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# PRECISION MASS DETERMINATION OF THE HIGGS BOSON AT PHOTON-PHOTON COLLIDERS**footnote *This work was supported in part by the U.S. Department of Energy under Contract No. DE-AC03-76SF00098.
## 1 Introduction
One of the most important tasks of the current and future collider experiments will be to detect and study Higgs boson(s). The accuracy of the measurement of the Higgs boson mass will impact precision tests of loop corrections, both in the standard model (SM) and in the extended models such as the minimal supersymmetric model (MSSM) .$`^{\text{“sevenrm}\text{?}\text{,}\text{?}\text{,}\text{?}}`$ Deviations of the total widths of the Higgs bosons from SM predictions can be directly compared to predictions of alternative models such as the MSSM, the non-minimal supersymmetric standard model, or the general two-Higgs-doublet model .$`^{\text{“sevenrm}\text{?}\text{,}\text{?}\text{,}\text{?}}`$ The total widths for the SM Higgs boson $`h_{SM}`$ and the three neutral Higgs bosons $`h^0,H^0,A^0`$ of the MSSM are shown in Fig. 1.
The interaction of high energy photons at a photon-photon collider $`^{\text{“sevenrm}\text{?}\text{,}\text{?}\text{,}\text{?}}`$ provides us with an unique opportunity to study Higgs boson, because the SM Higgs boson in $`s`$-channel resonance can be produced at photon-photon colliders .$`^{\text{“sevenrm}\text{?}\text{,}\text{?}\text{,}\text{?}\text{,}\text{?}\text{,}\text{?}}`$ In this paper we point out precision measurements of mass ($`m_h`$) and total width ($`\mathrm{\Gamma }_h`$) of the Higgs boson by the method of energy scanning, using the high energy edge of the photon spectrum.
The method of energy scanning at photon-photon colliders was first mentioned by V. Telnov .$`^{\text{“sevenrm}\text{?}}`$ The luminosity of the photon-photon collider has a very sharp edge at high energy, much narrower than the width of the luminosity peak. If the Higgs boson is a very narrow resonance, we will observe a rapid increase in the visible cross section of the Higgs production during energy scanning.
## 2 Luminosity of Photon-Photon Colliders
Figure 2 shows ten differential luminosities with the $`J_z=0`$ angular momentum state of initial photon collisions in a photon-photon collider for energy scanning at $`m_h=100`$ GeV. In this study, we have scanned the Higgs boson resonance from the left side to the right side in Fig. 2. The circles exhibit the luminosity points in contact with the Higgs boson and the rise of the luminosity at $`m_h=100`$ GeV is rapid at the threshold of energy scanning. Here we introduce the required parameters for the luminosity calculation. A laser photon of energy $`\omega _L`$ is scattered by an electron beam of energy $`E_e`$ in the conversion region of the photon-photon collider. The kinematics of Compton scattering is characterized by the dimensionless parameter $`^{\text{“sevenrm}\text{?}}`$
$`x{\displaystyle \frac{4E_e\omega _L}{m_e^2}}15.3\left[{\displaystyle \frac{E_e}{\mathrm{TeV}}}\right]\left[{\displaystyle \frac{\omega _L}{\mathrm{eV}}}\right],`$ (1)
where $`m_e`$ is electron mass. The maximum energy of the scattered photon $`\omega _{max}`$ is $`E_ex/(x+1)`$ given by $`x`$. The parameter $`x`$ is fixed to be 4.8, and we get $`\omega _{max}=100`$ GeV when $`E_e=121`$ GeV and $`\omega _L=2.6`$ eV. The combination of the polarizations of the electron $`P_e`$ and the laser $`P_L`$ should be $`P_LP_e=1`$ so that the generated photon spectrum peaks at its maximum energy.
The differential luminosity distribution depends on the variable $`\rho =b/(\gamma a)`$, where $`a`$ is the rms radius of the electron beam at the interaction point (IP), $`b`$ is the distance between the conversion point (CP) and the IP, and $`\gamma =E_e/m_e`$. The polarized luminosities with the $`J_z=0`$ and the $`J_z=\pm 2`$ in a photon-photon collider were used in Ref. ?. Here we assumed $`\rho =1`$ and the conversion coefficient $`k=0.6`$. It should be noted that the shape of the high energy edge and $`w_{max}`$ are influenced by nonlinear effects due to very strong focus of the laser field at the CP. Prior to the actual energy scan, we need to have a fairly good estimate for nonlinear effects including the polarization.
## 3 Higgs Boson and Backgrounds
Once the Higgs boson is observed at future $`e^+e^{}`$ colliders, we must determine its precise mass and width in order to reveal the fundamental properties of the Higgs boson. At a photon-photon collider, the feasibility of the measurement of the two-photon decay width of a Higgs boson has been studied in the mass range $`M_W<m_h<2M_W`$ .$`^{\text{“sevenrm}\text{?}\text{,}\text{?}\text{,}\text{?}\text{,}\text{?}\text{,}\text{?}}`$ For $`m_h<2M_W`$, the SM Higgs boson mainly decays into a $`b\overline{b}`$ pair and the daughter $`b`$-flavored hadrons will be easily identified due to their long lifetime; therefore, the $`b\overline{b}`$ events are the best signals. The cross section of the Higgs boson production can be described by the Breit-Wigner approximation:
$`\sigma _{\gamma \gamma hb\overline{b}}(\sqrt{s})=8\pi {\displaystyle \frac{\mathrm{\Gamma }(h\gamma \gamma )\mathrm{\Gamma }(hb\overline{b})}{(sm_h^2)^2+m_h^2\mathrm{\Gamma }_h^2}}(1+\lambda _1\lambda _2),`$ (2)
where $`\mathrm{\Gamma }(h\gamma \gamma )`$ and $`\mathrm{\Gamma }(hb\overline{b})`$ are the decay widths of the Higgs boson into two photons and a $`b\overline{b}`$ pair, $`\lambda _1`$ and $`\lambda _2`$ the initial photon helicities, respectively. The effective cross section of the signal events within $`m_h\delta <\sqrt{s}<m_h+\delta `$ is
$`\sigma _{\gamma \gamma hb\overline{b}}^{\mathrm{eff}}={\displaystyle _{m_h\delta }^{m_h+\delta }}16\pi {\displaystyle \frac{\mathrm{\Gamma }(h\gamma \gamma )\mathrm{\Gamma }(hb\overline{b})}{(\widehat{s}m_h^2)^2+m_h^2\mathrm{\Gamma }_h^2}}{\displaystyle \frac{1}{L_{\gamma \gamma }}}{\displaystyle \frac{dL_{\gamma \gamma }^{J_z=0}}{d\sqrt{\widehat{s}}}}𝑑\sqrt{\widehat{s}},`$ (3)
where $`\delta `$ expresses the effect of the detector resolution and we assumed $`\delta =5`$ GeV. Here we supposed that the total luminosity is $`L_{\gamma \gamma }=L_{\gamma \gamma }^{J_z=0}+L_{\gamma \gamma }^{J_z=\pm 2}`$.
The main background processes may be the continuum $`\gamma \gamma b\overline{b}`$, $`c\overline{c}`$ as well as the radiative processes $`\gamma \gamma b\overline{b}g`$, $`c\overline{c}g`$. The continuum backgrounds dominantly produced by initial photon collisions in the $`J_z=\pm 2`$ can be suppressed by controlling the polarization of the colliding photon beams. Several authors reported that the effect of QCD corrections to $`\gamma \gamma q\overline{q}`$ is large since the helicity suppression which affects the background $`q\overline{q}`$ events does not work due to a gluon emission .$`^{\text{“sevenrm}\text{?}\text{,}\text{?}}`$ Recently leading double-logarithmic QCD corrections for $`J_z=0`$ were resummed to all orders and the account of non-Sudakov form factor to higher orders makes the cross-section well defined and positive definite in all regions of the phase space .$`^{\text{“sevenrm}\text{?}}`$ In this study we take account of the one-loop QCD corrections of the soft gluon emission, hard gluon emission, and virtual correction, where higher order double logarithmic corrections are not taken into account .$`^{\text{“sevenrm}\text{?}}`$ The effective cross section of the background process $`\gamma \gamma b\overline{b}(g)`$ or $`c\overline{c}(g)`$ within $`m_h\delta <\sqrt{s}<m_h+\delta `$ is
$`\sigma _{\mathrm{bg}}^{\mathrm{eff}}={\displaystyle _{m_h\delta }^{m_h+\delta }}\sigma _{\mathrm{bg}}(\sqrt{\widehat{s}}){\displaystyle \frac{1}{L_{\gamma \gamma }}}{\displaystyle \frac{dL_{\gamma \gamma }}{d\sqrt{\widehat{s}}}}𝑑\sqrt{\widehat{s}},`$ (4)
where $`\sigma _{\mathrm{bg}}(\sqrt{s})`$ is the cross section of the background process.
Since the cross section of $`\gamma \gamma c\overline{c}`$ is larger than that of $`\gamma \gamma b\overline{b}`$ due to the large electric charge of the quark, we apply the $`b`$ tagging in order to eliminate the charm and the light quark backgrounds. By the topological vertexing method $`^{\text{“sevenrm}\text{?}}`$ and the LC Vertex Detector design ,$`^{\text{“sevenrm}\text{?}}`$ the efficiency and purity of $`b`$-quark jet identification are 70% and 99%, respectively. Therefore the tagging efficiencies of $`b\overline{b}(g)`$ and $`c\overline{c}(g)`$ events are assumed as 49% and 0.005% with double tagging, respectively.<sup>a</sup><sup>a</sup>aThe interaction region (IR) at photon-photon colliders is complicated, because there are the sweeping magnet for spent electrons and the optical mirror system for laser focusing around the vertex detector. We need to study the performance of the vertex detector at the IR. $`^{\text{“sevenrm}\text{?}}`$ We impose the following cuts to remove backgrounds: (1) the double $`b\overline{b}`$ tagging in the event; (2) $`|\mathrm{cos}\theta _{b,\overline{b}}|<0.95`$, where $`\theta _{b,\overline{b}}`$ is the scattering angle of the $`b(\overline{b})`$ quark; (3) $`|M_{b\overline{b}}m_h|<5`$ GeV.
## 4 Results and Discussion
Figure 4 shows an example of energy scan to determine $`m_h`$. Each energy point corresponds to 5 $`\mathrm{fb}^1`$ and the total luminosity of photon-photon collisions is 50 $`\mathrm{fb}^1`$ in the same distributions as with Fig. 2. The total width of the SM Higgs boson $`\mathrm{\Gamma }_{h\mathrm{SM}}`$ for $`m_h`$=100 GeV is 2.16 MeV, which is computed by the HDECAY program .$`^{\text{“sevenrm}\text{?}}`$ The partial widths $`\mathrm{\Gamma }(h\gamma \gamma )`$ and $`\mathrm{\Gamma }(hb\overline{b})`$ at the predicted SM value with $`m_h=`$ 100 GeV are fixed for energy scanning at $`m_h=`$ 99.8, 100, 100.2 GeV. The statistical errors in Fig. 4 indicate $`\sqrt{S+B}`$, where $`S`$ and $`B`$ are the numbers for signal and background events. From Fig. 4, we can understand that the method of energy scanning for $`m_h`$ is more effective than that of the measurement of a single point at the luminosity peak using the same total luminosity, because the statistical errors at the threshold of energy scanning are smaller than that at the luminosity peak and we can distinguish the mass difference of 200 MeV. With the energy scanning of 10 points, the attainable error in $`m_h`$ is about 110 MeV at the $`1\sigma `$ level.
The measurement for the determination of $`\mathrm{\Gamma }_h`$ by the method of energy scan is shown in Fig. 4. Each energy point corresponds to 5 $`\mathrm{fb}^1`$ and the total luminosity is 50 $`\mathrm{fb}^1`$. The partial widths $`\mathrm{\Gamma }(h\gamma \gamma )`$ and $`\mathrm{\Gamma }(hb\overline{b})`$ at the predicted SM value with $`m_h=100`$ GeV are fixed for energy scanning $`\mathrm{\Gamma }_h/\mathrm{\Gamma }_{h\mathrm{SM}}=`$ 0.9, 1.0, 1.1. The large difference between the total widths at the luminosity peak can be seen easily in Fig. 4. The statistical error in $`\mathrm{\Gamma }_h`$ is about 6% at the $`1\sigma `$ level. If there are additional invisible decay modes of Higgs boson, only the total decay width increases keeping the partial widths of two photons and a $`b\overline{b}`$ pair unchanged. In this study we find $`\mathrm{\Gamma }_h/\mathrm{\Gamma }_{h_{\mathrm{SM}}}>1`$. Of course, this deviation from the SM should have also been observed in the parent $`e^+e^{}`$ collider. However, this will be independent observation in gamma-gamma energy scan, which confirms the $`e^+e^{}`$ result.
Here we consider two cases for the photon-photon collider. First, we choose $`x=4.8`$ while tuning the energies of the laser photon and the electron beam while tuning the scan. Second, we fix the laser energy and only the energy of the electron beam is tuned during this scan. The two cases are called the tunable and fixed cases, respectively.
Table 4 lists the statistical errors of the SM Higgs boson mass at the $`1\sigma `$ level, using an integrated luminosity of 50/10 $`\mathrm{fb}^1`$. In this table, the mass errors of the tunable case are almost smaller than those of the fixed case. Since the background processes $`\gamma \gamma q\overline{q}(g)`$ are increasing at the lower Higgs mass and the branching ratio $`B(hb\overline{b})`$ is decreasing at the higher Higgs mass, the errors of the Higgs boson mass near 100 GeV are the smallest. The statistical errors $`\sqrt{S+B}/S`$ of the total width $`\mathrm{\Gamma }_h/\mathrm{\Gamma }_{h\mathrm{SM}}`$ of the SM Higgs boson with a 50 $`\mathrm{fb}^1`$ luminosity are listed in Table 4. The statistical errors of the total width for intermediate-mass Higgs bosons are almost within 8% in Table 4. Comparatively the results with the tagging efficiencies 70% and 3.5% of $`b\overline{b}(g)`$ and $`c\overline{c}(g)`$ events are listed in Tables 4 and 4.
At the future colliders, the expected precision for the mass of the Higgs boson with $`m_{h_{\mathrm{SM}}}`$=100 GeV is listed in Table 4. The NLC threshold result is at $`\sqrt{s}=m_Z+m_{h_{\mathrm{SM}}}+0.5`$ GeV including the initial state radiation and the beam energy spread .$`^{\text{“sevenrm}\text{?}}`$ The LHC error is for ATLAS+CMS including the statistical and systematic errors .$`^{\text{“sevenrm}\text{?}}`$ The error at the muon collider is devoted to the scan with beam energy resolution of $`0.01\%`$ .$`^{\text{“sevenrm}\text{?}}`$ From the table, the accuracy of the Higgs boson mass at the muon collider is the highest, however the systematic error at the muon collider is neglected assuming accurate beam energy determination. The accuracy at the photon-photon collider is 1.5 times lower than that at the NLC threshold case. Therefore we can perform the complementary measurement of Higgs boson mass at photon-photon colliders.
As for other origins of the errors, we need to know the systematic uncertainties on the luminosity distribution. The possibilities of the luminosity measurements at photon-photon colliders have been studied using the process $`\gamma \gamma l^+l^{}`$ or $`\gamma \gamma W^+W^{}`$ .$`^{\text{“sevenrm}\text{?}\text{,}\text{?}}`$ For energy scanning the measurement of the luminosity distribution at the high energy-edge is crucial and we need to study it further.
## 5 Summary
In this paper, we have shown that it is possible to determine the Higgs boson mass to a high precision by the method of energy scanning at photon-photon colliders, using the high energy edge of the photon spectrum.
Acknowledgments
I express sincere thanks to T. Takahashi, T. Tauchi, V. Telnov, I. Watanabe, M. Xie and K. Yokoya for useful discussions.
References
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# 1 Phase 𝛿₁-dependence of the oscillation probabilities 𝑃_{𝛼𝛽} and their differences Δ𝑃_{𝛼𝛽} for the long-baseline experiment. 𝑃_{𝛼𝛽}≡𝑃(𝜈_𝛼→𝜈_𝛽) and Δ𝑃_{𝛼𝛽}≡𝑃(𝜈_𝛼→𝜈_𝛽)-𝑃((𝜈_𝛼)̄→(𝜈_𝛽)̄). 𝛿₁ is in degree.
CP violation effect in long-baseline neutrino oscillation
in the four-neutrino model
Toshihiko Hattori$`,^{a),}`$<sup>1</sup><sup>1</sup>1e-mail: hattori@ias.tokushima-u.ac.jp Tsutom Hasuike$`,^{b),}`$<sup>2</sup><sup>2</sup>2e-mail: hasuike@anan-nct.ac.jp and Seiichi Wakaizumi<sup>c),</sup><sup>3</sup><sup>3</sup>3e-mail: wakaizum@medsci.tokushima-u.ac.jp
<sup>a)</sup>Institute of Theoretical Physics, University of Tokushima, Tokushima 770-8502, Japan
<sup>b)</sup>Department of Physics, Anan College of Technology, Anan 774-0017, Japan
<sup>c)</sup>School of Medical Sciences, University of Tokushima, Tokushima 770-8509, Japan
Abstract
We investigate CP-violation effect in the long-baseline neutrino oscillation in the four-neutrino model with mass scheme of the two nearly degenerate pairs separated with the order of 1 eV, by using the data from the solar neutrino deficit, the atmospheric neutrino anomaly and the LSND experiments along with the other accelerator and reactor experiments. By use of the most general parametrization of the mixing matrix with six angles and six phases, we show that the genuine CP-violation effect could attain as large as 0.3 for $`\mathrm{\Delta }P(\nu _\mu \nu _\tau )P(\nu _\mu \nu _\tau )P(\overline{\nu _\mu }\overline{\nu _\tau })`$ and that the matter effect is negligibly small such as at most 0.01 for $`\mathrm{\Delta }P(\nu _\mu \nu _\tau )`$ for $`\mathrm{\Delta }m^2=(15)\times 10^3`$ $`\mathrm{eV}^2`$, which is the mass-squared difference relevant to the long-baseline oscillation.
I Introduction
It has long been assumed that neutrinos are massless. However, since the atmospheric neutrino anomaly was discovered by several experimental Collaborations and was affirmatively confirmed by the Super-Kamiokande, people have come to think through the neutrino oscillation interpretation for the anomaly that neutrinos seem to have a certain amount of mass. Together with the solar neutrino deficit, the anomaly has been analyzed in the three-neutrino model and two typical mass scales have been derived for the neutrino mass-squared difference; $`\mathrm{\Delta }m_{\mathrm{atm}}^2=(0.56)\times 10^3`$ $`\mathrm{eV}^2`$ with a large mixing angle of $`\mathrm{sin}^22\theta _{\mathrm{atm}}>0.82`$ as the $`\nu _\mu \nu _\tau `$ oscillation from the atmospheric neutrino anomaly and $`\mathrm{\Delta }m_{\mathrm{solar}}^2=(10^{11}10^5)`$ $`\mathrm{eV}^2`$, a large range depending on the three solutions of the vacuum oscillation and the MSW solutions in the matter with small- and large-angle mixings from the solar neutrino deficit.
As in the quark sector, CP violation would be a characteristic feature in the three-neutrino model. It has been shown by using the constraints on the mixing matrix elements obtained from the analyses of these anomalies along with the results from the other accelerator and reactor experiments that the CP violation effect, defined as a difference of the oscillation probabilities between the neutrino and the antineutrino, is typically $`13\%`$ even in the long-baseline neutrino oscillations, depending on the assumed mass hierarchies.
On the other hand, sterile neutrinos were considered in the context of neutrino oscillations. After that, a four-neutrino model of the ordinary three active neutrinos and one sterile neutrino was introduced with a mass pattern of two nearly degenerate pairs separated with a mass gap of the order of 1eV motivated from the hot dark matter and then, by using the only one possible positive evidence from the terrestrial LSND experiments on the oscillations $`\nu _\mu \nu _e`$ and $`\overline{\nu _\mu }\overline{\nu _e}`$ , which suggest the mass scale of $`\mathrm{\Delta }m_{\mathrm{LSND}}^2=(0.32.2)`$ $`\mathrm{eV}^2`$, the four-neutrino model with the same mass pattern as the above is studied \[14-20\]. In this model, a sizable CP violation effect is shown to be possible in the long-baseline experiments, and different magnitudes of the probability difference between the CP-conjugate channels are expected in between the three-neutrino model and the four-neutrino model by using the most general parametrization of the mixing matrix. And, some features of CP asymmetry defined as the normalized probability difference are discussed in the long-baseline experiments at a neutrino factory.
We will investigate the CP violation effect in the long-baseline neutrino oscillations numerically in more detail in the four-neutrino model with mass scheme of the two nearly degenerate pairs separated with the order of 1 eV by using the most general parametrization of the mixing matrix, and in addition we will study the matter effect in the four-neutrino model.
The paper is organized as follows. In Sec. II the four-neutrino model we use here is presented and the expressions of the difference of oscillation probabilities between the CP-conjugate channels are given both in the exact form and in the approximate forms relevant to the short-baseline and the long-baseline neutrino oscillations for the neutrino mass scheme mentioned above. In Sect. III constraints on the neutrino mixing matrix are derived by using the solar neutrino deficit, atmospheric neutrino anomaly, Bugey reactor experiment, CHOOZ experiment, LSND experiments, CHORUS and NOMAD experiments and the other accelerator and reactor experiments. In Sect. IV the most general parametrization of the mixing matrix is adopted to obtain the constraints on the mixing angles and phases from the ones on the mixing matrix derived in Sect. III. And then, CP-violation in the long-baseline neutrino oscillation is investigated on the basis of these constraints. The behavior of the oscillation probability differences is analyzed in detail with respect to the two relevant phases of the mixing matrix and $`\mathrm{\Delta }m^2/E`$. The matter effect is shown to be negligibly small in the four-neutrino model with the mass scheme adopted here. Finally, Sect. V is devoted to the conclusion.
II The four-neutrino model
In order to consider the solar neutrino deficit, the atmospheric neutrino anomaly and the LSND experiment, we will take the four-neutrino model with the three ordinary active neutrinos and one sterile neutrino with three different scales of the neutrino mass-squared difference, $`\mathrm{\Delta }m_{\mathrm{solar}}^2=(10^{11}10^5)`$ $`\mathrm{eV}^2,\mathrm{\Delta }m_{\mathrm{atm}}^2=(10^310^2)`$ $`\mathrm{eV}^2`$ and $`\mathrm{\Delta }m_{\mathrm{LSND}}^2=(0.310)`$ $`\mathrm{eV}^2`$.
Under the neutrino oscillation hypothesis, the flavor eigenstates of neutrinos are the mixtures of mass eigenstates with masses $`m_i(i=1,2,3,4)`$ as follows,
$$\nu _\alpha =\underset{i=1}{\overset{4}{}}U_{\alpha i}\nu _i,\alpha =e,\mu ,\tau ,s$$
(1)
where $`\nu _e,\nu _\mu `$ and $`\nu _\tau `$ are the ordinary neutrinos and $`\nu _s`$ is the sterile neutrino, and $`U`$ is the unitary mixing matrix. The neutrino oscillation probability of $`\nu _\alpha \nu _\beta `$ in vacuum is given in the usual manner in the four-neutrino model by
$$P(\nu _\alpha \nu _\beta )=\delta _{\alpha \beta }4\underset{k>j}{}\mathrm{Re}(U_{\alpha k}^{}U_{\alpha j}U_{\beta j}^{}U_{\beta k})\mathrm{sin}^2\mathrm{\Delta }_{kj}+2\underset{k>j}{}\mathrm{Im}(U_{\alpha k}^{}U_{\alpha j}U_{\beta j}^{}U_{\beta k})\mathrm{sin}2\mathrm{\Delta }_{kj},$$
(2)
where $`\mathrm{\Delta }_{kj}\mathrm{\Delta }m_{kj}^2L/(4E)`$, $`L`$ being the distance from the neutrino source and $`E`$ the energy of neutrino. The oscillation probability for the antineutrinos is given by the exchange of $`UU^{}`$ in Eq.(2). And, the probability difference between CP-conjugate channels given by
$`\mathrm{\Delta }P_{\alpha \beta }`$ $``$ $`P(\nu _\alpha \nu _\beta )P(\overline{\nu _\alpha }\overline{\nu _\beta })`$ (3)
$`=`$ $`4{\displaystyle \underset{k>j}{}}\mathrm{Im}(U_{\alpha k}^{}U_{\alpha j}U_{\beta j}^{}U_{\beta k})\mathrm{sin}2\mathrm{\Delta }_{kj}`$
is a direct measure of the genuine CP-violation effect in the neutrino oscillation in vacuum.
The four neutrino masses should be devided into two pairs of close masses separated by a gap of about 1eV in order to accomodate with the solar and atmospheric neutrino deficits and the LSND experiments along with the other results from the accelerator and reactor experiments on the neutrino oscillation. There are the following two schemes for that mass pattern; (i) $`\mathrm{\Delta }m_{\mathrm{solar}}^2\mathrm{\Delta }m_{21}^2\mathrm{\Delta }m_{\mathrm{atm}}^2\mathrm{\Delta }m_{43}^2\mathrm{\Delta }m_{\mathrm{LSND}}^2\mathrm{\Delta }m_{32}^2`$, and (ii) $`\mathrm{\Delta }m_{\mathrm{solar}}^2\mathrm{\Delta }m_{43}^2\mathrm{\Delta }m_{\mathrm{atm}}^2\mathrm{\Delta }m_{21}^2\mathrm{\Delta }m_{\mathrm{LSND}}^2\mathrm{\Delta }m_{32}^2`$, where $`\mathrm{\Delta }m_{kj}^2m_k^2m_j^2`$. We will adopt the first scheme in the following analyses, and the second scheme can be attained only through the exchange of indices $`(1,2)(3,4)`$ in the following various expressions such as the oscillation probabilities. In the first scheme, the measure of CP violation in the neutrino oscillation in vacuum is given for the short-baseline experiment ( $`L/E1`$ \[km/GeV\] ) as follows,
$`\mathrm{\Delta }P_{\alpha \beta }`$ $``$ $`4[\mathrm{Im}(U_{\alpha 3}^{}U_{\alpha 2}U_{\beta 2}^{}U_{\beta 3})+\mathrm{Im}(U_{\alpha 3}^{}U_{\alpha 1}U_{\beta 1}^{}U_{\beta 3})`$ (4)
$`+\mathrm{Im}(U_{\alpha 4}^{}U_{\alpha 2}U_{\beta 2}^{}U_{\beta 4})+\mathrm{Im}(U_{\alpha 4}^{}U_{\alpha 1}U_{\beta 1}^{}U_{\beta 4})]\mathrm{sin}2\mathrm{\Delta }_{32},`$
since $`\mathrm{\Delta }_{21}`$ and $`\mathrm{\Delta }_{43}1`$, and $`\mathrm{\Delta }_{41},\mathrm{\Delta }_{42},\mathrm{\Delta }_{31},\mathrm{\Delta }_{32}1`$. $`\mathrm{\Delta }P_{\alpha \beta }`$ in Eq.(4) is zero due to the unitarity of the mixing matrix $`U`$. So, CP violation is negligibly small in the short-baseline oscillation experiments in the four-neutrino model.
On the other hand, for the long-baseline experiment ( $`L/E=1001000`$ \[km/GeV\]) the probability difference in vacuum is given as follows,
$$\mathrm{\Delta }P_{\alpha \beta }4\mathrm{I}\mathrm{m}(U_{\alpha 4}^{}U_{\alpha 3}U_{\beta 3}^{}U_{\beta 4})\mathrm{sin}2\mathrm{\Delta }_{43},$$
(5)
since $`\mathrm{\Delta }_{21}1`$, $`\mathrm{\Delta }_{41},\mathrm{\Delta }_{42},\mathrm{\Delta }_{31},\mathrm{\Delta }_{32}1`$, and $`\mathrm{\Delta }_{43}1`$. There are six $`\mathrm{\Delta }P_{\alpha \beta }`$’s, that is, $`\mathrm{\Delta }P_{\mu e},\mathrm{\Delta }P_{e\tau },\mathrm{\Delta }P_{\mu \tau },\mathrm{\Delta }P_{es},\mathrm{\Delta }P_{\mu s}`$, and $`\mathrm{\Delta }P_{\tau s}`$. Three of these six $`\mathrm{\Delta }P_{\alpha \beta }`$’s are independent due to the unitarity of $`U`$ for the approximate expression of $`\mathrm{\Delta }P_{\alpha \beta }`$ in Eq.(5) as well as for the exact expression in Eq.(3).
III Constraints on the mixing matrix $`U`$
In order to numerically calculate the oscillation probability differences $`\mathrm{\Delta }P_{\alpha \beta }`$, we will derive the constraints on the mixing matrix $`U`$ from the solar neutrino deficit, atmospheric neutrino anomaly, LSND experiments and the other terrestrial oscillation experiments using the accelerators and reactors.
(i) Solar neutrino deficit
Since $`\mathrm{\Delta }_{21}1`$ and all the other five $`\mathrm{\Delta }_{kj}`$’s are enormously larger than 1, the survival probability of $`\nu _e`$ is given from Eq.(2) by
$`P_{\mathrm{solar}}(\nu _e\nu _e)`$ $``$ $`14|U_{e1}|^2|U_{e2}|^2\mathrm{sin}^2\mathrm{\Delta }_{21}2|U_{e3}|^2(1|U_{e3}|^2|U_{e4}|^2)`$ (6)
$`2|U_{e4}|^2(1|U_{e4}|^2),`$
where the unitarity of $`U`$ is used. For the solar neutrino deficit, there are three different kinds of solutions, that is, the vacuum solution and the MSW solutions with small and large angle mixings, and a unique solution is not yet found, so that we will not use this deficit in order to obtain the constraints.
(ii) Atmospheric neutrino anomaly
Since $`\mathrm{\Delta }_{21}1,\mathrm{\Delta }_{43}1`$ and $`\mathrm{\Delta }_{41},\mathrm{\Delta }_{42},\mathrm{\Delta }_{31},\mathrm{\Delta }_{32}1`$, the survival probability of $`\nu _\mu `$ is given by
$$P_{\mathrm{atm}}(\nu _\mu \nu _\mu )14|U_{\mu 3}|^2|U_{\mu 4}|^2\mathrm{sin}^2\mathrm{\Delta }_{43}2(|U_{\mu 1}|^2+|U_{\mu 2}|^2)(1|U_{\mu 1}|^2|U_{\mu 2}|^2).$$
(7)
By using the data from the Super-Kamiokande experiments, that is, $`\mathrm{sin}^22\theta _{\mathrm{atm}}>0.82`$ for $`5\times 10^4<\mathrm{\Delta }m_{\mathrm{atm}}^2<6\times 10^3`$ $`\mathrm{eV}^2`$, and expecting from this data that $`|U_{\mu 1}|^2+|U_{\mu 2}|^21`$, the following constraint is obtained,
$$|U_{\mu 3}|^2|U_{\mu 4}|^2>0.205.$$
(8)
(iii) The Bugey experiment (including Krasnoyarsk, CDHS and CCFR experiments)
By being typically represented by the Bugey reactor experiment with $`L/E=320`$ \[m/MeV or km/GeV\], since $`\mathrm{\Delta }_{21}1,\mathrm{\Delta }_{43}1`$ and $`\mathrm{\Delta }_{41},\mathrm{\Delta }_{42},\mathrm{\Delta }_{31},\mathrm{\Delta }_{32}1`$, the survival probability of $`\overline{\nu _e}`$ is given by
$$P_{\mathrm{Bugey}}(\overline{\nu _e}\overline{\nu _e})14(|U_{e3}|^2+|U_{e4}|^2)(1|U_{e3}|^2|U_{e4}|^2)\mathrm{sin}^2\mathrm{\Delta }_{32}.$$
(9)
If we use the data from the Bugey experiment conservatively, that is, $`\mathrm{sin}^22\theta _{\mathrm{Bugey}}<0.1`$ for $`0.1<\mathrm{\Delta }m^2<1`$ $`\mathrm{eV}^2`$, the following constraint is obtained,
$$|U_{e3}|^2+|U_{e4}|^2<0.025.$$
(10)
(iv) The CHOOZ experiment
This experiment is the first long-baseline reactor experiment, since $`L1`$ km and $`E3`$ MeV so that $`L/E300`$ \[km/GeV\]. Therefore, $`\mathrm{\Delta }_{21}1,\mathrm{\Delta }_{43}1`$ and $`\mathrm{\Delta }_{41},\mathrm{\Delta }_{42},\mathrm{\Delta }_{31},\mathrm{\Delta }_{32}1`$, and the survival probability of $`\overline{\nu _e}`$ is given by
$$P_{\mathrm{CHOOZ}}(\overline{\nu _e}\overline{\nu _e})14|U_{e3}|^2|U_{e4}|^2\mathrm{sin}^2\mathrm{\Delta }_{43}2(|U_{e3}|^2+|U_{e4}|^2)(1|U_{e3}|^2|U_{e4}|^2).$$
(11)
By using the data from the CHOOZ experiment, that is, $`\mathrm{sin}^22\theta _{\mathrm{CHOOZ}}<0.12`$ for $`3\times 10^3<\mathrm{\Delta }m^2<1.0\times 10^2`$ $`\mathrm{eV}^2`$ and adopting Eq.(10), the following constraint is obtained,
$$4|U_{e3}|^2|U_{e4}|^2<0.12.$$
(12)
If we use, however, the constraint of Eq.(10) and the unequality of $`2|U_{e3}||U_{e4}||U_{e3}|^2+|U_{e4}|^2`$, a constraint $`4|U_{e3}|^2|U_{e4}|^2<6.3\times 10^4`$ is obtained so that Eq.(12) is included in the constraint from the Bugey experiment.
(v) The LSND experiments
This experiment is of the short baseline, $`L/E=0.51`$ \[m/MeV\]. Since $`\mathrm{\Delta }_{21}1,\mathrm{\Delta }_{43}1`$ and $`\mathrm{\Delta }_{41},\mathrm{\Delta }_{42},\mathrm{\Delta }_{31},\mathrm{\Delta }_{32}1`$, the oscillation probability of $`\nu _\mu \nu _e`$ is expressed as follows,
$`P_{\mathrm{LSND}}(\nu _\mu \nu _e)`$ $``$ $`4\mathrm{R}\mathrm{e}\left[(U_{\mu 3}^{}U_{e3}+U_{\mu 4}^{}U_{e4})(U_{\mu 1}U_{e1}^{}+U_{\mu 2}U_{e2}^{})\right]\mathrm{sin}^2\mathrm{\Delta }_{32}`$ (13)
$`=`$ $`4|U_{\mu 3}^{}U_{e3}+U_{\mu 4}^{}U_{e4}|^2\mathrm{sin}^2\mathrm{\Delta }_{32},`$
where the unitarity of $`U`$ is used. By using the data from the LSND experiments, that is, $`\mathrm{sin}^22\theta _{\mathrm{LSND}}=1.5\times 10^31.0\times 10^1`$ for $`0.3<\mathrm{\Delta }m_{\mathrm{LSND}}^2<2.2`$ $`\mathrm{eV}^2`$, the following constraint is obtained,
$$|U_{\mu 3}^{}U_{e3}+U_{\mu 4}^{}U_{e4}|=0.020.16.$$
(14)
(vi) The CHORUS and NOMAD experiments
These experiments are also the short baseline ones searching for the $`\nu _\mu \nu _\tau `$ oscillation, $`L/E=0.020.03`$ \[km/GeV\]. Since $`\mathrm{\Delta }_{21}1,\mathrm{\Delta }_{43}1`$ and $`\mathrm{\Delta }_{41},\mathrm{\Delta }_{42},\mathrm{\Delta }_{31},\mathrm{\Delta }_{32}10^210^1`$, the oscillation probability is given by
$$P_{\mathrm{CHORUS}/\mathrm{NOMAD}}(\nu _\mu \nu _\tau )4|U_{\mu 3}^{}U_{\tau 3}+U_{\mu 4}^{}U_{\tau 4}|^2\mathrm{sin}^2\mathrm{\Delta }_{32}.$$
(15)
By using the data from the latest NOMAD experiment, $`\mathrm{sin}^22\theta _{\mathrm{NOMAD}}<0.3`$ for $`\mathrm{\Delta }m^2<2.2`$ $`\mathrm{eV}^2`$, the following constraint is obtained,
$$|U_{\mu 3}^{}U_{\tau 3}+U_{\mu 4}^{}U_{\tau 4}|<0.28.$$
(16)
Among the above-mentioned six typical phenomena and experiments, the useful constraints are of Eqs. (8), (10), (14) and (16).
IV. CP violation in the neutrino oscillations
In this section, by using the constraints obtained in the previous section, we will numerically investigate the CP violation effects in the long-baseline neutrino oscillation experiments in the four-neutrino model described in Sect. II.
We adopt the most general parametrization of the mixing matrix $`U`$ for Majorana neutrinos, which includes six mixing angles and six phases. The expression of the matrix is too complicated to write it down here. So, we cite only the matrix elements which are useful for the following numerical analyses; $`U_{e1}=c_{01}c_{02}c_{03},U_{e2}=c_{02}c_{03}s_{d01}^{},U_{e3}=c_{03}s_{d02}^{},U_{e4}=s_{d03}^{},U_{\mu 3}=s_{d02}^{}s_{d03}s_{d13}^{}+c_{02}c_{13}s_{d12}^{},U_{\mu 4}=c_{03}s_{d13}^{},U_{\tau 3}=c_{13}s_{d02}^{}s_{d03}s_{d23}^{}c_{02}s_{d12}^{}s_{d13}s_{d23}^{}+c_{02}c_{12}c_{23}`$, and $`U_{\tau 4}=c_{03}c_{13}s_{d23}^{}`$, where $`c_{ij}\mathrm{cos}\theta _{ij}`$ and $`s_{dij}s_{ij}\mathrm{e}^{\mathrm{i}\delta _{ij}}\mathrm{sin}\theta _{ij}\mathrm{e}^{\mathrm{i}\delta _{ij}}`$ , and $`\theta _{01},\theta _{02},\theta _{03},\theta _{12},\theta _{13},\theta _{23}`$ are the six angles and $`\delta _{01},\delta _{02},\delta _{03},\delta _{12},\delta _{13},\delta _{23}`$ are the six phases. As stated in Sect. II, three of the six oscillation probability differences are independent so that only three of the six phases are determined by the measurements of the CP violation effect in the neutrino oscillations. In this sense, our analyses apply both to the Dirac and Majorana neutrinos.
On the basis of this parametrization, we obtain the constraints on the mixing angles and phases by using the constraints on the mixing matrix elements derived in the previous section. First, the constraint of Eq.(10) leads to
$$c_{03}^2s_{02}^2+s_{03}^2<0.025.$$
(17)
This unequality means at least $`s_{02}^2,s_{03}^2<0.025`$. The next constraint of Eq.(8) of
$$|s_{02}s_{03}s_{13}\mathrm{e}^{\mathrm{i}(\delta _{02}\delta _{03}+\delta _{13})}+c_{02}c_{13}s_{12}\mathrm{e}^{\mathrm{i}\delta _{12}}|^2c_{03}^2s_{13}^2>0.205$$
(18)
leads to
$$s_{12}^2c_{13}^2s_{13}^2>0.205$$
(19)
due to the smallness of $`s_{02}`$ and $`s_{03}`$. The third constraint of Eq.(14) gives the following expression,
$$|c_{02}s_{02}c_{03}s_{12}c_{13}+c_{02}^2c_{03}s_{03}s_{13}\mathrm{e}^{\mathrm{i}\delta _1}|=0.020.16,$$
(20)
where $`\delta _1\delta _{02}\delta _{03}\delta _{12}+\delta _{13}`$. This constraint proves not to bring any constraint on the phase $`\delta _1`$, if we use Eqs.(17) and (19). The fourth constraint of Eq.(16) is expressed as
$`|`$ $`c_{02}^2c_{12}s_{12}c_{13}c_{23}c_{02}s_{02}s_{03}s_{12}c_{13}^2s_{23}\mathrm{e}^{\mathrm{i}(\delta _1+\delta _2)}c_{02}s_{02}s_{03}c_{12}s_{13}c_{23}\mathrm{e}^{\mathrm{i}\delta _1}`$ (21)
$`+c_{13}s_{13}s_{23}(c_{03}^2c_{02}^2s_{12}^2+s_{02}^2s_{03}^2)\mathrm{e}^{\mathrm{i}\delta _2}|<0.28,`$
where $`\delta _2\delta _{12}\delta _{13}+\delta _{23}`$. By using Eqs. (17) and (19), no constraint on $`\delta _1`$, and the fact of the large angle mixing in $`\nu _\mu \nu _\tau `$ oscillation for the atmospheric neutrino anomaly which leads to the nearly maximal mixing in the angle $`\theta _{23}`$, the constraint of Eq. (21) gives no constraint to the phase $`\delta _2`$.
So, in summary, we derive the two constraints of Eqs. (17) and (19) on the mixing angles and no constraints on the two phases of $`\delta _1`$ and $`\delta _2`$.
Using these two constraints on the mixing angles, we will calculate the differences of the oscillation probabilities between the CP-conjugate channels for the long-baseline neutrino oscillations. As stated before, only three of the six probability differences among the four neutrinos are independent so that three of the six phases are relevant here. However, only two phases dominantly affect the differences as is shown by the leading terms relevant to the long-baseline oscillation, which are given in the following,
$`\mathrm{\Delta }P_{\mu e}`$ $``$ $`4c_{03}^2c_{02}s_{02}s_{03}s_{12}c_{13}s_{13}\mathrm{sin}\delta _1\mathrm{sin}2\mathrm{\Delta }_{43},`$
$`\mathrm{\Delta }P_{e\tau }`$ $``$ $`4c_{03}^2c_{02}s_{02}s_{03}c_{13}s_{23}\left[c_{12}c_{23}\mathrm{sin}(\delta _1+\delta _2)+s_{12}s_{13}s_{23}\mathrm{sin}\delta _1\right]\mathrm{sin}2\mathrm{\Delta }_{43},`$
$`\mathrm{\Delta }P_{\mu \tau }`$ $``$ $`4c_{03}^2c_{02}c_{13}s_{13}s_{23}[c_{02}c_{12}s_{12}c_{13}c_{23}\mathrm{sin}\delta _2+s_{02}s_{03}s_{12}s_{23}\mathrm{sin}\delta _1`$ (22)
$`s_{02}s_{03}c_{12}s_{13}c_{23}\mathrm{sin}(\delta _1+\delta _2)]\mathrm{sin}2\mathrm{\Delta }_{43},`$
$`\mathrm{\Delta }P_{es}`$ $``$ $`4c_{03}^2c_{02}s_{02}s_{03}c_{13}c_{23}\left[c_{12}s_{23}\mathrm{sin}(\delta _1+\delta _2)+s_{12}s_{13}c_{23}\mathrm{sin}\delta _1\right]\mathrm{sin}2\mathrm{\Delta }_{43},`$
$`\mathrm{\Delta }P_{\mu s}`$ $``$ $`4c_{03}^2c_{02}c_{13}s_{13}c_{23}[c_{02}c_{12}s_{12}c_{13}s_{23}\mathrm{sin}\delta _2s_{02}s_{03}s_{12}c_{23}\mathrm{sin}\delta _1`$
$`s_{02}s_{03}c_{12}s_{13}s_{23}\mathrm{sin}(\delta _1+\delta _2)]\mathrm{sin}2\mathrm{\Delta }_{43},`$
$`\mathrm{\Delta }P_{\tau s}`$ $``$ $`4c_{03}^2c_{02}c_{12}c_{13}^2c_{23}s_{23}\left[c_{02}s_{12}s_{13}\mathrm{sin}\delta _2+s_{02}s_{03}c_{13}\mathrm{sin}(\delta _1+\delta _2)\right]\mathrm{sin}2\mathrm{\Delta }_{43},`$
where $`\delta _1`$ and $`\delta _2`$ are the linear combinations of $`\delta _{ij}`$’s as stated before. We take the range of phases as $`0\delta _{ij}<2\pi `$ and the range of mixing angles as $`0\theta _{ij}\pi `$ so that $`s_{ij}`$’s can be taken only positive and $`c_{ij}`$’s can be taken both positive and negative. Since Eq. (17) means that the angles $`\theta _{02}`$ and $`\theta _{03}`$ are very small and Eq. (19) leads to $`s_{12}^2\mathrm{sin}^22\theta _{13}>0.82`$ which means that $`s_{12}`$ is in the range of $`0.9s_{12}1.0`$ and the angle $`\theta _{13}`$ is around $`\pi /4`$, $`\mathrm{\Delta }P_{\mu e}`$ and $`\mathrm{\Delta }P_{e\tau }`$ are expected from Eq.(22) to be very small and $`\mathrm{\Delta }P_{\mu \tau }`$ is to be able to take a sizable magnitude. In the following, we calculate the oscillation probabilities $`P(\nu _\alpha \nu _\beta )`$ and their differences $`\mathrm{\Delta }P_{\alpha \beta }`$ by using the rigorous expressions of Eqs. (2) and (3).
The probabilities of $`P(\nu _\mu \nu _e)`$ and $`P(\overline{\nu _\mu }\overline{\nu _e})`$ as functions of the phase $`\delta _1`$ with $`\delta _2=\pi /2`$ fixed are shown in Fig.1 and those of $`P(\nu _e\nu _\tau )`$ and $`P(\overline{\nu _e}\overline{\nu _\tau })`$ as functions of the phase $`\delta _2`$ with $`\delta _1=\pi /2`$ fixed are shown in Fig.2 for the values of the parameter set of angles and phases; $`s_{02}=s_{03}=0.11(c_{02}=c_{03}=0.994),s_{12}=0.91(c_{12}=0.415),s_{13}=0.67(c_{13}=0.742),s_{01}=s_{23}=1/\sqrt{2}(c_{01}=c_{23}=1/\sqrt{2})`$ and $`\delta _{01}=\delta _{02}=\delta _{03}=\delta _{12}=0`$, which are chosen so as to give the probability differences as large as possible within the parameter ranges allowed by the constaints of Eqs. (17) and (19). The magnitude of these probabilities is at most 0.04 as shown in Figs.1 and 2. Therefore, the probability differences $`\mathrm{\Delta }P(\nu _\mu \nu _e)`$ and $`\mathrm{\Delta }P(\nu _e\nu _\tau )`$ are at most $`\pm 0.02`$ as shown in Fig.3 for the same parameter values. On the other hand, $`P(\nu _\mu \nu _\tau )`$ and $`P(\overline{\nu _\mu }\overline{\nu _\tau })`$ can rise to as large as $`0.400.45`$ as shown in Fig.4 and $`\mathrm{\Delta }P(\nu _\mu \nu _\tau )`$ can attain as large as $`\pm 0.28`$ as shown in Fig.5 for the same parameter values. These facts agree with the above-mentioned expectations. The angle $`\theta _{23}`$-dependence of $`\mathrm{\Delta }P(\nu _\mu \nu _\tau )`$ and $`\mathrm{\Delta }P(\nu _e\nu _\tau )`$ is shown in Fig.6, where the phases $`\delta _1`$ and $`\delta _2`$ are taken as $`\pi /2`$ and the values of the other angles and phases are the same as the above. We display the phase $`\delta _1`$-dependence in Table 1 and the phase $`\delta _2`$-dependence in Table 2 of $`P(\nu _\mu \nu _e),P(\nu _e\nu _\tau ),P(\nu _\mu \nu _\tau ),\mathrm{\Delta }P(\nu _\mu \nu _e),\mathrm{\Delta }P(\nu _e\nu _\tau )`$ and $`\mathrm{\Delta }P(\nu _\mu \nu _\tau )`$.
Here we comment on the matter effect on the oscillation probability difference in the four-neutrino model. By using the Minakata-Nunokawa procedure, the probability difference with the matter effect is expressed for the long-baseline $`\nu _\alpha \nu _\beta `$ oscillation in the four-neutrino model with mass scheme of the two nearly degenerate pairs separated with the order of 1 eV as follows,
$`\mathrm{\Delta }P_{\alpha \beta }`$ $``$ $`4\mathrm{I}\mathrm{m}(U_{\alpha 4}^{}U_{\alpha 3}U_{\beta 3}^{}U_{\beta 4})\mathrm{cos}B_{34}\mathrm{sin}\left({\displaystyle \frac{\mathrm{\Delta }m^2L}{2E}}\right)`$ (23)
$`+4\mathrm{R}\mathrm{e}(U_{\alpha 4}^{}U_{\alpha 3}U_{\beta 3}^{}U_{\beta 4})\mathrm{sin}B_{34}\mathrm{sin}\left({\displaystyle \frac{\mathrm{\Delta }m^2L}{2E}}\right)`$
$`8{\displaystyle \underset{j>i}{}}\mathrm{Re}(UUU\delta V)_{\alpha \beta ;ij}\mathrm{cos}^2\left({\displaystyle \frac{B_{ij}}{2}}\right)\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m_{ij}^2L}{4E}}\right),`$
where
$$B_{ij}=(|U_{ei}|^2|U_{ej}|^2)aL+(|U_{si}|^2|U_{sj}|^2)a^{}L,$$
(24)
$`(UUU\delta V)_{\alpha \beta ;ij}`$ $`=`$ $`U_{\alpha i}^{}U_{\alpha j}U_{\beta j}^{}\delta V_{\beta i}+U_{\alpha i}^{}U_{\alpha j}\delta V_{\beta j}^{}U_{\beta i}`$ (25)
$`+U_{\alpha i}^{}\delta V_{\alpha j}U_{\beta j}^{}U_{\beta i}+\delta V_{\alpha i}^{}U_{\alpha j}U_{\beta j}^{}U_{\beta i}.`$
In Eq.(24), the quantity $`a`$ represents the matter effect for $`\nu _e`$ and we take $`a=1.04\times 10^{13}`$ eV for the constant matter density of 2.7 $`\mathrm{g}/\mathrm{cm}^3`$ , and $`a^{}`$ represents the one for $`\nu _s`$ and we take $`a^{}=a/2`$ . In Eq.(25), $`\delta V_{\alpha i}`$ is given by
$$\delta V_{\alpha i}=\underset{ji}{}\frac{2E}{\mathrm{\Delta }m_{ij}^2}U_{\alpha j}(U_{ej}^{}U_{ei}a+U_{sj}^{}U_{si}a^{}).$$
(26)
In Eq.(23), the first term represents the genuine CP-violation effect corrected by the matter effect, the second term does the CP-violation effect coming from the phase evolution of the neutrino wave function in the matter, and the third one results from the corrections to the mixing matrix $`U`$ due to the existence of matter.
We estimate these matter effects for the $`\nu _\mu \nu _\tau `$ oscillation. The first term of Eq.(23) is almost the genuine CP-violation effect, since the magnitude of the matter effect $`B_{34}`$ is at most $`1\times 10^3`$ for the above-mentioned parameter values of the mixing angles and phases. The second term is approximately $`0.4\times 10^3`$, since $`\mathrm{Re}(U_{\mu 4}^{}U_{\mu 3}U_{\tau 3}^{}U_{\tau 4})0.1`$ and $`\mathrm{sin}B_{34}1\times 10^3`$. This should be compared with the possible maximum value of the genuine CP-violation effect displayed in Fig.5 and Table 1, that is, $`|\mathrm{\Delta }P(\nu _\mu \nu _\tau )|0.3`$. The third term of Eq.(23) is expressed as
$$8\underset{i=1,2,j=3,4}{}\mathrm{Re}(UUU\delta V)_{\mu \tau ;ij}\mathrm{sin}^2\left(\frac{\mathrm{\Delta }M^2L}{4E}\right)8\mathrm{R}\mathrm{e}(UUU\delta V)_{\mu \tau ;34}\mathrm{sin}^2\left(\frac{\mathrm{\Delta }m^2L}{4E}\right),$$
(27)
since $`\mathrm{cos}^2\left(B_{ij}/2\right)1.0`$. The coefficient of the first term of Eq.(27) is given by
$`{\displaystyle \underset{i=1,2,j=3,4}{}}`$ $`\mathrm{Re}`$ $`(UUU\delta V)_{\mu \tau ;ij}`$ (28)
$`=`$ $`4{\displaystyle \frac{2Ea}{\mathrm{\Delta }M^2}}\mathrm{Re}[(U_{\mu 3}^{}U_{\tau 3}+U_{\mu 4}^{}U_{\tau 4})\{(U_{\mu 3}U_{e3}^{}+U_{\mu 4}U_{e4}^{})(U_{\tau 3}^{}U_{e3}+U_{\tau 4}^{}U_{e4})`$
$`+{\displaystyle \frac{1}{2}}(U_{\mu 3}U_{s3}^{}+U_{\mu 4}U_{s4}^{})(U_{\tau 3}^{}U_{s3}+U_{\tau 4}^{}U_{s4})\}],`$
where the relation $`a^{}=a/2`$ is used, and the terms with $`1/\mathrm{\Delta }m^2`$ and $`1/\mathrm{\Delta }m_{\mathrm{solar}}^2`$ do not appear due to the symmetry of the mass scheme of the four neutrinos adopted in our model. The coefficient of the second term of Eq.(27) is given by
$`\mathrm{Re}`$ $`(UUU\delta V)_{\mu \tau ;34}`$ (29)
$`={\displaystyle \frac{2Ea}{\mathrm{\Delta }M^2}}\mathrm{Re}[\{U_{\mu 3}U_{\mu 4}^{}(|U_{\tau 3}|^2+|U_{\tau 4}|^2)+U_{\tau 3}U_{\tau 4}^{}(|U_{\mu 3}|^2+|U_{\mu 4}|^2)\}(U_{e3}^{}U_{e4}`$
$`+{\displaystyle \frac{1}{2}}U_{s3}^{}U_{s4})+2U_{\mu 3}^{}U_{\mu 4}U_{\tau 3}U_{\tau 4}^{}(|U_{e3}|^2+|U_{e4}|^2+{\displaystyle \frac{1}{2}}|U_{s3}|^2+{\displaystyle \frac{1}{2}}|U_{s4}|^2)]`$
$`+{\displaystyle \frac{2Ea}{\mathrm{\Delta }m^2}}\mathrm{Re}[\{U_{\mu 3}U_{\mu 4}^{}(|U_{\tau 3}|^2|U_{\tau 4}|^2)+U_{\tau 3}U_{\tau 4}^{}(|U_{\mu 3}|^2|U_{\mu 4}|^2)\}(U_{e3}^{}U_{e4}`$
$`+{\displaystyle \frac{1}{2}}U_{s3}^{}U_{s4})].`$
The magnitude of Eq.(28) is estimated to be $`0.9\times 10^4`$, and the magnitude of Eq.(29) is to be $`1.1\times 10^2,0.51\times 10^2,1.3\times 10^3`$ for $`\mathrm{\Delta }m^2=(1.0,2.0,5.0)\times 10^3`$ $`\mathrm{eV}^2`$, respectively. So, the third term of Eq.(23), that is, Eq.(27) for $`\nu _\mu \nu _\tau `$ oscillation is again negligibly small as compared with the possible maximum value of the genuine CP-violation effect. So, the matter effect can be totally neglected in the $`\nu _\mu \nu _\tau `$ oscillation in the four-neutrino model with mass scheme of the two nearly degenerate pairs separated with the order of $`1\mathrm{e}\mathrm{V}`$, as was generally studied for any channels in ref..
As can be seen in Figs.5 and 6 and in Tables 1 and 2, CP violation could be observed as the probability difference between the $`\nu _\mu \nu _\tau `$ and $`\overline{\nu _\mu }\overline{\nu _\tau }`$ oscillations in the four-neutrino model. So, we show in Figs.7 and 8 the oscillation probabilities $`P(\nu _\mu \nu _\tau )`$ and $`P(\overline{\nu _\mu }\overline{\nu _\tau })`$, and their difference $`\mathrm{\Delta }P(\nu _\mu \nu _\tau )`$ as functions of $`\mathrm{\Delta }m^2/E`$ $`[\mathrm{eV}^2/\mathrm{GeV}]`$, respectively, for the long-baseline experiments of the MINOS and CERN-ICARUS types, where $`E`$ is the neutrino energy. In Figs.7 and 8, we have assumed the baseline length as $`L=730`$ km. We can observe from Fig.8 that if the beam energy is taken as 7 GeV, magnitude of the CP violation effect for the $`\nu _\mu \nu _\tau `$ channel could attain as large as $`|\mathrm{\Delta }P|0.22`$ in the case of $`\mathrm{\Delta }m^23.5\times 10^3`$ $`\mathrm{eV}^2`$. Incidentally, we display the probabilities $`P(\nu _\mu \nu _e)`$ and $`P(\overline{\nu _\mu }\overline{\nu _e})`$ in Fig.9 and $`\mathrm{\Delta }P(\nu _\mu \nu _e)`$ by a dashed curve in Fig. 8. We show in Table 3 the $`\mathrm{\Delta }m^2/E`$-dependence of $`P(\nu _\mu \nu _e),P(\nu _e\nu _\tau ),P(\nu _\mu \nu _\tau ),\mathrm{\Delta }P(\nu _\mu \nu _e),\mathrm{\Delta }P(\nu _e\nu _\tau )`$ and $`\mathrm{\Delta }P(\nu _\mu \nu _\tau )`$ for $`L=730`$ km.
V Conclusion
We have derived the constraints on the neutrino mixing matrix by using the data from the solar neutrino deficit, atmospheric neutrino anomaly, LSND oscillation experiments, Bugey experiment and the CHORUS and NOMAD experiments along with the other accelerator and reactor experiments in the four-neutrino model with mass scheme of the two nearly degenerate pairs separated with the order of 1 eV. We have used the most general parametrization of the mixing matrix with six mixng angles and six phases applicable to both Majorana and Dirac neutrinos and have obtained the two serious constraints about the four of the six mixing angles and no constraints on the phases.
By using these constraints, we have calculated the oscillation probabilities of $`P(\nu _\mu \nu _e),P(\nu _e\nu _\tau )`$ and $`P(\nu _\mu \nu _\tau )`$ and have investigated CP violation in the long-baseline neutrino oscillations of $`\mathrm{\Delta }P(\nu _\mu \nu _e),\mathrm{\Delta }P(\nu _e\nu _\tau )`$ and $`\mathrm{\Delta }P(\nu _\mu \nu _\tau )`$. The quantity $`\mathrm{\Delta }P(\nu _\mu \nu _\tau )`$ is found to be able to attain a value as large as $`\pm 0.28`$ due to the large mixing between $`\nu _\mu `$ and $`\nu _\tau `$ and the mass scheme of the four neutrinos and, therefore, it could be observed in the long-baseline experiments. We have shown that the contribution to $`\mathrm{\Delta }P(\nu _\mu \nu _\tau )`$ from the matter effect is at most 0.01, 0.005 and 0.001 in magnitude for $`\mathrm{\Delta }m^2=1.0\times 10^3,2.0\times 10^3`$ and $`5.0\times 10^3\mathrm{eV}^2`$, respectively. So,we can conclude that the matter effect is negligibly small in comparison with the possible maximum value of the genuine CP-violation effect of $`|\mathrm{\Delta }P(\nu _\mu \nu _\tau )|=0.28`$.
Figure captions
Fig.1. The oscillation probability of $`\nu _\mu \nu _e`$(solid curve) and $`\overline{\nu _\mu }\overline{\nu _e}`$(dashed curve) with respect to the phase $`\delta _1`$ of the mixing matrix for the long-baseline experiment. The other angles and phases are fixed as $`s_{02}=s_{03}=0.11(c_{02}=c_{03}=0.994),s_{12}=0.91(c_{12}=0.415),s_{13}=0.67(c_{13}=0.742),s_{01}=s_{23}=1/\sqrt{2}(c_{01}=c_{23}=1/\sqrt{2}),\delta _{01}=\delta _{02}=\delta _{03}=\delta _{12}=0`$ and $`\delta _2=\pi /2`$.
Fig.2. The oscillation probability of $`\nu _e\nu _\tau `$(solid curve) and $`\overline{\nu _e}\overline{\nu _\tau }`$(dashed curve) with respect to the phase $`\delta _2`$ of the mixing matrix for the long-baseline experiment. The other angles and phases are the same as in Fig.1 except for $`\delta _1=\pi /2`$ fixed.
Fig.3. The probability difference $`\mathrm{\Delta }P(\nu _\mu \nu _e)`$ (solid curve) and $`\mathrm{\Delta }P(\nu _e\nu _\tau )`$(dashed curve) with respect to the phase $`\delta _1`$ for the long-baseline experiment. The other angles and phases are the same as in Fig.1.
Fig.4. The oscillation probability of $`\nu _\mu \nu _\tau `$ (solid curve) and $`\overline{\nu _\mu }\overline{\nu _\tau }`$(dashed curve) with respect to the phase $`\delta _2`$ for the long-baseline experiment. The other angles and phases are the same as in Fig.2.
Fig.5. The probability difference $`\mathrm{\Delta }P(\nu _\mu \nu _\tau )`$ (solid curve) and $`\mathrm{\Delta }P(\nu _e\nu _\tau )`$ (dashed curve) with respect to the phase $`\delta _2`$ for the long-baseline experiment. The other angles and phases are the same as in Fig.2.
Fig.6. The probability difference $`\mathrm{\Delta }P(\nu _\mu \nu _\tau )`$ (solid curve) and $`\mathrm{\Delta }P(\nu _e\nu _\tau )`$ (dashed curve) with respect to the angle $`\theta _{23}`$ for the long-baseline experiment. The other angles and phases are the same as in Fig.1 except for $`\delta _1=\pi /2`$ fixed.
Fig.7. The oscillation probability of $`\nu _\mu \nu _\tau `$ (solid curve) and $`\overline{\nu _\mu }\overline{\nu _\tau }`$(dashed curve) with respect to $`\mathrm{\Delta }m^2/E`$ for the long-baseline experiment with the distance of $`L=730`$ km. The angles and phases are the same as in Fig.1 except for $`\delta _1=\pi /2`$ fixed.
Fig.8. The probability difference $`\mathrm{\Delta }P(\nu _\mu \nu _\tau )`$ (solid curve) and $`\mathrm{\Delta }P(\nu _\mu \nu _e)`$(dashed curve) with respect to $`\mathrm{\Delta }m^2/E`$ for the long-baseline experiment with the distance of $`L=730`$ km. The angles and phases are the same as in Fig.7.
Fig.9. The oscillation probability of $`\nu _\mu \nu _e`$(solid curve) and $`\overline{\nu _\mu }\overline{\nu _e}`$(dashed curve) with respect to $`\mathrm{\Delta }m^2/E`$ for the long-baseline experiment with the distance of $`L=730`$ km. The angles and phases are the same as in Fig.7.
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# 1 Introduction
## 1 Introduction
Let us consider the following classical hamiltonian model of a system of rotators:
$$H=\frac{1}{2}\underset{i=1}{\overset{N}{}}L_i^2+\frac{1}{2}\underset{i,j=1}{\overset{N}{}}\left[1\mathrm{cos}(\theta _i\theta _j)\right]=K+V.$$
(1)
The potential energy $`V`$ is not thermodynamically stable and the ensemble averaged energy density $`U=\frac{H}{N}`$ diverges in the thermodynamic limit (TL) . If the potential energy term is divided by $`N`$, then the energy density becomes intensive and it is bounded as $`N`$ goes to infinity.
Indeed, dynamics and thermodynamics of the $`1/N`$ rescaled model has been extensively investigated ; in particular, Ruffo and Antoni, who called it the hamiltonian mean field X-Y model (HMF), solved it in the canonical ensemble, and compared the theoretical caloric ($`TvsU`$) and magnetization ($`MvsU`$) curves with those obtained from a microcanonical simulation .
Here we consider a generalization of model (1):
$$H=\frac{1}{2}\underset{i=1}{\overset{N}{}}L_i^2+\frac{1}{2}\underset{ij}{\overset{N}{}}\frac{1\mathrm{cos}(\theta _i\theta _j)}{r_{ij}^\alpha }.$$
(2)
The rotators are placed at the sites of a lattice and the interaction between rotators $`i`$ and $`j`$ decays as the inverse of their distance to the power $`\alpha `$.
A onedimensional version of model (2) has been studied by Anteneodo and Tsallis , who have numerically measured the largest Lyapounov exponent, as a function of $`N`$ and $`\alpha `$. Through a rescaling factor $`N^{}=\frac{N^{1\alpha }1}{1\alpha }`$ Anteneodo and Tsallis showed that their results coincide with those previously obtained for the HMF ($`\alpha =0`$) model; this rescaling could then give a well defined TL to model (2).
In a recent paper Tamarit and Anteneodo, using a rescaling factor $`\stackrel{~}{N}=2^\alpha \frac{N^{1\alpha }1}{1\alpha }`$, have shown that the caloric and magnetization curves of model (2) in one dimension collapse onto the curves of the HMF model . This universality emerges plotting $`T/\stackrel{~}{N}`$ as a function of $`H/N\stackrel{~}{N}`$ and $`M`$ as a function of $`H/N\stackrel{~}{N}`$, from molecular dynamics simulation of model (2) for different $`N`$ and $`\alpha `$ values. These authors conjecture that the results they obtained in the onedimensional case might be general, valid in any dimension $`d`$ and for $`\alpha <d`$, as suggested also in .
## 2 Partition function
In this work, inspired by and , we analytically compute the partition function of an $`\stackrel{~}{N}`$-rescaled model (2) for any $`d`$ and $`\alpha <d`$. In formula (21) we give the right expression of the rescaling function $`\stackrel{~}{N}`$, to obtain universal state curves for all lattice models with long range ($`\alpha <d`$) interactions.
Let us now rewrite the rescaled version of Hamiltonian (2):
$`H`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}L_i^2+{\displaystyle \frac{1}{2\stackrel{~}{N}}}{\displaystyle \underset{i,j=1}{\overset{N}{}}}{\displaystyle \frac{1\mathrm{cos}(\theta _i\theta _j)}{r_{ij}^\alpha }}`$ (3)
$`h_x{\displaystyle \underset{i=1}{\overset{N}{}}}m_{ix}h_y{\displaystyle \underset{i=1}{\overset{N}{}}}m_{iy},`$
where we have introduced an external magnetic field $`𝐡=(h_x,h_y)`$ of modulus $`h`$, that makes possible to compute the magnetization. The indexes $`i,j`$ label the sites of a $`d`$-dimensional generic lattice; $`r_{ij}`$ is the distance between them, with periodic boundary conditions and nearest image convention (the definition of $`r_{ii}`$ will be given shortly); $`\alpha 0`$. At each site a classical rotator (X-Y spin) of unit momentum of inertia is represented by conjugate canonical coordinates $`(L_i,\theta _i)`$, where the $`L_i`$’s are angular momenta, and the $`\theta _i`$’s $`[0,2\pi )`$ are the angles of rotation on a family of parallel planes, each one defined at each lattice point; $`x`$ and $`y`$ refer to the components of boldface twodimensional vectors defined over these planes. To each lattice site a spin vector
$$𝐦_i=(m_{ix},m_{iy})=(\mathrm{cos}\theta _i,\mathrm{sin}\theta _i)$$
(4)
is associated, and the total magnetization is given by:
$$𝐌=(M_x,M_y)=\frac{1}{N}\underset{i=1}{\overset{N}{}}𝐦_i.$$
(5)
Note in (3) the rescaling factor $`\stackrel{~}{N}`$ in front of the potential energy term, now written as a free double sum over both indexes. $`\stackrel{~}{N}`$ should be regarded as an unknown function of $`N,\alpha ,d`$ and the geometry of the lattice, with the fundamental property of making
$$\frac{1}{\stackrel{~}{N}}\underset{j,ji}{}\frac{1}{r_{ij}^\alpha }$$
(6)
an intensive quantity; this guarantees the thermodynamic stability of the potential. We also note that the sum in (6) is independent of the origin $`i`$ because of periodic conditions. To reproduce the usual HMF it is also necessary that $`\stackrel{~}{N}(N,\alpha =0,d)=N`$. The constraint $`ij`$ over the double sum is removed defining $`r_{ii}^\alpha =1/b`$, a finite number. Since the numerator $`1\mathrm{cos}(\theta _i\theta _j)`$ is zero for $`i=j`$ the choice of $`b`$ is free. The removal of the constraint allows to introduce the distance matrix $`R_{ij}^{}=\frac{1}{r_{ij}^\alpha }`$; the diagonalization of such matrix is the key point to obtain, in the computation of the partition function, known integrals in the variables $`\theta _i`$.
As usual the partition function factorizes in a kinetic part $`Z_K=\left(\frac{2\pi }{\beta }\right)^{\frac{N}{2}}`$, where $`\beta =1/k_BT`$, and a potential part $`Z_V`$. After defining $`R_{ij}=\frac{\beta }{2\stackrel{~}{N}}R_{ij}^{}`$, $`𝐁=\beta 𝐡`$, $`C=\mathrm{exp}\left(\frac{\beta }{2\stackrel{~}{N}}_{ij}\frac{1}{r_{ij}^\alpha }\right)`$, the potential part can be written as:
$$Z_V=C_\pi ^\pi d^N\theta \mathrm{exp}\left[\underset{i,j,\mu }{}m_{i\mu }R_{ij}m_{j\mu }+\underset{i}{}B_\mu m_{i\mu }\right],$$
(7)
where $`\mu =x,y`$. Diagonalizing the symmetric matrix $`R=(R_{ij})`$ with the unitary matrix $`U`$ such that $`R=U^TDU`$, $`D=(R_i\delta _{ij})`$, where $`R_i`$ are the eigenvalues of $`R`$, we can write the first part of the exponent in (7) as:
$$\underset{ij}{}\left(m_{ix}R_{ij}m_{jx}+m_{iy}R_{ij}m_{jy}\right)=\underset{i}{}\left(n_{ix}^2R_i+n_{iy}^2R_i\right),$$
(8)
where $`n_{i\mu }=_jU_{ij}m_{j\mu }`$. In order to apply the gaussian transformation:
$$e^{aS^2}=\frac{1}{\sqrt{4\pi a}}_{\mathrm{}}^+\mathrm{}𝑑ze^{\frac{z^2}{4a}+Sz}a>0$$
(9)
to each term of the sum in the right hand side of (8), each $`R_i`$ must be positive. The spectrum can be explicitly computed using a $`d`$-dimensional Fourier transform of matrix $`R`$, the eigenvalues being labelled by vectors of the reciprocal lattice. These eigenvalues are trivially related to those of matrix $`R^{}`$. A study of the spectrum of $`R^{}`$ in the limit $`N\mathrm{}`$ and for $`b=0`$ shows that: when $`\alpha >d`$ each element of the spectrum converges to a finite quantity, the least eigenvalues being negative and of order one in modulus; when $`\alpha <d`$ a part of the spectrum converges to a finite quantity, another part diverges to $`+\mathrm{}`$, at most as $`\stackrel{~}{N}`$. However this last part consists of a fraction of the total number of eigenvalues which goes to zero in the limit $`N\mathrm{}`$. The least eigenvalue is still negative and of order one in modulus. Then part of the spectrum is negative, but it is easily seen that it is shifted by $`b`$. Thus calling $`p`$ the least eigenvalues of $`R^{}`$ for $`b=0`$ and choosing
$$b=p+ϵϵ>0,$$
(10)
we have that with this $`b`$ the whole spectrum of $`R^{}`$ (and therefore that of $`R`$) becomes positive. Then for each $`i=1,\mathrm{},N`$, $`\mu =x,y`$ we can apply (9) with the correspondence $`aR_i`$, $`Sn_{i\mu }`$, $`zz_{i\mu }`$. Performing the integrals over variables $`\theta _i`$ and using the transformation $`z_{i\mu }=2_j(UR)_{ij}\mathrm{\Psi }_{j\mu }`$ with Jacobian $`2^NdetR`$, we can rewrite the partition function as:
$`Z`$ $`=`$ $`CZ_K{\displaystyle \frac{detR}{\pi ^N}}{\displaystyle _{\mathrm{}}^+\mathrm{}}d^N\mathrm{\Psi }_xd^N\mathrm{\Psi }_y`$
$`e^{N\left[_{ij\mu }\mathrm{\Psi }_{i\mu }\frac{R_{ij}}{N}\mathrm{\Psi }_{j\mu }+\frac{1}{N}_l\mathrm{ln}\left(2\pi I_0\left(|2_jR_{lj}𝚿_j+𝐁|\right)\right)\right]}`$
where $`I_0`$ is the zeroth order modified Bessel function. The isolation of the $`N`$ factor in the exponential prepares the object for the use of the saddle point method. The quantity in square brackets is intensive. Double sums in the first two terms are compensated by $`R/N=(\beta /2N\stackrel{~}{N})R^{}`$ and the last sum has $`1/N`$ in front of it. The argument of $`I_0`$ is also intensive because involves a term of the form $`_jR_{lj}=(\beta /2\stackrel{~}{N})_jR_{lj}^{}`$. If we call $`f(w)`$ the function in square brackets, where $`w=(\mathrm{\Psi }_{1x},\mathrm{},\mathrm{\Psi }_{Nx},\mathrm{\Psi }_{1y},\mathrm{},\mathrm{\Psi }_{Ny})`$, then the application of the method requires the following three conditions: $`f(w)`$ admits a stationary point $`w_0`$; $`w_0`$ is a simple stationary point, i.e., $`detHef|_00`$, where $`Hef|_0`$ is the hessian matrix of $`f`$ in $`w_0`$; the path of integration can be deformed (generally going into $`𝒞^{2N}`$) into a path that passes through $`w_0`$ following the steepest descent of $`f(w)`$ and such that $`f(w)<f(w_0)`$ throughout the all path. If the point $`w_0`$ is a maximum no deformation is necessary and the method is also called the Laplace method. Since, as we show below, $`w_0`$ is indeed a real-valued maximum, we readily obtain for the free energy per particle $`F`$:
$`\beta F`$ $`=`$ $`\underset{N\mathrm{}}{lim}{\displaystyle \frac{\mathrm{ln}Z}{N}}=\underset{N\mathrm{}}{lim}\{{\displaystyle \frac{1}{2}}\mathrm{ln}\left({\displaystyle \frac{2\pi }{\beta }}\right){\displaystyle \frac{\beta }{2\stackrel{~}{N}}}{\displaystyle \underset{j}{}}{\displaystyle \frac{1}{r_{ij}^\alpha }}`$ (12)
$`+\underset{w}{\mathrm{max}}[f(w)]+{\displaystyle \frac{1}{N}}\mathrm{ln}{\displaystyle \frac{detR}{\sqrt{det\left(\frac{N}{2}Hef|_0\right)}}}\}.`$
The stationary point $`w_0`$ is given by the vector $`(\mathrm{\Psi }_x,\mathrm{},\mathrm{\Psi }_x,\mathrm{\Psi }_y,\mathrm{},\mathrm{\Psi }_y)`$, homogeneous on the lattice sites. Defining $`𝚿=(\mathrm{\Psi }_x,\mathrm{\Psi }_y)`$, its direction is that of $`𝐁`$, and its modulus $`\mathrm{\Psi }`$ is given by the solution of:
$$\mathrm{\Psi }=\frac{I_1}{I_0}\left(\beta \left[A\mathrm{\Psi }+h\right]\right),$$
(13)
with
$$A=\frac{1}{\stackrel{~}{N}}\underset{j}{}R_{ij}^{}=\frac{1}{\stackrel{~}{N}}\left[b+\underset{ji}{}\frac{1}{r_{ij}^\alpha }\right],$$
(14)
and where $`I_1`$ is the first order modified Bessel function. In (14) $`A`$ does not depend on $`i`$ because of the periodic boundary conditions. We note that when $`h=0`$ we have infinitely many degenerate solutions, since only the modulus $`\mathrm{\Psi }`$ is determined. Evaluation of the elements of the hessian matrix at the stationary point gives:
$$\frac{N}{2}\frac{^2f}{\mathrm{\Psi }_{i\mu }\mathrm{\Psi }_{j\nu }}|_0=R_{ij}\delta _{\mu \nu }(R^2)_{ij}g_{\mu \nu }(w_0)$$
(15)
where we do not give the explicit expression of $`g_{\mu \nu }(w_0)`$. As we will see shortly, the eigenvalues analysis of the hessian matrix (15) shows that the stationary point $`w_0`$ is a maximum. Then, Laplace method applies and Eq. (12) is valid. However, only in the long range case ($`\alpha <d`$) the last term in the rightmost side of (12) is zero; when $`\alpha >d`$ its expression does not appear to be manageable. We will comment on this point later. Restricting then to $`\alpha <d`$, and computing the derivative of (12) with respect to the magnetic field we find that the magnetization $`M=|𝐌|`$ is given by the solution $`\mathrm{\Psi }`$ of (13). Then the internal energy $`U`$ is given by:
$$U=\frac{(\beta F)}{\beta }=\frac{1}{2\beta }+\frac{A}{2}(1M^2)hM.$$
(16)
Equations (13) and (16) are the same as those of HMF, as soon as a proper $`\stackrel{~}{N}`$ rescaling gives
$$A=\frac{1}{\stackrel{~}{N}}\underset{j}{}R_{ij}^{}=\frac{1}{\stackrel{~}{N}}\left[b+\underset{ji}{}\frac{1}{r_{ij}^\alpha }\right]=1.$$
(17)
Now, from equations (15) and (17), and calling $`\lambda _n`$ the eigenvalues of $`R^{}`$, we find that, choosing $`𝐁`$ along one of the coordinate axes, the eigenvalues of the hessian matrix at the stationary point are given by:
$`\chi _n^{(1)}`$ $`=`$ $`{\displaystyle \frac{\beta }{2}}{\displaystyle \frac{\lambda _n}{\stackrel{~}{N}}}\left[1\left(\beta \mathrm{\Psi }^2\beta {\displaystyle \frac{\mathrm{\Psi }}{\mathrm{\Psi }+h}}\right){\displaystyle \frac{\lambda _n}{\stackrel{~}{N}}}\right]`$ (18)
$`\chi _n^{(2)}`$ $`=`$ $`{\displaystyle \frac{\beta }{2}}{\displaystyle \frac{\lambda _n}{\stackrel{~}{N}}}\left[1{\displaystyle \frac{\mathrm{\Psi }}{\mathrm{\Psi }+h}}{\displaystyle \frac{\lambda _n}{\stackrel{~}{N}}}\right]n=1,\mathrm{},N`$
Following our previous analysis we have that:
$$\frac{ϵ}{\stackrel{~}{N}}\frac{\lambda _n}{\stackrel{~}{N}}1.$$
(19)
Then we immediately see that $`\chi _n^{(2)}`$ are all positive for any $`\beta `$ and $`h`$; for $`\chi _n^{(1)}`$ we need to include $`\mathrm{\Psi }(\beta ,h)`$ from (13). We have checked numerically that the quantity in round brackets in (18) is always smaller than $`1`$, and therefore $`\chi _n^{(1)}`$ are also all positive. From (18) we can derive an expression for the determinant of matrix (15). It is given by:
$`{\displaystyle \frac{1}{N}}`$ $`\mathrm{ln}det\left({\displaystyle \frac{N}{2}}Hef|_0\right)={\displaystyle \frac{2}{N}}\mathrm{ln}detR`$
$`+`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{n=1}{\overset{N}{}}}\{\mathrm{ln}[1(\beta \mathrm{\Psi }^2\beta {\displaystyle \frac{\mathrm{\Psi }}{\mathrm{\Psi }+h}}){\displaystyle \frac{\lambda _n}{\stackrel{~}{N}}}]`$
$`+`$ $`\mathrm{ln}[1{\displaystyle \frac{\mathrm{\Psi }}{\mathrm{\Psi }+h}}{\displaystyle \frac{\lambda _n}{\stackrel{~}{N}}}]\}.`$ (20)
When $`\alpha <d`$ most of $`\frac{\lambda _n}{\stackrel{~}{N}}`$ go to zero for $`N\mathrm{}`$, then the sum in (2) is effectively constituted by the terms with the remaining $`\frac{\lambda _n}{\stackrel{~}{N}}`$. These terms are a fraction of $`N`$ that, as we already pointed out, goes to zero when $`N\mathrm{}`$. If we call $`N^{}(N)`$ this fraction, then the sum in (2) can be bounded from above by $`\frac{N^{}}{N}c0`$ for $`N\mathrm{}`$, where $`c`$ is a finite number. Therefore the last term in (12) is zero. When $`\alpha >d`$ all terms contribute to the sum in (2), and we can not give a meaningful expression for (12). At the end of the calculations we can let $`ϵ0`$ in (10).
Then we have shown that any model with $`\alpha <d`$ on any lattice is equivalent to HMF. From (17) we get an exact expression for $`\stackrel{~}{N}`$:
$$\stackrel{~}{N}=p+\underset{ji}{}\frac{1}{r_{ij}^\alpha }.$$
(21)
We have made a microcanonical simulation of Hamiltonian (3) on a threedimensional simple cubic lattice in zero magnetic field, using a fourth order simplectic algorithm with time step $`0.02`$, selected to have relative energy fluctuations not exceeding $`1/10^6`$. We have chosen a fixed $`N=343=7^3`$, and have simulated various energy densities $`H/N`$ and various $`\alpha <3`$. In Fig. 1 we show that the numerical caloric curves collapse onto the universal HMF curve. The kind of results shown in for a onedimensional lattice, where a slightly different $`\stackrel{~}{N}`$ has been used.
## 3 Conclusions
Going back to the beginning of our discussion: it is now clear that model (2) completely reduces to model (1) for $`\alpha =0`$. In model (1) the range of the interactions is infinite; each rotator interacts with all the others and with the same intensity. To get a well defined TL it is sufficient to divide $`V`$ in (1) by $`N`$, the total numbers of rotators. It is then possible to compute caloric and magnetization curves ; the spatial arrangement of the rotators has no effect on them since the intensity of the interaction is the same for each couple of rotators. In this work we have shown that, when considering model (2), it is possible to take into account the spatial $`d`$-dimensional arrangement of the rotators and the decaying of their mutual interaction through a factor $`\stackrel{~}{N}`$, which is computable for any periodic lattice and any $`\alpha <d`$. Dividing by $`\stackrel{~}{N}`$ the potential energy in (2), the model gets a well defined TL and it is possible to compute state curves which become those of the HMF model with a proper normalization of the constant $`A`$ in (17). The HMF ($`\alpha =0`$) model has revealed peculiar equilibrium and nonequilibrium properties , namely: ensemble inequivalence, metastability, collective oscillations, anomalous diffusion and interesting chaotic properties, both in the ferromagnetic and antiferromagnetic case. On the basis of the thermodynamical equivalence here established it would be interesting to investigate the $`\alpha `$ dependence of all these properties. The study of the Lyapounov exponents in is the first in this direction.
## 4 Aknowledgments
A. G. warmly thanks C. Tsallis for having suggested the study of the long range interacting rotators.
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# Preprint RM3-TH/00-2 Parton picture of inclusive quasi-elastic electron scattering off nuclei
## I Introduction
Inclusive quasi-elastic ($`QE`$) electron-nucleus scattering is widely recognized as an important tool for extracting unique information on the dynamics of nucleons in nuclei, in particular on the nature of short-range and tensor nucleon-nucleon ($`NN`$) correlations . In the $`QE`$ kinematical regions the dominant reaction mechanism is expected to be the elastic scattering from individual nucleons bound in the nucleus. Therefore, the concept of $`y`$-scaling was introduced in , showing that at sufficiently high values of the squared four-momentum transfer $`Q^2`$ ($`=qq=\nu ^2|\stackrel{}{q}|^2`$) the measured $`QE`$ cross section is predicted to scale in terms of the variable $`y`$, which represents the minimum momentum of the struck nucleon along the direction $`\stackrel{}{q}`$ of the virtual photon. Adopting the non-relativistic approximation for the nuclear wave function, the scaling function $`F(y)`$ was shown to be the nucleon longitudinal momentum distribution, $`f(y)`$, which depends only on the ground-state of the target nucleus.
Existing inclusive $`SLAC`$ measurements do not exhibit scaling at $`y>0`$, because of the relevance of competitive mechanisms, like meson-exchange currents at low $`Q^2`$ and nucleon inelastic channels at high $`Q^2`$. However, in the low-energy side of the $`QE`$ peak (corresponding to $`y<0`$) the $`QE`$ mechanism is dominant and the $`SLAC`$ data show indeed an approximate scaling for low negative values of $`y`$. At $`y\text{ }<0.3GeV/c`$, corresponding to the kinematical regions relevant for the investigation of $`NN`$ correlations, the $`y`$-scaling breaks down because of the effects from the final state interactions ($`FSI`$) among the struck nucleon and the residual nuclear system (cf. ). Thus, the standard approach is to extract the scaling function $`F(y)`$ by extrapolation from finite $`Q^2`$ to the limit $`Q^2\mathrm{}`$, where the $`FSI`$ are expected to vanish. However, a careful theoretical analysis has shown that $`F(y)f(y)`$ because of the effects of the nucleon removal energy distribution in the nucleus (binding effects). Since microscopic calculations of the latter are still lacking (except for very light nuclei) and only models are available (cf. ), the y-scaling does not allow to extract the nucleon longitudinal momentum distribution in a model-independent way .
The aim of this work is to present concisely a new approach to scaling for inclusive $`QE`$ electron-nucleus scattering, in which the individual nucleons are treated as the (non point-like) partons of the nucleus. All the target-mass effects, due to the non-vanishing masses of the struck particle (the nucleon mass $`M`$) and of the target system (the nuclear mass $`M_AAM`$), are evaluated and a new scaling variable $`\xi _{QE}`$ is obtained, viz.
$`\xi _{QE}x{\displaystyle \frac{1+\sqrt{1+4M^2/Q^2}}{1+\sqrt{1+4M^2x^2/Q^2}}}`$ (1)
where $`x=Q^2/2M\nu `$ is the Bjorken variable ($`0xM_A/MA`$). Note that the low-energy side of the $`QE`$ peak ($`x1`$) corresponds to $`\xi _{QE}1`$ with $`\xi _{QE}x`$. Moreover, $`\text{lim}_{Q^2\mathrm{}}\xi _{QE}=x`$. The variable $`\xi _{QE}`$ has the physical meaning of the light-cone fraction of the nucleus momentum carried by the struck nucleon and the resulting scaling function is proportional to the nucleon light-cone momentum distribution in the nucleus, which depends only on the target ground-state (cf., e.g., ). The concept of $`\xi _{QE}`$ scaling is then positively checked against the recent inclusive data from $`JLab`$ , which extend the $`Q^2`$ range of the previous $`SLAC`$ measurements up to $`Q^27(GeV/c)^2`$. Finally, the nucleon light-cone momentum distribution in $`{}_{}{}^{56}Fe`$ is extracted from the $`JLab`$ data in a model-independent way.
## II Parton picture of $`QE`$ scattering
As it is well known, the parton model successfully describes the deep inelastic electron-nucleon scattering process in terms of elastic processes on free (point-like) constituents, the partons. At sufficiently large values of $`Q^2`$ the nucleon structure function is predicted to scale in the Bjorken variable $`x`$ and to be directly related to the parton light-cone momentum distribution in the nucleon $`\rho _q(x)`$ as
$`F_2^N(x,Q^2)_{Bj}x{\displaystyle \underset{q}{}}e_q^2\rho _q(x)`$ (2)
where $`e_q`$ is the charge of the parton $`q`$. The predictions of the parton model are fully confirmed by the analysis of the light-cone singularities of the time-ordered product of quark currents. Moreover, using the Operator Product Expansion ($`OPE`$) formalism it has been shown that an approximate scaling can still hold at finite values of $`Q^2`$, provided target-mass corrections are properly taken into account. To this end a new scaling variable $`\xi x`$, known as the Nachtmann variable , has to be introduced, depending on $`Q^2`$ and the target mass (the nucleon mass), viz.
$`\xi ={\displaystyle \frac{2x}{1+\sqrt{1+4M^2x^2/Q^2}}}.`$ (3)
Violations of the $`\xi `$ scaling are related to dynamical higher-twist contributions, which manifest themselves as power suppressed terms (i.e., powers of $`1/Q^2`$), and to logarithmic $`pQCD`$ corrections (cf. ).
Let us now consider in a similar picture the $`QE`$ scattering off nuclei, being inspired by the fact that the $`QE`$ mechanism is an elastic process on the individual nucleons in the nucleus. Thus, at sufficiently large values of $`Q^2`$ we expect that the $`QE`$ contribution to the nuclear structure function (which will be simply denoted hereafter by $`F_2^A`$) can be written in terms of the nuclear Bjorken variable $`x_A=Q^2/2M_A\nu `$ as
$`F_2^A(x_A,Q^2)x_A{\displaystyle \underset{N=1}{\overset{A}{}}}G_N^2(Q^2)\rho _N(x_A)`$ (4)
where $`0x_A1`$, $`\rho _N(x_A)`$ is the nucleon light-cone momentum distribution in the nucleus (appearing also in the nucleonic contribution to the $`EMC`$ effect ), while $`G_N^2(Q^2)`$ generalizes the squared parton charge $`e_q^2`$ to the case of nucleons, viz.
$`G_N^2(Q^2){\displaystyle \frac{G_E^2(Q^2)+\tau G_M^2(Q^2)}{1+\tau }}`$ (5)
with $`\tau Q^2/4M^2`$. For point-like nucleons one has $`G_E(Q^2),G_M(Q^2)e_N`$, implying $`G_N^2(Q^2)e_N^2`$. Note that the distribution $`\rho _N(x_A)`$, appearing in Eq. (4), satisfies naturally both the baryon ($`_{N=1}^A_0^1𝑑x_A\rho _N(x_A)=A`$) and momentum sum rules ($`_{N=1}^A_0^1𝑑x_Ax_A\rho _N(x_A)=1`$), provided the nucleons are the only relevant degrees of freedom in the nucleus (cf. ). Let us point out that Eq. (4) does not exhibit a scaling property because of the $`Q^2`$-dependence of the nucleon form factors. Therefore, in what follows we will consider reduced nuclear structure functions, defined as
$`\widehat{F}_2^A(x_A,Q^2){\displaystyle \frac{F_2^A(x_A,Q^2)}{_{N=1}^AG_N^2(Q^2)}}.`$ (6)
Basically, in the kinematical regions of interest in this work the $`Q^2`$-dependence of $`\widehat{F}_2^A(x_A,Q^2)`$ drops out. Indeed, assuming the dipole-law for $`G_E^N(Q^2)`$ and $`G_M^N(Q^2)`$, one gets
$`\widehat{F}_2^A(x_A,Q^2)x_A{\displaystyle \underset{j=n,p}{}}\stackrel{~}{e}_j^2(Q^2)\rho _j(x_A)`$ (7)
where $`\stackrel{~}{e}_p^2(Q^2)=(1+\tau \mu _p^2)/[1+\tau (\mu _p^2+\mu _n^2N/Z)]`$ and $`\stackrel{~}{e}_n^2(Q^2)=1\stackrel{~}{e}_p^2(Q^2)`$ play the role of squared effective charges for protons and neutrons, respectively. For $`Q^2\text{ }>1(GeV/c)^2`$one has $`\stackrel{~}{e}_p^2(Q^2)1/(1+0.47N/Z)`$.
We have now to find the variable that generalizes the Nachtmann variable $`\xi `$ to the case of $`QE`$ scattering, where both the mass of the struck particle (the nucleon mass $`M`$) and the mass of the target system (the nuclear mass $`M_AAM`$) have to be considered. (For an investigation of nuclear scaling in terms of the variable $`\xi `$ see ). Let us consider first the effects of the target mass (the nuclear mass $`M_A`$), which lead to the following nuclear Nachtmann variable: $`\xi _A=2x_A/[1+\sqrt{1+4M_A^2x_A^2/Q^2}]`$. However, since $`M_Ax_A=Mx`$, the variable $`\xi _A`$ turns out to be simply a rescaled Nachtmann variable, i.e. $`\xi _A=\xi M/M_A\xi /A`$ where now $`0\xi A`$. The effects of the struck mass (the nucleon mass $`M`$) can be easily included in our partonic picture. Indeed, since partons are assumed to be free (and therefore on-mass-shell), one gets the conditions $`(k_N+q)^2=M^2`$ and $`k_N^2=M^2`$. In terms of light-cone variables we can put $`k_N=(k^+,M^2/k^+,\stackrel{}{0}_{})`$ and $`q=(M_A\xi _A,Q^2/M_A\xi _A,\stackrel{}{0}_{})`$, where we have used the relation $`|\stackrel{}{q}|\nu =M\xi =M_A\xi _A`$. Thus, one gets the polynomial equation $`k^{+2}M_A\xi _Ak^+M^2M_A^2\xi _A^2/Q^2=0`$, which implies $`k^+=(M_A\xi _A/2)[1+\sqrt{1+4M^2/Q^2}]`$. Introducing the (rescaled) light-cone fraction $`z=k^+/M`$ ($`0zA`$) we obtain $`z=\xi _{QE}`$ with $`\xi _{QE}`$ given by Eq. (1), representing therefore the appropriate variable in which we expect to observe approximate scaling for the $`QE`$ (reduced) nuclear structure function at finite values of $`Q^2`$. We have to mention that a scaling variable conceptually similar to $`\xi _{QE}`$ was introduced in (b).
Since the variable $`\xi _{QE}`$ allows to take into account kinematical (target-mass) corrections, violations of the $`\xi _{QE}`$ scaling are expected to be related to dynamical effects, like those due to $`FSI`$, which should appear as power suppressed terms, i.e. powers of $`1/Q^2`$ (see next Section). Finally note that the scaling variable $`\xi _{QE}`$ does not depend on the mass number $`A`$ and therefore it is the same for all nuclei.
We now derive the target-mass corrections to the nuclear $`QE`$ structure functions. Let us first consider the case of point-like nucleons and then insert the effects of the nucleon size where appropriate. For point-like nucleons we can adapt the procedure of , where the case of scattering from massive quarks in the nucleon was considered. Let us start from the nuclear forward Compton amplitude (cf. )
$`T_{\mu \nu }^A(P,q)`$ $`=`$ $`{\displaystyle \frac{i}{\pi }}{\displaystyle d^4ze^{iqz}P|𝒯[J_\mu (z)J_\nu (0)]|P}=`$ (8)
$``$ $`T_1^A(\nu ,Q^2)\left(g_{\mu \nu }+{\displaystyle \frac{q_\mu q_\nu }{Q^2}}\right)`$ (9)
$`+`$ $`{\displaystyle \frac{1}{M_A^2}}T_2^A(\nu ,Q^2)\stackrel{~}{P}_\mu \stackrel{~}{P}_\nu `$ (10)
where $`|P`$ is the nuclear ground-state with total four-momentum $`P`$, $`\stackrel{~}{P}_\mu =P_\mu +(qP)q_\mu /Q^2`$ and $`J_\mu (z)`$ is the nuclear current given by $`J_\mu (z)=\overline{\psi }(z)\gamma _\mu \psi (z)`$, with $`\psi (z)`$ being the nucleon field. The singularities generated in the product $`J_\mu (z)J_\nu (0)`$ as $`z0`$ can be treated within the framework of the $`OPE`$ and it has been shown that such an $`OPE`$ can be rewritten as an expansion over Gegenbauer polynomials, viz.
$`3T_1^A(1+{\displaystyle \frac{\nu ^2}{Q^2}})T_2^A`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}{\displaystyle \underset{n2}{}}C_n^1(i{\displaystyle \frac{\nu }{Q}})sin[{\displaystyle \frac{\pi }{2}}(n+1)]`$ (11)
$``$ $`\left({\displaystyle \frac{M_A}{Q}}\right)^n\mu _n^{(1)}(Q^2)`$ (12)
and
$`T_2^A`$ $`=`$ $`{\displaystyle \frac{8}{\pi }}{\displaystyle \underset{n2}{}}C_{n2}^3(i{\displaystyle \frac{\nu }{Q}})sin[{\displaystyle \frac{\pi }{2}}(n+1)]`$ (13)
$``$ $`\left({\displaystyle \frac{M_A}{Q}}\right)^n\mu _n^{(2)}(Q^2)`$ (14)
where $`C_n^\lambda `$ are Gegenbauer polynomials and $`\mu _n^{(i)}(Q^2)`$ are Nachtmann moments ($`i=1,2`$). Using standard dispersion relations for $`T_i^A(\nu ,Q^2)`$ (cf. the nucleon case in ) and orthogonality relations among Gegenbauer polynomials, from Eqs. (11-13) it follows
$`\mu _n^{(1)}={\displaystyle \frac{2M_A}{Q^2}}{\displaystyle 𝑑\nu \xi _A^{n+1}\left[3W_1^A\frac{\nu ^2+Q^2}{Q^2}W_2^A\right]}`$ (15)
and
$`\mu _n^{(2)}`$ $`=`$ $`{\displaystyle \frac{4M_A}{Q^4}}{\displaystyle }d\nu \xi _A^{n+1}W_2^A[(n^2+2n+3)\nu ^2`$ (16)
$`+`$ $`n(n+2)Q^2+3(n+1)\nu \sqrt{\nu ^2+Q^2}]`$ (17)
$`/`$ $`\left[(n+2)(n+3)\right]`$ (18)
where $`W_i^A(\nu ,Q^2)`$ are the structure functions of the absorptive part of $`T_{\mu \nu }^A`$ (i.e., of the nuclear tensor). As noticed in , the Nachtmann moments $`\mu _n^{(i)}(Q^2)`$ can be written as $`\mu _n^{(i)}(Q^2)=\stackrel{~}{\mu }_n^{(i)}(Q^2)P|\widehat{O}_n|P`$, where $`P|\widehat{O}_n|P`$ are the reduced matrix elements of the traceless operators appearing in the $`OPE`$ of the nuclear tensor and yielding the scaling function, while $`\stackrel{~}{\mu }_n^{(i)}(Q^2)`$ are the Nachtmann moments corresponding to the sum of the individual parton contributions in $`W_i^A(\nu ,Q^2)`$, namely: $`W_1^A\tau _{N=1}^Ae_N^2\delta (\nu Q^2/2M)`$ and $`W_2^A_{N=1}^Ae_N^2\delta (\nu Q^2/2M)`$. The resummation of the series over the Gegenbauer polynomials in Eqs. (11-13) can be done in the same way as in and the final result for the nuclear structure function $`\nu W_2^A(\xi _{QE},Q^2)`$ reads as
$`\nu W_2^A={\displaystyle \frac{x^2}{r^3}}\left[(1+{\displaystyle \frac{1}{\tau }}){\displaystyle \frac{F_2^A(\xi _{QE})}{\xi _{QE}^2}}+3T_2^A(\xi _{QE},Q^2)\right]`$ (19)
where $`r\sqrt{1+4M^2x^2/Q^2}`$ and
$`T_2^A`$ $`=`$ $`\sqrt{1+{\displaystyle \frac{1}{\tau }}}{\displaystyle \frac{2M^2x}{Q^2r}}{\displaystyle _{\xi _{QE}}^{\xi _{QE}^{max}}}𝑑\xi _{QE}^{}\left(1{\displaystyle \frac{1}{\xi _{QE}^2}}\right)`$ (20)
$``$ $`{\displaystyle \frac{F_2^A(\xi _{QE}^{})}{\xi _{QE}^2}}+{\displaystyle \frac{4M^4x^2}{Q^4r}}{\displaystyle _{\xi _{QE}}^{\xi _{QE}^{max}}}𝑑\xi _{QE}^{}\left(1{\displaystyle \frac{1}{\xi _{QE}^2}}\right)`$ (21)
$``$ $`{\displaystyle \frac{F_2^A(\xi _{QE}^{})}{\xi _{QE}^2}}\left[\xi _{QE}^{}+{\displaystyle \frac{1}{\xi _{QE}^{}}}\xi _{QE}{\displaystyle \frac{1}{\xi _{QE}}}\right]`$ (22)
with $`\xi _{QE}^{max}=\text{min}[A,Q(1+\sqrt{1+4M^2/Q^2})/2M]`$ (cf. ). The function $`F_2^A(\xi _{QE})`$ is the asymptotic limit of $`\nu W_2^A(\xi _{QE},Q^2)`$, i.e. at fixed $`\xi _{QE}`$ one has
$`\nu W_2^A_{Q^2\mathrm{}}F_2^A(\xi _{QE})=\xi _{QE}{\displaystyle \underset{N=1}{\overset{A}{}}}e_N^2\rho _N(\xi _{QE}).`$ (23)
As for the nuclear response $`W_1^A(\nu ,Q^2)`$, one can note that for the combination $`\{3W_1^AW_2^A[\nu ^2+Q^2]/Q^2\}`$ Eq. (15) is similar to the Nachtmann moment of a scalar current (cf. ). Thus, the target-mass-corrected nuclear response $`W_1^A(\xi _{QE},Q^2)`$ is given by
$`W_1^A={\displaystyle \frac{x}{2Mr}}\left[{\displaystyle \frac{1+1/\tau }{1+R_N}}{\displaystyle \frac{F_2^A(\xi _{QE})}{\xi _{QE}^2}}+T_2^A(\xi _{QE},Q^2)\right]`$ (24)
where $`R_N(Q^2)=1/\tau `$ is the longitudinal to transverse ($`L/T`$) cross section ratio for point-like nucleons.
The modifications of our basic equations (19-24) due to the nucleon size can be argued in the following way. From $`y`$-scaling analysis it is known (cf., e.g., ) that the quantity $`_{N=1}^AG_N^2(Q^2)`$ can be factorized out, so that the scaling function $`F(y)`$ describes the asymptotic nuclear response as the nucleons were point-like in the virtual photon coupling. Therefore, we expect that Eqs. (19-23) still hold by replacing $`\nu W_2^A`$, $`F_2^A`$ and $`T_2^A`$ with their reduced counterparts, defined as in Eq. (6). Correspondingly, in Eq. (24) we expect the replacements of $`W_1^A`$, $`F_2^A`$ and $`T_2^A`$ with their reduced counterparts and, furthermore, the replacement of $`L/T`$ ratio $`R_N(Q^2)=1/\tau `$ with
$`\widehat{R}_N(Q^2){\displaystyle \frac{1}{\tau }}{\displaystyle \frac{\underset{N^{}=1}{\overset{A}{}}[G_E^N^{}(Q^2)]^2}{_{N^{}=1}^A[G_M^N^{}(Q^2)]^2}}.`$ (25)
It is interesting to note that the above expectations may follow from the hypothesis that the Gegenbauer expansions (11-13) hold as well in case of non point-like nucleons. Indeed, within this hypothesis, in constructing the moments $`\stackrel{~}{\mu }_n^{(i)}(Q^2)`$ one has to consider in $`W_i^A(\nu ,Q^2)`$ the sum of the individual nucleon contributions given now by $`W_1^A\tau _{N=1}^A[G_M^N(Q^2)]^2\delta (\nu Q^2/2M)`$ and $`W_2^A_{N=1}^AG_N^2(Q^2)\delta (\nu Q^2/2M)`$. It is then easy to check that in Eqs. (15-16) the quantity $`_{N=1}^AG_N^2(Q^2)`$ does factorize out and that the $`L/T`$ ratio $`R_N(Q^2)=1/\tau `$ has to be replaced by $`\widehat{R}_N(Q^2)`$ given in Eq. (25). In the next Section we will show that our basic equations and the $`\xi _{QE}`$ variable are successful in describing the scaling in the $`QE`$ process, providing therefore a challenge to demonstrate that Eqs. (11-13) hold as well for non point-like nucleons. To this end a possibility (which can be investigated in future works) might be offered by the application of nuclear effective field theories.
Note that according to Eq. (19) the structure function $`\nu W_2^A(\xi _{QE},Q^2)`$ \[$`\nu \widehat{W}_2^A(\xi _{QE},Q^2)`$\] does not depend on the nucleon $`L/T`$ ratio $`R_N(Q^2)`$ \[$`\widehat{R}_N(Q^2)`$\]. This is at variance with the results of the impulse approximation obtained within the instant form of the dynamics (cf., e.g., ), while it agrees with the results of the light-cone impulse approximation (cf. ). This happens because our results are based on a partonic description which is naturally formulated in the light-cone form of the dynamics. As a matter of fact, the impulse approximations in the instant and the light-cone forms should coincide only asymptotically, while they can differ at finite values of $`Q^2`$.
A very relevant feature of the target-mass corrections is that the asymptotic function $`\widehat{F}_2^A(\xi _{QE})`$ can be extracted from the structure functions $`\widehat{W}_i^A(\xi _{QE},Q^2)`$ at finite values of $`Q^2`$, viz.
$`\widehat{F}_2^A(\xi _{QE})`$ $`=`$ $`{\displaystyle \frac{r^3\xi _{QE}^2}{x^2(1+1/\tau )}}{\displaystyle \frac{1+\widehat{R}_N(Q^2)}{2\widehat{R}_N(Q^2)}}`$ (26)
$``$ $`{\displaystyle \frac{2R_A(\xi _{QE},Q^2)}{1+R_A(\xi _{QE},Q^2)}}\nu \widehat{W}_2^A(\xi _{QE},Q^2)`$ (27)
where
$`R_A(\xi _{QE},Q^2){\displaystyle \frac{W_2^A(\nu ,Q^2)}{W_1^A(\nu ,Q^2)}}\left(1+{\displaystyle \frac{\nu ^2}{Q^2}}\right)1`$ (28)
is the nuclear $`QE`$ longitudinal to transverse cross section ratio. In other words Eq. (27) tells us how to cancel out all the target-mass corrections from the structure function $`\nu \widehat{W}_2^A(\xi _{QE},Q^2)`$.
## III Results
Let us now consider the experimental nuclear structure function $`\nu \overline{W}_2^A(\xi _{QE},Q^2)`$, which in terms of the inclusive nuclear cross section $`\sigma _A`$ is given by
$`\nu \overline{W}_2^A`$ $`=`$ $`{\displaystyle \frac{\nu \sigma _A}{\sigma _{Mott}[_{N=1}^AG_N^2(Q^2)]}}`$ (29)
$``$ $`{\displaystyle \frac{1}{1+2tg^2(\frac{\theta _e}{2})(1+\nu ^2/Q^2)/(1+R_A)}}`$ (30)
where $`\sigma _{Mott}`$ is the Mott cross section and $`\theta _e`$ is the electron scattering angle. We do not expect to observe scaling of Eq. (30) in terms of $`\xi _{QE}`$, at least because of the $`Q^2`$-dependence induced by target-mass corrections (see Eq. (19)). Significative scaling violations are indeed present in the iron data of at any $`\xi _{QE}`$, as it is illustrated in Fig. 1, where the reported uncertainties include the small impact of the variation of $`R_A`$ in Eq. (30) between $`0`$ and $`2\widehat{R}_N`$.
Inspired by Eq. (27) let us define the following experimental scaling function:
$`\overline{F}_2^A(\xi _{QE},Q^2)`$ $``$ $`{\displaystyle \frac{r^3\xi _{QE}^2}{x^2(1+1/\tau )}}{\displaystyle \frac{1+\widehat{R}_N(Q^2)}{2\widehat{R}_N(Q^2)}}`$ (31)
$``$ $`{\displaystyle \frac{2R_A(\xi _{QE},Q^2)}{1+R_A(\xi _{QE},Q^2)}}\nu \overline{W}_2^A(\xi _{QE},Q^2).`$ (32)
Therefore, from the results of the previous Section we expect that
$`\overline{F}_2^A(\xi _{QE},Q^2)=\xi _{QE}\overline{\rho }(\xi _{QE})+O({\displaystyle \frac{1}{Q^2}})`$ (33)
where $`\overline{\rho }(\xi _{QE})`$ is the nucleon light-cone momentum distribution averaged over protons and neutrons in the nucleus \[cf. Eq. (7)\], while the power-suppressed terms correspond to higher twists containing unique information about the structure of $`FSI`$ and/or other relevant mechanisms which violate the $`\xi _{QE}`$ scaling. Note that Eq. (33) implies that any $`Q^2`$-dependence of $`\overline{F}_2^A(\xi _{QE},Q^2)`$ signals the breakdown of the impulse approximation.
The results of Eq. (32) evaluated for the iron data of are shown in Fig. 2. It can clearly be seen that scaling violations are present for $`\xi _{QE}<1.1`$ because of the large contributions from nucleon inelastic channels (cf. ), but now the data free from target-mass effects scale nicely for $`1.1\text{ }<\xi _{QE}\text{ }<1.4`$, while for larger $`\xi _{QE}`$ the scaling is partially broken by $`FSI`$ effects.
Applying a power correction analysis to the data of Fig. 2 (after the subtraction of the nucleon inelastic channel contribution evaluated as in ), the asymptotic function $`\widehat{F}_2^A(\xi _{QE})=\xi _{QE}\overline{\rho }(\xi _{QE})`$ can be extracted in a model-independent way and the resulting nucleon light-cone momentum distribution $`\overline{\rho }(\xi _{QE})`$ is shown in Fig. 3. We point out that our procedure is fully relativistic and does not involve the non-relativistic approximation for the nuclear wave function; moreover, being formulated in the light-cone form of the dynamics, it does not suffer from any ambiguities related to the modelling of binding effects and to off-shell prescriptions.
Finally, since the variable $`\xi _{QE}`$ and the Nachtmann variable $`\xi `$ are simply related by $`\xi _{QE}=\xi [1+\sqrt{1+4M^2/Q^2}]/2`$, one could use $`\xi `$ as a scaling variable for the $`QE`$ process (cf. ). From the theoretical point of view it is however clear that the variable $`\xi `$ is appropriate only for elastic processes on massless constituents in a massive target. Consequently, all the machinery of $`QE`$ target-mass corrections developed in this work cannot be formulated in terms of $`\xi `$. Therefore, any approximate phenomenological scaling in $`\xi `$ should be simply reminiscent of the scaling in $`\xi _{QE}`$, which we stress is the appropriate variable for the $`QE`$ process because it includes the mass effects of both the struck particle and the target.
## IV Conclusions
In conclusion, a new approach to scaling in inclusive quasi-elastic electron-nucleus scattering has been concisely presented. It is based on a parton picture of the quasi-elastic process, where individual nucleons are treated as the (non point-like) partons of the nucleus. All the target-mass corrections to asymptotic scaling have been taken into account, leading to a new scaling variable $`\xi _{QE}`$ given by Eq. (1). The $`\xi _{QE}`$ scaling has been positively checked against recent inclusive iron data from Jefferson Lab (see Fig. 2) and the nucleon light-cone momentum distribution in iron (see Fig. 3) has been extracted in a model-independent way.
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# Fermion production during preheating after hybrid inflation
## 1 Introduction
Cosmological inflation is an extremely efficient mechanism in diluting any particle species or fluctuations. At the end of inflation, the universe is empty and extremely cold, dominated by the homogeneous coherent mode of the inflaton. Its potential energy density is converted into particles, as the inflaton field oscillates around the minimum of its potential. These particles are initially very far from equilibrium, but they interact among themselves (or decay) and thermal equilibrium is achieved at a very large temperature. From there on the universe expanded isoentropically, cooling down as it expanded, in the way described by the standard hot Big Bang model. Recent developments in the theory of reheating suggest that the decay of the inflaton energy could have been explosive , due to the coherent oscillations of the inflaton zero mode, which induce its stimulated decay. The result is a resonant production of particles in just a few inflaton oscillations, in a process known as preheating . The number of particles produced in this way is exponentially large, which may account for the extraordinarily large entropy, of order $`10^{89}`$, in our observable patch of the universe today.
Preheating is not generic, it may occur in different models of inflation but only under special circumstances. Preheating strongly depends on the inflaton couplings to other fields, as well as on the amplitude and frequency of the inflaton oscillations. This explosive production occurs through parametric resonance of the long wavelength modes of any field coupled to the inflaton, either bosonic or fermionic . Due to the parametric nature of the resonance the particle production occurs for well defined frequency bands. The difference between bosonic or fermionic preheating is spectral: bosons can have large occupation numbers for a single mode, while fermions saturate due of the Pauli exclusion principle, and therefore a larger fraction of their energy can be transferred to higher momentum modes. Also, the rate of growth of these two components is very different: while the boson energy density grows exponentially , that of fermions cannot . As a consequence, very soon bosons dominate the decay of the inflaton energy if both kinds of particles are resonantly produced. After a few inflaton oscillations, the energy density in fermions and bosons may be large enough to backreact on the inflaton oscillations and hence this will eventually stop their resonant production .
After preheating in chaotic inflation, the main part of the energy density of the universe is in the coherent modes of bosons: the inflaton, as well as any other parametrically amplified bosonic fields with large occupation numbers, while a subdominant part may be in the parametrically produced fermions. This state is very far from thermal equilibrium (characterized by a decoherent ensemble with small occupation numbers for all momentum modes) and, in fact, such a state could be used to generate the required baryon asymmetry during preheating, either via GUT baryogenesis , EW baryogenesis or leptogenesis . At the end, the final thermalization of the universe occurs through the rescattering of all the particles, which breaks the coherence of the bosonic fields. See, for instance, Ref. for thermalization via fermionic modes.
In this paper we study the resonant production of both fermions and bosons in a hybrid model of inflation . In these models it is possible to choose, without fine-tuning, the masses and couplings of the fields in such a way that the rate of expansion is negligible compared to the masses involved . In that case, preheating occurs in a few inflaton oscillations after the end of inflation, before the scale factor has grown by a single e-fold, and therefore we can ignore the universe expansion. This reduces the problem to resonant production of fermions in Minkowski space, for which there is a complete analytic treatment . We find that in the case of hybrid inflation, contrary to what happens in chaotic inflation, it is possible to produce, for certain couplings of the fermions to the inflaton field, a larger contribution of fermionic modes than bosonic ones to the energy density of the universe at the end of preheating.
We describe in Section 2 the hybrid model under consideration, with a Higgs-type field coupled to the inflaton, and discuss the solution of the classical equations of motion. In Section 3 we study the pair production of fermions in this model, for a non-zero coupling of fermions to both the inflaton and the Higgs field. We analyze the growth of the fermion energy density for different values of the couplings, as well as for different fermion masses. In Section 4 we study the boson production for both couplings to inflaton and Higgs, together with the growth of the bosonic energy density. In Section 5 we draw our conclusions.
## 2 Hybrid model
We will consider a model of inflation with two fields, one that slow-rolls down a potential, driving inflation, and another one with a symmetry breaking potential that triggers the end of inflation, which we will call the Higgs. Such a model is called hybrid inflation and was proposed by Linde long ago . There are many particle physics realizations of this general class of models . One particular potential for hybrid inflation is given by the tree level potential <sup>1</sup><sup>1</sup>1See Ref. for a discussion of loop-corrections to this potential during inflation.
$$V(\varphi ,\sigma )=\frac{M^4}{4\lambda }\frac{1}{2}M^2\varphi ^2+\frac{1}{4}\lambda \varphi ^4+\frac{1}{2}g^2\varphi ^2\sigma ^2+\frac{1}{2}m^2\sigma ^2.$$
(1)
During hybrid inflation, the inflaton field $`\sigma `$ evolves along the Higgs valley at $`\varphi =0`$. As soon as the Higgs acquires a negative mass term, it triggers the end of inflation. That is, in less than one e-fold the two fields start oscillating around their common absolute minimum $`(\varphi =v,\sigma =0)`$, where $`v=M/\sqrt{\lambda }`$ is the Higgs vacuum expectation value (vev). The end of inflation occurs when the mass of the Higgs vanishes, i.e. at $`\sigma =\sigma _cM/g`$.
By defining $`y\sigma /\sigma _c`$ and $`x\varphi /v`$, we can write the potential (1) as
$$V(x,y)=V_0[(1x^2)^2+2x^2y^2+2\gamma y^2],$$
where $`V_0=M^4/4\lambda `$, and $`\gamma =\lambda m^2/g^2M^2`$ is a constant that is constrained to be small by the amplitude of the microwave background (CMB) anisotropies. It does not play any significant role after inflation and we will neglect it here. Also, the rate of expansion can be made relatively small in hybrid inflation (as long as $`MgM_\mathrm{P}`$) and still satisfy the CMB constraints, see Ref. . Therefore, we will ignore the rate of expansion of the universe here. We will also concentrate in the regime $`\lambda g^2`$, for which the Higgs evolves along the minimum of the potential, following the inflaton oscillations.
We can then write the evolution equations after inflation, redefining the time unit as $`\tau \overline{M}t`$ (with $`\overline{M}=gM/\sqrt{\lambda }`$), as:
$`y^{\prime \prime }+(1y^2)y`$ $`=`$ $`0,`$ (2)
$`1y^2`$ $`=`$ $`x^2.`$ (3)
For initial conditions $`y(0)=y_0`$ and $`y^{}(0)=0`$, we find the solution
$$\begin{array}{cc}\hfill y=y_0\mathrm{cd}(u|m),& \\ \hfill u=\tau \sqrt{a_2/2},& m=a_1/a_2,\hfill \\ \hfill a_1=y_0^2,& a_2=2y_0^2,\hfill \end{array}$$
(4)
where $`\mathrm{cd}(u|m)`$ is the Jacobi elliptic function. This solution is periodic in the $`\tau `$ variable, with period
$$T=4\sqrt{\frac{2}{a_2}}K(m),$$
(5)
with $`K`$ the elliptic integral. We have plotted the evolution of both the normalized Higgs and the inflaton after inflation in Fig. 1. Throughout the paper we will take $`y_0=0.9999`$. The corresponding period of the inflaton oscillations is $`T29.97`$. Although this period increases with $`y_0`$ approaching one, the mean energy density transferred to fermions coupled to the inflaton is nearly unchanged (provided that $`y_0`$ is close to one). If the fermions are coupled to the Higgs, the results are more sensitive to the value of $`y_0`$ for strong coupling, as we will discuss later.
## 3 Pair production of fermions
The oscillations of the inflaton and Higgs fields at the end of inflation trigger the explosive production of particles that couple to either or both of these fields. The production of bosons in hybrid inflation has been studied in Ref. . Here we will analyze the fermion production in this model and in the next Section will elaborate further on the bosonic case in order to compare the two results.
We will consider first the coupling of fermions to the inflaton $`\sigma `$, with coupling $`h_1\sigma \overline{\psi }\psi `$, and a possible mass term $`m_\psi \overline{\psi }\psi `$. Then we will consider the coupling of fermions to the Higgs, with coupling $`h_2\varphi \overline{\psi }\psi `$. When the symmetry gets broken this last coupling will give the fermions a mass through the non vanishing vev of the Higgs.
### 3.1 Coupling to the inflaton
Let us consider here a fermionic field $`\psi `$ satisfying the Dirac equation
$$\left(i\gamma ^\mu _\mu h_1\sigma (t)m_\psi \right)\psi =0.$$
(6)
The solutions are more easily obtained using an auxiliary field $`X(\stackrel{}{x},t)`$, such that $`\psi =\left(i\gamma ^\mu _\mu +h_1\sigma (t)+m_\psi \right)X`$. Decomposing it as $`\mathrm{exp}(i\stackrel{}{K}\stackrel{}{x})X_K(t)R_r`$, with $`R_r`$ eigenvectors of $`\gamma ^0`$ with eigenvalue $`+1`$, we can write the equation of motion for fermion modes $`X_k`$ as
$$X_k^{\prime \prime }+\mathrm{\Omega }_k^2X_ki\sqrt{q_1}y^{}X_k=0,$$
(7)
where
$`\mathrm{\Omega }_k^2(\tau )`$ $`=`$ $`k^2+(\sqrt{q_1}y(\tau )+\overline{m}_\psi )^2,`$ (8)
$`q_1`$ $``$ $`{\displaystyle \frac{\lambda h_1^2}{g^4}},`$ (9)
$`\overline{m}_\psi `$ $``$ $`m_\psi /\overline{M}`$ (10)
and we have rescaled $`kK/\overline{M}`$. We will here display the results for $`q_1`$ between 1 and $`10^6`$. Notice that if we were to take e.g. $`\lambda =1`$ and $`g=0.01`$ this will lead to $`q_1=10^8h_1^2`$, and hence the values of $`q_1`$ considered would correspond to $`h_1`$ between $`10^4`$ and 0.1.
We take the initial conditions corresponding to positive frequency plane waves at $`\tau <0`$:
$`X_k(0)`$ $`=`$ $`[2\mathrm{\Omega }_k(\mathrm{\Omega }_k+\sqrt{q_1}y_0+\overline{m}_\psi )]^{1/2},`$ (11)
$`X_k^{}(0)`$ $`=`$ $`i\mathrm{\Omega }_kX_k(0).`$ (12)
The quantity of interest to us is the fraction of the total energy, $`\rho _{\mathrm{total}}=V_0`$, which is transferred into fermions,
$$\frac{\rho __F(\tau )}{\rho _{\mathrm{total}}}=h_1^2\frac{2N}{\pi ^2q_1}𝑑kk^2\mathrm{\Omega }_k(\tau )n_k(\tau ),$$
(13)
with the number of fermion-pair degrees of freedom $`N=1`$ or 2 for Majorana or Dirac fields, respectively. The occupation number (for fermion pairs), $`n_k`$, can be calculated as
$$n_k(\tau )=\frac{1}{2}\frac{k^2}{\mathrm{\Omega }_k}\mathrm{Im}[X_kX_{k}^{}{}_{}{}^{}]\frac{\sqrt{q_1}y+\overline{m}_\psi }{2\mathrm{\Omega }_k}.$$
(14)
Note that $`n_k(0)=0`$, thanks to the initial conditions (11) and (12). Also, it is easy to see that the occupation number of fermion pairs is always smaller than one,<sup>2</sup><sup>2</sup>2This is related to the fact that $`n_k=|\beta _k|^2=1|\alpha _k|^21`$, for fermions, in terms of the Bogoliubov transformation coefficients ($`\alpha _k,\beta _k`$), see Ref. . as expected.
Particle production through parametric resonance occurs at those moments when the adiabaticity condition d$`\mathrm{\Omega }_k/`$d$`\tau <\mathrm{\Omega }_k^2`$ is violated, see Ref. . For fermion production this will occur in general whenever the effective fermion mass $`\overline{m}_\psi +\sqrt{q_1}y`$ approaches zero. This is illustrated in Fig. 2, where the continuous lines show the evolution of the exact occupation number for massless fermions and two values of the $`q_1`$ parameter. The occupation numbers $`n_k`$ (solid lines) have jumps every quarter and three quarters of the inflaton period, corresponding to the times when the inflaton value vanishes, $`y=0`$ (see Fig. 1), and therefore when the effective fermion mass also vanishes.
To evaluate Eq. (13) using $`n_k(\tau )`$ from Eq. (14) and the numerical solutions of (7) is however quite demanding, so that it is convenient to use the approximate analytical method developed in . This method exploits the periodicity of $`y`$ and allows to obtain a smooth function
$$\overline{n}_k(\tau )=F_k\mathrm{sin}^2(\nu _k\tau ),$$
(15)
which coincides with $`n_k(\tau )`$ for $`\tau =nT`$. Hence, $`\overline{n}_k(\tau )`$ gives an approximate expression for $`n_k(\tau )`$ without its fine (spiky) details. The approximate solution (dashed lines in Fig. 2) follows the overall oscillations of the occupation number, matching the exact results at every inflaton period (and also at every half period in this case).
The advantage of using $`\overline{n}_k(\tau )`$ is that it has a simple temporal behaviour, while its $`k`$ dependence can be obtained from the knowledge of the functions $`F_k`$ and $`\nu _k`$, where
$`F_k`$ $`=`$ $`{\displaystyle \frac{k^2(\mathrm{Im}X_k^{(1)}(T))^2}{\mathrm{\Omega }_k^2(T)\mathrm{sin}^2(\nu _kT)}},`$ (16)
$`\mathrm{cos}(\nu _kT)`$ $`=`$ $`\pm \mathrm{Re}X_k^{(1)}(T).`$ (17)
and $`X_k^{(1)}`$ satisfies the same equation (7) with the initial condition $`X_k^{(1)}(0)=1`$, $`X_{k}^{(1)}{}_{}{}^{}(0)=0`$. Therefore, to obtain $`\overline{n}_k(\tau )`$ we only need to solve Eq. (7) during one inflaton period. The two signs in Eq. (17) correspond to two possible functions $`\overline{n}_k`$, which oscillate with the same amplitude but different frequency, and both match $`n_k`$ at every period of the inflaton oscillation. The best approximation to $`n_k(\tau )`$ is given by the one with smaller frequency $`\nu _k`$ for the $`k^2`$ values where $`F_k`$ is maximum (i.e. for the momenta contributing significantly to $`\rho __F`$), which in this case corresponds to taking the minus sign in Eq. (17). This is the function plotted in Fig. 4. Notice that the maximum value of $`\nu _k`$ is $`\pi /T0.105`$, as can be deduced from Eq. (17).
We show in Fig. 3 the instability chart for massless fermions coupled to the inflaton. It displays the contours in the $`(q_1,k^2)`$ plane of equal $`F_k`$ values. Fermions are mainly produced with momenta in the darker regions, corresponding to maxima of $`F_k`$. We see that the bands get narrower with increasing $`k^2`$ for a given $`q_1`$ value, and after several bands they shrink to a negligible width.
The upper panel in Fig. 4 shows the behaviour of $`F_k`$ for different $`q_1`$ values extending up to $`q_1=10^6`$. These correspond to cuts in the instability chart (Fig. 3) at a fixed value of $`q_1`$. The maximum momentum for which the bands are sizeable grows as $`k_{\mathrm{max}}q_1^{1/4}`$ for fermions coupled to the inflaton.<sup>3</sup><sup>3</sup>3Note that the same behaviour was found in Ref. for chaotic inflation. However, a different behaviour, $`k_{\mathrm{max}}q^{1/3}`$, was found in Ref. and was attributed, in their case, to the redshift of the modes during the particle production process. The lower panel of Fig. 4 shows $`\nu _k`$, as a function of $`k^2`$, for the same values of $`q_1`$. Note that the maxima of $`F_k`$ correspond to local minima of $`\nu _k`$.
The fraction of the total energy transferred to the fermions then finally results
$$\frac{\rho __F(\tau )}{\rho _{\mathrm{total}}}=h_1^2\frac{2N}{\pi ^2q_1}𝑑kk^2\mathrm{\Omega }_k(\tau )F_k\mathrm{sin}^2(\nu _k\tau ).$$
(18)
Its time evolution is shown in Fig. 5, for different values of the coupling parameter $`q_1`$. We see that after the first oscillation of the inflaton this fraction already reaches its asymptotic value and then fluctuates around it. The asymptotic value $`\rho __F/\rho _{\mathrm{total}}`$ scales as $`q_1^{1/4}h_1^2`$. Moreover, for large values of $`q_1`$ a significant fraction of the inflaton energy can be transferred into fermions (as long as $`h_1`$ is not too small). We have checked that the final density is insensitive to the initial value of the inflaton field $`y_0`$ as long as $`|y_01|10^2`$. This can be simply understood since, for $`m_\psi =0`$, the production of fermions takes place as the inflaton crosses through $`y=0`$.
If fermions are massive, the inflaton energy is transferred to the fermions with nearly the same efficiency as for massless fermions up to a maximum cut off value of the fermion’s mass, above which it drastically drops. Fig. 6 shows this behaviour for different values of the $`q_1`$ parameter. The cut off value of the mass goes like $`\overline{m}_{\psi ,\mathrm{cutoff}}q_1^{1/2}`$<sup>4</sup><sup>4</sup>4The same behaviour was found for chaotic inflation in Ref. . The reason is simple, particle production occurs whenever the effective fermion mass vanishes, $`m_{\mathrm{eff}}=\overline{m}_\psi +\sqrt{q_1}y(\tau )=0`$. For small values of $`q_1`$, even a small (positive) $`\overline{m}_\psi `$ will prevent $`m_{\mathrm{eff}}`$ from vanishing as the inflaton $`y(\tau )`$ oscillates. However, as we increase $`q_1`$, larger bare masses are still allowed for particle production. The largest (negative) value of $`y(\tau )1`$ gives the cutoff mass. Recalling that the actual mass of the fermion is $`m_\psi =(g/\sqrt{\lambda })\overline{m}_\psi M`$, and that we are working on the regime $`M<gM_\mathrm{P}`$, we see that this still allows for quite large fermion masses to be produced.
### 3.2 Coupling to the Higgs
We will consider now the coupling of fermions to the Higgs, and study their pair production. The associated Mathieu equation is analogous to (7), but with $`x(\tau )=+\sqrt{1y(\tau )^2}`$ in the place of $`y(\tau )`$, i.e.
$$X_k^{\prime \prime }+\mathrm{\Omega }_k^2X_ki\sqrt{q_2}x^{}X_k=0,$$
(19)
where $`\mathrm{\Omega }_k^2=k^2+q_2x^2`$, $`x^{}=y^{}y/x`$, and $`q_2=h_2^2/g^2`$.
The initial conditions are given by (11) with $`y_0x_0`$ and we take the bare mass $`\overline{m}_\psi =0`$, since we are assuming that the fermion acquires a mass through the Higgs mechanism.
The occupation number is given by
$$n_k(\tau )=\frac{1}{2}\frac{k^2}{\mathrm{\Omega }_k}\mathrm{Im}[X_kX_{k}^{}{}_{}{}^{}]\frac{\sqrt{q_2}x}{2\mathrm{\Omega }_k}.$$
(20)
Particle production again occurs periodically. In this case, the jumps in the occupation number take place for every period and half period of the inflaton, i.e. at $`\tau =nT/2`$. These correspond to the times when the Higgs approaches zero, and thus when the effective mass of the fermion is at a minimum. An approximate expression $`\overline{n}_k(\tau )`$ can be constructed as in the previous case of fermions coupled to the inflaton. However, a function that matches $`n_k(\tau )`$ at the times when it has the spikes would not generally be a good approximation to the $`n_k(\tau )`$. Thus we have performed a shift of a quarter of period in the initial time, so that $`\overline{n}_k(\tau )`$ matches the value of $`n_k(\tau )`$ at $`\tau =(2n+1)T/4`$, i.e. where $`n_k(\tau )`$ has the plateau, and this gives a very good approximation as it is shown in Fig. 7. Starting with $`n_k(\tau )=0`$ at $`\tau =T/4`$ does not affect the results after a few inflaton oscillations, since the density rapidly saturates to its asymptotic value.
The expressions for $`F_k`$ and $`\nu _k`$ are given by Eqs. (16) and (17), respectively. In this case the function with smaller frequency in the peaks of $`F_k`$ corresponds to the solution of Eq. (17) with the plus sign choice. We show in Fig. 8 the instability chart for fermions coupled to the Higgs. The resonance bands are narrower than those corresponding to the coupling to the inflaton.
The upper panel in Fig. 9 shows $`F_k`$ as a function of $`k^2`$ for several values of $`q_2`$, between 1 and $`10^4`$, while the lower panel shows the corresponding $`\nu _k`$. Note how quickly the bands become very narrow as $`k^2`$ is increased. As a consequence, the fermion production through the coupling to the Higgs is in general less efficient than through the coupling to the inflaton in hybrid inflation models.
The ratio of Higgs-coupled fermion energy density to total energy is given by
$$\frac{\rho __F(\tau )}{\rho _{\mathrm{total}}}=\frac{g^2h_2^2}{\lambda }\frac{2N}{\pi ^2q_2}𝑑kk^2\mathrm{\Omega }_k(\tau )\overline{n}_k(\tau ).$$
(21)
and it is shown in Fig. 10. If we take as a crude fit of the numerical results that $`(\lambda /g^2h_2^2)\rho __F/\rho _{\mathrm{total}}10^2/q_2`$, we see that $`\rho __F/\rho _{\mathrm{total}}10^2g^4/\lambda `$, and since we are working in the regime $`\lambda g^2`$, we see that the fraction of energy transferred to fermions here remains small (but not necessarily negligible).
Since the fermionic production takes place when $`x0`$, it is necessary to take the initial condition $`y_0`$ sufficiently close to one so that small Higgs values occur. Taking, as we did, $`y_0=0.9999`$, ensures that at least $`x10^2`$ is reached, and hence the parametric resonance is not inhibited by a mass gap for the values of $`q_2`$ considered.
Notice that when the symmetry is finally broken, the fermions get a mass $`m_\psi =h_2M/\sqrt{\lambda }`$, which also in this case can be quite large.
## 4 Boson production
The parametric resonant production of Higgs field particles has been shown to be quite inefficient in hybrid inflation models . However, other scalar particles coupled to the inflaton or the Higgs can can be abundantly produced. Here we will compute the parametric production of these bosons in order to compare it with the fermionic production.
We consider a scalar field $`\chi `$ with mass $`m_\chi `$ and coupling to $`\varphi `$ and $`\sigma `$ given by
$$V(\chi )=\frac{1}{2}m_\chi ^2\chi ^2+\frac{1}{2}g_1^2\chi ^2\sigma ^2+\frac{1}{2}g_2^2\chi ^2\varphi ^2.$$
(22)
We can write the equation of motion for the scalar field modes $`X_k`$ as
$$X_k^{\prime \prime }+\mathrm{\Omega }_k^2(\tau )X_k(\tau )=0,$$
(23)
where
$`\mathrm{\Omega }_k^2(\tau )`$ $`=`$ $`k^2+\overline{m}_\chi ^2+q_1^{}y^2(\tau )+q_2^{}(1y^2(\tau )),`$ (24)
$`q_1^{}`$ $`=`$ $`{\displaystyle \frac{\lambda g_1^2}{g^4}},q_2^{}={\displaystyle \frac{g_2^2}{g^2}},\overline{m}_\chi ={\displaystyle \frac{m_\chi }{\overline{M}}}.`$ (25)
We chose initial conditions of positive frequency plane waves at $`\tau <0`$
$`X_k(0)`$ $`=`$ $`[2\mathrm{\Omega }_k]^{1/2},`$ (26)
$`X_k^{}(0)`$ $`=`$ $`i\mathrm{\Omega }_kX_k(0).`$ (27)
The boson occupation number can then be calculated as
$$n_k(\tau )=\frac{1}{2\mathrm{\Omega }_k}|X_k^{}|^2+\frac{\mathrm{\Omega }_k}{2}|X_k|^2\frac{1}{2}.$$
(28)
Note that $`n_k(0)=0`$, thanks to Eqs. (26) and (27).
As in the case of the fermion production, we will use the method of Ref. to compute an approximate expression for the boson occupation number
$$\overline{n}_k(\tau )=2\mathrm{sinh}^2(\mu _k\tau ),$$
(29)
where the Floquet index or growth factor $`\mu _k`$ is determined by:
$$\mathrm{cosh}(\mu _kT)=\mathrm{Re}X_k^{(1)}(T),$$
(30)
with $`X_k^{(1)}`$ satisfying the same equation (23) with the initial condition $`X_k^{(1)}(0)=1`$, $`X_{k}^{(1)}{}_{}{}^{}(0)=0`$. The energy density of bosons can be obtained using this approximate solution as
$$\rho __B(\tau )=\frac{1}{2\pi ^2}𝑑kk^2\mathrm{\Omega }_k(\tau )\overline{n}_k(\tau ).$$
(31)
and therefore, the ratio of boson to total energies is
$$\frac{\rho __B(\tau )}{\rho _{\mathrm{total}}}=\frac{4g_{1(2)}^2}{\pi ^2q_{1(2)}}𝑑kk^2\mathrm{\Omega }_k(\tau )\mathrm{sinh}^2(\mu _k\tau ).$$
(32)
We show in Fig. 11 the instability chart for bosons coupled to the inflaton (upper panel) and to the Higgs (lower panel). Shaded areas correspond to the instability bands for boson production. The darker areas correspond to larger Floquet index $`\mu _k`$, and thus to a more efficient particle production. In the unshaded areas there is no exponential boson production. The instability bands are thinner in the case of the coupling to the Higgs compared to those corresponding to the coupling to the inflaton, and they shrink to zero for specific values $`q_2^{}=1,3,6,10,`$ etc. This reflects the fact that modes with $`q_2^{}=n(n+1)/2`$ have no instabilities .
In Fig. 12 we show the Floquet index $`\mu _k`$ for bosons coupled to the inflaton and different values of the parameters $`q_1^{}`$ up to $`10^6`$, corresponding to vertical slices in the instability charts of Fig. 11 at those $`q_1^{}`$ values. The instability bands become wider and extend up to larger $`k^2`$ with increasing $`q_1^{}`$ values. Fig. 13 shows the fraction of the total energy transferred to bosons coupled to the inflaton. Its mean value grows exponentially with time. The exponent is larger for larger $`q_1^{}`$ values and the boson production is more efficient. When this fraction approaches unity, the backreaction effect becomes important and the results will be modified . These plots have to be compared with those in Figs. 5 and 10 for fermions coupled to the inflaton and Higgs respectively. It is clear that during the first few oscillations of the inflaton, the energy transfer to fermions can be more important than that to bosons.
In the case that the bosons are coupled to the Higgs, as we noticed above, the particle production vanishes for $`q_2^{}=n(n+1)/2`$, and thus no energy is transferred to fermions for those $`q_2^{}`$ values. We show in Fig. 14 the fraction of the total energy transferred to bosons as a function of time for a set of $`q_2^{}`$ values spanning the range between the two zeroes, corresponding to $`n=3`$ and 4. The efficiency of bosonic production has an oscillatory behaviour with increasing $`q_2^{}`$: it is maximal for $`q_2^{}n^2/2`$ and minimal (zero) for $`q_2^{}=n(n+1)/2`$.
## 5 Conclusions
In this paper we have studied the parametric resonant production of fermions in hybrid inflation, with both fields, inflaton and Higgs coherently oscillating after inflation. We have assumed that fermions may couple to either the inflaton or the Higgs (or both). The behaviour in the two cases is very different. While fermion production is very important in the case of a coupling to the inflaton, even in the presence of a bare mass, the production of fermions coupled only to the Higgs is generically weak. This is related to the fact that the non-adiabaticity condition, $`d\mathrm{\Omega }_k/d\tau >\mathrm{\Omega }_k^2`$, is harder to achieve for the Higgs since, when $`\mathrm{\Omega }_k`$ is at a minimum, $`d\mathrm{\Omega }_k/d\tau `$ is also at a minimum (contrary to the inflaton case). When the bare fermion mass exceeds the value $`\overline{m}_\psi >\sqrt{q_1}`$ the fermion production by the coupling to the inflaton field is also suppressed.
We have studied the growth of the fermion energy density and seen that it very quickly saturates to an approximately constant value. For fermions coupled to the inflaton, the asymptotic value grows with the resonance parameter like $`h_1^2q_1^{1/4}`$. For natural values of the couplings, a significant fraction of the inflaton energy can be transferred to fermions. On the other hand, for fermions coupled to the Higgs, the fermion energy density is of the order $`\rho __F/\rho _{\mathrm{total}}10^2g^4/\lambda `$, which under our working assumptions ($`g^2\lambda `$) is quite small.
We have also studied the boson production for both couplings to inflaton and Higgs, and compared with the fermion production. While the energy density transferred to the parametrically produced fermions saturates after a few oscillations, the one in bosons grows exponentially with time. Hence, if both fermions and bosons have similar couplings, most of the particle production goes initially into fermions, while at late times the boson production is exponentially dominant. There are however some values of boson couplings to the Higgs, corresponding to $`q_2^{}=n(n+1)/2`$, which completely inhibit the parametric resonance of bosons.
When the energy density of the bosonic or fermionic particles produced becomes sizeable, they are expected to backreact on the inflaton, affecting its evolution and eventually suppressing the parametric production of particles. We have not considered this process in detail in this paper since the approach followed is not the most suited one for this purpose. For a proper discussion of this issue in the context of chaotic inflation, see Refs. .
As a summary, we have shown that the production of fermions in the preheating stage of hybrid inflation can be very important. For the range of model parameters assumed, hybrid inflation models lead to a more efficient fermion production than chaotic inflation models , without the need to go to extremely large values of the resonance parameter $`q`$. Depending on the relative size of the couplings, and on the backreaction process, the inflaton energy transferred to fermions may even be larger than that transferred to bosons. At any rate, the parametric production of out of equilibrium fermions could have interesting consequences for cosmological issues such as the generation of the baryon asymmetry through e.g. the leptogenesis mechanism.
Note added: while completing the writing of this work a related paper appeared where the gravitino production during preheating is computed in a supersymmetric model of hybrid inflation.
## Acknowledgements
The work of JGB has been supported in part by a CICYT project AEN/97/1678. SM and ER acknowledge financial support from CONICET, Fundación Antorchas and Agencia Nacional de Promoción Científica y Tecnológica. SM and ER thank Centro Atómico Bariloche for hospitality during the completion of this work.
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# 1 Introduction
## 1 Introduction
Many measurements of heavy quark production have been performed in $`\mathrm{e}^+\mathrm{e}^{}`$ collisions on the $`\mathrm{Z}^0`$ resonance . Among the measured parameters are the production cross-sections of bottom and charm quark pairs relative to the hadronic cross-section, $`R_\mathrm{b}`$ and $`R_\mathrm{c}`$, and the forward-backward asymmetries $`A_{\mathrm{FB}}^\mathrm{b}`$ and $`A_{\mathrm{FB}}^\mathrm{c}`$. In this paper, measurements of $`R_\mathrm{b}`$, $`A_{\mathrm{FB}}^\mathrm{b}`$, and $`A_{\mathrm{FB}}^\mathrm{c}`$ at energies above the $`\mathrm{Z}^0`$ resonance are presented for $`\mathrm{e}^+\mathrm{e}^{}\mathrm{Z}^0/\gamma ^{}\mathrm{q}\overline{\mathrm{q}}`$ events, where the effective centre-of-mass energy $`\sqrt{s^{}}`$ after initial-state radiation is required to satisfy $`\sqrt{s^{}/s}>0.85`$. Similar measurements have been performed previously by other collaborations .
The data collected with the OPAL detector at LEP at centre-of-mass energies between $`130`$ GeV and $`189`$ GeV are analysed. The basic techniques are similar to those adopted in previous OPAL measurements . The $`R_\mathrm{b}`$ measurement is based on the selection of a sample enriched in $`\mathrm{b}\overline{\mathrm{b}}`$ events obtained with a secondary vertex tagging technique. Both $`r\varphi `$ and $`rz`$ information<sup>1</sup><sup>1</sup>1The OPAL coordinate system is defined as a right-handed Cartesian coordinate system, with the $`x`$ axis pointing in the plane of the LEP collider towards the centre of the ring, the $`z`$ axis in the direction of the outgoing electrons, and $`\theta `$ and $`\varphi `$ defined as the usual spherical polar coordinates. from the silicon microvertex detector are used in the selection of $`\mathrm{b}\overline{\mathrm{b}}`$ events, improving on where only $`r\varphi `$ information was used. This provides higher efficiency at comparable purity, thus enhancing the statistical precision of the measurement.
Two measurements of forward-backward asymmetries are performed. In the events selected by secondary vertex tagging, the asymmetry $`A_{\mathrm{FB}}^\mathrm{b}`$ is measured using a hemisphere charge technique to identify the direction of emission of the primary quark. An independent measurement of both the bottom and charm forward-backward asymmetries is performed using leptons from semileptonic decays of heavy hadrons and pions from $`\mathrm{D}^+\mathrm{D}^0\pi ^+`$ decays. The asymmetries obtained with the hemisphere charge method and the lepton and slow pion methods are combined.
In Section 2, a brief description of the OPAL detector and the event selection is given. The $`R_\mathrm{b}`$ measurement is discussed in Section 3, followed by a description of the asymmetry measurements in Section 4. Systematic errors on both the $`R_\mathrm{b}`$ and asymmetry measurements are given in Section 5, and in Section 6, all results are summarised.
## 2 The OPAL Detector and Event Selection
A detailed description of the OPAL detector can be found elsewhere . For this analysis, the most relevant parts of the detector are the silicon microvertex detector, the tracking chambers, the electromagnetic and hadronic calorimeters, and the muon chambers. The microvertex detector is essential for the reconstruction of secondary vertices. The central detector provides precise measurements of the momenta of charged particles by the curvature of their trajectories in a magnetic field of $`0.435`$T. In addition, it allows an identification of charged particles through a combination of the measurement of the specific energy loss $`\mathrm{d}E/\mathrm{d}x`$ and the momentum. The electromagnetic calorimeter consists of approximately 12000 lead glass blocks, which completely cover the azimuthal range up to polar angles of $`|\mathrm{cos}\theta |<0.98`$. Nearly the entire detector is surrounded with four layers of muon chambers, after approximately one metre of iron from the magnet return yoke, which is instrumented as a hadron calorimeter.
Starting in 1995, the LEP experiments have collected data at increasing energies well above the $`\mathrm{Z}^0`$ peak. In this paper, the energy points are classified in five different sets, at centre-of-mass energies which will be generically called $`\sqrt{s}`$=133, 161, 172, 183, and 189 GeV. Table 1 shows the luminosity-weighted mean centre-of-mass energies at which data were taken, and the corresponding integrated luminosities. Additionally, calibration data taken at the $`\mathrm{Z}^0`$ peak during 1996, 1997, and 1998 are used to cross-check the analyses.
Hadronic events, $`\mathrm{e}^+\mathrm{e}^{}\mathrm{q}\overline{\mathrm{q}}`$, are selected based on the number of reconstructed charged tracks and the energy deposited in the calorimeters. The selection of the subsample of non-radiative hadronic events, defined by the requirement $`\sqrt{s^{}/s}>0.85`$, and the identification and rejection of $`\mathrm{W}^+\mathrm{W}^{}`$ background are described in detail in . The remaining contamination from radiative hadronic events with true effective centre-of-mass energy below $`0.85\sqrt{s}`$ is $`510\%`$, depending on the centre-of-mass energy. The residual contamination from four-fermion events (mainly W and Z pairs) is largest at $`\sqrt{s}=189\mathrm{GeV}`$, where it is about 8$`\%`$. These backgrounds to non-radiative hadronic $`\mathrm{q}\overline{\mathrm{q}}`$ events are accounted for in the measurements.
Jets are reconstructed using the JADE algorithm with the E0 recombination scheme , keeping the invariant mass cut-off $`x_{\mathrm{min}}=49\mathrm{GeV}^2`$ fixed at all centre-of-mass energies.
Hadronic events are simulated using the PYTHIA Monte Carlo generator . Heavy quark fragmentation is modelled according to the scheme by Peterson et al. with fragmentation parameters tuned according to the results in . Four-fermion background events are simulated with the grc4f generator . The events are passed through a detailed simulation of the OPAL detector before being analysed using the same procedure as for the data.
## 3 Measurement of $`𝑹_𝐛`$
The tagging of $`\mathrm{b}\overline{\mathrm{b}}`$ events is based on the long lifetime ($``$1.5 ps) and hard fragmentation of b-flavoured hadrons, which give rise to secondary vertices significantly displaced from the primary vertex. The secondary vertex tag described in Section 3.1 allows a clean and efficient reconstruction of $`\mathrm{b}\overline{\mathrm{b}}`$ events. In Section 3.2, the measurement of $`R_\mathrm{b}`$ with vertex tagged events is described.
### 3.1 Secondary Vertex Tag
The algorithm used for secondary vertex reconstruction is described in . For the results presented here, a three-dimensional vertex tagging algorithm is used, which takes advantage of the precise $`z`$ information provided by the OPAL microvertex detector.
The primary vertex in each event is reconstructed as described in , incorporating the average beam spot position determined from the measured tracks and the LEP beam-orbit measurements as a constraint. Although the beam spot is less precisely determined at energies above the $`\mathrm{Z}^0`$ resonance than at the $`\mathrm{Z}^0`$ peak, the resulting error on the primary vertex position is still small compared to the error on the reconstructed secondary vertex position. The angular acceptance is restricted to $`|\mathrm{cos}\theta _T|<0.9`$, where $`\theta _T`$ denotes the polar angle of the thrust axis of the event. Charged tracks used to reconstruct secondary vertices are selected as described in and least three of them are required to form a vertex.
Each hadronic event is divided into two hemispheres by the plane perpendicular to the thrust axis and containing the nominal interaction point. For each reconstructed secondary vertex the signed decay length $`L`$ is defined as the distance between the secondary and the primary vertex. $`L`$ is taken to be positive if the secondary vertex is in the hemisphere pointed at by the momentum vector of the jet which contains the vertex, and negative otherwise. The decay length significance $`L/\sigma _L`$ is defined as the ratio of the decay length and its error. Secondary vertices with $`L/\sigma _L>8`$ are used to tag $`\mathrm{b}\overline{\mathrm{b}}`$ events. This cut represents the best compromise between tagging efficiency and purity.
### 3.2 Determination of $`𝑹_𝐛`$ from Vertex-Tagged Events
For the measurement of $`R_\mathrm{b}`$, the number of events tagged by a secondary vertex is determined and corrected for tagging efficiency and background. Due to the limited statistics compared with the data collected at the $`\mathrm{Z}^0`$ peak, a double tag technique as e.g. in cannot be applied. In order to reduce the sensitivity of the analysis to the detector resolution, a folded tag technique is used: a hemisphere is assigned a tag if it contains a secondary vertex with a decay length significance $`L/\sigma _L>8`$, or an anti-tag<sup>2</sup><sup>2</sup>2In the literature, tagged and anti-tagged events are sometimes referred to as “forward” and “backward” tagged events. This convention is not used here to avoid confusion with the distinction between forward and backward event hemispheres in the asymmetry analyses. if it contains a vertex with a decay length significance $`L/\sigma _L<8`$. The number of anti-tagged hemispheres is then subtracted from the number of tagged hemispheres. After subtraction of the four-fermion background, the difference between the number of tagged and anti-tagged hemispheres, $`N\overline{N}`$, in a sample of $`N_{\mathrm{had}}`$ hadronic events can be expressed as
$$N\overline{N}=2N_{\mathrm{had}}[(ϵ_\mathrm{b}\overline{ϵ}_\mathrm{b})R_\mathrm{b}+(ϵ_\mathrm{c}\overline{ϵ}_\mathrm{c})R_\mathrm{c}+(ϵ_{\mathrm{uds}}\overline{ϵ}_{\mathrm{uds}})(1R_\mathrm{b}R_\mathrm{c})],$$
where $`(ϵ_\mathrm{b}\overline{ϵ}_\mathrm{b})`$, $`(ϵ_\mathrm{c}\overline{ϵ}_\mathrm{c})`$, and $`(ϵ_{\mathrm{uds}}\overline{ϵ}_{\mathrm{uds}})`$ are the differences between the tagging and anti-tagging efficiencies for a given quark flavour. The tagging efficiencies for $`\mathrm{u}\overline{\mathrm{u}}`$, $`\mathrm{d}\overline{\mathrm{d}}`$, and $`\mathrm{s}\overline{\mathrm{s}}`$ events are averaged (uds), since they are very similar. For $`R_\mathrm{c}`$, the prediction of ZFITTER is used. The four-fermion background is determined from Standard Model production cross-sections and from selection efficiencies estimated from Monte Carlo simulation. The b-purity for the sample with the folded tag is defined as the fraction of $`\mathrm{b}\overline{\mathrm{b}}`$ event hemispheres contained in the sample $`N\overline{N}`$, and is about 75$`\%`$. The distribution of the decay length significance $`L/\sigma _L`$ is shown in Figure 1 for events at 189 GeV centre-of-mass energy, together with the expectation from the Monte Carlo simulation. The relative difference between data and Monte Carlo is below 1$`\%`$ for the folded-tag rate and $`(8\pm 4)\%`$ for the anti-tag rate, which is more sensitive to modelling of the detector resolution. For $`L/\sigma _L<10`$, the number of events predicted by Monte Carlo differs by two standard deviations from the number of events in the data. This indicates an incomplete simulation of the detector resolution, which is taken into account in the systematic errors as described in Section 5.1.6 below. The agreement between data and simulation is similar at the other centre-of-mass energies.
The selected event sample contains a $`510\%`$ contamination of radiative hadronic events, with one or more energetic photons emitted in the initial state. In these events the effective centre-of-mass energy is reduced to values below $`\sqrt{s^{}/s}=0.85`$ where the predicted value of $`R_\mathrm{b}`$ is up to 30$`\%`$ larger. In addition, the selection of non-radiative events is around $`3\%`$ less efficient for $`\mathrm{b}\overline{\mathrm{b}}`$ final states than for other flavours, because of a larger missing energy due to neutrinos in semileptonic b hadron decays. The results are corrected for these effects, which are estimated from Monte Carlo simulation. The measured values of $`R_\mathrm{b}`$ are also corrected for interference between initial and final-state radiation as described in . This correction comes from the fact that Monte Carlo, which is used to model the data, does not contain interference between initial and final-state radiation.
The numbers of selected events and of tagged and anti-tagged hemispheres are given in Table 2. The differences in hemisphere tagging efficiency, also listed in Table 2, have been determined from Monte Carlo simulation. Their errors include all the systematic uncertainties that will be described in Section 5.1. No systematic uncertainties other than those due to Monte Carlo statistics and detector resolution are assigned to the efficiencies in $`\mathrm{u}\overline{\mathrm{u}}`$, $`\mathrm{d}\overline{\mathrm{d}}`$, and $`\mathrm{s}\overline{\mathrm{s}}`$ events, as they represent a small fraction of the tagged sample. The systematic error is dominated by the uncertainties from the event selection, the modelling of b and c fragmentation and decay, and from the simulation of the detector resolution.
The dependence of the result on the assumed value of $`R_\mathrm{c}`$ can be parametrised as
$$\mathrm{\Delta }R_\mathrm{b}=b\left(R_\mathrm{c}R_\mathrm{c}^{\mathrm{SM}}\right).$$
(1)
The parameter $`b`$ has been determined separately for each centre-of-mass energy. Its values are given in Table 3.
As a cross-check, the analysis is repeated on calibration data collected at the $`\mathrm{Z}^0`$ peak. A value of $`R_\mathrm{b}(\sqrt{s}=m_{\mathrm{Z}^0})=0.221\pm 0.002(\mathrm{stat}.)\pm 0.010(\mathrm{syst}.)`$ is obtained. The systematic error has been determined as for the high-energy data samples. This result agrees within the errors with the LEP1 combined value of $`R_\mathrm{b}^0=0.21664\pm 0.00076`$ and can be regarded as a check of the evaluation of systematic errors for the measurements at energies above the $`\mathrm{Z}^0`$ peak. Note that the $`4\%`$ systematic error which is assigned to the $`R_\mathrm{b}`$ measurement at $`\sqrt{s}=189`$ GeV is larger than the $`2\%`$ difference with respect to the LEP1 average which is observed at $`\sqrt{s}=m_{\mathrm{Z}^0}`$.
## 4 Measurement of Forward-Backward Asymmetries
For the measurement of heavy quark forward-backward asymmetries, it is necessary to distinguish the event hemispheres of the primary quark and antiquark in addition to the quark flavour tagging. Two complementary techniques are used. The first analysis provides a measurement of $`A_{\mathrm{FB}}^\mathrm{b}`$ for the events that have been tagged by the presence of a secondary vertex. A hemisphere charge method is used to distinguish between quark and anti-quark hemispheres.
The second technique is used for a simultaneous measurement of $`A_{\mathrm{FB}}^\mathrm{b}`$ and $`A_{\mathrm{FB}}^\mathrm{c}`$. It is based on the identification of leptons from semileptonic decays of heavy hadrons (“prompt leptons”) and pions from $`\mathrm{D}^+\mathrm{D}^0\pi ^+`$ decays (“slow pions”). The charge of these particles provides a clean distinction between the primary quark and anti-quark hemispheres.
In Section 4.1, the measurement of $`A_{\mathrm{FB}}^\mathrm{b}`$ with the hemisphere charge technique is described. After a discussion of the lepton and slow pion identification in Section 4.2 and the flavour separation of the tagged samples (Section 4.3), the asymmetry measurement based on prompt leptons and slow pions is presented in Section 4.4. Finally, the combination of the two measurements is treated in Section 4.5.
### 4.1 Measurement of $`𝑨_{\mathrm{𝐅𝐁}}^𝐛`$ with a Hemisphere Charge Technique
To obtain a sample enriched in $`\mathrm{b}\overline{\mathrm{b}}`$ events suitable for the $`A_{\mathrm{FB}}^\mathrm{b}`$ measurement, the secondary vertex algorithm used for the $`R_\mathrm{b}`$ analysis is employed. The analysis is limited to the range $`|\mathrm{cos}\theta _T|<0.9`$. Only tagged events are used, since Monte Carlo studies show that with the current statistics, the folded tag would result in a larger error on the forward-backward asymmetry. Thus, an event is considered if it contains a secondary vertex with a decay length significance $`L/\sigma _L>8`$ and no other secondary vertex with $`L/\sigma _L<8`$. This cut is chosen to minimize the overall error.
Each non-radiative hadronic event is divided into two hemispheres by the plane perpendicular to the thrust axis that contains the nominal interaction point. Monte Carlo simulation shows that the direction of the thrust axis is a good approximation of the direction of emission of the initial $`\mathrm{q}\overline{\mathrm{q}}`$ pair. For each hemisphere, the hemisphere charge $`Q_{\mathrm{hem}}`$ is computed as
$$Q_{\mathrm{hem}}=\frac{_{i=1}^n|p_i|^\kappa Q_i}{_{i=1}^n|p_i|^\kappa },$$
where the sum runs over all $`n`$ tracks in the hemisphere, $`p_i`$ is the momentum component of track $`i`$ along the thrust axis, $`Q_i`$ denotes its charge, and $`\kappa =0.4`$ is a parameter tuned on simulated events for an optimal charge identification of a primary b or $`\overline{\mathrm{b}}`$. To ensure a good hemisphere charge reconstruction, only events with more than three charged tracks per hemisphere are used. Events are classified according to the sign of the difference $`Q_FQ_B`$ between the forward ($`Q_F`$) and backward ($`Q_B`$) hemisphere charges, where the forward hemisphere is defined as the one that contains the momentum vector of the incoming electron.
In order to ensure that each event is used at most once in this analysis and the measurement based on prompt leptons and slow pions described below, every event is assigned a figure of merit $`𝒫_{\mathrm{sig}}^{(\mathrm{vtx})}`$ defined as
$$𝒫_{\mathrm{sig}}^{(\mathrm{vtx})}=\stackrel{~}{F}_\mathrm{b}(2\stackrel{~}{P}_\mathrm{b}1),$$
(2)
where $`\stackrel{~}{F}_\mathrm{b}`$ denotes the estimated $`\mathrm{b}\overline{\mathrm{b}}`$ purity as a function of decay length significance and $`\stackrel{~}{P}_\mathrm{b}`$ stands for the estimated probability of correct charge identification<sup>3</sup><sup>3</sup>3The charge identification probability $`\stackrel{~}{P}_\mathrm{b}`$ does not depend significantly on the decay length significance of the tagged vertex. as a function of $`|Q_FQ_B|`$, both determined from the simulation. Events are rejected from the vertex-tagged sample if they are also tagged by the presence of a prompt lepton or slow pion with a corresponding figure of merit $`𝒫_{\mathrm{sig}}^{(\mathrm{}/\pi _s)}>𝒫_{\mathrm{sig}}^{(\mathrm{vtx})}`$, where $`𝒫_{\mathrm{sig}}^{(\mathrm{}/\pi _s)}`$ is determined according to Equation 8 (see Section 4.3 below). At $`\sqrt{s}=189\mathrm{GeV}`$, the events that contain both a tagged secondary vertex and a prompt lepton or slow pion correspond to $`56\%`$ of the vertex tagged and $`28\%`$ of the lepton or slow pion tagged samples, respectively. Of these common events, $`58\%`$ are assigned to the secondary vertex tagged event sample by the procedure described above. It has been checked that systematic cross-dependences between the two asymmetry measurements, which may in principle be introduced by this method, are negligibly small. Note that the quantity $`𝒫_{\mathrm{sig}}^{(\mathrm{vtx})}`$ is not used in the fit which determines $`A_{\mathrm{FB}}^\mathrm{b}`$, but only to define the selected sample of events.
For the final vertex-tagged event sample, the tagging efficiencies for each flavour are determined from Monte Carlo and are shown in Table 4. The uncertainties include the systematic errors, which will be discussed in detail in Section 5.1. They are dominated by uncertainties in the bottom and charm physics modelling and detector resolution, as in the $`R_\mathrm{b}`$ analysis. The b purity is about 60$`\%`$, with a fraction of four-fermion background up to 5$`\%`$. The lower b purity with respect to the folded tag purity does not limit the precision of the asymmetry measurement with the present statistics. Because of the rejection of some of the events that are also lepton or slow pion tagged, the final $`\mathrm{b}\overline{\mathrm{b}}`$ efficiencies and purities are lower than those which could otherwise be obtained for a vertex-tagged event sample.
The quantities $`P_\mathrm{q}`$ are the probabilities for the hemisphere charge method to correctly identify the event hemisphere into which the primary quark $`\mathrm{q}`$ was emitted. They are determined from Monte Carlo simulation and are given in Table 4 for the different quark flavours at $`\sqrt{s}=189`$ GeV. Similar values are obtained at different centre-of-mass energies. Their errors include systematic uncertainties which will be discussed in detail in Section 5.1. The largest contribution to the systematic error of the b and c quark charge identification probabilities arises from the modelling of heavy flavour fragmentation and decay. For light flavours the uncertainties on fragmentation are expected to have a small effect on the total systematic error, and only Monte Carlo statistics and detector resolution are considered. Possible detector biases in the charge identification probability for positive and negative quarks have been investigated using $`\mathrm{Z}^0`$ calibration data, and have been found to be negligible.
The quantity
$$x=\mathrm{sign}(Q_FQ_B)|\mathrm{cos}\theta _T|$$
(3)
is computed for each event, where $`\theta _T`$ denotes the polar angle of the thrust axis. Its observed distribution at 189 GeV centre-of-mass energy is compared with the Monte Carlo prediction in Figure 2. When only vector and axial vector couplings of the quarks to a gauge boson exchanged in the s-channel are allowed, the observed angular distribution of the primary quark can be expressed as
$$\frac{\mathrm{d}\sigma ^{obs}}{\mathrm{d}x}=𝒞ϵ(x)(1+x^2+\frac{8}{3}A_{\mathrm{FB}}^{obs}x),$$
(4)
where the quark masses have been neglected. The constant $`𝒞`$ is for normalization, and $`ϵ(x)`$ is the tagging efficiency as a function of $`\mathrm{cos}\theta _T`$ for an event. It is assumed that the efficiencies are symmetric functions of $`x`$, and it has been checked in the simulation that their dependence on $`x`$ is the same for all primary flavours. For other event types (e.g. four-fermion events), the predicted differential cross-section is not a second-order polynomial, but the resulting effects are negligible within the precision of the measurements presented here.
Using Equation 4, the observed asymmetry $`A_{\mathrm{FB}}^{obs}`$ is obtained by maximising the log likelihood
$$\mathrm{ln}=\underset{j=1}{\overset{N}{}}\mathrm{ln}\left\{𝒞ϵ(x_j)\right\}+\underset{j=1}{\overset{N}{}}\mathrm{ln}\left\{1+x_j^2+\frac{8}{3}A_{\mathrm{FB}}^{obs}x_j\right\},$$
(5)
where the sum is over all $`N`$ selected events, and $`A_{\mathrm{FB}}^{obs}`$ is the only free parameter in the fit. The first term is a constant for a given set of events.
After four-fermion background subtraction, $`A_{\mathrm{FB}}^\mathrm{b}`$ is determined from the relation
$$A_{\mathrm{FB}}^{obs}=\underset{\mathrm{q}=\mathrm{u},\mathrm{d},\mathrm{s},\mathrm{c},\mathrm{b}}{}s_\mathrm{q}F_\mathrm{q}(2P_\mathrm{q}1)A_{\mathrm{FB}}^\mathrm{q}.$$
Here, $`s_\mathrm{q}`$ is $`1`$ ($`+1`$) for up-type (down-type) quarks, and $`A_{\mathrm{FB}}^\mathrm{q}`$ is the forward-backward asymmetry for flavour $`\mathrm{q}`$. The asymmetries for non-b events are fixed to their Standard Model values as calculated by ZFITTER. The fractions $`F_\mathrm{q}`$ of events of flavour $`\mathrm{q}`$ in the sample are determined as
$$F_\mathrm{q}=\frac{R_\mathrm{q}ϵ_\mathrm{q}}{_jR_jϵ_j},$$
where $`R_\mathrm{q}`$ is the ratio of the cross-section of quark type $`\mathrm{q}`$ to the total hadronic cross-section, determined from ZFITTER, and $`ϵ_\mathrm{q}`$ is the tagging efficiency determined from Monte Carlo. The factor $`(2P_\mathrm{q}1)`$ is to account for charge misassignment. The small contamination from four-fermion events is evaluated as for the $`R_\mathrm{b}`$ measurement and subtracted from the sample. Its observed asymmetry is found to be consistent with zero within the available Monte Carlo statistics.
The numbers of tagged events, the observed asymmetries $`A_{\mathrm{FB}}^{obs}`$, and the corrected asymmetries $`A_{\mathrm{FB}}^\mathrm{b}`$ with their statistical and systematic errors are given in Table 5, together with the Standard Model expectations at the different centre-of-mass energies. Most of the systematic errors are in common with the $`R_\mathrm{b}`$ analysis, as discussed below in Section 5.1. The largest systematic errors arise from uncertainties in the detector resolution and the event selection procedure. Uncertainties in the fragmentation of light quarks are assumed to be negligible and are not considered. For all centre-of-mass energies, the statistical error is dominant.
While the observed asymmetries are well within the physical range of $`0.75<A_{\mathrm{FB}}^{obs}<0.75`$, values of $`A_{\mathrm{FB}}^\mathrm{b}`$ outside this range are possible because of the corrections due to sample composition and charge identification probability. All corrected $`A_{\mathrm{FB}}^\mathrm{b}`$ values are compatible with a value inside the physical range. No constraint is applied to force the corrected $`A_{\mathrm{FB}}^\mathrm{b}`$ values to lie within this range, in order to facilitate the combination with the values determined in the lepton and slow pion analysis and with measurements by other experiments.
### 4.2 Identification of Leptons and Slow Pions
Prompt leptons from semileptonic decays of heavy hadrons provide a means of tagging both $`\mathrm{b}\overline{\mathrm{b}}`$ and $`\mathrm{c}\overline{\mathrm{c}}`$ events that is largely independent of the secondary vertex tag. In addition, slow pions from $`\mathrm{D}^+\mathrm{D}^0\pi ^+`$ decays are used for tagging heavy flavour events. Both prompt leptons and slow pions allow a clean identification of the event hemisphere that contains the primary quark.
#### 4.2.1 Electron Identification
Electron candidates are required to have a momentum of at least $`2\mathrm{GeV}`$. Tracks with less than 20 d$`E`$/d$`x`$ samplings in the tracking chamber are rejected to ensure a good measurement of the specific energy loss. The difference between the measured energy loss and that expected for an electron, divided by the measurement error, is required to be between $`2`$ and $`4`$.
In this sample of preselected tracks, electrons are identified with the help of an artificial neural network, which is described in detail in . In addition to the electron preselection, a network output $`𝒩_{\mathrm{el}}>0.9`$ is required. At $`189\mathrm{GeV}`$, this selection has an efficiency for prompt electrons of approximately $`25\%`$, defined with respect to the total number of prompt electrons that are reconstructed as tracks in the detector. The resulting sample is $`75\%`$ pure in electrons.
After this selection, electrons from photon conversions are an important background in the sample. A separate artificial neural network is used to identify pairs of conversion electrons . The contribution from photon conversions is reduced by requiring a network output of $`𝒩_{\mathrm{cv}}<0.4`$. In the Monte Carlo simulation at $`189\mathrm{GeV}`$ centre-of-mass energy, $`89\%`$ of the electrons from photon conversions are rejected by this cut, while $`90\%`$ of the prompt electrons are kept. In Figure 3, the $`𝒩_{\mathrm{el}}`$ and $`𝒩_{\mathrm{cv}}`$ output distributions are shown for tracks at $`189\mathrm{GeV}`$ centre-of-mass energy.
#### 4.2.2 Muon Identification
The muon selection proceeds in two steps. First, muon track segments are formed from the hits in the muon chambers. Tracks from the central tracking chambers with a momentum greater than $`2\mathrm{GeV}`$ are extrapolated to the muon chambers. For each track segment in the muon chambers, only the “best matching track” is considered for use in the asymmetry fit. It is defined as the extrapolated track that has the smallest angular separation $`\alpha `$ to the muon track segment in question.
In a second step, an artificial neural network trained for muon identification is used to enhance the purity of the muon sample. The network uses the following eleven inputs:
* Information from the matching:
+ The square root of the $`\chi ^2`$ for the position match in $`\theta `$ and $`\varphi `$ between the extrapolated track and the associated muon track segment in the muon chambers, as described in ;
+ the ratio of distances $`R_{\mathrm{mis}}=\alpha ^{(1)}/\alpha ^{(2)}`$ of the best and second best matching track to the muon segment; this is a measure of how ambiguous the choice of the best matching track was in the preselection;
+ the $`\chi ^2`$ probability for the matching computed using both position and direction information for the track in the central detector and the associated muon track segment.
* Information from the hadron calorimeter:
+ The number of calorimeter layers in the cluster associated with the central track;
+ the number of the outermost such layer;
+ the $`\chi ^2`$ probability for the match in $`\theta `$ and $`\varphi `$ between the track (extrapolated to the hadron calorimeter) and the associated cluster.
* Specific energy loss:
+ The muon $`\mathrm{d}E/\mathrm{d}x`$ weight for the track, which is a measure of the probability that the track is compatible with a muon hypothesis;
+ $`\sigma _{\mathrm{d}E/\mathrm{d}x}`$, the error on the $`\mathrm{d}E/\mathrm{d}x`$ measurement;
+ the momentum of the track.
* Geometrical information:
+ The position in $`|\mathrm{cos}\theta |`$ and $`\varphi `$ where the extrapolated track enters the muon chambers.
The distribution of the neural network output $`𝒩_\mu `$ is shown in Figure 4 for “best matching” tracks according to the definition above. Muon candidates are retained if $`𝒩_\mu `$ is larger than $`0.65`$. In Monte Carlo simulated events at 189 GeV centre-of-mass energy, the muon selection results in an efficiency of $`43\%`$ for prompt muons, defined with respect to all prompt muons that are reconstructed as tracks, and a muon purity of $`73\%`$.
#### 4.2.3 Preselection of Slow Pion Candidates
Pions from $`\mathrm{D}^+\mathrm{D}^0\pi ^+`$ decays, denoted $`\pi _s`$ in the following, are selected based on the kinematic properties of this decay. Due to the low momentum, $`p^{}=39\mathrm{MeV}`$ , of the decay products in the $`\mathrm{D}^+`$ rest frame, pions from this decay have momenta smaller than
$$p_{\pi _s}^{\mathrm{max}}=\frac{\sqrt{s}}{2m_{\mathrm{D}^+}}\left(E^{}+p^{}\right)=0.0458\sqrt{s}$$
(6)
in the laboratory frame, where $`E^{}=\sqrt{(p^{})^2+m_{\pi ^+}^2}`$, and $`m_{\mathrm{D}^+}`$ and $`m_{\pi ^+}`$ denote the $`\mathrm{D}^+`$ and $`\pi ^+`$ masses, respectively. In addition, slow pions have a transverse momentum with respect to the $`\mathrm{D}^+`$ flight direction of at most $`p^{}`$, and are thus dominantly found in the core of the jet containing the $`\mathrm{D}^+`$ meson.
Slow pion candidates are required to have a momentum between $`1.0\mathrm{GeV}`$ and $`p_{\pi _s}^{\mathrm{max}}`$ and, if at least 20 $`\mathrm{d}E/\mathrm{d}x`$ samplings are available, a specific energy loss $`\mathrm{d}E/\mathrm{d}x`$ whose probability compatibility for the pion hypothesis exceeds $`2\%`$. Tracks that form single charged particle jets are rejected. The $`\mathrm{D}^+`$ flight direction is estimated by the jet direction, which is recalculated in an iterative procedure similar to the one described in , based on the rapidities of the tracks and clusters in the jet containing the slow pion candidate. If the jet mass exceeds $`2.3\mathrm{GeV}`$, the track or calorimeter cluster in the jet with the smallest rapidity with respect to the jet axis is removed from the calculation, and the direction is recomputed. The transverse momentum $`p_t`$ of the slow pion candidate is calculated with respect to this jet direction. The $`p_t^2`$ distribution is shown in Figure 5 for tracks at $`189\mathrm{GeV}`$ centre-of-mass energy. Slow pion candidates are accepted if $`p_t^2<0.02\mathrm{GeV}^2`$.
In Monte Carlo simulated events at $`189\mathrm{GeV}`$, this preselection is $`56\%`$ efficient and yields a sample that contains $`5.7\%`$ of slow pions. This sample is further enriched with a cut that is described in Section 4.3 below.
### 4.3 Flavour Separation of the Lepton and Slow Pion Samples
Three different sources of prompt leptons are considered: $`\mathrm{b}\mathrm{}`$, meaning leptons from semileptonic decays of b-flavoured hadrons; cascade bottom decays, which include the contributions from $`\mathrm{b}\mathrm{c}\mathrm{}`$ and $`\mathrm{b}\overline{\mathrm{c}}\mathrm{}`$ processes; and $`\mathrm{c}\mathrm{}`$, leptons from semileptonic decays of charm hadrons. The background can be classified as “non-prompt” leptons, i.e. all other leptons that are not produced in the decay of b- or c-flavoured hadrons, and particles that are mis-identified as electrons or muons.
For electrons and muons, separate artificial neural networks have been constructed with the aim of separating prompt $`\mathrm{b}\mathrm{}`$ decays, cascade $`\mathrm{b}[\mathrm{c},\overline{\mathrm{c}}]\mathrm{}`$ decays, prompt $`\mathrm{c}\mathrm{}`$ decays, and all other contributions. The technique used is similar to the one described in , but the inputs and training have been re-optimized for centre-of-mass energies above the $`\mathrm{Z}^0`$ resonance; also an artificial neural network for the identification of cascade decays is included. Candidates that form single charged particle jets are rejected since they are expected to be dominantly produced in leptonic W decays.
The first network, denoted $`𝒩_\mathrm{b}`$, is designed to separate $`\mathrm{b}\mathrm{}`$ decays from all other contributions. Two networks, $`𝒩_{\mathrm{bc}}`$ and $`𝒩_\mathrm{c}`$, have been trained on a Monte Carlo simulation that does not contain $`\mathrm{b}\mathrm{}`$ decays in order to classify lepton candidates that are background to the $`𝒩_\mathrm{b}`$ net. All three networks use the following input variables, where jet variables are defined with the lepton candidate included in the jet:
* $`p`$, the lepton track momentum;
* $`p_t`$, the transverse momentum of the lepton track calculated relative to the jet which contains the track;
* $`L/\sigma _L`$, the decay length significances of the secondary vertices (if existing) in the jet containing the lepton and the most energetic jet in the hemisphere not containing the lepton, where secondary vertices are reconstructed with the same algorithm as described in Section 3.1;
* the jet charges of the jets containing the lepton and the most energetic jet in the hemisphere not containing the lepton, each multiplied by the lepton charge, where the jet charge is defined as in Equation 4.1 with $`\kappa =0.4`$, but using only tracks associated to the jet;
* the forward multiplicity in the lepton jet, defined as the number of tracks with an impact parameter significance with respect to the primary vertex larger than $`2`$. For each track, the impact parameter is defined as the distance between the primary vertex and the track at its the point of closest approach; the impact parameter significance is defined as this distance divided by its error;
* the $`|\mathrm{cos}\theta |`$ of the jet momentum vector, where $`\theta `$ is the jet polar angle;
* the outputs $`𝒩_{\mathrm{el}}`$ and $`𝒩_{\mathrm{cv}}`$ of the electron identification network and the conversion finder network, respectively, in the case of electrons;
* the output $`𝒩_\mu `$ of the muon identification network, in the case of muons.
From the $`𝒩_\mathrm{b}`$, $`𝒩_{\mathrm{bc}}`$, and $`𝒩_\mathrm{c}`$ network outputs, the following quantities are computed which are related to the probabilities of a lepton candidate to come from one of the three sources:
$`𝒫_{\mathrm{sig}}^{(\mathrm{b}\mathrm{})}`$ $`=`$ $`𝒩_\mathrm{b}`$
$`𝒫_{\mathrm{sig}}^{(\mathrm{b}[\mathrm{c},\overline{\mathrm{c}}]\mathrm{})}`$ $`=`$ $`\left(1𝒩_\mathrm{b}\right){\displaystyle \frac{\left[𝒩_{\mathrm{bc}}\left(1𝒩_\mathrm{c}\right)\right]}{\left[𝒩_{\mathrm{bc}}\left(1𝒩_\mathrm{c}\right)\right]+\left[𝒩_\mathrm{c}\left(1𝒩_{\mathrm{bc}}\right)\right]+\left[\left(1𝒩_\mathrm{c}\right)\left(1𝒩_{\mathrm{bc}}\right)\right]}}`$ (7)
$`𝒫_{\mathrm{sig}}^{(\mathrm{c}\mathrm{})}`$ $`=`$ $`\left(1𝒩_\mathrm{b}\right){\displaystyle \frac{\left[𝒩_\mathrm{c}\left(1𝒩_{\mathrm{bc}}\right)\right]}{\left[𝒩_{\mathrm{bc}}\left(1𝒩_\mathrm{c}\right)\right]+\left[𝒩_\mathrm{c}\left(1𝒩_{\mathrm{bc}}\right)\right]+\left[\left(1𝒩_\mathrm{c}\right)\left(1𝒩_{\mathrm{bc}}\right)\right]}}.`$
The difference in the treatment of the $`𝒩_\mathrm{b}`$ output from $`𝒩_{\mathrm{bc}}`$ and $`𝒩_\mathrm{c}`$ is due to the fact that $`\mathrm{b}\mathrm{}`$ decays were omitted in the training of the latter two networks. Only candidates that satisfy the condition
$$𝒫_{\mathrm{sig}}^{(\mathrm{})}=\sqrt{\left(𝒫_{\mathrm{sig}}^{(\mathrm{b}\mathrm{})}\right)^2+\left(𝒫_{\mathrm{sig}}^{(\mathrm{b}[\mathrm{c},\overline{\mathrm{c}}]\mathrm{})}\right)^2+\left(𝒫_{\mathrm{sig}}^{(\mathrm{c}\mathrm{})}\right)^2}>0.1$$
(8)
are used in the subsequent analysis. Candidates with lower values of $`𝒫_{\mathrm{sig}}^{(\mathrm{})}`$ are expected to be dominantly background and to have a negligible contribution to the overall result. Note that the quantity $`𝒫_{\mathrm{sig}}^{(\mathrm{})}`$, although used to define the selected sample of events, is not itself used in the fit that determines the bottom and charm asymmetries.
Similarly to the lepton case, slow pion candidates are classified as cascade $`\mathrm{b}[\mathrm{c},\overline{\mathrm{c}}]\pi _s`$ decays, $`\mathrm{c}\pi _s`$, and background. Another two artificial neural networks, $`𝒩_{\mathrm{bc}}`$ and $`𝒩_\mathrm{c}`$, have been trained for the separation of the three components in the preselected slow pion sample. The following inputs are used by both networks:
* $`p`$, the slow pion track momentum;
* $`p_t^2`$, the transverse momentum squared of the slow pion track calculated relative to the jet direction obtained with the same rapidity based algorithm that is used in Section 4.2.3;
* $`E_{\pi _s\mathrm{jet}}`$, the total energy of the jet containing the slow pion;
* $`E_{\pi _s\mathrm{s}ubjet}`$, the energy of the sub-jet containing the slow pion: Each jet containing a slow pion is split into two sub-jets, where the slow pion sub-jet is seeded by the slow pion track. In an iterative procedure, any particle that forms a smaller opening angle with the slow pion sub-jet than with the remainder of the jet is then assigned to the slow pion sub-jet;
* $`L/\sigma _L`$, the decay length significance of the vertex in the jet containing the slow pion, if a vertex is found;
* the jet charge of the jet containing the slow pion, calculated with $`\kappa =0.4`$, multiplied by the slow pion charge; and
* the $`|\mathrm{cos}\theta |`$ of the jet momentum vector.
In addition, the network for identification of $`\mathrm{b}[\mathrm{c},\overline{\mathrm{c}}]\pi _s`$ decays uses:
* $`L`$, the decay length of the vertex (if existing) in the jet containing the slow pion;
* the jet charge of the jet containing the slow pion, calculated with a different parameter $`\kappa =2.0`$, multiplied by the slow pion charge;
* the forward multiplicity in the jet containing the slow pion, defined as above;
* $`(p_t)_{\mathrm{jet}}`$, the scalar sum of the transverse momenta relative to the jet axis of all tracks in the jet; and
* the maximum longitudinal momentum component of any track in the jet containing the slow pion, measured relative to the jet direction.
The quantities $`𝒫_{\mathrm{sig}}^{(\mathrm{b}[\mathrm{c},\overline{\mathrm{c}}]\pi _s)}`$ and $`𝒫_{\mathrm{sig}}^{(\mathrm{c}\pi _s)}`$ are computed in analogy to Equation 4.3 for prompt lepton candidates, while $`𝒫_{\mathrm{sig}}^{(\mathrm{b}\pi _s)}`$ is set to zero. The same requirement on $`𝒫_{\mathrm{sig}}^{(\pi _s)}`$ (computed according to Equation 8) is imposed as for prompt leptons.
If more than one lepton or slow pion candidate per event passes the selection, the one with the highest $`𝒫_{\mathrm{sig}}^{(\mathrm{}/\pi _s)}`$ value is taken. Furthermore, this $`𝒫_{\mathrm{sig}}^{(\mathrm{}/\pi _s)}`$ value is required to be larger than the value of $`𝒫_{\mathrm{sig}}^{(\mathrm{vtx})}`$ as determined for the hemisphere charge measurement (see Section 4.1) if the event is also tagged by the presence of a secondary vertex.
A breakdown of the composition of the samples of electron, muon, and slow pion tagged events at 189 GeV centre-of-mass energy together with the efficiencies is given in Table 6. In Figure 6, the output distributions of the flavour separation networks are shown for the events that pass all cuts.
### 4.4 Measurement of $`𝑨_{\mathrm{𝐅𝐁}}^𝐛`$ and $`𝑨_{\mathrm{𝐅𝐁}}^𝐜`$ with Leptons and Slow Pions
The forward-backward asymmetries for bottom and charm, $`A_{\mathrm{FB}}^\mathrm{b}`$ and $`A_{\mathrm{FB}}^\mathrm{c}`$, are extracted from the data using an unbinned maximum likelihood fit to the charge signed polar angle distribution of the thrust axis $`\stackrel{}{T}`$ in lepton and slow pion tagged events. It is assumed that $`\stackrel{}{T}`$ is the axis along which the primary quark-antiquark pair is emitted. The quantity $`y=q\mathrm{cos}\theta _T`$ is computed event by event, where $`q`$ is the charge of the lepton or slow pion, and the thrust direction $`\stackrel{}{T}`$ is defined such that $`\stackrel{}{T}\stackrel{}{p}>0`$, with $`\stackrel{}{p}`$ being the momentum of the jet containing the lepton or slow pion. The inclusive $`y`$ distributions for electron, muon, and slow pion tagged events at $`189\mathrm{GeV}`$ are presented in Figure 7.
In the fit, both the bottom and charm asymmetries are to be determined. Therefore, the fit uses for each event its probabilities to be a correctly tagged $`\mathrm{b}\overline{\mathrm{b}}`$ or $`\mathrm{c}\overline{\mathrm{c}}`$ event or background, as determined from the simulation: The events are divided into several subsamples according to their $`𝒩_\mathrm{b}`$, $`𝒩_{\mathrm{bc}}`$, $`𝒩_\mathrm{c}`$, and $`|y|`$ values. Three bins are used for each quantity, making a total of $`81`$ subsamples separately for both electron and muon tagged events and $`27`$ subsamples for the slow pion tagged events. These subsamples have different bottom and charm purities and are fitted simultaneously.
The cross-section for producing a $`\mathrm{q}\overline{\mathrm{q}}`$ pair is assumed to depend on $`y`$ according to
$$\frac{\mathrm{d}\sigma _{\mathrm{q}\overline{\mathrm{q}}}}{\mathrm{d}y}1+y^2+\frac{8}{3}A_{\mathrm{FB}}^\mathrm{q}y.$$
(9)
The likelihood $`_{\mathrm{sub}}`$ for one subsample is given by
$$\mathrm{ln}_{\mathrm{sub}}=\underset{\mathrm{candidates}}{}\mathrm{ln}\left(1+y^2+\frac{8}{3}A_{\mathrm{FB}}^{obs}y\right),$$
(10)
where the sum is taken over all candidates in the subsample. The total likelihood is then given by the product of the likelihoods of all subsamples. The expected observed asymmetry $`A_{\mathrm{FB}}^{obs}`$ in each subsample is computed as
$$A_{\mathrm{FB}}^{obs}(𝒩_\mathrm{b},𝒩_{\mathrm{bc}},𝒩_\mathrm{c},|y|)=\underset{i=1}{\overset{5}{}}f_i(𝒩_\mathrm{b},𝒩_{\mathrm{bc}},𝒩_\mathrm{c},|y|)A_{\mathrm{FB}}^i.$$
(11)
In this equation $`f_i`$ denotes the predicted fraction of leptons or slow pions from source $`i`$, and $`A_{\mathrm{FB}}^i`$ is the corresponding asymmetry:
$$\{\begin{array}{cccc}\hfill A_{\mathrm{FB}}^1=& \hfill (12\overline{\chi }_1^{\mathrm{eff}})& A_{\mathrm{FB}}^\mathrm{b}& \mathrm{for}\mathrm{b}\mathrm{},\hfill \\ \hfill A_{\mathrm{FB}}^2=& \hfill (12\overline{\chi }_2^{\mathrm{eff}})& A_{\mathrm{FB}}^\mathrm{b}& \mathrm{for}\mathrm{b}\mathrm{c}\mathrm{}/\pi _s,\hfill \\ \hfill A_{\mathrm{FB}}^3=& \hfill (12\overline{\chi }_3^{\mathrm{eff}})& A_{\mathrm{FB}}^\mathrm{b}& \mathrm{for}\mathrm{b}\overline{\mathrm{c}}\mathrm{}/\pi _s,\hfill \\ \hfill A_{\mathrm{FB}}^4=& \hfill & A_{\mathrm{FB}}^\mathrm{c}& \mathrm{for}\mathrm{c}\mathrm{}/\pi _s,\mathrm{and}\hfill \\ \hfill A_{\mathrm{FB}}^5=& & 0& \mathrm{for}\mathrm{background},\hfill \end{array}$$
(12)
where $`\overline{\chi }_i^{\mathrm{eff}}`$ is the effective $`\mathrm{B}\overline{\mathrm{B}}`$ mixing parameter. The fractions $`f_i`$ have been calculated from the Monte Carlo simulation and depend on the mis-identification probability in each bin of $`𝒩_\mathrm{b}`$, $`𝒩_{\mathrm{bc}}`$, $`𝒩_\mathrm{c}`$, and $`|y|`$, on the production rates of bottom and charm quarks, on the semileptonic branching ratios of heavy hadrons, and on the hadronisation fractions $`f(\mathrm{b}[\mathrm{c},\overline{\mathrm{c}}]\mathrm{D}^\pm )`$ and $`f(\mathrm{c}\mathrm{D}^+)`$. Variations in the sample composition with $`|\mathrm{cos}\theta _T|`$ are taken into account since the $`f_i`$ are binned in $`|y|`$. As described in Section 4.1, this likelihood fit has the advantage that the $`|\mathrm{cos}\theta _T|`$ dependence of the efficiency for identifying leptons and slow pions is not needed explicitly. The effective mixing parameters $`\overline{\chi }_1^{\mathrm{eff}}`$, $`\overline{\chi }_2^{\mathrm{eff}}`$, and $`\overline{\chi }_3^{\mathrm{eff}}`$ are determined from the simulation for each selected subsample of events. The mixing parameter $`\overline{\chi }=0.118\pm 0.006`$ used in the simulation for inclusive $`\mathrm{b}\overline{\mathrm{b}}`$ events is taken as an external input. The small non-zero contributions to the observed asymmetry from prompt leptons and slow pions from radiative and four-fermion events are accounted for, but left out of the above list for simplicity. The assumption that the backgrounds from mis-identified leptons and slow pions do not contribute to the observed asymmetry has been checked and will be discussed in Section 5.2.2.
The fit is done separately for the data taken at $`133\mathrm{GeV}`$, 161 $`\mathrm{GeV}`$, 172 $`\mathrm{GeV}`$, 183 $`\mathrm{GeV}`$, and 189 $`\mathrm{GeV}`$. In Table 7, the numbers of lepton and slow pion tagged events are given, and the results of the fit for the bottom and charm asymmetries are summarised for each energy point together with the errors and correlations. The systematic errors have been evaluated as described in Section 5.2.
As a cross-check, the fit is also performed on the calibration data taken at the $`\mathrm{Z}^0`$ peak in the years 1996, 1997, and 1998. The results of this fit are consistent with the average of LEP1 and SLD measurements given in .
### 4.5 Combination of the $`𝑨_{\mathrm{𝐅𝐁}}^𝐛`$ Measurements
By construction, the samples used in the hemisphere charge analysis do not have any events in common with the lepton or slow pion tagged samples. The bottom asymmetry as measured with the vertex tag depends on the assumed value of the charm asymmetry, while the fit to the lepton and slow pion tagged events yields results for the bottom and charm asymmetries with a non-zero correlation. This dependence on the charm asymmetry has to be taken into account when combining the two bottom asymmetry measurements. The vertex tag measurement is applied as a constraint in the fit to the lepton and slow pion tagged events that is described in Section 4.4. Since the bottom and charm asymmetries are correlated in the fit to the lepton and slow pion tagged events, both the fitted bottom and charm asymmetries are expected to change in the constrained fit; but the error on the charm asymmetry is essentially unchanged.
The results for $`A_{\mathrm{FB}}^\mathrm{b}`$ and $`A_{\mathrm{FB}}^\mathrm{c}`$ are listed in Table 8 and are also shown in Figure 8 together with the Standard Model expectations as a function of the centre-of-mass energy. Good agreement is observed between the measurements and the predictions from ZFITTER.
In all asymmetry measurements, $`R_\mathrm{c}`$ and $`R_\mathrm{b}`$ are fixed to the Standard Model values as predicted by ZFITTER. The dependence of the measured values on the assumed values of $`R_\mathrm{b}`$ and $`R_\mathrm{c}`$ is parametrised as
$$\begin{array}{cc}\hfill \mathrm{\Delta }A_{\mathrm{FB}}^\mathrm{b}=a^\mathrm{b}(R_\mathrm{b})\frac{\mathrm{\Delta }R_\mathrm{b}}{R_\mathrm{b}^{\mathrm{SM}}}+a^\mathrm{b}(R_\mathrm{c})\frac{\mathrm{\Delta }R_\mathrm{c}}{R_\mathrm{c}^{\mathrm{SM}}}& \text{}\hfill \\ \hfill \mathrm{\Delta }A_{\mathrm{FB}}^\mathrm{c}=a^\mathrm{c}(R_\mathrm{b})\frac{\mathrm{\Delta }R_\mathrm{b}}{R_\mathrm{b}^{\mathrm{SM}}}+a^\mathrm{c}(R_\mathrm{c})\frac{\mathrm{\Delta }R_\mathrm{c}}{R_\mathrm{c}^{\mathrm{SM}}}& ,\text{}\hfill \end{array}$$
(13)
where $`\mathrm{\Delta }R_\mathrm{q}=R_\mathrm{q}R_\mathrm{q}^{\mathrm{SM}}`$ for q=b, c. The Standard Model values for $`R_\mathrm{b}`$ and $`R_\mathrm{c}`$ are given in Table 9 together with the values of the coefficients $`a^\mathrm{q}(R_\mathrm{q}^{})`$.
Any systematic error from a common source (see Section 5 for the description of systematic uncertainties) is treated as fully correlated between the two measurements. However, there are large contributions to the systematic error that affect only one of the analyses, such that systematic errors are clearly no limitation to the combination of the two results.
## 5 Systematic Errors
In this section, the evaluation of the systematic errors for the analyses presented here is discussed. Both the secondary vertex and the lepton and slow pion tags depend crucially on the knowledge of the detector resolution for reconstructing charged particle tracks. Most of the other main systematic errors are independent for the two measurements.
### 5.1 Systematic Errors on $`𝑹_𝐛`$ and $`𝑨_{\mathrm{𝐅𝐁}}^𝐛`$ with the Secondary Vertex Tag
In Table 10, a breakdown of the systematic error is given for both the $`R_\mathrm{b}`$ and $`A_{\mathrm{FB}}^\mathrm{b}`$ measurements with the vertex tag at 189 GeV centre-of-mass energy. The systematic errors considered for these measurements are described in the following paragraphs.
#### 5.1.1 Event Selection
The bias on the measurement of $`R_\mathrm{b}`$ from the selection of non-radiative events has already been discussed in Section 3.2. A correction has been applied, and an uncertainty of 100$`\%`$ is assigned to the correction. The uncertainty on the corrections for the interference between initial and final state radiation results in a systematic error of 0.5% on $`R_\mathrm{b}`$. The contamination from radiative events has the effect of decreasing the measured asymmetry by $`3\%`$. The value is corrected accordingly, and a $`3\%`$ systematic error is assigned. The likelihood fit is based on the assumption that the shape of the tagging efficiency as a function of $`|\mathrm{cos}\theta _T|`$ is the same for all flavours, which might not be true at the edges of the acceptance. The effect has been estimated by dividing the sample in bins of $`|\mathrm{cos}\theta _T|`$, over which the above assumption is valid, determing the asymmetry in each bin independently and comparing their average with the reference result. It has been found to affect the asymmetry by less than 1$`\%`$, and is thus neglected.
#### 5.1.2 Final State QCD Corrections
Final state QCD corrections are included in the calculation of the predictions with which the asymmetry measurements are compared. However, the experimental event selection is less efficient for events with hard gluon radiation due to the cuts on the decay length significance (for the vertex tag) and the lepton or slow pion momentum (for the lepton and slow pion analysis). The fraction of the QCD correction that has to be applied to the measurements has been determined previously at LEP1 and is typically between $`0.3`$ and $`0.6`$ . The overall QCD corrections at $`\sqrt{s}=189\mathrm{GeV}`$ have been determined from ZFITTER to be 0.015 for $`A_{\mathrm{FB}}^\mathrm{b}`$ and 0.022 for $`A_{\mathrm{FB}}^\mathrm{c}`$, respectively. Half of this correction is assigned as a systematic error.
#### 5.1.3 Bottom and Charm Physics Modelling
Uncertainties in bottom and charm fragmentation and decay properties have been treated as follows:
* b fragmentation: Although the mean scaled energy $`2E_\mathrm{b}/\sqrt{s}`$ of weakly decaying b hadrons is expected to change from LEP1 to LEP2 energies, the Peterson fragmentation parameter $`ϵ_P^\mathrm{b}`$, which describes one step in the fragmentation process, is assumed not to vary with energy. Simulated $`\mathrm{b}\overline{\mathrm{b}}`$ events are reweighted within the range of $`0.0030<ϵ_P^\mathrm{b}<0.0048`$ , which corresponds to a variation of the mean scaled energy $`2E_\mathrm{b}/\sqrt{s}`$ of weakly decaying b hadrons in $`\mathrm{Z}^0`$ decays in the range of $`2E_\mathrm{b}/\sqrt{s}=0.702\pm 0.008`$ .
* b decay multiplicity: The charged decay multiplicity of hadrons containing a b quark is varied in the Monte Carlo simulation according to .
* b hadron composition: The tagging efficiency differs for the various b hadron species. The fractions of b hadrons and their errors have been taken from . The fractions $`f(\mathrm{B}^0+\mathrm{B}^+)`$ and $`f(\mathrm{B}_\mathrm{s}^0)`$ are varied independently within their errors, and their variation is compensated by the b baryon fraction.
* b hadron lifetimes: The lifetimes of the different b hadrons are varied in the Monte Carlo by their errors according to .
* c fragmentation: $`\mathrm{c}\overline{\mathrm{c}}`$ Monte Carlo events are reweighted by varying the Peterson fragmentation parameter $`ϵ_P^\mathrm{c}`$ in the range of $`0.025<ϵ_P^\mathrm{c}<0.031`$ . This corresponds to a variation of the mean scaled energy $`2E_\mathrm{c}/\sqrt{s}`$ of weakly decaying c hadrons in $`\mathrm{Z}^0`$ decays in the range of $`2E_\mathrm{c}/\sqrt{s}=0.484\pm 0.008`$.
* c decay multiplicity: The average charged track multiplicities of $`\mathrm{D}^+`$, $`\mathrm{D}^0`$ and $`\mathrm{D}_\mathrm{s}^+`$ decays are varied in the Monte Carlo within the ranges of the experimental measurements .
* c hadron composition: The $`\mathrm{D}^0`$ fraction is written as $`f(\mathrm{D}^0)=1\mathrm{f}(\mathrm{D}^+)\mathrm{f}(\mathrm{D}_\mathrm{s}^+)\mathrm{f}(\mathrm{c}_{\mathrm{baryon}})`$. The last three parameters are varied independently by their errors according to to evaluate the uncertainty on the charm efficiency.
* c hadron lifetimes: Charmed hadron lifetimes are varied within their experimental errors according to .
#### 5.1.4 Four-fermion Background
Above the WW production threshold, the four-fermion background is largely dominated by W pairs. The uncertainty in the W pair production cross-section is taken into account and has been found to have a negligible systematic effect. The background from W pairs has the highest probability to be accepted in the tagged sample when one or both W bosons decay into a final state containing a charm quark. The systematic error on the W pair tagging efficiency is estimated by varying the charm physics modelling parameters as described above. The effect of the detector resolution is also taken into account, as described in Section 5.1.6.
#### 5.1.5 Monte Carlo Statistics
Tagging efficiencies and charge identification probabilities are varied by the statistical error arising from the finite number of Monte Carlo simulated events.
#### 5.1.6 Track Reconstruction
The effect of the detector resolution on the track parameters is estimated by degrading or improving the resolution of all tracks in the Monte Carlo simulation. This is done by applying a single multiplicative scale factor to the difference between the reconstructed and true track parameters. A $`\pm 10\%`$ variation is applied independently to the $`r\varphi `$ and $`rz`$ track parameters.
In addition, the matching efficiency for assigning measurement points in the silicon microvertex detector to the tracks is varied by $`1\%`$ in $`r\varphi `$ and $`3\%`$ in $`rz`$. The systematic errors resulting from the individual variations are summed in quadrature.
### 5.2 Systematic Errors on $`𝑨_{\mathrm{𝐅𝐁}}^𝐛`$ and $`𝑨_{\mathrm{𝐅𝐁}}^𝐜`$ with the Lepton and Slow Pion Tag
Three different groups of systematic errors have been considered: those from detector effects, those related to physics models and external inputs used in the analysis, and effects introduced by the limited Monte Carlo statistics and the fitting procedure. A list of all systematic errors is given in Table 11.
#### 5.2.1 Detector Systematics
The lepton and slow pion identification relies on the proper modelling of the detector response in the Monte Carlo simulation.
* Track Reconstruction: The influence of the simulated detector resolution and matching efficiency for hits in the microvertex detector is estimated as described in Section 5.1.6.
* Lepton Identification: The fractions of misidentified electrons and muons are taken from Monte Carlo simulation. The modelling of the input variables to the artificial neural network for electron identification has been studied on $`\mathrm{Z}^0`$ calibration data in a manner similar to that described in . Differences between the data and the modelling in the Monte Carlo simulation have been determined using a pure sample of electrons from photon conversions and an inclusive sample of tracks depleted in conversion electrons. The dependence of the efficiency and background contamination of the selected electron sample on these differences has been studied, and the resulting uncertainties for each input have been added in quadrature. In addition, it has been shown that there is no large dependence of the resulting systematic error on the position of the $`𝒩_{\mathrm{el}}`$ cut. The systematic error has been evaluated separately for each year of data taking at energies above the $`\mathrm{Z}^0`$ peak; the uncertainty on the efficiency is around $`10\%`$, and that on the background contamination around $`20\%`$.
The performance of the muon tagging network has been studied on $`\gamma \gamma \mu ^+\mu ^{}`$ events recorded at LEP energies above the $`\mathrm{Z}^0`$ peak and on an inclusive sample of tracks from the $`\mathrm{Z}^0`$ calibration data that fail the “best match” preselection criterion. The muon and background rates have been compared between data and Monte Carlo simulation in bins of $`𝒩_\mu `$, and the simulation has been reweighted to match the data distribution. The reweighting factors in each bin have been varied independently by their statistical errors, and the resulting variations of the measured asymmetries added in quadrature. For the inclusive muon sample selected with a cut at $`𝒩_\mu >0.65`$, the reweighting factors for muon signal and background are $`0.95\pm 0.05`$ and $`0.97\pm 0.04`$, respectively.
Conversion candidates are explicitly removed from the sample of events used in the fit. The $`𝒩_{\mathrm{cv}}`$ output distribution is well described in the simulation. The conversion finding efficiency is tested in the data using an algorithm which does not need the $`𝒩_{\mathrm{el}}`$ outputs for the two tracks as input. The accuracy of this test is $`18\%`$, and the conversion rate is varied by this amount to calculate the systematic error. The rate of non-prompt muons from pion and kaon decays in flight has been studied previously in , and has been found to be modelled to within $`9\%`$.
The four above systematic errors are added in quadrature to yield the value given in Table 11.
* Slow Pion Efficiency: Because of the large backgrounds in the slow pion sample, it is crucial to check that the slow pion reconstruction is modelled correctly in the Monte Carlo simulation. The efficiency has been studied by reconstructing slow pions in the $`\mathrm{Z}^0`$ calibration data using the same cuts as described in Section 4.2.3, except for the upper momentum cut which has been lowered to $`p_{\pi _s}^{\mathrm{max}}(\sqrt{s}=m_{\mathrm{Z}^0})=4.17\mathrm{GeV}`$. The slow pion content in this sample has been measured using a fit to the $`p_t^2`$ spectrum. Functions which are found to describe the shape of the background distribution well in the simulation are fitted to the $`p_t^2`$ sideband from $`0.03`$ to $`0.10\mathrm{GeV}^2`$. The number of signal slow pion events is then determined by extrapolating the background estimate to $`p_t^2=0\mathrm{GeV}^2`$. In the $`\mathrm{Z}^0`$ calibration data, this fit has a relative statistical precision of $`6\%`$, but an additional systematic error of $`23\%`$ is assigned to cover biases resulting from the extrapolation procedure and the particular choice of fit function.
Combining this result with the OPAL measured yield of $`\mathrm{D}^+`$ mesons in hadronic $`\mathrm{Z}^0`$ decays of $`\overline{n}_{\mathrm{Z}^0\mathrm{D}^+\mathrm{X}}=0.1854\pm 0.0041\pm 0.0059\pm 0.0069`$ , the efficiency of the slow pion reconstruction in $`\mathrm{Z}^0`$ decays is measured to be $`(37.3\pm 9.0)\%`$, where the relative errors of $`6\%`$ and $`23\%`$ have been added in quadrature. This is consistent with the value in the simulation of $`(31.9\pm 1.5(\mathrm{stat}.))\%`$. The fractional error on the slow pion efficiency is then assigned to the slow pion efficiencies at energies above the $`\mathrm{Z}^0`$ peak, which are taken from the simulation. In order to determine the systematic error on the asymmetry measurements, the expected number of slow pions in each of the different subsamples is varied according to the error on the efficiency, and balanced by the expected number of background events so that the total number of tagged events is kept constant.
No additional systematic error is assigned for the modelling of the shape of the slow pion signal $`p_t^2`$ distribution since the signal shape in $`\mathrm{Z}^0`$ decays has been found to be well described in the simulation.
* Modelling of Artificial Neural Network Inputs: Each of the input distributions used in the flavour separation networks has been compared between data and Monte Carlo simulation. The simulated distributions are reweighted for each input variable in turn to agree with the corresponding data distributions, and the analysis is repeated with the weighted events. The observed differences from the original fit result are added in quadrature to yield the systematic uncertainty due to the modelling of the input variables.
#### 5.2.2 Physics Systematics
The fragmentation of charm and bottom quarks and the momentum spectra of leptons emitted in their semileptonic decay are described using phenomenological models tuned to experimental data. Systematic errors introduced by models of heavy hadron semileptonic decays, by the use of externally measured inputs, and from the assumed asymmetries of background candidates are assessed as follows:
* Lifetimes: The lifetimes of weakly decaying b- and c-flavoured hadrons are varied as described in Section 5.1.3. The resulting changes in measured asymmetries are negligible.
* Fragmentation: Bottom and charm fragmentation uncertainties are estimated as described in Section 5.1.3.
* Semileptonic Decay Models: Systematic effects due to the modelling of semileptonic decays of heavy hadrons are studied following the recommendations in . The lepton momentum spectra in the Monte Carlo simulation are reweighted to different theoretical models, with ranges of parameters chosen such that the experimental errors are covered.
* Branching Ratios and Hadronisation Fractions: The values for the semileptonic branching ratios $`B(\mathrm{b}\mathrm{})`$, $`B(\mathrm{b}\mathrm{c}\mathrm{})`$, $`B(\mathrm{b}\overline{\mathrm{c}}\mathrm{})`$, and $`B(\mathrm{c}\mathrm{})`$ are taken from , and they are varied within their errors. Similarly, the values for the hadronisation fractions $`f(\mathrm{b}\mathrm{c}\mathrm{D}^+)`$ and $`f(\mathrm{c}\mathrm{D}^+)`$ are taken from and varied within their errors.
In the absence of any measurements, the hadronisation fraction $`f(\mathrm{b}\overline{\mathrm{c}}\mathrm{D}^{})`$ is taken from the simulation and varied by $`100\%`$ in order to assess the systematic error.
* $`𝐁\mathbf{}\overline{𝐁}`$ Mixing: The $`\mathrm{B}\overline{\mathrm{B}}`$ mixing parameter is taken as an external input. The value of $`\overline{\chi }=0.118\pm 0.006`$ is used, and the variation within its error yields the systematic uncertainty listed in Table 11.
* Background Asymmetries: Prompt leptons and slow pions from four-fermion background in the selected event sample can affect the measured asymmetries. They may lead to a rather large observed forward-backward asymmetry, but their contribution to the overall sample is small. The asymmetry from these events is varied within $`\pm 0.5`$ to determine the systematic error. The error due to the uncertainties on the bottom and charm asymmetries in radiative events is negligible.
The asymmetry in the $`y`$ distribution of mis-identified lepton and slow pion candidates is assumed to be zero for any event type. For lepton candidates, this assumption is verified within an accuracy of $`2\%`$ on an inclusive sample of tracks that pass the lepton momentum cuts and do not pass the cut on $`𝒫_{\mathrm{sig}}`$. For slow pions, a sample of tracks that pass the cuts on the slow pion momentum and transverse momentum, but do not pass the cut on $`𝒫_{\mathrm{sig}}`$, is used. The asymmetry of slow pion background is found to be zero with an error of $`4\%`$. The asymmetries from lepton and slow pion backgrounds are then varied within these limits to assess the associated systematic errors.
All the above errors are added in quadrature to yield the uncertainty given in Table 11.
* Final State QCD Corrections: As described in Section 5.1.2, half the QCD correction to $`A_{\mathrm{FB}}^\mathrm{b}`$ and $`A_{\mathrm{FB}}^\mathrm{c}`$ as computed using ZFITTER is assigned as systematic error.
#### 5.2.3 Monte Carlo Statistics
An error arises from the limited statistics in the Monte Carlo simulation that is used to predict the fractions $`f_i`$ for each subsample (see Section 4.4). This uncertainty has been evaluated by varying for each subsample in turn the contribution from each source by its statistical error, redoing the fit, and adding all observed differences in quadrature.
#### 5.2.4 Fitting Procedure
The fitting procedure has been studied in a large number of simulated experiments with the same statistics as in the actual measurements. It has been found to be essentially unbiased for the large data samples recorded at 183 GeV and 189 GeV. For the datasets at 133 GeV, 161 GeV, and 172 GeV, however, the statistical error from the fit has been found to be underestimated by a factor of up to $`1.16`$ (bottom) or $`1.05`$ (charm asymmetry), and biases of up to $`0.12\sigma `$ and $`0.06\sigma `$ have been found for the fitted bottom and charm asymmetry, respectively, where in each case $`\sigma `$ denotes the statistical error. The statistical errors are adjusted, the measurements are corrected according to the observed biases, and the full bias is treated as an additional systematic error.
#### 5.2.5 Cross-Checks
Consistent results are observed when the fit is repeated with the numbers of bins in $`𝒩_\mathrm{b}`$, $`𝒩_{\mathrm{bc}}`$, $`𝒩_\mathrm{c}`$, and $`|y|`$ varied independently between 2 and 4. In addition, all selection cuts have been varied, with no significant deviations in the results. Consistent results have been found when the slow pion tagged events are not used in the fit. In the 189 GeV data sample, the statistical error of the charm asymmetry increases from $`0.18`$ to $`0.21`$ without the information from slow pion tagged events, while the error on the bottom asymmetry stays the same.
## 6 Conclusions
Using data collected at centre-of-mass energies between $`130\mathrm{GeV}`$ and $`189\mathrm{GeV}`$ with the OPAL detector at LEP, the relative $`\mathrm{e}^+\mathrm{e}^{}\mathrm{b}\overline{\mathrm{b}}`$ production rate and the forward-backward asymmetries in $`\mathrm{e}^+\mathrm{e}^{}\mathrm{b}\overline{\mathrm{b}}`$ and $`\mathrm{e}^+\mathrm{e}^{}\mathrm{c}\overline{\mathrm{c}}`$ production have been measured to be
$$\begin{array}{cccccc}& & & & & \\ \mathrm{Energy}& R_\mathrm{b}& \multicolumn{3}{c}{A_{\mathrm{FB}}^\mathrm{b}}& A_{\mathrm{FB}}^\mathrm{c}\text{}\\ & & & & & \\ & & & & & \\ 133\mathrm{GeV}& 0.190\pm 0.023\pm 0.007& \hfill 0.19& \pm \mathrm{\hspace{0.17em}0.30}& \pm \mathrm{\hspace{0.17em}0.12}\hfill & 0.50_{0.32}^{+0.31}\pm 0.13\text{}\\ 161\mathrm{GeV}& 0.195\pm 0.035\pm 0.007& \hfill 0.03& {}_{0.42}{}^{}{}_{}{}^{+0.45}& \pm 0.11\hfill & 0.87_{0.60}^{+0.58}\pm 0.12\text{}\\ 172\mathrm{GeV}& 0.091\pm 0.034\pm 0.005& \hfill 0.82& {}_{0.72}{}^{}{}_{}{}^{+0.67}& \pm 0.14\hfill & 0.69_{0.53}^{+0.49}\pm 0.12\text{}\\ 183\mathrm{GeV}& 0.213\pm 0.020\pm 0.009& \hfill 0.77& {}_{0.24}{}^{}{}_{}{}^{+0.23}& \pm 0.10\hfill & 0.55_{0.28}^{+0.27}\pm 0.11\text{}\\ 189\mathrm{GeV}& 0.158\pm 0.012\pm 0.007& \hfill 0.63& {}_{0.16}{}^{}{}_{}{}^{+0.15}& \pm 0.10\hfill & 0.50_{0.19}^{+0.18}\pm 0.11\text{}\end{array}$$
where in each case, the first error is statistical and the second systematic. These values are illustrated in Figure 8 together with the dependences on the centre-of-mass energy as predicted in the Standard Model. For all measurements, good agreement is observed with the Standard Model expectation. The measurements of $`\mathrm{b}\overline{\mathrm{b}}`$ production presented in this paper supersede the previously published values of references .
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# The Dark Side of the Solar Neutrino Parameter Space[1]
## Abstract
Results of neutrino oscillation experiments have always been presented on the $`(\mathrm{sin}^22\theta ,\mathrm{\Delta }m^2)`$ parameter space for the case of two-flavor oscillations. We point out, however, that this parameterization misses the half of the parameter space $`\frac{\pi }{4}<\theta \frac{\pi }{2}`$ (“the dark side”), which is physically inequivalent to the region $`0\theta \frac{\pi }{4}`$ (“the light side”) in the presence of matter effects. The MSW solutions to the solar neutrino problem can extend to the dark side, especially if we take the conservative attitude to allow higher confidence levels, ignore some of the experimental results in the fits, or relax theoretical predictions. Furthermore, even the so-called “vacuum oscillation” solution distinguishes the dark and the light sides. We urge experimental collaborations to present their results on the entire parameter space.
preprint: UCB-PTH-00/03, LBNL-45023
In the Standard Model of particle physics, neutrinos are strictly massless. Recently, however, the Super-Kamiokande collaboration studied atmospheric neutrinos and reported a strong evidence for neutrino oscillations , and hence a finite neutrino mass. The most likely interpretation of their data is the oscillation between $`\nu _\mu `$ and $`\nu _\tau `$. This made it also natural to interpret another long-standing issue in neutrino physics, the deficit of the solar $`\nu _e`$ flux , in terms of neutrino oscillations. However, the solar neutrino deficit has not been regarded as convincing evidence for neutrino oscillations in the community. The reason is probably multifold but two main objections are the following. (1) Neutrino experiments are so difficult that it is possible that some of the data are not entirely correct. (2) The physics of the Sun is so complex that the neutrino flux calculations in the Standard Solar Model (SSM) may have underestimated the theoretical uncertainties.
To resolve this situation, a new generation of solar neutrino experiments, such as Super-Kamiokande, SNO, Borexino, GNO, KamLAND, etc, is looking for an evidence for solar neutrino oscillations without relying on the SSM in well-understood experimental environments. They aim not only at establishing oscillations but also at overdetermining the solution in the next few years. Such data will eventually supersede data from the past experiments. It is therefore important to analyze the future data without too much prejudice based on the past data.
In this letter, we point out that the study of neutrino oscillations on the $`(\mathrm{\Delta }m^2,\mathrm{sin}^22\theta )`$ parameter space done traditionally is incomplete, since it covers only the range $`0\theta \frac{\pi }{4}`$ (“the light side”). Indeed, some of the solutions to the solar neutrino puzzle extend to the other half of the parameter space $`\frac{\pi }{4}<\theta \frac{\pi }{2}`$, which we call “the dark side,” and hence it is phenomenologically necessary to include both halves of the parameter space. This is especially true once one employs a more conservative attitude which either allows higher confidence levels, ignores some of the experimental data (especially Homestake ), or relaxes the theoretical prediction on the <sup>8</sup>B flux.
Neutrino oscillations occur if neutrino mass eigenstates are different from neutrino weak eigenstates. Assuming that only two neutrino states mix, the relation between mass eigenstates ($`\nu _1`$ and $`\nu _2`$) and flavor eigenstates (for example $`\nu _e`$ and $`\nu _\mu `$) is simply given by
$`|\nu _1=\mathrm{cos}\theta |\nu _e\mathrm{sin}\theta |\nu _\mu ,`$ (1)
$`|\nu _2=\mathrm{sin}\theta |\nu _e+\mathrm{cos}\theta |\nu _\mu ,`$ (2)
where $`\theta `$ is the vacuum mixing angle. The mass-squared difference is defined as $`\mathrm{\Delta }m^2m_2^2m_1^2`$. We are interested in the range of parameters that encompasses all physically different situations. First, observe that Eq. (1) is invariant under $`\theta \theta +\pi `$, $`|\nu _e|\nu _e`$, $`|\nu _\mu |\nu _\mu `$, and hence the ranges $`[\frac{\pi }{2},\frac{\pi }{2}]`$ and $`[\frac{\pi }{2},\frac{3\pi }{2}]`$ are physically equivalent. Next, note that it is also invariant under $`\theta \theta `$, $`|\nu _\mu |\nu _\mu `$, $`|\nu _2|\nu _2`$, hence it is sufficient to only consider $`\theta [0,\frac{\pi }{2}]`$. Finally, it can also be made invariant under $`\theta \frac{\pi }{2}\theta `$, $`|\nu _\mu |\nu _\mu `$ by relabeling the mass eigenstates $`|\nu _1|\nu _2`$, i.e. $`\mathrm{\Delta }m^2\mathrm{\Delta }m^2`$. Thus, we can take ($`\mathrm{\Delta }m^2>0`$) without loss of generality. All physically different situations are obtained by allowing $`0\theta \frac{\pi }{2}`$.
For the case of oscillations in the vacuum, the survival probability is given by
$$P(\nu _e\nu _e)=1\mathrm{sin}^22\theta \mathrm{sin}^2\left(1.27\frac{\mathrm{\Delta }m^2}{E}L\right).$$
(3)
Here, $`\mathrm{\Delta }m^2`$ is given in eV$`{}_{}{}^{2}/c^4`$, $`E`$ in GeV, and $`L`$ in km. In this case the oscillation phenomenon can be parameterized by $`\mathrm{\Delta }m^2`$ and $`\mathrm{sin}^22\theta `$, since $`\theta `$ and $`\frac{\pi }{2}\theta `$ yield identical physics. Therefore we can restrict ourselves to $`0\theta \frac{\pi }{4}`$, and use the parameter space $`(\mathrm{\Delta }m^2,\mathrm{sin}^22\theta )`$ without any ambiguity. This is indeed an adequate parameterization for reactor antineutrino oscillation experiments, short-baseline accelerator neutrino oscillation experiments, and $`\nu _\mu \nu _\tau `$ atmospheric neutrino oscillation experiments.
In certain cases, however, neutrino-matter interactions can dramatically change the oscillation probability . These matter effects are particularly important in explaining the solar $`\nu _e`$ flux deficit in terms of neutrino oscillations. In the presence of matter effects, Eq. (3) is modified to
$`P(\nu _e\nu _e)=P_1\mathrm{cos}^2\theta +(1P_1)\mathrm{sin}^2\theta `$ (4)
$`\sqrt{P_c(1P_c)}\mathrm{cos}2\theta _M\mathrm{sin}2\theta \mathrm{cos}\left(2.54{\displaystyle \frac{\mathrm{\Delta }m^2}{E}}L+\delta \right),`$ (5)
where $`P_c`$ is the hopping probability, $`\theta _M`$ is the mixing angle at the production point, $`P_1=P_c\mathrm{sin}^2\theta _M+(1P_c)\mathrm{cos}^2\theta _M`$, and $`\delta `$ is a phase induced by the matter effects, which is not important for our purposes. See Ref. for notation. Because $`P_1`$ depends on $`\mathrm{\Delta }m^2\mathrm{cos}2\theta `$, the half of the parameter space traditionally considered $`0\theta \frac{\pi }{4}`$ (the light side) is physically inequivalent to the other half $`\frac{\pi }{4}<\theta \frac{\pi }{2}`$ (the dark side). However, all data analysis have been reported on the $`(\mathrm{\Delta }m^2,\mathrm{sin}^22\theta )`$ plane with positive $`\mathrm{\Delta }m^2`$ only for solar neutrino experiments and hence only half of the parameter space has been analyzed. Even though the dark side has been studied in the context of three-flavor and four-flavor neutrino oscillations, the importance of studying both halves for the simplest case of two-flavor oscillations has been largely ignored in the literature. There also appeared to be a misconception in the literature that physics was discontinuous at maximal mixing $`\theta =\frac{\pi }{4}`$. For instance, matter effects in the Earth were once thought to disappear as the mixing approached maximal. However, the authors of Ref. emphasized that the matter effects remain important even for the maximal mixing, and the present authors further showed that physics is completely continuous beyond $`\theta =\frac{\pi }{4}`$ . One can still retain the dark side with only $`0\theta \frac{\pi }{4}`$ if a separate parameter space with $`\mathrm{\Delta }m^2<0`$ is added. This is indeed what Super-Kamiokande did in the case of $`\nu _\mu \nu _s`$ oscillations of atmospheric neutrinos . However, as argued in , it is more natural to use $`0\theta \frac{\pi }{2}`$ with the fixed sign of $`\mathrm{\Delta }m^2`$ to exhibit the continuity of the physics between the two halves of the parameter space.
Part of the reason why the dark side has been neglected in the literature is that it is impossible to obtain $`\nu _e`$ survival probabilities less than one half when the two mass eigenstates are incoherent, i.e., when the last term in Eq. (5) is absent. (This occurs in the so-called “MSW region” $`10^8\mathrm{\Delta }m^210^3\mathrm{eV}^2`$ .) Indeed, the data from the Homestake experiment used to be about a quarter of the SSM prediction, and this could have been an argument for dropping the dark side entirely in the MSW region. However, the change from BP95 to BP98 calculations increased the Homestake result to about a third of the SSM with a relatively large theoretical uncertainty. Therefore it is quite possible that the “MSW solutions” extend to the dark side as well. Moreover, some people question the SSM and/or the Homestake experiment, and perform fits by ignoring either (or both) of them . We show below that some of the MSW solutions indeed extend to the dark side and hence it is necessary to explore the dark side experimentally. If we further relax the theoretical prediction on the <sup>8</sup>B solar neutrino flux and/or ignore one of the solar neutrino experiments in the global fit, the preferred regions extend even deeper into the dark side.
Another possible reason for disregarding the dark side is that the so-called “vacuum oscillation region” ($`\mathrm{\Delta }m^210^9`$ eV<sup>2</sup>) was believed to be the same in the light and dark sides. This is because $`P_1`$ approaches $`\mathrm{cos}^2\theta `$ for $`\mathrm{\Delta }m^210^9\mathrm{eV}^2(E/\mathrm{MeV})`$ and Eq. (5) reduces to Eq. (3). It is remarkable, however, that low-energy (especially pp) neutrinos do not reach this limit for $`\mathrm{\Delta }m^210^{10}\mathrm{eV}^2`$ and hence the preferred regions are different in the light and the dark sides . This observation also implies that the separation of the MSW region and the vacuum oscillation region as traditionally done in the global fits is artificial and misleading. It is important to study the entire range of $`\mathrm{\Delta }m^2`$ continuously.
If $`\mathrm{sin}^22\theta `$ is not a good parameter, what is the alternative? Two suggestions have been made in the literature. One is $`\mathrm{sin}^2\theta `$, which is natural since the matter effect depends directly on $`\mathrm{sin}^2\theta `$ . If plotted on the linear scale, pure vacuum oscillations would yield physics reflection-symmetric around $`\mathrm{sin}^2\theta =0.5`$. If plotted on the log scale, the reflection symmetry is lost, but it is still a useful parameterization as physics is completely continuous and smooth from the light to the dark side. Another possible parameterization is $`\mathrm{tan}^2\theta `$, which retains the reflection symmetry for pure vacuum oscillation around $`\mathrm{tan}^2\theta =1`$ if plotted on the log scale . We employ $`\mathrm{tan}^2\theta `$ for the analysis below because we would like to use the log scale to present the MSW solutions as well as the importance of the matter effect on the “vacuum oscillation” region at the same time. Note that the Jacobian from $`\mathrm{sin}^2\theta `$ or $`\mathrm{tan}^2\theta `$ to $`\mathrm{sin}^22\theta `$ is singular at $`\theta =\frac{\pi }{4}`$ and plots with $`\mathrm{sin}^22\theta `$ will display unphysical singular behavior there .
We next present the results of global fits to the current solar neutrino data from water Cherenkov detectors (Kamiokande and Super-Kamiokande) , a chlorine target (Homestake) and gallium targets (GALLEX and SAGE) on the full parameter space. We do not include the spectral data from Super-Kamiokande as it appears to be still evolving with time. The fit is to the event rates measured at these experiments only. In computing the rates we include not only the pp, <sup>7</sup>Be, and <sup>8</sup>B neutrinos, but also the <sup>13</sup>N, <sup>15</sup>O, and pep neutrinos. We use Eq. (5) with $`P_c`$ computed in the exponential approximation for the electron number density profile in the Sun, and properly account for neutrino interactions in the Earth during the night with a realistic Earth electron number density profile by numerically solving Schrödinger equation as described in . Since the mixing angle at the production point in the Sun’s core depends on the electron number density, we integrate over the production region numerically. We treat the correlations between the theoretical uncertainties at different experiments following Ref. . To insure a smooth transition between the MSW and the vacuum oscillation region, we integrate over the energy spectrum (including the thermal broadening of the <sup>7</sup>Be neutrino “line”) for $`\mathrm{\Delta }m^210^8`$ eV<sup>2</sup> and average the neutrino fluxes over the seasons. For $`\mathrm{\Delta }m^2>10^8`$ eV<sup>2</sup> we treat the two mass eigenstates as incoherent. Results are completely smooth at $`\mathrm{\Delta }m^2=10^8`$ eV<sup>2</sup>, as expected. This allows us to fit the data from $`\mathrm{\Delta }m^2=10^{11}`$$`10^3`$ eV<sup>2</sup> all at once, unlike previous analyses which separate out the “vacuum oscillation region” from the rest.
As was mentioned earlier, we take the global fit to the currently available data only as indicative of the ultimate result because we expect much better data to be collected in the near future to eventually supersede the current data set. We would like to keep our minds open to surprises such as the possibility that one of the earlier experiments was not entirely correct or that the theoretical uncertainty in the flux prediction was underestimated. In this spirit, we employ more conservative attitudes in the global fit than most of the analyses in the literature in the following three possible ways. (1) We allow higher confidence levels, such as 3 $`\sigma `$. (2) We relax the theoretical prediction on the neutrino flux. (3) We ignore some of the experimental data in the fit.
The global fit results are presented in Fig. 1 at the 2 $`\sigma `$ (95% CL) and 3 $`\sigma `$ (99.7% CL) levels defined by $`\chi ^2\chi _{\mathrm{𝑚𝑖𝑛}}^2`$ for two degrees of freedom. It is noteworthy that both the LMA and LOW solutions (we use the nomenclature introduced in ) extend to the dark side at the 3 $`\sigma `$ level. At 99% CL, however, the LMA solution is confined to the light side. This result is consistent with the two-flavor limit of the three-flavor analysis in and the four-flavor analysis in , where the spectral data is included and the LOW solution extends into the dark side at 99% CL. Another interesting fact is that the LOW solution is smoothly connected to the VAC solution, where the preferred region is clearly asymmetric between the light and the dark sides. Note that, at $`\mathrm{\Delta }m^210^9\mathrm{eV}^2`$, the allowed region is bigger in the dark side. The region $`10^9<\mathrm{\Delta }m^2<10^8\mathrm{eV}^2`$ was, to the best of the authors’ knowledge, never studied fully in the literature and this result demonstrates the need to study the entire $`\mathrm{\Delta }m^2`$ region continuously without the artificial separation of the “MSW region” and “vacuum oscillation region,” as traditionally done in the literature.
We next present a fit where the theoretical prediction of the <sup>8</sup>B flux is relaxed. Even though the helioseismology data constraints the sound speed down to about 5% of the solar radius , the core region where <sup>8</sup>B neutrinos are produced is still not constrained directly. Given the sensitive dependence of the <sup>8</sup>B flux calculation on the core temperature $`\mathrm{\Phi }_{{}_{}{}^{8}B}T^{22}`$, we may consider it as a free parameter in the fit. This can be done within the formalism of Ref. by formally sending the error in $`C_{\mathrm{Be}}`$ to infinity. The result is presented in Fig. 2. The preferred region extends more into the dark side than the previous fit. Even though the LMA and LOW solutions are connected in this plot, the lack of a large day-night asymmetry at Super-Kamiokande would eliminate the range $`3\times 10^7\mathrm{\Delta }m^210^5\mathrm{eV}^2`$ for $`0.2\mathrm{tan}^2\theta <1`$ . It is important for Super-Kamiokande to report their exclusion region on the dark side.
Finally, we present a fit where the event rate measured at the Homestake experiment is not used in Fig. 3. This may be a sensible exercise given that the neutrino capture efficiency was never calibrated in this experiment. The preferred region extends into the dark side even at the 95% CL. Note also the asymmetry between the dark and the light sides even for $`\mathrm{\Delta }m^2<10^9\mathrm{eV}^2`$.
We expect the data of the current and next generation of solar neutrino experiments, such as Super-Kamiokande, SNO, GNO, Borexino, KamLAND, to eventually supersede the current data set. Therefore we regard the above global fits only as estimates of the ultimate results. The most important point is that all experimental collaborations should report their results, both exclusion and measurements, on both sides of the parameter space, without unnecessary theoretical bias towards the light side. We strongly urge the experimental collaborations to consider this point.
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# Paraxial propagation of a quantum charge in a random magnetic field
## I Introduction
The paper addresses the problem of quantum transport of a charge in an inhomogeneous static random magnetic field. In recent years, this or related problems have been met in a number of contexts in physics of 2-dimensional systems. For instance, in the composite fermion model of the Fractional Quantum Hall Effect, a (fictitious) random magnetic field is the environment which controls dynamics of effective charge carriers . One meets the fluctuating gauge fields in some models of high temperature superconductors , where the gauge field is the tool to impose the constrain of no double occupancy in the $`tJ`$ model . Besides, stochastically inhomogeneous magnetic field can be experimentally created by various ways. For example, the field is irregular near the surface of a superconductor in an external magnetic field if the Abrikosov flux lattice is disordered; depending on the experimental conditions, the magnetic field inhomogeneities may be weak and smooth, or the field may be concentrated in an irregular array of flux tubes. Various aspects of transport in the magnetic field of the Abrikosov vortices, the weak localization in particular, have been studied in Refs. . In recent years, the random magnetic field problem has been an active subject area Refs. .
The formulation of the problem is as follows. A particle with the electric charge $`e`$ and the mass $`m`$ moves on the $`xy`$ plane subject to a vector potential potential $`𝑨(A_x,A_y)`$ generated by a magnetic field $`b(x,y)=\left(\mathrm{𝐫𝐨𝐭}𝑨\right)_z`$. Two random fields models are considered in the paper. In the first one, the magnetic field $`b(𝒓)`$ is a random Gaussian variable with zero average, $`b=0`$, specified by the correlator
$$b(𝒓)b(𝒓^{})=\left(\frac{\mathrm{\Phi }_0}{2\pi }\right)^2\frac{1}{^2}\delta (𝒓𝒓^{}),$$
(1)
where $`\mathrm{\Phi }_0=\frac{hc}{e}`$ is the flux quantum; the strength of the random magnetic field is characterized via the length $``$ the meaning of which is that the magnetic flux through the area $`^2`$ is typically of order of $`\mathrm{\Phi }_0`$. The random field is assumed to be weak in the sense that $``$ much exceeds the wave length $`\mathrm{\lambda ̄}\mathrm{}/p`$, $`p`$ being the particle momentum. In another model which is motivated by fractional statistics theories, the gauge potential is created by a random array of the Aharonov-Bohm flux lines. A system of the Abrikosov vortices (e.g. in the gate of a MOSFET transistor ) may serve as an experimental realization of the Aharonov-Bohm array if the particle wave length much exceeds the vortex (magnetic) size.
In the random magnetic field case, the traditional approach of the theory of disordered systems meets difficulties on the very first steps. Indeed, the simplest object that is the single particle Green’s function $`G(1,2)=\psi (𝒓_1)\psi ^{}(𝒓_2)`$, is not gauge invariant and the physical meaning of its averaging with respect to the vector potential $`𝑨`$ generated by the random field is not clear. One may define a gauge invariant combination $`\stackrel{~}{G}(1,2)=\psi (𝒓_1)\psi ^{}(𝒓_2)\mathrm{exp}[i\frac{e\mathrm{}}{c}_{𝒞_{12}}𝑑𝒍\mathbf{}𝑨]`$ where the path $`𝒞_{12}`$ connects the points $`𝒓_1`$ and $`𝒓_2`$. Albeit gauge independent, $`\stackrel{~}{G}(1,2)`$ essentially depends on the choice of the path $`𝒞_{12}`$. With the point $`𝒓_{1,2}`$ connected by the straight line, the field averaged $`\stackrel{~}{G}(1,2)`$ has been found in Ref..
Another problem is the diverging scattering rate $`\frac{1}{\tau }`$. For small scattering angles $`\varphi `$, the differential cross-section behaves like $`\frac{1}{\varphi ^2}`$ so that the scattering total cross-section is infinite. In other words, the life time of a state with the definite momentum is zero. The conventional diagram technique where $`\frac{\mathrm{}}{\tau }`$ is assumed to be small compared with the kinetic energy, becomes questionable. On the other hand, it is known that in gauge-invariant response functions the self energy enters in combination with the vertex corrections and the divergence cancels out. In Ref., it has been attempted to introduce a physically sensible gauge-invariant “single-particle time” $`\tau `$ as a parameter entering the Landau level broadening.
The main purpose of this paper is to develop a scheme which allows one to study the most singular part of interaction with random magnetic field, that is the near forward scattering.
To pinpoint the physics behind the theoretical difficulties, consider first propagation of a plane wave. It is common in wave mechanics to analyze propagation in terms of wave fronts, i.e. the surfaces (lines in 2D) of constant phase. The property of the wave front line is that the probability current is locally perpendicular to the line. In the magnetic field, the phase of the wave function is ill-defined because of the gauge freedom. Nevertheless, one can construct a gauge invariant quantity $`\chi `$ defined on a line, which in a limited sense plays the role of the phase: Given the wave function $`\psi (𝒓)`$ and the vector potential $`𝑨`$, the phase $`\chi (𝒔)`$ for the points $`𝒔`$ on a line $`S`$ is defined through its differential as
$$d\chi =\frac{m}{\mathrm{}|\psi |^2}𝒋\mathbf{}d𝒔$$
(2)
where $`𝒋=\frac{1}{m}\mathrm{}\left(\psi ^{}\left(\frac{\mathrm{}}{i}\mathbf{}\frac{e}{c}𝑨\right)\psi \right)`$ is the probability current density. Provided the line $`S`$ does not have self-intersections, $`\chi `$ is an unique function of $`𝒔`$.
If $`\chi `$ is a constant, i.e. $`𝒋\mathbf{}d𝒔=0`$ and $`𝒋d𝒔`$, the local current and the normal to $`S`$ are parallel, so that $`S`$ is a wave front. If $`\chi `$ is a slowly varying function, $`\mathrm{\lambda ̄}\left(\chi (𝒔)\chi (\mathrm{𝟎})\right)`$ gives the local distance from $`S`$ to the wave front passing through the point $`𝒔=0`$.
Consider now how the random magnetic field affects the wave front upon its propagation. Take a state $`\psi `$, for which the line $`x=0`$ is a wave front corresponding to the propagation in the positive $`x`$-direction ($`j_x`$ ¿0). To satisfy the requirement $`d\chi =0`$, the wave function is
$`\psi (x=0,y)=\mathrm{exp}\left[i{\displaystyle \frac{e\mathrm{}}{c}}{\displaystyle \underset{\mathrm{}}{\overset{y}{}}}𝑑yA_y(x=0,y)\right],`$along the wave front $`x=0`$ (choice of lower limit of integration is not important).
To find the profile of the wave front having advanced from $`x=0`$ to a finite $`x`$, one can apply the usual eikonal-type approximation where the field affects only the phase of the wave function through the factor $`\mathrm{exp}[\frac{ie}{\mathrm{}c}𝑑𝒍\mathbf{}𝑨]`$, $`d𝒍\widehat{x}`$ being along the direction of propagation:
$`\psi (x,y)=\mathrm{exp}\left[{\displaystyle \frac{i}{\mathrm{}}}\left(px+{\displaystyle \frac{e}{c}}{\displaystyle \underset{\mathrm{}}{\overset{y}{}}}𝑑yA_y(0,y)+{\displaystyle \frac{e}{c}}{\displaystyle \underset{0}{\overset{x}{}}}𝑑xA_x(x,y)\right)\right].`$Integrating Eq.(2), one finds the phase $`\chi (y;x)`$ as a function of $`y`$ for fixed $`x`$. For the phase difference $`\mathrm{\Delta }\chi (y_1,y_2;x)=\chi (y_2;x)\chi (y_1;x)`$, one gets after simple calculations:
$$\mathrm{\Delta }\chi (y_1,y_2;x)=2\pi \frac{\mathrm{\Phi }(y_1,y_2;x)}{\mathrm{\Phi }_0}$$
(3)
where $`\mathrm{\Phi }(y_1,y_2;x)`$ is the magnetic flux through the area enclosed by the path $`(0,y_1)(0,y_2)(x,y_2)(x,y_1)(0,y_1)`$ (see fig. 2).
A non-local character of the interaction with magnetic filed is clearly seen from Eq.(3): the phase difference is controlled by the flux rather than the magnetic field in the vicinity of the particle trajectories. The non-locality is obviously of the Aharonov-Bohm type.
Averaging with respect of the random magnetic field Eq.(1), one gets the variation of the phase difference
$`(\mathrm{\Delta }\chi )^2(y_1,y_2;x)={\displaystyle \frac{1}{2}}{\displaystyle \frac{|x||y_1y_2|}{^2}}.`$Most notable feature here is that $`(\mathrm{\Delta }\chi )^2`$ grows with the separation $`|y_1y_2|`$ (cf. Ref.). For points separated by the distance $`\mathrm{\Delta }y`$, the random phase difference is of order of $`1`$ when the wave front advances to $`\mathrm{\Delta }x^2/\mathrm{\Delta }y`$. One sees that an infinite plane wave, for which $`\mathrm{\Delta }y\mathrm{}`$, looses its coherence immediately, whatever small is the propagation distance $`\mathrm{\Delta }x`$. (See Section V for a more formal derivation.) These qualitative arguments explain the actual meaning of the zero life time and show that it is not an artifact arising e.g. due to a violation of the gauge invariance.
To handle the anomalously intensive forward scattering, one needs a method suitable for a nonperturbative analysis of the small-angle multiple scattering. For this, the paraxial (parabolic) approximation to the Schrödinger equation is chosen in the paper. The paraxial theory is applicable when the particle moves mainly in the direction of an “axis” and the momentum transverse to the axis remains always small. The paraxial approximation to the wave equation is most popular in optics where it gives a convenient description of light beams propagating in optical systems, their diffraction, focusing etc . Taking scattering and diffraction broadening on equal footing, the paraxial approximation is more generally applicable then the eikonal one .
To make the paper self-contained, the derivation of the paraxial approximation is outlined in Sect.II. The case of magnetic field is considered in Sect.II A where a scheme for description of scattering by magnetic field is suggested. The scheme is in a sense gauge invariant, gauge freedom revealing itself only in the overall phases. As a limiting case, one recovers the well-known eikonal approximation (see Sect. II B).
To illustrate usage of the paraxial approximation, a simple problem of scattering of a charged particle by the Aharonov-Bohm magnetic flux line is considered in Sect.III. This (or equivalent) problem is of interest in a broad variety of contexts extending from the cosmic string theory to superfluids (the Iordanskii force) and superconductors where the scattering of excitations by quantized vortex lines controls the vortex dynamics. Although the exact solution to the scattering problem has been known since the original paper of Aharonov and Bohm (see also review and references therein), certain controversy in the analysis and the interpretation of the solution still remains. Different opinions exist in the literature about the existence of the transverse force exerted by the Aharonov-Bohm line or a superfluid vortex. On the basis of the left-right symmetry in the Aharonov-Bohm differential cross-section, the authors of Refs. and Ref. have come to the conclusion that the line does not exert any Lorentz-like force (translated as the Iordanskii force in a superfluid). Other authors, predict a finite force. Due to its simplicity, the paraxial solution allows one to perform a detailed analysis and resolve the controversy.
It is shown in In Sect.IV, that the paraxial scattering theory becomes manifestly gauge invariant if formulated in terms of by-linear in $`\psi `$ and $`\psi ^{}`$ object that is the density matrix $`\rho `$. The evolution of the density matrix is given by a gauge invariant two-particle Green function.
As discussed in Sect.IV A, the paraxial 2D stationary equation with inhomogeneous magnetic field can be written as a non-stationary 1D Schrödinger equation for a particle in a time-dependent electric field. This mapping allows one to present stationary solutions to the 2D magnetic field problem as the Feynman path integral for the effective 1D problem.
In Sect.V, the paraxial theory is applied to the model of $`\delta `$-correlated random magnetic field. It turns out to be possible to evaluate the Feynman path integral and by this to find a (paraxially) exact expression for the two-particle Green’s function averaged with respect to the magnetic field fluctuations. It is shown that the density matrix evolution can be mapped to the Boltzmann kinetic equation.
In Sect.VI, another model is considered where the random gauge filed is generated by a random array of Aharonov-Bohm fluxes. The flux lines are randomly distributed in the plane, the flux of a line, $`\mathrm{\Phi }`$, is distributed with the probability $`p(\mathrm{\Phi })`$. The Aharonov-Bohm array may create an effective magnetic field $`\stackrel{~}{B}`$ if $`p(\mathrm{\Phi })`$ is asymmetric, $`p(\mathrm{\Phi })p(\mathrm{\Phi })`$.
In Sect.VI A, the Boltzmann equation for charge subject to a Gaussian random magnetic field or field of AB-array is derived. With the help of the Boltzmann equation, the resistivity tensor is found. Finally, in Sect.VI B we discuss the density of states of the levels due to the quantization of motion in the field $`\stackrel{~}{B}`$.
The results are summarized in Section VII
## II Paraxial approximation
The paraxial approximation allows one to construct a family of solutions to the Schrödinger equation which are close to to the plane wave with a certain momentum $`𝒑_0`$. The wave function $`\mathrm{\Psi }_{Sch}`$, a solution to the stationary Schrödinger equation, is presented as
$$\mathrm{\Psi }_{Sch}(𝒓)=\mathrm{\Psi }(𝒓)e^{\frac{i}{\mathrm{}}𝒑_0\mathbf{}𝒓},$$
(4)
where the envelope paraxial function $`\mathrm{\Psi }(𝒓)`$ is supposed to be slowly varying at the distances of order of the wave length $`\mathrm{\lambda ̄}=\frac{\mathrm{}}{p_0}`$.
The Schrödinger equation reads
$$\left(E(\widehat{𝒑})+U\right)\mathrm{\Psi }_{Sch}=E\mathrm{\Psi }_{Sch}$$
(5)
$`E(𝒑)=\frac{1}{2m}𝒑^2`$ and $`U(𝒓)`$ being the kinetic and potential energy respectively. Given $`𝒑_0`$, the family of solutions in Eq.(4) corresponds to the eigen-energy $`E=E(𝒑_0)`$ and the velocity $`𝒗=\frac{E(𝒑_0)}{𝒑_0}`$. Inserting Eq.(4) into Eq.(5), one gets equation for $`\mathrm{\Psi }`$,
$`\left(\stackrel{~}{E}\left(\widehat{𝒑}\right)E(𝒑_0)+U\right)\mathrm{\Psi }=0,\stackrel{~}{E}\left(\widehat{𝒑}\right)e^{\frac{i}{\mathrm{}}𝒑_0\mathbf{}𝒓}E\left(\widehat{𝒑}\right)e^{\frac{i}{\mathrm{}}𝒑_0\mathbf{}𝒓},`$which is still exact. The operator $`\stackrel{~}{E}(\widehat{𝒑})=E(𝒑_0+\frac{\mathrm{}}{i}\mathbf{})`$ acting on the slowly varying function $`\mathrm{\Psi }`$ is approximated in the paraxial theory as
$`\stackrel{~}{E}(\widehat{𝒑})E(𝒑_0)+{\displaystyle \frac{\mathrm{}}{i}}𝒗\mathbf{}\mathbf{}+{\displaystyle \frac{1}{2m}}\left({\displaystyle \frac{\mathrm{}}{i}}\mathbf{}_{}\right)^2,`$here $`\mathbf{}_{}`$ denotes the gradient in the direction perpendicular to $`𝒗`$.
The paraxial approximation to the Schrödinger equation reads
$$\left(i\mathrm{}𝒗\mathbf{}\mathbf{}+\frac{\mathrm{}^2}{2m}\mathbf{}_{}^2U\right)\mathrm{\Psi }=0.$$
(6)
(Condition of applicability are discussed later). The main feature of the paraxial approximation Eq.(6), is that it is of first order differential equation relative to the coordinate in the direction of the propagation $`x=𝒓\mathbf{}𝒗/v`$.
Introducing formally “time” $`\tau =x/v`$, Eq.(6) takes the form, of the time dependent Schrödinger equation in a reduced space dimension:
$$i\mathrm{}\frac{}{\tau }\mathrm{\Psi }=\left(\frac{\mathrm{}^2}{2m}\mathbf{}_{}^2+U\right)\mathrm{\Psi },$$
(7)
This formal analogy allows one to discuss the stationary solutions in terms of the wave moving in the direction $`𝒗`$, and call $`x`$ the current coordinate of the wave. The Feynman path integral equivalent to Eq.(7) gives an alternative method of solving the equation.
The probability current $`𝑱`$ is derived from the standard expression $`𝑱=\frac{\mathrm{}}{m}\mathrm{}\mathrm{\Psi }_{Sch}^{}\mathbf{}\mathrm{\Psi }_{Sch}`$ and the definition Eq.(4). In the main approximation, the components parallel, $`𝑱_{||}`$, and perpendicular, $`𝑱_{}`$, relative to $`𝒗`$, are
$$𝑱_{||}=𝒗|\mathrm{\Psi }|^2,𝑱_{}=\frac{\mathrm{}}{m}\mathrm{}\mathrm{\Psi }^{}\mathbf{}_{}\mathrm{\Psi }.$$
(8)
These expressions are consistent with the current conservation and Eq.(6) or Eq.(7): Indeed, the continuity equation which follows from Eq.(7),
$$\frac{|\mathrm{\Psi }|^2}{\tau }+\text{div }𝑱_{}=0$$
(9)
is equivalent to $`\text{div }𝑱=0`$ with $`𝑱`$ from Eq.(8). The continuity in Eq.(9) means that the paraxial wave function can be normalized:
$`{\displaystyle 𝑑𝒓_{}|\mathrm{\Psi }|^2}=1`$fixing the total flux in the beam to $`𝒗`$.
The required solution to the paraxial equation can be chosen by imposing a proper boundary condition to Eq.(6) (“initial” condition in the case of Eq.(7))
#### a Paraxial approximation: conditions of applicability
The paraxial theory is based on the approximation $`E(𝒑)E(𝒑_0)v\delta p_{||}+\frac{1}{2m}(\delta 𝒑_{})^2`$, $`\delta 𝒑𝒑𝒑_0`$ where the term $`\frac{1}{2m}(\delta 𝒑_{||})^2`$ is neglected. This is justifiable if the angle, $`\theta \frac{\delta p_{}}{p_0}`$, between $`𝒑`$ and $`𝒑_0`$ is small, $`\theta 1`$. If the motion is free, $`\delta p_{}\frac{\mathrm{}}{w}`$ where $`w`$ is the width of the beam (defined by the boundary conditions). Paraxial approximation is therefore applicable if the beam is wide in the scale of the wave length $`\mathrm{\lambda ̄}`$,
$$w\frac{\delta p_{}}{p_0}\mathrm{\lambda ̄}.$$
(10)
The scattering by the potential $`U`$ changes the angle by $`\theta _{\text{scat}}U/E`$ . The paraxial approximation requires the small angle scattering to be dominant, $`\theta \theta _{\text{scat}}1`$, so that the theory is applicable only for fast particles: $`E|U|`$. It is important however, that, as in the case of the eikonal approximation , the theory is applicable beyond the Born approximation: the phase shift $`\delta \phi \frac{Ua}{\mathrm{}v}\frac{U}{E}\frac{pa}{\mathrm{}}`$, $`a`$ being the thickness of the layer where $`U0`$, may be large even for fast particles.
Unlike the eikonal approximation where only the phase variations are taken into account, the paraxial theory allows also for the change of the profile of the beam, i.e. $`|\mathrm{\Psi }(𝒓)|`$ due to diffraction (or, equivalently, to broadening of the wave packet in the language of Eq.(7)). If $`w`$ is a typical size of the transverse structure defined by either the initial condition or scattering, the “diffraction blurring of the image” happens when the beam travels the distance $`x_{\text{diff}}`$,
$$x_{\text{diff}}\frac{w^2}{\mathrm{\lambda ̄}}.$$
(11)
The diffraction length $`x_{\text{diff}}`$ is the typical distance for the paraxial approximation while region $`xx_{\text{diff}}`$ is described by the eikonal approximation.
Applicability of the approximation at large distances requires further analysis. The paraxial relation $`\delta p_{||}=\frac{1}{2p_0}(\delta 𝒑_{})^2`$ is valid up to a small correction $`(\delta p_{||})_2`$ due to the neglected quadratic term: $`(\delta p_{||})_2\frac{1}{2p_0}(\delta p_{||})^2`$. In the main approximation, $`(\delta p_{||})_2p_{}^4/p_0^3`$. Although small in comparison with $`\delta p_{||}`$, the correction is important at long enough distances; it can be ignored only if $`(\delta p_{||})_2x/\mathrm{}<1`$. Thus, the paraxial approximation is reliable if the distance traveled by the beam is not not too large,
$$x<\frac{w^4}{\mathrm{\lambda ̄}^3}\left(\frac{w}{\mathrm{\lambda ̄}}\right)^2x_{\text{diff}}.$$
(12)
If the main condition Eq.(10) of applicability of the paraxial approximation is met, this requirement is distances large compared with the typical diffraction length $`x_{\text{diff}}`$.
### A Paraxial approximation: magnetic scattering
In this section the 2D paraxial theory is applied to the case of an external magnetic field; for simplicity $`U=0`$, generalization to $`U0`$ is straightforward.
The paraxial wave equation, the gauge covariant form of Eq.(6) with $`U=0`$, reads
$$\left(i\mathrm{}v_x+\frac{\mathrm{}^2}{2m}_y^2\right)\psi =0$$
(13)
where $`_{x,y}\frac{}{x,y}i\frac{e}{\mathrm{}c}A_{x,y}`$ $`𝑨`$ being the vector potential; the $`x`$axis is chosen in the direction of the propagation $`𝒗`$. The current density is given by Eq.(8) if modified by the standard diamagnetic term .
It is convenient to consider first the situation when the field is present only in a finite region. Divide the space into three regions: $`x<x_{in}`$, incoming (I); $`x_{in}<x<x_{out}`$ scattering (II); and outgoing region (III), $`x>x_{out}`$ (see Fig. 3). In regions I and II a magnetic field is absent. Present the wave function in I as
$$\mathrm{\Psi }(𝒓)=e^{i\frac{e}{\mathrm{}}\underset{𝑹_𝑰}{\overset{𝒓}{}}𝑑𝒓𝑨}\psi _{\text{in}}(𝒓),xx_{in},$$
(14)
and by this define $`\psi _{\text{in}}`$ . The integral in Eq.(14) does not depend on the path of the integration if the latter is in the magnetic free region I. The overall phase of $`\psi _{\text{in}}`$ depends on the choice of $`𝑹_I`$ and the gauge of $`𝑨`$. It is convenient to put $`𝑹_I`$ on the I-II interface, $`𝑹_I=(x=x_{in},y=y_{})`$, with somehow chosen $`y_{}`$.
The new function $`\psi _{\text{in}}`$ has the properties of the wave function in the gauge $`𝑨=0`$: It obeys the free equation,
$$\left(i\mathrm{}v\frac{}{x}+\frac{\mathrm{}^2}{2m}\frac{^2}{y^2}\right)\psi _{in}=0$$
(15)
and the probability current is given by Eq.(8) (without any diamagnetic term).
The wave function of the incoming beam is the input to the scattering problem: It is assumed that the problem of a free propagation in I is solved, and the incoming wave at the I-II interface, $`\psi _{\text{in}}(𝒓)|_{x=x_{in}}=\psi _{\text{in}}(y)`$, is known. The normalization condition
$`{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑y|\psi _{\text{in}}(y)|^2=1`$makes the total (conserving) flux in the $`x`$-direction equal to the velocity $`v`$. With the above choice of $`𝑹_I`$, the wave function Eq.(14) at the boundary of the region I reads
$$\mathrm{\Psi }(x,y)|_{x=x_{in}}=e^{i\frac{e}{\mathrm{}}\underset{y_{}}{\overset{y}{}}𝑑y_1A_y(x_{in},y_1)}\psi _{\text{in}}(y).$$
(16)
where the integration is performed along a piece of the straight line $`x=x_{in}`$.
The wave at $`x>x_{in}`$ has to be found from the paraxial equation Eq.(13), solved with Eq.(16) as the boundary condition. The solution in the both, scattering (II) and outgoing (III), regions may be generally written as,
$$\mathrm{\Psi }(x,y)=\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑y^{}G^R(x,y;x_{in},y^{})e^{i\frac{e}{\mathrm{}}\underset{y_{}}{\overset{y^{}}{}}𝑑y_1A_y(x_{in},y_1)}\psi _{\text{in}}(y^{}),x>x_{in}$$
(17)
where the Green function, $`G^R(𝒓,𝒓^{})`$ solves
$$\left(i\mathrm{}v_x+\frac{\mathrm{}^2}{2m}_y^2\right)G^R(𝒓,𝒓^{})=i\mathrm{}v\delta (𝒓𝒓^{}),G^R=0\text{ for }x<x^{}.$$
(18)
Similar to the above consideration of I (see Eq.(14)), one defines in III a new function, $`\psi _{\text{out}}`$, by
$$\mathrm{\Psi }(𝒓)=e^{i\frac{e}{\mathrm{}}\underset{𝑹_{𝑰𝑰𝑰}}{\overset{𝒓}{}}𝑑𝒓𝑨}\psi _{\text{out}}(𝒓),xx_{out};$$
(19)
again, the whole path of integration must be in the field free region III; choose the initial point of the integration path $`𝑹_{III}=(x_{out},y_{})`$ (or any other point at the $`IIIII`$ interface). Analogously to $`\psi _{\text{in}}`$ in I, $`\psi _{\text{out}}(𝒓)`$ obeys the free Eq.(15) and is fully defined in the whole region III by $`\psi _{\text{out}}(y)`$ that is the boundary value at the II-III interface, $`\psi _{\text{out}}(y)\psi (x_{out},y)`$. From Eq.(19) and (17), we see that
$$\psi _{\text{out}}(y)=\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑y^{}𝒢(x_{out},y;x_{in},y^{})\psi _{\text{in}}(y^{})$$
(20)
where the Green function $`𝒢`$ is
$$𝒢^R(x,y;x^{},y^{})e^{i\frac{e}{\mathrm{}}\underset{y_{}}{\overset{y}{}}𝑑y_2A_y(x,y_2)}G^R(x,y;x^{},y^{})e^{i\frac{e}{\mathrm{}}\underset{y_{}}{\overset{y^{}}{}}𝑑y_1A_y(x^{},y_1)}.$$
(21)
Eq.(20) combined with Eq.(21) relates the outgoing wave amplitude $`\psi _{\text{out}}`$ to the incoming wave $`\psi _{\text{in}}`$, and thus gives a general solution to the magnetic scattering problem in the paraxial approximation.
Known $`\psi _{\text{out}}`$, one finds $`\psi _{\text{out}}(x,y)`$ in the outgoing region $`III`$ as
$$\psi _{\text{out}}(x,y)=\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑y^{}G_0^R(xx_{out},yy^{})\psi _{\text{out}}(y^{}),xx_{out}$$
(22)
with
$$G_0^R(x,y)=\theta (x)\frac{1}{\sqrt{2\pi i\mathrm{\lambda ̄}x}}e^{\frac{i}{2\mathrm{\lambda ̄}x}y^2},$$
(23)
being the free propagator .
### B The eikonal approximation
If the scattering region II is narrow enough compared with the diffraction length Eq.(11), one can neglect the transverse derivatives in Eq.(13). This limit corresponds to the well-known eikonal approximation . The wave function obeys the eikonal equation ($`U=0`$)
$$i\mathrm{}_x\psi =0.$$
(24)
Solution to Eq.(18) reads in the eikonal limit
$$G^R(x,y;x^{},y^{})=e^{i\frac{e}{\mathrm{}}\underset{x^{}}{\overset{x}{}}𝑑x^{\prime \prime }A_x(x^{\prime \prime },y)}\delta (yy^{})\theta (xx^{}).$$
(25)
Substituting Eq.(25) into Eqs.(21) and (20), one gets the eikonal solution to the magnetic scattering problem,
$$\psi _{\text{out}}(y)=e^{i\frac{e}{\mathrm{}c}\mathrm{\Phi }(y)}\psi _{\text{in}}(y)$$
(26)
where $`\mathrm{\Phi }(y)`$ is the total magnetic flux through the directed area (see Fig. 4) restricted by the path $`(x_{in},y_{})(x_{in},y)(x_{out},y))(x_{out},y_{})`$; choice of $`y_{}`$ is arbitrary affecting only the overall phase of $`\psi _{\text{out}}`$ .
## III Illustrations: Aharonov-Bohm scattering
To illustrate the usage of the paraxial theory, the scattering on the Aharonov-Bohm magnetic line is considered in this section as an example. Some results presented in this section has been published in a short communication Ref..
The scattering problem is formulated as follows. The incident wave $`\psi _{in}=\psi _{\text{in}}(y)`$ comes from the negative side of the $`x`$-axis and the charge sees a magnetic line (extending in the $`z`$-direction), which creates the magnetic field $`b(x,y)=\mathrm{\Phi }_0\delta (x)\delta (y)`$. The final goal is to characterize the outgoing wave at $`x>0`$.
As the magnetic field is finite only at the origin, the scattering region II in the terminology of Section II A can be shrunk to just the line $`y=0`$, so that $`x_{in}=0`$ and $`x_{out}=+0`$. Seeing that scattering region is narrow, the eikonal approximation discussed in Section II B is applicable, except, perhaps, in the immediate vicinity of the singular magnetic line point $`y=0`$. The flux function, $`\mathrm{\Phi }(y)`$, in the eikonal expression Eq.(26) is easily found to be $`\mathrm{\Phi }(y)=\mathrm{\Phi }\theta (y)`$ (if $`y_{}`$ is chosen at $`\mathrm{}`$) . The outcoming wave reads
$$\psi _{\text{out}}(y)=\mathrm{exp}\left(i\pi \stackrel{~}{\mathrm{\Phi }}sign(y)\right)\psi _{\text{in}}(y).$$
(27)
where $`\stackrel{~}{\mathrm{\Phi }}\frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}`$ . Same result can readily obtained by solving the first order differential equation Eq.(24) with the vector potential in the gauge
$`A_x(x,y)={\displaystyle \frac{1}{2}}\mathrm{\Phi }sign(y)\delta (x),A_y(x,y)=0.`$
The freely propagating outgoing wave is found from Eq.(22) with $`x_{out}=0`$ and $`\psi _{\text{out}}`$ from Eq.(27) Ref.,
$$\psi _{\text{out}}(x,y)=\psi _{\text{in}}(x,y)\mathrm{cos}\pi \stackrel{~}{\mathrm{\Phi }}+iV(x,y)\mathrm{sin}\pi \stackrel{~}{\mathrm{\Phi }},x>0,$$
(28)
where $`\psi _{\text{in}}(x,y)`$ is the incoming wave continued to the region $`x>0`$,
$$\psi _{\text{in}}(x,y)=\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑y^{}G_0^R(x,yy^{})\psi _{\text{in}}(y^{}),$$
(29)
i.e. the wave in the absence of the Aharonov-Bohm line, and
$$V(x,y)=\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑y^{}G_0^R(x,yy^{})\widehat{y^{}}\psi _{\text{in}}(y^{}).$$
(30)
Eqs.(28 30) give the solution of the Aharonov-Bohm problem in the paraxial approximation for an arbitrary incoming wave $`\psi _{in}(y)`$.
Using Eq.(23), one can check that
$`{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑y|\psi _{\text{in}}(x,y)|^2=1`$ , $`{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑y|V(x,y)|^2=1,`$ (31)
$`\mathrm{}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}dy\psi _{\text{in}}^{}(x,y`$ $`)`$ $`V(x,y)=0;`$ (32)
for any $`x>0`$ and arbitrary (normalized) $`\psi _{\text{in}}`$; these relations guarantee the current conservation .
The solution in Eq.(28) gives a convenient tool for studying interaction of beams (waves packets) with the Aharonov-Bohm flux line. Some examples are considered below.
### A Plane incident wave
Plane wave corresponds to $`\psi _{\text{in}}(y)=1`$. From Eqs.(29) and (23), one gets $`\psi _{\text{in}}(x,y)=1`$, and from Eq.(30) Ref.,
$$V(x,y)=K(\stackrel{~}{y}),\stackrel{~}{y}=\frac{y}{\sqrt{2\mathrm{\lambda ̄}x}}$$
(33)
where $`K(t)`$ is a familiar function,
$$K(t)=\frac{2}{\sqrt{i\pi }}\underset{0}{\overset{t}{}}𝑑t^{}e^{it^2},$$
(34)
related to the Fresnel integrals. The outgoing wave Eq.(28) reads
$$\psi _{\text{out}}(x,y)=\mathrm{cos}\pi \stackrel{~}{\mathrm{\Phi }}+iK(\stackrel{~}{y})\mathrm{sin}\pi \stackrel{~}{\mathrm{\Phi }}.$$
(35)
The complex wave function $`\psi _{\text{out}}(x,y)`$ is a function of the real parameter $`\stackrel{~}{y}`$ only. This means, that $`\psi _{out}(x,y)`$ for any point of the half-plane $`x>0`$ spans a line on the complex plane $`\mathrm{}\psi \mathrm{}\psi `$. This line (Fig. 5) is the Cornu spiral well-known in the diffraction theory . The property of the Cornu spiral is that $`(dt)^2`$ is proportional to $`(d\mathrm{})^2`$ where $`d\mathrm{}`$ is the distance between the points corresponding to $`t`$ and $`t+dt`$. From Eq.(33) and Eq.(34), the element $`d\mathrm{}`$ of the length along the spiral (the arc-length) and $`d\stackrel{~}{y}`$ are related as $`(d\mathrm{})^2=\frac{4}{\pi }\mathrm{sin}^2(\pi \stackrel{~}{\mathrm{\Phi }})(d\stackrel{~}{y})^2`$. Since the wave function Eq.(35) is real for $`\stackrel{~}{y}=0`$, the arc-length counted from the spiral point with $`\mathrm{}\psi =0`$ is proportional to $`\stackrel{~}{y}`$
The singular Aharonov-Bohm line problem gives a rather tough test for the validity of the paraxial approximation. The evaluation of its accuracy has been made simple by the recent observation Ref. that the exact solution can be presented in the form Eq.(35) with $`\stackrel{~}{y}\stackrel{~}{Y}`$ with $`\stackrel{~}{Y}`$ found from
$$\stackrel{~}{Y}^2=\frac{(rx)}{\mathrm{\lambda ̄}}+\phi \left(\frac{1}{2}\stackrel{~}{\mathrm{\Phi }}\right)+\mathrm{},$$
(36)
$`r`$ and $`\phi `$ being the cylindrical coordinates; the terms not shown in Eq.(36) are of order $`𝒪\left(\frac{\mathrm{\lambda ̄}}{r}\right)`$ . (Unexpectedly, the Aharonov-Bohm wave function on the $`xy`$ plane can be mapped on the Cornu spiral not only in the forward direction but for any scattering angle (and $`r\mathrm{\lambda ̄}`$), see Ref. for details.) The paraxial approximation gives correctly the leading term $`\sqrt{\frac{rx}{\mathrm{\lambda ̄}}}\stackrel{~}{y}`$ in the vicinity of the forward direction and short wave length $`\mathrm{\lambda ̄}r`$.
### B Finite size beam
The scattering of a plane wave by the Aharonov-Bohm line is highly singular: Eq.(35) shows that in the forward direction, $`\phi \sqrt{\frac{\mathrm{\lambda ̄}}{x}}`$ i.e. $`\stackrel{~}{y}1`$, the wave function equals to $`\mathrm{cos}\pi \stackrel{~}{\mathrm{\Phi }}`$ and does not converge at large distances to the plane incident wave as assumed in the standard scattering theory. Known from the exact solution , this anomaly is the reason why the text-book scattering theory fails: The incoming plane wave and scattered wave cannot be separated and the scattering amplitude cannot be introduced as the object carrying the complete information about scattering .
The singular behaviour is due to the combination of two factors: (i) the infinitely long range of the interaction with the line; (ii) the infinite extension of the plane wave. Since any physical state has a finite transverse extension $`W`$, the behaviour of the potential beyond the width $`W`$ is irrelevant, and the singularities are expected to be regularized.
To show this, suppose that the incoming wave is beam-like with the profile
$`\psi _{in}(y)=e^{\frac{y}{|W|}},`$where $`W`$ is the beam width; it is assumed that $`W\mathrm{\lambda ̄}`$ so that the paraxial approximation is applicable. The beam has a small but finite angular width $`\phi _0=\mathrm{\lambda ̄}/W`$. One can easily see from Eq.(28) that at large distances $`xW^2/\mathrm{\lambda ̄}`$, the outgoing wave behaves like a spherical wave i.e. $`|\psi _{out}(x,y)|^2P(\phi )/x`$ where $`\phi =y/x`$ is the scattering angle and $`P(\phi )`$,
$$P(\phi )=\frac{2\mathrm{\lambda ̄}}{\pi }\frac{(\phi \mathrm{sin}\pi \stackrel{~}{\mathrm{\Phi }}\phi _0\mathrm{cos}\pi \stackrel{~}{\mathrm{\Phi }})^2}{(\phi ^2+\phi _0^2)^2},$$
(37)
has the meaning of the angular distribution of the intensity in the outgoing wave. As expected, the distribution shown in Fig. 6 is perfectly smooth.
For the angles larger than the beam angular width, i.e. $`|\phi |\phi _0`$, one gets from Eq.(37) that
$$P(\phi )\frac{2\mathrm{\lambda ̄}}{\pi }\frac{\mathrm{sin}^2\pi \stackrel{~}{\mathrm{\Phi }}}{\phi ^2},$$
(38)
recovering the small angle asymptotics of the Aharonov-Bohm scattering cross-section
$$\left(\frac{d\sigma }{d\phi }\right)_{\text{AB}}=\frac{\mathrm{\lambda ̄}}{2\pi }\frac{\mathrm{sin}^2\pi \stackrel{~}{\mathrm{\Phi }}}{\mathrm{sin}^2\frac{\varphi }{2}}.$$
(39)
It would be rather trivial if the angular broadening of the incident wave led just to a regularization of the forward scattering singularity. Most important is that a qualitatively new feature becomes seen: Unlike the Aharonov-Bohm cross-section Eq.(39) the angular distribution in Eq.(37) is left-right asymmetric. The antisymmetric part is concentrated in the forward direction $`|\phi |\phi _0`$, i.e. within the angular width of the incoming wave i.e. in the region where the scattered and incident waves cannot be separated . The asymmetry means that the beam is deflected by the Aharonov-Bohm line as a whole . The order of magnitude of the deflection is the initial angular width $`\phi _0\mathrm{\lambda ̄}/W`$. One may say that the Aharonov-Bohm line not only scatters the incident wave but also modifies the unscattered part of the wave. Later, it will be shown that in a system of many Aharonov-Bohm lines, the deflections by individual lines coherently add together, and the lines act as an effective magnetic field
The deflection of the beam $`\mathrm{\Delta }\phi `$ can be presented via the momentum $`\mathrm{\Delta }p_{}=\frac{\mathrm{}}{\mathrm{\lambda ̄}}\mathrm{\Delta }\phi `$ transfered to the charge in the direction perpendicular to its initial velocity. A calculation details of which are collected in Appendix A gives the following result Ref.:
$$\mathrm{\Delta }p_{}=\mathrm{}|\psi _{\text{in}}(0)|^2\mathrm{sin}2\pi \frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}$$
(40)
here the momentum transfer is expressed via the value of the normalized incoming wave at the position of the line. By comparison with the exact theory, the validity of this paraxial result has been confirmed by Berry .
## IV Density matrix
The theory of magnetic scattering presented in Section II A is gauge invariant only in a limited sense. Although Eq.(20) holds in arbitrary gauge, each of the objects there, $`\psi _{\text{in}}`$, $`\psi _{\text{out}}`$, and $`𝒢`$ is gauge dependent, although only through the overall phase. For example, under $`𝑨𝑨+\mathbf{}\chi `$ simultaneously with $`\mathrm{\Psi }e^{i\frac{e}{\mathrm{}c}\chi }`$, the incoming wave $`\psi _{\text{in}}`$ Eq.(14) is modified as $`\psi _{\text{in}}e^{i\frac{e}{\mathrm{}}\chi (𝑹_I)}\psi _{\text{in}}`$; we see also that the overall phase depends on the arbitrarily chosen $`𝑹_I`$. Of course, the observables are independent from the global phase and the above gauge dependence does not create any problem.
A truly gauge invariant theory can be formulated in terms of the by-linear in $`\mathrm{\Psi }`$ and $`\mathrm{\Psi }^{}`$ “density matrix” $`\rho (y_1,y_2;x)`$ defined as
$$\rho (y_1,y_2;x)e^{i\frac{e}{\mathrm{}}\underset{y_2}{\overset{y_1}{}}𝑑y^{}A_y(x,y^{})}\mathrm{\Psi }(x,y_1)\mathrm{\Psi }^{}(x,y_2).$$
(41)
The “density matrix” $`\rho `$ carries the full quantum information needed to find observables and is gauge invariant. In the field free regions, the density matrix is built from $`\psi _{\text{in}}`$ and $`\psi _{\text{in}}^{}`$ or $`\psi _{\text{out}}`$ and $`\psi _{\text{out}}^{}`$; these combinations depend on neither gauge nor $`𝑹_{I,II}`$.
The current density the $`x`$\- and $`y`$\- directions are (cf. Eq.(8))
$`J_x(x,y)=v\rho (y,y;x),J_y(x,y)={\displaystyle \frac{\mathrm{}}{2im}}\left({\displaystyle \frac{}{y_1}}{\displaystyle \frac{}{y_2}}\right)\rho (y_1,y_2;x)\text{ }_{y_1=y_2=y}`$
Seeing that the evolution of the wave function $`\mathrm{\Psi }`$ in Eq.(41) is given by the propagator $`G^R`$ Eq.(18), the density matrix evolves from $`x^{}`$ to $`x`$ ($`x>x^{}`$) as
$$\rho (y_1,y_2;x)=\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑y_1^{}𝑑y_2^{}𝒢^R(y_1,x;y_1^{},x^{})\rho (y_1^{},y_2^{};x^{})𝒢^A(y_2^{},x^{};y_2,x)$$
(42)
where $`𝒢^R`$ is defined by (21) and the advanced Green function $`𝒢^A`$
$`𝒢^A(𝒓_1,𝒓_2)=\left(𝒢^R(𝒓_2,𝒓_1)\right)^{}.`$
Introducing the two-particle Green function
$$𝒦(y_1,y_2;x|y_1^{},y_2;,x^{})=𝒢^R(y_1,x|y_1^{},x^{})𝒢^A(y_2^{},x^{}|y_2,x),$$
(43)
Eq.(42) can be written as
$$\rho (y_1,y_2;x)=\underset{\mathrm{}}{\overset{\mathrm{}}{}}dy_1^{}dy_2^{}𝒦(y_1,y_2;x|y_1^{},y_2;,x^{})\rho (y_1^{},y_2^{};x^{})$$
(44)
As before, the incoming wave enters the scattering problem as the boundary condition at $`x=x_{in}`$: $`\rho (y_1,y_2;x_{in})=\rho _{in}(y_1,y_2)`$. Further propagation of the incoming beam is given by Eq.(44) with $`x^{}=x_{in}`$.
For future references, we note the following property of the Green function:
$`{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑yG^R(x,y;x^{},y_1)G^A(x^{},y_2;x,y)`$ $`=`$ $`\theta (xx^{})\delta (y_1y_2),`$ (45)
$`{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑yG^R(x,y_1;x^{},y)G^A(x^{},y;x,y_2)`$ $`=`$ $`\theta (xx^{})\delta (y_1y_2),`$ (46)
this relations expresses the current conservation .
In particular, from Eq.(45), one gets the conservation of the total current in the beam:
$`{\displaystyle 𝑑y\rho (y,y;x)}={\displaystyle 𝑑y^{}\rho (y^{},y^{};x=0)}`$
### A Path integral representation
Similar to Eq.(7), the stationary paraxial equation in 2D may be mapped to a time dependent 1D problem:
$$i\mathrm{}\frac{}{\tau }\psi =\left(\frac{\mathrm{}^2}{2m}\left(\frac{}{y}i\frac{e}{\mathrm{}c}a\right)^2+e\phi \right)\psi =0$$
(47)
where “time” $`\tau =\frac{x}{v}`$, $`a=A_y`$, and $`\phi =\frac{v}{c}A_x`$. In the effective 1D problem, the particle moves in the “electric field”,
$$F=\frac{1}{c}\dot{a}\phi ,$$
(48)
defined by the “vector potential”, $`a`$, and the “scalar potential”, $`\phi `$. The gauge transformation, $`𝑨𝑨+\mathbf{}\chi `$, translates to $`aa+\chi `$, $`\phi \phi \frac{1}{c}\dot{\chi }`$, so that the effective electric field Eq.(48) is indeed gauge invariant. Locally, the electric field is related to the magnetic field of the original problem $`b(x,y)`$ as $`F=vb`$.
The mapping to the effective 1D problem allows one to use a convenient path integral representation for the paraxial propagators. Obviously, $`G^R`$ is just the retarded Green function for nonstationary equation Eq.(47) and in the Feynman path integral representation
$`G^R(x,y;x^{},y^{})={\displaystyle \underset{y^{}}{\overset{y}{}}}𝒟[y(x)]e^{iS[y(x)]}`$where the action of the effective 1D problem $`S=\frac{1}{\mathrm{}}𝑑\tau (m\dot{y}^2/2+ea\dot{y}/ce\phi )`$ translates as
$$S[y(x)]=\frac{1}{2\mathrm{\lambda ̄}}\underset{x^{}}{\overset{x}{}}𝑑xy_x^2+\frac{e}{\mathrm{}c}\underset{y=y(x)}{}𝑑𝒓\mathbf{}𝑨,$$
(49)
where $`y_x=\frac{dy}{dx}`$. The action corresponding to $`𝒢^R`$ differs from Eq.(49) in the path of integration which should be extended in an obvious way to include the additional exponential factors in the definition of $`𝒢^R`$ Eq.(21).
With the help of Eq.(49), the two particle Green function Eq.(43) can be presented as
$$𝒦(y_1,y_2;x|y_1^{},y_2;,x^{})=𝒟[y_1(x)]𝒟[y_2(x)]e^{i𝒮[y_1(x),y_2(x)]},$$
(50)
where
$`𝒮=𝒮_0+{\displaystyle \frac{2\pi }{\mathrm{\Phi }_0}}\mathrm{\Phi }([y_1],[y_2]),`$$`𝒮_0`$ being the free motion contribution,
$$𝒮_0[y_1(x),y_2(x)]=\frac{1}{2\mathrm{\lambda ̄}}\underset{x^{}}{\overset{x}{}}𝑑x\left(y_{1x}^2y_{2x}^2\right),$$
(51)
and $`\mathrm{\Phi }`$ is the flux
$$\mathrm{\Phi }([y_1],[y_2])=\underset{C([y_1],[y_2])}{}𝑑𝒓\mathbf{}𝑨$$
(52)
threading the (oriented) area bounded by the paths $`y_1(x)`$ and $`y_2(x)`$ and the vertical lines at $`x`$ and $`x^{}`$ (see Fig.7).
The path integral representation is used below to perform averaging with respect to the gauge field.
## V Gaussian Random magnetic field
This Section concerns the averaging the two-particle Green’s function Eq.(43) with respect to the Gaussian random magnetic field Eq.(1). In the path integral representation Eq.(50), the random field enters via the flux $`\mathrm{\Phi }([y_1],[y_2])`$ Eq.(52). For the Gaussian field Eq.(1), the averaged value
$`\mathrm{exp}\left(2\pi i{\displaystyle \frac{\mathrm{\Phi }([y_1],[y_2])}{\mathrm{\Phi }_0}}\right)=\mathrm{exp}\left({\displaystyle \frac{𝒜_{\text{no}}([y_1y_2])}{2^2}}\right)`$where $`𝒜_{\text{no}}`$
$`𝒜_{\text{no}}([y_1y_2])={\displaystyle \underset{x^{}}{\overset{x}{}}}𝑑x|y_1y_2|`$is the non-oriented area bounded by the paths $`y_1(x)`$ and $`y_2(x)`$ and lines $`x=x^{}`$ and $`x=x`$ (see Fig.7). Further calculations are rather simple thanks to the fact that $`e^{i2\pi \frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}}`$ is a functional only of $`y_1y_2`$. This important simplification is a property of the models with a $`\delta `$-correlated magnetic field.
In the variables
$`y=y_1y_2,Y={\displaystyle \frac{1}{2}}\left(y_1+y_2\right),`$the kinetic energy contribution $`𝒮_0`$ Eq.(51) reads after integration by parts,
$`𝒮_0[y_1(x),y_2(x)]={\displaystyle \frac{1}{2\mathrm{\lambda ̄}}}y_xY|_x^{}^x{\displaystyle \frac{1}{2\mathrm{\lambda ̄}}}{\displaystyle \underset{x^{}}{\overset{x}{}}}𝑑xy_{xx}Y.`$Since $`Y(x)`$ enters only the $`𝒮_0`$, the integration $`e^{i𝒮}`$ with respect to $`Y`$ gives $`\delta (y_{xx})`$. This means, that the integration with respect to $`y(x)`$ is limited to the path with $`y_{xx}\frac{d^2y}{dx^2}=0`$ that is the straight line connecting the initial and final points. After this, the integral is easily calculated.
Finally, $`𝒦_{\text{av}}`$ that is the paraxial two-particle Green function Eq.(43) averaged with respect to the fluctuation magnetic field, reads
$$𝒦_{\text{av}}(y_1,y_2;x|y_1^{},y_2;,x^{})=𝒦_\text{0}(y_1,y_2;x|y_1^{},y_2;,x^{})\mathrm{exp}(\frac{𝒜_{\text{no}}}{2^2}),$$
(53)
here $`𝒦_0`$ is the free two-particle propagator, and $`𝒜`$ is the (non-oriented) area formed by the straight line trajectories Fig.8,
$`𝒜`$ $`=`$ $`{\displaystyle \frac{1}{2}}(xx^{})\times \{\begin{array}{ccc}|y_1+y_1^{}y_2y_2^{}|\hfill & ,& \hfill \text{if }(y_1y_2)(y_1^{}y_2^{})>0\text{}\\ \frac{(y_1y_2)^2+(y_1^{}y_2^{})^2}{y_1y_2+y_1y_2}\hfill & ,& \hfill \text{if }(y_1y_2)(y_1^{}y_2^{})<0\end{array}`$ (56)
In the variables $`y,Y`$ and $`(xx^{})x`$, an expanded version of Eq.(53) reads
$`K_{\text{av}}(y,Y;y^{},Y^{};x)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi x\mathrm{\lambda ̄}}}\mathrm{exp}\left[{\displaystyle \frac{i}{x\mathrm{\lambda ̄}}}(YY^{})(yy^{}){\displaystyle \frac{1}{8}}{\displaystyle \frac{x}{^2}}\left(y+y^{}+{\displaystyle \frac{(y+y^{})^2}{y+y^{}}}\right)\right]`$ (57)
The two-particle Green function in Eq.(57) allows one to find averaged over the random field “evolution” of the density matrix and thus describes correlations in the gauge invariant observables like the density or the current. For instance, it describes transmission through a slab with a random magnetic field. In what follows, Eq.(57) is applied to some simple cases: (i) focusing of a coherent wave; (ii) the scattering of the partially coherent spatially uniform incoming wave.
### A Coherent propagation: Focusing
Let the incident wave $`\psi _{in}(y)`$ be a converging Gaussian beam i.e.
$$\psi _{in}(y)=\frac{1}{(\pi )^{1/4}\sqrt{w}}e^{\frac{y^2}{2}(\frac{1}{w^2}+\frac{i}{\mathrm{\lambda ̄}f})},$$
(58)
where $`f`$ is the distance to the focal point and $`wf`$ is the width of the beam. It is well-known that the “spherical” wave like that in Eq.(58) will converge at the focal point $`x=f`$ producing a diffraction limited spot with the waist $`\mathrm{\lambda ̄}/\theta `$ where $`\theta =w/f1`$ is the angular size of the beam as seen from the focal point.
The distribution of the averaged intensity for the beam propagating in a random magnetic field can be found with the help of Eq.(57). Taking for simplicity only the points on the beam axis $`y=0`$, the intensity $`I(x)\rho (0,0;x)`$, reads
$`I(x)={\displaystyle \frac{1}{\sqrt{\pi }w}}{\displaystyle \frac{1}{\lambda x}}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑Y𝑑ye^{i\frac{1}{\mathrm{\lambda ̄}}\left(\frac{1}{x}\frac{1}{f}\right)yY\frac{Y^2}{w^2}\frac{y^2}{4w^2}\frac{x}{4^2}y}.`$At the focal point $`x=f`$,
$`I(f)={\displaystyle \frac{1}{\lambda f}}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑ye^{\frac{y^2}{4w^2}\frac{f}{4^2}y},`$In the limiting cases,
$$I(f)=\{\begin{array}{ccc}\frac{2\sqrt{\pi }}{\lambda f}\times w\hfill & ,& fw^2\hfill \\ \frac{8}{\lambda f}\times \frac{^2}{f}\hfill & ,& fw^2\hfill \end{array}.$$
(59)
For the conditions when the magnetic field is not important, the upper line gives the usual diffraction limited value of the intensity; the width of the spot at the focal plane (line) is of order of $`1/I(f)\mathrm{\lambda ̄}/\theta `$, where $`\theta =w/f`$. The larger the aperture $`w`$, the larger the intensity in the focus and the smaller the size of the spot. However, in the presence of the random magnetic field, the intensity saturates when the aperture $`w^2/f`$ and $`\theta \left(/f\right)^2`$.
This behaviour is very different from that when the scattering is due to a random scalar potential with a short correlation length. In this case, the relevant parameter characterizing the disorder is the focal length $`f`$ in units of the mean free path $`l`$ rather then the size of the aperture: On the background created by incoherent scattering, one would see a spot with the disorder insensitive profile and the integral intensity $`e^{f/l}`$ (the exponential factor is probability that the wave does not experience any scattering).
### B Incoherent wave: the Boltzmann equation
Consider the initial density matrix of the form $`\rho _{in}(y_1,y_2)=\rho _0(y_1y_2)`$, $`y=y_1y_2`$, which corresponds to a partially coherent spatially homogeneous state. At $`x>0`$, the density matrix $`\rho (y_1,y_2;x)=\rho (y;x)`$ is found from Eq.(44) with the two-particle propagator from Eq.(57). The Fourier transform,
$$\rho (y;x)=\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑\phi e^{i\phi \frac{y}{\mathrm{\lambda ̄}}}n_\phi (x),$$
(60)
defines $`n_\phi `$ which has the meaning of the distribution function with respect of the transverse momentum, $`q=\phi \frac{\mathrm{}}{\mathrm{\lambda ̄}}`$, and $`\phi `$ is the angle between the velocity of the particle and the $`x`$-axis. In the paraxial situation, the distribution is concentrated at small angles and the integration in Eq.(60) can be taken in infinite limits.
From, Eqs.(44) and (57) we easily get the density matrix of the wave having traveled the distance $`x`$:
$$\rho (y;x)=\rho _0(y)\mathrm{exp}\left(\frac{xy}{2^2}\right).$$
(61)
As required by the current conservation, $`\rho (0;x)`$ does not depend on $`x`$ whereas the non-diagonal elements of the density matrix $`y0`$ decay to zero, the faster, the more “distance to the diagonal” $`|y|`$. In other words, random magnetic field is very effective in destroying a long range coherence, the longer the coherence, the faster it decays.
Considering the limiting case of plane infinite incident wave $`n_\phi (x=0)=\delta (\phi )`$ i.e. $`\rho _{in}(y)=1`$, Eq.(61) gives
$$n_\phi (x)=\frac{1}{\pi }\frac{\mathrm{\Delta }}{\phi ^2+\mathrm{\Delta }^2},\mathrm{\Delta }=x\frac{\mathrm{\lambda ̄}}{2^2}$$
(62)
The evolution is nonperturbative in the sense that the plane wave looses its shape immediately at any $`x0`$ transforming into the Lorentz distribution with the width of the distribution proportional to $`x`$ and the strength of the field . In particular, it means that the plane wave is not a good basis for the perturbation theory.
More insight can be gained if the evolution of the density matrix is mapped to a Boltzmann-type kinetic equation. For this, note that the density matrix in Eq.(61) satisfies the equation,
$$v\frac{}{x}\rho +\widehat{I}\rho =0,\widehat{I}=\frac{v}{2^2}|y|.$$
(63)
Written for the distribution function $`n_\phi `$ introduced in Eq.(60), the equation acquires the familiar Boltzmann form,
$$v\frac{n_\phi }{x}+\widehat{I}n_\phi =0,$$
(64)
where $`\widehat{I}`$ is the collision integral i.e. operator $`\widehat{I}`$ in the $`\phi `$-representation.
One may present the collision integral in the standard form
$$\widehat{I}n_\phi =𝑑\varphi w(\varphi )(n_\phi n_{\phi +\varphi }),$$
(65)
where $`w(\varphi )`$ is the scattering rate for the process $`\phi \phi +\varphi `$. From the condition that the operators in Eq.(64) and Eq.(65) have same eigenvalues corresponding to the common eigenfunctions, $`n_\phi e^{i\frac{y_0}{\mathrm{\lambda ̄}}\phi }`$, $`w`$ must satisfy the requirement that Ref.
$$\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑\varphi w(\varphi )(1e^{i\frac{y_0}{\mathrm{\lambda ̄}}\varphi })=\frac{v}{2^2}|y_0|.$$
(66)
From here,
$$w(\varphi )=\frac{2}{\pi \tau _0}\frac{1}{\varphi ^2},\frac{1}{\tau _0}=\frac{v\mathrm{\lambda ̄}}{4^2}$$
(67)
Usually, one can split the collision integral into the in- and out-scattering pieces. In the case of random magnetic field, the scattering-out rate is ill-defined as $`𝑑\varphi w(\varphi )`$ diverges at small angles, and the split hardly makes sense. On the other hand, the collision integral, as an operator acting on the distribution function, is well defined and the transport is not singular.
Treating the random magnetic field in the Born approximation, Aronov et al. found the the scattering rate $`W(\varphi )`$ to be
$$W(\varphi )=\frac{1}{2\pi \tau _0}\mathrm{cot}^2\frac{\varphi }{2}.$$
(68)
Because of divergence at small angles, one may doubt the validity of the Born approximation. However, the small $`\varphi `$ asymptotics of $`W(\varphi )`$ agrees with Eq.(67). This means that Eq.(68) is actually valid for arbitrary $`\varphi `$ if used for constructing the collision integral. The other way around, the collision integral Eq.(65) is expected to correctly describe scattering with arbitrary scattering angles if $`W(\varphi )`$ Eq.(68) is used instead of paraxial $`w(\varphi )`$ in Eq.(67). This allows one to generalize the paraxial kinetic equation, including large angle scattering.
The kinetic equation for the distribution function $`n_\phi `$ reads
$$𝒗\mathbf{}\mathbf{}n_\phi +\widehat{I}n_\phi =0,\mathrm{\hspace{0.33em}\hspace{0.33em}0}<\phi <2\pi ,$$
(69)
where the collision integral
$`\widehat{I}n_\phi ={\displaystyle \frac{1}{\tau _0}}{\displaystyle \underset{0}{\overset{2\pi }{}}}{\displaystyle \frac{d\varphi }{2\pi }}\mathrm{cot}^2{\displaystyle \frac{\varphi }{2}}\left(n_\phi n_{\phi +\varphi }\right).`$As before, the spilt of the collision integral in the in- and out- scattering parts leads to divergences and, therefore, has a very limited sense. At the same time, the collision operator is well-defined: If $`n_\phi `$ is presented as the sum,
$`n_\phi ={\displaystyle \underset{m}{}}n_me^{im\phi },`$over the eigenfunctions of the collision operator $`e^{im\phi }`$, $`m=0,\pm 1,\mathrm{}`$, the collision operator acts as
$`\widehat{I}n_\phi ={\displaystyle \frac{1}{\tau _0}}{\displaystyle \underset{m0}{\overset{\mathrm{}}{}}}(2|m|1)n_me^{im\phi },`$The parameter $`\tau _0`$ has the meaning of the relaxation time for the first harmonics $`m=\pm 1`$ i.e. the transport relaxation time.
One concludes that transport of a charge in a random magnetic field can be described by the Boltzmann equation, and is not anomalous in spite of the fact that the total scattering rate is infinite.
## VI Random array of Aharonov-Bohm lines
This section deals with the model of a random gauge field where the gauge field is created by an array of Aharonov-Bohm lines. It is assumed that the lines in the array take random space positions and the flux of a line $`\mathrm{\Phi }`$ may be random. The model is specified by the averaged density of the lines $`d_{_{\text{AB}}}`$ and the probability distribution $`p(\mathrm{\Phi })`$ for the magnetic flux $`\mathrm{\Phi }`$ in a line. Our primarily goal is to average the paraxial two-particle Green function Eq.(43) over the distribution of the lines.
The calculations turn out to be very similar to those in Section V so that we only outline them. Repeating the arguments from Section V one comes to an expression similar to Eq.(53):
$`𝒦_{\text{av}}(y_1,y_2;x|y_1^{},y_2;,x^{})=𝒦_\text{0}(y_1,y_2;x|y_1^{},y_2;,x^{})\mathrm{exp}\left(2\pi i{\displaystyle \frac{\mathrm{\Phi }^{(t)}(y_1,y_2;x|y_1^{},y_2;,x^{})}{\mathrm{\Phi }_0}}\right)`$where $`\mathrm{\Phi }^{(t)}(y_1,y_2;x|y_1^{},y_2;,x^{})`$ is the flux through the oriented area bounded by the straight (directed) lines connecting the initial and finite points (see Fig.8). Given the configuration of the Aharonov-Bohm array, the flux through the area is the sum over the lines piercing the area. The $`k`$’s line with the flux $`\mathrm{\Phi }_k`$ contributes to the total flux as $`\sigma _k\mathrm{\Phi }_k`$ where $`\sigma _k=+1`$ or $`1`$ depending on the orientation, positive or negative, of the area the line is situated in. Let $`N_+`$ ($`N_{}`$) be the (random) number of lines in the area with positive (negative) orientation; The variables $`\mathrm{\Phi }_k`$’s are independent in the model, and the averaging $`\mathrm{exp}\left(2\pi i\frac{\mathrm{\Phi }^{(t)}}{\mathrm{\Phi }_0}\right)`$ over the configurations with fixed $`N_\pm `$ is simple:
$`\mathrm{exp}\left(i{\displaystyle \frac{2\pi }{\mathrm{\Phi }_0}}{\displaystyle \underset{k=1}{\overset{N_++N_{}}{}}}\sigma _k\mathrm{\Phi }_k\right)=\mathrm{exp}\left(2\pi i{\displaystyle \frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}}\right)^{N_+}\mathrm{exp}\left(2\pi i{\displaystyle \frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}}\right)^N_{}`$where $`<\mathrm{exp}(\pm \frac{2\pi i\mathrm{\Phi }}{\mathrm{\Phi }_0})>`$ implies averaging with the distribution function $`p(\mathrm{\Phi })`$.
The random numbers $`N_\pm `$ obey the Poisson distribution $`P_N=e^{\overline{N}}\overline{N}^N/N!`$ with $`\overline{N}`$ either $`\overline{(N_++N_{})}=d_{_{\text{AB}}}𝒜_{\text{no}}`$ or $`\overline{N_+N_{}}=d_{_{\text{AB}}}𝒜_\text{o}`$, $`𝒜_{\text{no}}`$ and $`𝒜_\text{o}`$ being non-oriented and oriented area, respectively.
Finally,
$`\mathrm{exp}\left(2\pi i{\displaystyle \frac{\mathrm{\Phi }(y_1,y_2;x|y_1^{},y_2;,x^{})}{\mathrm{\Phi }_0}}\right)=\mathrm{exp}\left({\displaystyle \frac{𝒜_{\text{no}}}{2_{AB}^2}}+i{\displaystyle \frac{2\pi }{\mathrm{\Phi }_0}}\stackrel{~}{B}𝒜_\text{o}\right),`$where $`𝒜_{\text{no}}`$ is the non-oriented area Eq.(56), $`𝒜_\text{o}`$ is the oriented area,
$`𝒜_\text{o}={\displaystyle \frac{1}{2}}(xx^{})(y_1+y_1^{}y_2y_2^{}),`$
$$\frac{1}{_{AB}^2}=2d_{_{\text{AB}}}1\mathrm{cos}2\pi \frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}$$
(70)
and the effective magnetic field
$$\stackrel{~}{B}=d_{_{\text{AB}}}\frac{\mathrm{\Phi }_0}{2\pi }\mathrm{sin}2\pi \frac{\mathrm{\Phi }}{\mathrm{\Phi }_0};$$
(71)
in Eqs.(70) and (71), the averaging is performed with the distribution function of the flux in the line $`p(\mathrm{\Phi })`$.
Collecting the results together, the two-particle Green function reads
$`K_{\text{av}}(y,Y;y^{},Y^{};x)`$ $`=`$ (73)
$`={\displaystyle \frac{1}{2\pi x\mathrm{\lambda ̄}}}\mathrm{exp}\left[{\displaystyle \frac{i}{x\mathrm{\lambda ̄}}}(YY^{})(yy^{}){\displaystyle \frac{1}{8}}{\displaystyle \frac{x}{_{_{\text{AB}}}^2}}\left(y+y^{}+{\displaystyle \frac{(y+y^{})^2}{y+y^{}}}\right)i{\displaystyle \frac{\pi }{\mathrm{\Phi }_0}}\stackrel{~}{B}x(y+y^{})\right]`$
where $`y=y_1y_2`$, $`y^{}=y_1^{}y_2^{}`$, $`Y=\frac{1}{2}(y_1+y_2)`$ $`Y^{}=\frac{1}{2}(y_1^{}+y_2^{})`$ and $`x(xx^{})`$.
If compared with the Green function derived in Sect.V Eq.(57), the propagator Eq.(73) contains an additional term proportional $`\stackrel{~}{B}`$, finite to the extent the distribution $`p(\mathrm{\Phi })`$ is asymmetric. As we will see later, $`\stackrel{~}{B}`$ creates the Lorentz force and plays in dynamics the role of an effective magnetic field. A finite Lorentz force in an Aharonov-Bohm array is not readily obvious: A classical Lorentz force is absent since magnetic field is locally zero, whereas the quantum Aharonov-Bohm cross-section Eq.(39) is left-right and $`\mathrm{\Phi }\mathrm{\Phi }`$ symmetric and cannot explain $`\stackrel{~}{B}`$. Obviously, $`\stackrel{~}{B}`$ and the associated force is directly related to the transverse momentum transfer considered in Sect.III (see Eq.(A3)).
In the limit of dense array with small typical flux, $`\mathrm{\Phi }0`$, $`d_{_{\text{AB}}}\mathrm{}`$, the effective field $`\stackrel{~}{B}`$ reduces to to the macroscopic mean-field magnetic induction $`B=d_{_{\text{AB}}}<\mathrm{\Phi }>`$. In this limit, the Aharonov-Bohm array model is equivalent the $`\delta `$correlated field model Eq.(1) with $`^2d_{_{\text{AB}}}<\mathrm{\Phi }^2>`$, in the external homogeneous magnetic field $`\stackrel{~}{B}`$.
In general, however, $`\mathrm{\Phi }`$-periodic $`\stackrel{~}{B}`$ Eq.(71) is very different from the mean-field expectations. Depending on the flux distribution function $`p(\mathrm{\Phi })`$, the magnetic induction and $`\stackrel{~}{B}`$ may be in any relation. For instance, the mean-field induction can be always compensated to zero by adding some properly oriented lines of flux $`\mathrm{\Phi }_N=\frac{N}{2}\mathrm{\Phi }_0,N=1,2,\mathrm{}`$. However, the added lines do not affect the effective magnetic field $`\stackrel{~}{B}`$ seen by the particles (since $`\mathrm{sin}\left(2\pi \frac{\mathrm{\Phi }_N}{\mathrm{\Phi }_0}\right)=0`$). In an Aharonov-Bohm array the Lorentz force may be finite even when the macroscopic magnetic induction is zero.
### A Kinetic equation
As before, in Sect.V B, consider the spatially uniform situation when the density matrix of the partially coherent wave at $`x=0`$ is of the form $`\rho _{in}(y_1,y_2)=\rho _0(y)`$, $`y=y_1y_2`$. The density matrix $`\rho (y;x)`$ of the wave at distance $`x`$ can be found from Eqs.(44) and Eq.(73),
$`\rho (y;x)=\rho _0(y)\mathrm{exp}\left({\displaystyle \frac{xy}{2_{AB}^2}}i{\displaystyle \frac{2\pi }{\mathrm{\Phi }_0}}\stackrel{~}{B}xy\right).`$Again, as in Sect.V B, the density matrix obeys the following kinetic equation (compare with Eq.(63)):
$`v{\displaystyle \frac{}{x}}\rho +i{\displaystyle \frac{2\pi }{\mathrm{\Phi }_0}}\stackrel{~}{B}y\rho +\widehat{I}_{_{AB}}\rho =0,\widehat{I}_{_{AB}}={\displaystyle \frac{v}{2_{AB}^2}}|y|.`$Performing the Fourier transform, one gets the equation for the distribution function $`n_\phi `$ in Eq.(60):
$$v\frac{n_\phi }{x}+\frac{e\stackrel{~}{B}}{mc}\frac{n_\phi }{\phi }+\widehat{I}_{_{AB}}n_\phi =0.$$
(74)
The equation has the Boltzmann form and $`\stackrel{~}{B}`$ enters kinetics as a magnetic field. The collision integral $`\widehat{I}_{_{AB}}`$ has the form of Eq.(65) with the scattering rate
$`W(\varphi )={\displaystyle \frac{v\mathrm{\lambda ̄}}{2\pi _{_{AB}}^2}}{\displaystyle \frac{1}{\varphi ^2}}`$From here one can conclude that the relaxation is governed by an incoherent scattering by the flux lines: Indeed, $`W`$ is proportional to the density of lines and the contribution of a line is given by the small limit of the Aharonov-Bohm cross-section Eq.(39). Since even most dangerous small angle scattering fits this simple picture, it seems plausible that Eq.(74) can be generalized to arbitrary scattering angle using Eq.(39) as the probability scattering. Similar to Eq.(69), the kinetic equation reads
$`𝒗\mathbf{}\mathbf{}n_\phi {\displaystyle \frac{e\stackrel{~}{B}}{mc}}{\displaystyle \frac{n_\phi }{\phi }}+\widehat{I}_{_{AB}}n_\phi =0,\mathrm{\hspace{0.33em}\hspace{0.33em}0}<\phi <2\pi ,`$where the collision integral
$`\widehat{I}_{_{AB}}n_\phi ={\displaystyle \frac{1}{2\tau _{_{AB}}}}{\displaystyle \underset{0}{\overset{2\pi }{}}}{\displaystyle \frac{d\varphi }{2\pi }}{\displaystyle \frac{1}{\mathrm{sin}^2\frac{\varphi }{2}}}\left(n_\phi n_{\phi +\varphi }\right).`$with
$$\frac{1}{\tau _{_{AB}}}=\frac{\mathrm{}}{m}d_{_{\text{AB}}}1\mathrm{cos}2\pi \frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}.$$
(75)
As before, the collision integral is a regular linear operator regardless its singular scattering-out term. Its action is defined by the following relation
$`\widehat{I}_{_{AB}}e^{im\phi }={\displaystyle \frac{|m|}{\tau _{_{AB}}}}e^{im\phi }`$From here, one sees that $`\tau _{_{AB}}`$ has the meaning of the transport scattering time.
Applicability of the Boltzmann equation requires the mean free path $`lv\tau _{_{AB}}`$ to be large on the scale of the wave length $`\mathrm{\lambda ̄}`$.
$$\frac{\mathrm{\lambda ̄}}{l}\mathrm{\lambda ̄}^2d_{_{\text{AB}}}\mathrm{sin}^2\pi \frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}1$$
(76)
If typically $`\mathrm{\Phi }/\mathrm{\Phi }_01`$, the density of the lines must not be too high: $`d_{_{\text{AB}}}\mathrm{\lambda ̄}^21`$. In case of lines with small $`\mathrm{\Phi }`$, the condition is milder: $`\mathrm{\Phi }^2d_{_{\text{AB}}}\mathrm{\lambda ̄}^21`$.
As an illustration, the Drude conductivity tensor can be readily derived from the Boltzmann equation Eq.(75):
$`\sigma _{xx}`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^2N_0v^2\tau _{_{AB}}{\displaystyle \frac{1}{1+(\mathrm{\Omega }_c\tau _{_{AB}})^2}}`$ (77)
$`\sigma _{xy}`$ $`=`$ $`(\mathrm{\Omega }_c\tau _{_{AB}})\sigma _{xx}`$ (78)
where $`N_0`$ is the density of states and $`\mathrm{\Omega }_c=\frac{e\stackrel{~}{B}}{mc}`$ plays the role of the Larmor frequency. One sees that the Hall angle $`\theta _H`$, $`\mathrm{tan}\theta _H=\sigma _{xy}/\sigma _{xx}`$,
$$\mathrm{tan}\theta _H=\frac{\mathrm{sin}2\pi \frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}}{1\mathrm{cos}2\pi \frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}},$$
(79)
does not depend on the density of lines and parameters of the system. If the Aharonov-Bohm lines in array have same flux $`\mathrm{\Phi }`$, the Hall angle is simply
$$\mathrm{tan}\theta _H=\mathrm{cot}\left(\pi \frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}\right).$$
(80)
Counter intuition, the Hall effect is strongest at $`\mathrm{\Phi }0`$.
### B Landau quantization
The Hall angle Eq.(80) in an Aharonov-Bohm array may be large which means that the particle orbit tends to be a circle, similar to Larmor orbits in an external magnetic field. The periodic classical motion is expected to be quantized (the Landau quantization). To check this possibility, one should calculate the density of states of a particle moving in an Aharonov-Bohm array. In the quasiclassical region, when the energy of the particle $`E\mathrm{}\mathrm{\Omega }_c`$, the problem can be generally solved by the method used in . In the present paper, only the most promising case of small flux array is considered. Not to repeat the calculation in , qualitative arguments which lead to the same result, are presented.
In the linear in $`\mathrm{\Phi }`$ approximation, when scattering $`\mathrm{\Phi }^2`$ can be neglected, the transverse momentum Eq.(40) is transfered to the particle, and it undergoes periodic motion with the angular frequency $`\mathrm{\Omega }_c=\frac{e\stackrel{~}{B}}{mc}`$ along the Larmor circle with the radius $`R_L=v/\mathrm{\Omega }_c`$. In the small $`\mathrm{\Phi }`$-limit, the distinction between $`\stackrel{~}{B}`$ and magnetic induction is immaterial, and the Bohr-Sommerfeld quantization condition can be formulated as the requirement of the flux quantization,
$$\mathrm{\Phi }^{\text{orbit}}(E_N)=\mathrm{\Phi }_0(N+\frac{1}{2}),$$
(81)
where $`\mathrm{\Phi }^{\text{orbit}}(E)`$ is the total flux encircled by an orbit of the particle with energy $`E`$. Due to the randomness in the flux line positions, the number of the encircled lines fluctuates, and so does the total flux. The total flux $`\mathrm{\Phi }^{\text{orbit}}=\mathrm{\Phi }^{\text{orbit}}+\delta \mathrm{\Phi }`$ is a random Gaussian variable which fluctuates around the average $`\mathrm{\Phi }^{\text{orbit}}=\stackrel{~}{B}S(E)`$, $`S(E)=\pi R_L^2`$ being the Larmor circle area. For the flux lines with uncorrelated positions, the total flux fluctuates with the variance
$$\left(\delta \mathrm{\Phi }\right)^2=d_{_{\text{AB}}}S(E)\mathrm{\Phi }^2.$$
(82)
Driven by the flux enclosed by the orbit, the energy of the level fluctuates, $`E_N=E_N+\delta E`$. Neglecting the flux fluctuations in Eq.(81), $`\mathrm{\Phi }^{\text{orbit}}\mathrm{\Phi }^{\text{orbit}}`$, one gets the average energy of the N-th level $`E_N=\mathrm{}\mathrm{\Omega }_c(N+\frac{1}{2})`$. To preserve the quantization Eq.(81), the level energy acquires a shift $`\delta E`$ under the flux variations $`\delta \mathrm{\Phi }`$: From the condition
$`\delta \mathrm{\Phi }+\delta E_N{\displaystyle \frac{\mathrm{\Phi }}{E}}|_{E_N}=0,`$one gets the energy shift caused by the change of the flux $`\delta \mathrm{\Phi }`$:
$$\delta E_N=\mathrm{}\mathrm{\Omega }_c\frac{\delta \mathrm{\Phi }}{\mathrm{\Phi }_0}.$$
(83)
Combining Eq.(83) and Eq.(82), we see that the Landau levels acquire the Gaussian distribution form with variance (the width of the level)
$$\left(\delta E_N\right)^2=\frac{1}{\pi }\frac{\mathrm{}}{\tau _{_{AB}}}E$$
(84)
Physics here is similar to the inhomogeneous broadening: The Larmor circle sees different realizations of the random relief in different places of the “sample”, and the energy levels adjust their positions to the local conditions.
As already noticed in , it is rather unusual that importance of disorder (measured by the Landau level broadening Eq.(84)) increases with the energy $`E`$. The physical reason for this is that the larger is proportional to $`E`$ the area under the Larmor circle, the bigger are the absolute fluctuations of the number of the Aharonov-Bohm lines encircled by the orbit and the flux fluctuation $`\delta \mathrm{\Phi }`$.
The density of states is given by the sum over the discrete levels. The position of the level is the Gaussian of the width in Eq.(84) and centered at $`E_N=\mathrm{}\mathrm{\Omega }_c(N+\frac{1}{2})`$. The overlapping Landau levels create oscillating density of states. Similar to Ref., one applies the Poisson summation formula and finds the first harmonics of the oscillations:
$`\rho ^{osc}(E)={\displaystyle \frac{m}{\pi \mathrm{}^2}}\mathrm{exp}(\gamma )\mathrm{cos}\left(2\pi {\displaystyle \frac{E}{\mathrm{}\mathrm{\Omega }_c}}\right),`$where the damping of the oscillations is controlled by
$`\gamma ={\displaystyle \frac{2\pi }{\mathrm{\Omega }_c\tau _{_{AB}}}}{\displaystyle \frac{E}{\mathrm{}\mathrm{\Omega }_c}}.`$Therefore, the Shubnikov-De-Haas oscillations in the Aharonov-Bohm array are strongly suppressed: even when $`\mathrm{\Omega }_c\tau _{_{AB}}=\mathrm{tan}\theta _H1`$ and the Larmor circling is well pronounced, quasiclassical quantization at high Landau levels $`E\mathrm{}\mathrm{\Omega }_c`$ may not be seen because of the large damping $`\gamma `$.
The damping parameter $`\gamma `$ can be also presented as
$$\gamma =\frac{\pi }{2}\frac{1}{\mathrm{\lambda ̄}^2d_{_{\text{AB}}}}\frac{\mathrm{\Phi }^2}{\mathrm{\Phi }^2}\frac{\pi }{2}\frac{1}{\mathrm{\lambda ̄}^2d_{_{\text{AB}}}}$$
(85)
where $`\mathrm{\lambda ̄}`$ ($`=\mathrm{}/\sqrt{2mE}`$) is the wave length corresponding to the energy $`E`$; the equality sign realizes in the case when the lines have equal flux. One concludes from Eq.(85) that as a necessary condition, the Larmor motion may be quantized only if the wave length exceed the distance between Aharonov-Bohm lines. As one could expect, this requirement tends to be complementary to the condition of applicability of the Boltzmann equation in Eq.(76).
## VII Conclusions
The main goal of the paper has been to understand specific features of the transport properties of a quantum charge subject to a random magnetic field, namely, the features related to the long range correlations of the gauge potential and the anomalous forward scattering. The non-perturbative method used in the paper to handle the divergence is based on the paraxial approximation to the stationary Schrödinger equation (Section II).
The paraxial theory of magnetic scattering is presented in Section II A. To show usage of the theory, it is applied to the Aharonov-Bohm line problem in Section III. Being calculationally simple, the paraxial approximation proves to be rather efficient. The paraxial solution reproduces the small angle asymptotics of the Aharonov-Bohm exact solution for the plane incident wave. The wave packet solution (see Eq.(37)) allows one to resolve an old controversy discussed in the Introduction concerning the transverse force exerted by the Aharonov-Bohm line. One sees that the angular distribution in the outgoing wave is indeed left-right asymmetric, so that there is a finite momentum transfer in the transverse direction. However, the asymmetry is concentrated within the angular width of the incident wave and it cannot be described in terms of the differential cross-section. For an arbitrary incident wave, the transverse momentum transfered to the charge can be found by Eq.(40). By comparison with the exact solution, the validity of this formula has been recently confirmed by Berry .
The main result of the paper is the paraxial two-particle Green’s function averaged with respect to the random magnetic field. In the paraxial theory, a stationary 2D problem becomes equivalent to a non-stationary 1D one, and one can use the standard Feynman representation for the propagators. It turns out that the corresponding path integral can be evaluated exactly. The expressions for the Green’s function is given by Eqs.(57) and Eq.(73) for the Gaussian random magnetic field and the Aharonov-Bohm array model, respectively.
The paraxial two-particle Green’s function solves the quantum problem of the near forward multiple scattering by random gauge potential: for a given incident wave, one is able to find correlators $`\psi (𝒓_1)\psi ^{}(𝒓_2)`$, where averaging is performed with respect to the random field. To draw physical conclusions, two cases are analyzed: (i) (de)focusing of a coherent converging wave in the random magnetic field environment, and (ii) propagation of spatially homogeneous partially coherent beam.
The non-local character of interaction with magnetic field is clearly seen from the analysis in Section V A of defocusing of a converging beam. The lost of coherence measured by defocusing is controlled by the size of the entire area “occupied “ by the system that is the region where the wave function is finite: Indeed, it follows from Eq.(59) that the wave cannot be focussed if the random flux threading the area (aperture size)$`\times `$(focus length) is of order of the flux quantum $`\mathrm{\Phi }_0`$. Same conclusion follows from Eq.(62): having traveled a distance $`x`$, a perfect plane wave becomes a mixture of waves with the transverse momenta $`\mathrm{\Delta }p_y\mathrm{}x/^2`$ where the spatial coherence survives only within the region $`\mathrm{\Delta }y^2/x`$. Again, supporting the qualitative arguments presented in the Introduction, the coherence exists only within the area $`\mathrm{\Delta }x\times \mathrm{\Delta }y^2`$; the area is defined by the condition that the random flux is typically not bigger than the flux quantum $`\mathrm{\Phi }_0`$.
On the other hand, the evolution in the momentum space is rather ordinary: It is described by the usual Boltzmann equation derived in Sections V B and VI A. The only uncommon feature of the Boltzmann equation is that the collision integral being a well-defined operator nevertheless cannot be split in the scattering -in and -out regular parts: An attempt of the split produces two singular pieces, the two infinities canceling when combined. In the diagrammatic language, this can be rephrased as the cancellation of the divergences in the self energy and the vertex correction as observed in Ref..
The condition of applicability of the paraxial approximation can be derived from Eq.(12): At the distance $`x`$, the typical transverse size $`w`$ is of order $`min[w_0,^2/x]`$ where $`w_0`$ is the characteristic length in the initial distribution. For large enough $`x`$, $`w^2/x`$. Substituting this value into Eq.(12), one gets
$`x<x_{max},x_{max}\left({\displaystyle \frac{}{\mathrm{\lambda ̄}}}\right)^{\frac{3}{5}}.`$One sees that the theory is applicable in the non-perturbative region $`x`$ in the paraxial limit $`\mathrm{\lambda ̄}`$.
In the paraxial picture, the particle always moves in (almost) same direction. Therefore, any effect related to the Anderson localization is beyond the paraxial approximation. Although the Boltzmann equation allows for large angle scattering events, it is, of course, unable to describe the quantum localization either. Localization in random magnetic field remains a controversial issue, see Refs. and .
Two models of the random field has been considered in the paper: the Gaussian random magnetic field with zero average and the array of Aharonov-Bohm magnetic flux lines with arbitrary distribution of the line fluxes. Comparing the Green’s function in Eqs.(57), and (73) (with $`\stackrel{~}{B}0`$), one sees that the models are paraxially equivalent. (For the Gaussian model, Eq.(57) has been derived assuming that the field is zero on average, $`b=0`$. In a more general case, the Gaussian model with a finite magnetic induction $`B=b`$ has the Green’s function of the form in Eq.(73) with $`\stackrel{~}{B}`$ substituted for $`B`$).
The origin of the effective magnetic field in the Aharonov-Bohm array can be traced back to the left-right asymmetry in the scattering by an isolated Aharonov-Bohm line (see Section III). One sees that the transverse momenta $`\mathrm{\Delta }p_{}`$ Eq.(40) gained as a result of collisions with Aharonov-Bohm lines, add together giving rise to the Lorentz force and the effective magnetic field $`\stackrel{~}{B}`$ Eq.(71). Since any integer flux can be gauged out, $`\mathrm{\Delta }p_{}`$ and $`\stackrel{~}{B}`$ are periodic functions of the fluxes. Note that the magnetic induction $`B`$, which by definition equals to $`\mathrm{\Phi }`$, and the effective field $`\stackrel{~}{B}`$ Eq.(71) are same quantities only if lines flux is small. Generally, they may be in any relation. In particular, one may have finite $`\stackrel{~}{B}0`$ even if the magnetic induction $`B=0`$: A complex of 3 lines with the fluxes: $`(+\frac{\mathrm{\Phi }_0}{4},+\frac{\mathrm{\Phi }_0}{4},\frac{\mathrm{\Phi }_0}{2}`$) gives an example.
In Section VI A, the kinetic equation for a charge in an Aharonov-Bohm line array is solved to find the Drude conductivity tensor Eq.(77-78) and the Hall angle Eq.(79). These results for the array of lines are in agreement with Ref. where the Hall effect due to a single line has been studied. As discussed in , the Aharonov-Bohm periodicity, $`\stackrel{~}{B}_{\mathrm{\Phi }+\mathrm{\Phi }_0}=\stackrel{~}{B}_\mathrm{\Phi }`$ combined with the time reversal symmetry, $`\stackrel{~}{B}_\mathrm{\Phi }=\stackrel{~}{B}_\mathrm{\Phi }`$, requires that the half-integer flux lines, $`\stackrel{~}{\mathrm{\Phi }}=\frac{N}{2}`$ do not generate any Lorentz force. The Abrikosov vortex carries the flux $`\frac{1}{2}\frac{hc}{e}`$, therefore, does not exert the transverse force (of course, only in the limit when the particle wave length much exceeds the vortex size). As noticed before, this is in a qualitative agreement with the experimental observation of the reduced Hall effect exhibited by 2D electrons in the magnetic field of the Abrikosov vortices .
In an array of Aharonov-Bohm lines with small fluxes $`\mathrm{\Phi }\mathrm{\Phi }_0`$, the transverse momentum transfer $`\mathrm{\Phi }`$ (the “Lorentz” force exerted by $`\stackrel{~}{B}`$) is more efficient than relaxation of momentum the rate of which is $`\frac{1}{\tau _{\text{AB}}}\mathrm{\Phi }^2`$. Therefore, the “Lorentz” force significantly bends the particle trajectory leading to a large Hall angle Eq.(80) and, probably, to Landau quantization. As shown in Section VI B, these expectation are almost never met. Even if the Hall angle is large and the Larmor circling is well pronounced, the Landau levels are very broad. The inhomogeneous in its nature broadening is due to the fluctuations in the total flux threading the Larmor orbit.
Summarizing, the analysis of the forward scattering anomalous scattering in a random magnetic field (Gaussian and the Aharonov-Bohm array) has been presented using the paraxial approximation to the stationary Schrödinger equation. The gauge-invariant two-particle Green’s function has been found by exact evaluation of the corresponding Feynman integral. The propagation of coherent (defocusing) and incoherent (Boltzmann equation) waves has been analyzed, as well as Landau quantization in the Aharonov-Bohm array.
## VIII Acknowledgements
I am grateful to A. Mirlin and P. Wölfle for discussions, and to M. Ozana for the help with preparation of the paper. This work was supported by SFB 195 der Deutschen Forschungsgemeinschaft and partly Swedish Natural Science Research Council.
## A The transverse force
The transverse force that is the transfered momentum in the direction perpendicular to the velocity, can be easily found in the paraxial theory. The expectation value of the transverse momentum in the outgoing wave is
$$\widehat{p}_y_{out}=\frac{\mathrm{}}{i}\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑y\psi _{\text{out}}^{}(x,y)\frac{}{y}\psi _{\text{out}}(x,y)$$
(A1)
To calculate this integral, one exploits the fact that the free propagation at $`x>0`$ conserves the momentum, and, therefore, the average in Eq.(A1) does not depend on $`x>0`$. The goal now is to present it as an integral at $`x=0`$ where the $`\psi _{\text{out}}`$ in Eq.(27) is simplest possible; it cannot be done directly because of the eikonal discontinuity in $`\psi _{\text{out}}`$ at $`y=0`$
First, note that $`\psi _{\text{out}}(x,y)`$ in Eq.(22), is a well-behaved function of $`y`$ at any finite $`x>0`$. Thus, the derivate $`\frac{}{y}`$ in Eq.(A1) can be safely taken as the limit: $`\frac{}{y}f(y)=\frac{1}{2\eta }\left(f(y+\eta )f(y\eta )\right),\eta 0`$, and Eq.(A1) transforms to
$$\widehat{p}_y_{out}=\frac{1}{4\eta }\left(P_{out}(\eta )P_{out}(\eta )\right),\eta 0$$
(A2)
where
$`P_{out}(\eta )={\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑y\psi _{\text{in}}^{}(x,y_{})\psi _{\text{in}}(x,y_+),y_\pm =y\pm {\displaystyle \frac{1}{2}}\eta .`$Substituting everywhere the subscript $`outin`$, one gets the transverse momentum in the incoming wave $`\widehat{p}_y_{in}`$ expressed via the corresponding $`P_{in}(\eta )`$.
Identically, $`P_{out}(\eta )=e^{\frac{\mathrm{}}{i}\eta \widehat{p}_y}`$. Since $`P_{out}`$ is the expectation value of the conserving variable $`e^{i\eta \widehat{p}_y}`$, its value does not depend on $`x`$. Choose the point $`x=+0`$ where $`\psi _{\text{out}}`$ is given by Eq.(27) to evaluate $`P_{out}`$. After simple calculations,
$`P_{out}(\eta )=P_{in}(\eta )+|\eta |\left(e^{2\pi \widehat{\eta }i\frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}}1\right)\overline{|\psi _{\text{in}}|_\eta ^2},`$where
$`\overline{|\psi _{\text{in}}|_\eta ^2}={\displaystyle \frac{1}{\eta }}{\displaystyle \underset{\frac{\eta }{2}}{\overset{\frac{\eta }{2}}{}}}\psi _{\text{in}}^{}(y_{})\psi _{\text{in}}(y_+)`$Now, we see that the limit in Eq.(A2) is well defined, and we get for the transfered momentum $`\mathrm{\Delta }p_y=\widehat{p_y}_{out}\widehat{p_y}_{in}`$
$$\mathrm{\Delta }p_y=\mathrm{}|\psi _{\text{in}}(0)|^2\mathrm{sin}2\pi \frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}$$
(A3)
This the final result for the transfered momentum.
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# The Morphological and Structural Classification of Planetary Nebulae
## 1. Introduction
For decades now, it has been well established that PNe are the result of the evolution of low and intermediate mass stars (M$`<`$ 10 M). The actual nebular formation process has been well understood since Kwok, Purton & Fitzgerald (1978) explained the formation of PN as the result of the interaction of a low density fast wind with a high density slow wind. The only drawback of this two wind model is that it can not explain well the formation of asymmetric PNe. Since it has been long observed that most PNe do not have a round, regular shape, some other mechanism has to be invoked to produce asymmetry.
Mellema & Frank (1995) implemented an interacting wind model with an equatorial density enhancement. The fast low density wind interacts with an azimuthal dependent wind forming an asymmetrical PN. Rotation as a way to produce asymmetric PNe has been proposed by different authors (e.g. Calvet & Peimbert 1983; Ignace, Cassinelli & Bjorkman 1996; García-Segura et al. 1999). The presence of a magnetic field is also able to convey an asymmetrical nebular evolution (e.g. Pascoli 1992; Chevalier & Luo 1994; Soker 1998; García-Segura et al. 1999). The common envelope evolutionary phase, typical of a close binary star, also produces the appropriate equatorial density enhancement to make the PN ejecta asymmetric; a similar effect can be obtained by the presence of a substellar object in the system (e.g. Soker 1997).
By correlating the morphological class with the different nebular and stellar parameters, it may be possible to disclose the predominant mechanism responsible for the observed morphology.
Morphological studies of PNe have lacked a univocal classification scheme. Since the pioneer work by Curtis (1918), who discovered large structured ”haloes” around some PNe, there have been several studies all aimed at the same goal: a better understanding on the physical meaning of nebular shapes. Greig (1971) classified PNe into 15 morphological classes, ultimately grouped in two main classes: binebulous and circular; he found that binebulous PNe have lower galactic height distribution than circular PNe (Greig 1972). Zuckerman & Aller (1986) studied a sample of 108 PNe, classifying them into 16 morphological types, then regrouped the many classes into bipolar, elliptical, round, irregular, and other shapes; 50 % of their PN sample was bipolar, 30 % elliptical, 15 % round, and the rest irregular or other shapes. Zuckermann & Aller (1986) could not find any correlation between the morphological class and the C/O abundance. Balick (1987) divided the morphological classes in round, elliptical, and butterfly. He proposed an evolutionary sequence within each morphological type. Chu, Jacoby & Arendt (1987), studied a sample of 126 extended PNe. They found that the frequence of multiple-shell planetary nebulae (MSPN) was 50 %. Schwarz, Corradi & Stanghellini (1993) classified the Schwarz, Corradi & Melnick (1992) sample of southern Galactic PNe into elliptical, bipolar, pointsymmetric, irregular, and stellar shapes. On a subsample of the same catalog of PNe, Stanghellini, Corradi & Schwarz (1993) found that the central star distribution was different for bipolar and elliptical PNe. Corradi & Schwarz (1995), using a large PN sample ($``$ 400 PNe), found a different galactic height distribution for elliptical (z=320 pc) and for bipolar (z=130 pc) PNe.
Most of these classification schemes are based on incomplete or inhomogeneous samples. On this basis, Manchado et al. (1996) presented a complete set of northern Galactic PNe, to be analyzed for their morphological properties. In $`\mathrm{\S }`$2 we will illustrate this sample; we also discuss the completeness and the morphological classification based on the sample. $`\mathrm{\S }`$ 3 presents some of the relations found between the morphology and other nebular and stellar parameters, and includes a discussion on the possible evolutionary scheme for the different types of pointsymmetric PNe. Conclusions are in $`\mathrm{\S }`$ 4.
## 2. The PN sample and its completeness
The selection criteria for our homogeneous sample of northern Galactic PNe includes: (1) all northern PNe with declination larger than –11 in the Acker et al. (1992) catalog; (2) all the PNe larger than 4 arc-second; (3) images must be obtained in the narrow band filters (e.g. H$`\alpha `$, \[N II\] or \[O III\]).
There are 255 PNe that fulfill these selection criteria, 205 from the survey by Manchado et al. (1996), 28 from Balick (1987) and 22 from Schwarz, Corradi & Melnick (1992). In Figure 1 we show a selection of these images, representative of the various morphologies.
After a thorough analysis of the whole sample we decided to revise the morphological classification by Manchado et al. (1996). In fact, to make the individual morphological classes statistically meaningful, we decided to make only three major morphological classes: round (63 cases), elliptical (149 cases), and bipolar (43 cases) PNe. The quadrupolar PNe (7 cases) were included in the bipolar class, because the formation mechanisms could be very similar (Manchado, Stanghellini & Guerrero 1996). Pointsymmetry can be defined as a sub-class of elliptical and bipolar PNe: in fact, most pointsymmetric PNe have either bipolar or elliptical main shapes. A typical case of an elliptical pointsymmetric PNe is a PN with FLIERS (e.g. Balick et al. 1993). Figure 2 shows a diagram with the new morphological scheme.
Although the sample is complete as far as known PNe are concerned, there may be observational biases due to a different surface brightness limit for each morphological class. In order to investigate this possible bias we compare the statistical distribution of each morphological class taking into account the distance (Cahn, Kaler & Stanghellini 1992) and the extinction (Cahn, Kaler & Stanghellini 1992; Tylenda et al. 1992). The overall distribution of morphology in our PN sample is 58 % elliptical, 25 % round and 17 % bipolar. However, we realized that the sample is only complete up to a distance of 7 Kpc for all the morphological classes. If we were to limit the statistical studies to those PNe within a distance of 7 Kpc, the morphological distribution would be like 61 % elliptical, 26 % round and 13 % bipolar.
It can be argued that the statistical distances are not correct, so we used the extinction to infer completeness. If the sample is confined to the galactic plane, we can assume an extinction distance relationship of c = 0.2 per Kpc. Therefore, if we limit the sample to PNe with $`|b|<`$ 4 , c must be $`<`$ 1. In this newly defined space volume, we find 59 % elliptical, 28 % round, and 13 % bipolar PNe. Therefore, the results obtained using the statistical distance and the extinction rule are very similar, which means that the completeness within this space volume is sound.
## 3. Relations across morphological types
Each morphological class was correlated with a set of nebular and stellar parameters from the literature (for a complete reference list, see Manchado et al. 2000).
It was found that electronic density has different values for each morphological class; the median value is 1500 cm<sup>-3</sup> for elliptical, 400 cm<sup>-3</sup> for round, and 1000 cm<sup>-3</sup> for bipolar.
Dust temperatures were derived using the IRAS 25 and 60 $`\mu m`$ fluxes and dust emissivities taken from Draine & Lee (1984). Both elliptical and round PNe have a median dust temperature of 82 K, while bipolar temperature is 69 K.
The \[N II\]/H$`\alpha `$ ratio is higher for bipolar than for elliptical and round PNe. The N/O and He abundances of bipolar PNe are consistent with type I PNe (as defined by Peimbert & Torres-Peimbert 1983) and in MSPN they are consistent with type II PNe.
The galactic latitude distribution is different for each morphological class. The median of the galactic latitude is $`|\mathrm{b}|=7`$ for elliptical, $`|\mathrm{b}|=`$12 for round and $`|\mathrm{b}|=3`$ for bipolar PNe. Figure 3 shows the galactic distribution of these three morphological classes. The median values of the Galactic height are $`<\mathrm{z}>=`$308 pc for elliptical, $`<\mathrm{z}>=`$753 pc for round, and $`<\mathrm{z}>=`$179 pc for bipolar.
By studying the pointsymmetric PNe, we find that for elliptical pointsymmetric the scale height is $`<\mathrm{z}>=`$310 pc, while the elliptical PNe without pointsymmetry have a $`<\mathrm{z}>=`$ 308 pc. Bipolar with pointsymmetric structure have $`<\mathrm{z}>=`$248 pc, while bipolar without pointsymmetry $`<\mathrm{z}>=`$110 pc.
The different galactic height for the various morphological classes may imply a different stellar population. According to Miller & Scalo (1979) $`<\mathrm{z}>=`$300 pc implies that the progenitor star has mass $`<`$ 1.0 M. For $`<\mathrm{z}>=`$150 pc, the mass is $`>`$ 1.5 M, while for $`<\mathrm{z}>=`$230 pc and $`<\mathrm{z}>=`$ 110 pc masses will be $`>`$ 1.2 M and $`>`$ 1.9 M. Therefore elliptical and bipolar PNe might have different distribution masses for their progenitor stars ($`<`$ 1.0 M and $`>`$ 1.5 M). In the bipolar class there is also a mass segregation. In fact, according to to their scale height on the Galactic plane, PNe with pointsymmetric structure evolve from stellar masses $`>`$ 1.2 M, while those without the pointsymmetric structure evolve from stellar masses $`>`$ 1.5 M. These results are consistent with the other results from our statistical analysis, as bipolar PNe have higher N/O and helium abundances.
The fact that two different mass distributions can be inferred for bipolar PNe, depending on the presence of pointsymmetry, can be explained with two evolutionary schemes for the two types: a single, high mass star would form a bipolar PNe, due possibly to rotation and magnetic field (e.g. García-Segura et al. 1999), while a bipolar pointsymmetric PN could be due to magnetic collimation around a precessing star (e.g. García-Segura 1997).
## 4. Conclusions
A proper statistical analysis of a complete sample of PNe has allowed us to classify them into elliptical, round, and bipolar, with the sub-classes of multiple-shell and pointsymmetric PNe. It was found that 60 % of our PN sample present an elliptical shape, while 26 % are round, and 13 % bipolar.
We use statistical distances that appear to be sound for the task. If they are indeed correct, the different scale heights that characterize each morphological class hint of different progenitor mass distribution for each class. Two evolutionary schemes are proposed for bipolar PNe and bipolar PNe with pointsymmetric structure.
## Acknowledgment
The work of EV and AM is supported by a grant of the Spanish DGES PB97-1435-C02-01. MAG is supported by the Spanish Ministerio de Educación y Cultura.
## References
Acker, A., Ochsenbein, F., Stenholm, B., Tylenda, R., Marcout, J., Schohn, C.: 1992, Strasbourg–ESO catalogue of Galactic planetary nebulae, ESO
Balick, B. 1987, AJ 94, 671
Balick, B., Rugers, M., Terzian, Y., Chengalur, J.N. 1993, ApJ 411, 778
Cahn, J.H., Kaler, J., & Stanghellini, L. 1992, A&AS 94, 399
Calvet, N., & Peimbert, M. 1983, Rev. Mex. Astron. Astrofis. 5, 319
Chevalier, R. A. & Luo, D. 1994, ApJ, 421, 225
Chu, Y.-H., Jacoby, G., & Arendt, R. 1987, ApJSS, 64, 529
Corradi, R. L. M., & Schwarz H. E. 1995, A&A 293, 871
Curtis, H.D. 1918, Pub. Lick Obs XIII, 55
Draine B.T. & Lee, H.M. 1984, ApJ, 285, 89
García-Segura, G. 1997, ApJ 489, L189
García-Segura, G., Langer, N., Różyczka, M., Mac-Low, M. Franco, J. 1999, ApJ, 517, 767
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Kwok, S., Purton, C. R., & Fitzgerald, P. M. 1978, ApJ 219, L 125
Manchado, A., Guerrero, M., Stanghellini, L., & Serra–Ricart, M. 1996, The IAC Morphological Catalog of Northern Galactic planetary nebulae , (La Laguna: IAC)
Manchado, A., Stanghellini, L., & Guerrero, M., 1996, ApJ, 466, L95
Manchado, A., Villaver, E, Stanghellini, L., & Guerrero, M., 2000, ApJS (in preparation)
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Pascoli, G. 1992, PASP, 104, 350
Peimbert, M. & Torres-Peimbert, S. 1983, IAU Symp 103, p. 233 (Reidel:Dordrecht)
Schwarz, H. E., Corradi, R., & Melnick, J. 1992, A&AS, 96, 23
Schwarz, H. E., Corradi, R. & Stanghellini L. 1993, IAU Symp 155, p. 214, eds. Weinberger and Acker, (Kluwer:Dordrecht)
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# Extended Quintessence: imprints on the cosmic microwave background spectra
## 1 Introduction
One of the most interesting novelty in modern cosmology is the observational trend for an accelerating Universe, as suggested by distance measurements to type Ia Supernovae . These results astonishingly indicate that almost two thirds of the energy density today is vacuum energy.
It has been thought that this vacuum energy could be mimicked by a minimally-coupled scalar field , considered as a ”Quintessence” (Q). The main features of such a vacuum energy component, that could also allow to distinguish it from a cosmological constant, are its time-dependence as well as its capability to develop spatial perturbations.
Theoretically, Quintessence models are attractive, since they offer a valid alternative explanation of the smallness of the present vacuum energy density instead of the cosmological constant; indeed, we must have $`|\rho _{vac}|<10^{47}`$ GeV<sup>4</sup> today, while quantum field theories would predict a value for the cosmological constant which is larger by more than 100 orders of magnitude . Instead, the vacuum energy associated to the Quintessence is dynamically evolving towards zero driven by the evolution of the scalar field. Furthermore, in the Quintessence scenarios one can select a subclass of models, which admit ”tracking solutions” : here a given amount of scalar field energy density today can be reached starting from a wide set of initial conditions.
The effects of possible couplings of this new cosmological component with the other species have been explored in recent works, both for what regards matter and gravity . Here we review some of the results obtained in a recent paper , for what concerns the effects on the Cosmic Microwave Background (CMB) anisotropy: this scenario has been named ‘Extended Quintessence’ (EQ), by meaning that the scalar field coupled with the Ricci scalar $`R`$ has been proposed as the Quintessence candidate, in analogy with Extended Inflation models .
## 2 Cosmological dynamics in scalar-tensor theories of gravity
The action $`S=d^4x\sqrt{g}[F(\varphi )R\varphi ^{;\mu }\varphi _{;\mu }2V(\varphi )+L_{fluid}]`$ represents scalar-tensor theories of gravity, where $`R`$ is the Ricci scalar and $`L_{fluid}`$ includes matter and radiation.
We assume a standard Friedman-Robertson-Walker (FRW) form for the unperturbed background metric, with signature $`(,+,+,+)`$, and we restrict ourselves to a spatially flat universe. The FRW and Klein Gordon equations are
$$^2=\frac{a^2\rho _{fluid}}{3F}+\frac{\dot{\varphi }^2}{6F}++\frac{a^2V}{3F}\frac{\dot{F}}{F},\ddot{\varphi }+2\dot{\varphi }=\frac{a^2F_{,\varphi }R}{2}a^2V_{,\varphi },$$
(1)
where the overdot denotes differentiation with respect to the conformal time $`\tau `$ and $`=\dot{a}/a`$. Furthermore, the continuity equations for the individual fluid components are $`\dot{\rho }_i=3(\rho _i+p_i)`$.
For what concerns our treatment of the perturbations , we give here only the very basic concepts. A scalar-type metric perturbation in the synchronous gauge is parameterized as
$$ds^2=a^2[d\tau ^2+(\delta _{ij}+h_{ij})dx^idx^j];$$
(2)
by linearly perturbing the Einstein and Klein Gordon equations above, the equation for the metric perturbing quantities can be derived; these equations are linked to the fluid perturbed quantities, from any species including $`\varphi `$, obeying the perturbed continuity equations.
Let us define now the gravitational sector of the Lagrangian. We require that $`F`$ has the correct physical dimensions of $`1/G`$. Note that all this fixes the link between the value of $`F`$ today and the Newtonian gravitational constant $`G`$: $`F_0=F(\varphi _0)=1/8\pi G`$. Different forms of $`F(\varphi )`$ can be considered . In Induced Gravity (IG) models, that we treat here, the gravitational constant is directly linked to the scalar field itself, as originally proposed in the context of the Brans-Dicke theory:
$$F(\varphi )=\xi \varphi ^2,$$
(3)
where $`\xi `$ is the IG coupling constant. Note that solar system experiments already offer constraints to the viable values of $`\xi `$, that may be easily obtained by integrating the background equations . The dynamics of $`\varphi `$ is determined by its coupling with $`R`$, as well as by its potential, that is responsible for the vacuum energy today; we take the simplest inverse power potential, $`V(\varphi )=M^5/\varphi `$, where the mass-scale $`M`$ is fixed by the level of energy contribution today from the Quintessence. In our integrations, we adopt adiabatic initial conditions. We require that the present value of $`\mathrm{\Omega }_\varphi `$ is $`0.6`$, with Cold Dark Matter at $`\mathrm{\Omega }_{CDM}=0.35`$, three families of massless neutrinos, baryon content $`\mathrm{\Omega }_b=0.05`$ and Hubble constant $`H_0=50`$ Km/sec/Mpc; the initial kinetic energy of $`\varphi `$ is taken equal to the potential one at the initial time $`\tau =0`$.
## 3 Effects on the CMB
The phenomenology of CMB anisotropies in EQ models is rich and possesses distinctive features. In Fig.1, the effect of increasing $`\xi `$ on the power spectrum of COBE-normalized CMB anisotropies is shown. The rise of $`\xi `$ makes substantially three effects: the low $`\mathrm{}`$’s region is enhanced, the oscillating one attenuated, and the location of the peaks shifted to higher multipoles. Let us now explain these effects. The first one is due to the integrated Sachs-Wolfe effect, arising from the change from matter to Quintessence dominated era occurred at low redshifts. This occurs also in ordinary Q models, but in EQ this effect is enhanced. Indeed, in ordinary Q models the dynamics of $`\varphi `$ is governed by its potential; in the present model, one more independent dynamical source is the coupling between the Q-field and the Ricci curvature $`R`$. The dynamics of $`\varphi `$ is boosted by $`R`$ together with its potential $`V`$. As a consequence, part of the COBE normalization at $`\mathrm{}=10`$ is due to the Integrated Sachs-Wolfe effect; thus the actual amplitude of the underlying scale-invariant perturbation spectrum gets reduced. In addition, it can be seen that the Hubble length was smaller in the past in EQ than in Q models. This has the immediate consequence that the horizon crossing of a given cosmological scale is delayed, making the amplitude the acoustic oscillations slightly decreasing since the matter content at decoupling is increased.
Finally, note how the location of the acoustic peaks in term of the multipole $`\mathrm{}`$ at which the oscillation occurs, is shifted to the right. Again, the reason is the time dependence of the Hubble length, which at decoupling subtended a smaller angle on the sky. It can indeed verified that the ratio of the peak multipoles in Fig.1 coincides numerically with the the ratio of the values of the Hubble lengths at decoupling in EQ and Q models .
We have used here values of $`\xi `$ large in order to clearly show the CMB effects. It can be seen these values do not satisfy the solar system experimental constraints; however, a smaller $`\xi `$ produces the same spectral features, reduced but still detectable by the future generation of CMB experiments, able to bring the accuracy on the CMB power spectrum at percent level up to $`\mathrm{}1000`$ .
## References
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# 1 Introduction
## 1 Introduction
The production of $`W\gamma `$ and $`Z\gamma `$ in hadronic collisions has been studied extensively since the Born cross sections have been computed . In particular, these processes allow to study the triple gauge boson couplings $`WW\gamma `$, $`ZZ\gamma `$ and $`Z\gamma \gamma `$. The study of these couplings is mainly motivated by the hope that some new physics may modify them. If the new physics occurs at an energy scale well above that being probed experimentally, it is possible to integrate it out. The result is an effective theory which might result in non standard triple gauge boson couplings.
Both collaborations at the Tevatron have studied the production of $`W\gamma `$ and $`Z\gamma `$ pairs. The bounds on the anomalous couplings obtained at the Tevatron tend to be less constraining than those obtained at LEP . However it has to be kept in mind that these analyses are complementary. At the hadron colliders a whole range in the center of mass energy is tested, whereas at LEP the center of mass energy is fixed by the collider. Furthermore, with Run II the expected number of events at the Tevatron increases substantially. Assuming a data sample of 2 fb<sup>-1</sup>, more than 3000 $`W\gamma \mathrm{}\nu \gamma `$ events and 700 $`Z\gamma \mathrm{}\mathrm{}\gamma `$ events are expected for each experiment . Of course, the expected number of events is even bigger for the LHC.
Anomalous triple gauge boson couplings lead to deviations from Standard Model predictions. Obviously, observables for which these deviations are enhanced offer better chances to find new physics or get tighter constraints on anomalous couplings. There are basically two classes of such observables. Either we consider observables which are strongly suppressed in the Standard Model or observables with large transverse momentum (or center of mass energy). In both cases, the inclusion of next-to-leading order (NLO) QCD corrections is mandatory.
A prominent example of an observable that is suppressed in the Standard Model is the so called radiation zero for $`W\gamma `$ production. At leading order (LO) there exist some kinematic configurations for which the amplitude vanishes . This is manifest in some observables as a dip in the rapidity distributions. Since anomalous coupling contributions fill in the dips, there seemed to be excellent prospects to obtain accurate limits for them from experimental data. Unfortunately, next-to-leading order QCD corrections strongly affect the LO analysis. They have the same effect as the anomalous coupling contributions. The dips are filled in, making the extraction of anomalous couplings quite more difficult.
For processes with large transverse momentum or center of mass energy, the NLO corrections are particularly large. This is due to the fact that the cross sections in these cases get large contributions from gluon induced partonic subprocesses, which only enter in a next-to-leading order description of the cross section. Thus, even though the anomalous contributions are enhanced in these regions, a calculation at NLO in $`\alpha _s`$ is required to reliably exclude (or establish) physics beyond the Standard Model.
The relevance of NLO corrections was first shown for the production of real (spin-summed) $`W`$ and $`Z`$ bosons with Standard Model couplings and without considering lepton decays and spin correlations . These calculations were later extended, in order to include the leptonic decays and anomalous couplings . However, the full one-loop amplitudes including leptonic decays became available only very recently . Therefore, refs. included decay correlations everywhere except for the finite part of the virtual contributions.
In this paper we present order $`\alpha _s`$ results for the production of $`W\gamma `$ and $`Z\gamma `$ in hadronic collisions, including the full leptonic correlations. We work in the narrow-width approximation, where only ‘single-resonant’ Feynman diagrams have to be considered. The simplicity of the helicity method allows to take into account anomalous couplings as well and present for the first time analytical expressions for the corresponding amplitudes.
For the case of $`WW,ZZ`$ and $`WZ`$ production at hadron colliders, some results beyond the narrow-width approximation are known. The narrow-width approximation requires only the calculation of ‘doubly-resonant’ Feynman diagrams. However, for these processes, also the amplitudes including ‘single-resonant’ diagrams have been computed and implemented into a Monte Carlo program .
To perform the phase space integration we use the subtraction method discussed in ref. . This allows for a straightforward implementation of the one-loop $`q\overline{q}^{}V\gamma \mathrm{}\mathrm{}^{}\gamma `$ and bremsstrahlung $`q\overline{q}^{}gV\gamma g\mathrm{}\mathrm{}^{}\gamma `$ amplitudes, presented in ref. $`(V\{Z,W\})`$. The constructed Monte Carlo code allows the computation of any infrared-safe observable<sup>1</sup><sup>1</sup>1The corresponding Fortran codes are available upon request.
A brief overview of the calculation is given in section 2, were we summarize the input parameters used and the cuts implemented to obtain our phenomenological results. In section 3 we first present some benchmark cross section numbers for both $`W\gamma `$ and $`Z\gamma `$ production at the LHC and study the typical scale dependence of some observables at NLO in the Standard Model. Since many distributions have been studied in the past, we refrain from doing a detailed analysis. However, as soon as more precise data becomes available such an analysis can easily be done.
In section 4 we concentrate on anomalous triple gauge boson couplings. We describe the parameterization of the triple gauge boson vertex in terms of anomalous coupling parameters and search for the kinematical region where its effect is amplified, namely at large transverse momentum for both photon and leptons. We also analyze the possibility of seeing the effect of approximate radiation zeros in the $`W\gamma `$ process, i.e. by looking for ’dips’ in rapidity distributions. In order to avoid the arbitrariness introduced by form factors, we propose to analyze the anomalous couplings as a function of the squared partonic center of mass energy $`\widehat{s}`$. This has been suggested previously for the $`Z\gamma `$ case , where such an analysis is straightforward. We extend this idea to the $`W\gamma `$ production. This case is more involved, since a complete reconstruction of $`\widehat{s}`$ is impossible, due to the appearance of a non observed neutrino in the $`W`$ decay. Particularly, we present an observable quantity which is highly correlated to $`\widehat{s}`$ and, therefore, allows such an analysis even for $`W\gamma `$ production. Finally, in section 5 we give our conclusions and in the appendix we present analytical expressions for the amplitudes relevant for the processes under consideration.
## 2 Formalism
The helicity amplitudes needed for the computation of the NLO corrections to $`W\gamma `$ and $`Z\gamma `$ production in the Standard Model were presented in ref. . The amplitudes relevant for the inclusion of anomalous couplings are presented in the appendix. In order to cancel analytically the soft and collinear singularities coming from the bremsstrahlung and one loop parts, we have used the version of the subtraction method presented in ref. . The amplitudes are therefore implemented into a numerical Monte Carlo style program, which allows to calculate any infrared-safe physical quantity with arbitrary cuts.
Obviously, the Monte Carlo program can be used for the Tevatron and the LHC. However, in this paper we will mainly concentrate on results for the LHC collider, which corresponds to $`pp`$ scattering at $`\sqrt{s}=14`$ TeV. Unless otherwise stated, the results are obtained using the following cuts: we make a transverse momentum cut of $`p_T^{\mathrm{}}>25`$ GeV for the charged leptons and the rapidity is limited to $`|\eta |<2.4`$ for all detected particles. The photon transverse momentum cut is $`p_T^\gamma >50(100)`$ GeV for $`W\gamma `$ ($`Z\gamma `$) production. For the $`W\gamma `$ case we require a minimum missing transverse momentum carried by the neutrinos $`p_T^{\mathrm{miss}}>50`$ GeV. Additionally, charged leptons and the photons must be separated in the rapidity-azimuthal angle by $`\mathrm{\Delta }R_\mathrm{}\gamma =\sqrt{(\eta _\gamma \eta _{\mathrm{}})^2+(\varphi _\gamma \varphi _{\mathrm{}})^2}>0.7`$. Moreover, since our calculation is done in the narrow-width approximation and, therefore, ignores the radiation of photons from the final state leptons, we apply an additional cut to suppress the contribution from the off-resonant diagrams. For that purpose, we require the transverse mass $`M_T>90`$ GeV for $`W\gamma `$ production and the invariant mass of the $`\mathrm{}\mathrm{}\gamma `$ system $`M_\mathrm{}\mathrm{}\gamma >100`$ GeV for the $`Z\gamma `$ case.
Finally, photons can also be significantly produced at LHC from the fragmentation of a final state parton<sup>2</sup><sup>2</sup>2This contribution is also known in the literature as ‘bremsstrahlung’. Unfortunately, fragmentation functions of partons into photons are not very well determined and the NLO calculation for such a contribution is not available yet. In principle, a full NLO calculation should include it however, since only the sum of the ‘direct’ plus ‘fragmentation’ components is physically well defined at NLO (only in the sum all collinear singularities cancel out). In order to circumvent this problem, we include the LO component of the fragmentation part but using NLO fragmentation distributions, where we can factorize the final state $`q\gamma `$ collinear singularities. Since the ‘fragmentation’ component can be further suppressed implementing certain cuts (see below) the lack of its NLO calculation is not expected to affect the final result beyond the few percent level.
The fragmentation contribution constitutes a background to the search of anomalous couplings, since it does not involve any triple gauge boson coupling. Fortunately, there is a way to suppress its contribution by requiring the photons to be isolated from hadrons. In this paper we require the transverse hadronic momentum in a cone of size $`R_0=0.7`$ around the photon to be smaller than a small fraction of the transverse momentum of the photon
$$\underset{\mathrm{\Delta }R<R_0}{}p_T^{\mathrm{had}}<0.15p_T^\gamma $$
(1)
This completes the definition of our ’standard’ cuts.
When indicated, we also apply a jet-veto, which means that we reject any event where a jet of $`p_T^{\mathrm{jet}}>50`$ GeV and $`|\eta _{\mathrm{jet}}|<2.5`$ is observed.
In our results we do not include the branching ratios of the vector boson into leptons. They can be taken into account by simply multiplying our final results with the corresponding branching ratio. For both the LO and NLO results we use the latest (corrected) set of parton distributions of MRST(cor01) and the two loop expression (with $`n_f=5`$ for the typical scales of these processes) for the strong coupling constant with $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}(n_f=4)=300`$ MeV which corresponds to $`\alpha _s(M_Z)=0.1175`$. For the fragmentation component we use the fragmentation functions from ref. .
Since we are particularly interested in the large $`p_T`$ tail, which is more sensitive to the anomalous coupling contributions, we use (unless otherwise stated)
$$\mu ^2=\mu _{\mathrm{st}}^2M_V^2+\frac{1}{2}\left[(p_T^V)^2+(p_T^\gamma )^2\right].$$
(2)
as the ’standard’ scale for both the factorization and renormalization scales.
Contributions from $`b`$ and $`t`$ quark initial states have been neglected and, consistently, the following values have been used for the Cabibbo-Kobayashi-Maskawa (CKM) matrix elements in the case of $`W\gamma `$ production: $`|V_{ud}|=|V_{cs}|=0.975`$ and $`|V_{us}|=|V_{cd}|=0.222`$.
The masses of the vector bosons have been set to $`M_Z=91.187`$ GeV and $`M_W=80.41`$ GeV. We do not include any QED or electroweak corrections but choose the coupling constants $`\alpha `$ and $`\mathrm{sin}^2\theta _W`$ in the spirit of the ‘improved Born approximation’ , with $`\mathrm{sin}^2\theta _W=0.230`$. Notice that the observable is order $`\alpha ^2`$; within the same spirit we use the running $`\alpha =\alpha (M_Z)=1/128`$ for the coupling between the vector boson and the quarks (to effectively take into account the EW corrections) whereas we keep $`\alpha =1/137`$ for the photon coupling. It is worth noticing that this modification results in a 7% change in the normalization of the cross section with respect to the standard approach of using both running coupling constants.
## 3 Standard Results
We begin the presentation of our results with some total cross section numbers in Table 1, which can be useful for future checks and for an estimate of the number of events to be observed at the LHC. The first three results were obtained by imposing only the cut on the transverse momentum of the photon, i.e. $`p_T^\gamma >50(100)`$ GeV for $`W\gamma `$ ($`Z\gamma `$) production. Apparently, the NLO corrections as well as the fragmentation contribution are very large. As discussed above, the relative importance of the fragmentation contribution can be reduced substantially by applying the isolation cut prescription. This can be seen from the results for the total cross section obtained after the implementation of our standard cuts, which are also presented in Table 1. Unfortunately, most of the previous publications on the subject do not present cross sections numbers. Nevertheless, we have compared many of the plots shown in refs. , specially for the case of real (spin-summed) $`W/Z\gamma `$ production (which is not affected by lepton correlations). Within the precision that can be reached in such a comparison, we found good agreement.
In what follows we will estimate the theoretical uncertainty of our results by analyzing the variation of various distributions when changing the scale $`\mu `$ by a factor of two in both directions $`1/2\mu _{\mathrm{st}}\mu 2\mu _{\mathrm{st}}`$. Since many observables already have been studied in the past and in order to avoid the proliferation of plots, we refrain from presenting a detailed analysis here. We simply concentrate on a couple of typical examples in order to give a general picture and illustrate the importance of NLO corrections in both $`W\gamma `$ and $`Z\gamma `$ production.
In Figure 1 we show the scale dependence of the $`p_T`$ distribution of the photon in $`W^+\gamma `$ production with the standard cuts applied (upper curves) and also with the additional requirement of a jet-veto (lower curves). As can be observed, the scale dependence is still large ($`\pm 10\%`$) as long as only the standard cuts are applied. However, it is considerably reduced when the jet-veto is applied. The situation is similar to what has been observed in the case of $`WW`$ production and is caused by the suppression of the contribution from the $`qg`$ initial state appearing for the first time at NLO. Since this initial state dominates the cross section, the NLO result behaves, regarding the scale dependence, effectively like a LO one.
In the inset plot we present the ratio between the NLO and LO results (with the standard scale), which remains larger than 3 and increases with the photon transverse momentum. This clearly shows that the LO calculation is not even sufficient for an understanding of the shape of the distribution, since the NLO effect goes beyond a simple change in the normalization. As is well known , the relevance of the NLO corrections for this process is mainly due to the breaking of the radiation zero appearing at LO and to the large $`qg`$ initial state parton luminosity at the LHC. We also show the ratio of the NLO jet-veto and the LO result. As expected, this ratio is closer to 1, again due to the fact that most of the contributions coming from the new subprocesses appearing at NLO are suppressed by the jet-veto.
It is worth mentioning that the scale dependence of the LO result turns out to be very small. This is an artificial effect and illustrates that a small scale dependence is by no means a guaranty for small NLO corrections. In fact, there is no renormalization scale dependence at all at LO. The only scale dependence comes from the factorization scale dependence of the parton distribution functions. Furthermore, we would like to mention that the situation concerning the scale dependence is slightly more favorable at the Tevatron. This is simply due to the fact that the gluon initiated process is less important.
In Figure 2 we study the lepton correlation in the azimuthal angle for $`Z\gamma `$ production, $`\mathrm{\Delta }\varphi _{\mathrm{}\mathrm{}}=|\varphi _{\mathrm{}^{}}\varphi _\mathrm{}^+|`$. Notice that this observable can be studied at NLO since we take fully into account the spin correlations between the leptons in our implementation of the one-loop corrections. The NLO corrections are rather sizeable and increase the cross section by $`50\%`$ for small $`\mathrm{\Delta }\varphi _{\mathrm{}\mathrm{}}`$. The region $`\mathrm{\Delta }\varphi _{\mathrm{}\mathrm{}}>2`$ (with our standard cuts) is kinematically forbidden unless a jet with a high transverse momentum is produced. Therefore, the cross section vanishes at LO and it is strongly suppressed for the NLO calculation with jet-veto. In this region, the full NLO calculation is effectively only a LO calculation and its scale dependence becomes larger, as expected.
Because there is no radiation zero appearing at LO for $`Z\gamma `$ production, the NLO corrections are under better control in the kinematical region where the LO cross section does not vanish. Nevertheless, for large transverse momentum, the process with a $`qg`$ initial state again dominates the NLO contribution and the corrections increase considerably .
Finally we mention that we also considered a more stringent photon isolation prescription, introduced by Frixione . This prescription completely eliminates the fragmentation contribution. We have checked that the main features of all studied distributions remain unchanged when it is imposed.
## 4 Sensitivity to Anomalous Couplings
The study of triple gauge boson couplings is motivated by the hope that some physics beyond the Standard Model leads to a modification of these couplings which eventually could be detected. In order to quantify the effects of the new physics an effective Lagrangian is introduced which in principle contains all Lorentz invariant terms. The prefactor of these operators are the anomalous couplings. A general approach is impractical, since it would lead to a proliferation of new couplings. Therefore, some additional constraints have to be imposed.
The usual choice for $`W\gamma `$ production is to insist on electromagnetic gauge invariance and on $`C`$ and $`P`$ invariance. Also, only operators of dimension six or less are considered. This leads to a momentum-space vertex $`W_\alpha ^{}(q)W_\beta ^+(\overline{q})\gamma _\mu (p)`$ (where all momenta are outgoing $`p+q+\overline{q}=0`$ ) which can be written as
$`\mathrm{\Gamma }_{WW\gamma }^{\alpha \beta \mu }(q,\overline{q},p)`$ $`=`$ $`\overline{q}^\alpha g^{\beta \mu }\left(2+\mathrm{\Delta }\kappa ^\gamma +\lambda ^\gamma {\displaystyle \frac{q^2}{M_W^2}}\right)q^\beta g^{\alpha \mu }\left(2+\mathrm{\Delta }\kappa ^\gamma +\lambda ^\gamma {\displaystyle \frac{\overline{q}^2}{M_W^2}}\right)`$ (3)
$`+\left(\overline{q}^\mu q^\mu \right)\left[g^{\alpha \beta }\left(1+{\displaystyle \frac{1}{2}}p^2{\displaystyle \frac{\lambda ^\gamma }{M_W^2}}\right)+{\displaystyle \frac{\lambda ^\gamma }{M_W^2}}p^\alpha p^\beta \right],`$
where the overall coupling has been chosen to be $`|e|`$. Note that in the Feynman rule for this vertex there is also a factor $`i`$ that is conventionally not included in $`\mathrm{\Gamma }^{\alpha \beta \mu }`$. In the Standard Model we have $`\mathrm{\Delta }\kappa ^\gamma =\lambda ^\gamma =0`$.
For $`Z\gamma `$ production, we consider operators up to dimension 8 (all of them $`C`$ odd) resulting in $`ZZ\gamma `$ and $`Z\gamma \gamma `$ couplings. The non-standard $`Z_\alpha (q_1)\gamma _\beta (q_2)Z_\mu (p)`$ momentum-space vertex is given by
$`\mathrm{\Gamma }_{Z\gamma Z}^{\alpha \beta \mu }(q_1,q_2,p)={\displaystyle \frac{i(p^2q_1^2)}{M_Z^2}}(`$ $`h_1^Z\left(q_2^\mu g^{\alpha \beta }q_2^\alpha g^{\mu \beta }\right)+{\displaystyle \frac{h_2^Z}{M_Z^2}}p^\alpha \left(Pq_2g^{\mu \beta }q_2^\mu p^\beta \right)`$ (4)
$``$ $`h_3^Z\epsilon ^{\mu \alpha \beta \nu }q_{2\nu }{\displaystyle \frac{h_4^Z}{M_Z^2}}\epsilon ^{\mu \beta \nu \sigma }p^\alpha p_\nu q_{2\sigma })`$
where the overall coupling has been chosen to be $`|e|`$ (and $`ϵ^{0123}=+1`$). This vertex is absent altogether in the Standard Model. The non-standard $`Z_\alpha (q_1)\gamma _\beta (q_2)\gamma _\mu (p)`$ momentum-space vertex can be obtained from eq. (4) by setting $`q_1^20`$ and replacing $`h_i^Zh_i^\gamma `$. Notice that the vertex differs from the one implemented in ref. by an overall factor $`i`$, which ensures the hermiticity of the corresponding effective Lagrangian<sup>3</sup><sup>3</sup>3There is also a sign difference in the $`h_{3,4}`$ contributions coming from the different definition of $`ϵ^{0123}`$ in ref. . Furthermore, the $`i`$ factor modifies the interference pattern of anomalous coupling and Standard Model amplitudes: $`CP`$-violating $`h_{1,2}`$ contributions do not interfere with the Standard Model ones, whereas the $`CP`$-conserving $`h_{3,4}`$ do . Therefore, with the corrected vertex, there are contributions linear in $`h_{3,4}`$ and the cross section is generally not invariant anymore under a change of sign of $`h_{3,4}`$. This must be considered in a future precise analysis of anomalous couplings from experimental data since the limits for $`h_{3,4}`$ will not be symmetric at variance with present analyses.
The anomalous couplings spoil the gauge cancellation in the high energy limit and, therefore, will lead to violation of unitarity for increasing partonic center of mass energy $`\sqrt{\widehat{s}}`$. Usually, in an analysis of anomalous couplings from experimental data in hadronic collisions this problem is circumvented by supplementing the anomalous couplings, $`\alpha _{\mathrm{AC}}`$, with form factors. A common choice for the form factor is
$$\alpha _{\mathrm{AC}}\frac{\alpha _{\mathrm{AC}}}{(1+\frac{\widehat{s}}{\mathrm{\Lambda }^2})^n}$$
(5)
where $`n`$ has to be large enough to ensure unitarity and $`\mathrm{\Lambda }`$ is interpreted as the scale for new physics. Obviously, this procedure is rather ad hoc and introduces some arbitrariness . Also, it is not really consistent with the effective theory approach. Increasing anomalous contributions would require the inclusion of even higher dimensional operators. At the end of this section section we will address the question on how to avoid this arbitrariness in an analysis of anomalous couplings at hadron colliders.
Anomalous couplings mainly affect the events with large $`\widehat{s}`$ or large $`p_T`$. Since the total cross section is dominated by low $`p_T`$ events this is not a good observable to get tight constraints on anomalous couplings. A more promising possibility is to consider a double binned cross section. We therefore consider the total cross section binned in $`p_T^{\mathrm{}}`$ and $`p_T^\gamma `$ for the process $`ppW^+\gamma \mathrm{}^+\nu \gamma `$ at the LHC. In Figure 3 we show the ratio of $`\sigma _{\mathrm{AC}}`$ over $`\sigma _{\mathrm{SM}}`$ where $`\sigma _{\mathrm{SM}}`$ is the Standard Model cross section and $`\sigma _{\mathrm{AC}}`$ is the cross section obtained with $`\mathrm{\Delta }\kappa =0.08`$ and $`\lambda =0.02`$ (both within the present experimental limits from LEP and Tevatron ) and a form factor as defined in eq. (5) with $`n=2`$ and $`\mathrm{\Lambda }=`$ 2 TeV. As expected, the ratio is large for the high $`p_T`$ bins, whereas it is very close to one for the low $`p_T`$ bins. We checked that the uncertainty coming from the scale variation is much smaller than the effect of the anomalous couplings for the high $`p_T`$ bins.
Another possibility to get a large effect due to anomalous couplings is to consider the approximate radiation zeros present in the $`W\gamma `$ process . At tree-level in the Standard Model, the $`\mathrm{\Delta }\eta _{W\gamma }\eta _W\eta _\gamma `$ distribution has a radiation zero. This dip is filled by next-to-leading order corrections and anomalous effects. In order to get an observable quantity, we do not consider $`\mathrm{\Delta }\eta _{W\gamma }`$ but rather $`\mathrm{\Delta }\eta _\gamma \mathrm{}`$ . This will wash out the dip as the rapidity of the lepton is not equal to the rapidity of the $`W`$. However, requiring the energy (or the transverse momentum) of the lepton to be large enough forces the lepton to follow closely the $`W`$ direction. Also, applying a jet-veto reduces the effects of the next-to-leading order corrections. Thus, for larger $`p_T^{\mathrm{}}`$ even the next-to-leading order $`\mathrm{\Delta }\eta _\gamma \mathrm{}`$ distribution shows a clear dip in the Standard Model. This is illustrated in Figure 4, where we show the $`\mathrm{\Delta }\eta _\gamma \mathrm{}`$ distribution for the standard cuts and the additional cuts $`p_T^{\mathrm{}}>100`$ GeV and $`p_T^{\mathrm{}}>200`$ GeV respectively. In all figures we apply our jet-veto. Also shown are the three distributions with anomalous couplings, which are chosen as for Figure 3. Clearly, the effect is dramatic for high energy leptons but of course, the disadvantage of applying such a cut is a big loss in statistics.
We now turn to the question whether it is possible to avoid using form factors in the analysis of anomalous couplings at hadron colliders. This would bring these analyses more into line with those at $`e^+e^{}`$ colliders. In order to do so one should analyze the data at fixed values of $`\widehat{s}`$, as it is done at LEP. This results in limits for the anomalous parameters which are a function of $`\widehat{s}`$.
Obviously, it is possible to do such analysis for the production of $`Z\gamma `$ when both leptons are detected. Since the center of mass partonic energy can be reconstructed from the kinematics of the final state particles<sup>4</sup><sup>4</sup>4To simplify the discussion we assume that all final state particles are detected, including the jets. the cross section can be measured for different bins of fixed $`\widehat{s}`$ . As an example, we show in Figure 5 the cross section as a function of $`\widehat{s}`$. In order to enhance the effect of the anomalous couplings we do not only apply our standard cuts but we also require $`p_T^\gamma >200`$ GeV and $`p_T^Z>200`$ GeV. We show four curves: the curve (a) is the Standard Model result; curve (b) includes anomalous couplings in the standard way, that is with a form factor as defined in eq. (5) with $`n=3`$ ($`n=4`$) for $`h_3^V(h_4^V)`$ and $`\mathrm{\Lambda }=2`$ TeV and we have set $`h_3^\gamma =h_3^Z=0.01,h_4^\gamma =h_4^Z=10^4`$; curve (c) uses the same values for $`h_i^{Z/\gamma }`$ but does not include any form factor; finally curve (d) does also not include any form factors but the anomalous couplings are smaller $`h_3^\gamma =h_3^Z=0.001,h_4^\gamma =h_4^Z=10^5`$. For all of them we set $`h_1^V=h_2^V=0`$. Of course, for large $`\widehat{s}`$ the effects are much more dramatic if the form factor is omitted and at some point the corresponding curves would violate unitarity. This simply reflects the breakdown of the effective theory approach in this region. The idea behind Figure 5 is that such an analysis can be done for suitably defined bins in $`\widehat{s}`$. As a result, for each bin, i.e. each value of $`\widehat{s}`$ a bound on the anomalous couplings can be obtained.
In the inset plot we present again curve (b) and the one corresponding to the parameters of (b) but with the opposite sign for values of the anomalous couplings $`h_3^\gamma =h_3^Z=0.01,h_4^\gamma =h_4^Z=10^4`$. From there the effect of the interference between AC and SM results can be observed, i.e. the contribution of linear terms in $`h_3^V`$ and $`h_4^V`$ appearing in the squared amplitudes due to the correct treatment of the $`i`$ factor in the $`Z\gamma V`$ vertex . For this particular configuration, the interference effects are mostly relevant at $`\sqrt{\widehat{s}}<1`$ TeV and modify the cross section by more than 10%. Clearly, they must be taken into account in a precise extraction of anomalous couplings from future experimental data. Results for $`h_{1,2}^V`$ couplings are similar to those obtained for the same values of $`h_{3,4}^V`$. The only difference comes from the fact that the $`CP`$-violating couplings do not interfere with the SM. For a configuration like (b) with $`h_1^\gamma =h_1^Z=0.01,h_2^\gamma =h_2^Z=10^4`$, the cross section is given by the average of both curves in the inset plot of Figure 5.
The situation is more complicated for $`W\gamma `$ production since the neutrino is not observed. Nevertheless, by identifying the transverse momentum of the neutrino with the missing transverse momentum, and assuming the $`W`$ boson to be on shell, it is possible to reconstruct the neutrino kinematics (particularly the longitudinal momentum) with a twofold ambiguity. In the case of the Tevatron, since it is a $`p\overline{p}`$ collider, it is possible to choose the ‘correct’ neutrino kinematics 73% of the times by selecting the maximum (minimum) of the two reconstructed values for the longitudinal momentum of the neutrino for $`W^+\gamma `$ ($`W^{}\gamma `$.
This is not true at the LHC, where due to the symmetry of the colliding beams both reconstructed kinematics have equal chances to be correct. Fortunately, in the case of anomalous couplings, we are interested in a efficient way to reconstruct the $`\widehat{s}`$ rather than the full kinematics. Again there are two possible values of $`\widehat{s}`$. It turns out that there is a simple method to choose the ‘correct’ one 66% of the times at the LHC (73% of the times at Tevatron) by selecting the minimum $`\widehat{s}_{\mathrm{min}}`$ of the two reconstructed values. This applies to both $`W^+\gamma `$ and $`W^{}\gamma `$ production. Furthermore, we checked that the selected value $`\widehat{s}_{\mathrm{min}}`$ differs in almost 90% of the events by less than 10% from the exact value $`\widehat{s}`$. This is likely to be precise enough, since the data will be collected in sizeable bins of $`\widehat{s}`$ and the anomalous parameters are not expected to change very rapidly as a function of the energy. To quantify the advantage of the method, we show in Figures 7 and 7 the correlations of $`\sqrt{\widehat{s}_{\mathrm{min}}}`$ with $`\sqrt{\widehat{s}}`$. The cross section drops very rapidly for increasing $`\sqrt{\widehat{s}}\sqrt{\widehat{s}_{\mathrm{min}}}`$. This correlation also holds in the particularly interesting large $`\sqrt{\widehat{s}}`$ region and also for the anomalous contribution. To investigate the latter point, we show in Figure 7 the correlation for (already experimentally ruled out) huge values of $`\mathrm{\Delta }\kappa =0.8`$ and $`\lambda =0.2`$. For this figure we still use an ordinary form factor but in order to increase the anomalous contribution further we set $`\mathrm{\Lambda }=1`$ TeV.
Finally, in Figure 8 we show the same correlation for $`W\gamma `$ production at the Tevatron (Run II), which corresponds to $`p\overline{p}`$ scattering at $`\sqrt{s}=2`$ TeV. We impose the same kinematical cuts as for the LHC, with the following exceptions: for the transverse momentum of the photon we require $`p_T^\gamma >10`$ GeV and the rapidities of the observed lepton has to be in the range $`|\eta |<1.5`$ (instead of 2.4). As for the LHC, there is a strong correlation between the ’true’ center of mass energy $`\sqrt{\widehat{s}}`$ and the ’reconstructed’ one $`\sqrt{\widehat{s}_{\mathrm{min}}}`$. Again, this correlation also holds in the large $`\sqrt{\widehat{s}}`$ region. In Figure 8 no anomalous couplings are included, but we have checked that also for the Tevatron the inclusion of large anomalous couplings does not spoil the correlations.
As a result of these investigations we conclude that even in the case of $`W\gamma `$ production reliable bounds for anomalous couplings as a function of $`\widehat{s}`$ can be obtained. To this end one merely has to do an analysis as in the $`Z\gamma `$ case but with $`\widehat{s}_{\mathrm{min}}`$ replacing $`\widehat{s}`$. This has the advantage that the anomalous effects can be quantized without introducing the ambiguity of form factors. Such a procedure would certainly facilitate a comparison of various bounds from different experiments. Finally, we note that in principle any quantity which has a very strong correlation with $`\widehat{s}`$ can be used. However, we could not find any better candidate than $`\widehat{s}_{\mathrm{min}}`$. In particular, the correlations of $`\widehat{s}`$ with the cluster mass and transverse mass respectively is not quite as strong .
## 5 Conclusions
In this work, we have presented a general purpose Monte Carlo program for the calculation of any infrared safe observable in $`W\gamma `$ and $`Z\gamma `$ production at hadron colliders at next-to-leading order in $`\alpha _s`$. The leptonic decay of the $`W`$ and $`Z`$-boson respectively has been included in the narrow-width approximation. We retained all spin information via decay angle correlations and thereby generalized previous calculations . We also included anomalous triple gauge boson couplings at NLO in $`\alpha _s`$ and presented the analytical expressions for the corresponding amplitudes.
As an illustration of the usefulness of the program, we have studied several observables for the LHC. Generally we find that the NLO corrections are relevant for all of them, confirming results of refs. .
Moreover, we searched for the kinematical regions were the effect of anomalous couplings is amplified and proposed an alternative way to study its energy dependence. Using the strong correlations between the partonic center of mass energy and a measurable variable, $`\widehat{s}_{\mathrm{min}}`$, makes it possible to extract anomalous couplings from the data without need to introduce ad hoc form factors. Such an analysis is possible even for the $`W\gamma `$ process and, in our view, should be undertaken at the Tevatron and LHC.
## Acknowledgments
It is a pleasure to thank Z. Kunszt for his participation on the initial stages of this work and for helpful discussions. D. de F. would like to thank the Department of Physics of the University of Durham for its hospitality while part of this work was carried out.
This work was partly supported by the EU Fourth Framework Programme ‘Training and Mobility of Researchers’, Network ‘Quantum Chromodynamics and the Deep Structure of Elementary Particles’, contract FMRX-CT98-0194 (DG 12 - MIHT).
## 6 Appendix
The helicity amplitudes that are needed for the calculation of $`W\gamma `$ and $`Z\gamma `$ production at next-to-leading order can be found in ref. . In this appendix we list the additional amplitudes that are needed in the presence of anomalous couplings which lead to the non-standard vertices as given in eqs. (3) and (4). We use the notation and conventions of refs. .
We start with the $`W\gamma `$ amplitudes. In order to maintain electromagnetic gauge invariance we set $`g_1^\gamma =0`$ and only allow anomalous couplings $`\mathrm{\Delta }\kappa ^\gamma `$ and $`\lambda ^\gamma `$. The only diagram that gets modified by these anomalous couplings is the diagram with a $`WW\gamma `$ coupling. Fortunately, this diagram does not contribute to the rather complicated finite pieces $`F_\gamma ^a,F_\gamma ^b`$ of the amplitudes, given in eqs. (4.6) and (4.7) of ref. . Therefore, only the tree-level amplitudes have to be computed with anomalous couplings.
We do not have to list the explicit results for all possible helicities of the photon and the gluon. In order to reverse the helicities of all gauge bosons, i.e. the photon and the gluon (if the latter is present), we merely have to apply a ‘flip’ operation, defined as
$$\mathrm{flip2}:12;34;ab\left[ab\right];\mathrm{and}A^aA^b(\mathrm{for}W\gamma )$$
(6)
The amplitudes $`A_\gamma ^{\mathrm{tree},\mathrm{a}},A_\gamma ^{\mathrm{tree},\mathrm{b}}`$ given in eqs. (4.4) and (4.5) of ref. are modified as follows in the presence of anomalous couplings:
$`A_{\mathrm{AC},\gamma }^{\mathrm{tree},\mathrm{a}}`$ $`=`$ $`A_\gamma ^{\mathrm{tree},\mathrm{a}}+A_{\mathrm{AC}}^5`$ (7)
$`A_{\mathrm{AC},\gamma }^{\mathrm{tree},\mathrm{b}}`$ $`=`$ $`A_\gamma ^{\mathrm{tree},\mathrm{b}}A_{\mathrm{AC}}^5`$ (8)
where
$$A_{\mathrm{AC}}^5=\frac{i\left[45\right]}{2s_{34}(s_{12}s_{34})\left[34\right]}\left((\mathrm{\Delta }\kappa ^\gamma +\lambda ^\gamma )13\left[25\right]\left[34\right]+\lambda ^\gamma 1|5|2\left[45\right]\right)$$
(9)
As usual, this will also lead to a modification of the one-loop amplitudes. The corresponding divergent pieces now read $`c_\mathrm{\Gamma }VA_{\mathrm{AC},\gamma }^{\mathrm{tree},\mathrm{a}}`$ and $`c_\mathrm{\Gamma }VA_{\mathrm{AC},\gamma }^{\mathrm{tree},\mathrm{b}}`$ respectively, where
$$V=\frac{1}{ϵ^2}\left(\frac{\mu ^2}{s_{12}}\right)^ϵ\frac{3}{2ϵ}\left(\frac{\mu ^2}{s_{12}}\right)^ϵ\frac{7}{2}.$$
(10)
In the case of the bremsstrahlung amplitudes we have to consider two cases. The additional gluon can have positive or negative helicity. The corresponding Standard Model amplitudes are given in eqs. (4.9) to (4.12) of ref. . In the presence of anomalous couplings $`\mathrm{\Delta }\kappa `$ and $`\lambda `$ they are modified as follows.
$`A_{\mathrm{AC},6,\gamma }^{\mathrm{tree},\mathrm{a}}(1^{},2^+,3^{},4^+,5^+,6^+)`$ $`=`$ $`A_{6,\gamma }^{\mathrm{tree},\mathrm{a}}(1^{},2^+,3^{},4^+,5^+,6^+)+A_{\mathrm{AC}}^{6+}`$ (11)
$`A_{\mathrm{AC},6,\gamma }^{\mathrm{tree},\mathrm{a}}(1^{},2^+,3^{},4^+,5^{},6^+)`$ $`=`$ $`A_{6,\gamma }^{\mathrm{tree},\mathrm{a}}(1^{},2^+,3^{},4^+,5^{},6^+)+A_{\mathrm{AC}}^6`$ (12)
$`A_{\mathrm{AC},6,\gamma }^{\mathrm{tree},\mathrm{b}}(1^{},2^+,3^{},4^+,5^+,6^+)`$ $`=`$ $`A_{6,\gamma }^{\mathrm{tree},\mathrm{b}}(1^{},2^+,3^{},4^+,5^+,6^+)A_{\mathrm{AC}}^{6+}`$ (13)
$`A_{\mathrm{AC},6,\gamma }^{\mathrm{tree},\mathrm{b}}(1^{},2^+,3^{},4^+,5^{},6^+)`$ $`=`$ $`A_{6,\gamma }^{\mathrm{tree},\mathrm{b}}(1^{},2^+,3^{},4^+,5^{},6^+)A_{\mathrm{AC}}^6`$ (14)
where we defined
$`A_{\mathrm{AC}}^{6+}`$ $``$ $`i{\displaystyle \frac{\left[45\right]1|2+6|5}{2s_{34}^2(t_{126}s_{34})1662}}\left(\lambda ^\gamma 341545+(\mathrm{\Delta }\kappa ^\gamma +\lambda ^\gamma )13s_{34}\right)`$ (15)
$`A_{\mathrm{AC}}^6`$ $``$ $`i{\displaystyle \frac{3515}{2s_{34}^2(t_{126}s_{34})1662}}\times `$ (16)
$`\left(\lambda ^\gamma 351|2+6|5\left[34\right](\mathrm{\Delta }\kappa ^\gamma +\lambda ^\gamma )1|2+6|4s_{34}\right)`$
We now turn to the amplitudes for $`Z\gamma `$ production. In the Standard Model, there are no diagrams with a triple gauge boson coupling and the corresponding amplitudes $`A_\gamma ^{\mathrm{tree},\mathrm{s}}`$ can simply be obtained as the symmetric combination of the $`W\gamma `$ amplitudes
$$A_\gamma ^{\mathrm{tree},\mathrm{s}}=A_\gamma ^{\mathrm{tree},\mathrm{a}}+A_\gamma ^{\mathrm{tree},\mathrm{b}}$$
(17)
(see eq. (4.15) of ref. ). These combinations can be simplified somewhat but we refrain from listing the simplified versions.
The amplitudes related to an anomalous $`ZZ\gamma `$ or $`Z\gamma \gamma `$ coupling will be denoted by $`A_{5/6,\mathrm{AC}}^{(Z\gamma )}`$. In the former case, the intermediate vector boson is a $`Z`$-boson whereas in the latter case it is a $`\gamma `$. This results in different couplings of the intermediate vector boson to the initial state quarks. Apart from this difference, the amplitudes with an intermediate $`Z`$ and $`\gamma `$ are the same. The anomalous couplings always appear in the combination
$$\stackrel{~}{h}_1^{Z/\gamma }\frac{h_1^{Z/\gamma }}{M_Z^2};\stackrel{~}{h}_2^{Z/\gamma }\frac{h_2^{Z/\gamma }}{M_Z^4};\stackrel{~}{h}_3^{Z/\gamma }\frac{h_3^{Z/\gamma }}{M_Z^2};\stackrel{~}{h}_4^{Z/\gamma }\frac{h_4^{Z/\gamma }}{M_Z^4}.$$
(18)
We start with the tree-level amplitude for a positive helicity photon
$$A_{5,\mathrm{AC}}^{(Z\gamma )}=\frac{i}{4s_{34}}\left((i\stackrel{~}{h}_1^{Z/\gamma }+\stackrel{~}{h}_3^{Z/\gamma })213\left[25\right]\left[45\right]+(i\stackrel{~}{h}_2^{Z/\gamma }+\stackrel{~}{h}_4^{Z/\gamma })123|5|4\left[25\right]^2\right)$$
(19)
In order to obtain the amplitudes for a negative helicity photon we have to apply the ‘flip2’ operation defined in eq. (6) in addition to
$$\stackrel{~}{h}_1^{Z/\gamma }\stackrel{~}{h}_1^{Z/\gamma };\stackrel{~}{h}_2^{Z/\gamma }\stackrel{~}{h}_2^{Z/\gamma };$$
(20)
The anomalous bremsstrahlung amplitude with an additional positive helicity gluon reads
$`A_{6,\mathrm{AC}}^{(Z\gamma )}`$ $`=`$ $`{\displaystyle \frac{i}{4s_{34}1626}}\times `$
$`\left((i\stackrel{~}{h}_1^{Z/\gamma }+\stackrel{~}{h}_3^{Z/\gamma })213\left[45\right]1|2+6|5+(i\stackrel{~}{h}_2^{Z/\gamma }+\stackrel{~}{h}_4^{Z/\gamma })3|5|41|2+6|5^2\right)`$
whereas for a negative helicity gluon
$`A_{6,\mathrm{AC}}^{(Z\gamma )}`$ $`=`$ $`{\displaystyle \frac{i}{4s_{34}\left[16\right]\left[26\right]}}\times `$
$`\left((i\stackrel{~}{h}_1^{Z/\gamma }+\stackrel{~}{h}_3^{Z/\gamma })2\left[45\right]\left[25\right]3|1+6|2+(i\stackrel{~}{h}_2^{Z/\gamma }+\stackrel{~}{h}_4^{Z/\gamma })t_{126}3|5|4\left[25\right]^2\right)`$
Again, the operation ‘flip2’ reverses the helicities of the photon and the gluon. Finally we mention that in order to get the amplitudes with a positive helicity lepton, $`3^+`$, we simply have to exchange $`34`$ in the amplitudes presented above. Correspondingly, the amplitudes with opposite helicities for the partons are obtained by a simple $`12`$ crossing.
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# Polarizations on abelian varieties
## 1. Introduction
If $`F`$ is a field and $`\mathrm{}`$ is a prime number, let $`F^{(\mathrm{})}`$ denote the extension of $`F`$ obtained by adjoining all $`\mathrm{}`$-power roots of unity. In §6 we prove a result (Theorem 6.1) which implies:
###### Theorem 1.1.
Suppose $`F`$ is a number field and $`r`$ is a positive integer. Then there exists an abelian variety $`A`$ over $`F`$ with the following properties:
1. every $`F`$-polarization on every abelian variety $`F`$-isogenous to $`A`$ has degree divisible by $`r`$;
2. if $`\mathrm{}`$ is a prime and $`\mathrm{}^a`$ divides $`r`$, then every $`F^{(\mathrm{})}`$-polarization on every abelian variety $`F^{(\mathrm{})}`$-isogenous to $`A`$ has degree divisible by $`\mathrm{}^{2a}`$.
We deduce Theorems 6.1 and 1.1 from Theorem 5.2, which gives a general method for obtaining isogeny classes of abelian varieties all of whose polarizations have degree divisible by a given prime power. The proof uses some results on elliptic curves (see Corollary 4.7) that may be of independent interest.
It is known that every isogeny class over an algebraically closed field contains a principally polarized abelian variety (, Cor. 1 to Thm. 4 in §23).
Howe (; see also ) gave examples of isogeny classes of abelian varieties over finite fields with no principal polarizations (but not with the degrees of all the polarizations divisible by a given non-zero integer, as in Theorem 1.1 above).
In we obtained, for all odd primes $`\mathrm{}`$, isogeny classes of abelian varieties in positive characteristic, all of whose polarizations have degree divisible by $`\mathrm{}^2`$. (We gave results in the more general context of invertible sheaves; see also Theorems 6.1 and 5.2 below.) Our results gave the first examples for which all the polarizations of the abelian varieties in an isogeny class have degree divisible by a given prime.
Inspired by our results, Howe recently obtained, for all odd primes $`\mathrm{}`$, examples of isogeny classes of abelian varieties over fields of arbitrary characteristic (including over number fields), all of whose polarizations have degree divisible by $`\mathrm{}^2`$. The results of this paper were inspired by an early version (§2 of the current version) of . In particular, we employ Howe’s idea of considering the Jordan-Hölder factors of an appropriate module.
Zarhin was partially supported by EPSRC grant GR/M98135. Silverberg would like to thank NSF, NSA, MSRI, the Mathematics Institute of the University of Erlangen, and the Alexander von Humboldt Foundation.
## 2. Notation
If $`T`$ is a $`𝐙_{\mathrm{}}`$-lattice and $`e:T\times T𝐙_{\mathrm{}}`$ is a non-degenerate pairing, let
$$T^{}=\{yT_𝐙_{\mathrm{}}𝐐_{\mathrm{}}e(T,y)𝐙_{\mathrm{}}\},$$
the dual lattice of $`T`$ with respect to $`e`$.
Suppose $`F`$ is a field. Let $`F^s`$ be a separable closure of $`F`$ and let $`G_F=\mathrm{Gal}(F^s/F)`$. If $`A`$ is an abelian variety over $`F`$ and $`n`$ is a positive integer, let $`A_n`$ denote the kernel of multiplication by $`n`$ in $`A(F^s)`$. Let $`T_{\mathrm{}}(A)=\underset{}{lim}A_\mathrm{}^m`$, let $`V_{\mathrm{}}(A)=T_{\mathrm{}}(A)_𝐙_{\mathrm{}}𝐐_{\mathrm{}}`$, let $`\overline{A}=A\times _FF^s`$, let $`A^{}`$ denote the dual of $`A`$, and let
$$\rho _{A,n}:G_F\mathrm{Aut}(A_n)\mathrm{GL}_{2d}(𝐙/n𝐙)$$
denote the mod $`n`$ representation (where $`d=\mathrm{dim}(A)`$). Let
$$F^{(\mathrm{})}=F(𝝁_{\mathrm{}^{\mathrm{}}})\text{and}F_A^{(\mathrm{})}=F(A_{\mathrm{}^{\mathrm{}}}).$$
## 3. Algebra
###### Lemma 3.1.
Suppose $`N`$ is a group, $`S`$ is a $`𝐙_{\mathrm{}}[N]`$-module which is free of finite rank $`m`$ over $`𝐙_{\mathrm{}}`$, and $`W=S_𝐙_{\mathrm{}}𝐐_{\mathrm{}}`$. Then the following are equivalent:
1. the natural homomorphism $`𝐙_{\mathrm{}}[N]\mathrm{End}_𝐙_{\mathrm{}}(S)`$ is surjective;
2. the $`N`$-module $`S/\mathrm{}S`$ is absolutely simple;
3. the image of the natural homomorphism $`𝐙_{\mathrm{}}[N]\mathrm{End}_𝐐_{\mathrm{}}(W)`$ is isomorphic to the matrix ring $`\mathrm{M}_m(𝐙_{\mathrm{}})`$.
If these equivalent conditions hold, then:
1. the $`N`$-module $`W`$ is absolutely simple;
2. every $`N`$-stable $`𝐙_{\mathrm{}}`$-lattice in $`W`$ is of the form $`\mathrm{}^iS`$ for some integer $`i`$;
3. if $`f:S\times S𝐙_{\mathrm{}}`$ is an $`N`$-invariant pairing such that $`f(S,S)=𝐙_{\mathrm{}}`$ then $`f`$ is perfect.
###### Proof.
Assume (ii). By Theorem 9.2 of , the natural homomorphism $`𝐅_{\mathrm{}}[N]\mathrm{End}(S/\mathrm{}S)`$ is surjective. Now (i) follows from Nakayama’s Lemma.
Clearly, (i) implies (iii).
Assume (iii). Then the natural homomorphism $`𝐅_{\mathrm{}}[N]\mathrm{End}(S/\mathrm{}S)`$ factors through a surjective homomorphism $`𝐅_{\mathrm{}}[N]\mathrm{M}_m(𝐅_{\mathrm{}})`$, i.e., $`S/\mathrm{}S`$ is an $`\mathrm{M}_m(𝐅_{\mathrm{}})`$-module. Since $`\mathrm{M}_m(𝐅_{\mathrm{}})`$ is a simple algebra and $`m=\mathrm{dim}_𝐅_{\mathrm{}}(S/\mathrm{}S)`$, the map $`\mathrm{M}_m(𝐅_{\mathrm{}})\mathrm{End}(S/\mathrm{}S)`$ is injective and therefore surjective. Thus the composition $`𝐅_{\mathrm{}}[N]\mathrm{M}_m(𝐅_{\mathrm{}})\mathrm{End}(S/\mathrm{}S)`$ is surjective, and (ii) follows.
If (i) holds, then the natural map $`𝐐_{\mathrm{}}[N]\mathrm{End}_𝐐_{\mathrm{}}(W)`$ is surjective by dimension arguments, and (a) follows.
Part (b) is an immediate corollary of (ii) and Exercise 15.3 of .
The form $`f`$ in (c) induces an $`N`$-invariant pairing $`\overline{f}:S/\mathrm{}S\times S/\mathrm{}S𝐅_{\mathrm{}}`$. The left and right kernels of $`\overline{f}`$ are $`N`$-submodules of the simple $`N`$-module $`S/\mathrm{}S`$, and thus are $`0`$, since $`\overline{f}0`$. It follows that $`f`$ is perfect. ∎
###### Definition 3.2.
We will say a $`𝐙_{\mathrm{}}[N]`$-module $`S`$ is well-rounded if $`S`$, $`N`$, and $`\mathrm{}`$ satisfy the hypotheses and the equivalent conditions (i)-(iii) of Lemma 3.1.
###### Remark 3.3.
Clearly, if $`S_1`$ is a well-rounded $`𝐙_{\mathrm{}}[N_1]`$-module and $`S_2`$ is a well-rounded $`𝐙_{\mathrm{}}[N_2]`$-module, then $`S_1_𝐙_{\mathrm{}}S_2`$ is a well-rounded $`𝐙_{\mathrm{}}[N_1\times N_2]`$-module.
###### Lemma 3.4.
Suppose $`S`$ is a well-rounded $`𝐙_{\mathrm{}}[N]`$-module, $`W=S_𝐙_{\mathrm{}}𝐐_{\mathrm{}}`$, and $`f`$ is a non-degenerate $`N`$-invariant $`𝐐_{\mathrm{}}`$-valued alternating pairing on $`W`$. Suppose $`D`$ is a group, and $`V`$ is a finite-dimensional $`𝐐_{\mathrm{}}`$-vector space and a simple $`D`$-module. Then every ($`N\times D`$)-stable $`𝐙_{\mathrm{}}`$-lattice in $`WV`$ is of the form $`S\mathrm{\Gamma }`$ where $`\mathrm{\Gamma }`$ is a $`D`$-stable $`𝐙_{\mathrm{}}`$-lattice in $`V`$, and every ($`N\times D`$)-invariant $`𝐐_{\mathrm{}}`$-valued pairing on $`WV`$ is of the form $`fh`$ where $`h`$ is a $`D`$-invariant $`𝐐_{\mathrm{}}`$-valued pairing on $`W`$.
###### Proof.
Let $`r=\mathrm{dim}(V)`$ and $`2d=\mathrm{dim}(W)`$. Let $`M`$ denote the image of $`𝐙_{\mathrm{}}[N]\mathrm{End}_𝐐_{\mathrm{}}(W)`$. By Lemma 3.1iii, $`M=\mathrm{End}_𝐙_{\mathrm{}}(S)\mathrm{M}_{2d}(𝐙_{\mathrm{}})`$. Suppose $`\mathrm{\Lambda }`$ is an ($`N\times D`$)-stable lattice in $`WV`$. Then $`\mathrm{\Lambda }`$ is a $`𝐙_{\mathrm{}}[N]`$-module and thus an $`M`$-module. Since every finitely generated $`M`$-module is a direct sum of copies of $`𝐙_{\mathrm{}}^{2d}`$ ($`S`$) with the standard action, therefore $`\mathrm{\Lambda }S^r=S_𝐙_{\mathrm{}}𝐙_{\mathrm{}}^r`$, as $`M`$-modules. Since $`D`$ commutes with $`N`$, it commutes with the image of $`𝐙_{\mathrm{}}[N]`$ in $`\mathrm{End}(\mathrm{\Lambda })`$, i.e., with $`\mathrm{End}(S)1`$. Therefore the image of $`D`$ in $`\mathrm{Aut}(\mathrm{\Lambda })`$ is contained in $`1\mathrm{Aut}(𝐙_{\mathrm{}}^r)`$, and the isomorphism $`\mathrm{\Lambda }S_𝐙_{\mathrm{}}𝐙_{\mathrm{}}^r`$ is also an isomorphism of $`D`$-modules. Therefore $`\mathrm{\Lambda }(𝐙_{\mathrm{}}^r)^{2d}`$ as $`D`$-modules, so $`V^{2d}WV\mathrm{\Lambda }𝐐_{\mathrm{}}(𝐐_{\mathrm{}}^r)^{2d}`$ as $`D`$-modules. Since $`V`$ is a simple $`D`$-module, therefore $`V^{2d}`$ and $`(𝐐_{\mathrm{}}^r)^{2d}`$ are semisimple. It follows that $`V𝐐_{\mathrm{}}^r`$ as $`D`$-modules, and we can identify $`𝐙_{\mathrm{}}^r`$ with a $`D`$-stable lattice $`\mathrm{\Gamma }`$ in the $`D`$-module $`V`$.
If $`g`$ is an ($`N\times D`$)-invariant $`𝐐_{\mathrm{}}`$-valued pairing on $`WV`$, then $`WVV^{2d}`$ is self-dual with respect to $`g`$, since $`V`$ is simple as a $`D`$-module. Therefore $`V`$ is self-dual, i.e., admits a non-degenerate $`D`$-invariant $`𝐐_{\mathrm{}}`$-valued pairing $`h_0`$. Since $`W`$ is absolutely simple as an $`N`$-module, every $`N`$-invariant $`𝐐_{\mathrm{}}`$-valued pairing on $`W`$ is a multiple of $`f`$. Similarly, since $`\mathrm{End}_{N\times D}(WV)=\mathrm{End}_D(V)`$, thus every ($`N\times D`$)-invariant $`𝐐_{\mathrm{}}`$-valued pairing on $`WV`$ is of the form $`fh`$, where for some $`u\mathrm{Aut}_D(V)`$ we have $`h(x,y)=h_0(ux,y)`$ for every $`x,yV`$. Since $`h_0`$ is $`D`$-invariant, so is $`h`$. ∎
The next result extends Theorem 6.2 of . In that result, $`N`$ was restricted to being a finite group of order not divisible by $`\mathrm{}`$, $`\mathrm{}`$ was odd, and $`D`$ was $`𝝁_{\mathrm{}}`$.
###### Theorem 3.5.
Suppose $`S`$ is a well-rounded $`𝐙_{\mathrm{}}[N]`$-module, $`W=S_𝐙_{\mathrm{}}𝐐_{\mathrm{}}`$, and $`f`$ is a non-degenerate $`N`$-invariant $`𝐐_{\mathrm{}}`$-valued alternating pairing on $`W`$. Suppose $`D\mathrm{GL}_m(𝐐_{\mathrm{}})`$ is an irreducible subgroup. Assume that there does not exist a $`D`$-stable $`𝐙_{\mathrm{}}`$-lattice $`\mathrm{\Gamma }`$ in $`𝐐_{\mathrm{}}^m`$ with a perfect $`D`$-invariant $`𝐙_{\mathrm{}}`$-valued symmetric pairing. Suppose $`T`$ is an ($`N\times D`$)-stable $`𝐙_{\mathrm{}}`$-lattice in $`W^m`$ and $`e`$ is a non-degenerate ($`N\times D`$)-invariant $`𝐙_{\mathrm{}}`$-valued alternating pairing on $`T`$. Then $`e`$ is not perfect, and $`\mathrm{\#}(T^{}/T)=\mathrm{}^{2dt}`$ for some positive integer $`t`$, where $`T^{}`$ is the dual lattice of $`T`$ with respect to $`e`$ and $`2d=\mathrm{dim}_𝐐_{\mathrm{}}(W)`$.
###### Proof.
Replacing $`f`$ by a suitable multiple, we may suppose that $`f(S,S)=𝐙_{\mathrm{}}`$. By Lemma 3.4, we have $`T=S\mathrm{\Gamma }`$ where $`\mathrm{\Gamma }`$ is a $`D`$-stable lattice in $`𝐐_{\mathrm{}}^m`$, and $`e=fh`$ where $`h`$ is a non-degenerate $`D`$-invariant $`𝐐_{\mathrm{}}`$-valued pairing on $`𝐐_{\mathrm{}}^m`$. Since $`e`$ is alternating, $`h`$ is symmetric. Since $`f(S,S)=𝐙_{\mathrm{}}`$, we have $`h(\mathrm{\Gamma },\mathrm{\Gamma })𝐙_{\mathrm{}}`$. Since $`h:\mathrm{\Gamma }\times \mathrm{\Gamma }𝐙_{\mathrm{}}`$ is not perfect, there exists $`z\mathrm{\Gamma }\mathrm{}\mathrm{\Gamma }`$ such that $`h(z,\mathrm{\Gamma })\mathrm{}𝐙_{\mathrm{}}`$. If $`xS\mathrm{}S`$, then $`xzT\mathrm{}T`$ but $`e(xz,T)\mathrm{}𝐙_{\mathrm{}}.`$ Therefore $`e`$ is not perfect, so $`T^{}T`$. The Jordan-Hölder factors of the $`N`$-module $`T/\mathrm{}T(S/\mathrm{}S)^m`$ are all isomorphic to $`S/\mathrm{}S`$. Since $`TT^{}\mathrm{}^rT`$ for some positive integer $`r`$, the Jordan-Hölder factors of the $`N`$-module $`T^{}/T`$ are all isomorphic to $`S/\mathrm{}S`$. Therefore $`\mathrm{\#}(T^{}/T)`$ is a power of $`\mathrm{\#}(S/\mathrm{}S)=\mathrm{}^{2d}`$ as desired. ∎
The symmetric group $`S_n`$ acts on $`𝐐_{\mathrm{}}^n`$ with standard basis $`\{e_1,\mathrm{},e_n\}`$. Let
$$U=\{\underset{i=1}{\overset{n}{}}c_ie_i\underset{i=1}{\overset{n}{}}c_i=0\},$$
a hyperplane in $`𝐐_{\mathrm{}}^n`$. Write $`h`$ for the restriction to $`U`$ of the standard symmetric pairing on $`𝐐_{\mathrm{}}^n`$. Then $`h`$ is non-degenerate. There are two well-known examples of $`S_n`$-stable $`𝐙_{\mathrm{}}`$-lattices in $`U`$. First, the root lattice $`Q`$ generated by “simple roots” $`\alpha _i=e_ie_{i+1}`$ for $`1in1`$. Second, its dual with respect to $`h`$, the weight lattice $`P`$ generated by $`\omega _i=e_i(e_1+\mathrm{}+e_n)/n`$ for $`1in`$. We have $`h(\omega _i,\alpha _j)=\delta _{ij}`$. If $`t_iS_n`$ is the transposition $`(i,i+1)`$ then $`t_i(x)=xh(x,\alpha _i)\alpha _i`$ for $`xU`$.
###### Lemma 3.6 (see and Theorem 1 of ).
1. $`QP`$.
2. $`P/Q`$ is a cyclic group whose order is the largest power of $`\mathrm{}`$ dividing $`n`$.
3. If $`L`$ is an $`S_n`$-stable $`𝐙_{\mathrm{}}`$-lattice in $`U`$, then there exists an integer $`t`$ such that $`Q\mathrm{}^tLP`$.
4. If $`n>2`$ and $`n`$ is divisible by $`\mathrm{}`$, then there are no $`S_n`$-invariant perfect $`𝐙_{\mathrm{}}`$-valued pairings on $`P`$ or on $`Q`$.
5. If $`n>2`$, and $`\mathrm{}`$ divides $`n`$ but $`\mathrm{}^2`$ does not divide $`n`$, then there does not exist an $`S_n`$-stable $`𝐙_{\mathrm{}}`$-lattice in $`U`$ with a perfect $`S_n`$-invariant $`𝐙_{\mathrm{}}`$-valued pairing.
6. If $`\mathrm{}`$ does not divide $`n`$, then $`Q=P`$ and the $`𝐙_{\mathrm{}}[S_n]`$-module $`P`$ is well-rounded.
###### Proof.
Since $`\alpha _i=\omega _i\omega _{i+1}`$, we have (i). The group isomorphism $`P/Q𝐙_{\mathrm{}}/n𝐙_{\mathrm{}}`$ defined by $`\omega _i1`$ gives (ii). For (iii), multiplying $`L`$ by a suitable power of $`\mathrm{}`$, we may assume that $`L`$ contains $`Q`$ but does not contain $`\mathrm{}^1Q`$. Assume that $`L`$ is not contained in $`P`$. Then there exist $`xL`$ and $`j\{1,\mathrm{},n1\}`$ such that $`h(x,\alpha _j)𝐙_{\mathrm{}}`$. Multiplying $`x`$ by a suitable $`\mathrm{}`$-adic integer, we may assume that $`h(x,\alpha _j)=\mathrm{}^1`$. Since $`xL`$ and $`t_j(x)=x\mathrm{}^1\alpha _j`$, therefore $`\alpha _j/\mathrm{}L`$. Since the orbit $`S_n(\alpha _j)`$ contains all the roots and $`L`$ is $`S_n`$-stable, thus $`\mathrm{}^1QL`$. This contradiction gives (iii). The $`𝐐_{\mathrm{}}[S_n]`$-module $`U`$ is absolutely simple, so every $`S_n`$-invariant $`𝐐_{\mathrm{}}`$-valued pairing on $`U`$ is of the form $`ch(x,y)`$ for some $`c𝐐_{\mathrm{}}`$. We have $`h(Q,Q)=𝐙_{\mathrm{}}`$, since $`h(\alpha _1,\alpha _2)=1`$. Further, $`ch(Q,Q)=𝐙_{\mathrm{}}`$ if and only if $`c𝐙_{\mathrm{}}^{}`$. Suppose that $`n`$ is divisible by $`\mathrm{}`$. Since the dual of $`Q`$ with respect to $`h`$ is $`PQ`$, thus $`ch(x,y)`$ is never perfect. Let $`\mathrm{}^r`$ be the largest power of $`\mathrm{}`$ dividing $`n`$. Since $`h(\omega _1,\omega _1)=(n1)/n`$, we have $`\mathrm{}^r(P,P)=𝐙_{\mathrm{}}`$. Further, $`ch(P,P)=𝐙_{\mathrm{}}`$ if and only if $`c\mathrm{}^r𝐙_{\mathrm{}}^{}`$. The dual of $`P`$ with respect to $`\mathrm{}^rh`$ is $`\mathrm{}^rQ`$. If $`n>2`$, then $`\mathrm{}^rQP`$ so $`ch(x,y)`$ is perfect for no $`c`$, and we have (iv). Part (v) follows from (ii), (iii), and (iv). Suppose that $`n`$ is not divisible by $`\mathrm{}`$. Let bars denote reduction mod $`\mathrm{}`$. By (ii), $`P=Q`$. By (iii), every $`S_n`$-stable $`𝐙_{\mathrm{}}`$-lattice in $`U`$ is of the form $`\mathrm{}^tP`$. Therefore the $`S_n`$-module $`\overline{P}=P/\mathrm{}P`$ is simple, by Exercise 15.3 of . It thus suffices to show $`\mathrm{End}_{S_n}(\overline{P})=𝐅_{\mathrm{}}`$. Let $`\overline{P}_1`$ denote the set of elements of $`\overline{P}`$ invariant under $`\{\sigma S_n\sigma (1)=1\}`$. Then $`\overline{\omega }_1`$ generates the one-dimensional $`𝐅_{\mathrm{}}`$-vector space $`\overline{P}_1`$. Suppose $`\gamma \mathrm{End}_{S_n}(\overline{P})`$. Then $`\overline{P}_1`$ is $`\gamma `$-stable, so there is a $`c𝐅_{\mathrm{}}`$ such that $`\gamma (\overline{\omega }_1)=c\overline{\omega }_1`$. Since $`\gamma `$ commutes with $`S_n`$ and the orbit $`S_n(\overline{\omega }_1)`$ contains all the $`\overline{\omega }_i`$’s, we have $`\gamma (\overline{\omega }_i)=c\overline{\omega }_i`$. Since $`\overline{P}`$ is generated by the $`\overline{\omega }_i`$’s, we have $`\gamma =c𝐅_{\mathrm{}}`$. ∎
## 4. $`\mathrm{}`$-full abelian varieties
Suppose $`A`$ is an abelian variety over a field $`F`$ of characteristic $`\mathrm{}`$.
###### Definition 4.1.
We will say that $`A`$ is $`\mathrm{}`$-full if the $`𝐙_{\mathrm{}}[G_{F^{(\mathrm{})}}]`$-module $`T_{\mathrm{}}(A)`$ is well-rounded (see Definition 3.2).
###### Remarks 4.2.
1. Since $`T_{\mathrm{}}(A)/\mathrm{}T_{\mathrm{}}(A)=A_{\mathrm{}}`$, the abelian variety $`A`$ is $`\mathrm{}`$-full if and only if the $`G_{F^{(\mathrm{})}}`$-module $`A_{\mathrm{}}`$ is absolutely simple.
2. Clearly, $`\mathrm{}`$-fullness is stable under $`F^{(\mathrm{})}`$-isogeny.
3. If $`A`$ is $`\mathrm{}`$-full then $`\mathrm{End}_{F^{(\mathrm{})}}(A)=𝐙`$.
4. If $`A`$ is $`\mathrm{}`$-full then it has a polarization defined over $`F`$ of degree prime to $`\mathrm{}`$ (this follows from Thm. 3 in §23 of and Lemma 3.1c).
5. Suppose $`A`$ is $`\mathrm{}`$-full and $`\varphi :AB`$ is an $`F`$-isogeny whose degree is a power of $`\mathrm{}`$. Then $`A`$ and $`B`$ are $`F`$-isomorphic (since $`\mathrm{ker}(\varphi )=A_\mathrm{}^m`$ for some $`m`$).
Suppose now that $`A`$ is an elliptic curve $`E`$.
1. The image of $`G_{F^{(\mathrm{})}}`$ in $`\mathrm{Aut}(T_{\mathrm{}}(E))\mathrm{GL}_2(𝐙_{\mathrm{}})`$ is the intersection of the image of $`G_F`$ and $`\mathrm{SL}(T_{\mathrm{}}(E))\mathrm{SL}_2(𝐙_{\mathrm{}})`$. In particular, if $`G_F\mathrm{Aut}(T_{\mathrm{}}(E))`$ is surjective then the image of $`G_{F^{(\mathrm{})}}`$ in $`\mathrm{Aut}(T_{\mathrm{}}(E))`$ is $`\mathrm{SL}(T_{\mathrm{}}(E))\mathrm{SL}_2(𝐙_{\mathrm{}})`$ and $`E`$ is $`\mathrm{}`$-full.
2. Suppose $`\rho _{E,\mathrm{}}(G_{F^{(\mathrm{})}})`$ contains $`\mathrm{SL}_2(𝐅_{\mathrm{}})`$ (and therefore is $`\mathrm{SL}_2(𝐅_{\mathrm{}})`$). Then $`E`$ is $`\mathrm{}`$-full (see (ii) of Lemma 3.1).
3. By Lemma 3 on p. IV-23 of , if $`\mathrm{}5`$ and $`F`$ is a number field, then the $`\mathrm{}`$-adic representation $`G_{F^{(\mathrm{})}}\mathrm{SL}_2(𝐙_{\mathrm{}})`$ is surjective if and only if the mod $`\mathrm{}`$ representation $`\rho _{E,\mathrm{}}:G_{F^{(\mathrm{})}}\mathrm{SL}_2(𝐅_{\mathrm{}})`$ is surjective.
###### Example 4.3.
Let $`E`$ be the elliptic curve $`y^2+y=x^3x`$ over $`𝐐`$. The map $`G_𝐐\mathrm{Aut}(T_{\mathrm{}}(E))`$ is surjective for all primes $`\mathrm{}`$ (see , 5.5.6). In particular, $`E`$ is $`\mathrm{}`$-full for all primes $`\mathrm{}`$.
###### Proposition 4.4.
Suppose $`\mathrm{}`$ is an odd prime, and $`E`$ is an elliptic curve over a field $`F`$ of characteristic $`\mathrm{}`$. Suppose $`\rho _{E,\mathrm{}}(G_F)`$ contains $`\mathrm{SL}_2(𝐅_{\mathrm{}})`$. Then $`E`$ is $`\mathrm{}`$-full.
###### Proof.
The commutator subgroup $`[G_F,G_F]`$ lies in $`G_{F^{(\mathrm{})}}`$. First suppose $`\mathrm{}>3`$. Then $`\rho _{E,\mathrm{}}(G_{F^{(\mathrm{})}})`$ contains $`[\mathrm{SL}_2(𝐅_{\mathrm{}}),\mathrm{SL}_2(𝐅_{\mathrm{}})]=\mathrm{SL}_2(𝐅_{\mathrm{}})`$, and we are done by Remark 4.2vii. Suppose $`\mathrm{}=3`$. We have $`[\mathrm{GL}_2(𝐅_3):\mathrm{SL}_2(𝐅_3)]=2`$. Let $`H=\rho _{E,3}(G_{F^{(3)}})`$. By Remark 4.2vii, we may suppose that $`H`$ is a proper subgroup of $`\mathrm{SL}_2(𝐅_3)`$. If $`\rho _{E,3}(G_F)=\mathrm{GL}_2(𝐅_3)`$, then $`H[GL_2(𝐅_3),\mathrm{GL}_2(𝐅_3)]=\mathrm{SL}_2(𝐅_3)`$. Thus, we may assume $`\rho _{E,3}(G_F)=\mathrm{SL}_2(𝐅_3)`$. Then $`F`$ contains a primitive cube root of unity and therefore $`G_F/G_{F^{(3)}}`$ is a pro-cyclic $`3`$-group. Then $`\mathrm{SL}_2(𝐅_3)/H`$ is a cyclic $`3`$-group. Since $`\mathrm{\#}\mathrm{SL}_2(𝐅_3)=24`$, $`H`$ has order $`8`$ which is prime to $`3`$ and therefore the $`H`$-module $`E_3𝐅_3^2`$ is semisimple. Since $`H`$ is non-commutative and $`\mathrm{dim}_{𝐅_3}E_3=2`$, it follows that the $`H`$-module $`E_3`$ is absolutely simple. ∎
###### Proposition 4.5.
Suppose $`E`$ is an elliptic curve over a number field $`F`$. Suppose $`j>2`$ is an integer such that $`\zeta _{2^j}F(\zeta _{2^{j1}})`$, and suppose $`\rho _{E,2^j}(G_F)`$ contains $`\mathrm{SL}_2(𝐙/2^j𝐙)`$. Then $`\rho _{E,2}(G_{F^{(2)}})=\mathrm{SL}_2(𝐅_2)`$, and therefore $`E`$ is $`2`$-full.
###### Proof.
Write $`F^{}=F(\zeta _{2^{j1}})`$ and $`F^{\prime \prime }=F(\zeta _{2^j})`$. The mod $`2`$ representation
$$\rho _{E,2}:G_F\mathrm{Aut}(E_2)\mathrm{SL}_2(𝐅_2)S_3$$
is surjective, since $`\mathrm{SL}_2(𝐙/2^j𝐙)\mathrm{SL}_2(𝐅_2)`$ is surjective. Let $`G=\rho _{E,2}(G_{F^{(2)}})`$. Since $`G_F/G_{F^{(2)}}`$ is a pro-$`2`$-group, the index of $`G`$ in $`\mathrm{SL}_2(𝐅_2)S_3`$ is a power of $`2`$. Thus either $`G\mathrm{SL}_2(𝐅_2)`$ and we are done, or $`GA_3`$. Assume the latter. We have
(1)
$$\rho _{E,2^j}(G_F^{})\rho _{E,2^j}(G_{F^{\prime \prime }})=\rho _{E,2^j}(G_F)\mathrm{SL}_2(𝐙/2^j𝐙)=\mathrm{SL}_2(𝐙/2^j𝐙).$$
Let $`\sigma `$ be the composition
$$G_F^{}\stackrel{\rho _{E,2}}{}\mathrm{GL}_2(𝐅_2)=S_3S_3/A_3=\{1,1\},$$
and let $`F_\sigma `$ be the fixed field of $`\mathrm{ker}(\sigma )`$. Then $`F_\sigma `$ is a quadratic extension of $`F^{}`$. By (1), $`F_\sigma F^{\prime \prime }`$. Since $`j>2`$, $`F^{\prime \prime }`$ is the only quadratic extension of $`F^{}`$ in $`F^{(2)}`$. Therefore $`F_\sigma `$ is not contained in $`F^{(2)}`$, so $`\sigma (G_{F^{(2)}})1`$. This contradicts the assumption that $`G=A_3`$. ∎
###### Theorem 4.6.
Suppose $`F`$ is a number field and $`n`$ is a positive integer. Then there are infinitely many elliptic curves $`E`$ over $`F`$, non-isomorphic over $`𝐂`$, for which $`\mathrm{SL}_2(𝐙/n𝐙)\rho _{E,n}(G_F)`$.
###### Proof.
As shown on pp. 145–6 of , there is an elliptic curve $`E(t)`$ over the function field $`𝐐(t)`$, with $`j`$-invariant $`t`$, such that
$$\mathrm{Gal}(𝐐(t,E(t)_n)/𝐐(\zeta _n)(t))=\mathrm{SL}_2(𝐙/n𝐙)$$
and $`𝐐(\zeta _n)`$ is algebraically closed in $`𝐐(t,E(t)_n)`$. Thus
$$F(\zeta _n)(t)𝐐(t,E(t)_n)=𝐐(\zeta _n)(t),$$
and we have
$$\mathrm{Gal}(F(t,E(t)_n)/F(\zeta _n)(t))=\mathrm{SL}_2(𝐙/n𝐙).$$
By Prop. 2 on p. 123 of and Hilbert’s Irreducibility Theorem (p. 130 of ), there are infinitely many specializations $`t_0F`$ such that
$$\mathrm{Gal}(F(E(t_0)_n)/F(\zeta _n))=\mathrm{Gal}(F(t,E(t)_n)/F(\zeta _n)(t))=\mathrm{SL}_2(𝐙/n𝐙),$$
and therefore $`\mathrm{Gal}(F(E(t_0)_n)/F)\mathrm{SL}_2(𝐙/n𝐙)`$. ∎
###### Corollary 4.7.
Let $`F`$ be a number field and $`r`$ a positive integer. Then there are infinitely many elliptic curves over $`F`$, non-isomorphic over $`𝐂`$, that are $`\mathrm{}`$-full for all prime divisors $`\mathrm{}`$ of $`r`$.
###### Proof.
Choose $`j>2`$ such that $`\zeta _{2^j}F(\zeta _{2^{j1}})`$ (such a $`j`$ exists since $`[F(\zeta _{2^k}):F]`$ is unbounded as $`k\mathrm{}`$). Let $`n^{}`$ be the product of the prime divisors of $`r`$. Let $`n=n^{}`$ if $`r`$ is odd and let $`n=2^jn^{}`$ if $`r`$ is even. By Theorem 4.6, there are infinitely many elliptic curves $`E`$ over $`F`$, non-isomorphic over $`𝐂`$, for which $`\mathrm{SL}_2(𝐙/n𝐙)\rho _{E,n}(G_F)`$. Thus for all odd prime divisors $`\mathrm{}`$ of $`r`$ we have $`\mathrm{SL}_2(𝐅_{\mathrm{}})\rho _{E,\mathrm{}}(G_F)`$, so $`E`$ is $`\mathrm{}`$-full by Proposition 4.4. If $`r`$ is even, then $`E`$ is $`2`$-full by Proposition 4.5. ∎
###### Proposition 4.8 (see ).
Suppose $`Y`$ is an abelian variety over a field $`F`$ of characteristic $`\mathrm{}`$. Then there exists a finite Galois extension $`K_{\mathrm{}}/F`$ such that every finite extension of $`F`$ linearly disjoint from $`K_{\mathrm{}}`$ is linearly disjoint from $`F_Y^{(\mathrm{})}`$.
###### Proof.
Let $`H=\mathrm{Gal}(F_Y^{(\mathrm{})}/F)`$, an $`\mathrm{}`$-adic Lie group. By the Proposition and Example 1 in § 10.6 of , there exists an open normal subgroup $`\mathrm{\Phi }H`$ of finite index such that whenever $`\stackrel{~}{H}`$ is a closed normal subgroup in $`H`$ with $`H=\stackrel{~}{H}\mathrm{\Phi }`$, then $`\stackrel{~}{H}=H`$. Let $`K_{\mathrm{}}`$ be the fixed field of $`\mathrm{\Phi }`$. Suppose $`K`$ is a finite extension of $`F`$, let $`\stackrel{~}{H}=\mathrm{Gal}(K_Y^{(\mathrm{})}/K)`$, and view $`\stackrel{~}{H}`$ as a closed normal subgroup of $`H`$. If $`K`$ and $`K_{\mathrm{}}`$ are linearly disjoint, then the restriction to $`\stackrel{~}{H}`$ of $`HH/\mathrm{\Phi }`$ is surjective, i.e., $`H=\stackrel{~}{H}\mathrm{\Phi }`$. Therefore $`\stackrel{~}{H}=H`$, so $`K`$ and $`F_Y^{(\mathrm{})}`$ are linearly disjoint. ∎
## 5. A general method
Suppose $`A`$ is an abelian variety over a field $`F`$ of characteristic $`\mathrm{}`$. There is a canonical $`𝐙_{\mathrm{}}`$-bilinear $`G_F`$-equivariant perfect pairing
$$e_{\mathrm{}}:T_{\mathrm{}}(A)\times T_{\mathrm{}}(A^{})𝐙_{\mathrm{}}(1)𝐙_{\mathrm{}}$$
where $`𝐙_{\mathrm{}}(1)`$ is the projective limit of the groups of $`\mathrm{}^m`$-th roots of unity (, §20; , Ch. VII, §2). Since $`G_{F^{(\mathrm{})}}`$ acts trivially on $`𝐙_{\mathrm{}}(1)`$, the pairing $`e_{\mathrm{}}`$ is $`G_{F^{(\mathrm{})}}`$-invariant. There exists an $`F`$-isogeny $`AA^{}`$, and therefore a surjection $`A_{\mathrm{}^{\mathrm{}}}A_{\mathrm{}^{\mathrm{}}}^{}`$ defined over $`F`$. Thus, $`F_A^{}^{(\mathrm{})}F_A^{(\mathrm{})}`$. Since $`e_{\mathrm{}}`$ is perfect, it follows that
$$F^{(\mathrm{})}F_A^{(\mathrm{})}.$$
Suppose $``$ is an invertible sheaf on $`\overline{A}`$. Let $`\chi ()`$ denote the Euler characteristic, and let $`\varphi _{}:\overline{A}\overline{A^{}}`$ be the associated natural homomorphism defined in §13 (see p. 131) of . Then $`\mathrm{deg}(\varphi _{})=\chi ()^2`$ (see §16 of ), and $`\varphi _{}`$ is an isogeny if and only if $`\chi ()0`$. If $``$ is ample, then $`\chi ()>0`$ and the isogeny $`\varphi _{}`$ is called a polarization on $`\overline{A}`$.
Suppose now that $`\varphi _{}`$ is defined over $`F`$ (i.e., is obtained by extension of scalars from a morphism $`AA^{}`$). Then $`\varphi _{}`$ induces a homomorphism of $`G_F`$-modules $`f_{}:T_{\mathrm{}}(A)T_{\mathrm{}}(A^{})`$ which is injective if and only if $`\varphi _{}`$ is an isogeny. If $`f_{}`$ is injective, then the index of the image of $`T_{\mathrm{}}(A)`$ in $`T_{\mathrm{}}(A^{})`$ is the exact power of $`\mathrm{}`$ dividing $`\mathrm{deg}(\varphi _{})`$ (, Ch. VII, §1, Thm. 1 and its proof). Extending $`𝐐_{\mathrm{}}`$-linearly, the induced homomorphism $`f_{}:V_{\mathrm{}}(A)V_{\mathrm{}}(A^{})`$ is an isomorphism of $`𝐐_{\mathrm{}}[G_F]`$-modules. The homomorphism $`f_{}`$ gives rise to an alternating $`G_F`$-equivariant pairing
$$E^{}:T_{\mathrm{}}(A)\times T_{\mathrm{}}(A)𝐙_{\mathrm{}}(1)𝐙_{\mathrm{}},x,ye_{\mathrm{}}(x,f_{}(y))$$
(see §20 of ). The form $`E^{}`$ is non-degenerate if and only if $`\varphi _{}`$ is an isogeny. It is perfect if and only if $`\mathrm{deg}(\varphi _{})`$ is not divisible by $`\mathrm{}`$.
Assume that $`\varphi _{}`$ is an isogeny. Let $`T_{\mathrm{}}(A)^{}`$ denote the dual lattice of $`T_{\mathrm{}}(A)`$ with respect to $`E^{}`$. Then
$$T_{\mathrm{}}(A)T_{\mathrm{}}(A)^{}=f_{}^1(T_{\mathrm{}}(A^{}))V_{\mathrm{}}(A),$$
and $`\mathrm{\#}(T_{\mathrm{}}(A)^{}/T_{\mathrm{}}(A))=\mathrm{\#}(T_{\mathrm{}}(A^{})/f_{}(T_{\mathrm{}}(A))`$ is the exact power of $`\mathrm{}`$ dividing $`\mathrm{deg}(\varphi _{})`$ (see the proof of Thm. 3 in §20 of ).
###### Lemma 5.1.
Suppose $`Y`$ is an abelian variety over a field $`F`$ of characteristic $`\mathrm{}`$, and $`s`$ is a positive integer. Suppose $`K/F`$ is a finite Galois extension with $`\mathrm{Gal}(K/F)=D\mathrm{GL}_s(𝐙)`$, such that $`1D`$. Let $`A`$ be the $`K/F`$-form of $`Y^s`$ attached to the inclusion $`\mathrm{Gal}(K/F)=D\mathrm{GL}_s(𝐙)\mathrm{Aut}(Y^s)`$. Let $`M`$ denote the image of $`G_F`$ in $`\mathrm{Aut}(T_{\mathrm{}}(Y))`$. Then:
1. the natural map $`f:M\times D\mathrm{Aut}(T_{\mathrm{}}(Y)_𝐙_{\mathrm{}}𝐙_{\mathrm{}}^s)=\mathrm{Aut}(T_{\mathrm{}}(A))`$ is injective,
2. $`F_A^{(\mathrm{})}=KF_Y^{(\mathrm{})}`$.
Suppose now that $`K`$ and $`F_Y^{(\mathrm{})}`$ are linearly disjoint over $`F`$. Then:
1. the image of $`G_F`$ in $`\mathrm{Aut}(T_{\mathrm{}}(A))`$ is $`f(M\times D)`$,
2. if $`Y`$ is $`\mathrm{}`$-full and the $`𝐙_{\mathrm{}}[D]`$-module $`𝐙_{\mathrm{}}^s`$ is well-rounded, then $`A`$ is $`\mathrm{}`$-full.
###### Proof.
The kernel of
$$\mathrm{Aut}(T_{\mathrm{}}(Y))\times \mathrm{GL}_s(𝐙_{\mathrm{}})\mathrm{Aut}(T_{\mathrm{}}(Y)_𝐙_{\mathrm{}}𝐙_{\mathrm{}}^s)=\mathrm{GL}_s(\mathrm{End}_𝐙_{\mathrm{}}(T_{\mathrm{}}(Y))$$
is $`\{(a,a^1)a𝐙_{\mathrm{}}^{}\}`$. Since $`1D`$, we have
$$1=D\mathrm{Aut}(T_{\mathrm{}}(Y))\mathrm{GL}_s(\mathrm{End}_𝐙_{\mathrm{}}(T_{\mathrm{}}(Y)),$$
and (a) follows. From the natural injections
$$\mathrm{Gal}(KF_Y^{(\mathrm{})}/F)M\times D\stackrel{f}{}\mathrm{Aut}(T_{\mathrm{}}(A))$$
we have (b). Suppose $`K`$ and $`F_Y^{(\mathrm{})}`$ are linearly disjoint over $`F`$. Then the natural injection $`\mathrm{Gal}(KF_Y^{(\mathrm{})}/F)\mathrm{Gal}(F_Y^{(\mathrm{})}/F)\times \mathrm{Gal}(K/F)M\times D`$ is surjective, and (c) follows. Now (d) follows from Remark 3.3, by applying (c) and (a) with $`F^{(\mathrm{})}`$ in place of $`F`$. ∎
###### Theorem 5.2.
Suppose $`F`$ is a field of characteristic different from $`\mathrm{}`$, $`Y`$ is an $`\mathrm{}`$-full $`d`$-dimensional abelian variety over $`F`$, and $`s`$ is a positive integer. Suppose $`D`$ is an irreducible subgroup of $`\mathrm{GL}_s(𝐙)`$ such that $`1D`$ and such that there does not exist a $`D`$-stable $`𝐙_{\mathrm{}}`$-lattice in $`𝐐_{\mathrm{}}^s`$ with a perfect $`D`$-invariant $`𝐐_{\mathrm{}}`$-valued symmetric pairing. Suppose $`K/F`$ is a Galois extension with $`\mathrm{Gal}(K/F)=D`$ and with $`K`$ and $`F_Y^{(\mathrm{})}`$ linearly disjoint over $`F`$. Let $`A`$ be the twist of $`Y^s`$ via
$$\mathrm{Gal}(K/F)=D\mathrm{GL}_s(𝐙)\mathrm{Aut}(Y^s),$$
and suppose $`B`$ is an abelian variety $`F^{(\mathrm{})}`$-isogenous to $`A`$. Then $`\mathrm{}^{2d}`$ divides the degree of every $`F^{(\mathrm{})}`$-polarization on $`B`$, and $`\mathrm{}^d`$ divides the Euler characteristic of every invertible sheaf $``$ on $`\overline{B}`$ such that $`\varphi _{}`$ is defined over $`F^{(\mathrm{})}`$. If $`\chi ()0`$ and $`\mathrm{}^r`$ is the highest power of $`\mathrm{}`$ dividing $`\chi ()`$, then $`\mathrm{}^r`$ is a power of $`\mathrm{}^d`$.
###### Proof.
If $`\varphi _{}`$ is not an isogeny then $`\chi ()=0`$ and we are done. Suppose now that $`\varphi _{}`$ is an isogeny. Let $`N`$ denote the image of $`G_{F^{(\mathrm{})}}`$ in $`\mathrm{Aut}(T_{\mathrm{}}(Y))`$. By Lemma 5.1, the image of $`G_{F^{(\mathrm{})}}`$ in $`\mathrm{Aut}(V_{\mathrm{}}(A))`$ is $`N\times D`$. Therefore the $`G_{F^{(\mathrm{})}}`$-stable $`𝐙_{\mathrm{}}`$-lattice $`T_{\mathrm{}}(B)`$ is also ($`N\times D`$)-stable. Since $`E^{}`$ is $`G_{F^{(\mathrm{})}}`$-invariant, it is also ($`N\times D`$)-invariant. Applying Theorem 3.5 with $`S=T_{\mathrm{}}(Y)`$, $`T=T_{\mathrm{}}(B)`$, and $`e=E^{}`$, then $`\mathrm{\#}(T^{}/T)=\mathrm{}^{2dt}`$ for some positive integer $`t`$, and the result follows. ∎
## 6. Application
###### Theorem 6.1.
Suppose $`F`$ is a number field and $`r`$ is a positive integer. Then there exists an abelian variety $`A`$ over $`F`$ with the following properties:
1. if $`B`$ is an abelian variety $`F`$-isogenous to $`A`$ and $``$ is an invertible sheaf on $`\overline{B}`$ such that $`\varphi _{}`$ is defined over $`F`$, then $`\chi ()`$ is divisible by $`r`$;
2. if $`\mathrm{}`$ is a prime, $`\mathrm{}^a`$ divides $`r`$, $`B`$ is an abelian variety $`F^{(\mathrm{})}`$-isogenous to $`A`$, and $``$ is an invertible sheaf on $`\overline{B}`$ such that $`\varphi _{}`$ is defined over $`F^{(\mathrm{})}`$, then $`\chi ()`$ is divisible by $`\mathrm{}^a`$.
###### Proof.
We may assume that $`r>1`$. Let $`m`$ be a positive integer relatively prime to $`r`$ and such that for every prime divisor $`\mathrm{}`$ of $`r`$, $`\mathrm{}^{m1}`$ does not divide $`r`$. Let $`n=6`$ if $`r`$ is a power of $`2`$, and otherwise let $`n`$ be the product of the prime divisors of $`r`$. Identify the symmetric group $`S_t`$ in the standard way as an irreducible subgroup of $`\mathrm{GL}_{t1}(𝐙)`$. By Lemma 3.6, for all prime divisors $`\mathrm{}`$ of $`r`$ the $`𝐙_{\mathrm{}}[S_m]`$-module $`𝐙_{\mathrm{}}^{m1}`$ is well-rounded. Note that $`m>2`$ and $`n>2`$, so for $`t=m`$ and $`n`$ we have $`1S_tGL_{t1}(𝐙)`$. By Lemma 3.6v, no $`S_n`$-stable $`𝐙_{\mathrm{}}`$-lattice in $`U𝐐_{\mathrm{}}^{n1}`$ has a perfect $`S_n`$-invariant $`𝐙_{\mathrm{}}`$-valued pairing. By Corollary 4.7, there exists an elliptic curve $`E`$ over $`F`$ that is $`\mathrm{}`$-full for all prime divisors $`\mathrm{}`$ of $`r`$. By Prop. 4.8, for every prime $`\mathrm{}`$ there exists a finite Galois extension $`K_{\mathrm{}}/F`$ such that every finite extension of $`F`$ linearly disjoint from $`K_{\mathrm{}}`$ is linearly disjoint from $`F_E^{(\mathrm{})}`$. There exists a finite Galois extension $`L`$ of $`F`$ with $`\mathrm{Gal}(L/F)=S_m\times S_n`$ which is linearly disjoint from $`K_{\mathrm{}}`$ over $`F`$ for all prime divisors $`\mathrm{}`$ of $`r`$ (see , especially Prop. 3.3.3). We can write $`L=K^{}K`$ with $`\mathrm{Gal}(K^{}/F)=S_m`$ and $`\mathrm{Gal}(K/F)=S_n`$. Then $`K^{}`$ and $`K`$ are linearly disjoint over $`F`$, and $`K^{}`$ (resp., $`K`$) and $`F_E^{(\mathrm{})}`$ are linearly disjoint over $`F`$. Let $`Y`$ be the $`K^{}/F`$-form of $`E^{m1}`$ attached to
$$\mathrm{Gal}(K^{}/F)=S_m\mathrm{GL}_{m1}(𝐙)\mathrm{Aut}(E^{m1}),$$
and let $`A`$ be the $`K/F`$-form of $`Y^{n1}`$ attached to
$$\mathrm{Gal}(K/F)=S_n\mathrm{GL}_{n1}(𝐙)\mathrm{Aut}(Y^{n1}).$$
By Lemma 5.1b, $`F_Y^{(\mathrm{})}=K^{}F_E^{(\mathrm{})}`$, so $`K`$ and $`F_Y^{(\mathrm{})}`$ are linearly disjoint over $`F`$. By Lemma 5.1d, $`Y`$ is $`\mathrm{}`$-full. Part (b) now follows from Theorem 5.2 with $`D=S_n`$. Part (a) follows from (b) for all prime divisors $`\mathrm{}`$ of $`r`$. ∎
Applying Theorem 6.1 when $``$ is ample yields Theorem 1.1.
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# Space VLBI Observations of 3C 279 at 1.6 and 5 GHz
## 1 Introduction
The quasar 3C 279 ($`z`$=0.536) is one of the most intensively studied quasars for several reasons. It was the first radio source observed to exhibit superluminal motion (Knight et al. 1971; Whitney et al. 1971; Cohen et al. 1971). It was the first blazar — and remains one of the brightest — detected in high-energy $`\gamma `$-rays by the EGRET instrument on the Compton Gamma Ray Observatory (Hartman et al. 1992). The EGRET detection prompted several large multiwavelength studies of this source. Results of these studies are presented by Maraschi et al. (1994), Grandi et al. (1996), Hartman et al. (1996), and Wehrle et al. (1998). The radio flux density of 3C 279 has been monitored by the Michigan group since 1965 at frequencies of 4.8, 8.0, and 14.5 GHz (Aller et al. 1985). The observations presented in this paper occurred at the beginning of a total flux density flare that would later reach the highest flux densities yet recorded in this program for 3C 279. <sup>4</sup><sup>4</sup>4http://www.astro.lsa.umich.edu/obs/radiotel/gif/1253-055.gif
Following the discovery of superluminal motion, the parsec scale structure of 3C 279 has been monitored using the VLBI technique. Cotton et al. (1979) measured a speed of 15 $`c`$ for the original superluminal jet component. (Throughout the paper we assume $`H_0`$=70 km s<sup>-1</sup> Mpc<sup>-1</sup> and $`q_0`$=0.1, and component speeds measured by others have been expressed in these terms.) Unwin et al. (1989) and Carrara et al. (1993) describe VLBI monitoring of 3C 279 at 5, 11, and 22 GHz throughout the 1980s. These authors observed the motions of several new superluminal components, and found that the speeds of these components were only one-quarter to one-third (3-5 $`c`$) of that measured for the original superluminal component during the 1970s. VLBI monitoring of 3C 279 has been undertaken at 22 and 43 GHz during the 1990s. Initial results of this high-frequency monitoring are reported by Wehrle, Unwin, & Zook (1994) and Wehrle et al. (1996), and final results will be reported by Wehrle et al. (2000). High-resolution VLBI polarimetric images of 3C 279 have been made by Leppänen, Zensus, & Diamond (1995), Lister, Marscher, & Gear (1998) and Wardle et al. (1998). The detection of circularly polarized radio emission by Wardle et al. provides some of the first direct evidence that electron-positron pairs are an important component of the jet plasma. One notable feature of the VLBI observations of this source has been the differing speeds and position angles of the VLBI components. Carrara et al. (1993) and Abraham & Carrara (1998) claim that these can be explained by ejection of components by a precessing jet.
3C 279 was observed during the TDRSS space VLBI experiments at 2.3 GHz (Linfield et al. 1989) and 15 GHz (Linfield et al. 1990), with source frame brightness temperatures between 1.6 and 2.0 $`\times `$ 10<sup>12</sup> K being measured for this source. A brightness temperature of 1.9 $`\times `$ 10<sup>12</sup> K (translated from observed frame to source frame) was also measured for 3C 279 in the 22 GHz VSOP Pre-Launch Survey (Moellenbrock et al. 1996). This source has been detected in VLBI observations up to frequencies of 215 GHz (Krichbaum et al. 1997). This paper reports on the VSOP observations of 3C 279 made during the first Announcement of Opportunity period (AO1). Hirabayashi et al. (1999) and Edwards et al. (1999) presented preliminary analyses of these 5 and 1.6 GHz VSOP observations respectively. Here we present higher dynamic range images together with analysis and interpretation of model fits and a spectral index map.
## 2 Observations
The quasar 3C 279 was observed on two consecutive days during the AO1 phase of the VSOP mission: on 1998 January 9 at 1.6 GHz, and on 1998 January 10 at 5 GHz. The VSOP mission uses the Japanese HALCA satellite as an element in a changeable VLBI array in order to obtain visibility measurements on baselines larger than the Earth’s diameter. HALCA was launched on 1997 February 12 and carries an 8 meter antenna through an elliptical orbit with an apogee height of 21,400 km (yielding baselines up to 2.6 Earth diameters), a perigee height of 560 km, and an orbital period of 6.3 hours. HALCA has operational observing bands at 1.6 and 5 GHz (18 and 6 cm). The data from the satellite are recorded by a network of ground tracking stations and subsequently correlated with the data from the participating ground telescopes. The VSOP system and initial science results are discussed by Hirabayashi et al. (1998).
The observations of 1998 January 9–10 were conducted using the standard VSOP observing mode: two 16 MHz intermediate frequency bands, each 2-bit sampled at the Nyquist rate in left circular polarization, for a total data rate of 128 Mbps. The ground telescope arrays were made up of nine elements of the NRAO Very Long Baseline Array<sup>5</sup><sup>5</sup>5The National Radio Astronomy Observatory is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc. (VLBA) — Hancock did not observe because of a power failure — with the addition of the 70 m telescopes at Goldstone, California, U.S.A. and Tidbinbilla, Australia on January 9, and the 64 m telescope at Usuda, Japan on January 10. The total observing time on January 9 was 6 hours, including 3.5 hours of HALCA data from the tracking station at Green Bank, West Virginia, covering the portion of HALCA’s orbit from near perigee to near apogee. The data from Goldstone were not used in the final image as it observed for only a short time and there were significant calibration uncertainties with this data. The total observing time on January 10 was 10 hours, although for the final 1.5 hours only Mauna Kea, Usuda and HALCA observed and this data was also excluded from the final image (as four telescopes are required for amplitude self-calibration). This observation included two HALCA tracking passes of 3 hours each (separated by 4 hours) by the tracking stations at Robledo, Spain and Goldstone, California, U.S.A. Each tracking pass covered the portion of HALCA’s orbit from near perigee to near apogee. The data from both observations were correlated at the VLBA correlator in Socorro. Through the remainder of this paper, we use “VSOP” to refer to the HALCA+ground-array combination.
Calibration and fringe-fitting were done with the AIPS software package. 3C 279 is a strong source, and good fringes were found to HALCA during all tracking passes. The antenna gain for Saint Croix had to be adjusted upward by a factor of $``$ 2 from its nominal value at 5 GHz since the VSOP observations were conducted at 4800–4832 MHz because HALCA’s performance is better at these frequencies. Figure 1 shows the $`(u,v)`$ plane coverages of the two observations. These $`(u,v)`$ plane coverages result in highly elliptical beams: the major to minor axis ratio of the 5 GHz beam is 8:1. The addition of Tidbinbilla to the 1.6 GHz observation improves the north-south coverage and reduces this ratio to 4:1. Plots of correlated flux density vs. $`(u,v)`$ distance projected along a position angle of $`115^{}`$ (the position angle of the brightest structure) are shown in Figure 2. The addition of HALCA extends the projected $`(u,v)`$ distances compared to the ground-only baselines from 120 to 340 M$`\lambda `$ at 5 GHz, and from 65 to 120 M$`\lambda `$ at 1.6 GHz. The beating evident in the correlated flux densities indicates that on milliarcsecond scales the morphology is dominated by two components of similar flux density.
## 3 Results
### 3.1 Images
Images from these datasets were produced using standard CLEAN and self-calibration procedures from the Difmap software package (Shepherd, Pearson, & Taylor 1994). Figure 3 shows two images from the 5 GHz observation on 1998 January 10: the full-resolution space VLBI image and the VLBA-only image made from the same dataset with the space baselines removed. Figure 4 shows the full-resolution space VLBI and the VLBA-only images from the 1.6 GHz observation on the previous day (Tidbinbilla baselines were also removed from this ground-only image to reduce the size of the $`(u,v)`$ holes). The VSOP images are displayed using uniform weighting (although cleaning was done with both uniform and natural weighting to model the extended structure), while the VLBA images are displayed with natural weighting. We stress the importance of using uniform weighting with VSOP datasets. A naturally weighted VSOP image degrades the resolution of the ground-space array because of the higher weighting of the larger ground antennas and the much denser sampling of the $`(u,v)`$ plane on the ground baselines.
The VSOP and VLBA-only images differ in scale by about a factor of four, showing the greatly increased resolution provided by the space baselines. The VSOP images presented in Figures 3 and 4 are the highest resolution images yet produced of 3C 279 at these frequencies. The 1.6 and 5 GHz VSOP images have formal dynamic ranges (peak:rms) of 1,500:1 and 2,000:1 respectively, demonstrating that high dynamic range images can be made despite the small size of the HALCA orbiting antenna. For comparison, the VLBA-only images have dynamic ranges of 5,000:1 and 15,000:1 respectively. The 5 GHz VSOP image shows that on small angular scales 3C 279 is dominated by a double structure. This structure consists of the compact core and inner jet region (the feature to the east) and a bright jet component about 3 mas from the core along a position angle of $`115^{}`$ (the feature to the west). This bright jet component, ‘C4’, is a well known feature first observed in 1985. This component will be discussed further in $`\mathrm{\S }`$ 4.1.
The structure on slightly larger angular scales is quite different. The 5 GHz VLBA image and the 1.6 GHz VSOP image both show the double structure mentioned above (although in the 1.6 GHz image the jet component is brighter than the core which results in the jet component being placed at the phase center), as well as a more extended jet to the southwest along a position angle of approximately $`140^{}`$. This position angle is similar to that seen in older VLBI images (see $`\mathrm{\S }`$ 4.1) as well as VLA and MERLIN images (de Pater & Perley 1983; Pilbratt, Booth, & Porcas 1987; Akujor et al. 1994). The resolutions of the 5 GHz VLBA image and the 1.6 GHz VSOP image are roughly equal, a fact that allows a spectral index map to be made from these images ($`\mathrm{\S }`$ 4.2). The jet emission in the 1.6 GHz VSOP image is quite complex, and it appears in Figure 4 that the jet may be limb brightened. However, we caution against over-interpreting these complex jet features because the CLEAN striping produced by the holes in the $`(u,v)`$ plane coverage runs parallel to the jet, and because simulations indicate that, due to the lack of complete $`(u,v)`$ plane coverage, space VLBI images may show such knotty structure when the actual brightness distribution is smoother (D. Murphy 1999, private communication). The 1.6 GHz VLBA image shows structure extending out to $``$ 100 mas from the core, all the way out to the smallest size scales sampled by the 22 GHz VLA images of de Pater & Perley (1983).
### 3.2 Model Fits
The Difmap model-fitting routine was used to fit elliptical Gaussian components to the visibility data for each image in Figures 3 and 4. When model fitting VSOP data, it is important to increase the weight of HALCA over the default weighting used by Difmap in order for the space baseline data to have any effect on the model. We increased the weight of HALCA during model fitting by a factor equal to the product of the ratios of the average ground baseline sensitivity to the average space baseline sensitivity ($``$ 50 for these observations) and the number of ground visibilities to the number of space visibilities ($``$ 10 for these observations). This effectively achieves a “uniform” weighting during model fitting, and causes the space visibilities to have an effect on the model fitting equal to that of the ground visibilities.
The results of the model fitting are given in Table 1. Component numbers are given only for ease of later reference, and are not meant to identify the same component between images. Tentative component identifications obtained from the discussion in $`\mathrm{\S }`$ 4.1 are given in the third column. Note that the lower resolution images (e.g. the 1.6 GHz VLBA image) may not properly split the flux between the core and the closest component. In each case we have taken the far northeastern component to be the core, and have defined the other component positions relative to the position of the core. In the 1.6 GHz VSOP image, the southwestern jet is too complex to be fit with simple Gaussian components, and we left the CLEAN components in this region during model fitting. When fitting elliptical components, the model fitting chi-squared statistic is frequently minimized by an ellipse of zero axial ratio. This is unphysical in the sense that these components have formally infinite brightness temperatures. In these cases an upper limit to the size of the component can be used instead of the best-fit value to find a lower limit to the brightness temperature. Since we use an error analysis method to find these limits in $`\mathrm{\S }`$ 4.3, we have left these zero axial ratio components in the models if they minimize the chi-squared for that model. Table 1 also gives the source frame brightness temperatures for the VSOP models, where the maximum brightness temperature of a Gaussian component is given by
$$T_B=1.22\times 10^{12}\frac{S(1+z)}{ab\nu ^2}\mathrm{K},$$
(1)
where $`S`$ is the flux density of the component in Janskys, $`a`$ and $`b`$ are the FWHMs of the major and minor axes respectively in mas, $`\nu `$ is the observation frequency in GHz, and $`z`$ is the redshift.
## 4 Discussion
### 4.1 Identification of Historical Components
In this section we consider the identification of components seen in these images with previously published VLBI components, using these components’ published positions and velocities. We work from the innermost components outward, starting with the components fit to the 5 GHz VSOP image. We caution that any such identifications are highly speculative, particularly for the older components where much time has elapsed since the last published image. Note that prior to completion of the NRAO VLBA (1995), global VLBI network sessions at 22 GHz occurred only twice per year and used less than half a dozen antennas.
A total of six elliptical Gaussian components are required to fit the 5 GHz VSOP data. The first three of these components are interior to 1 mas, and represent the core and two components of the inner jet. The region interior to 1 mas has been studied by Wehrle et al. (2000) at 22 and 43 GHz and Rantakyrö et al. (1998) at 86 GHz, and they find it to be a complex region with multiple components. Attempts to name components in this region have resulted in some confusion. The component named ‘C5’ in three 11 GHz maps from 1989.3-1991.1 by Carrara et al. (1993) and Abraham & Carrara (1998) is not the same component referred to as either ‘C5’ by Wehrle et al. (1994) at 22 GHz, Leppänen et al. (1995) at 22 GHz, and Lister et al. (1998) at 43 GHz or ‘the stationary 1 mas feature’ by Wehrle et al. (1996) at 22 GHz. Wehrle et al. (1994) and Wehrle et al. (1996) identify a component between C5 and the core, and Leppänen et al. (1995) identify two components in this region which they name ‘C6’ and ‘C7’. Lister et al. (1998) do not detect these components but record another new component named ‘C8’. Clarification of the components in this region must await completion of the Wehrle et al. (2000) analysis.
The situation beyond $``$ 1 mas is easier to interpret. The bright feature at $``$ 3 mas is the component C4 that has been seen by many other authors (Unwin et al. 1989; Carrara et al. 1993; Wehrle et al. 1994; Leppänen et al. 1995; Wehrle et al. 1996; Lister et al. 1998; Wardle et al. 1998; Kellermann et al. 1998). This component has been moving along a position angle of $`115^{}`$ for over 10 years. The brightness distribution of C4 is asymmetric, with the leading edge being sharper than the trailing edge, and two model components are required to represent it (components 5 and 6 in the 5 GHz VSOP model). The 43 GHz data of Wehrle et al. (2000) also require two components to represent the structure of C4. The sharp leading edge of this component is suggestive of a working surface or shock front. The polarization observations of Leppänen et al. (1995), Lister et al. (1998), and Wardle et al. (1998) show that C4 has a magnetic field transverse to the jet, also indicative of a shock front. Carrara et al. (1993) determined a motion of 0.15$`\pm `$0.01 mas/yr (4 $`c`$) for C4, however the position of C4 in our 5 GHz image is inconsistent with a simple extrapolation at this speed. However, we note that at the time of the observations of Carrara et al. C4 was $``$1 mas from the core and in the region noted above as being difficult to interpret. More recent data from Wehrle et al. (2000) and Kellermann (1999, private communication) indicate a speed of $``$$`c`$ for C4 between 1991 and 1999, which is consistent with the position seen in the VSOP images in this paper. Although the larger scale structure has a much different position angle than C4, C4 shows no signs of altering its path to follow the larger scale structure, and appears to be continuing along a position angle of $`115^{}`$. Component C4 dominates the emission in the 1.6 GHz VSOP image and it is represented by component 4 in the 1.6 GHz VSOP model.
Older VLBI measurements (Cotton et al. 1979; Unwin et al. 1989; Carrara et al. 1993) followed a series of components (C1–C3) moving along position angles of $`130`$ to $`140^{}`$. Since this is also the position angle of the larger scale structure seen in our images, it is reasonable that this string of components may form the southwestern jet in our images. The 5 GHz VSOP image shows a fainter, more diffuse component at a position angle of $`133^{}`$ located $``$2.6 mas from the core (component 4 in the 5 GHz model, component 3 in the 1.6 GHz model). Emission is also seen at this position in the images of Kellermann et al. (1998) and Wehrle et al. (2000). An attempt to identify this component with the most recent component ejected along this position angle (C3) would imply a drastic deceleration for C3; a straightforward extrapolation of the motion of C3 given by Carrara et al. (1993) would place it at a separation of 4 mas in 1998 and so we consider such an identification unlikely. An extrapolation of the motion of C2 estimated from the positions given by Unwin et al. (1989) would place C2 between 6 and 10 mas from the core at the epoch of our VSOP observations, and conceivably identify it with either component 5 or 6 in the 5 GHz VLBA model (component 3 in the 1.6 GHz VLBA model). Cotton et al. (1979) reported a large proper motion of 0.5 mas/yr during the early 1970s; their data was re-examined and their interpretation judged to be correct by Unwin et al. (1989). An extrapolation the motion of the Cotton et al. (1979) component (which could also be the component C1 of Unwin et al. located $``$8 mas from the core at 1984.25) would place it about 15 mas from the core in 1998, and mean that it could be identified with the relatively bright component 7 of the 5 GHz VLBA model (component 4 in the 1.6 GHz VLBA model). The uncertainties of extrapolating component motions makes this analysis highly speculative; the validity of the assumption of constant projected speed with time is unclear, and enough time has elapsed that the fate of these older components will probably never be certain.
A series of VLA maps of 3C 279 at varying resolution has been published by de Pater & Perley (1983). Their highest resolution map shows a component (component ‘C’) 95 mas from the core in position angle $`145^{}`$. (This map has a higher resolution than the MERLIN maps of Pilbratt et al. and Akujor et al. because it is at a higher frequency). Our 1.6 GHz VLBA image extends out to this distance from the core, and we see a 140 mJy component at 88 mas in position angle $`141^{}`$ (component 6 in the 1.6 GHz VLBA model). We can tentatively identify this with component C of de Pater & Perley (1983), given the large errors implied by their VLA beam (60 mas resolution) and the similar size of our model-fit component. These errors make it impossible to infer anything about the motion of this component between 1982 and 1998. We have, however, matched the largest scale VLBI structures with the smallest scale VLA structures, and established a continuous connection between the parsec and kiloparsec scales in this source.
### 4.2 Spectral Index Map
Construction of spectral index maps is often hindered by the differing resolutions of the images at different frequencies; however, a unique capability of the VSOP mission is that it can provide matched resolution images to ground-based images at higher frequencies, enabling the construction of a spectral index map from two images of approximately equal resolution. We have used this capability to produce a spectral index map from the 1.6 GHz VSOP image and the 5 GHz VLBA-only image; this spectral index map is shown in Figure 5. These two datasets were taken only 1 day apart, so the errors due to component motions are negligible. To produce the spectral index map the two images were restored with their average beam of 3.12 $`\times `$ 1.14 mas with a major axis position angle of 15.9, each image was restored from the clean components without residuals (the flux left over after the clean components convolved by the dirty beam have been subtracted from the data), and no spectral index was calculated for pixels where the flux densities were less than 3 times the rms noise level in the images at both frequencies. (A spectral index map using images restored from clean components with residuals was also produced and gave essentially the same results. We prefer to use the images without residuals as the residuals can then be used to assess the possible errors in the spectral index map.) In the following discussion we use $`S\nu ^{+\alpha }`$.
A major difficulty in making spectral index maps lies in correctly registering the two images. We investigated several alignments of these two images, including aligning the peak core pixels and aligning the peak pixels in the bright jet component (C4). We doubled the number of pixels across each image in order to measure the required shifts as accurately as possible. Aligning the peak core pixels produced unphysical results, including a highly inverted spectral index along the right edge of component C4. Of the different alignments tried, aligning the peak pixels in the jet component C4 produced the most physically reasonable results. The reason for this can be seen a posteriori from Figure 5. The spectral index is constant across component C4, meaning that the brightest pixels in C4 at 1.6 and 5 GHz will represent the same physical location. On the other hand, there are steep spectral index gradients across the core region, so the peak pixel in the core region will be at different locations at 1.6 and 5 GHz.
We also constructed a map of the error in the spectral index; this error map is shown in Figure 6. The error was calculated by standard propagation of errors, using the fluxes at each pixel in each residual map as the flux errors. This method actually gives a lower limit to the error at each pixel, because it does not take into account calibration errors or errors in imaging the source structure caused by the holes in the $`(u,v)`$ plane coverage. In the region comprising the southwestern jet the errors in the spectral index $`\alpha `$ range from $`\pm 0.05`$ in the brighter parts of the jet (the knot at $``$ 18 mas) to $`\pm 0.3`$ in the fainter parts. The formal errors in the core and C4 regions are quite low and so the errors in these regions will be dominated by the other effects mentioned above.
The core of 3C 279 has an inverted spectrum with steep spectral index gradients. Such inverted spectra are commonly interpreted as being due to self-absorption of the radio synchrotron emission. The calculated spectral index in the core region ranges from $``$ 1.0 at the western edge to the theoretical limiting value for synchrotron self-absorption of 2.5 (assuming a constant magnetic field) over a small region at the eastern edge. Spectral indices approaching this theoretical value are almost never seen; the flatter spectra usually observed are commonly interpreted as being due to an inhomogeneous source made up of a number of synchrotron components with differing turnover frequencies (e.g. Cotton et al. 1980). The highly inverted spectrum at the eastern edge may imply the detection of a homogeneous compact component in this region, which should be an efficient producer of inverse Compton gamma-rays. The spectral index gradients in the core imply that components will have different measured separations at different frequencies. We do indeed observe this frequency-dependent separation, the measured separation between the core and C4 model components is 3.3 mas at 5 GHz and only 2.8 mas at 1.6 GHz. The fact that the apparent position of the core is a function of wavelength is an important verification of the twin exhaust model, which argues that the observed core is that position in the throat of a nozzle where the opacity is of the order of unity.
The jet component C4 has a flat spectrum with $`\alpha `$ approximately 0.25. This is unusual, as jet components usually have steeper spectra ($`\alpha <0`$). The southwestern jet also has structure in its spectral index distribution, with the edges of the jet appearing to have a steeper spectrum than the center. This spectral index structure is related to the apparent limb brightening at 1.6 GHz which, as noted above, should be interpreted with care. AO2 VSOP observations of 3C 279 at 1.6 GHz in which the baselines to the orbiting antenna have a different orientation in the $`(u,v)`$ plane will allow a consistency check on the brightness and spectral index structures transverse to the jet.
### 4.3 High Brightness Temperatures
Space VLBI observations have a major advantage over ground-based observations because they are able to measure higher brightness temperatures. This is because the smallest measurable major and minor axes in equation are proportional to the resolution of the interferometer, which is proportional to 1/(baseline $`\times `$ frequency), so the denominator in equation depends only on baseline length. The brightness temperature limit for ground-based VLBI is $`10^{11}S(1+z)/f^2`$ and that for space VLBI with HALCA is $`10^{12}S(1+z)/f^2`$, where $`S`$ is the flux density in Janskys, $`z`$ is the redshift, and $`f`$ is the smallest size that can be measured expressed as a percentage of the beam size. The improvement gained by space VLBI thus covers the interesting transition region around $`10^{12}`$ K, the nominal inverse Compton brightness temperature limit (Kellermann & Pauliny-Toth 1969).
Observed brightness temperatures are often used to calculate Doppler beaming factors by assuming an intrinsic brightness temperature and using the fact that $`T_{B,obs}=\delta T_{B,int}`$, where $`T_{B,obs}`$ is the observed source frame brightness temperature, $`T_{B,int}`$ is the intrinsic brightness temperature, and $`\delta `$ is the Doppler factor. The intrinsic brightness temperature depends on the physical mechanism imposing the limitation. Kellermann & Pauliny-Toth (1969) showed that inverse Compton losses limit the intrinsic brightness temperatures to $`5\times 10^{11}1\times 10^{12}`$ K, otherwise the source can radiate away most of its energy on a timescale of days. Readhead (1994) proposed that the limiting mechanism is equipartition of energy between the particles and magnetic field, and that intrinsic brightness temperatures are limited to $`5\times 10^{10}1\times 10^{11}`$ K.
To be complete, brightness temperature measurements must be presented with associated errors. Since the brightness temperature depends on the product of the major and minor axes of the model-fit component, and these axis sizes can have large errors, the measured brightness temperature can also have large errors. Error analysis of VLBI model-fit parameters has historically been problematic. In this paper we use the “Difwrap” program<sup>6</sup><sup>6</sup>6http://halca.vsop.isas.ac.jp./survey/difwrap/ (Lovell 2000) to analyze the upper and lower limits on our brightness temperature measurements. This program uses the method described by Tzioumis et al. (1989) in which the parameter of interest is varied in steps around the best-fit value, allowing the other parameters to relax at each step, and the resultant model is then visually compared with the data to determine whether or not the fit remains acceptable. The brightness temperature of a component depends on the flux density and size of the component, and the size depends on three of the Difmap model-fit parameters: the major axis length, the axial ratio, and indirectly on the position angle of the major axis (since different major axis lengths and axial ratios may be allowed at different position angles). The size error analysis therefore searches a three-dimensional cube in parameter space, varying the major axis length and position angle and the axial ratio over all possible combinations given input search ranges and step sizes. A visual inspection is done to determine the goodness of the fit instead of using a numerical cutoff in the chi-squared because the true number of degrees of freedom is not well known. Using the actual number of measured visibilities ($`10^5`$ for these observations) to determine the degrees of freedom gives errors that are unrealistically small, and methods used by other authors did not have a clear physical motivation (e.g. one degree of freedom per antenna per hour \[Biretta, Moore, & Cohen 1986\]).
In Table 2 we show our brightness temperature error analysis for the six components in Table 1 that have best-fit brightness temperatures over $`10^{12}`$ K. For each of these components we searched a $`7\times 7\times 7`$ cube in major axis length, axial ratio, and major axis position angle. Initially we searched major axis lengths from zero to twice the best-fit length, axial ratios from zero to one, and a range of $`\pm 90^{}`$ in position angle; and then refined the search to a smaller grid if necessary. The parameter values yielding the maximum and minimum area that still gave an acceptable fit to the data were recorded. A similar error analysis was done for the flux density, and the extreme allowed values of area and flux density were used to determine the maximum and minimum brightness temperatures. Since errors in flux density and size are searched for separately, the flux density was held constant during the size error analysis and vice-versa. The position of the component was also held constant to avoid it ‘trading identities’ with another model component. All other model components were allowed to vary.
Inspection of Table 2 shows that in all but one of the cases investigated the component shrinking to zero area and infinite brightness temperature (in all cases caused by a valid fit with zero axial ratio at some position angle) produced acceptable results, and therefore it appears that many measured brightness temperatures, even those measured by space VLBI, may have error bars that extend to infinity in the positive direction. The measured lower limits also indicate a considerable error in the best-fit brightness temperature values: three of the six components with best-fit brightness temperatures over $`10^{12}`$ K have minimum brightness temperatures under this value. The other three components have minimum brightness temperatures over $`10^{12}`$ K, with minimum brightness temperatures of 1.2, 2.0, and 4.9 $`\times 10^{12}`$ K being measured for components 1 and 2 of the 5 GHz VSOP model fit (the core and first jet component) and component 4 of the 1.6 GHz VSOP model fit (C4) respectively. Bower & Backer (1998) and Shen et al. (1999) report brightness temperatures of $`3\times 10^{12}`$ K from VSOP observations of NRAO 530 and PKS 1921$``$293 respectively, but without accompanying error analyses.
If a brightness distribution other than a Gaussian is used in the model fitting, the derived values of the brightness temperature will be different. For example, the brightness temperature of a homogeneous optically thick component is given by
$$T_B=1.77\times 10^{12}\frac{S(1+z)}{ab\nu ^2}\mathrm{K},$$
where the constant in front is different from that in equation (1), and $`a`$ and $`b`$ are the lengths of the major and minor axes respectively rather than the FWHMs. The visibility of a homogeneous optically thick component drops to 50% at the same baseline length as a Gaussian when its diameter equals 1.6 times the Gaussian’s FWHM (Pearson 1995), so we expect the homogeneous optically thick brightness temperature to be about 0.6 of the Gaussian brightness temperature (see also Hirabayashi et al. 1998 and in particular the correction in the erratum to this paper). We have fit homogeneous optically thick components to the data, and for the two components in Table 2 where neither component type goes to zero size (components 2 and 5 of the 5 GHz model fit) we measure brightness temperatures of 1.9 and 1.5 $`\times 10^{12}`$ K respectively for homogeneous optically thick components rather than 2.9 and 2.6 $`\times 10^{12}`$ K for Gaussian components, so we see about the expected decrease. Since the true brightness distribution is not known, and Gaussian components are the standard for VLBI model fitting and provide a somewhat better fit for these observations, we remained with Gaussian components.
Our highest brightness temperature lower limit of $`5\times 10^{12}`$ K for component C4 at 1.6 GHz implies Doppler factor lower limits of 5 and 50 for the inverse Compton and equipartition brightness temperature limits respectively. A Doppler factor of 50 is at the upper end of the Doppler factor distributions expected for flux-limited samples of flat-spectrum radio sources (Lister & Marscher 1997) and gamma-ray sources (Lister 1998). Bower & Backer (1998) found similar values for the Doppler factor of NRAO 530 under these same two limiting conditions. If VSOP observations reveal a brightness temperature much higher than $`5\times 10^{12}`$ K, or many brightness temperatures around $`5\times 10^{12}`$ K, it may be difficult to reconcile the high Doppler factors implied by the equipartition brightness temperature limit with beaming statistics and with the relatively slow speeds measured in studies of apparent velocity distributions (e.g. Vermeulen 1995).
Using an estimated speed for C4 of 7 $`c`$ (see $`\mathrm{\S }`$ 4.1), and assuming the pattern speed observed with VLBI equals the bulk fluid speed in the jet, bulk Lorentz factors and angles to the line-of-sight can be calculated for the jet. For $`\delta =5`$, $`\mathrm{\Gamma }`$=7.5 and $`\theta `$=11; for $`\delta =50`$, $`\mathrm{\Gamma }`$=25.5 and $`\theta `$=0.3. Again, the equipartition brightness temperature limit implies values for 3C 279 near the extremes of expected EGRET source properties (Lister 1998). 3C 279 and NRAO 530 have both been detected by EGRET, and Bower & Backer (1998) speculate that blazars detected by EGRET may be those where the equipartition brightness temperature limit is briefly (on a timescale of years) superseded by the inverse Compton catastrophe limit. The observations presented in this paper occurred at the beginning of a total flux density flare at 5 GHz recorded by the Michigan monitoring program that would later reach the highest flux density yet recorded in this program for 3C 279 at 5 GHz. Measurements of the variability brightness temperature of this flare (Lähteenmäki, Valtaoja, & Wiik 1999) together with VSOP brightness temperatures measured during AO2 should allow calculation of the intrinsic brightness temperature and Doppler factor and allow us to estimate any departures from equipartition in this source.
## 5 Conclusions
We have presented the first space VLBI images of 3C 279, which are the highest resolution images yet obtained of this source at 5 and 1.6 GHz. The parsec-scale emission is dominated by the core and the jet component C4 which has been visible in VLBI images since 1985. The 1.6 GHz VSOP image and the 5 and 1.6 GHz VLBA-only images show emission from a jet extending to the southwest. The 1.6 GHz VLBA-only image has structure that matches that seen in the highest resolution VLA images, connecting the parsec and kiloparsec scale structures in this source.
We have exploited two of the main strengths of VSOP: the ability to obtain matched-resolution images to ground-based images at higher frequencies and the ability to measure high brightness temperatures. The spectral index map constructed from the 1.6 GHz VSOP image and the 5 GHz VLBA-only image has an unusually inverted spectral index in the core region, approaching the limiting value for synchrotron self-absorption of $`+2.5`$. An extensive error analysis conducted on the model-fit brightness temperatures reveals brightness temperature lower limits as high as $`5\times 10^{12}`$ K. This lower limit is significantly above both the nominal inverse Compton and equipartition brightness temperature limits. The derived Doppler factor, Lorentz factor, and angle to the line-of-sight in the case of the equipartition limit are at the upper end of the range of expected values for EGRET blazars.
Part of the work described in this paper has been carried out at the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration. A.E.W. acknowledges support from the NASA Long Term Space Astrophysics Program. We gratefully acknowledge the VSOP Project, which is led by the Japanese Institute of Space and Astronautical Science in cooperation with many organizations and radio telescopes around the world. The National Radio Astronomy Observatory is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc. This research has made use of data from the University of Michigan Radio Astronomy Observatory which is supported by the National Science Foundation and by funds from the University of Michigan, and the NASA/IPAC extragalactic database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration.
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# New Dark Matter Physics: Clues from Halo Structure
## I How Cold and How Collisionless is the Dark Matter?
The successful concordance of predictions and observations of large scale structure and microwave anisotropy vindicates many assumptions of standard cosmology, in particular the hypothesis that the dark matter is composed of primordial particles which are cold and collisionless. At the same time, there are hints of discrepancies observed in the small-scale structure within galaxy haloes, which we explore as two related but separate issues, namely the predictions of excessive substructure and sharp central cusps in dark matter halos.
The first “substructure problem” is that CDM predicts excessive relic substructure: much of the mass of a CDM halo is not smoothly distributed but is concentrated in many massive sublumps, like galaxies in galaxy clusters. The model predicts that galaxy halos should contain many dwarf galaxies which are not seen, and which would disrupt disks even if they are invisible. The substructure problem appears to be caused by the “bottom-up” hierarchical clustering predicted by CDM power spectra; fluctuations on small scales collapse early and survive as dense condensations. Its absence hints that the small scale power spectrum is filtered to suppress early collapse on subgalactic scales.
The second “cusp problem” is that CDM also predicts a universal, monotonic increase of density towards the center of halos which is not seen in close studies of dark-matter-dominated galaxies (although the observational issue is far from settled). The formation of central cusps has been observed for many years in simulations of collapse of cold matter in a wide variety of circumstances; it may be thought of as low-entropy material sinking to the center during halo formation. Simulations suggest that dynamical “pre-heating” of CDM by hierarchical clustering is not enough to prevent a cusp from forming— that some material is always left with a low entropy and sinks to the center. If this is right, the central structure of halos might provide clues to the primordial entropy which are insensitive to complicated details of nonlinear collapse.
It may be possible to explain these discrepancies in a CDM framework, for example by using various baryonic contrivances. It is also possible that the observations can be interpreted more sympathetically for CDM; we explore this possibility in more detail in a separate paper. However it is also possible that the problems with halo structure are giving specific quantitative clues about new properties of the dark matter particles. By examining halo structure and stability, in this paper we make a quantitative assessment of the effects of modifications of the two main properties of CDM— the addition of primordial velocity dispersion, and/or the addition of particle self-interactions. In particular we focus on aspects of halo structure which provide the cleanest “laboratories” for studying dark matter properties. The ultimate goal of this exercise is to measure and constrain particle properties from halo structure.
Endowing the particles with non-zero primordial velocity dispersion produces two separate effects: a filter in the primordial power spectrum which limits small-scale substructure, and a phase packing or Liouville limit which produces halo cores. Both effects depend on the same quantity, the “phase density” which we choose to define using the most observationally accessible units, $`Q\rho /v^2^{3/2}`$, where $`\rho `$ is the density and $`v^2`$ is the velocity dispersion. The definitions of these quantities depend on whether we are discussing fine-grained or coarse-grained $`Q`$.<sup>*</sup><sup>*</sup>*For a uniform monatomic ideal thermal gas, $`Q`$ is related in a straightforward way to the usual thermodynamic entropy; for $`N`$ particles, $`S=kN[\mathrm{ln}(Q)+\mathrm{constant}].`$ For collisionless particles, the fine-grained $`Q`$ does not change but the coarse-grained $`Q`$ can decrease as the sheet occupied by particles folds up in phase space. The coarse-grained $`Q`$ can be estimated directly from astronomical observations, while the fine-grained $`Q`$ relates directly to microphysics of dark matter particles. For particles which decouple when still relativistic, the initial microscopic phase density $`Q_0`$, which for nondissipationless collisionless particles is the maximum value for all time, can be related to the particle mass and type, with little reference to the cosmology. The most familiar examples are the standard neutrinos, but we include in our discussion the more general case which yields different numerical factors for bosons and for particles with a significant chemical potential.
The physics of both the filtering and the phase packing in the collisionless case closely parallels that of massive neutrinos, the standard form of “hot” dark matter. Dominant hot dark matter overdoes both of these effects— the filtering scale is too large to agree with observations of galaxy formation (both in emission and quasar absorption) and the phase density is too low to agree with observations of giant-galaxy halos. However one can introduce new particles with a lower velocity dispersion (“warm” dark matter, ), which is the option we consider here.
Although warm dark matter has most often been invoked as a solution to fixing apparent (and no longer problematic) difficulties with predictions of the CDM power spectrum for matching galaxy clustering data, a spectrum filtered on smaller scales may also solve several other classic problems of CDM on galactic and subgalactic scales which are sometimes attributed to baryonic effects. The main effect in warm models is that the first nonlinear objects are larger and form later, suppressing substructure and increasing the angular momentum of galaxies. This improves the predictions for dwarf galaxy populations, baryon-to-dark-matter ratio, disk size and angular momentum, and quiet flows on the scale of galaxy groups. If the filtering is confined to small scales the predictions are likely to remain acceptable for Lyman-$`\alpha `$ absorption during the epoch of galaxy formation at $`z3`$.
Liouville’s theorem tells us that dissipationless, collisionless particles can only decrease their coarse-grained phase density, and we conjecture that halo cores on small scales approximately preserve the primordial phase density. The universal character of the phase density allows us to make definite predictions for the scaling of core density and core radius with halo velocity dispersion. These relations are analogous to those governing nonrelativistic degenerate-dwarf stars: more tightly bound (i.e. massive) halos should have smaller, denser cores. A survey of available evidence on the phase density of dark matter cores on scales from dwarf spheroidal galaxies to galaxy clusters shows that the phase density needed to create the cores of rotating dwarf galaxies is much lower than that apparently present in dwarf spheroidal galaxies— so at least one of these populations is not probing primordial phase density. Translating into masses of neutrino-like relics, the spheroidals prefer masses of about 1 keV (unless the observed stars occupy only a small central portion of an implausibly large, massive and high-dispersion halo), and the disks prefer about 200 eV.
The larger phase density is also preferred from the point of view of filtering. If we take $`\mathrm{\Omega }0.3`$ (instead of 1 as in most of the original warm scenarios— which reduces the scale for a given mass, because it lowers the temperature and therefore the number of the particles), the filtering scale for 1 keV particles is at about $`k=3\mathrm{M}\mathrm{p}\mathrm{c}^1`$— small enough to preserve the successful large-scale predictions of CDM but also large enough to impact the substructure problem. Galaxy halo substructure therefore favors a primordial phase density corresponding to collisionless thermal relics with a mass of around 1 keV. In this scenario the densest dwarf spheroidals might well preserve the primordial phase density and in principle could allow a measurement of the particle mass.This raises another unresolved issue: whether the filtering actually prevents systems as small as dwarf spheroidals from forming at all. The predictions of warm dark models are not yet worked out enough to answer this question. (Conversely, a mass as large as 1 keV can only solve the core problem in disks with additional nonlinear dynamical heating, so that the central matter no longer remains on the lowest adiabat, or with the aid of baryonic effects.)
To have the right mean density and phase density today, relativistically-decoupling particles of this phase density must have separated out at least as early as the QCD era, when the number of degrees of freedom was much larger than at classical weak decoupling. Their interactions with normal Standard Model particles must therefore be “weaker than weak,” ruling out not only standard neutrinos but many other particle candidates. The leading CDM particle candidates, such as WIMPs and axions, form in standard scenarios with much higher phase densities, although more elaborate mechanisms are possible to endow these particles with the velocities to dilute $`Q`$. We review briefly some of the available options for making low-$`Q`$ candidates, such as particles decaying out of equilibrium.
A new wrinkle on this story comes if we endow the particles with self-interactions. We consider a simple parametrized model of particle self-interactions based on massive intermediate particles of adjustable mass and coupling, and explore the constraints on these parameters from halo structure. Self-interactions change the filtering of the power spectrum early on, and if they are strong enough they qualititatively change the global structure and stability of halos.
In the interacting case, linear perturbations below the Jeans scale oscillate as sound waves instead of damping by free streaming— analogous to a baryon plasma rather than a neutrino gas. This effect introduces a filter which is sharper in $`k`$ than that from streaming, and also on a scale about ten times smaller than the streaming for the same rms particle velocity— about right to reconcile the appropriate filtering scale with the $`Q`$ needed for phase-density-limited disk cores. These self-interactions could be so weak that the particles are effectively collisionless today as in standard CDM.
On the other hand stronger self-interactions have major effects during the nonlinear stages of structure formation and on the structure of galaxy halos. We consider this possibility in some detail, using Lane-Emden polytropes as fiducial models for collisional halos. Their structures are close analogs of degenerate dwarf stars and we call them “giant dwarfs”. We find that these structures are subject to an instability caused by heat conduction by particle diffusion.Degenerate dwarf stars are not subject to this instability because they are supported without a temperature gradient; the same stabilization could occur in halo cores only if the dark matter is fermionic and degenerate (e.g., ). The instability we discuss here is essentially what happens in a thermally-supported star with no nuclear reactions, except that the conduction is by particle diffusion rather than by radiation. This effect may have already been observed numerically. Although a little of this might be interesting (e.g. leading to the formation of central black holes or to high-density, dwarf spheroidal galaxies), typical halos can only be significantly collisional if they last for a Hubble time; for this to be the case, the particle interactions must be so strong that diffusion is suppressed, which in turn requires a fluid behavior for all bound dark matter structures. This option is not very attractive from a phenomenological point of view; for example, dwarf galaxies or galaxies in clusters tend to sink like rocks instead of orbiting like satellites, and the collapse of cores occurs most easily in those low-dispersion halos where we seek to stabilize them.
## II Particle Properties and Phase Densities
We adopt the hypothesis that some dark matter cores are real and due to dark matter rather than baryonic physics or observational artifacts. At present this interpretation is suggested rather than proven by observations. We also conjecture that the heating which sets the finite central phase density is primordial, part of the physics of the particle creation rather than some byproduct of hierarchical clustering. At present this is a conjecture suggested rather than proven by simulations.
In the clustering hierarchy, more higher-entropy material is created as time goes on, but numerical experiments indicate that this heated material tends to end up in the outer halo. This is the basic reason why CDM halos always have central cusps: there is always a little bit of material which remembers the low primordial entropy and sinks to the center. The halo center contains the lowest-entropy material, which we conjecture is a relic of the original entropy of the particles— or equivalently, their original phase density, which is most directly related to measurable properties of halo dynamics. We begin by relating the phase density to particle properties in some simple models.
### A Phase Density of Relativistically-Decoupled Relics
Consider particles of mass $`m`$ originating in equilibrium and decoupling at a temperature $`T_D>>m`$ or chemical potential $`\mu >>m`$. The original distribution function is
$$f(\stackrel{}{p})=(e^{(E\mu )/T_D}\pm 1)^1(e^{(p\mu )/T_D}\pm 1)^1$$
(1)
with $`E^2=p^2+m^2`$ and $`\pm `$ applies to fermions and bosons respectively. The number density and pressure of the particles are
$$n=\frac{g}{(2\pi )^3}fd^3p$$
(2)
$$P=\frac{g}{(2\pi )^3}\frac{p^2}{3E}fd^3p$$
(3)
where $`g`$ is the number of spin degrees of freedom. Unless stated otherwise, we adopt units with $`\mathrm{}=c=1`$.
With adiabatic expansion this distribution is preserved with momenta of particles varying as $`pR^1`$, so the density and pressure can be calculated at any subsequent time. For thermal relics $`\mu =0`$, we can derive the density and pressure in the limit when the particles have cooled to be nonrelativistic:
$$n=\frac{gT_0^3}{(2\pi )^3}\frac{d^3p}{e^p\pm 1}$$
(4)
$$P=\frac{gT_0^5}{(2\pi )^33m}\frac{p^2d^3p}{e^p\pm 1}$$
(5)
where the pseudo-temperature $`T_0=T_D(R_D/R_0)`$ records the expansion of any fluid element relative to its initial size and temperature at decoupling $`R_D,T_D`$.
It is useful to define a “phase density” $`Q\rho /v^2^{3/2}`$ proportional to the inverse specific entropy for nonrelativistic matter, which is preserved under adiabatic expansion and contraction. For nondissipative particles $`Q`$ cannot increase, although it can decrease due to shocks (in the collisional case) or coarse-graining (in the collisionless case, e.g. in “violent relaxation” and other forms of dynamical heating.) Combining the above expressions for density and pressure and using $`v^2=3P/nm`$, we find
$$Q_X=q_Xg_Xm_X^4.$$
(6)
The dimensionless coefficient for the thermal case is
$$q_T=\frac{4\pi }{(2\pi )^3}\frac{[𝑑p(p^2/e^p\pm 1)]^{5/2}}{[𝑑p(p^4/e^p\pm 1)]^{3/2}}=0.0019625,$$
(7)
where the last equality holds for thermal fermions. An analogous calculation for the degenerate fermion case ($`T=0,\mu _D>>m_X`$) yields the same expression for $`Q`$ but with a different coefficient,
$$q_d=\frac{4\pi }{(2\pi )^3}\frac{[_0^1p^2𝑑p]^{5/2}}{[_0^1p^4𝑑p]^{3/2}}=0.036335.$$
(8)
To translate from $`\mathrm{}=c=1`$ into more conventional astronomers’ units,
$$(100eV)^4/c^5=12,808\frac{(M_{}/\mathrm{kpc}^3)}{(\mathrm{km}\mathrm{s}^1)^3}=12.808\frac{(M_{}/\mathrm{pc}^3)}{(100\mathrm{km}\mathrm{s}^1)^3}$$
(9)
The phase density in this situation depends on the particle properties but not at all on the cosmology; the decoupling temperature, the current temperature and density do not enter. The numerical factors just depend on whether the particles are thermal or degenerate, bosons or fermions, which makes the quantity $`Q`$ a potentially precise tool for measuring particle properties. Many scenarios envision thermal relics so we adopt this as a fiducial reference in quoting phase densities in $`m^4`$ units—bearing in mind that the actual mass may be different in cases such as degnerate sterile neutrinos, and that for the astrophysical effects discussed below, it is the phase density that matters. For a neutrino-like ($`g=2`$), thermal relic,
$$Q_T=5\times 10^4\frac{(M_{}/\mathrm{pc}^3)}{(\mathrm{km}\mathrm{s}^1)^3}(m_X/1\mathrm{k}\mathrm{e}\mathrm{V})^4.$$
(10)
### B Space Density of Thermal Relics
For a standard, relativistically-decoupled thermal relic, the mean density of the particles can be estimated from the number of particle degrees of freedom at the epoch $`T_D`$ of decoupling, $`g_D`$; the ratio to the critical density is
$$\mathrm{\Omega }_X=78.3h^2[g_{eff}/g_D](m_X/1\mathrm{k}\mathrm{e}\mathrm{V})=2.4h_{70}^2(m_X/1\mathrm{k}\mathrm{e}\mathrm{V})(g_{eff}/1.5)(g_D/100)^1$$
(11)
where $`g_{eff}`$ is the number of effective photon degrees of freedom of the particle ($`=1.5`$ for a two-component fermion). For standard neutrinos which decouple at around 1MeV, $`g_D=10.75`$.
Current observations suggest that the dark matter density $`\mathrm{\Omega }_X0.3`$ to $`0.5`$, hence the mass density for a warm relic with $`m_X`$ 200 eV clearly requires a much larger $`g_D`$ than the standard value for neutrino decoupling. Above about 200 MeV, the activation of the extra gluon and quark degrees of freedom (24 and 15.75 respectively including $`uds`$ quarks) give $`g_D50`$; activation of heavier modes of the Standard Model above $`200`$GeV produces $`g_d100`$; this gives a reasonable match for $`m_X200`$ eV and $`\mathrm{\Omega }_X0.5`$, as suggested by current evidence. Masses of the order of 1 keV can be accomodated by somewhat earlier decoupling ($``$ TeV) and including many extra (e.g., supersymmetric or extra-dimensional) degrees of freedom. Alternatively a degenerate particle can be introduced via mixing of a sterile neutrino, combined with a primordial chemical potential adjusted to give the right density. In any of these cases, the particle must interact with Standard Model particles much more weakly than normal weak interactions, which decouple at $`1`$ MeV.
Note that warm dark matter particles have low densities compared with photons and other species at 1 MeV so they do not strongly affect nucleosynthesis. However, their effect is not entirely negligible since they are relativistic at early times and add considerably more to the mean total density in the radiation era than standard CDM particles. They add the equivalent of $`(T_X/T_\nu )^3=10.75/g_d`$ of an effective extra neutrino species, which leads to a small increase in the predicted primordial helium abundance for a given $`\eta `$. Because the phase density fixes the mean density at which the particles become relativistic, it also fixes this effect on nucleosynthesis (independent of the other particle properties, thermal or degenerate etc.) This effect might eventually become detectable with increasingly precise measurements of cosmic abundances.
### C Decaying WIMPs and Other Particle Candidates
Thermally decoupled relics are the simplest way to obtain the required finite phase density, but they are not the only way. Heavier particles can be produced with a kinetic temperature higher than the radiation, accelerated by some nonthermal process. Weakly interacting massive particles, including the favored Lightest Supersymmetric Particles, can reduce their phase density if they form via out-of-equilibrium particle decay. A small density of heavy unstable particles (X1) can separate out in the standard way, then later decay into the present-day (truly stable) dark matter particles (X2). In a supersymmetric scheme one can imagine for example a gravitino separating out and decaying into neutralino dark matter.
In the normal Lee-Weinberg scenario for WIMP generation, the particle density is in approximate thermal equilibrium until $`Tm_X/20`$. The particles thin out by annihilation until their relic density freezes out when the the annihilation rate matches the Hubble rate, $`n_X\sigma _{ann}vH`$. The density today is then
$$\mathrm{\Omega }_XT_{\gamma 0}^3H_0^2m_{Planck}^3\sigma _{ann}v^1(m_W/100\mathrm{G}\mathrm{e}\mathrm{V})^2(m_W/m_X)^2$$
(12)
where we have used the typical weak annihilation cross section $`\sigma _{ann}\alpha ^2m_X^2/m_W^4`$ determined by the mass of the $`W`$. The kinetic temperature of the WIMPs freezes out at about the same time as the abundance, so they are very cold today, with typical velocities $`v\sqrt{20}T_0/m_X10^{14}(m_X/100GeV)^1`$. This of course endows them with small velocities and an enormous phase density.
A smaller phase density can be produced if these particles decay at some point into the particles present today. If the secondary particles are much lighter than the first, they can be generated with relativistic velocities at relatively late times as we require. Suppose the primary $`X1`$ particles decay into secondary $`X2`$ particles at a temperature $`T_{decay}`$. To produce particles with the velocity $`0.4`$ km/sec today (characteristic of a fiducial 200 eV thermal relic phase density), or $`vc`$ at $`T/300\mathrm{e}\mathrm{V}`$,
$$m_{X1}m_{X2}T_{decay}/300\mathrm{e}\mathrm{V}.$$
(13)
We also want to get the right density of $`X2`$ particles. Suppose the density of $`X1`$ is determined by a Lee-Weinberg freezeout, such that $`n_{X1}(T_\gamma m_{X1}/20)\sigma _{ann}vH`$. In order to have $`\rho _{X2}\rho _{rad}/600`$ at $`z_{nr}10^6`$, $`\rho _{X1}\rho _{rad}/600`$ at $`T_{decay}`$, and then
$$m_{X2}^2T_{decay}\frac{600\times 20m_W^4}{\alpha ^2m_{Planck}}(100\mathrm{M}\mathrm{e}\mathrm{V})^3.$$
(14)
Thus we obtain
$$m_{X1}m_{X2}(30\mathrm{G}\mathrm{e}\mathrm{V})^2.$$
(15)
A simple example would be a more or less standard 50 GeV WIMP primary which decays at $`T_{decay}1`$keV into marginally relativistic 20 GeV secondaries. Alternatively the primary could be heavier than this and the secondary lighter. Such scenarios have to be crafted to be consistent with various constraints, such as the long required lifetime for $`X1`$ (in the example just given, a week or so) and the decay width of the $`Z`$ (which must not notice the existence of $`X2`$); although not compelling, they are not all ruled out.<sup>§</sup><sup>§</sup>§It is also possible to reduce the scale of filtering of linear perturbations for a given phase density by arranging for the decay relatively late, and for the decay products to be nonrelativistic. This option seems even more contrived and we will not pursue it in detail here.
The other perennial favorite dark matter candidate is the axion. The usual scenario is to produce these by condensation, which if homogeneous produces dark matter even colder than the WIMPs— indeed, as bosons in a macroscopic coherent state. However, it is natural to contemplate modifications to this picture where the condensing fields are not uniform but have topological defects or Goldstone excitations, produced by the usual Kibble mechanism during symmetry breaking (e.g. ). In this case the axions are produced with relativistic velocities and could in principle lead to the desired velocity dispersion.
## III Cores from Finite Primordial Phase Density
We have shown several examples of how particle properties determine primordial phase density. Here we explore how the phase density affects the central structure of dark matter halos.
### A Core Radius of an Isothermal Halo
Consider the evolution of classical dissipationless, collisionless particles in phase space. Truly Cold Dark Matter is formed with zero velocity dispersion occupying a three dimensional subspace (determined by the Hubble flow $`\stackrel{}{v}=H\stackrel{}{r}`$) of six dimensional phase space. Subsequent nonlinear evolution wraps up the phase sheet so that a coarse-grained average gives a higher entropy and a lower phase density. In general a small amount of cold material remains which naturally sinks to the center of a system. There is in principle no limit to the central density; the phase sheet can pack an arbitrary number of phase wraps into a small volume.
By contrast, with warm dark matter the initial phase sheet has a finite thickness. The particles do not radiate so the phase density can never exceed this initial value. In the nonlinear formation of a halo, the phase sheet evolves as an incompressible fluid in phase space. The outer parts of a halo form in the same way as CDM by wraps of the sheet whose thickness is negligible, but in the central parts the finite thickness of the sheet prevents arbitrarily close packing— it reaches a “phase packing” limit. For a given velocity dispersion at any point in space, the primordial phase density of particles imposes an upper limit on their density $`\rho `$, corresponding to adiabatic compression. Thus warm dark matter halos cannot form the singular central cusps predicted by Cold Dark Matter but instead form cores with a maximum limiting density at small radius, determined by the velocity dispersion.
We estimate the structure of the halo core by conjecturing that the matter in the central parts of the halo lies close to the primordial adiabat defined by $`Q`$. This will be a good assumption for cores which form quietly without too much dynamical heating. Simulations indicate this to be the case in essentially all CDM halos, although in principle it could be that warm matter typically experiences more additional dynamical heating than cold matter, in which case the core could be larger. This question can be resolved with warm simulations, including a reasonable sampling of the particle distribution function during nonlinear clustering; for the present we derive a rigorous upper limit to the core density for a given velocity dispersion, and conjecture that this will be close the actual central structure.
A useful model for illustration and fitting is a standard isothermal sphere model for the halo. The spherical case with an isotropic distribution of velocities maximizes the central density compatible with the phase density limit. The conventional definition of core size in an isothermal sphere is the “King radius”
$$r_0=\sqrt{9\sigma ^2/4\pi G\rho _0}$$
(16)
where $`\sigma `$ denotes the one-dimensional velocity dispersion, and $`\rho `$ denotes the central density. Making the adiabatic assumption, $`\rho _0=Q(3\sigma ^2)^{3/2}`$, we find
$$r_0=\sqrt{9\sqrt{2}/4\pi 3^{3/2}}(QGv_c\mathrm{})^{1/2}=0.44(QGv_c\mathrm{})^{1/2}$$
(17)
where $`v_c\mathrm{}=\sqrt{2}\sigma `$ denotes the asymptotic circular velocity of the halo’s flat rotation curve. (Note that aside from numerical factors this is the same mass-radius relation as a degenerate dwarf star; the galaxy core is bigger than a Chandrasekhar dwarf of the same specific binding energy by a factor $`(m_{proton}/m_X)^2`$. The collisional case treated below is even closer to a scaled version of a degenerate dwarf star.)
For the thermal and degenerate phase densities derived above,
$$r_{0,thermal}=5.5\mathrm{kpc}(m_X/100\mathrm{e}\mathrm{V})^2(v_c\mathrm{}/30\mathrm{k}\mathrm{m}\mathrm{s}^1)^{1/2}$$
(18)
$$r_{0,degenerate}=1.3\mathrm{kpc}(m_X/100\mathrm{e}\mathrm{V})^2(v_c\mathrm{}/30\mathrm{k}\mathrm{m}\mathrm{s}^1)^{1/2},$$
(19)
where we have set $`g=2`$. The circular velocity in the central core displays the harmonic behavior $`v_cr`$; it reaches half of its asymptotic value at a radius $`r_{1/2}0.4r_0`$.
Instead of fitting an isothermal sphere to an entire rotation curve, in some situations we might opt to measure the central density directly by fitting the linear inner portion of a rotation curve if it is well-resolved in the core:
$$v_c/r=\sqrt{4\pi G\rho /3}=2.77G^{1/2}Q^{1/2}v_c\mathrm{}^{3/2}=6.71\mathrm{km}\mathrm{s}^1\mathrm{kpc}^1(m_X/100\mathrm{e}\mathrm{V})^2(v_c\mathrm{}/30\mathrm{k}\mathrm{m}\mathrm{s}^1)^{3/2}.$$
(20)
### B Comparison with galaxy and cluster data
In a separate paper we review the current relevant data in more detail, including a consideration of interpretive ambiguities. Here we offer a summary of the situation.
The relationship of core radius or central density with halo velocity dispersion is a simple prediction of the primordial phase density hypothesis, which can be in principle be tested on a cosmic population of halos. In particular if phase packing is the explanation of dwarf galaxy cores, the dark matter cores of giant galaxies and galaxy clusters are predicted to be much smaller than for dwarfs, unobservably hidden in a central region dominated by baryons. There is currently at least one well-documented case of a galaxy cluster with a large core ($`30`$kpc) as measured by a lensing fit, which cannot be explained at all by phase packing with primordial phase density. On the other hand more representative samples of relaxed clusters do not show evidence of cores.
The favorite laboratories for finding evidence of dark matter cores are dwarf disk galaxies which display a central core even at radii where the baryonic contribution is negligible. Rotation curves allow a direct estimate of the enclosed density as a function of radius, right out to a fairly flat portion which allows an estimate of the dark matter velocity dispersion— all the information we need to estimate a phase density for a core. Three of the best-resolved cases yield estimates of $`Q10^710^6(M_{}/\mathrm{pc}^3)/(\mathrm{km}\mathrm{s}^1)^3`$. The sensitive dependence of $`Q`$ on particle mass means that $`m_X`$ is reasonably well bounded even from just from a handful of such cases; a thermal value of $`m_X`$300 eV does not produce large enough cores to help at all (that is, one must seek unrelated explanations of the data), while values $`m_X`$100 eV produce such large cores that they conflict with observed rotation curves of normal giant galaxies and LSB galaxies. This is why we adopt a fiducial reference value of 200 eV for dwarf disk cores.
Dwarf spheroidal galaxies do not have gas on circular orbits so their dynamics is studied with stellar velocity dispersions. Here we have an estimate of the mean density in the volume encompassed by the stellar test particles, but we do not know the velocity dispersion of the dark matter halo particles (which may larger than that of the stars if the latter occupy only the harmonic central portion of a large dark matter core) so estimates of the phase density are subject to other assumptions and modeling constraints. If we assume that the stars are not much more concentrated than the dark matter, we get the largest estimateThis is the largest value of the mean phase density of material in the region enclosed by the stellar velocity tracers; there is no real observational upper limit for the maximum phase density. Without the rotation curve information, these systems are consistent with singular isothermal spheres or other cuspy profiles for the dark matter of the phase density, which in the largest case is about $`Q2\times 10^4(M_{}/\mathrm{pc}^3)/(\mathrm{km}\mathrm{s}^1)^3`$ corresponding to a thermal relic of mass $`m_X800`$eV. The apparent phase densities estimated for dwarf spheroidals are thus much larger than for dwarf disks, even at the same radius. The mass-to-light ratio in the most extreme of these systems is about 100 in solar units, an order of magnitude more than that found for purely baryonic, old stellar populations in elliptical galaxies, so there is little doubt that they are dominated by dark matter. The CDM prediction is that there are other, more weakly bound halos in which gas was unable to cool and form stars, and which therefore have an even higher mass-to-light ratio.
## IV Filtering of Small-Scale Fluctuations
The non-zero primordial velocity dispersion naturally leads to a filtering of the primordial power spectrum. The transfer function of Warm Dark Matter is almost the same as Cold Dark Matter on large scales, but is filtered by free-streaming on small scales. The characteristic wavenumber for filtering at any time is given by $`k_X=H/v^2^{1/2}`$, the inverse distance travelled by a particle at the rms velocity in a Hubble time. The detailed shape of the transfer function depends on the detailed evolution of the Boltzmann equation, and in particular whether the particles are free-streaming or collisional.
In the current application, we are concerned with $`H`$ during the radiation-dominated era ($`z10^4`$), so that $`H^2=8\pi G\rho _{rel}/3(1+z)^4`$, where $`\rho _{rel}`$ includes all relativistic degrees of freedom. For constant $`Q`$, $`v^2^{1/2}=(\rho _X/Q)^{1/3}(1+z)`$ as long as the $`X`$ particles are nonrelativistic. For particles with a small velocity dispersion today, the comoving filtering scale is thus approximately independent of redshift over a considerable interval of redshift (see Figure 1). The “plateau” scale is independent of $`H_0`$:
$$k_{X,comov}=H_0\mathrm{\Omega }_{rel}^{1/2}v_{X0}^1=0.65\mathrm{Mpc}^1(v_{X0}/1\mathrm{k}\mathrm{m}\mathrm{s}^1)^1$$
(21)
where $`\mathrm{\Omega }_{rel}=4.3\times 10^5h^2`$ is the density in relativistic species and $`v_{X0}=(Q/\overline{\rho }_{X0})^{1/3}`$ is the rms velocity of the particles at their present mean cosmic density $`\overline{\rho }_{X0}`$. For the thermal case, in terms of particle mass, we have
$$v_{X0,thermal}=0.93\mathrm{km}\mathrm{s}^1h_{70}^{2/3}(m_X/100\mathrm{e}\mathrm{V})^{4/3}(\mathrm{\Omega }_X/0.3)^{1/3}(g/2)^{1/3},$$
(22)
and hence
$$k_{X,comov}=15\mathrm{Mpc}^1h_{70}^{2/3}(m_X/1\mathrm{k}\mathrm{e}\mathrm{V})^{4/3}(\mathrm{\Omega }_X/0.3)^{1/3}(g/2)^{1/3}.$$
(23)
In the case of free-streaming, relativistically-decoupled thermal particles, the transfer function has been computed precisely; the characteristic wavenumber where the square of the transfer function falls to half the CDM value is about $`k_{1/2,stream}k_{X,comov}/5.5`$. The simple streaming case only works for high phase densities $`m_X1`$keV, that is, comparable to that observed in dwarf spheroidals. For example, to produce an acceptable number of galaxies at a dwarf galaxy scale without invoking disruption, Press-Schechter theory implies a spectral cutoff at about $`k=3h_{70}\mathrm{Mpc}^1`$, requiring a thermal relic mass of about 1100 eV. Hydrodynamic simulations show that the same cutoff scale preserves the large scale success of CDM and probably improves the CDM situation on galaxy scales in ways mentioned previously. Although the typical uncertainty on the phenomenologically best filtering scale is at least a factor of two, it is clear that the smallest phase density compatible with standard streaming filtering is too large to have a direct impact on the core problem in dwarf disk galaxies.
On the other hand the discrepancy is only a factor of a few in mass, less than an order of magnitude in linear damping scale. We have already mentioned two modifications which could reconcile these scales. It could be that warm models turn out to be sometimes more effective at producing smooth cores than we have guessed from the minimal phase-packing constraint, due to more efficient dynamical heating than CDM; this would produce a nonlinear amplifier of the primordial velocities, probably with a large variation depending on dynamical history (an especially good option if cores turn out to be common in galaxy clusters.) Another possibility is that the primordial velocities are introduced relatively late (nonrelativistically) by particle decay.
Still another possibility is a different relationship of $`k_{1/2}`$ and $`k_X`$ from the standard collisionless streaming behavior. For example, if the particles are self-interacting, then the free streaming is suppressed and the relevant scale is the standard Jeans scale dividing growing behavior from acoustic oscillations, $`4\pi G\rho _{total}k_J^2c_S^2=0`$. This comes out to $`k_J=\sqrt{3}H/c_S=\sqrt{27/5}k_X`$, 13 times shorter than $`k_{1/2,stream}`$ at a fixed phase density. (An intuitive view of the this numerical factor is that during the long period when $`k_X`$ is flat, streaming particles continue to move and damp on larger scales, whereas the comoving Jeans scale just remains fixed, sharply dividing oscillating from growing behavior.) The acoustic case is similar to the behavior of fluctuations in high-density, baryon-dominated models, which have a sharp cutoff at the Jeans scale. We conclude that some particle self-interactions may be desirable to reconcile the scale of the transfer function of primordial perturbations with the phase packing effect on disk cores.
## V Collisional Dark Matter
We now turn to the case where the dark matter particles are not collisionless, but scatter off of each other via a new intermediate force. Self-interactions of dark matter have been motivated from both an astrophysical and a particle physics point of view. Our goal here is again to relate the properties of the new particles to the potentially observable properties of dark matter halos. In addition to the single parameter $`Q`$ considered for the collisionless case, we can use halo properties to constrain fundamental parameters of the particles— the masses of the dark matter particles and intermediate bosons carrying the interactions, as well as a coupling constant.
Such self-interactions lead to modifications in several of the previous arguments. As we have seen, self-interactions can have observable effects via the transfer function even if they are negligible today. Stronger self-interactions also affect the structure and stability of halos; collisional matter has a fluid character leading to equilibrium states of self-gravitating halos much like those of stars. These systems are quite different from collisionless systems. Although entropy must increase outwards for stability against convection (which naturally happens due to shocks in the hierarchy), it cannot increase too rapidly and remain hydrostatically stable; in particular, stable solutions have a minimum nonneglible temperature gradient, and the isothermal case is no longer a stable static solution as it is for collisionless matter. Since collisional matter conducts heat between fluid elements, these solutions are all unstable on some timescale.
### A Particles and Interactions
We now apply several simple physical arguments to constrain properties of the dark matter candidate and its interactions. Some of these have been considered previously. The most important constraints are summarized in figure 2.
Suppose that the dark matter $`X`$ particles with mass $`m_X`$, which may be either fermions or bosons, interact via massive bosons $`Y`$ whose mass $`m_Y`$ determines the range of the interactions, and a coupling constant $`e`$. These may be considered analogous to strong interaction scatterings where we regard pions as Yukawa scalar intermediates, or electroweak interactions with $`W,Z`$ as vector intermediates. The interactions must be elastic scatterings to avoid a net energy loss, although “dissipative” three-body encounters are permitted as long as the energy does not leave the $`XY`$ subsystem nor travel far in space. For most purposes even the sign of the interaction does not matter— it may be attractive or repulsive, as happens with vectors and like charges. The $`Y`$ particles at tree level interact only with $`X`$, although the $`X`$ may (as is usual with dark matter candidates) be allowed some much weaker interactions with ordinary matter. In this model the collision cross section for strong scattering is about
$$\sigma m_Y^2\mathrm{min}[e^4\left(\frac{m_Y}{m_Xv^2}\right)^2,e^4\left(\frac{m_X}{m_Y}\right)^2,1]$$
(24)
where the first case is coupling-limited (and depends on the particle velocity and coupling strength, like electromagnetic scattering of electrons), the second case holds for $`m_Y>m_X`$ (like neutrino neutral-current interactions) and the third is the range-limited, strong interaction limit (like neutron scattering).
There are several simple constraints on the particle masses. If the dark matter is collisional, the rate of net annihilations of $`X`$ must be highly suppressed compared to the scattering rate, or the mass of the halo would quickly radiate away as $`Y`$ particles. Either there is a primordial asymmetry (so the number of $`\overline{X}`$ is negligible), or
$$m_Y>2m_X,$$
(25)
suppressing what would otherwise be a rapid channel for $`X`$ to annihilate and radiate $`Y`$. (Recall that in this model, there is no direct route to annihilate into anything else). In any case the $`Y`$ must not be too light or the typical inelastic collisions will radiate them; for particles with relative velocities $`v10^3`$ typical of dark matter in galaxies, we must have
$$m_Y>m_Xv^210^6m_X,$$
(26)
so that the energy of collisions is typically insufficient to create a real $`Y`$. In addition, if attractive, the range of the interactions must be less than the “Bohr radius” for these interactions, requiring
$$m_Y>e^2m_X,$$
(27)
in order not to form bound “atoms”. The close analogy with $`Y`$ is the pion, which is just light enough to allow a bound state of deuterium. Bound states would be a disaster since they would behave like nuclear reactions in stars. Such states would add an internal source of energy in the halos, creating winds or other energy flows which would unbind large amounts of matter. All of these constraints eliminate the upper left region of figure 2, with details depending on the coupling strength and halo velocity.
### B Parameters for Collisional Behavior
The properties of interacting particles define a characteristic column density, $`m_X/\sigma `$; a slab of $`X`$ at this column is one mean free path thick. This is the quantity that specifies the degree of collisional or collisionless behavior of a system. In order to connect the halo astrophysics with dark matter properties we convert from units with $`\mathrm{}=c=1`$:
$$(1\mathrm{GeV})^3=4.6\times 10^3\mathrm{g}\mathrm{cm}^2=2.2\times 10^7\mathrm{M}_{}\mathrm{pc}^2$$
(28)
For comparison, the average mass column density within radius $`r_{kpc}\mathrm{kpc}`$ for a halo with a circular velocity $`v_{30}\times 30\mathrm{k}\mathrm{m}\mathrm{sec}^1`$ is
$$\mathrm{\Sigma }_h=\frac{v^2}{\pi Gr}=0.014v_{30}^2r_{kpc}^1\mathrm{g}\mathrm{cm}^2=(15\mathrm{M}\mathrm{e}\mathrm{V})^3v_{30}^2r_{kpc}^1.$$
(29)
A halo therefore enters the strongly-collisional regime— qualitatively different from classical CDM— if
$$m_Y^4e^4(15\mathrm{M}\mathrm{e}\mathrm{V})^3v_{30}^2r_{kpc}<m_X<\mathrm{min}[(15\mathrm{M}\mathrm{e}\mathrm{V})^3v_{30}^2r_{kpc}^1m_Y^2,15\mathrm{M}\mathrm{e}\mathrm{V}(e/v)^{4/3}(v_{30}^2r_{kpc}^1)^{1/3}].$$
(30)
This criterion is shown in figure 2 as the right boundary of the “unstable cores” region; indeed this marginally-collisional case maximizes the rate of thermal conduction instability, as discussed below.
We also compute the criterion for non-streaming behavior in the early universe— the amount of self-interaction needed to affect the transfer function as discussed above. It is significantly less than that required for collisional behavior today:
$$\frac{\sigma }{m_X}H(t_{eq})/n_X(t_{eq})v_X(t_eq)m_X=\mathrm{\Sigma }_0^1\mathrm{\Omega }_{rel}^{5/2}\mathrm{\Omega }_X^1v_X(t_0),$$
(31)
where $`eq`$ refers to the epoch of equal densities in dark matter and relativistic species, and
$$\mathrm{\Sigma }_0c\rho _{crit}/H_0=0.1213h_{70}^1\mathrm{g}\mathrm{cm}^2$$
(32)
is the characteristic cosmic column density today. Using the units conversion above we have
$$\frac{\sigma }{m_X}(600\mathrm{M}\mathrm{e}\mathrm{V})^3(v_{X0}/1\mathrm{k}\mathrm{m}\mathrm{s}^1)^1(\mathrm{\Omega }_X/0.3)^2h_{70}^4,$$
(33)
corresponding to a mass column for one expected scattering of $`2\times 10^4\mathrm{g}\mathrm{cm}^2`$. Particles scattering off of each other more strongly than this no longer have streaming behavior at high redshift but support acoustic oscillations, much like baryons but with only their own pressure (that is, without the interaction with radiation pressure and without decoupling from it). We should bear in mind that a somewhat larger cross section is needed to avoid diffusive (“Silk”) damping, but even at this level of interaction the scale of damping is is significantly reduced from the streaming case. This criterion is shown in figure 2 as the right boundary of the “Jeans” region (although some acoustic behavior before $`t_{eq}`$ occurs even to the right of this).
### C Polytropes
The equilibrium configurations of collisional dark matter correspond to those of classical self-gravitating fluids. The simplest cases to consider and general enough for our level of precision are classical polytrope solutions— stable configurations of a classical, self-gravitating, ideal gas with a polytropic equation of state. In the absence of shocks or conduction, the pressure and density of a fluid element obey an equation of state $`p=K_1\rho ^{\gamma _1}`$. For an adiabatic, classical, nonrelativistic, monatomic gas, or for nonrelativistic degenerate particles, the adiabatic index $`\gamma _1=5/3`$ and different values of $`K_1`$ correspond to different entropy. If the entropy varies radially as a power-law, equilibrium self-gravitating configurations are given by classical Lane-Emden polytrope solutions. The radial variation of pressure and density obey $`p(r)=K_2\rho ^{\gamma _2}(r)`$; the second index $`\gamma _2`$ tracks the radial variation between different fluid elements in some particular configuration (that is, including variations in entropy). For gas on the same adiabat everywhere, $`\gamma _1=\gamma _2`$; for the case of nonrelativistic degenerate or adiabatic matter, $`\gamma _1=\gamma _2=5/3`$ applies and is a good model of degenerate dwarfs. If the entropy is increasing with radius, as would be expected if assembled in a cosmological hierarchy, then $`\gamma _2<\gamma _1`$, conferring stability against convection.
The character of the solutions is well known. As long as $`\gamma _2>6/5`$ the halo structure is like a star, with a flat-density core in the center, falling off in the outer parts to vanishing density at a boundary. If it is rotating, the structure is similar but rotationally flattened. These solutions describe approximately the structure of stars, especially degenerate dwarfs, and halos of highly collisional dark matter.It is worth commenting on some differences and similarities with collisionless halos with finite phase density material. The polytrope solutions are for collisional matter with an isotropic pressure and local balance of pressure gradient and gravity. Collisionless particles can fill phase space more sparsely, but this just means that at a given mass density they must have a larger maximum velocity; the collisional solution saturates the phase density limit and has the largest mass density for a given coarse-grained phase density. In this sense, once one is solving the cusp problem with finite phase density, nothing further is gained by making the particles collisional. Collisionless particles allow anisotropy in the momentum distribution function, and therefore a wider range of ellipsoidal figures, but cannot pack into tighter cores. For the same reason, the inner phase-density-limited core is expected to be close to spherical except for rotational support, whether the particles are collisional or not. The phase space is fully occupied and therefore the velocity distribution is close to isotropic wherever the local entropy approaches the primordial value. If $`\gamma _2<6/5`$ (and in particular for the isothermal case $`\gamma _2=1`$) there is a dynamical instability and no stable solution; the system runs away on a gravitational timescale, with the center collapsing and the outer layers blowing off.
### D Giant Dwarfs
At zero entropy the equilibrium configuration is the exactly soluble $`\gamma =5/3`$ polytrope, which we adopt as an illustrative example. That is, we model a dwarf galaxy core as a degenerate dwarf star, the only difference being a particle mass much smaller than a proton allowing a halo mass much bigger than a star. For total mass $`M`$ and radius $`R`$, the Lane-Emden solution gives a central pressure $`p_c=0.770GM^2/R^4`$ and a central density $`\rho _c=5.99\overline{\rho }=1.43M/R^3`$. Using the above relation for the equation of state we obtain the standard degenerate dwarf solution, which has
$$R=4.5m_X^{8/3}M^{1/3}m_{Planck}^2=0.98\mathrm{kpc}\left(\frac{m_X}{100\mathrm{e}\mathrm{V}}\right)^{8/3}\left(\frac{M}{10^{10}M_{}}\right)^{1/3},$$
(34)
where $`m_{Planck}=\sqrt{\mathrm{}c/G}`$ and $`M_{}=9.48\times 10^{37}m_{Planck}`$. This “giant dwarf” configuration is stable even at zero temperature up to the Chandrasekhar limit for $`X`$ particles.<sup>\**</sup><sup>\**</sup>\**Defined analogously to the Chandrasekhar limit for standard dwarfs (with $`Z=A`$ because there is just one kind of particle providing both mass and pressure, similar to a neutron star),
$$M_{CX}=3.15\frac{m_{Planck}^3}{m_X^2}=4.95\times 10^{14}M_{}(m_X/100\mathrm{e}\mathrm{V})^2.$$
(35)
Since the mass is not directly observable, it is more useful to consider the velocity of a circular orbit at the surface, $`v_c=(GM/R)^{1/2}`$. We then obtain the relation for a degenerate system,
$$m_X=4.5^{3/8}v_c^{1/4}(r_cm_{Planck})^{1/2}m_{Planck},$$
(36)
or in more conventional astrophysical units,
$$m_X=87eV(v_c/30\mathrm{k}\mathrm{m}/\mathrm{s})^{1/4}(r_c/1kpc)^{1/2}.$$
(37)
Note that as in the collisionless case, no cosmological assumptions or parameters have entered into this expression.
For any adiabatic nonrelativistic matter the solution is similar. The absolute scale of the giant dwarf, determined by $`K_2`$, is fixed by the phase density $`Q`$. In general there is a range of entropy but once again the the lowest-entropy material (which is densest at a given pressure) sinks to the center of a halo and forms an approximately adiabatic core. The rest of the halo forms a thermally-supported atmosphere above it. Once again cores are the places to look for signs of a primordial ceiling to phase density. However, as we see below the behavior changes if conduction or radiation are not negligible. As we know, a thermally supported star which conducts heat and has no nuclear or other source of energy is unstable.
### E Heat Conduction Time and Halo Stability
If the collision rates are not very high we must consider heat and momentum transport between fluid elements by particle diffusion. The most serious consideration for radial stability is the transport of heat. In all stable thermally supported solutions the dense inner parts are hotter; if conduction is allowed, heat is transported outwards. The entropy of the central material decreases, the interior is compressed to higher density and the outer layers spread to infinity, a manifestation of the gravothermal catastrophe. With conduction the inner gas falls in and the outer gas drifts out on a diffusion timescale, attempting to approach a singular isothermal sphere.
Consider the scenario where the dark matter cross section is small enough to remain essentially noninteracting on large scales, preserving the successes of CDM structure formation simulations, but large enough to become collisional in the dense central regions of galaxies. Although this scenario was introduced to help solve the cusp problem, we will see that the conductive instability makes matters worse. If stable cores are to last for a Hubble time, the dark matter halos must either be effectively collisionless (standard dark matter), or very strongly interacting, so that the inevitable conduction is slow (or made of degenerate fermions so there is no temperature gradient.)
Elementary kinetic theory yields an estimate for the the conduction of heat by particle diffusion; the ratio of energy flux to temperature gradient is the classical conduction coefficient $`\kappa \sigma ^1\sqrt{T/m}`$. Assuming a halo in approximate virial equilibrium and profile $`v(r)`$, this yields a timescale for heat conduction,
$$t_{cond}\frac{v\sigma }{2Gm_X}\frac{d\mathrm{log}r}{d\mathrm{log}v}$$
(38)
where $`v`$ is the typical particle velocity (which is about the virial velocity of the halo independent of the mass of the particles $`m_X`$). The first factor is essentially the time it takes a particle to random walk a distance $`r`$, $`t_{diffuse}r^2n\sigma /v`$. The last factor characterizes the temperature and entropy gradient; dynamical stability prevents it from being very large, and in most of the matter it typically takes a value not much larger than unity.<sup>††</sup><sup>††</sup>††The conductive destabilization probably happens faster than Spergel and Steinhardt estimated. They used the Spitzer formula describing core collapse in globular clusters, which takes about 300 times longer than the two-particle relaxation time. However, the large factor arises because in the globular cluster case the relaxation is entirely gravitational and is dominated by very long-range interactions with distant stars. In the present situation the interactions are strong and short-range, leading to significant exchange of both energy and momemtum in each scattering. The transport of heat takes place on the same timescale as the diffusion of particles, with numerical factors of the order of unity as in standard solutions of the Boltzmann equation for gases.
A halo with conduction therefore forms a kind of cooling flow, with the core collapsing and the envelope expanding. If it is hydrostatically quasi-stable (that is, if the core collapse is slow and regulated by the particle diffusion), we can use the Lane-Emden solutions to set bounds on the numerical factor $`d\mathrm{log}r/d\mathrm{log}v`$ governing the instability. The equation of state tells us that $`v\rho ^{1/2n_2}`$ where $`n_2=(\gamma _21)^1`$. The largest value of $`n_2`$ which corresponds to a quasi-stable solution is $`n_2=5`$. The density profile is steeper than isothermal ($`n_2=\mathrm{}`$), $`\rho r^2`$; therefore $`|d\mathrm{log}r/d\mathrm{log}v|n_25`$. In the rough estimates here we set these factors to unity.<sup>‡‡</sup><sup>‡‡</sup>‡‡Another interesting limit is that of small but nonzero self-interactions. The halo is essentially collisionless, but occasional scatterings still take place. The collisionless isothermal sphere, singular or not, is then an approximate solution, but still subject to a slow secular instability from heat conduction. It is also possible to set up situations where halos are evaporated by a hot external environment, heated from outside by collapse of the cosmic web.
Conduction can be suppressed if the scattering is very frequent. For nondegenerate $`X`$, stable cores require that the conduction time exceeds the Hubble time $`H_0^1`$. For stability over a Hubble time, the column density of a halo with velocity $`v`$ must exceed $`m_X/\sigma =Hv/G`$; therefore the particles must satisfy
$$\frac{m_X}{\sigma }1.0\times 10^4\mathrm{g}\mathrm{cm}^2h_{70}v_{stable,30}.$$
(39)
where $`v_{stable,30}\times 30\mathrm{k}\mathrm{m}\mathrm{s}^1`$ denotes the velocity in the lowest-velocity stable halo. Perhaps surprisingly, the mass and radius of the halo do not enter explicitly.
This condition constrains the particles to be highly interactive. Galaxy halos have slow conduction compared to $`H`$ only above a critical velocity dispersion $`v_{crit}(G/H)(m_X/\sigma )`$. Halos below this threshold should have collapsed cores, and above the threshold the core radius/mass relation is determined as before by the giant dwarf sequence for the the particle’s phase density. The existence of stable bound 30 km/s halos of highly-collisional dark matter requires
$$\frac{m_X}{\sigma }(2.8\mathrm{MeV})^3h_{70}v_{stable,30},$$
(40)
shown in figure 2 as the right boundary of the “fluid” region.
The “thickness” of a halo with velocity $`v_{30}\times 30\mathrm{k}\mathrm{m}\mathrm{s}^1`$, in units of particle pathlengths, is
$$\frac{\mathrm{\Sigma }_h\sigma }{m_X}10^2v_{30}^2r_{kpc}^1h_{70}^1v_{stable,30}^1$$
(41)
so it is clear that all dark-matter-dominated structures, from small galaxies to galaxy clusters ($`v_{30}130`$, $`r_{kpc}0.11000`$), are highly collisional and their dark matter behaves as a fluid. Even for very diffuse matter at the mean cosmic density ($`\mathrm{\Omega }_X=0.3`$), the particle mean free path is at most $`12v_{core,30}h_{70}^1`$Mpc, about the same as the scale of nonlinear clustering, so all bound dark matter structures act like fluids.
Are other data consistent with the idea that essentially all dark matter acts like a fluid? This option has been considered previously and while it is perhaps not definitively ruled out, it is not phenomenogically compelling. Serious problems arise for example from satellite galaxies which are thought to have had several orbits without stopping and sinking as they would in a fluid, or from galaxies in clusters, at least some of which appear (from lens reconstruction mass maps) to have retained some of their dark matter halos. An intriguing possibility is that a small collision rate might contribute to enough instability to feed the formation of black holes. However the rate of the instability is greatest in the lowest mass, lowest density galaxy cores, a trend not conspicuous in the demography of central black holes of galaxies.<sup>\**</sup><sup>\**</sup>\**We have to take note of another possibility: perhaps the dwarf spheroidals, which have the lowest velocity dispersions of all galaxies and are also the densest, have already collapsed by heat conduction. In this way we could use phase packing to give the cores of the dwarf disk galaxies and still explain why the dwarf spheroidals have such a large phase density. Note that this scheme also gives the right filtering scale since the particles are collisional at early times. The dwarf spheroidals need not of course collapse all the way to black holes, but they may well have singular dark matter profiles.
We conclude that dark matter self-interactions are likely to be negligible in galaxy halos, and that this places significant constraints on the particles. Figure 2 summarizes the constraints on the parameters $`m_X,m_Y`$ of this interacting-particle model from the various constraints considered here.
## VI Conclusions
We have found that some halos might preserve in their inner structure observable clues to new dark matter physics, and that indeed some current observations already hint that the dark matter might be warm rather than cold. We conclude with a summary:
1. Halo cores can be created by a “phase-packing limit” depending on finite initial phase density. They may provide a direct probe of primordial velocity dispersion in dissipationless dark matter.
2. For relativistically-decoupled thermal relics, the phase density depends on the particle mass and spin but not on cosmological parameters.
3. Rotation curves in a few dwarf disk galaxies indicate cores with a phase density corresponding to that of a 200 eV thermal relic or an rms velocity of about 0.4 km/sec at the current cosmic mean density. Velocity dispersions in dwarf spheroidal galaxies indicate a higher phase density, corresponding to a thermal relic mass of about 1 keV. At most one of these populations can be tracing the primordial phase density.
4. Thermal relics in this mass range can match the mean cosmic density with a plausible superweak decoupling from Standard Model particles before the QCD epoch.
5. Other very different particles are consistent with the halo data, provided they have the about the same mean density and phase density. Examples include WIMPs from particle decay and axions from defect decay.
6. Cores due to phase packing limited by primoridial $`Q_0`$ predict a universal relation between core radius and halo velocity dispersion. The relation is not found in a straightforward interpretation of the data.
7. Primordial velocity dispersion also suppresses halo substructure (and solves some other difficulties with CDM) by filtering primordial adiabatic perturbations. Estimates based on luminosity functions prefer filtering on a scale of about $`k3\mathrm{M}\mathrm{p}\mathrm{c}^1`$; for collisionless particles, this scale corresponds to a filter caused by streaming of about a 1keV thermal relic.
8. Weak self-interactions change from streaming to acoustic behavior, reducing the damping scale and sharpening the filter.
9. Stronger self-interactions destabilize halos by thermal conduction, making the cusp problem worse (unless they are very strong— too strong for satellite-galaxy kinematics— or particles are degenerate, eliminating the central temperature gradient).
10. A simulation which samples a warm distribution function reasonably well is strongly motivated, to determine whether primordial $`Q`$ is preserved in the centers of halos, or whether nonlinear effects can amplify dynamical heating in such models to explain cores on all scales.
## ACKNOWLEDGMENTS
We are grateful for useful discussions of these issues with F. van den Bosch, A. Dolgov, G. Fuller, B. Moore, J. Navarro, T. Quinn, J. Stadel, J. Wadsley, and S. White. JD gratefully acknowledges the hospitality of the Institute for Theoretical Physics at UC Santa Barbara, which is supported in part by the National Science Foundation under Grant No. PHY94-07194. JD was partially supported by NSF Grant AST-990862. CJH thanks the Max-Planck-Institute für Astrophysik, the Isaac Newton Institute for Mathematical Sciences and the Ettore Majorana Centre for Scientific Culture for hospitality. His work was supported at the University of Washington by NSF and NASA, and at the Max-Planck-Institute für Astrophysik by a Humboldt Research Award.
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# \title — ?????
(January 21, 2000) (January 21, 2000)
## Abstract
We use the Popov-Fedotov representation of spin operators to construct an effective action for a Kondo lattice model with quenched disorder at finite temperatures. We study the competition between the Kondo effect and frozen spin order in Ising-like spin glass. We present the derivation of new mean-field equations for the spin-glass order parameter and analyze the effects of screening of localized spins by conduction electrons on the spin-glass phase transition.
PACS: 71.27.+a, 75.20.Hr
One of the most interesting questions of physics of heavy-fermion (HF) compounds is the competition between Kondo screening of localized spins by conduction electrons (CE) and ordering of these spins due to Ruderman-Kittel-Kasuya-Yosida (RKKY) interaction (see, e.g. ). The screening is attributed to the Kondo effect \- the resonance scattering of an electron on a magnetic atom with simultaneous change of the spin projection. In dilute alloys such a scattering results in the sharp resonance at Fermi level with characteristic energy width $`ϵT_Kϵ_F\mathrm{exp}(\alpha ^1)`$, where $`T_K`$ is Kondo temperature, $`J`$ is a coupling constant, $`\rho `$ is a density of states of CE on the Fermi level and $`\alpha =\rho J`$. As it was recently discussed (see, e.g. -), such a competition can be responsible for the Non-Fermi-liquid (NFL) behaviour observed in some heavy fermion compounds. Most of such a materials share two characteristics: proximity to magnetic region of appropriate phase diagram (usually temperature vs. pressure or chemical composition), and disorder due to chemical substitution. In many respects the concentrated Kondo systems, namely the lattice of magnetic atoms interacting with CE ”bath” (Kondo lattice (KL)) show striking similarities with dilute Kondo systems. The Kondo temperature in these systems is a characteristic crossover temperature at which spins transform their local properties to some itinerant Fermi-liquid behaviour determining low temperature regime of HF compounds. NFL behaviour in HF system is mainly attributed then to reducing the Kondo temperature possibly even suppressing it to zero. In turn, magnetic or spin glass (SG) transition can also be suppressed due to interplay between Kondo scattering and spin-spin interaction. Thus, such an interplay can result in quantum phase transition (QPT) when both Kondo and magnetic temperatures are equal to zero at some finite doping. The role of chemical substitution in this case is to ”tune” the Fermi level of metallic system providing sharp Kondo resonance.
The problem of competition between RKKY and Kondo interaction in clean system was studied for the first time by Doniach in the ”Kondo necklace” model. The transition typically takes place between a paramagnetic metal and magnetic (usually AFM) metal. In this case there are two possibilities: the compound will have long range magnetic order when the RKKY interaction is sufficiently large compared with the Kondo interaction, or the compound will be paramagnetic due to the quenching of magnetic moments of the rare earth atoms and the ground state has the features of Kondo-singlet state. Nevertheless, in the region $`T_{RKKY}^MT_K`$ the competition between magnetic and Kondo interaction results in dramatically change of ”naive” Doniach diagram (see ). Namely, both Kondo and magnetic temperatures are strongly suppressed and spin-liquid state (e.g of Resonance Valence Bond type ) occurs.
The goal of this letter is to present some results concerning the competition between Kondo effect and Ising-like SG transition which is in many aspects similar to the magnetic instability. We study mechanisms of suppressing the SG transition and effects of screening in disordered environment. In this paper we consider the high temperature regime of KL model. We leave aside the issue of the ground state properties and especially the question whether NFL behaviour is a generic feature of vicinity to QPT for future publication.
The Hamiltonian of KL model with additional quenched randomness of exchange interaction between localized spins is given by
$$H_{KL}=\underset{k\sigma }{}\epsilon _kc_{k\sigma }^{}c_{k\sigma }+J\underset{i}{}\left(\stackrel{}{s}_i\stackrel{}{S}_i+\frac{1}{4}n_iN_i\right)\underset{ij}{}I_{ij}\left(S_i^zS_j^z+\lambda S_i^+S_j^{}\right)$$
(1)
The system under consideration is a periodic lattice of magnetic atoms modeled by $`f`$ \- orbitals interacting with metallic background spin density operator $`\stackrel{}{s}_i=\frac{1}{2}c_{i\alpha }^{}\stackrel{}{\sigma }_{\alpha \alpha ^{}}c_{i\alpha ^{}}`$. The first term in the Hamiltonian (1) describes kinetic energy of CE, the second stands for Kondo coupling ($`J>0`$). We denote $`n_i=_\sigma c_{i,\sigma }^{}c_{i,\sigma }`$ as the CE density operator. The identity $`N_i=1`$ describes the half-filled f-electron shell. Quenched independent random variables $`I_{ij}`$ with distribution $`P(I_{ij})\mathrm{exp}(I_{ij}^2N/(2I^2))`$ stand for direct spin-spin interaction . We assume that this random interaction is of RKKY origin <sup>1</sup><sup>1</sup>1It has been pointed out in , that the presence of nonmagnetic impurities makes RKKY interaction a random interaction even in the case of regular arrangement of magnetic moments., namely, for $`d`$-dimensional system $`I\alpha ^2ϵ_Fl^d`$, $`l`$ is a lattice constant in magnetic sublattice. The magnetic effects can also be included in our approach by introducing nonzero standard deviation $`\mathrm{\Delta }I=\overline{I}_{RKKY}`$ into the distribution $`P(I_{ij})`$, which, in turn, can result in the additional competition between SG and AFM (or, rarely FM) states. For simplicity we neglect these effects in present letter concentrating on the interplay between Kondo interaction and effects of bond disorder. Since indirect RKKY interaction through CE is mostly determined by ”fast” electrons with characteristic energies $`ϵϵ_FT_K`$ we neglect also the Kondo renormalizations of RKKY exchange.
As it well-known for a long time, the spin $`S=1/2`$ matrices can be exactly replaced by bilinear combination of Fermi operators $`S_i^z=\frac{1}{2}(f_i^{}f_if_i^{}f_i)`$, $`S_i^+=f_i^{}f_i`$, $`S_i^{}=f_i^{}f_i.`$ Nevertheless, most of fermionic representations of spin are not free of constraint problem. For this reason, the dimensionality of space in which these operators act is always greater than the dimensionality of the spin matrices. Elimination of unphysical states is a serious problem which makes the diagrammatic techniques quite complicated. Moreover, in most cases, the analytical continuation of Feynman diagrams becomes extremely uneasy. To avoid the main part of difficulties related to constraint, the new representation for spin operators was proposed in well-forgotten paper of Popov and Fedotov . In this representation the partition function of the problem containing spin operators ($`H_S`$) then can be easily expressed in terms of new fermions with imaginary chemical potential ($`H_S^f`$): $`Z_S=\mathrm{𝚃𝚛}e^{\beta H_S}=i^N\mathrm{𝚃𝚛}e^{\beta (H_S^f+i\pi N_f/(2\beta ))}`$, $`N_f=_{i\sigma }f_{i\sigma }^{}f_{i\sigma }`$ and $`\beta =1/T`$. As a result, there is no constraint, the unphysical states are eliminated and standard Matsubara-Abrikosov-Gor’kov diagrammatic technique is obtained ,,.
We sketch our derivation of the effective action and of resulting mean field (MF) equations for KL model in order to explicit the approximations made and the physics underlying these approximations. To construct the path integral representation for the partition function, the new Grassmann variables $`c_{i\sigma }^{}\overline{\mathrm{\Psi }}_{i\sigma },c_{i\sigma }\mathrm{\Psi }_{i\sigma }`$ for CE with chemical potential $`\mu `$ and $`f_{i\alpha }^{}\overline{a}_{i\alpha },f_{i\alpha }a_{i\alpha }`$ for Popov-Fedotov (PF) spin operators ($`S=1/2`$) are introduced. The Euclidean action for the KL model is given by
$$𝒜=_0^\beta 𝑑\tau \left(\underset{i\alpha }{}\left[\overline{\mathrm{\Psi }}_{i\alpha }(\tau )(_\tau +\mu )\mathrm{\Psi }_{i\alpha }(\tau )+\overline{a}_{i\alpha }(\tau )(_\tau i\pi T/2)a_{i\alpha }(\tau )\right]H_{int}(\tau )\right).$$
(2)
where the generalized Grassmann fields satisfy the following boundary conditions: $`\mathrm{\Psi }_{i\alpha }(\beta )=\mathrm{\Psi }_{i\alpha }(0)`$, $`\overline{\mathrm{\Psi }}_{i\alpha }(\beta )=\overline{\mathrm{\Psi }}_{i\alpha }(0)`$, $`a_{i\alpha }(\beta )=ia_{i\alpha }(0)`$, $`\overline{a}_{i\alpha }(\beta )=i\overline{a}_{i\alpha }(0)`$.
In this paper we consider $`\lambda =0`$ which corresponds to Sherrington- Kirkpatrik (SK) spin-glass model. Such an anisotropy of RKKY interaction can be associated e.g. with lattice geometry. In the case of Ising-like model the dynamical fluctuations in spin subsystem appear only due to interaction with conduction electrons and in high temperature regime $`TT_{SG}`$ can be neglected. To study the influence of Kondo-scattering on SG transition temperature $`T_{SG}`$ we use standard replica trick $`\mathrm{\Psi }_i(\tau )\upsilon _i^a(\tau ),a_i(\tau )\phi _i^a(\tau ),a=1,\mathrm{}n.`$ Then, the free energy of the model can be calculated (see, e.g. ) by taking the formal limit $`n0`$ in
$$<Z^n>_{av}=𝑑I_{ij}P(I_{ij})D[\phi _{i,\sigma }^a\upsilon _{i,\sigma }^a]\mathrm{exp}\left(𝒜_0[\upsilon ^a,\phi ^a]_0^\beta 𝑑\tau H_{int}(\tau )\right)$$
(3)
where $`𝒜_0`$ is corresponding to noninteracting fermions.
As we already mentioned, for considering the competition between Kondo scattering and trend of disorder we assume that the magnetic temperature $`T_{RKKY}^MT^{}`$, where $`T^{}`$ stands for characteristic temperature corresponding to the Kondo temperature in a lattice. This assumption allows one to decouple the Kondo interaction term $`H_i^K=\frac{J}{2}\overline{\upsilon }_{i,\sigma }^a\phi _{i,\sigma }^a\overline{\phi }_{i,\sigma ^{}}^a\upsilon _{i,\sigma ^{}}^a`$ in each site by the replica-dependent Hubbard-Stratonovich field $`\psi _i^a`$ . Performing the average over random potential in Eq.(3) results in
$$\begin{array}{c}<Z^n>_{av}=D[\upsilon ^a,\phi ^a,\psi ^a]\mathrm{exp}(𝒜_0+\frac{I^2}{4N}\mathrm{𝚃𝚛}[X^2]+\hfill \\ \hfill +_0^\beta d\tau \underset{i,a,\sigma }{}\{\psi _i^a\overline{\upsilon }_{i\sigma }^a\phi _{i\sigma }^a+\psi _i^a\overline{\phi }_{i\sigma }^a\upsilon _{i\sigma }^a\frac{2}{J}|\psi _i^a|^2\})\end{array}$$
(4)
with $`X^{ab}(\tau ,\tau ^{})=_i_{\sigma ,\sigma ^{}}\overline{\phi }_{i,\sigma }^a(\tau )\sigma \phi _{i,\sigma }^a(\tau )\overline{\phi }_{i,\sigma ^{}}^b(\tau ^{})\sigma ^{}\phi _{i,\sigma ^{}}^b(\tau ^{})`$. The next step is to perform the Gaussian integration over the replica-dependent Grassmann field $`\upsilon ^a`$ describing CE and to decouple the eight-fermion term $`\mathrm{𝚃𝚛}[X^2]`$ with the help of $`Q`$ matrices (see details in ). As a result, the partition function is given by:
$$\begin{array}{c}<Z^n>_{av}=D[Q]\mathrm{exp}(\frac{1}{4}(\beta I)^2N\mathrm{𝚃𝚛}[Q^2]+\hfill \\ \hfill +\underset{i}{}\mathrm{ln}\{D[\phi ^a,\psi ^a]\mathrm{exp}[\underset{a}{}\underset{\{\omega \}}{}\overline{\phi }_{i\sigma }^a𝒢_a^1\phi _{i\sigma }^a+\frac{1}{2}(\beta I)^2\mathrm{𝚃𝚛}[QX]]\})\end{array}$$
(5)
where $`𝒢^1`$ is inverse Green function for PF fermions depending on Matsubara frequences $`\omega _n=2\pi T(n+1/4)`$ (see details in )
$$𝒢_a^1=i\omega _{n_1}\delta _{\omega _{n_1},\omega _{n_2}}T\underset{ϵ}{}\psi _i^a(ϵ_l+\omega _{n_1})G_0(i_i,ϵ_l)\psi _i^a(ϵ_l+\omega _{n_2})$$
(6)
and $`G_0(i,ϵ_l)=(iϵ_l\epsilon (i)+\mu )^1`$ stands for CE Greens function, $`ϵ_l=2\pi T(l+1/2)`$.
We are still left with a term of fourth order residing in $`\mathrm{𝚃𝚛}[QX]`$ and can not evaluate the Grassmann integral directly. Consequently, the second decoupling is needed. To perform it, we stress that we do not intend to deal with dynamical behaviour here confining ourselves by high temperature regime in the vicinity of SG transition such that the lowest Matsubara frequency is sufficient. Assuming this and recalling that the spatial fluctuations are suppressed by the choice of infinite range interaction , one can consider $`Q`$ as a constant saddle point matrix under condition $`Q=Q^T`$. The elements of this matrix will later be determined self-consistently from the saddle point condition. Assuming that the elements of $`Q`$ are $`Q_{SP}^{aa}=\stackrel{~}{q}`$ and $`Q_{SP}^{ab}=q`$ one can decouple the $`\mathrm{𝚃𝚛}[QX]`$ term by introducing replica-independent $`z`$ and replica-dependent $`y^a`$ fields and map KL problem with disorder onto effective one-site interacting spin system coupled with external local replica-dependent magnetic field:
$$\begin{array}{c}<Z^n>_{av}=\mathrm{exp}(\frac{1}{4}(\beta I)^2N(n\stackrel{~}{q}^2+n(n1)q^2)+\hfill \\ \hfill +\underset{i}{}\mathrm{ln}[D[\phi ^a,\psi ^a]_z^G_{y^a}^G\mathrm{exp}(𝒜[\phi ^a,\psi ^a,y^a,z])])\end{array}$$
(7)
where $`_z^Gf(z)`$ denotes $`_{\mathrm{}}^{\mathrm{}}𝑑z/\sqrt{2\pi }\mathrm{exp}(z^2/2)f(z)`$,
$$𝒜[\phi ^a,\psi ^a,y^a,z]=\underset{a,\sigma }{}\overline{\phi }_\sigma ^a\left[𝒢_a^1\sigma H(y^a,z)\right]\phi _\sigma ^a\frac{2}{J}\underset{\omega }{}|\psi ^a(\omega )|^2$$
(8)
and $`H(y^a,z)=I\sqrt{q}z+I\sqrt{\stackrel{~}{q}q}y^a`$ is effective local magnetic field. Note, that the variable $`q=<S_i^aS_i^b>`$ corresponds to Edwards-Anderson SG order parameter when the limit $`n0`$ is taken. Nevertheless, the diagonal element $`\stackrel{~}{q}`$ can be set neither zero nor one, in contrast to the classical Ising glass theory because of dynamical effects due to the interaction with CE ”bath”. To take into account this interaction we include replica dependent magnetic field into bare Green’s function $`𝒢_{0\sigma }^a=(i\omega _n\sigma H(y^a,z))^1`$ and perform the integration over PF Grassmann variables with the help of the expression
$$\mathrm{𝚃𝚛}\mathrm{ln}\left(𝒢_a^1\sigma H\right)=\mathrm{ln}\left(2\mathrm{cosh}(\beta H)\right)+\mathrm{𝚃𝚛}\underset{m=1}{\overset{\mathrm{}}{}}\frac{(1)^{m+1}}{m}\left(𝒢_{0\sigma }^a(H)\mathrm{\Sigma }(\psi ^a)\right)^m$$
(9)
where $`\mathrm{\Sigma }(\psi ^a)=T_ϵ\psi _i^a(ϵ+\omega _{n_1})G_0(i_i,ϵ)\psi _i^a(ϵ+\omega _{n_2})`$ depends on the variable $`\psi `$ ”responsible” for Kondo interaction. Calculating the first term in expansion (9) one gets the following expression for the effective ”bosonic” action in the one loop approximation
$$𝒜[\psi ^a,H]=\mathrm{ln}\left(2\mathrm{cosh}(\beta H(y^a,z))\right)\frac{2}{J}\underset{n}{}\left[1J\mathrm{\Pi }(i\mathrm{\Omega }_n,H(y^a,z))\right]|\psi ^a|^2O(|\psi ^a|^4)$$
(10)
The polarization operator $`\mathrm{\Pi }`$ in the limit $`T,Hϵ_F`$ is given by:
$$\begin{array}{c}\mathrm{\Pi }(i\mathrm{\Omega }_n,H)=\beta ^1\underset{n,\stackrel{}{k},\sigma }{}G_0(\stackrel{}{k},iϵ_n+i\mathrm{\Omega }_n)𝒢_{0\sigma }(iϵ_n,H)\stackrel{\mathrm{\Omega }_n=0}{}\hfill \\ \hfill \stackrel{\mathrm{\Omega }_n=0}{}\rho (0)\left[\mathrm{ln}\left(\frac{ϵ_F}{\sqrt{H^2+\pi ^2\beta ^2/4}}\right)+\frac{\pi }{2\mathrm{cosh}(\beta H)}+O\left(\frac{H^2}{ϵ_F^2}\right)\right]\end{array}$$
(11)
When $`H=0`$, the coefficient in front of $`|\psi ^a|^2`$ in (10) changes its sign at $`T^{}ϵ_F\mathrm{exp}(\alpha ^1)`$. This is a manifestation of single-impurity Kondo effect (see, e.g. ,).
One can now perform the Gaussian integration over $`\psi ^a`$ fields in (7) by stationary phase method $`D[\psi ^a]\mathrm{exp}(\delta 𝒜[\psi ^a])=\mathrm{exp}\left(\mathrm{𝚃𝚛}\mathrm{ln}\left[1J\mathrm{\Pi }(i\mathrm{\Omega }_n,H(y^a,z))\right]\right)`$. After the last step, namely integration over replica dependent field $`y^a`$ the limit $`n0`$ can be taken. The free energy per site $`f=\beta ^1lim_{n0}(1<Z^n>_{av})/(nN)`$ is given by
$$\beta f(\stackrel{~}{q},q)=\frac{1}{4}(\beta I)^2\left(\stackrel{~}{q}^2q^2\right)_z^G\mathrm{ln}\left(_y^G\frac{2\mathrm{cosh}(\beta H(y,z))}{1J\mathrm{\Pi }(0,H(y,z))}\right).$$
(12)
New equations for $`q`$, $`\stackrel{~}{q}`$ are determined by conditions $`f(\stackrel{~}{q},q)/\stackrel{~}{q}=0`$, $`f(\stackrel{~}{q},q)/q=0`$:
$$\frac{1}{2}(\beta I)^2\stackrel{~}{q}=_z^G\frac{\mathrm{ln}}{\stackrel{~}{q}},\frac{1}{2}(\beta I)^2q=_z^G\frac{\mathrm{ln}}{q},=_y^G\frac{2\mathrm{cosh}(\beta H(y,z))}{1J\mathrm{\Pi }(0,H(y,z))}$$
(13)
The equations (12 \- 13) contain the key result of the paper. They represent the solution of KL problem with quenched disorder on a replica symmetrical level. To demonstrate some interesting physical effects, described by these equations let us consider the case $`TT_{SG}`$$``$$`T^{}`$ (Kondo high temperature limit). Since $`H(y^a,z)`$ is dynamical variable, we break the parametrical region of $`H`$ to several pieces. First, when $`H`$$``$ $`T`$,$`T^{}`$, the logarithm in (11) is cut by $`H`$ and there are no temperature dependent Kondo corrections to the MF equations. This corresponds to the limit $`T^{}`$ $``$ $`I`$ providing frozen spins and preventing them from resonance scattering <sup>2</sup><sup>2</sup>2We also note, that when $`T^{}`$ $``$ $`I`$ the SG transition does not happen.. Nevertheless, when $`T^{}I`$, the region $`H`$ $``$ $`T`$ becomes very important. We calculate $``$ expanding the r.h.s of (12) up to $`(H/T)^2`$
$$\mathrm{ln}(C_{z,\stackrel{~}{q},q})=\frac{1}{2}\mathrm{ln}\left(1+\gamma u^2r^2\right)+\frac{u^2}{2}\frac{r^2q\gamma z^2}{1+\gamma u^2r^2}+\mathrm{ln}\left[\mathrm{cosh}\left(\frac{uz\sqrt{q}}{1+\gamma u^2r^2}\right)\right]$$
(14)
We use the following short-hand notations: $`u=\beta I`$, $`\gamma =2c/\mathrm{ln}(T/T^{})`$, $`r^2=\stackrel{~}{q}q`$, $`C=2c\alpha /\gamma `$ with $`c=(\pi /4+2/\pi ^2)1`$. We note again that when $`J=0`$, which corresponds to the absence of Kondo interaction, $`(z,\stackrel{~}{q},q)=\mathrm{exp}\left(\frac{1}{2}(\beta I)^2(\stackrel{~}{q}q)\right)\mathrm{cosh}(\beta Iz\sqrt{q})`$ and standard SK equation takes place, providing, e.g. an exact identity $`\stackrel{~}{q}=1`$.
In the vicinity of the phase transition point Eq.(13) reads:
$$\begin{array}{c}\stackrel{~}{q}=1\frac{2c}{\mathrm{ln}(T/T^{})}+O\left(\frac{1}{\mathrm{ln}^2(T/T^{})}\right),\hfill \\ \hfill q=_z^G\mathrm{tanh}^2\left(\frac{\beta Iz\sqrt{q}}{1+2c(\beta I)^2(\stackrel{~}{q}q)/\mathrm{ln}(T/T^{})}\right)+O\left(\frac{q}{\mathrm{ln}^2(T/T^{})}\right)\end{array}$$
(15)
These equations describe a second-order SG transition in Ising-like SK <sup>3</sup><sup>3</sup>3When an Ising system described by (1) with nearest neighbor interaction is treated with MF theory, equations identical to (13) are obtained with $`\sqrt{Z}I`$ replacing $`I`$, where $`Z`$ is average number of neighbors. system coupled with CE ”bath” in the presence of Kondo scattering. Taking the limit $`q0`$ we estimate the temperature of SG transition $`(T_{SG}/I)^2=14c/\mathrm{ln}(T_{SG}/T^{})\mathrm{}<1`$. Thus, the Kondo scattering resonance results in depressing of SG transition temperature due to the screening effects in the same way as magnetic moments and one-site susceptibility are screened in single-impurity Kondo problem . This screening shows up at large time scale $`t`$ $``$ $`1/T^{}`$ and affects both diagonal and nondiagonal elements of $`Q`$ matrix. Moreover, $`\stackrel{~}{q}`$ becomes partially screened well above the SG transition point. Recalling that $`H`$ $``$ $`Iy\sqrt{\stackrel{~}{q}}`$ one can see that our assumption $`H/T`$ $``$ $`1`$ is consistent with Eq.(15) even if $`T`$ $``$ $`T_{SG}`$. It is necessary to note, that a growing SG order parameter in Eq.(10-11) suppresses Kondo effect as well as providing a broader validity domain for equations (15). We leave the self-consistent analysis of Eq.(12,15) for a future detailed publication.
In conclusion, we have considered the Kondo high temperature limit (in a sence of $`T>T^{}`$) of a KL model with quenched disorder. We derived new MF equations for SG transition in the presence of strong Kondo scattering and have shown that the partial screening of both diagonal and nondiagonal elements of $`Q`$ matrix takes place. As a result, the temperature of SG transition is strongly suppressed when Ising and Kondo interactions are of the same order of magnitude.
We thank F.Bouis, B.Coqblin, K.Kikoin, and P.Pfeuty for useful discussions. This work is supported by the SFB410 (II-VI semiconductors). One of us (MNK) is grateful to Alexander von Humboldt Foundation for support during his stay in Germany.
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# Group Cohomology and Gauge Equivalence of some Twisted Quantum Doubles
## 1 Introduction
The purpose of the present paper is to investigate relationships that exist between certain quasi-Hopf algebras (namely, twisted quantum doubles of a finite abelian group) and their module categories. This is closely related to aspects of group cohomology concerning the relationship of the bar complex to some other complexes defined and studied by Eilenberg and MacLane \[Mac52\], \[EM50a\] and \[EM50b\]. Although our results are purely homological and algebraic in nature, the motivation for studying these questions derives mainly from connections with conformal field theory (CFT) and vertex operator algebras \[DW90\], \[DPR92\], \[MS89\], \[DL93\]. We will thus review background motivation (although it is hardly necessary for an understanding of our main results) and then give a more detailed overview of the paper.
It is not necessary for the reader to appreciate what a vertex operator algebra (VOA) is (cf. \[Bor86\], \[FLM88\], \[FHL93\]), but only to understand this: VOAs have a representation theory, hence a (linear) module category, and it is important to understand the nature of this category. It is expected that the module category $`V\text{-}\text{Mod}`$ for a VOA $`V`$ has the structure of a braided monoidal category \[JS93\]. Roughly, this means that there is a notion of tensor product of $`V`$-modules, and that the tensor product of three modules is associative by an isomorphism that depends on the particular modules. Moreover, tensor product of modules is commutative, again via an isomorphism that depends on the particular modules. Particularly nice VOAs called rational have the additional property that $`V\text{-}\text{Mod}`$ is a semi-simple category, so that $`V`$ has only a finite number of isomorphism classes of simple modules, and every object in $`V\text{-}\text{Mod}`$ is a direct sum of simple objects. The simplest VOAs of all from this perspective are called holomorphic: they are rational and have but one simple module, namely the adjoint module $`V`$ itself. Thus if $`V`$ is a holomorphic VOA, then $`V\text{-}\text{Mod}`$ is equivalent (as a semi-simple braided monoidal category) to the category of vector spaces over $``$. Examples of holomorphic VOAs include the famous Moonshine Module and VOAs attached to positive-definite, even, self-dual lattices \[FLM88\].
It is a very difficult problem to establish the nature of $`V\text{-}\text{Mod}`$, even for relatively well-understood VOAs. However, if $`V\text{-}\text{Mod}`$ is indeed a braided monoidal category, then reconstruction theory (cf. \[HO97\], \[Maj92\] ) suggests that there might be a quasi-Hopf algebra $`Q`$ with the property that $`Q\text{-}\text{Mod}`$ is equivalent ( as braided monoidal category) to $`V\text{-}\text{Mod}`$. (Recall that the module category of a quasi-Hopf algebra automatically carries a braided monoidal structure, cf. \[Kas95\], \[Dri90\].) In this context, Dijkgraaf, Pasquier and Roche \[DPR92\] put forward a remarkable proposal which essentially identifies such a $`Q`$ in the case that $`V`$ is a so-called holomorphic orbifold. This means that there is a holomorphic VOA $`W`$ and a finite group of automorphisms $`G`$ of $`W`$ such that $`V=W^G`$, i.e., $`V`$ is the subVOA of $`G`$-fixed-points. The DPR ansatz is that one can take $`Q`$ to be a twisted quantum double $`D^\omega (G)`$ of $`G`$. This is a certain quasi-Hopf algebra, first constructed in \[DPR92\], which is a version of the Drinfeld double of the group algebra $`[G]`$ \[Dri87\], but which is twisted by a certain 3-cocycle $`\omega `$ in $`Z^3(G,^{})`$. For more details see \[DPR92\], \[Mas95\], and section 2 below.
One of the problems of vertex operator algebras is that of classification. There are far too many VOAs to make the problem of classification up to isomorphism at all practical, at least for the foreseeable future. A meaningful alternative is to classify VOAs up to equivalence of their module categories. Thus we say that two VOAs $`V`$ and $`W`$ are physically equivalent if $`V\text{-}\text{Mod}`$ and $`W\text{-}\text{Mod}`$ are equivalent as linear, braided monoidal categories. For example, all holomorphic VOAs are equivalent in this sense because they have identical module categories. In the present paper we are interested in the case of holomorphic orbifolds: given two holomorphic VOAs $`W_1`$ and $`W_2`$, and finite automorphism groups $`G_1`$ and $`G_2`$ of $`W_1`$ and $`W_2`$ respectively, when are the corresponding holomorphic orbifolds physically equivalent? Because of the DPR ansatz, one expects that this is equivalent to the following purely algebraic question:
(\*) given finite groups $`G_1`$ and $`G_2`$ and normalized 3-cocycles $`\omega _1Z^3(G_1,^{})`$ and $`\omega _2Z^3(G_2,^{})`$, when are the module categories $`D^{\omega _1}(G_1)\text{-}\text{Mod}`$ and $`D^{\omega _2}(G_2)\text{-}\text{Mod}`$ equivalent as (linear) braided monoidal categories?
As far as we know, the question in this generality has not been discussed. In the present paper we will limit ourselves to the case in which the relevant groups and algebras are commutative. Note that it is well-known \[Kas95\] that a sufficient condition for the equivalence of the categories in (\*) is that the corresponding quasi-bialgebras are gauge equivalent. We will prove (Theorem 10.5) that this is also a necessary condition, at least when the twisted quantum doubles are commutative. In practice, however, this is not very helpful in trying to decide the answer to (\*) in a given case. The main results and methods of this paper show essentially that there is a method that allows one to answer (\*) in any given case. We also establish results that give a complete answer for various classes of groups.
Let us begin by considering a finite abelian group $`G`$. With each normalized 3-cocycle $`\omega `$ in $`Z^3(G,^{})`$ we consider the twisted quantum double $`D^\omega (G)`$ together with its group of group-like elements $`\mathrm{\Gamma }^\omega (G)`$. Now $`\mathrm{\Gamma }^\omega (G)`$ is an abelian group, and we will be particularly interested in the case in which $`\mathrm{\Gamma }^\omega (G)`$ spans $`D^\omega (G)`$: this happens precisely when $`D^\omega (G)`$ is isomorphic as a bialgebra to the group algebra $`[\mathrm{\Gamma }^\omega (G)]`$, in which case $`\mathrm{\Gamma }^\omega (G)`$ is a central extension of $`G`$ by its group of characters $`\widehat{G}`$.
$$1\widehat{G}\mathrm{\Gamma }^\omega (G)G1$$
(1)
However, there is an essential difference between the two bialgebras which generally means that they are not isomorphic as quasi-bialgebras: the associator of $`D^\omega (G)`$ carries information about the cocycle $`\omega `$, whereas the associator of the group algebra is trivial and carries no cohomological information.
We call the normalized 3-cocycle $`\omega `$ abelian in case $`\mathrm{\Gamma }^\omega (G)`$ spans $`D^\omega (G)`$, and denote by $`Z^3(G,^{})_{ab}`$, resp $`H^3(G,^{})_{ab}`$ the 3-cocycles, resp. 3-cohomology classes which are abelian in this sense. They are groups, and via (1) we get a group homomorphism
$$\mathrm{\Lambda }:H^3(G,^{})_{ab}H_{ab}^2(G,\widehat{G})$$
(2)
where the subscript ab on the second cohomology group refers to abelian extensions of $`G`$ by $`\widehat{G}`$.
It transpires that many of our results hinge on a close analysis of this map. Though we have described $`\mathrm{\Lambda }`$ in terms of twisted quantum doubles, it turns out to have a purely homological interpretation. This arises as follows: Eilenberg and MacLane in \[Mac52\], \[EM50a\] and \[EM50b\] described a sequence of complexes $`A_i(G)`$ associated to an abelian group $`G`$ with a view to computing the cohomology of the Eilenberg-MacLane spaces $`K(G,n)`$. We are only interested in the complex $`A_0(G)`$, which is just the bar complex, and the second complex $`A_1(G)`$ which contains $`A_0(G)`$ as a subcomplex. This containment yields a canonical short exact sequence of complexes
$$0\text{Hom}(B(G),^{})\text{Hom}(A_1(G),^{})\text{Hom}(A_0(G),^{})0$$
($`B(G)`$ is the quotient complex $`A_1(G)/A_0(G)`$), and thereby a long exact sequence in cohomology. It turns out that the map $`\mathrm{\Lambda }`$ in (2) is the restriction to $`H^3(G,^{})_{ab}`$ of the connecting homomorphism which maps $`H^3(G,^{})`$ to $`H^4(B(G),^{})`$. There is naturally a close connection between $`\mathrm{\Lambda }`$ and the cohomology of the complex $`A_1(G)`$, which is sometimes called abelian cohomology. We observe that this latter cohomology already plays a role in \[MS89\] and in \[JS93\], and is quite essential in the work of Dong and Lepowsky \[DL93\] on the generalized Jacobi identity satisfied by the vertex operators associated to lattice elements in a lattice VOA. But the context of these works is quite different to our own.
Our analysis of the long exact sequence also reveals a fundamental difference between groups of odd order and groups of even order : we will see that $`\mathrm{\Lambda }`$ is injective if, and only if, $`G`$ has odd order, and that in general $`\text{ker }\mathrm{\Lambda }`$ is isomorphic to the group<sup>2</sup><sup>2</sup>2For a finite abelian group $`G`$, $`\mathrm{\Omega }_p(G)`$ is the subgroup of elements of order dividing $`p`$. $`\mathrm{\Omega }_2(G)`$ of elements in $`G`$ of order at most 2. We also give a precise description of $`\text{Im }\mathrm{\Lambda }`$ : in the easier case of groups of odd order it is precisely the group of invariants of the duality map $`ϵ`$ : $`H_{ab}^2(G,\widehat{G})H_{ab}^2(G,\widehat{G})`$ which maps a short exact sequence of abelian groups to the dual short exact sequence obtained by applying the functor $`\text{Hom}(\mathrm{?},^{})`$. In case $`G`$ is an abelian 2-group, it will turn out from the structure of $`\text{Im }\mathrm{\Lambda }`$ that the extension $`\mathrm{\Gamma }^\omega (G)`$ is necessarily a product of an even number of cyclic factors. This and other results in the paper have consequences for the theory of VOAs and CFT that we hope to pursue elsewhere.
The main consequence of our analysis of the long exact sequence is that we are able to characterize the pair $`(G,\omega )`$ for $`\omega `$ abelian in terms that are more amenable to calculation. Namely, we show that it is equivalent to the existence of a certain non-degenerate quadratic form on the abelian group $`\mathrm{\Gamma }^\omega (G)`$, call it $`q`$, with the property that $`q`$ has a metabolizer $`\widehat{G}`$. In general we say that a non-degenerate quadratic space $`(E,q)`$ for a finite abelian group $`E`$ has a metabolizer $`M`$, if the restriction of $`q`$ to $`M`$ vanishes identically and if the order of $`E`$ is the square of the order of $`M`$. So here we are following the definition of metabolizer in \[tD79\] rather than the classical definition \[Kne70\], \[MH73\], where $`M`$ is also required to be a direct summand of $`E`$. In the case we are interested in, this only happens if the extension (1) splits, and we say that $`M`$ is a split metabolizer in this case. For us, metabolizers are not necessarily split and the most interesting metabolizers are not. In this context we may speak of metabolic triples $`(E,q,G)`$, and we can try to construct suitable Witt-type groups. We essentially carry this out when $`E`$ is a homogeneous (or homocyclic) abelian $`p`$-group, however we do not use the language of Witt groups. We will see (Theorem 10.4) that the equivalence (\*) for abelian groups $`G_1`$, $`G_2`$ and abelian 3-cocycles $`\omega _1,\omega _2`$ is the same as the equivalence of the associated quadratic spaces $`(\mathrm{\Gamma }^{\omega _1}(G_1),q_1)`$ and $`(\mathrm{\Gamma }^{\omega _2}(G_2),q_2)`$.
We point out that quadratic forms play a prominent role in several of the papers we have already quoted concerning abelian cohomology and/or VOAs (\[DL93\], \[JS93\], \[MS89\]). But in \[DL93\] and \[MS89\], for example, conformal field theories more general than holomorphic orbifolds are considered, so that the existence of a metabolizer, so crucial for us, plays no role in these works.
Once the connection with quadratic forms with a metabolizer has been forged, a number of further results flow. Consider a finite abelian group $`G`$ together with an abelian 3-cocycle $`\omega `$. By analogy with the standard decomposition of $`G`$ into a direct product of cyclic groups, one may ask whether $`D^\omega (G)`$ is gauge equivalent to a tensor product
$$D^\omega (G)D^{\omega _1}(C_1)\mathrm{}D^{\omega _r}(C_r)$$
(3)
with each $`C_i`$ a cyclic group. One may assume that $`G`$ is a $`p`$-group for some prime $`p`$, in which case the question has an affirmative answer in the following cases: (a) $`p`$ satisfies the congruence $`p1(mod4)`$; (b) $`G`$ is a homogeneous $`p`$-group and $`p`$ is odd. However, the result is generally false if $`p=2`$, and $`p3(mod4)`$, and $`G`$ is not homogeneous. Furthermore, if there is an equivalence of the form (3) then it is generally not unique, and the direct product $`C_1\times \mathrm{}\times C_r`$ of the groups $`C_i`$ is not necessarily isomorphic to $`G`$. It may seem curious that one has to impose congruence conditions on $`p`$ in order to establish (3) for all abelian $`p`$-groups. The reason is that via the quadratic forms we can use the technique of Gauss sums, and odd primes which are sums of two squares are better behaved. The proof of (3) for homogeneous p-groups is established using Witt-group type techniques. The paper of C.T.C. Wall \[Wal63\], which classifies the quadratic forms on finite abelian groups (at least if they are of odd order) proves to be invaluable in many of our calculations.
By a rational lattice we mean a free abelian group $`M`$ of finite rank together with a non-degenerate bilinear form $`,:M\times M`$. We do not necessarily assume that $`,`$ is positive-definite. The lattice is called integral in case it takes values in $``$, and even if, in addition, one has $`x,x2`$. Such lattices play a prominent role in \[DL93\] because they naturally give rise to vertex algebras, indeed to examples of what Dong-Lepowsky call “abelian intertwining algebras.” Via the connection with quadratic forms, we are able to relate the elements of $`H^3(G,^{})_{ab}`$ to certain pairs of lattices $`ML`$, essentially because such a pair (together with the bilinear form $`,`$ ) yields a quadratic form on $`L/M`$. Moreover, a theorem of Wall \[Wal63\] allows us to see that the existence of a metabolizer implies that the lattice $`M`$ is necessarily self-dual (i.e., unimodular). As we will show elsewhere, this has consequences for the theory of lattice VOAs. It also permits us to see that the following pieces of data are essentially equivalent<sup>3</sup><sup>3</sup>3a vague phrase which will be clarified in section 11 for a finite abelian group $`G`$:
1. $`[\omega ]H^3(G,^{})_{ab}`$.
2. a metabolic triple $`(E,q,\widehat{G})`$.
3. a pair of rational lattices $`ML`$ with $`M`$ even and self-dual and $`L/M=G`$.
We will also prove (Theorem 13.8) that if $`G`$ is a homogeneous abelian $`p`$-group of odd order then two twisted quantum doubles $`D^\omega (G)`$ and $`D^\omega ^{}(G)`$ for abelian cocycles $`\omega `$ and $`\omega ^{}`$ are gauge equivalent if, and only if, the corresponding cohomology classes $`[\omega ]`$ and $`[\omega ^{}]`$ are equivalent under the action of the automorphism group $`\text{Aut}(G)`$. In some sense this result is the quintessential goal that one seeks for any abelian $`p`$-group $`G`$; but the result is false in general due to “hidden” gauge equivalences which are hard to enumerate, but which can be detected via the metabolic triples.
Finally, the connection with quadratic forms allows us to construct various kinds of dualities and symmetries (other than categorical equivalence) between the module categories in (\*). We give various examples in section 14, including the following: for positive integers $`n`$, $`k`$ and odd prime $`p`$, let $`G_1`$ and $`G_2`$ denote the homogeneous groups $`\left(_{p^n}\right)^k`$ and $`\left(_{p^k}\right)^n`$ respectively. Then there are precisely $`\left(\genfrac{}{}{0pt}{}{n+k}{k}\right)`$ equivalence classes of monoidal categories (i.e., tensor categories excluding braiding) both of the form $`D^{\omega _1}(G_1)\text{-}\text{Mod}`$ and $`D^{\omega _2}(G_2)\text{-}\text{Mod}`$, for some $`\omega _1`$ or $`\omega _2`$ respectively. Moreover there is a canonical bijection (duality) between these two sets of equivalence classes of tensor categories. Each such tensor category can be labeled by a partition, and at the level of partitions the duality is simply that which maps a partition to its dual (or conjugate) partition. Moreover, duality induces a bijection between the braided monoidal categories associated to each group.
The paper is organized as follows: we give some background in section 2, including the observation that for a finite abelian group and normalized 3-cocycle $`\omega Z^3(G,^{})`$, the object $`D^\omega (G)`$–when equipped with the trivial associator–is a semi-simple self-dual Hopf algebra which is generally neither commutative nor cocommutative. We study the group-like elements of $`D^\omega (G)`$ in section 3, giving rise to the central extension (1). After a short section 4 concerned with tensor products of twisted quantum doubles, we take up in sections 5 \- 8 the structure of the long exact sequence alluded to above. Much of the resulting homological information is encoded in a 7 term exact sequence which we record in Remark 8.6. In section 9 we consider the category $`D^\omega (G)\text{-}\text{Mod}`$ as a monoidal category i.e., forgetting the braiding. For example, we show (Lemma 9.2 and Theorem 9.4) that for finite abelian groups $`G`$ of odd order and normalized 3-cocycles $`\omega _1`$ and $`\omega _2`$ in $`Z^3(G,^{})_{ab}`$, the monoidal categories $`D^{\omega _1}(G)\text{-}\text{Mod}`$ and $`D^{\omega _2}(G)\text{-}\text{Mod}`$ are tensor equivalent if, and only if, the corresponding groups of fusion rules $`\mathrm{\Gamma }^{\omega _1}`$ and $`\mathrm{\Gamma }^{\omega _2}`$ are isomorphic, and that in this case each $`D^{\omega _i}(G)`$ is gauge equivalent to the group algebra $`[\mathrm{\Gamma }^{\omega _i}]`$. This answers in the affirmative–at least for groups of odd order–a question posed in \[DPR92\], namely can $`D^\omega (G)`$ be obtained by twisting a Hopf algebra? We show by example that the answer is “no” in general. In section 10 we study the connections among gauge equivalence, braided monoidal categories and quadratic forms and prove the results already mentioned. Section 11 explains the connections with lattices, while section 12 presents some results involving Gauss sums and applies them to the proof of the gauge equivalence (3) in the case $`p1(mod4)`$. In section 13 we consider metabolic triples $`(E,q,G)`$ when $`G`$ is a homogeneous $`p`$-group and–in all but name–construct the corresponding Witt-type group. This is used (Theorem 13.8) to establish both (3) in the case that $`G`$ is a homogeneous abelian $`p`$-group, and other results concerning homogeneous $`p`$-groups already discussed. Finally, we discuss symmetries and dualities in section 14.
The authors thank Chongying Dong for useful discussions.
## 2 Twisted Quantum Double Of A Finite Group
Let $`G`$ be a finite group and $`\omega :G\times G\times G^{}`$ be a normalized 3-cocycle<sup>4</sup><sup>4</sup>4All cocycles will take values in a trivial $`G`$-module.. For any $`x,y,gG`$, define
$`\theta _g(x,y)`$ $`=`$ $`{\displaystyle \frac{\omega (g,x,y)\omega (x,y,(xy)^1gxy)}{\omega (x,x^1gx,y)}}`$ (4)
$`\gamma _g(x,y)`$ $`=`$ $`{\displaystyle \frac{\omega (x,y,g)\omega (g,g^1xg,g^1yg)}{\omega (x,g,g^1yg)}}`$ (5)
The twisted quantum double $`D^\omega (G)`$ of $`G`$ with respect to $`\omega `$ is the quasi-triangular quasi-Hopf algebra with underlying vector space $`(G)^{}G`$ and multiplication,
$$(e(g)x)(e(h)y)=\theta _g(x,y)e(g)e(xhx^1)xy$$
(6)
$$\mathrm{\Delta }(e(g)x)=\underset{hk=g}{}\gamma _x(h,k)e(h)xe(k)x$$
(7)
$$\mathrm{\Phi }=\underset{g,h,kG}{}\omega (g,h,k)^1e(g)1e(h)1e(k)1$$
(8)
where $`\{e(g)|gG\}`$ is the dual basis of the canonical basis of $`G`$ (cf. \[DPR92\]). The counit and antipode are given by
$$ϵ(e(g)x)=\delta _{g,1}\text{and}S(e(g)x)=\theta _{g^1}(x,x^1)^1\gamma _x(g,g^1)^1e(x^1g^1x)x^1$$
where $`\delta _{g,1}`$ is the Kronecker delta. The corresponding elements $`\alpha `$ and $`\beta `$ are $`1_{D^\omega (G)}`$ and $`\underset{gG}{}\omega (g,g^1,g)e(g)1`$ respectively. For the definition and more details about quasi-Hopf algebras, see \[Dri90\], \[Kas95\] or \[CP95\]. Verification of the detail involves the following identities, which result from the 3-cocycle identity for $`\omega `$:
$$\theta _z(a,b)\theta _z(ab,c)=\theta _{a^1za}(b,c)\theta _z(a,bc),$$
(9)
$$\theta _y(a,b)\theta _z(a,b)\gamma _a(y,z)\gamma _b(a^1ya,a^1za)=\theta _{yz}(a,b)\gamma _{ab}(y,z),$$
(10)
$$\gamma _z(a,b)\gamma _z(ab,c)\omega (z^1az,z^1bz,z^1cz)=\gamma _z(b,c)\gamma _z(a,bc)\omega (a,b,c),$$
(11)
for all $`a,b,c,y,zG`$. The universal $`R`$-matrix is given by
$$R=\underset{g,hG}{}e(g)1e(h)g.$$
(12)
###### Remark 2.1
1. Any Hopf algebra $`H`$ can be viewed as a quasi-Hopf algebra with the trivial associator $`\mathrm{\Phi }=1_H1_H1_H`$ and $`\alpha =\beta =1_H`$.
2. If $`w`$ and $`w^{}`$ are cohomologous 3-cocycles, then $`D^\omega (G)`$ and $`D^\omega ^{}(G)`$ are gauge equivalent as quasi-triangular quasi-bialgebras or simply gauge equivalent (cf. \[Kas95\]) with the algebra isomorphism $`\mathrm{\Theta }:D^\omega (G)D^\omega ^{}(G)`$ defined by
$$\mathrm{\Theta }(e(g)x)=\frac{b(g,x)}{b(x,x^1gx)}e(g)x$$
and the gauge transform $`FD^\omega ^{}(G)D^\omega ^{}(G)`$ given by
$$F=\underset{g,h}{}b(g,h)^1e(g)1e(h)1$$
where $`\omega ^{}=\omega \delta b`$ (cf. \[DPR92\]).
3. Let $`\sigma `$ be a group automorphism of $`G`$. Then $`\widehat{\sigma }\omega `$ defined by
$$\widehat{\sigma }\omega (g,h,k)=\omega (\sigma ^1(g),\sigma ^1(h),\sigma ^1(k))$$
is a 3-cocycle. Moreover, $`\sigma `$ induces an isomorphism $`\stackrel{~}{\sigma }`$ of quasi-triangular quasi-Hopf algebras from $`D^\omega (G)`$ to $`D^{\sigma \omega }(G)`$, namely
$$\stackrel{~}{\sigma }(e(g)x)=e(\sigma (g))\sigma (x).$$
For $`gG`$, denote by $`C_G(g)`$ the centralizer of $`g`$ in $`G`$. One can easily see that $`\theta _g`$ is a 2-cocycle of $`C_G(g)`$ with coefficient in $`^{}`$. Moreover, the map $`D_g:\omega \theta _g`$ induces a group homomorphism from $`H^3(G,^{})`$ to $`H^2(C_G(g),^{})`$. We will denote by $`H^3(G,^{})_{ab}`$ the group $`\underset{gG}{}\text{ker }D_g`$ and $`Z^3(G,^{})_{ab}`$ the subgroup of normalized 3-cocycles whose cohomology classes are in $`H^3(G,^{})_{ab}`$.
When $`G`$ is abelian, $`\theta _g=\gamma _g`$ and $`C_G(g)=G`$ for all $`gG`$ and we will simply write $`\theta _g`$ as $`\omega _g`$. It follows from equation (10) that we have
$$\frac{\omega _g(x,y)\omega _h(x,y)}{\omega _{gh}(x,y)}=\frac{\omega _{xy}(g,h)}{\omega _x(g,h)\omega _y(g,h)}.$$
Hence the map $`GH^2(G,^{})`$, $`g[\omega _g]`$ is a group homomorphism. Let us denote this map by $`\mathrm{\Omega }(\omega )`$. Notice that $`\mathrm{\Omega }(\omega )H^1(G,H^2(G,^{}))`$ and $`\mathrm{\Omega }(\omega \omega ^{})=\mathrm{\Omega }(\omega )\mathrm{\Omega }(\omega ^{})`$ for any other 3-cocycle $`\omega ^{}`$. Moreover, $`\mathrm{\Omega }(\omega )`$ is independent of the choice of representative of the cohomology class of $`\omega `$. Therefore, $`\mathrm{\Omega }`$ induces a group homomorphism $`\overline{\mathrm{\Omega }}`$ from $`H^3(G,^{})`$ to $`H^1(G,H^2(G,^{}))`$, and $`H^3(G,^{})_{ab}=\text{ker }\overline{\mathrm{\Omega }}`$.
###### Proposition 2.2
Let $`G`$ be a finite abelian group and $`\omega `$ a normalized 3-cocycle.
1. When equipped with the trivial associator, $`D^\omega (G)`$ is a self-dual Hopf algebra.
2. Let $`b:G\times G^{}`$ be a normalized 2-cochain and $`\omega ^{}=\omega \delta b`$. Then the map $`\mathrm{\Theta }:D^\omega (G)D^\omega ^{}(G)`$ defined by
$$\mathrm{\Theta }(e(g)x)=\frac{b(g,x)}{b(x,g)}e(g)x$$
is an isomorphism of Hopf algebras.
Proof. (i) Since $`G`$ is abelian, $`e(g)1`$ is in the center of $`D^\omega (G)`$. Hence, $`\beta \text{Cent}(D^\omega (G))`$ and $`\mathrm{\Phi }\text{Cent}(D^\omega (G)^3)`$. Therefore, $`\mathrm{\Delta }`$ is coassociative and $`S`$ satisfies the antipode conditions.
To show that $`D^\omega (G)`$ is self-dual, let us denote by $`\{f_{x,g}\}_{g,xG}`$ the dual basis of $`\{e(g)x\}_{g,xG}`$. Then,
$`f_{x,g}f_{y,h}(e(k)z)`$ $`=`$ $`(f_{x,g}f_{y,h})\left({\displaystyle \underset{ab=k}{}}\omega _z(a,b)e(a)ze(b)z\right)`$
$`=`$ $`\delta _{x,z}\delta _{y,z}\delta _{gh,k}\omega _x(g,h)`$
$`=`$ $`\delta _{x,y}\omega _x(g,h)f_{x,gh}(e(k)z).`$
Therefore,
$$f_{x,g}f_{y,h}=\delta _{x,y}\omega _x(g,h)f_{x,gh}.$$
Similar, one can derive that
$$\mathrm{\Delta }(f_{x,g})=\underset{uv=x}{}\omega _g(u,v)f_{u,g}f_{v,g}.$$
The identity, $`1_{D^\omega (G)^{}}`$, and the counit, $`ϵ`$, of $`D^\omega (G)^{}`$ are given by
$$1_{D^\omega (G)^{}}=\underset{xG}{}f_{x,1},ϵ(f_{x,g})=\delta _{x,1}.$$
respectively, Therefore, the linear map $`\phi :D^\omega (G)D^\omega (G)^{}`$ defined by
$$\phi (e(x)g)=f_{x,g}$$
(13)
is a Hopf algebra isomorphism.
(ii) Notice that for any $`x,y,gG`$
$$\omega _g^{}(x,y)=\omega _g(x,y)\delta b_g(x,y)$$
where $`b_g(x)=\frac{b(x,g)}{b(g,x)}`$. The result follows immediately from this identity. $`\mathrm{}`$
In the sequel, we will denote by $`D^\omega (G)_0`$ the Hopf algebra structure on $`D^\omega (G)`$, and we use $`D^\omega (G)`$ for the quasi-Hopf algebra structure with respect to the associator $`\mathrm{\Phi }`$ given by the equation (8).
## 3 Group-like Elements of $`D^\omega (G)`$
###### Definition 3.1
Let $`G`$ be a finite group and $`\omega `$ a normalized 3-cocycle. A nonzero element $`u`$ in $`D^\omega (G)`$ is called group-like if $`\mathrm{\Delta }(u)=uu`$. We will denote by $`\mathrm{\Gamma }^\omega (G)`$ the set of all group-like elements of $`D^\omega (G)`$. We simply write $`\mathrm{\Gamma }^\omega `$ for $`\mathrm{\Gamma }^\omega (G)`$ when the context is clear.
As in the case of group-like elements in a coalgebra (cf. \[Swe69\]), $`\mathrm{\Gamma }^\omega `$ is a linearly independent set and any group-like element $`u`$ is invertible in $`D^\omega (G)`$ with inverse $`S(u)`$. Since $`\mathrm{\Delta }:D^\omega (G)D^\omega (G)D^\omega (G)`$ is an algebra map, $`\mathrm{\Gamma }^\omega `$ is a subgroup of the group of units of $`D^\omega (G)`$. Moreover, $`\mathrm{\Gamma }^\omega `$ can be characterized by the following proposition.
###### Proposition 3.2
A nonzero element $`u`$ in $`D^\omega (G)`$ is a group-like element if, and only if, $`u=\underset{gG}{}\alpha (g)e(g)x`$ for some $`xG`$ and a map $`\alpha :G^{}`$ such that
$$\gamma _x(g,h)=\frac{\alpha (g)\alpha (h)}{\alpha (gh)}$$
for any $`g,hG`$.
Proof. The result follows from direct computation. $`\mathrm{}`$
Let $`\widehat{G}`$ be the character group of $`G`$. Since $`\gamma _1=1`$, $`\underset{gG}{}\alpha (g)e(g)1\mathrm{\Gamma }^\omega `$ for any $`\alpha \widehat{G}`$ by Proposition 3.2. The map $`\widehat{G}\mathrm{\Gamma }^\omega `$, $`\alpha \underset{gG}{}\alpha (g)e(g)1`$ is an injective group homomorphism. In the sequel, we often identify $`\widehat{G}`$ with the image of $`\widehat{G}`$ under this map. On the other hand, the map
$$D^\omega (G)G,\underset{g,xG}{}\alpha (g,x)e(g)x\underset{xG}{}\alpha (1,x)x$$
(14)
is an algebra map. Let $`B^\omega `$ be the image of $`\mathrm{\Gamma }^\omega `$ under this map. One can easily see that
$$B^\omega =\{xG|\gamma _x\text{ is a 2-coboundary.}\}.$$
###### Lemma 3.3
Let $`G`$ be a finite group. Then $`\widehat{G}`$ is in the center of $`D^\omega (G)`$. Moreover, $`\mathrm{\Gamma }^\omega `$ is a central extension
For each $`xB^\omega `$, let $`\gamma _x=\delta \tau _x`$ for a 1-cochain $`\tau _x:G^{}`$. The 2-cocycle $`\beta `$ associated to this central extension is given by
$$\beta (x,y)(g)=\frac{\tau _x(g)\tau _y(x^1gx)}{\tau _{xy}(g)}\theta _g(x,y).$$
Proof. For any $`\alpha \widehat{G}`$ and $`e(h)xD^\omega (G)`$,
$`(e(h)x)({\displaystyle \underset{gG}{}}\alpha (g)e(g)1)`$ $`=`$ $`{\displaystyle \underset{gG}{}}\alpha (g)e(h)e(x^1gx)x`$
$`=`$ $`{\displaystyle \underset{gG}{}}\alpha (xgx^1)e(h)e(g)x`$
$`=`$ $`{\displaystyle \underset{gG}{}}\alpha (g)e(h)e(g)x`$
$`=`$ $`({\displaystyle \underset{gG}{}}\alpha (g)e(g)1)(e(h)x).`$
Therefore, the first statement follows.
An element $`u\text{ker }\left(\mathrm{\Gamma }^\omega B^\omega \right)`$ if, and only if, $`u=\underset{gG}{}\alpha (g)e(g)1`$ and $`\gamma _1=\delta \alpha `$. Since $`\gamma _1=1`$, $`\alpha \widehat{G}`$. Hence,
$$\widehat{G}=\text{ker }\left(\mathrm{\Gamma }^\omega B^\omega \right).$$
Since $`\gamma _x=\delta \tau _x`$ for $`xB^\omega `$,
$$x\underset{gG}{}\tau _x(g)e(g)x$$
is a section of the map $`\mathrm{\Gamma }^\omega B^\omega `$. One can check directly that
$$\left(\underset{gG}{}\tau _x(g)e(g)x\right)\left(\underset{gG}{}\tau _y(g)e(g)y\right)=\underset{gG}{}\tau _x(g)\tau _y(x^1gx)\theta _g(x,y)e(g)xy.$$
Hence, the formula for the 2-cocycle associated to the central extension follows. $`\mathrm{}`$
###### Remark 3.4
In general, $`B^\omega `$ and the extension $`\mathrm{\Gamma }^\omega `$ depend on the individual cocycle. However, both of them are independent of the representative of the cohomology class $`[\omega ]`$ when $`G`$ is abelian.
###### Proposition 3.5
Let $`G`$ be a finite abelian group and $`\omega ,\omega ^{}`$ normalized 3-cocycles of $`G`$.
1. $`\mathrm{\Gamma }^\omega `$ is abelian.
2. If $`\omega `$ and $`\omega ^{}`$ are cohomologous, then $`B^\omega =B^\omega ^{}`$ and the central extensions
$$1\widehat{G}\mathrm{\Gamma }^\omega B^\omega 1\text{and}1\widehat{G}\mathrm{\Gamma }^\omega ^{}B^\omega ^{}1$$
are equivalent.
Proof. (i) By Proposition 2.2, $`D^\omega (G)_0`$ is a self-dual Hopf algebra and hence $`\mathrm{\Gamma }^\omega `$ is isomorphic to the group of group-like elements of $`D^\omega (G)_0^{}`$. Let $`\sigma _1`$, $`\sigma _2`$ be group-like elements of $`D^\omega (G)_0^{}`$. Then, $`\sigma _1`$, $`\sigma _2:D^\omega (G)`$ are algebra maps. Let $`V_1`$, $`V_2`$ be the 1-dimensional representations associated to $`\sigma _1`$ and $`\sigma _2`$ respectively. Then, $`V_1V_2`$ and $`V_2V_1`$ are the 1-dimensional representations associated to $`\sigma _1\sigma _2`$ and $`\sigma _2\sigma _1`$ respectively. Since $`V_1V_2V_2V_1`$ as $`D^\omega (G)`$ modules, $`\sigma _1\sigma _2=\sigma _2\sigma _1`$. Therefore, $`\mathrm{\Gamma }^\omega `$ is abelian.
(ii) Let $`\omega ^{}=\omega \delta b`$ for some normalized 2-cochain $`b:G\times G^{}`$. Then
$$\omega _g^{}(x,y)=\omega _g(x,y)\delta b_g(x,y)$$
where $`b_g(x)=b(x,g)/b(g,x)`$. Hence, $`\omega _g^{}`$ is a 2-coboundary if, and only if, $`\omega _g`$ is a 2-coboundary. Therefore, $`B^\omega ^{}=B^\omega `$. Since the map $`\mathrm{\Theta }:D^\omega (G)_0D^\omega ^{}(G)_0`$ defined in Proposition 2.2 is an Hopf algebra isomorphism, $`\mathrm{\Theta }(\mathrm{\Gamma }^\omega )=\mathrm{\Gamma }^\omega ^{}`$. Moreover, $`\mathrm{\Theta }`$ satisfies the commutative diagram :
$`\mathrm{}`$
###### Corollary 3.6
Let $`G`$ be a finite group and $`\omega `$ a normalized 3-cocycle. Then the following statements are equivalent :
1. $`D^\omega (G)`$ is spanned by $`\mathrm{\Gamma }^\omega `$.
2. $`G`$ is abelian and $`B^\omega =G`$.
3. $`G`$ is abelian and $`\omega Z^3(G,^{})_{ab}`$
4. $`D^\omega (G)`$ is a commutative algebra.
Proof. (i) $``$ (ii). It follows from Lemma 3.3 that
$$|\mathrm{\Gamma }^\omega |=|\widehat{G}||B^\omega |.$$
Therefore, $`D^\omega (G)`$ is spanned by $`\mathrm{\Gamma }^\omega `$ if and only if $`|\widehat{G}|=|G|`$ and $`|B^\omega |=|G|`$ which is equivalent to the statement (ii).
The equivalence of (ii) and (iii) follows directly from the definition of $`Z^3(G,^{})_{ab}`$ and equation (6). (ii) $``$ (iv). Since $`G`$ is abelian and $`B^\omega =G`$, $`\omega _g`$ is a coboundary for any $`gG`$. In particular, $`\omega _g(x,y)=\omega _g(y,x)`$ for any $`x,yG`$. Therefore, the multiplication on $`D^\omega (G)`$ is commutative.
(iv) $``$ (ii) Assume $`D^\omega (G)`$ is commutative. By the surjectivity of the algebra map defined in (14), $`G`$ is abelian and hence $`\theta _g=\gamma _g=\omega _g`$. Moreover, the commutative multiplication in $`D^\omega (G)`$ implies that
$$\omega _g(x,y)=\omega _g(y,x)$$
for any $`x,y,gG`$. Therefore, $`\omega _g`$ is a 2-coboundary for any $`gG`$. $`\mathrm{}`$
In the sequel, if $`G`$ is a finite abelian group we will denote by $`\mathrm{\Lambda }_G[\omega ]`$ the cohomology class in $`H^2(B^\omega ,\widehat{G})`$ which corresponds to the central extension
In particular, if $`[\omega ]H^3(G,^{})_{ab}`$, then $`B^\omega =G`$ and so $`\mathrm{\Lambda }_G[\omega ]H^2(G,\widehat{G})`$. We simply write $`\mathrm{\Lambda }`$ instead of $`\mathrm{\Lambda }_G`$ when there is no ambiguity.
###### Remark 3.7
If one of the conditions (i)–(iv) of Corollary 3.6 holds then we have an isomorphism of Hopf algebras $`D^\omega (G)_0[\mathrm{\Gamma }^\omega ]`$
###### Proposition 3.8
The map $`\mathrm{\Lambda }:H^3(G,^{})_{ab}H^2(G,\widehat{G})`$, $`[\omega ]\mathrm{\Lambda }[\omega ]`$ is a group homomorphism.
Proof. Let $`[\omega ]H^3(G,^{})_{ab}`$. Then for each $`xG`$, there is a normalized 1-cochain $`\tau _x`$ such that $`\delta \tau _x=\omega _x`$. Then $`\mathrm{\Lambda }[\omega ]`$ can be represented by $`\beta `$ where $`\beta :G\times G\widehat{G}`$ is given by
$$\beta (x,y)(z)=\frac{\tau _x(z)\tau _y(z)}{\tau _{xy}(z)}\omega _z(x,y).$$
(15)
The result follows easily from this formula. $`\mathrm{}`$
Let $`G`$ be a finite abelian group and $`\omega Z^3(G,^{})_{ab}`$. Denote by $`T(\omega )`$ the set of all normalized 2-cochains $`\tau `$ on $`G`$ such that
$$\omega _x=\delta \tau _x$$
for $`xG`$ where $`\tau _x(y)=\tau (x,y)`$. Pick any $`\tau T(\omega )`$. Denote by $`\sigma (\alpha ,x)`$ the group-like element
$$\underset{gG}{}\alpha (g)\tau _x(g)e(g)x$$
for any $`\alpha \widehat{G}`$ and $`xG`$. We will write $`\sigma _\tau (\alpha ,x)`$ if we wish to emphasize the dependence of $`\sigma `$ on $`\tau `$. By Lemma 3.3, for any $`u\mathrm{\Gamma }^\omega `$, there exist unique $`\alpha \widehat{G}`$ and $`xG`$ such that $`u=\sigma (\alpha ,x)`$. It follows from Proposition 3.8 that
$$\sigma (\alpha ,x)\sigma (\lambda ,y)=\sigma (\alpha \lambda \beta (x,y),xy)$$
(16)
for any $`\alpha ,\lambda \widehat{G}`$ and $`x,yG`$ where $`\beta `$ is given by equation (15). We will simply write $`\mathrm{\Lambda }(\omega )`$ for the 2-cocycle given in equation (15) whenever the context is clear.
## 4 Tensor Products
We consider here the structure of the tensor product of two quantum doubles. Let $`G`$, $`H`$ be finite groups and $`\omega ,\omega ^{}`$ normalized 3-cocycles on $`G`$ and $`H`$ respectively (with $`^{}`$ coefficients as always). Of course there are canonical projections $`G\times HG`$, $`G\times HH`$, and we let $`infl\omega `$, $`infl\omega ^{}`$ denote the 3-cocycles in $`Z^3(G\times H,^{})`$ obtained by inflating $`\omega `$, $`\omega ^{}`$ using these projections. Let $`\zeta =(infl\omega )(infl\omega ^{})`$.
###### Proposition 4.1
There is a natural isomorphism of quasi-triangular quasi-Hopf algebras
given by $`\iota :e(g)xe(h)ye(g,h)(x,y)`$ for any $`g,xG`$ and $`h,yH`$.
This means that $`\iota `$ is an isomorphism of bialgebras, and that it maps the Drinfeld associators, $``$-matrices and elements $`\alpha `$, $`\beta `$ (as given in section 2) to the corresponding objects associated with $`D^\zeta (G\times H)`$. The calculations needed to prove the proposition are essentially routine and we omit the proof.
The group-like elements in $`D^\omega (G)D^\omega ^{}(H)`$ are of the form
$$\left(\underset{gG}{}\lambda _x(g)e(g)x\right)\left(\underset{hH}{}\lambda _y^{}(h)e(h)y\right)$$
where $`xB^\omega `$, $`yB^\omega ^{}`$ and $`\lambda _x`$ and $`\lambda _y^{}`$ are functions on $`G`$ and $`H`$ respectively such that $`\delta \lambda _x=\gamma _x`$, $`\delta \lambda _y^{}=\gamma _y^{}`$ and $`\gamma `$, $`\gamma ^{}`$ are those $`\gamma `$’s associated to $`\omega `$ and $`\omega ^{}`$ respectively (cf.(5)). It is easy to see that the group of all group-like elements in $`D^\omega (G)D^\omega ^{}(H)`$ is naturally isomorphic to $`\mathrm{\Gamma }^\omega (G)\times \mathrm{\Gamma }^\omega ^{}(H)`$. As $`\iota `$ induces an isomorphism from the group of group-like elements of $`D^\omega (G)D^\omega ^{}(H)`$ to $`\mathrm{\Gamma }^\zeta (G\times H)`$,
$$\mathrm{\Gamma }^\omega (G)\times \mathrm{\Gamma }^\omega ^{}(H)\stackrel{\iota }{}\mathrm{\Gamma }^\zeta (G\times H).$$
###### Proposition 4.2
Let $`G`$ be a finite group given as a direct product of groups $`G=H\times K`$ with $`|H|`$ and $`|K|`$ coprime. Then for any normalized 3-cocycle $`\omega `$ of $`G`$, there exist normalized 3-cocycles $`\eta `$, $`\eta ^{}`$ on $`K`$, $`H`$ respectively such that $`D^\omega (G)`$ and $`D^\eta (H)D^\eta ^{}(K)`$ are equivalent as quasi-triangular quasi-bialgebras.
Proof. It is well-known \[Bro82\] that $`H^n(G,^{})=infl(H^n(H,^{}))infl(H^n(K,^{}))`$ (written multiplicatively), and in particular, for any normalized 3-cocycle $`\omega `$ of $`G`$, $`\omega `$ is cohomologous to $`\zeta =(infl\eta )(infl\eta ^{})`$ for some normalized 3-cocycles $`\eta `$, $`\eta ^{}`$ on $`K`$, $`H`$ respectively. Hence, $`D^\omega (G)`$ and $`D^\eta (H)D^\eta ^{}(K)`$ are equivalent as quasi-triangular quasi-bialgebras by Proposition 4.1 and Remark 2.1 (ii). $`\mathrm{}`$
We state formally what obtains when $`G`$ is abelian since we will need this case.
###### Proposition 4.3
Let $`G`$ be a finite abelian group with decomposition $`G=P_1\times \mathrm{}\times P_k`$ into the product of its Sylow subgroups. Let $`\omega Z^3(G,^{})`$. Then there are cocycles $`\eta _iZ^3(P_i,^{})`$ such that $`D^\omega (G)`$ and $`D^{\eta _1}(P_1)\mathrm{}D^{\eta _k}(P_k)`$ are equivalent as quasi-triangular quasi-bialgebras.
Because of this result, many of the questions we address regarding $`D^\omega (G)`$ ($`G`$ abelian) can be reduced to the case that $`G`$ is a $`p`$-group for some prime $`p`$.
## 5 Eilenberg-Maclane Cohomology
Let $`G`$ be a group. The bar complex of $`G`$, denoted by $`A_0(G)`$, is the chain complex of abelian group
$$\mathrm{}\stackrel{}{}C_3\stackrel{}{}C_2\stackrel{}{}C_1$$
where $`C_n`$ is the free abelian group generated by all $`n`$-tuples $`(x_1,\mathrm{},x_n)`$ of elements $`x_i`$ of $`G`$ and $``$ is a $``$-linear map defined by
$$(x_1,\mathrm{},x_n)=(x_2,\mathrm{},x_n)+\underset{i=1}{\overset{n1}{}}(1)^i(x_1,\mathrm{},x_{i1},x_ix_{i+1},\mathrm{},x_n)+(1)^n(x_1,\mathrm{},x_{n1}).$$
We will call $`(x_1,\mathrm{},x_n)C_n`$ a $`n`$-dimensional cell of $`A_0(G)`$. For any abelian group $`\mathrm{\Pi }`$, $`\text{Hom}(A_0(G),\mathrm{\Pi })`$ is a cochain complex. We will denote by $`C^n(G,\mathrm{\Pi })`$, $`Z^n(G,\mathrm{\Pi })`$, $`B^n(G,\mathrm{\Pi })`$ and $`H^n(G,\mathrm{\Pi })`$ the dimension $`n`$-cochains, cocycles, coboundaries and cohomology classes of $`\text{Hom}(A_0(G),\mathrm{\Pi })`$ respectively.
For any two cells $`(x_1,\mathrm{},x_n)`$ and $`(y_1,\mathrm{},y_m)`$ in $`A_0(G)`$, we can define a “shuffle” of the $`n`$ letters $`x_1,\mathrm{},x_n`$ through the $`m`$ letters $`y_1,\mathrm{},y_m`$ to be the $`n+m`$-tuple in which the order of the $`x`$’s and the order of $`y`$’s are preserved. The sign of the shuffle is the sign of the permutation required to bring the shuffled letters back to the standard shuffle $`(x_1,\mathrm{},x_n,y_1,\mathrm{},y_m)`$. Then, we can define the “star” product of $`(x_1,\mathrm{},x_n)(y_1,\mathrm{},y_m)`$ to be the signed sum of the shuffles of the letters $`x`$ through the letters $`y`$.
Let $`G`$ be an abelian group. In the paper \[Mac52\], the author described a complex $`A_1(G)`$ in which the cells are symbols $`\sigma =[\alpha _1|\alpha _2|\mathrm{}|\alpha _p]`$, with each $`\alpha _i`$ a cell of $`A_0(G)`$. The dimension of $`\sigma `$ is $`p1`$ plus the sum of the dimensions of the $`\alpha _i`$, and the boundary of $`\sigma `$ is
$$\sigma =\underset{i=1}{\overset{p}{}}(1)^{ϵ_{i1}}[\alpha _1|\mathrm{}|\alpha _i|\mathrm{}|\alpha _p]+\underset{i=1}{\overset{p}{}}(1)^{ϵ_i}[\alpha _1|\mathrm{}|\alpha _i\alpha _{i+1}|\mathrm{}|\alpha _p],$$
where $`ϵ_i=1+dim[a_1|\mathrm{}|\alpha _i]`$. For any abelian $`\mathrm{\Pi }`$, $`\text{Hom}(A_1(G),\mathrm{\Pi })`$ is a cochain complex. We denote by $`C_{ab}^n(G,\mathrm{\Pi })`$, $`Z_{ab}^n(G,\mathrm{\Pi })`$, $`B_{ab}^n(G,\mathrm{\Pi })`$ and $`H_{ab}^n(G,\mathrm{\Pi })`$ the dimension $`n`$-cochains, cocycles, coboundaries and cohomology classes of $`\text{Hom}(A_1(G),\mathrm{\Pi })`$ respectively. Until the end of this section, we will tacitly assume that $`\mathrm{\Pi }`$ is a divisible abelian. The underlying product in $`\mathrm{\Pi }`$ will be written additively.
Notice that $`A_0(G)`$ is a subcomplex of $`A_1(G)`$. We have the exact sequence of complexes
$$0A_0(G)A_1(G)B(G)0$$
with $`B(G)`$ the quotient complex $`A_1(G)/A_0(G)`$. Then
$$0\text{Hom}(B(G),\mathrm{\Pi })\text{Hom}(A_1(G),\mathrm{\Pi })\text{Hom}(A_0(G),\mathrm{\Pi })0.$$
is exact and we have the long exact sequence
$$\mathrm{}\stackrel{\delta }{}H^{}(B(G),\mathrm{\Pi })H_{ab}^{}(G,\mathrm{\Pi })H^{}(G,\mathrm{\Pi })\stackrel{\delta }{}H^{+1}(B(G),\mathrm{\Pi })H_{ab}^{+1}(G,\mathrm{\Pi })\mathrm{}$$
where $`H^n(B(G),\mathrm{\Pi })`$ is the $`n`$-th cohomology group of the cochain complex $`\text{Hom}(B(G),\mathrm{\Pi })`$.
The low dimensional cells in $`B(G)`$ are given in the following table :
| dimension | cells |
| --- | --- |
| 1 | none |
| 2 | none |
| 3 | $`[x|y]`$ |
| 4 | $`[x,y|z],[x|y,z]`$ |
| 5 | $`[x|y|z],[x,y,z|u],[u|x,y,z],[x,y|u,v]`$ |
Since $`H_{ab}^2(G,\mathrm{\Pi })=\text{Ext}_{}^1(G,\mathrm{\Pi })=0`$, we have the exact sequence
$$0H^2(G,\mathrm{\Pi })\stackrel{\delta }{}H^3(B(G),\mathrm{\Pi })H_{ab}^3(G,\mathrm{\Pi })H^3(G,\mathrm{\Pi })\stackrel{\delta }{}H^4(B(G),\mathrm{\Pi }).$$
(17)
We are going to analyze the sequence (17) in some detail. Let us denote by $`B_n(G)`$ the group of $`n`$-chains of $`B(G)`$. For any cocycle $`b:B^3(G)\mathrm{\Pi }`$ and $`x,y,zG`$,
$`b[(x,y)|z]`$ $`=`$ $`0`$ (18)
$`b[x|(y,z)]`$ $`=`$ $`0.`$ (19)
Hence, $`b`$ is a bicharacter of $`G`$. Conversely, it is also easy to see that any bicharacter on $`G`$ defines an element in $`H^3(B(G),\mathrm{\Pi })`$ uniquely. Hence, we have the following lemma.
###### Lemma 5.1
$`H^3(B(G),\mathrm{\Pi })=\text{Hom}(GG,\mathrm{\Pi })`$.
Let $`cZ^2(G,\mathrm{\Pi })`$ and $`[c]`$ the corresponding cohomology class. There exists a natural isomorphism $`\theta :H^2(G,\mathrm{\Pi })\text{Hom}(^2G,\mathrm{\Pi })`$ defined by
$$\theta ([c])(x,y)=c(x,y)c(y,x)$$
for any $`x,yG`$(cf.\[Bro82\]). On the other hand, the connecting homomorphism
$`\delta :H^2(G,\mathrm{\Pi })H^3(B(G),\mathrm{\Pi })`$ is given by
$$\delta ([c])[x|y]=c((x)(y))$$
for $`x,yG`$. One can easily check that the following diagram commutes :
$$\begin{array}{ccc}H^2(G,\mathrm{\Pi })& \stackrel{\delta }{}& H^3(B(G),\mathrm{\Pi })\\ \theta & & & & \\ \text{Hom}(^2G,\mathrm{\Pi })& & \text{Hom}(GG,\mathrm{\Pi }).\end{array}$$
(20)
The bottom horizontal map in the diagram is the natural embedding. The exact sequence (17) is thus equivalent to the following exact sequence :
$$0\text{Hom}(^2G,\mathrm{\Pi })\text{Hom}(GG,\mathrm{\Pi })H_{ab}^3(G,\mathrm{\Pi })H^3(G,\mathrm{\Pi })\stackrel{\delta }{}H^4(B(G),\mathrm{\Pi }).$$
(21)
The 3-dimensional Eilenberg-Maclane cohomology group $`H_{ab}^3(G,\mathrm{\Pi })`$ was explicitly described in \[Mac52\]. The elements in $`Z_{ab}^3(G,\mathrm{\Pi })`$ are pairs $`(f,d)`$ where $`fZ^3(G,\mathrm{\Pi })`$ and $`dC^2(G,\mathrm{\Pi })`$ satisfying the following conditions :
$`d(xy|z)d(x|z)d(y|z)+f(x,y,z)f(x,z,y)+f(z,x,y)`$ $`=`$ $`0`$ (22)
$`d(x|yz)d(x|y)d(x|z)f(x,y,z)+f(y,x,z)f(y,z,x)`$ $`=`$ $`0.`$ (23)
The 3-dimensional cocycle $`(f,d)`$ is a coboundary if there exists $`hC^2(G,\mathrm{\Pi })`$ such that
$`f(x,y,z)`$ $`=`$ $`\delta h(x,y,z)`$ (24)
$`d(x|y)`$ $`=`$ $`h(x,y)h(y,x).`$ (25)
Thus, the maps $`\text{Hom}(GG,\mathrm{\Pi })H_{ab}^3(G,\mathrm{\Pi })`$ and $`H_{ab}^3(G,\mathrm{\Pi })H^3(G,\mathrm{\Pi })`$ in (21) are given by
$`b[(0,b)]`$ and $`[(f,d)][f]`$ (26)
respectively, where $`[x]`$ means the cohomology class associated to the cocycle $`x`$.
To any $`(f,d)Z_{ab}^3(G,\mathrm{\Pi })`$, one can assign the function $`t(x)=d(x|x)`$, called its trace. Any trace is a quadratic function–a function $`t:G\mathrm{\Pi }`$ such that
1. $`t(ax)=a^2t(x)`$ for $`a`$, and
2. $`b_t(x,y)=t(x+y)t(x)t(y)`$ defines a bilinear function on $`G`$.
###### Theorem 5.2 (Eilenberg-MacLane)
Let $`Q(G,\mathrm{\Pi })`$ consist of all quadratic functions from $`G`$ to $`\mathrm{\Pi }`$. The function assigning to each cocycle its trace induces an isomorphism
$$H_{ab}^3(G,\mathrm{\Pi })\stackrel{}{}Q(G,\mathrm{\Pi }).$$
A map $`F:B^4(G)\mathrm{\Pi }`$ is a cocycle if it satisfies the following:
$`F[(x,y,z)|u]`$ $`=`$ $`0`$ (27)
$`F[u|(x,y,z)]`$ $`=`$ $`0`$ (28)
$`F[(x,y)|u,v]`$ $`=`$ $`F[x,y|(u,v)]`$ (29)
$`F[x|(y)(z)]`$ $`=`$ $`F[(x)(y)|z]`$ (30)
$`F`$ is a coboundary if there is a function $`d:B^3(G)\mathrm{\Pi }`$ such that
$`d[(x,y)|z]`$ $`=`$ $`F[x,y|z]`$ (31)
$`d[x|(y,z)]`$ $`=`$ $`F[x|y,z].`$ (32)
Let $`fZ^3(G,\mathrm{\Pi })`$. The connecting homomorphism $`\delta :H^3(G,\mathrm{\Pi })H^4(B(G),\mathrm{\Pi })`$ in (21) is given by $`\delta [f]=[F]`$ where $`F:B^4(G)\mathrm{\Pi }`$ is defined by
$`F[x,y|z]`$ $`=`$ $`f((x,y)(z))`$ (34)
$`F[x|y,z]`$ $`=`$ $`f((x)(y,z))`$ (35)
Let $`F:B^4(G)\mathrm{\Pi }`$ be a cocycle. Define
$$\overline{F}(x,y,z)=F[x,y|z]F[y,x|z]$$
(36)
for $`x,y,zG`$.
###### Lemma 5.3
The map $`\mathrm{\Psi }:[F]\overline{F}`$ defines a group homomorphism from $`H^4(B(G),\mathrm{\Pi })`$ onto $`\text{Hom}(^3G,\mathrm{\Pi })`$. The map splits if $`G`$ is finitely generated.
Proof. Since $`(x,y)=(y,x)`$, one can check that $`\mathrm{\Psi }`$ is a well-defined. Let $`F:B^4(G)\mathrm{\Pi }`$ be a cocycle. Then, for any $`x,y,u,vG`$,
$`\overline{F}((u,v),x,y)`$ $`=`$ $`F[(u,v),x|y]F[x,(u,v)|y]`$
$`=`$ $`F[(u,v)|x,y]+F[(u,v)|y,x]`$
$`=`$ $`F[u,v|(x,y)]+F[u,v|(y,x)]`$
$`=`$ $`0.`$
Therefore, $`\overline{F}`$ is linear at the first entry. Clearly,
$$\overline{F}(x,y,z)=\overline{F}(y,x,z)\text{ and }\overline{F}(x,x,z)=0.$$
In particular, $`\overline{F}`$ is also linear at the second entry. Notice that
$`\overline{F}(x,y,z)`$ $`=`$ $`F[x,y|z]F[y,x|z]`$
$`=`$ $`F[x|z,y]F[x|y,z]`$
$`=`$ $`(F[x,z|y]F[z,x|y])`$
$`=`$ $`\overline{F}(x,z,y).`$
Therefore, $`\overline{F}`$ is also linear at the third entry and hence $`\overline{F}\text{Hom}(^3G,\mathrm{\Pi })`$.
Note that the map $`\psi ^{}:H^3(G,\mathrm{\Pi })\text{Hom}(^3G,\mathrm{\Pi })`$ given by
$$\psi ^{}[\omega ](x,y,z)=\underset{\sigma }{}\text{sgn}(\sigma )\omega (\sigma (x),\sigma (y),\sigma (z))$$
(37)
with $`\sigma `$ running through the permutations of $`x,y,z`$ is a surjective homomorphism (cf. \[Bro82\]). It follows directly from (34) and (35) that the diagram
commutes. Hence, $`\mathrm{\Psi }`$ is surjective. Moreover, if $`G`$ is finitely generated, the map $`\psi ^{}`$ has a linear section $`s`$. Then, $`\delta s`$ is a linear section of $`\mathrm{\Psi }`$. $`\mathrm{}`$
For any $`\beta Z_{ab}^2(G,\text{Hom}(G,\mathrm{\Pi }))`$, define $`F_\beta :B^4(G)\mathrm{\Pi }`$ by
$$F_\beta [x,y|z]=0\text{and}F_\beta [x|y,z]=\beta (y,z)(x)$$
(38)
for any $`x,y,zG`$. One can easily check that $`F_\beta Z^4(B(G),\mathrm{\Pi })`$. Moreover, $`\mathrm{\Xi }:[\beta ][F_\beta ]`$ defines a linear map from $`H_{ab}^2(G,\text{Hom}(G,\mathrm{\Pi }))`$ to $`H^4(B(G),\mathrm{\Pi })`$.
###### Lemma 5.4
The sequence
is exact. In particular, if $`G`$ is finitely generated, the sequence is split exact.
Proof. If $`[F_\beta ]=0`$, then there exists a function $`d:B_3(G)\mathrm{\Pi }`$ such that
$`0`$ $`=`$ $`d[(x,y)|z]`$
$`\beta (y,z)(x)`$ $`=`$ $`d[x|(y,z)].`$
Hence, $`dC^1(G,\text{Hom}(G,\mathrm{\Pi }))`$ and $`\beta =\delta d`$. Therefore, $`\mathrm{\Xi }`$ is injective.
Obviously, $`\mathrm{\Psi }\mathrm{\Xi }[\beta ]=0`$ for $`\beta Z_{ab}^2(G,\text{Hom}(G,\mathrm{\Pi }))`$. Let $`[F]\text{ker }\mathrm{\Psi }`$. Then, we have
$$F[x,y|z]=F[y,x|z]$$
for any $`x,y,zG`$. Since $`\mathrm{\Pi }`$ is divisible, there exists a map $`\tau C^2(G,\mathrm{\Pi })`$ such that
$$\tau [(x,y)|z]=F[x,y|z].$$
Then the map $`F^{}:B^4(G)\mathrm{\Pi }`$ defined by
$$F^{}[x,y|z]=\tau [(x,y)|z]\text{and}F^{}[x|y,z]=\tau [x|(y,z)]$$
is a coboundary and hence $`[F]=[FF^{}]`$. Let $`\beta (y,z)(x)=F[x|y,z]F^{}[x|y,z]`$. Then,
$`\beta (y,z)(u)\beta (y,z)(uv)+\beta (y,z)(v)`$ $`=`$ $`F[(u,v)|y,z]\tau ((u,v)|(y,z))`$
$`=`$ $`F[u,v|(y,z)]\tau ((u,v)|(y,z))`$
$`=`$ $`\tau ((u,v)|(y,z))\tau ((u,v)|(y,z))=0.`$
Moreover, $`\beta ((x,y,z))(u)=F[u|(x,y,z)]F^{}[u|(x,y,z)]=0`$. Therefore, $`\beta Z_{ab}^2(G,\text{Hom}(G,\mathrm{\Pi }))`$ and $`F_\beta =FF^{}`$. Thus, $`\mathrm{\Xi }[\beta ]=[F]`$. The second statement follows directly from Lemma 5.3. $`\mathrm{}`$
## 6 The Kernel Of $`\mathrm{\Lambda }`$
Let us return to the exact sequence (21) with $`\mathrm{\Pi }=^{}`$. We then have the exact sequence
$$0\text{Hom}(^2G,^{})\text{Hom}(GG,^{})H_{ab}^3(G,^{})H^3(G,^{})\stackrel{\delta }{}H^4(B(G),^{}).$$
(39)
Notice that for any $`(f,d)Z_{ab}^3(G,^{})`$, $`fZ^3(G,^{})`$ and $`f_z`$ is a coboundary for $`zG`$ because
$$f_z(x,y)=\frac{f(z,x,y)f(x,y,z)}{f(x,z,y)}=\frac{d(z|x)d(z|y)}{d(z|xy)}.$$
Therefore, $`[f]H^3(G,^{})_{ab}`$ and hence the image of the map $`H_{ab}^3(G,^{})H^3(G,^{})`$ is contained in $`H^3(G,^{})_{ab}`$.
Let $`fZ^3(G,^{})_{ab}`$. Then $`\delta [f]`$ admits the representative $`FZ^4(B(G),^{})`$ defined by
$$F[x,y|z]=f((x,y)(z))^1\text{and}F[x|y,z]=f((x)(y,z)).$$
On the other hand, $`\mathrm{\Xi }\mathrm{\Lambda }[f]`$ has the representative $`F^{}Z^4(B(G),^{})`$ given by
$$F^{}[x,y|z]=1\text{and}F^{}[x|y,z]=\frac{\tau _y(x)\tau _z(x)}{\tau _{yz}(x)}f_x(y,z).$$
where $`f_x(y,z)=\frac{\tau _x(y)\tau _x(z)}{\tau _x(yz)}`$. Let
$$d[x|y]=\tau _y(x).$$
Then $`F^{}/F=\delta d`$. Hence $`F^{}`$ and $`F`$ represent the same cohomology class in $`H^4(B(G),^{})`$ and the diagram
$$\begin{array}{ccc}H^3(G,^{})& \stackrel{\delta }{}& H^4(B(G),^{})\\ incl& & \mathrm{\Xi }& & \\ H^3(G,^{})_{ab}& \stackrel{\mathrm{\Lambda }}{}& H_{ab}^2(G,\widehat{G})\end{array}$$
commutes. Moreover, we obtain an exact sequence
(40)
Therefore, $`\text{ker }\mathrm{\Lambda }=\text{coker }\left(\text{Hom}(GG,^{})H_{ab}^3(G,^{})\right)`$. Since the map $`\text{Hom}(GG,^{})H_{ab}^3(G,^{})`$ is given by $`b[(1,b)]`$,
$$\text{Im }\left(\text{Hom}(GG,^{})H_{ab}^3(G,^{})\right)=\left\{[(1,b)]H_{ab}^3(G,^{})\right|b\text{ is a bicharacter}\}.$$
Let $`K(G,^{})`$ be the group of quadratic forms which are the trace of some bicharacter. Then, by Theorem 5.2, we see that the following holds:
###### Proposition 6.1
For any abelian group $`G`$, $`\text{ker }\mathrm{\Lambda }=Q(G,^{})/K(G,^{})`$. $`\mathrm{}`$
###### Lemma 6.2
Let $`G`$ be a finite abelian group.
1. Let $`G=HL`$ for subgroups $`H,L`$ of $`G`$. Then $`tK(G,^{})`$ if, and only if, $`t|_HK(H,^{})`$ and $`t|_LK(L,^{})`$.
2. Let $`H`$ be the Sylow 2-subgroup of $`G`$. Then
$$Q(G,^{})/K(G,^{})Q(H,^{})/K(H,^{})$$
as abelian groups.
3. $`Q(G,^{})^2K(G,^{})`$.
Proof. (i) The necessity of the statement is obvious. Let us assume $`tK(G,^{})`$ is such that $`t|_HK(H,^{})`$ and $`t|_LK(L,^{})`$. Let $`b_H`$ and $`b_L`$ be bicharacters of $`H`$ and $`L`$ respectively with
$$t|_H(x)=b_H(x,x)\text{and}t|_L(y)=b_L(y,y)$$
for $`xH`$ and $`yL`$. Since $`tQ(G,^{})`$,
$$c(u,v)=\frac{t(uv)}{t(u)t(v)}$$
defines a bicharacter on $`G`$. Consider the map $`b:G\times G^{}`$ given by
$$b(x_1y_1,x_2y_2)=b_H(x_1,x_2)c(x_1,y_2)b_L(y_1,y_2)$$
for any $`x_1,x_2H`$ and $`y_1,y_2L`$. One can check directly that $`b`$ is a bicharacter on $`G`$. Moreover,
$$b(xy,xy)=b_H(x,x)c(x,y)b_L(y,y)=t(xy).$$
Therefore, $`tK(G,^{})`$.
(ii) It suffices to show that $`tK(G,^{})`$ if, and only if, $`t|_HK(H,^{})`$. Let $`L`$ be the subgroup of odd order elements of $`G`$. Then $`G=HL`$. Since $`|L|`$ is odd, $`K(L,^{})=Q(L,^{})`$. Hence, by (i), $`tQ(G,^{})`$ if, and only if, $`t|_HQ(H,^{})`$.
(iii) Since
$$c(x,y)=\frac{t(xy)}{t(x)t(y)}$$
is a bicharacter and $`c(x,x)=t(x)^2`$, the result follows. $`\mathrm{}`$
###### Proposition 6.3
Let $`G`$ be a finite abelian group. Then $`Q(G,^{})/K(G,^{})\mathrm{\Omega }_2(G)`$.In particular, $`\text{ker }\mathrm{\Lambda }_G\mathrm{\Omega }_2(G)`$.
Proof. By virtue of Lemma 6.2(ii), we may assume that $`G`$ is a finite abelian 2-group. Let us consider the case that $`G`$ is cyclic of order $`2^n`$ with the generator $`y`$. Then
$$1=t(1)=t(y^{2^n})=t(y)^{2^{2n}}$$
Since $`y^1=y^{2^n1}`$,
$$t(y)=t(y^{2^n1})=t(y)^{2^{2n}2^{n+1}+1}.$$
Hence, $`t(y)`$ is a $`2^{n+1}`$th root of unity. Conversely, for any $`2^{n+1}`$th root of unity $`\xi `$, $`t(y^r)=\xi ^{r^2}`$ defines a quadratic function. Hence $`Q(G,^{})_{2^{n+1}}`$. Obviously, $`K(G,^{})`$ is a proper subgroup $`Q(G,^{})`$ and $`Q(G,^{})^2K(G,^{})`$ by Lemma 6.2(iii). Since $`Q(G,^{})^2`$ is the largest proper subgroup of $`Q(G,^{})`$, $`Q(G,^{})^2=K(Q,^{})`$. Hence, $`Q(G,^{})/K(G,^{})_2\mathrm{\Omega }_2(G)`$.
Now, we consider the general case. Let $`G=C_1\mathrm{}C_l`$ where $`C_1,\mathrm{},C_l`$ are cyclic subgroups of $`G`$. Consider the map $`p:Q(G,^{})Q(C_1,^{})/K(C_1,^{})\times \mathrm{}\times Q(C_l,^{})/K(C_l,^{})`$ defined by
$$p(t)=(t_1,\mathrm{},t_l)$$
where $`t_i`$ is the coset $`t|_{C_i}K(C_i,^{})`$. Obviously, $`p`$ is an epimorphism. By Lemma 6.2(i), $`\text{ker }p=K(G,^{})`$. Hence,
$$\frac{Q(G,^{})}{K(G,^{})}\frac{Q(C_1,^{})}{K(C_1,^{})}\times \mathrm{}\times \frac{Q(C_l,^{})}{K(C_l,^{})}.$$
As $`{\displaystyle \frac{Q(C_i,^{})}{K(C_i,^{})}}_2`$, $`{\displaystyle \frac{Q(G,^{})}{K(G,^{})}}_2^l\mathrm{\Omega }_2(G)`$. $`\mathrm{}`$
###### Corollary 6.4
If $`G`$ is a finite abelian group of odd order, then $`\mathrm{\Lambda }`$ is injective. $`\mathrm{}`$
###### Corollary 6.5
Let $`G`$ be a finite abelian group and $`H`$ the Sylow 2-subgroup of $`G`$. Let $`H=_{i=1}^lC_i`$ be a cyclic subgroup decomposition of $`H`$. Let $`p_i:HC_i`$ be the natural projection associated to the decomposition and $`p_i^{}:H^3(C_i,^{})H^3(G,^{})`$ the associated inflation of $`p_i`$. Then
$$\text{ker }\mathrm{\Lambda }_G=\underset{i=1}{\overset{l}{}}p_i^{}(\text{ker }\mathrm{\Lambda }_{C_i})$$
Proof. Note that the exact sequence (40) is natural in $`G`$. Therefore, we have the commutative diagram
with exact rows. Therefore, $`p_i^{}(\text{ker }\mathrm{\Lambda }_{C_i})\text{ker }\mathrm{\Lambda }_G`$ for $`i=1,\mathrm{},l`$. Hence, $`_{i=1}^lp_i^{}(\text{ker }\mathrm{\Lambda }_{C_i})\text{ker }\mathrm{\Lambda }_G`$. Since $`p_i^{}`$ is injective and $`\{\text{Im }p_i^{}\}`$ is $``$-linearly independent in $`H^3(G,^{})_{ab}`$,
$$\underset{i=1}{\overset{l}{}}p_i^{}(\text{ker }\mathrm{\Lambda }_{C_i})\underset{i=1}{\overset{l}{}}\text{ker }\mathrm{\Lambda }_{C_i}_2^l.$$
Therefore, the equality in the statement follows. $`\mathrm{}`$
## 7 The Cohomology Group $`H^3(G,^{})_{ab}`$
###### Definition 7.1
Let $`G`$ be an abelian group and $`M`$ be a left $`G`$-module. For $`gG`$, the group homomorphism $`m_g:C_1(G)C_{}(G)`$ given by $`yyg`$ is a chain map. We denote by $`D_{g,M}^{}:H^{}(G,M)H^1(G,M)`$ the map induced by $`m_g`$ and define $`D_{g,M}^0`$ to be the trivial map. We will also denote $`H^{}(G,M)_{ab}=_{gG}\text{ker }D_{g,M}^{}`$. When $`M`$ is the trivial module $`^{}`$, $`D_{g,M}`$ will simply written as $`D_g`$. In this case, one can easily see that $`H^3(G,M)_{ab}`$ coincides with $`H^3(G,^{})_{ab}`$ defined previously.
###### Lemma 7.2
For any positive integer $`n1`$,
$$H^n(G,^{})_{ab}H^{n+1}(G,)_{ab}$$
and
$$\frac{H^n(G,^{})}{H^n(G,^{})_{ab}}\frac{H^{n+1}(G,)}{H^{n+1}(G,)_{ab}}.$$
Proof. Recall that, for $`n1`$, $`H^n(G,^{})\stackrel{\delta }{}H^{n+1}(G,)`$ under the connecting maps of the long exact sequence associated to the exact sequence of trivial $`G`$-modules
Since the diagram
commutes for $`gG`$, the result follows. $`\mathrm{}`$
###### Remark 7.3
Let $`G`$ be an abelian group. For any $`gG`$, the map $`D_{g,}=_{n0}D_{g,}^n`$ is a graded derivation on $`\underset{n0}{}H^n(G,)`$ with respect to the cup product. Hence, $`\text{ker }D_{g,}`$ is a graded subring of $`H^{}(G,)`$. Since $`H^{}(G,)_{ab}=\underset{gG}{}\text{ker }D_{g,}`$, $`H^{}(G,)_{ab}`$ is a graded subring of $`H^{}(G,)`$. In particular, if $`G`$ is a finite cyclic group, $`H^{}(G,)`$ is generated by $`H^2(G,)`$. As $`H^1(G,)=0`$, $`H^2(G,)=H^2(G,)_{ab}`$. Therefore, $`H^{}(G,)_{ab}=H^{}(G,)`$ in case $`G`$ is cyclic.
###### Lemma 7.4
For any abelian group $`G`$, $`H^2(G,^{})_{ab}`$ is trivial and
$$\frac{H^3(G,^{})}{H^3(G,^{})_{ab}}\text{Hom}(^3G,^{}).$$
Proof. For $`[f]H^2(G,^{})_{ab}`$, $`f(x,y)/f(y,x)=1`$. Hence, $`f`$ is a coboundary. Recall that the map $`\phi ^{}:H^3(G,^{})\text{Hom}(^3G,^{})`$ defined by equation (37) is a split epimorphism. Moreover, for any normalized 3-cocycle $`\omega `$,
$$\phi ([\omega ])(x,y,z)=\omega _x(y,z)/\omega _x(z,y)$$
for any $`x,y,zG`$. Therefore, $`\text{ker }\phi ^{}=H^3(G,^{})_{ab}`$. $`\mathrm{}`$
Let $`G`$ be a finite abelian group with a cyclic subgroup decomposition $`G=\underset{i=1}{\overset{l}{}}C_i`$ . Let $`\{p_i\}`$ be the projections of $`G`$ associated to the decomposition and $`P_i^{}:H^{}(C_i,)H^{}(G,)`$ the induced graded ring homomorphism with respect to cup product. Let $`P_{C_1,\mathrm{},C_l;G}^{}:H^{}(C_1,)\mathrm{}H^{}(C_l,)H^{}(G,)`$ be the graded ring monomorphism defined by
$$P_{C_1,\mathrm{},C_l;G}^{}(f_1\mathrm{}f_l)=P_1(f_1)\mathrm{}P_l(f_l).$$
###### Proposition 7.5
Let $`G`$ be a finite abelian group with a cyclic subgroup decomposition $`G=\underset{i=1}{\overset{l}{}}C_i`$ . Then
$$H^4(G,)_{ab}=\text{Im }P_{C_1,\mathrm{},C_l;G}^4.$$
In particular,
$$H^4(G,)_{ab}=\underset{n_1+\mathrm{}+n_l=4}{}H^{n_1}(C_1,)\mathrm{}H^{n_l}(C_l,).$$
Proof. We proceed by induction on $`l`$. The case $`l=1`$ follows from Remark 7.3. Assume $`l>1`$ and $`H=\underset{i=2}{\overset{l}{}}C_i`$. Then $`\text{Im }P_{C_1,\mathrm{},C_l;G}^{}\text{Im }P_{C_1,H;G}^{}`$. Moreover, by the K$`\ddot{\text{u}}`$nneth theorem,
$$\frac{H^4(G,)}{\text{Im }P_{C_1,\mathrm{},C_l;G}^4}\frac{\text{Im }P_{C_1,H;G}^4}{\text{Im }P_{C_1,\mathrm{},C_l;G}^4}\text{Tor}(H^3(C_1,),H^2(H,)).$$
Note that
$$\frac{\text{Im }P_{C_1,H;G}^4}{\text{Im }P_{C_1,\mathrm{},C_l;G}^4}\frac{H^4(G,)}{P_{C_2,\mathrm{},C_l;H}^4}.$$
By induction and Lemma 7.4,
$$\frac{\text{Im }P_{C_1,H;G}^4}{\text{Im }P_{C_1,\mathrm{},C_l;G}^4}\frac{H^4(G,)}{H^4(G,)_{ab}}^3G.$$
Since
$$\text{Tor}(H^2(C_1,),H^3(H,))C_2H^3(H,)C_2H^2(H,^{})C_2^2H,$$
$$\frac{H^4(G,)}{\text{Im }P_{C_1,\mathrm{},C_l;G}^4}\left(C_2^2H\right)^3H^3G.$$
Hence, by Lemma 7.4 and Lemma 7.2,
$$\frac{H^4(G,)}{\text{Im }P_{C_1,\mathrm{},C_l;G}^4}\frac{H^4(G,)}{H^4(G,)_{ab}}.$$
Thus, $`|\text{Im }P_{C_1,\mathrm{},C_l;G}^4|=|H^4(G,)_{ab}|`$. Notice that $`\text{Im }P_{C_1,\mathrm{},C_l;G}^{}`$ is the graded subring of $`H^{}(G,)`$ generated by $`H^2(G,)`$. As $`H^2(G,)=H^2(G,)_{ab}`$, $`\text{Im }P_{C_1,\mathrm{},C_l;G}^{}H^{}(G,)_{ab}`$. $`\mathrm{}`$
## 8 The Image of $`\mathrm{\Lambda }_G`$
Let $`G`$ be a finite abelian group and
$$E:1\widehat{G}\mathrm{\Gamma }G1$$
an abelian central extension of $`G`$ by $`\widehat{G}`$. By applying the functor $`\widehat{\mathrm{?}}=\text{Hom}(\mathrm{?},^{})`$ to this exact sequence, we can obtain another abelian central extension of $`G`$ by $`\widehat{G}`$, namely
$$ϵ(E):1\widehat{G}\widehat{\mathrm{\Gamma }}G1$$
where $`G`$ and $`\widehat{\widehat{G}}`$ are naturally identified. Obviously, if $`E_1`$ and $`E_2`$ are equivalent abelian central extensions of $`G`$ by $`\widehat{G}`$, and so are $`ϵ(E_1)`$ and $`ϵ(E_2)`$. Therefore, $`ϵ`$ induces a map on $`H_{ab}^2(G,\widehat{G})`$. We will denote this map by the same symbol $`ϵ`$.
Let $`\beta `$ be a normalized abelian cocycle in $`Z_{ab}^2(G,\widehat{G})`$. Then the set $`\mathrm{\Gamma }=\widehat{G}\times G`$ equipped with the multiplication
$$(\alpha ,x)(\lambda ,y)=(\alpha \lambda \beta (x,y),xy)$$
(41)
is an abelian group. Moreover, a central extension of $`G`$ by $`\widehat{G}`$ associated to $`[\beta ]`$ is given by
(42)
where $`i(\alpha )=(\alpha ,1)`$ and $`p(\alpha ,x)=x`$ for $`\alpha \widehat{G}`$ and $`xG`$. Let $`\tau _x`$ be a normalized 1-cochain such that
$$\delta \tau _x(g,h)=\beta (g,h)(x)$$
(43)
for any $`g,hG`$. For any $`xG`$, denoted by $`\overline{x}`$ be a fixed element in $`\widehat{\mathrm{\Gamma }}`$ such that $`\widehat{i}(\overline{x})=x`$ and $`\overline{1}`$ the identity element of $`\widehat{\mathrm{\Gamma }}`$. Then, there exist a normalized cocycle $`\beta ^{}Z^2(G,\widehat{G})`$ such that
$$\overline{x}\overline{y}=\widehat{p}(\beta ^{}(x,y))\overline{x}\overline{y}$$
for any $`x,yG`$. Obviously, $`ϵ[\beta ]=[\beta ^{}]`$.
Let $`\tau _x^{}(g)=\overline{x}(1,g)`$ for any $`g,xG`$. Then
$$\beta ^{}(x,y)(g)=\frac{\tau _x^{}(g)\tau _y^{}(g)}{\tau _{xy}^{}(g)}$$
for $`x,y,gG`$. Since $`\overline{x}\widehat{\mathrm{\Gamma }}`$,
$$\tau _x^{}(g)\tau _x^{}(h)=\tau _x^{}(gh)\beta (g,h)(x).$$
Therefore, $`\tau _x^{}=\tau _x\lambda _x`$ for some $`\lambda _x\widehat{G}`$. Define
$$b_1^{}(x,y)(g)=\frac{\tau _x(g)\tau _y(g)}{\tau _{xy}(g)}\text{and}\lambda (x)(g)=\lambda _x(g).$$
(44)
Then $`\beta ^{}=\beta _1^{}\delta \lambda `$ and hence $`\beta _1^{}`$ is a normalized 2-cocycle in $`Z_{ab}^2(G,\widehat{G})`$ and
$$ϵ([\beta ])=[\beta _1^{}].$$
(45)
By the formulae (44) and (45), we obtain
###### Proposition 8.1
The map $`ϵ:H_{ab}^2(G,\widehat{G})H_{ab}^2(G,\widehat{G})`$ is a group homomorphism. $`\mathrm{}`$
Let us identify $`\widehat{G}`$ and $`H^2(G,)`$ and consider the map $`\mathrm{{\rm Y}}:\widehat{G}\widehat{G}H_{ab}^2(G,\widehat{G})`$, which assigns to $`\alpha \lambda `$ the cohomology class $`[\beta ]`$ with
$$\beta (x,y)(g)=\alpha (g)^{\stackrel{~}{\lambda }(x,y)}$$
where $`\stackrel{~}{\lambda }`$ is a normalized 2-cocycle in $`Z^2(G,)`$ corresponding to $`\lambda `$. Since $`G`$ is a finite, one can easily see that $`\mathrm{{\rm Y}}`$ is an isomorphism.
###### Lemma 8.2
Let $`T:\widehat{G}\widehat{G}\widehat{G}\widehat{G}`$ be the transposition automorphism. Then,
$$\mathrm{{\rm Y}}T=ϵ\mathrm{{\rm Y}}.$$
Proof. Let $`\alpha _1,\alpha _2\widehat{G}`$ and $`\stackrel{~}{\alpha _i}=\delta f_i`$ for some $`f_i:G`$ such that $`\text{Im }\delta f_i`$ and $`f_i(1)=0`$. Then $`\mathrm{{\rm Y}}(\alpha _1\alpha _2)`$ is the cohomology class represented by $`\beta `$ given by
$$\beta (x,y)(g)=\alpha _1(g)^{f_2(x)+f_2(y)f_2(xy)}.$$
Let
$$\tau _g(x)=\alpha _1(g)^{f_2(x)}\alpha _2(x)^{f_1(g)}.$$
Then $`\beta (x,y)(g)=\delta \tau _g(x,y)`$. By formulae (44) and (45), $`ϵ[\beta ]=[\beta ^{}]`$ where $`\beta ^{}`$ is given by
$$\beta ^{}(x,y)(g)=\frac{\tau _x(g)\tau _y(g)}{\tau _{xy}(g)}=\alpha _2(g)^{f_1(x)+f_1(y)f_1(xy)}.$$
Hence, $`[\beta ^{}]=\mathrm{{\rm Y}}(\alpha _2\alpha _1)`$. $`\mathrm{}`$
###### Lemma 8.3
For any normalized $`\beta Z_{ab}^2(G,\widehat{G})`$, $`\xi (\beta )(x,y,z)=\beta (x,y)(z)`$ defines a 3-cocycle in $`Z^3(G,^{})`$ and $`[\xi (\beta )]H^3(G,^{})_{ab}`$. Moreover, $`\mathrm{\Lambda }_G[\xi (\beta )]=[\beta ]ϵ[\beta ]`$.
Proof. For any normalized $`\beta Z_{ab}^2(G,\widehat{G})`$, one can easily check that $`\omega =\xi (\beta )`$ satisfies the 3-cocycle identity. For any $`zG`$, let $`\tau _z:G^{}`$ such that $`\beta (x,y)(z)=\delta \tau _z(x,y)`$. Then
$$\omega _z(x,y)=\delta \tau _z(x,y)$$
and so $`[\omega ]H^2(G,^{})_{ab}`$. Moreover, $`\mathrm{\Lambda }_G[\omega ]`$ can be represented by the 2-cocycle $`\beta _1`$ given by
$$\beta _1(x,y)(z)=\frac{\tau _z(x)\tau _z(y)}{\tau _z(xy)}\frac{\tau _x(z)\tau _y(z)}{\tau _{xy}(z)}.$$
However, from the foregoing we have $`[\beta _1]=[\beta ]ϵ[\beta ]`$. $`\mathrm{}`$
###### Theorem 8.4
Let $`G`$ be a finite abelian group. Then
$$\text{Im }\mathrm{\Lambda }_G=\text{Im }S$$
where $`S:H^2(G,\widehat{G})H^2(G,\widehat{G})`$ is defined by $`S[\beta ]=[\beta ]ϵ[\beta ]`$ for $`[\beta ]H_{ab}^2(G,\widehat{G})`$.
Proof. It follows from Lemma 8.3 that $`\text{Im }\mathrm{\Lambda }_G\text{Im }S`$. By virtue of Lemma 8.2, $`\text{Im }S\text{Im }(id+T)`$. Let $`G=C_1\mathrm{}C_l`$ be a cyclic subgroup decomposition of $`G`$. As $`\widehat{G}G`$, one can easily see that
$$\text{Im }(id+T)\left(\underset{i<j}{}C_iC_j\right)\left(\underset{i=1}{\overset{l}{}}\frac{C_i}{(C_i)_2}\right).$$
By Proposition 7.5 and Proposition 6.3
$$H^3(G,^{})_{ab}\left(\underset{i<j}{}C_iC_j\right)\left(\underset{i=1}{\overset{l}{}}C_i\right)\text{and}\text{ker }\mathrm{\Lambda }_G\mathrm{\Omega }_2(G).$$
Therefore, $`|\text{Im }\mathrm{\Lambda }_G|=|\text{Im }S|`$ and hence $`\text{Im }\mathrm{\Lambda }_G=\text{Im }S`$. $`\mathrm{}`$
###### Theorem 8.5
Let $`G`$ be a finite abelian 2-group and $`\omega `$ a normalized 3-cocycle such that $`[\omega ]H^2(G,^{})_{ab}`$. Then $`\mathrm{\Gamma }^\omega (G)`$ is a direct sum of an even number of cyclic subgroups.
Proof. By Theorem 8.4, $`\mathrm{\Lambda }_G[\omega ]=[\beta ]ϵ[\beta ]`$ for some normalized $`\beta Z_{ab}^2(G,\widehat{G})`$. For $`zG`$, let $`\tau _z:G^{}`$ such that
$$\beta (x,y)(z)=\delta \tau _z(x,y).$$
Then, $`\mathrm{\Lambda }_G[\omega ]`$ contains the normalized 2-cocycle $`\sigma `$ given by
$$\sigma (x,y)(z)=\frac{\tau _z(x)\tau _z(y)}{\tau _z(xy)}\frac{\tau _x(z)\tau _y(z)}{\tau _{xy}(z)}.$$
Consider the central extension associated to $`\sigma `$
(46)
as defined in (42) where $`\mathrm{\Gamma }`$ the group with the underlying set $`\widehat{G}\times G`$ endowed with the multiplication
$$(\alpha ,x)(\lambda ,y)=(\alpha \lambda \sigma (x,y),xy).$$
To show that $`\mathrm{\Gamma }`$ is a direct sum of an even number of cyclic subgroups, it is enough to show that $`\mathrm{\Gamma }^{}=p^1(\mathrm{\Omega }_2(G))/i(\widehat{G}^2)`$ is a direct sum of an even number of cyclic subgroups. For this, it suffices to show that $`\mathrm{\Gamma }^{}`$ admits a nonsingular alternating bilinear form valued in $`^{}`$ \[Wal63\]. Define $`b:\mathrm{\Gamma }^{}\mathrm{\Gamma }^{}^{}`$ by
$$b((\alpha \widehat{G}^2,x),(\lambda \widehat{G}^2,y))=\frac{\alpha (y)}{\lambda (x)}\frac{\tau _y(x)}{\tau _x(y)}.$$
for $`\alpha ,\lambda \widehat{G}`$ and $`x,y\mathrm{\Omega }_2(G)`$. One can easily see that $`b`$ is well-defined. The linearity of $`b`$ follows from the fact that
$$\beta (x,y)=\frac{\tau _x\tau _y}{\tau _{xy}}$$
is a character for any $`x,yG`$. The routine verification that $`b`$ is also non-degenerate and alternating will be left to the reader. $`\mathrm{}`$
###### Remark 8.6
We summarize some of our results.
We have an exact sequence
$$1\text{Hom}(^2G,^{})\text{Hom}(GG,^{})H_{ab}^3(G,^{})H^3(G,^{})_{ab}\stackrel{\mathrm{\Lambda }_G}{}\text{Im }S1$$
where $`S`$ is as in Theorem 8.4.
If $`|G|`$ is odd then $`\mathrm{\Lambda }`$ induces an isomorphism
where $`H_{ab}^2(G,\widehat{G})^+`$ is the group of $`ϵ`$-invariants. (In this case, $`\mathrm{\Lambda }`$ is an injection (Cor. 6.4) and $`\text{Im }S=H_{ab}^2(G,\widehat{G})^+`$ as follows from Theorem 8.4 and the fact that $`ϵ`$ is an involution on the group $`H_{ab}^2(G,\widehat{G})`$.)
## 9 The Monoidal Category $`D^\omega (G)\text{-}\text{Mod}`$
Let $`G`$ be a finite abelian group and $`\omega Z^3(G,^{})_{ab}`$. Adopt the notation introduced in section 3, and pick $`\tau T(\omega )`$. Then,
$$\mathrm{\Gamma }^\omega =\{\sigma (\alpha ,x)|\alpha \widehat{G},xG\}$$
where $`\sigma (\alpha ,x)`$ is given by (16). By Proposition 2.2, for any $`xG`$ and $`\alpha \widehat{G}`$, the element
$$\chi (\alpha ,x)=\underset{gG}{}\alpha (g)\tau _x(g)f_{g,x}$$
is an algebra map from $`D^\omega (G)`$ to $``$ where $`\{f_{x,g}\}`$ is the dual basis of $`\{e(g)x\}`$. Denote by $`V(\alpha ,x)`$ the irreducible representation of $`D^\omega (G)`$ associated to the algebra map $`\chi (\alpha ,x)`$. Then,
$$S=\left\{V(\alpha ,x)\right|\alpha \widehat{G},xG\}$$
is a complete set of irreducible representations of $`D^\omega (G)`$ up to isomorphism and $`e(h)y`$ acts on $`V(\alpha ,x)`$ as scalar multiplication by $`\alpha (y)\tau _x(y)\delta _{h,x}`$. Hence, as reflected in equation (8), the associativity map
$$a_{V(\alpha ,x),V(\lambda ,y),V(\mu ,z)}:(V(\alpha ,x)V(\lambda ,y))V(\mu ,z)V(\alpha ,x)(V(\lambda ,y)V(\mu ,z))$$
(47)
in $`D^\omega (G)\text{-}\text{Mod}`$ is the scalar $`\omega (x,y,z)^1`$. Notice that the group structure on $`S`$ induced by the tensor product of $`D^\omega (G)\text{-}\text{Mod}`$ is isomorphic to $`\mathrm{\Gamma }^\omega `$ with the isomorphism $`\sigma (\alpha ,x)V(\alpha ,x)`$. Hence, by the results of \[TY98\], we have the following lemma.
###### Lemma 9.1
Let $`G`$ be a finite abelian group and $`\omega `$ be a normalized 3-cocycle such that $`[\omega ]H^3(G,^{})_{ab}`$. Then, $`D^\omega (G)\text{-}\text{Mod}`$ is monoidally equivalent to $`D^\omega (G)_0\text{-}\text{Mod}`$ if, and only if, $`infl[\omega ]`$ is trivial, where $`infl:H^3(G,^{})H^3(\mathrm{\Gamma }^\omega ,^{})`$ is inflation along the projection $`\mathrm{\Gamma }^\omega G`$. $`\mathrm{}`$
In fact, we have
###### Lemma 9.2
Suppose that $`G`$ is finite and abelian and that $`infl[\omega ]`$ is trivial. Then $`D^\omega (G)`$ and $`D^\omega (G)_0`$ are gauge equivalent as quasi-bialgebras.
Proof. Since $`infl[\omega ]`$ is trivial, there is a normalized 3-cochain $`f`$ on $`\mathrm{\Gamma }^\omega `$ such that
$$infl(\omega )=\delta f.$$
(48)
Let $`\tau T(\omega )`$ and let $`E_{\sigma (\alpha ,x)}=\frac{1}{|G|}e(x)_{gG}\frac{g}{\alpha (g)\tau _x(g)}`$. One can easily see that $`E_uE_v=\delta _{u,v}E_v`$ and $`\mathrm{\Delta }E_u=_{s,t\mathrm{\Gamma }^\omega }E_sE_t`$. Set
$$F=\underset{u,v\mathrm{\Gamma }^\omega }{}f(u,v)^1E_uE_v.$$
One can easily check that $`F`$ is invertible with
$$F^1=\underset{u,v\mathrm{\Gamma }^\omega }{}f(u,v)E_uE_v.$$
As $`f`$ is normalized,
$$(\epsilon id)F=1_{D^\omega (G)}=(id\epsilon )(F).$$
Moreover, by equation (48),
$$\mathrm{\Phi }=F_{23}(id\mathrm{\Delta })(F)\mathrm{\Phi }_0(\mathrm{\Delta }id)(F^1)F_{12}^1.$$
where $`\mathrm{\Phi }_0=1_{D^\omega (G)}1_{D^\omega (G)}1_{D^\omega (G)}`$. Since $`D^\omega (G)_0`$ is a commutative, $`D^\omega (G)`$ is a twist of $`D^\omega (G)_0`$ by $`F`$, that is, $`D^\omega (G)`$ and $`D^\omega (G)_0`$ are gauge equivalent. $`\mathrm{}`$
###### Lemma 9.3
$`infl[\omega ]\text{ker }\mathrm{\Lambda }_{\mathrm{\Gamma }^\omega }`$ .
Proof. Set $`\overline{\omega }=infl\omega `$. Define $`d:\mathrm{\Gamma }^\omega \times \mathrm{\Gamma }^\omega ^{}`$ by
$$d[\sigma (\alpha ,x)|\sigma (\lambda ,y)]=\frac{1}{\alpha (y)\tau _x(y)}.$$
One can check directly that $`(\overline{\omega },d)Z_{ab}^3(\mathrm{\Gamma }^\omega ,^{})`$. By the exact sequence (40), we have the commutative diagram
with the rows exact, where the vertical homomorphisms are induced by the epimorphism $`\mathrm{\Gamma }^\omega G`$. As
$$[(\overline{\omega },d)][\overline{\omega }]=infl[\omega ],$$
it follows from the diagram that $`infl[\omega ]\text{ker }\mathrm{\Lambda }_{\mathrm{\Gamma }^\omega }`$. $`\mathrm{}`$
###### Theorem 9.4
Let $`G`$ be a finite group of odd order and $`\omega _1,\omega _2`$ normalized 3-cocycles such that $`[\omega _1],[\omega _2]H^3(G,^{})_{ab}`$. Then, $`D^{\omega _1}(G)\text{-}\text{Mod}`$ and $`D^{\omega _2}(G)\text{-}\text{Mod}`$ are monoidally equivalent if, and only if, $`\mathrm{\Gamma }^{\omega _1}\mathrm{\Gamma }^{\omega _2}`$ as groups.
Proof. Since the group structure, induced by the underlying tensor product, on the isomorphism classes of irreducible representations of $`D^{\omega _i}(G)`$ is isomorphic to $`\mathrm{\Gamma }^{\omega _i}`$, $`\mathrm{\Gamma }^{\omega _1}\mathrm{\Gamma }^{\omega _2}`$ as groups if $`D^{\omega _1}(G)\text{-}\text{Mod}`$ and $`D^{\omega _2}(G)\text{-}\text{Mod}`$ are monoidal equivalent. Conversely, assume that $`\mathrm{\Gamma }^{\omega _1}\mathrm{\Gamma }^{\omega _2}`$ as groups. Then, $`\mathrm{\Gamma }^{\omega _1}`$ and $`\mathrm{\Gamma }^{\omega _2}`$ are isomorphic as Hopf algebras. By Lemma 9.3, $`infl[\omega _i]\text{ker }\mathrm{\Lambda }_{\mathrm{\Gamma }^{\omega _i}}`$. As $`G`$ is of odd order, $`|\mathrm{\Gamma }^{\omega _i}|`$ is odd, and it follows from Corollary 6.4 that $`\text{ker }\mathrm{\Lambda }_{\mathrm{\Gamma }^{\omega _i}}`$ is trivial. Hence, by Lemma 9.2, $`D^{\omega _i}(G)`$ is equivalent to $`\mathrm{\Gamma }^{\omega _i}`$ and hence $`D^{\omega _i}(G)\text{-}\text{Mod}`$ is monoidally equivalent to the tensor category $`\mathrm{\Gamma }^{\omega _i}\text{-}\text{Mod}`$ with the usual associativity constraint. $`\mathrm{}`$
###### Example 9.5
Theorem 9.4 is not true for groups of even order. For instance, take $`G=_2`$. By Proposition 6.3, $`\text{ker }\mathrm{\Lambda }_G_2`$ and hence $`\text{ker }\mathrm{\Lambda }_G=H^3(G,^{})`$. Let $`\omega _1`$ be a normalized 3-cocycle on $`G`$ whose cohomology class is nontrivial and let $`\omega _0=1`$. As $`\mathrm{\Lambda }_G[\omega _i]`$ is trivial $`(i=0,1)`$, $`\mathrm{\Gamma }^{\omega _i}_2\times _2`$ as groups. Therefore, inflation of cohomology along $`\mathrm{\Gamma }^{\omega _i}G`$ is injective for $`i=0,1`$. Thus, $`infl[\omega _1]`$ is non-trivial, while $`infl[\omega _0]`$ is obviously trivial. Thanks to lemma 9.1, if $`[_2\times _2]\text{-}\text{Mod}`$ equipped with the usual associativity constraint then $`D^{\omega _0}(G)\text{-}\text{Mod}`$ is monoidally equivalent to $`[_2\times _2]\text{-}\text{Mod}`$ but $`D^{\omega _1}(G)\text{-}\text{Mod}`$ is not. Note that $`[_4]`$ and $`[_2\times _2]`$ are the only 4-dimensional semisimple Hopf algebras up to isomorphism. As the fusion rule of $`D^{\omega _1}(G)\text{-}\text{Mod}`$ is isomorphic to $`_2\times _2`$, $`D^{\omega _1}(G)\text{-}\text{Mod}`$ and $`[_4]\text{-}\text{Mod}`$ are not monoidally equivalent. Hence, $`D^{\omega _1}(G)`$ cannot be obtained from any Hopf algebra by a twist.
###### Remark 9.6
In the paper \[DPR92\], p92, the authors asked whether $`D^\omega (G)`$ can be obtained by twisting a Hopf algebra. Theorem 9.4 gives an affirmative answer to the question under the conditions stated in the theorem. However, in general, Example 9.5 gives a negative answer to the question.
## 10 Gauge Equivalence and Quadratic Forms
Let $`G`$ be a finite abelian group, $`\omega Z^3(G,^{})_{ab}`$ and $`\tau T(\omega )`$. Then, as shown in the proof of Lemma 9.3, $`(infl\omega ^1,d_\omega )Z_{ab}^3(\mathrm{\Gamma }^\omega (G),^{})`$ where $`d_\omega :\mathrm{\Gamma }^\omega (G)\times \mathrm{\Gamma }^\omega (G)^{}`$ is defined by
$$d_\omega [\sigma (\alpha ,x)|\sigma (\lambda ,y)]=\alpha (y)\tau _x(y).$$
Define $`q_\omega :\mathrm{\Gamma }^\omega (G)^{}`$ to be the trace of the abelian 3-cocycle $`(infl\omega ^1,d_\omega )`$. Obviously, $`q_\omega `$ is a quadratic map given by
$$q_\omega (\sigma (\alpha ,x))=\alpha (x)\tau _x(x).$$
(49)
Denote by $`b_\omega `$ the symmetric bicharacter on $`\mathrm{\Gamma }^\omega (G)`$ associated to $`q_\omega `$, namely $`b_\omega (x,y)=\frac{q_\omega (xy)}{q_\omega (x)q_\omega (y)}`$. One can easily show that $`b_\omega `$ is given by
$$b_\omega (\sigma (\alpha ,x),\sigma (\lambda ,y))=\alpha (y)\lambda (x)\tau _x(y)\tau _y(x)$$
for any $`\sigma (\alpha ,x),\sigma (\lambda ,y)\mathrm{\Gamma }^\omega (G)`$. Obviously, $`b_\omega `$ is non-degenerate on $`\mathrm{\Gamma }^\omega (G)`$. Let $`,_\omega `$ be the $``$-linear extension of $`b_\omega `$ on $`D^\omega (G)`$.
###### Lemma 10.1
With the previous notation, the map $`\phi _\omega :D^\omega (G)_0D^\omega (G)_0^{}`$, $`\phi _\omega :xx,_\omega `$ is identical to the Hopf algebra isomorphism $`\phi :D^\omega (G)_0D^\omega (G)_0^{}`$ defined in (13). In particular, $`,_\omega `$ is a non-degenerate symmetric bilinear form on $`D^\omega (G)`$. Moreover, for any $`u,v\mathrm{\Gamma }^\omega `$,
$$(\phi _\omega (u)\phi _\omega (v))R=d_\omega [v|u]$$
(50)
and
$$q_\omega (u)=(\phi _\omega (u)\phi _\omega (u))R$$
(51)
where $`R`$ is the $``$-matrix of $`D^\omega (G)`$.
Proof. To show that $`\phi _\omega =\phi `$, it suffices to show that $`\phi _\omega (u)=\phi (u)`$ for $`u\mathrm{\Gamma }^\omega `$. Hence, it is enough to show that
$$\phi (u)(v)=\phi _\omega (u)(v)=b_\omega (u,v)$$
for any $`u,v\mathrm{\Gamma }^\omega `$. However, it follows from (13) that
$$\phi (\sigma (\alpha ,x))(\sigma (\lambda ,y))=\alpha (y)\lambda (x)\tau _x(y)\tau _y(x)=b_\omega (\sigma (\alpha ,x),\sigma (\lambda ,y)).$$
Now, $`\phi _\omega (\sigma (\alpha ,x))(e(g)y)=\alpha (y)\tau _x(y)\delta _{g,x}`$. Thus,
$$(\phi _\omega (\sigma (\alpha ,x))\phi _\omega (\sigma (\lambda ,y)))R=\lambda (x)t_y(x)=d_\omega [\sigma (\lambda ,y)|\sigma (\alpha ,x)]$$
and hence
$$(\phi _\omega (u)\phi _\omega (u)R=d_\omega [u|u]=q_\omega (u)$$
for $`u\mathrm{\Gamma }^\omega `$. $`\mathrm{}`$
###### Lemma 10.2
Let $`\omega _1`$ and $`\omega _2`$ be normalized 3-cocycles in $`Z^3(G_1,^{})_{ab}`$ and $`Z^3(G_2,^{})_{ab}`$ respectively. If $`j:D^{\omega _1}(G_1)_0D^{\omega _2}(G_2)_0`$ is a bialgebra isomorphism, then $`j^t:D^{\omega _2}(G_2)_0D^{\omega _1}(G_1)_0`$, defined by
$$j^t(u),v_{\omega _1}=u,j(v)_{\omega _2}$$
for any $`uD^{\omega _2}(G_2)`$ and $`vD^{\omega _1}(G_1)`$, is also a bialgebra isomorphism.
Proof. Consider the map $`j^{}:D^{\omega _2}(G_2)_0D^{\omega _1}(G_1)(G)_0`$ given by
$$j^{}=\phi _{\omega _1}^1j^{}\phi _{\omega _2}.$$
Since $`j:D^{\omega _1}(G_1)_0D^{\omega _2}(G_2)_0`$ is a bialgebra isomorphism, and so is $`j^{}:D^{\omega _2}(G_2)_0^{}D^{\omega _1}(G_1)_0^{}`$. It follows from Lemma 10.1 that $`j^{}`$ is also a bialgebra isomorphism. Moreover, for any $`uD^{\omega _2}(G_2)`$ and $`vD^{\omega _1}(G_1)`$,
$$j^{}(u),v_{\omega _1}=\left(\phi _{\omega _1}j^{}(u)\right)(v)=\left(j^{}\phi _{\omega _2}(u)\right)(v)=\phi _{\omega _2}(j(u))(v)=u,j(v)_{\omega _2}.$$
Hence, $`j^{}=j^t`$. $`\mathrm{}`$
###### Definition 10.3
Let $`\mathrm{\Gamma }_1`$, $`\mathrm{\Gamma }_2`$ be abelian groups, and $`b_1`$, $`b_2`$ bicharacters on $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ respectively. The pairs $`(\mathrm{\Gamma }_1,b_1)`$ and $`(\mathrm{\Gamma }_2,b_2)`$ are said to be equivalent if there is a group isomorphism $`j:\mathrm{\Gamma }_1\mathrm{\Gamma }_2`$ such that
$$b_2(j(x),j(y))=b_1(x,y)$$
for $`x,y\mathrm{\Gamma }_1`$. Similarly, let $`q_1:\mathrm{\Gamma }_1^{}`$ and $`q_2:\mathrm{\Gamma }_2^{}`$ be quadratic maps. The pairs $`(\mathrm{\Gamma }_1,q_1)`$ and $`(\mathrm{\Gamma }_2,q_2)`$, which we call quadratic spaces, are said to be equivalent if there is a group isomorphism $`j:\mathrm{\Gamma }_1\mathrm{\Gamma }_2`$ such that
$$q_2j=q_1.$$
The orthogonal sum of two quadratic spaces $`(\mathrm{\Gamma }_1,q_1)`$ and $`(\mathrm{\Gamma }_2,q_2)`$ is the quadratic space $`(\mathrm{\Gamma }_1\times \mathrm{\Gamma }_2,q)`$ where $`q(x,y)=q_1(x)q_2(y)`$.
###### Theorem 10.4
Let $`G_1`$, $`G_2`$ be a finite abelian groups and $`\omega _1`$, $`\omega _2`$ normalized 3-cocycles in $`Z^3(G_1,^{})_{ab}`$ and $`Z^3(G_2,^{})_{ab}`$ respectively. Then, $`D^{\omega _1}(G_1)`$ and $`D^{\omega _2}(G_2)`$ are gauge equivalent if, and only if, $`(\mathrm{\Gamma }^{\omega _1},q_{\omega _1})`$ and $`(\mathrm{\Gamma }^{\omega _2},q_{\omega _2})`$ are equivalent.
Proof. Suppose that $`D^{\omega _1}(G_1)`$ and $`D^{\omega _2}(G_2)`$ are equivalent as quasi-triangular quasi Hopf algebras. Then there exists a gauge transform $`FD^{\omega _2}(G_2)D^{\omega _2}(G_2)`$ such that $`D^{\omega _1}(G_1)`$ and $`D^{\omega _2}(G_2)_F`$ are isomorphic as quasi-triangular quasi-bialgebras. Let $`j:D^{\omega _1}(G_1)D^{\omega _2}(G_2)_F`$ be such an isomorphism. Then,
$$(jj)R_1=F_{21}R_2F^1$$
(52)
where $`R_1`$ and $`R_2`$ are the $``$-matrices of $`D^{\omega _1}(G_1)`$ and $`D^{\omega _2}(G_2)`$ respectively. Since $`\mathrm{\Delta }_2=(\mathrm{\Delta }_2)_F`$, $`D^{\omega _2}(G_2)_F`$ and $`D^{\omega _2}(G_2)_0`$ are identical as bialgebras. By Lemma 10.2, $`j^t:D^{\omega _2}(G_2)_0D^{\omega _1}(G_1)_0`$ is a bialgebra isomorphism. In particular, $`j^t(\mathrm{\Gamma }^{\omega _2})=\mathrm{\Gamma }^{\omega _1}`$. Let $`\{E_u\}_{u\mathrm{\Gamma }^{\omega _2}}`$ be the dual basis of $`\mathrm{\Gamma }^{\omega _2}`$ with respect to the pairing $`,_{\omega _2}`$. Let $`F=_{u,v\mathrm{\Gamma }^{\omega _2}}f(u,v)E_uE_v`$. For any $`u\mathrm{\Gamma }^{\omega _2}`$, apply $`\phi _{\omega _2}(u)\phi _{\omega _2}(u)`$ to equation (52). The left hand side of the equation becomes
$$(\phi _{\omega _2}(u)\phi _{\omega _2}(u))(jj)R_1=(\phi _{\omega _1}(j^tu)\phi _{\omega _1}(j^tu))R_1=q_{\omega _1}(j^t(u))$$
and the right hand side becomes
$$(\phi _{\omega _2}(u)\phi _{\omega _2}(u))(F_{21}R_2F^1)=q_{\omega _2}(u)f(u,u)f(u,u)^1=q_{\omega _2}(u).$$
Therefore, $`(\mathrm{\Gamma }^{\omega _1},q_{\omega _1})`$ and $`(\mathrm{\Gamma }^{\omega _2},q_{\omega _2})`$ are equivalent. Conversely, assume that there exists an isomorphism $`j:\mathrm{\Gamma }^{\omega _2}\mathrm{\Gamma }^{\omega _1}`$ such that $`q_{\omega _1}j=q_{\omega _2}`$. Set
$$\widehat{j}(infl\omega _1)(u,v,w)=infl\omega _1(j(u),j(v),j(w))$$
and
$$\widehat{j}d_{\omega _1}[u|v]=d_{\omega _1}[j(u)|j(v)]$$
for any $`u,v\mathrm{\Gamma }^{\omega _2}`$. Then, $`(\widehat{j}(infl\omega _1^1),\widehat{j}(d_{\omega _1}))`$ is also an abelian 3-cocycle of $`\mathrm{\Gamma }^{\omega _2}`$. The trace of this cocycle is $`q_{\omega _1}j`$. By Theorem 5.2, there exists a normalized 2-cochain $`f`$ on $`\mathrm{\Gamma }^{\omega _2}`$ such that
$`infl\omega _1^1(j(u),j(v),j(w))`$ $`=`$ $`\omega _2^1(u,v,w)\delta f(u,v,w)`$ (53)
$`d_{\omega _1}[j(u)|j(v)]`$ $`=`$ $`d_{\omega _2}[u|v]f(u,v)/f(v,u)`$ (54)
for any $`u,v,w\mathrm{\Gamma }^{\omega _2}`$. Let
$$F=\underset{u,v\mathrm{\Gamma }^{\omega _2}}{}f(u,v)E_uE_v.$$
Notice that $`ϵ(E_u)=E_u,\sigma (1,1)=\delta _{1,u}`$ and $`E_uE_v=\delta _{u,v}E_u`$ for any $`u,v\mathrm{\Gamma }^{\omega _2}`$. This implies that $`F`$ is invertible in $`D^{\omega _2}(G_2)D^{\omega _2}(G_2)`$ and
$$(idϵ)F=1_{D^{\omega _2}(G_2)}=(ϵid)F.$$
For simplicity, we denote the linear extension of $`j`$ from $`D^{\omega _2}(G_2)_0`$ to $`D^{\omega _1}(G_1)_0`$ by the same symbol $`j`$. Obviously, $`j`$ is a bialgebra isomorphism. By Lemma 10.2, $`j^t`$ is also a bialgebra isomorphism. For any $`u,v,w\mathrm{\Gamma }^{\omega _2}`$,
$`(\phi _{\omega _2}(u)\phi _{\omega _2}(v)\phi _{\omega _2}(w))(j^tj^tj^t)\mathrm{\Phi }_1`$ $`=`$ $`(\phi _{\omega _1}(j(u))\phi _{\omega _1}(j(v))\phi _{\omega _1}(j(w)))\mathrm{\Phi }_1`$
$`=`$ $`infl\omega _{1}^{}{}_{}{}^{1}(j(u),j(v),j(w))`$
and
$$(\phi _{\omega _2}(u)\phi _{\omega _2}(v)\phi _{\omega _2}(w))(F_{23}(id\mathrm{\Delta })(F)\mathrm{\Phi }_2(\mathrm{\Delta }id)(F^1)F_{12}^1)=(infl\omega _2)\delta f(u,v,w).$$
Hence, by equation (53), we obtain
$$(j^tj^tj^t)\mathrm{\Phi }_1=F_{23}(id\mathrm{\Delta })(F)\mathrm{\Phi }_2(\mathrm{\Delta }id)(F^1)F_{12}^1).$$
It follows from equation (50) that
$$(\phi _{\omega _2}(u)\phi _{\omega _2}(v))(j^tj^t)R_1=(\phi _{\omega _1}(j(u))\phi _{\omega _1}(j(v)))R_1=d_{\omega _1}[j(v)|j(u)]$$
and
$$(\phi _{\omega _2}(u)\phi _{\omega _2}(v))F_{21}R_2F^1=f(v,u)f^1(u,v)d_{\omega _2}[v|u].$$
Hence, by equation (54),
$$(j^tj^t)R_1=F_{21}R_2F^1.$$
Therefore, $`j^t:D^{\omega _1}(G_1)D^{\omega _2}(G_2)_F`$ is a quasi-bialgebra isomorphism. $`\mathrm{}`$
###### Theorem 10.5
Let $`G_1`$, $`G_2`$ be finite abelian groups and $`\omega _1Z^3(G_1,^{})_{ab}`$, $`\omega _2Z^3(G_2,^{})_{ab}`$. Then, $`D^{\omega _1}(G_1)`$ and $`D^{\omega _2}(G_2)`$ are gauge equivalent if, and only if $`D^{\omega _1}(G_1)\text{-}\text{Mod}`$ and $`D^{\omega _2}(G_2)\text{-}\text{Mod}`$ are equivalent additive braided tensor categories; that is, there exists a braided tensor equivalence $`(,\phi _0,\phi _2)`$ between $`D^{\omega _1}(G_1)\text{-}\text{Mod}`$ and $`D^{\omega _2}(G_2)\text{-}\text{Mod}`$ with $``$ being additive.
Proof. The “only if” part of the statement is well-known (cf. \[Kas95\]). Conversely, let $`(,\phi _0,\phi _2)`$ be an additive braided tensor equivalence from $`K\text{-}\text{Mod}`$ to $`S\text{-}\text{Mod}`$ where $`K=D^{\omega _1}(G_1)`$ and $`S=D^{\omega _2}(G_2)`$. By Morita theory (cf. \[AF92\], \[Lam99\]), $``$ induces an algebra isomorphism $`\alpha `$ of $`\text{End}_S((K))`$ with $`K`$ and so $`(K)`$ becomes a $`S`$-$`K`$-bimodule. Moreover, $``$ is equivalent to the tensor functor $`(K)_K\mathrm{?}`$. Note that both $`S`$ and $`K`$ are semisimple and their irreducible modules are 1-dimensional. Therefore, $`(K)={}_{S}{}^{}S`$ and the right $`K`$-module action is given by the isomorphism $`\alpha ^1:KS`$ where $`S`$ is naturally identified with $`\text{End}({}_{S}{}^{}S)`$. For any $`MK\text{-}\text{Mod}`$, define $`\alpha ^{}(M)`$ to be the left $`S`$-module with the underlying space $`M`$ and the left $`S`$-action given by
$$sm=\alpha (s)m$$
for any $`sS`$ and $`mM`$. It is easy to see that $`\text{Hom}_K(M,N)=\text{Hom}_S(\alpha ^{}(M),\alpha ^{}(N))`$ for any $`M,NK\text{-}\text{Mod}`$. Hence, $`\alpha ^{}`$ defines an additive functor from $`K\text{-}\text{Mod}`$ to $`S\text{-}\text{Mod}`$. It is straightforward to show that $`\alpha ^{}`$ and $`(K)_K\mathrm{?}`$ are equivalent functors. Without loss of generality, we may assume $`=\alpha ^{}`$.
Now $`\stackrel{\phi _0}{}\alpha ^{}()`$ is scalar multiplication by a nonzero complex number $`\kappa `$. For any $`MS\text{-}\text{Mod}`$, define $`\eta _M`$ to be scalar multiplication $`\kappa ^1`$. Then $`\eta :\alpha ^{}\alpha ^{}`$ is a natural isomorphism. Hence, $`(\alpha ^{},id,\kappa \phi _2)`$ is a braided tensor equivalence from $`K\text{-}\text{Mod}`$ to $`S\text{-}\text{Mod}`$. Consider $`\kappa \phi _2(K,K)`$. Then $`\kappa \phi _2(K,K)(1_K1_K)`$ is an invertible in $`KK`$ and we set $`F^1=\kappa \phi _2(K,K)(1_K1_K)`$. Then, $`F`$ is a gauge transform of $`K`$ (cf. \[Kas95\],p381) and
$`\alpha (\mathrm{\Phi }_S)`$ $`=`$ $`F_{23}(id\mathrm{\Delta }_K)(F)\mathrm{\Phi }_K(\mathrm{\Delta }id)(F^1)F_{12}^1`$
$`\alpha (R_S)`$ $`=`$ $`F_{21}R_KF^1`$
where $`\mathrm{\Phi }_S`$, $`R_S`$ and $`\mathrm{\Phi }_K`$, $`R_K`$ are the associators and $``$-matrices of $`S`$ and $`K`$ respectively. $`\mathrm{}`$
###### Definition 10.6
Let $`\mathrm{\Gamma }`$ be a finite abelian group and $`b:\mathrm{\Gamma }\times \mathrm{\Gamma }^{}`$ a symmetric bicharacter. For any subset $`M`$ of $`G`$, denote by $`M^{}`$ the subgroup
$$\{x\mathrm{\Gamma }|b(x,M)=1\}$$
of $`\mathrm{\Gamma }`$. A subgroup $`N`$ of $`\mathrm{\Gamma }`$ is called a metabolizer of $`(\mathrm{\Gamma },b)`$ if $`N=N^{}`$. A quadratic map $`q:\mathrm{\Gamma }^{}`$ is called non-degenerate if the associated bicharacter $`b_q(x,y)=\frac{q(xy)}{q(x)q(y)}`$ is non-degenerate. A subgroup $`N`$ of $`\mathrm{\Gamma }`$ is called a metabolizer of $`(\mathrm{\Gamma },q)`$ if $`q|_N=1`$ and $`N`$ is a metabolizer of $`(\mathrm{\Gamma },b_q)`$. A metabolizer $`N`$ of $`(\mathrm{\Gamma },q)`$ or $`(\mathrm{\Gamma },b)`$ is called split if $`N`$ is a summand of $`\mathrm{\Gamma }`$.
###### Remark 10.7
1. Let $`G`$ be a finite abelian group and $`\omega `$ a normalized 3-cocycle in $`Z^3(G,^{})_{ab}`$. Then, $`q_\omega (\widehat{G})=1`$. Moreover, if $`b_\omega (\sigma (\alpha ,x),\sigma (\lambda ,1))=1`$ for all $`\lambda \widehat{G}`$, then $`\lambda (x)=1`$ for all $`\lambda \widehat{G}`$. This implies $`x=1`$. Therefore, $`\widehat{G}^{}=\widehat{G}`$ and hence $`\widehat{G}`$ is a metabolizer of $`(\mathrm{\Gamma }^\omega ,q_\omega )`$.
2. If $`(\mathrm{\Gamma },q)`$ is non-degenerate with metabolizer $`G`$, then $`|G|^2=|\mathrm{\Gamma }|`$.
Let $`G`$, $`G^{}`$ be finite abelian groups, $`\omega Z^3(G,^{})_{ab}`$ and $`\omega ^{}Z^3(G^{},^{})_{ab}`$. Let $`\zeta Z^3(G_1\times G_2,^{})_{ab}`$ be the product of the inflations of $`\omega `$ and $`\omega ^{}`$, namely
$$\zeta =(infl\omega )(infl\omega ^{})$$
where the inflations are induced by the natural surjections $`G\times G^{}G`$ and $`G\times G^{}G^{}`$. By the remarks following Proposition 4.1, there is an isomorphism of groups $`\iota :\mathrm{\Gamma }^\omega (G)\times \mathrm{\Gamma }^\omega ^{}(G)\mathrm{\Gamma }^\zeta (G\times G^{})`$ given by
$$\iota (\underset{gG}{}\lambda _x(g)e(g)x,\underset{hG^{}}{}\lambda _y^{}(h)e(h)y)=\underset{(g,h)G\times G^{}}{}\lambda _x(g)\lambda _y^{}(h)e(g,h)(x,y)$$
where $`\delta \lambda _x=\omega _x`$, $`\delta \lambda _y^{}=\omega _y^{}`$ and $`xG`$, $`yG^{}`$. Since, for any $`u=_{(g,h)G\times G^{}}\lambda _{(x,y)}(g,h)e(g,h)(x,y)\mathrm{\Gamma }^\zeta (G\times G^{})`$, $`q_\zeta (u)=\lambda _{(x,y)}(x,y)`$, $`\iota `$ defines an equivalence of the quadratic forms $`(\mathrm{\Gamma }^\zeta ,q_\zeta )`$ and $`(\mathrm{\Gamma }^\omega ,q_\omega )(\mathrm{\Gamma }^\omega ^{},q_\omega ^{})`$.
Conversely, suppose that $`H`$ is a finite abelian group and $`\eta Z^3(H,^{})_{ab}`$ such that $`(\mathrm{\Gamma }^\eta ,q_\eta )`$ is equivalent to the orthogonal sum $`(\mathrm{\Gamma }^\omega ,q_\omega )(\mathrm{\Gamma }^\omega ^{},q_\omega ^{})`$. Then, $`(\mathrm{\Gamma }^\eta (H),q_\eta )`$ and $`(\mathrm{\Gamma }^\zeta (G\times G^{}),q_\zeta )`$ are equivalent quadratic forms. By virtue of Proposition 10.4, $`D^\eta (H)`$ is equivalent to $`D^\omega (H)D^\omega ^{}(H^{})`$ as quasi-triangular quasi-bialgebras. This proves
###### Proposition 10.8
Let $`G`$ be a finite abelian group and $`\eta Z^3(G,^{})_{ab}`$. Then, $`D^\eta (G)`$ is equivalent to $`D^\omega (H)D^\omega ^{}(H^{})`$ as quasi-triangular quasi-bialgebras for some abelian groups $`H`$, $`H^{}`$ and $`\omega Z^3(H,^{})_{ab}`$, $`\omega ^{}Z^3(H^{},^{})_{ab}`$ if, and only if, $`(\mathrm{\Gamma }^\eta ,q_\eta )`$ is equivalent to $`(\mathrm{\Gamma }^\omega ,q_\omega )(\mathrm{\Gamma }^\omega ^{},q_\omega ^{})`$.
## 11 Lattices
We use the following notation: $`M`$ is a self-dual, even, lattice with respect to the nonsingular bilinear form
$$,:M\times M.$$
Thus $`xx,`$ is an isomorphism of $`M`$ with $`\text{Hom}_{}(M,)`$, moreover $`x,x`$ is an even integer for $`xM`$. Note that $`,`$ is not necessarily positive definite.
Let $`E`$ be the space $`_{}M`$ equipped with the $``$-linear extension of $`,`$, let $`MLE`$ with $`|L:M|<\mathrm{}`$, and let $`L_0=\{xE|x,L\}`$ be the $``$-dual of $`L`$. We set $`G=L/M`$. There is a short exact sequence
$$0M/L_0L/L_0L/M0$$
(55)
and we pick a section $`s:L/ML`$, such that $`s(0)=0`$, which naturally defines a section $`\overline{s}:L/ML/L_0`$.
Because $`M`$ is self-dual, the pairing $`,:L\times L`$ induces a perfect pairing
$$p:M/L_0\times L/MS^1,(x+L_0,y+M)e^{2\pi ix,y}$$
(56)
and so there is a natural identification of $`M/L_0`$ with $`\text{Hom}(L/M,S^1)=\widehat{G}`$. Thus the sequence (55) is of the type we have been considering. That is, the triple $`L_0ML`$ defines an element of $`H_{ab}^2(G,\widehat{G})`$ with $`G=L/M`$, and it is well-known that a 2-cocycle $`\beta Z_{ab}^2(G,\widehat{G})`$ which corresponds to the triple is defined, using the section $`\overline{s}`$, via
$$\beta (x,y)=\overline{s}(x)+\overline{s}(y)\overline{s}(x+y)$$
(57)
for $`x`$, $`yG`$.
Following Dong-Lepowsky \[DL93\], we pick an alternating, bilinear map $`c:L\times LS^1`$ with the property that
$$c(x,y)=(1)^{x,y}$$
(58)
for $`x,yM`$. Such $`c`$ always exists (cf. \[DL93\], Remark 12.18).
###### Proposition 11.1
With the previous notation, set
$$\omega (g,x,y)=c(s(g),s(x)+s(y)s(x+y))e^{\pi is(g),s(x)+s(y)s(x+y)}$$
(59)
$$\tau _g(x)=c(s(g),s(x))e^{\pi is(g),s(x)}.$$
(60)
Then $`\omega Z^3(G,^{})_{ab}`$, $`\omega _g=\delta \tau _g`$, and $`\mathrm{\Lambda }_G([\omega ])=[\beta ]`$.
Proof. After some computation, we find that for $`g,h,x,yG`$,
$$\begin{array}{cc}\hfill \delta \omega (g,h,x,y)=& c(s(g)+s(h)s(g+h),s(x)+s(y)s(x+y))\hfill \\ & e^{\pi is(g)+s(h)s(g+h),s(x)+s(y)s(x+y)}.\hfill \end{array}$$
(61)
Note that both $`a=s(g)+s(h)s(g+h)`$ and $`b=s(x)+s(y)s(x+y)`$ lie in $`M`$. Using equation (58), the expression of (61) is equal to 1 and, as a result, $`\omega `$ is a normalized 3-cocycle.
The identity $`\omega _g=\delta \tau _g`$ follows immediately from (59) and (60). Finally, using the expression (15) together with (60) shows that $`\mathrm{\Lambda }_G([\omega ])`$ is the cohomology class represented by $`\beta _1Z_{ab}^2(G,\widehat{G})`$ where
$`\beta _1(x,y)(g)`$ $`=`$ $`{\displaystyle \frac{c(s(g),s(x))c(s(x),s(g))c(s(g),s(y))c(s(y),s(g))}{c(s(g),s(x+y))c(s(x+y),s(g))}}e^{2\pi is(g),s(x)+s(y)s(x+y)}`$
for $`x,y,gG`$. Since $`c`$ is alternating, equations (56) and (57) yields
$$\beta _1(x,y)(g)=p(\beta (x,y),g).$$
So by definition of the pairing $`p`$, $`\beta _1=\beta `$. $`\mathrm{}`$
We now have two quadratic forms associated to this situation: $`(L/L_0,q_L)`$ where $`q_L(x+L_0)=e^{\pi ix,x}`$ for $`xL`$; and $`(\mathrm{\Gamma }^\omega ,q_\omega )`$ canonically associated to the 3-cocycle $`\omega `$ (cf. (49)).
###### Lemma 11.2
These two quadratic spaces are equivalent.
Proof. Let notation be as before. From section 3, $`\mathrm{\Gamma }^\omega `$ is a central extension of $`G`$ by $`\widehat{G}`$. The additive version of the corresponding 2-cocycle $`\beta `$ is given by (57). Here, as before, $`G=L/M`$ and $`\widehat{G}=M/L_0`$ under the identification given in (56). There is an isomorphism $`j:L/L_0\mathrm{\Gamma }^\omega `$ defined for $`xL`$ via
$$x+L_0\sigma (x+L_0\overline{s}(x+M),x+M).$$
(62)
Now use equations (49), (60) to see that
$$q_\omega (j(x+L_0))=p(x+L_0\overline{s}(x+M),x+M)e^{\pi is(x),s(x)}.$$
(63)
Write $`s(x+M)=x+m`$ for some $`mM`$. Then we see that
$`q_\omega (j(x+L_0))`$ $`=`$ $`e^{\pi i(2x(x+m),x+x+m,x+m)}`$
$`=`$ $`e^{\pi ix,x}=q_L(x+L_0).`$
This completes the proof of the lemma. $`\mathrm{}`$
We now begin again, this time with a non-degenerate quadratic form $`q`$ on the finite abelian group $`\mathrm{\Gamma }`$. It is convenient to write $`q`$ additively, that is $`q:\mathrm{\Gamma }/`$. We also consider $`q_1=2q`$ as a quadratic form with values in $`/2`$. According to \[Wal63\], Theorem (6), there is a rational lattice $`(H^{},,)`$ i.e. $`,:H^{}\times H^{}`$ non-degenerate, an even (integral) sublattice $`HH^{}`$, and an equivalence of quadratic forms
where $`\overline{q}`$ is the quadratic form arising from $`,`$. Thus, for $`x,yH^{}`$ we have $`\overline{q}(x+H)=x,x(mod2)`$ and
$`q_1(j(x+H))`$ $``$ $`x,x(mod2)`$ (64)
$`b(j(x+H),j(y+H))`$ $``$ $`x,y(mod)`$ (65)
where $`b`$ is the bilinear form on $`\mathrm{\Gamma }`$ associated to $`q`$.
Now assume that $`(\mathrm{\Gamma },q)`$ has a metabolizer $`\widehat{G}`$, and write $`j^1(\widehat{G})=M/H`$ for a sublattice $`M`$ with $`HMH^{}`$. Since $`\widehat{G}^{}=\widehat{G}`$, it follows from (65) that $`M`$ is a self-dual, integral sublattice of $`H^{}`$. Similarly, equation (65) shows that $`M`$ is an even lattice. We now have
###### Theorem 11.3
1. Let $`ML`$ be a pair of rational lattices with $`M`$ even and self-dual, and let $`L_0M`$ be the $``$-dual of $`L`$. Then the pair $`(L/L_0,q_L)`$ is a non-degenerate quadratic space with metabolizer $`M/L_0`$. Furthermore,
2. There is $`\omega Z^3(G,^{})_{ab}`$ such that the corresponding quadratic space $`(\mathrm{\Gamma }^\omega ,q_\omega )`$ is equivalent to $`(L/L_0,q_L)`$. Moreover, $`\mathrm{\Lambda }(\omega )`$ is cohomologous to the element of $`Z_{ab}^2(L/M,M/L_0)`$ defined by (57).
3. Conversely, suppose that $`(\mathrm{\Gamma },q)`$ is a non-degenerate quadratic space with metabolizer. Then there are lattices $`ML`$ as in (i) such that $`(\mathrm{\Gamma },q)`$ and $`(L/L_0,q_L)`$ are equivalent quadratic spaces.
4. Every non-degenerate quadratic space $`(\mathrm{\Gamma },q)`$ with a metabolizer $`\widehat{G}`$ is equivalent to one of the form $`(\mathrm{\Gamma }^\omega ,q_\omega )`$ for suitable $`\omega Z^3(G,^{})_{ab}`$.
Proof. (i) and (ii) follow from Proposition 11.1 and Lemma 11.2, while (iii) follows from the discussion following the proof of the Lemma. Part (iv) follows from (i)–(iii). $`\mathrm{}`$
###### Remark 11.4
¿From Theorem 11.3 and our earlier work, we see that the following pieces of data are more or less equivalent:
1. A pair of rational lattices $`ML`$ with $`M`$ even and self-dual, $`L/MG`$.
2. A non-degenerate quadratic space $`(\mathrm{\Gamma },q)`$ with metabolizer $`\widehat{G}`$.
3. A cohomology class $`[\omega ]H^3(G,^{})_{ab}`$.
4. A cohomology class $`[\beta ]Z_{ab}^2(G,\widehat{G})`$ which is in the image of $`S`$ (cf. Theorem 8.4).
## 12 Gauss Sums
We begin with a discussion of Gauss sums on pairs $`(M,q)`$ consisting of a finite abelian group $`M`$ and a non-degenerate quadratic form $`q:M/`$. Thus, $`q(x+y)=q(x)+q(y)+b(x,y)`$ where $`b:M\times M/`$ is a non-degenerate bilinear form. In our applications, $`(M,q)`$ will be the pair $`(\mathrm{\Gamma }^\omega ,q_\omega )`$ that we have already considered, or something related to it. Our discussion is closely related to that of tom Dieck (\[tD79\], 2.2) adapted to our present needs.
The Gauss sum $`G(M,q)`$ is defined via
$$G(M,q)=\underset{mM}{}e^{2\pi iq(m)}.$$
The Gauss map<sup>5</sup><sup>5</sup>5we say map because $`\gamma `$ can be considered as function on a suitable Witt group. $`\gamma `$ is defined via
$$\gamma (M,q)=\frac{1}{\sqrt{|M|}}G(M,q).$$
###### Lemma 12.1
Let $`NM`$ be a subgroup of $`M`$ such that the restriction of $`q`$ to $`N`$ vanishes identically. Then $`q`$ induces a non-degenerate quadratic form $`\overline{q}`$ on $`N^{}/N`$, and we have
$$\gamma (M,q)=\gamma (N^{}/N,\overline{q}).$$
Proof. Note that $`\overline{q}`$ is defined via $`\overline{q}(m+N)=q(m)`$ for $`mN^{}`$. It is easy to see that this is well-defined and yields a quadratic form on $`N^{}/N`$ which is non-degenerate. Then we calculate:
$`\gamma (M,q)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{|M|}}}{\displaystyle \frac{1}{|N|}}{\displaystyle \underset{nN}{}}{\displaystyle \underset{mM}{}}e^{2\pi iq(m+n)}`$
$`=`$ $`{\displaystyle \frac{1}{\sqrt{|M|}}}{\displaystyle \frac{1}{|N|}}{\displaystyle \underset{mM}{}}e^{2\pi iq(m)}{\displaystyle \underset{nN}{}}e^{2\pi ib(m,n)}`$
where we used $`q(m+n)=q(m)+q(n)+b(m,n)`$ and $`q(n)=0`$ for $`nN`$. Now the inner sum vanishes whenever $`mN^{}`$ by the orthogonality of group characters, and is otherwise equal to $`|N|`$. So in fact
$`\gamma (M,q)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{|M|}}}{\displaystyle \underset{mN^{}}{}}e^{2\pi iq(m)}`$
$`=`$ $`{\displaystyle \frac{|N|}{\sqrt{|N||N^{}|}}}{\displaystyle \underset{mN^{}/N}{}}e^{2\pi i\overline{q}(m)}`$
$`=`$ $`\gamma (N^{}/N,\overline{q}).\mathrm{}`$
###### Corollary 12.2
If $`(M,q)`$ has a metabolizer then $`\gamma (M,q)=1`$.
###### Lemma 12.3
Let $`p`$ be a prime and $`(M,q)`$ a non-degenerate quadratic space on the non-trivial cyclic $`p`$-group $`M`$. One of the following holds:
1. $`p`$ is odd, $`|M|`$ is a square, and $`(M,q)`$ has a metabolizer.
2. $`p1(mod4)`$, $`|M|`$ is not a square, and $`\gamma (M,q)=\pm 1`$ .
3. $`p3(mod4)`$, $`|M|`$ is not a square, and $`\gamma (M,q)=\pm i`$ .
4. $`p=2`$ and $`\gamma (M,q)`$ is a primitive 8th root of unity.
Proof. Pick $`NM`$ maximal subject to the condition that the restriction of $`q`$ to $`N`$ is identically zero. Set $`V=N^{}/N`$, so that $`\gamma (M,q)=\gamma (V,\overline{q})`$ as in Lemma 12.1. Now the maximality of $`N`$ implies that $`(V,\overline{q})`$ is anisotropic, that is if $`vV`$ and $`\overline{q}(v)=0`$ then in fact $`v=0`$. It is then easy to see that $`|V|p^2`$, and even $`|V|p`$ if $`p`$ is odd. First assume that $`p`$ is odd. Then (i) holds if $`V=0`$, and otherwise $`|V|=p`$ and $`|M|=|N|^2|V|`$ is not a square. Then (ii) or (iii) holds by the classical calculation of Gauss sums (cf. \[tD79\]).
Now take $`p=2`$. If $`|V|=2`$ or 4 it is an easy direct calculation that (iv) holds. To complete the lemma, we must show that, in fact, $`V0`$ if $`p=2`$. Assume, by way of contradiction, that indeed $`V=0`$. Then $`|M|=|N|^2=2^{2n}`$, say. Let $`x`$ be a generator of $`M`$. Then $`2^nx`$ generates $`N`$ and therefore we have $`q(2^nx)=0`$. As $`b`$ is non-degenerate, $`b(x,x)=k/2^{2n}`$ for some odd integer $`k`$. It follows from the equality
$$q(2x)2q(x)=b(x,x)$$
that $`q(x)=k^{}/2^{n+1}`$ for some odd integer $`k^{}`$. Hence $`q(2^nx)=k^{}/20`$, contradiction. $`\mathrm{}`$
Now consider a non-degenerate pair $`(M,q)`$ which has a metabolizer, and where $`M`$ is a $`p`$-group with $`p1(mod4)`$. We may write $`M`$ as an orthogonal direct sum of pairs $`(M_i,q_i)`$ with $`M_i`$ cyclic (\[Wal63\], Theorem 4). By Lemma 12.3, if $`|M_i|`$ is a square then $`(M_i,q_i)`$ has a metabolizer and $`\gamma (M_i,q_i)=1`$. Since $`|M|`$ is a square (by Remark 10.7 (ii)) and $`\gamma (M,q)=1`$, there must be a even number of $`M_i`$ whose order is not a square and satisfy $`\gamma (M_i,q_i)=1`$, and an even number of $`M_i`$ whose order is not a square and satisfy $`\gamma (M_i,q_i)=1`$. So we see that now we may write (with a change of notation)
$$(M,q)=(M_1,q_1)\mathrm{}(M_r,q_r)$$
(66)
with $`\gamma (M_i,q_i)=1`$ for each $`i`$, and $`M_i`$ is either cyclic of square order or the product of two cyclic groups, each of non-square order.
Before we proceed further with our discussion, let us recall Wall’s nomenclature for non-degenerate symmetric bilinear forms on a finite cyclic group. Let $`G`$ be a cyclic group of order $`p^n`$ with generator $`x`$, where $`p`$ is an odd prime. A bilinear form $`b:G\times G/`$ with $`b(x,x)=ϵp^n`$ is of type A if $`ϵ=1`$, and it is of type B if $`ϵ`$ is a quadratic non-residue modulo $`p`$ (cf. \[Wal63\]).
We assert that each $`(M_i,q_i)`$ in (66) has a cyclic metabolizer. This has already been explained if $`M_i`$ is cyclic. Otherwise, $`M_i=HK`$ with $`H`$, $`K`$ both non-zero cyclic groups and $`\gamma (H,q_i)\gamma (K,q_i)=1`$. Then $`\gamma (H,q_i)=\gamma (K,q_i)`$ and from section 5 of \[Wal63\] we see that $`H`$ and $`K`$ are either both of type A or both of type B. More precisely, there are generators $`x,y`$ of $`H,K`$ respectively (and of order $`p^h`$, $`p^k`$ respectively, say) such that $`b_i(x,x)=ϵp^h`$ and $`b_i(y,y)=ϵp^k`$ where $`ϵ`$ is either 1 or a non-residue $`(modp)`$. (Here we have used the fact that $`\left(\frac{1}{p}\right)=1`$ for $`p1(mod4)`$.) Now if $`h=ru`$ and $`k=r+u`$ for nonnegative integers $`r,u`$ we find that the subgroup generated by $`x+p^uy`$ is the metabolizer we require.
We can rewrite (66) in terms of twisted quantum doubles using Proposition 10.8. As each $`(M_i,q_i)`$ has a cyclic metabolizer $`G_i`$, say, then we know (cf. Remark 11.4) that the pair $`(M_i,q_i)`$ can be realized via $`D^{\omega _i}(G_i)`$ for suitable $`\omega _iZ^3(G,^{})_{ab}`$. That is, $`(M_i,q_i)`$ is equivalent to $`(\mathrm{\Gamma }^{\omega _i},q_{\omega _i})`$. The orthogonal sum (66) corresponds to the tensor product of twisted doubles. So if we start with a pair $`(G,\omega )`$ which gives rise to $`(M,q)`$, we can conclude
###### Theorem 12.4
Suppose that $`p1(mod4)`$ is a prime, $`G`$ an abelian $`p`$-group, and $`\omega Z^3(G,)_{ab}`$. Then there are cyclic $`p`$-groups $`G_i`$ and cocycles $`\omega _iZ^3(G_i,)_{ab}`$, $`1ir`$, such that the two quasi-triangular, quasi-Hopf algebras $`D^\omega (G)`$ and $`\underset{i=1}{\overset{r}{}}D^{\omega _i}(G_i)`$ are gauge equivalent. $`\mathrm{}`$
###### Remark 12.5
It is evident from our proof that the choice of the isomorphism types of the $`G_i`$ is far from unique. For example, suppose that $`M`$ is the orthogonal sum of four cyclic group of orders $`p`$, $`p^3`$, $`p^5`$, $`p^7`$, and each of type A. Then our proof shows that $`D^\omega (G)`$ is equivalent to $`D^{\omega _1}(G_1)D^{\omega _2}(G_2)`$, as quasi-triangular quasi-bialgebras, where $`G_1`$, $`G_2`$ can be chosen to be have orders $`p^2`$, $`p^6`$; $`p^3`$, $`p^5`$; or $`p^4`$, $`p^4`$. Of course, the number $`r`$ of tensor factors that occur depends only on $`(G,\omega )`$, and we always have $`|G|=\underset{i=1}{\overset{r}{}}|G_i|`$.
Theorem 12.4 fails for primes not congruent to $`1(mod4)`$. To treat this case we need
###### Lemma 12.6
Let $`(M,q)`$ be a non-degenerate quadratic space on the $`p`$-group $`M`$. Suppose that $`|M|`$ is a square. Then $`\gamma (M,q)=1`$ if, and only if, $`(M,q)`$ has a metabolizer.
Proof. Let $`N`$ and $`V`$ be as in the proof of Lemma 12.3. After Corollary 12.2 we may assume $`\gamma (M,q)=1`$ and $`V0`$ and try to reach a contradiction.
We see that $`V`$ is not cyclic by Lemma 12.3, and since $`V`$ is anisotropic then $`V_p_p`$, as is easily verified. From Lemma 12.3 (and its proof), we see that $`\gamma (M,q)=1`$ if either $`p`$ is odd, or if $`p=2`$ and $`V`$ is not an orthogonal sum of two cyclic groups. Otherwise, $`p=2`$, $`V=x,y`$, $`b(x,y)=0`$, $`q(x)=q(y)=\pm 1/4`$, $`q(x+y)=1/2`$. Hence $`\gamma (M,q)=\pm i`$ in this case. $`\mathrm{}`$
Now consider the analogue of (66) in the case $`p3(mod4)`$. The arguments preceding (66) together with Lemma 12.3 show that we may still write $`(M,q)`$ in the form of equation (66), with each $`\gamma (M_i,q_i)=1`$, and $`|M_i|`$ a square. However, $`M_i`$ may now be a product of 1, 2 or 4 cyclic subgroups. By Lemma 12.6, each $`(M_i,q_i)`$ has a metabolizer. Thus there is an analogue of Theorem 12.4 for $`p3(mod4)`$, except that the groups $`G_i`$ cannot necessarily be chosen to be cyclic. Rather, we have the property that the group $`\mathrm{\Gamma }^{\omega _i}(G_i)`$ is a product of 1, 2 or 4 cyclic groups.
Finally, take $`p=2`$, and assume again that $`\gamma (M,q)=1`$. We can decompose $`(M,q)`$ as follows (\[Wal63\], section 5,6)
$$(M,q)=(P_1,q_1)\mathrm{}(P_t,q_t)$$
(67)
with each $`P_i`$ homogeneous and the product of at most 2 cyclic groups. From Lemma 12.3, and the proof of Lemma 12.6, we see that the number of cyclic factors in (67) is necessarily even. Hence $`M`$ is the product of an even number of cyclic groups, which thus gives a different proof of Theorem 8.5. The analogue of Theorem 12.4 for $`p=2`$ is the following: $`D^\omega (G)`$ is gauge equivalent to $`\underset{i=1}{\overset{r}{}}D^{\omega _i}(G_i)`$ and for each $`i`$, $`\mathrm{\Gamma }^{\omega _i}(G_i)`$ is the product of 1,2,4 or 8 cyclic groups.
## 13 Metabolic Forms And Homogeneous Groups
A metabolic form is a triple $`(\mathrm{\Gamma },b,G)`$ with $`b:\mathrm{\Gamma }\times \mathrm{\Gamma }/`$ a non-degenerate symmetric bilinear form and $`G`$ a metabolizer of $`(\mathrm{\Gamma },b)`$. If $`\mathrm{\Gamma }^{}G`$ and $`G^{}G`$ such that $`(\mathrm{\Gamma }^{},b,G^{})`$ also a metabolic form, then $`(\mathrm{\Gamma }^{},b,G^{})`$ is called a sub-metabolic form of $`(\mathrm{\Gamma },b,G)`$. Obviously, $`(\{0\},b,\{0\})`$ is a trivial sub-metabolic form of $`(\mathrm{\Gamma },b,G)`$. A metabolic form is called simple if $`(\mathrm{\Gamma },b,G)`$ is the only nontrivial sub-metabolic form.
A group homomorphism $`j:\mathrm{\Gamma }_1\mathrm{\Gamma }_2`$ is called a morphism mapping $`(\mathrm{\Gamma }_1,b_1,G_1)`$ to $`(\mathrm{\Gamma }_2,b_2,G_2)`$ if $`j(G_1)G_2`$ and
$$b_2(j(u),j(v))=b_1(u,v)$$
for any $`u,v\mathrm{\Gamma }_1`$. If the morphism $`j`$ is bijective, then we call the two metabolic forms equivalent and denote this by
$$(\mathrm{\Gamma }_1,b_1,G_1)(\mathrm{\Gamma }_2,b_2,G_2).$$
###### Proposition 13.1
Let $`G`$, $`G^{}`$ be finite abelian groups of odd order, $`\omega Z^3(G,^{})_{ab}`$ and $`\omega ^{}Z^3(G^{},^{})_{ab}`$. Then $`(\mathrm{\Gamma }^\omega (G),b_\omega ,\widehat{G})(\mathrm{\Gamma }^\omega ^{}(G^{}),b_\omega ^{},\widehat{G}^{})`$ if, and only if, $`\omega ^{}`$ and $`\widehat{\rho }w`$ are cohomologous for some isomorphism $`\rho :GG^{}`$ where $`\widehat{\rho }\omega `$ is as in Remark 2.1(iii).
Proof. Let $`\rho :GG^{}`$ be an isomorphism such that $`\widehat{\rho }\omega =\omega ^{}\delta b^{}`$ for some normalized 2-cochain $`b^{}`$ on $`G^{}`$. Let $`\tau T(\omega )`$. Then $`\tau ^{}T(\omega ^{})`$ where
$$\tau _x^{}^{}(y^{})=\tau _{\rho ^1x^{}}(\rho ^1y^{})\frac{b^{}(y^{},x^{})}{b^{}(x^{},y^{})}$$
for any $`x^{},y,G^{}`$. The map $`j:D^\omega (G)D^\omega ^{}(G^{})`$ defined by
$$j:e(g)x\frac{b^{}(\rho g,\rho x)}{b^{}(\rho x,\rho g)}e(\rho g)\rho x$$
is a bialgebra isomorphism. In particular, $`j(\mathrm{\Gamma }^\omega (G))=\mathrm{\Gamma }^\omega ^{}(G^{})`$. Moreover, for any $`\sigma (\alpha ,x)\mathrm{\Gamma }^\omega (G)`$,
$$j(\sigma (\alpha ,x))=\underset{gG}{}\alpha (g)\tau _x(g)\frac{b^{}(\rho g,\rho x)}{b^{}(\rho x,\rho g)}e(\rho g)\rho x$$
and so
$$q_\omega ^{}(j(\sigma (\alpha ,x))=\alpha (x)\tau _x(x)=q_\omega (\sigma (\alpha ,x)).$$
It is obvious that $`j(\widehat{G})=\widehat{G^{}}`$ and hence $`j`$ is an isomorphism of $`(\mathrm{\Gamma }^\omega (G),b_\omega ,\widehat{G})`$ with $`(\mathrm{\Gamma }^\omega ^{}(G^{}),b_\omega ^{},\widehat{G}^{})`$.
Conversely, assume that $`j`$ is an isomorphism of $`(\mathrm{\Gamma }^\omega (G),b_\omega ,\widehat{G})`$ with $`(\mathrm{\Gamma }^\omega ^{}(G^{}),b_\omega ^{},\widehat{G}^{})`$. For $`xG`$,
$$j\sigma _\tau (1,x)=\sigma _\tau ^{}(b_x,\rho x)$$
for some $`b_x\widehat{G}^{}`$ and $`\rho xG`$. As $`j(\widehat{G})=\widehat{G}^{}`$, the map $`\rho :x\rho x`$ defines an isomorphism from $`G`$ to $`G^{}`$. As
$$b_\omega ^{}(\sigma _\tau ^{}(b_x,\rho x),\sigma _\tau ^{}(b_y,\rho y))=b_\omega (\sigma _\tau (1,x),\sigma _\tau (1,y)),$$
we have the equality
$$b_x(y)b_y(x)\tau _{\rho x}^{}(\rho y)\tau _{\rho y}^{}(\rho x)=\tau _x(y)\tau _y(x).$$
(68)
Hence,
$`{\displaystyle \frac{\tau _x(z)\tau _y(z)}{\tau _{xy}(z)}}\omega _z(x,y)`$ $`=`$ $`{\displaystyle \frac{\tau _x(z)\tau _z(x)\tau _y(z)\tau _z(y)}{\tau _{xy}(z)\tau _z(xy)}}`$
$`=`$ $`{\displaystyle \frac{\tau _{\rho x}^{}(\rho z)\tau _{\rho z}^{}(\rho x)\tau _{\rho y}^{}(\rho z)\tau _{\rho z}^{}(\rho y)}{\tau _{\rho (xy)}^{}(\rho z)\tau _{\rho z}^{}(\rho (xy))}}{\displaystyle \frac{b_x(z)b_y(z)}{b_{xy}(z)}}`$
$`=`$ $`{\displaystyle \frac{\tau _{\rho x}^{}(\rho z)\tau _{\rho y}^{}(\rho z)}{\tau _{\rho (xy)}^{}(\rho z)}}\omega _{\rho z}^{}(\rho x,\rho y)(\delta b)(x,y)(z)`$
where $`bC^1(G,\widehat{G})`$ given by $`b(x,y)=b_x(y)`$. By equation (68) and (15), we see that $`\mathrm{\Lambda }(\widehat{\rho }\omega )`$ and $`\mathrm{\Lambda }(\omega ^{})`$ are cohomologous. As $`|G|`$ is odd, it follows from Corollary 6.4 that $`[\widehat{\rho }\omega ]=[\omega ^{}]`$. $`\mathrm{}`$
For two metabolic forms $`(\mathrm{\Gamma }_1,b_1,G_1)`$ and $`(\mathrm{\Gamma }_2,b_2,G_2)`$, one can define the orthogonal sum, $`(\mathrm{\Gamma }_1,b_1,G_1)(\mathrm{\Gamma }_2,b_2,G_2)`$ to be the metabolic form $`(\mathrm{\Gamma }_1\times G_2,b,G_1\times G_2)`$ where
$$b((u_1,v_1),(u_2,v_2))=b_1(u_1,v_1)+b_2(u_2,v_2)$$
for any $`u_1,u_2G_1`$ and $`v_1,v_2G_2`$.
###### Proposition 13.2
Let $`(\mathrm{\Gamma },b,G)`$ be a metabolic form and $`(K,b|_K,H)`$ a sub-metabolic form of $`(\mathrm{\Gamma },b,G)`$. Then, $`(K^{},b|_K^{},GK^{})`$ is also a sub-metabolic form of $`(\mathrm{\Gamma },b,G)`$. Moreover,
$$(\mathrm{\Gamma },b,G)(K,b|_K,H)(K^{},b|_K^{},GK^{}).$$
Proof. As $`b`$ is non-degenerate on $`K`$, there exists a subgroup $`Q<K^{}`$ such that $`K+G=K+Q`$. So for any $`xG`$, $`x=u+v`$ for some $`uK`$ and $`vQ`$. Now $`HG`$ and $`b(v,H)=0`$. This implies $`b(u,H)=0`$. Since $`H`$ is a metabolizer of $`(K,b|_K)`$, $`uH`$, it follows that $`vG`$ and hence $`G=H+(K^{}G)`$. Moreover,
$$(GK^{})^{}K^{}=(G+K)K^{}=(K+(K^{}G))K^{}=GK^{}.$$
Hence $`GK^{}`$ is a metabolizer of $`(K^{},b|_K^{})`$. It is straightforward to show that
$$(K,b|_K,H)(\mathrm{\Gamma },b|_\mathrm{\Gamma },G)(K^{},b|_K^{},GK^{})(\mathrm{\Gamma },b,G)$$
with the morphism $`j:K\times K^{}\mathrm{\Gamma }`$ given by $`j(u,v)=u+v`$ for any $`(u,v)K\times K^{}`$. $`\mathrm{}`$
###### Proposition 13.3
Every metabolic form is equivalent to an orthogonal sum of simple metabolic forms.
Proof. The result follows directly from Proposition 13.2 by induction. $`\mathrm{}`$
###### Remark 13.4
It follows easily from Proposition 13.2 that a metabolic form $`(\mathrm{\Gamma },b,G)`$ is simple if $`G`$ is cyclic.
We will denote by $`(K,b|_K,H)^{}`$ the metabolic form $`(K^{},b|_K^{},K^{}G)`$ for any sub-metabolic form $`(K,b|_K,H)`$ of $`(\mathrm{\Gamma },b,G)`$.
###### Lemma 13.5
Let $`(\mathrm{\Gamma },b,G)`$ be a metabolic form with $`G`$ homogeneous of exponent $`p^n`$. For any non-degenerate cyclic subgroup $`D`$ of $`\mathrm{\Gamma }`$ of order greater than $`p^n`$, there exists a cyclic subgroup $`HG`$ such that $`(H+D,b|_{H+D},H)`$ is a sub-metabolic form of $`(\mathrm{\Gamma },b,G)`$. Moreover, for any generator $`y`$ of $`D`$, there is a generator $`x`$ of $`H`$ satisfying
$$x=x^{}+p^ay$$
for some $`x^{}D^{}`$ where $`|D|=p^{n+a}`$.
Proof. Let $`D`$ be as in the statement of the lemma. Set $`G_0=GD^{}`$. Then $`G/G_0`$ is cyclic since the map $`g+G_0b(g,)`$ defines an embedding from $`G/G_0`$ into $`\widehat{D}`$. If $`|G:G_0|<p^n`$, then $`\mathrm{\Omega }_p(G)G_0`$. Therefore, $`\mathrm{\Omega }_p(G)D`$ is trivial and so is $`GD`$. But this implies $`|D|p^n`$ since
$$D/(GD)(G+D)/G\mathrm{\Gamma }/G\widehat{G}.$$
Hence $`|G:G_0|=p^n`$. Then, $`G=G_0H`$ for some $`HG`$ is isomorphic to $`_{p^n}`$. $`H`$ can be chosen to contain $`GD`$ as $`G`$ is homogeneous. In particular, we have
$$GD=HD.$$
(69)
Set $`K=H+D`$. Then $`G_0K^{}`$. Therefore,
$$K^{}K=(K+K^{})^{}(D+H+G_0)^{}=G^{}D^{}=G_0.$$
As $`GK=H+(GD)=H`$ by equation (69), $`G_0K`$ is trivial. Hence, $`K^{}K=\{0\}`$, i.e. $`(K,b|_K)`$ is non-degenerate. Moreover,
$$H^{}K=H+(H^{}D)=H+((G_0+H)^{}D)=H+(GD),.$$
By (69), $`H^{}K=H`$ and hence $`H`$ is a metabolizer of $`(K,b|_K)`$.
Let $`x=H`$ and $`y=D`$. Then, $`x=x^{}+sp^ay`$ for some integer $`s`$ and $`x^{}D^{}`$. If $`p|s`$, $`b(p^{n1}x,y)=0`$ and so $`p^{n1}xG_0H`$ which contradicts $`\text{ord}(x)=p^n`$. So there is $`s^{}`$ such that $`ss^{}1(modp)`$. The result follows if one replace $`x`$ by $`s^{}x`$. $`\mathrm{}`$
###### Lemma 13.6
Let $`p`$ be an odd prime and $`(\mathrm{\Gamma },b,G_1)`$, $`(\mathrm{\Gamma },b,G_2)`$ metabolic forms with $`G_1`$ and $`G_2`$ isomorphic to $`(_{p^n})^k`$ for some positive integer $`k`$. Then, there exist non-trivial sub-metabolic forms $`(K_i,b|_{K_i},H_i)`$ of $`(\mathrm{\Gamma },b,G_i)`$, $`i=1,2`$ such that
$$(K_1,b|_{K_1},H_1)(K_2,b|_{K_2},H_2).$$
Proof. Pick $`D\mathrm{\Gamma }`$ cyclic of maximal order such that $`b|_D`$ is non-degenerate (cf. \[Wal63\]). Then $`|D|p^n`$. If $`|D|=p^n`$, $`G_1`$ and $`G_2`$ are split metabolizers of $`(\mathrm{\Gamma },b)`$. It is straightforward to show that there exist $`x_iG_i`$ and $`u_i,v_i\mathrm{\Gamma }`$ such that $`b(u_i,u_i)=1/p^n`$, $`b(v_i,v_i)=1/p^n`$, $`b(u_i,v_i)=0`$ and $`x_i=u_i+v_i`$ (cf. \[MH73\]). In particular, $`(K_i,b|_{K_i},x_i)`$ is a sub-metabolic form of $`(\mathrm{\Gamma },b,G_i)`$ where $`K_i`$ is the subgroup generated by $`u_i,v_i`$. Moreover, the map $`K_1K_2`$, $`u_1u_2`$ and $`v_1v_2`$ defines an isomorphism of $`(K_1,b|_{K_1},x_1)`$ with $`(K_2,b|_{K_2},x_2)`$.
Let $`y`$ be a generator of $`D`$ and $`\text{ord}(y)=p^{n+a}`$ for some positive integer $`a`$. By Lemma 13.5, for $`i=1,2`$ there exists $`H_iG_i`$ such that $`(K_i,b|_{K_i},H_i)`$ is a sub-metabolic form of $`(\mathrm{\Gamma },b,G_i)`$ where $`K_i=H_i+D`$. Moreover $`H_i`$ admits a generator $`x_i`$ satisfying
$$x_i=x_i^{}+p^ay$$
for some $`x_i^{}D^{}`$. The map $`K_1K_2`$, given by $`yy`$ and $`x_1^{}x_2^{}`$, defines an isomorphism of $`(K_1,b|_{K_1},H_1)`$ with $`(K_2,b|_{K_2},H_2)`$. $`\mathrm{}`$
###### Theorem 13.7
Let $`p`$ be an odd prime and $`(\mathrm{\Gamma },b,G)`$ a metabolic form with $`G(_{p^n})^k`$.
1. $`(\mathrm{\Gamma },b,G)`$ is equivalent to an orthogonal sum of sub-metabolic forms with metabolizers isomorphic to $`_{p^n}`$.
2. Let $`(\mathrm{\Gamma }^{},b^{},G^{})`$ be a metabolic form such that $`G^{}G`$. Then $`(\mathrm{\Gamma }^{},b^{},G^{})(\mathrm{\Gamma },b,G)`$ if, and only if $`(\mathrm{\Gamma }^{},b^{})`$ and $`(\mathrm{\Gamma },b)`$ are equivalent.
Proof. Statement (i) follows easily from Lemma 13.6 and Proposition 13.2 by induction on $`k`$. For statement (ii), it is obvious that the equivalence of two metabolic form implies the equivalence of the underlying bilinear forms. Conversely, we proceed by induction on $`k`$. Let $`j`$ be an isomorphism of $`(\mathrm{\Gamma }^{},b^{},G^{})`$ with $`(\mathrm{\Gamma },b,G)`$. Then, $`(\mathrm{\Gamma },b,j(G^{}))`$ is a metabolic form. By Lemma 13.6, there exist a nontrivial sub-metabolic form $`(K^{},b^{}|_K^{},H^{})`$ of $`(\mathrm{\Gamma }^{},b^{},G^{})`$ and a sub-metabolic form $`(K,b|_K,H)`$ of $`(\mathrm{\Gamma },b,G^{})`$ such that
$$(K^{},b|_K^{},H^{})(K,b|_K,H).$$
Thus, $`(K_{}^{}{}_{}{}^{},b^{}|_K_{}^{}{}_{}{}^{})`$ and $`(K^{},b|_K^{})`$ are equivalent bilinear form. By induction assumption $`(K^{},b|_K^{},H^{})^{}(K,b|_K,H)^{}`$ and hence
$$(\mathrm{\Gamma }^{},b^{},G,)(K^{},b|_K^{},H^{})(K^{},b|_K^{},H^{})^{}(K,b|_K,H)(K,b|_K,H)^{}(\mathrm{\Gamma },b,G).$$
$`\mathrm{}`$
###### Theorem 13.8
Let $`p`$ be an odd prime with $`G`$ isomorphic to $`(_{p^n})^k`$.
1. For any $`\omega Z^3(G,^{})_{ab}`$, there exists $`\eta _1,\mathrm{},\eta _kZ^3(_{p^n},^{})_{ab}`$ such that $`D^\omega (G)`$ is gauge equivalent to $`\underset{i=1}{\overset{k}{}}D^{\eta _i}(_{p^n})`$ as quasi-triangular quasi-bialgebras.
2. For any $`\omega ,\omega ^{}Z^3(G,^{})_{ab}`$, $`D^\omega (G)`$ and $`D^\omega ^{}(G)`$ are gauge equivalent as quasi-triangular quasi-bialgebras if, and only if, there exists $`\rho \text{Aut}(G)`$ such that $`[\widehat{\rho }\omega ]=[\omega ^{}]`$.
Proof. (i) Consider a metabolic form $`(\mathrm{\Gamma }^\omega ,b_\omega ,\widehat{G})`$ associated to $`\omega `$. By Theorem 13.7, $`(\mathrm{\Gamma }^\omega ,b_\omega ,\widehat{G})`$ equivalent to the orthogonal sum
$$(\mathrm{\Gamma }_1,b_1,_{p^n})\mathrm{}(\mathrm{\Gamma }_k,b_k,_{p^n}).$$
(70)
By Theorem 11.3 (iv), there exist $`\eta _iZ^3(_{p^n},^{})`$ such that $`(\mathrm{\Gamma }_i,b_i)`$ is equivalent to $`(\mathrm{\Gamma }^{\eta _i},b_{\eta _i})`$. Hence, by equation (70), the result follows.
(ii) If $`D^\omega (G)`$ and $`D^\omega ^{}(G)`$ are equivalent as quasi-triangular quasi-bialgebras, by Theorem 10.4, $`(\mathrm{\Gamma }^\omega ,b_\omega )`$ and $`(\mathrm{\Gamma }^\omega ^{},b_\omega ^{})`$ are equivalent bilinear forms. It follows from Theorem 13.7(ii) that $`(\mathrm{\Gamma }^\omega ,b_\omega ,\widehat{G})`$ and $`(\mathrm{\Gamma }^\omega ^{},b_\omega ^{},\widehat{G})`$ are equivalent metabolic forms. By Proposition 13.1, $`[\widehat{\rho }\omega ]=[\omega ^{}]`$ for some $`\rho \text{Aut}(G)`$. The “only if” part follows from Remark 2.1 (iii). $`\mathrm{}`$
## 14 Duality and Symmetry
In this section we discuss various kinds of relations that exist between suitably chosen $`D^\omega (G)`$ and $`D^\omega ^{}(G^{})`$. Our discussion is meant to illustrate the possibilities, and is by no means exhaustive. We make use of Wall’s classification of symmetric bilinear forms \[Wal63\] as well as results previously established in the present paper.
The first type relation we discuss–symmetry–arises from the possibility that a quadratic space $`(\mathrm{\Gamma },q)`$ may have metabolizers $`G`$, $`G^{}`$ which are not conjugate in the corresponding orthogonal group. Indeed, they may not even be isomorphic. This is, in fact, a rather common phenomena.
As an example, take $`\mathrm{\Gamma }_{p^2}\times _{p^2}`$ ( and for convenience, $`p`$ an odd prime). By \[Wal63\], the possible non-degenerate bilinear forms on $`\mathrm{\Gamma }`$ can be taken to be the following: $`A_{p^2}A_{p^2}`$, $`A_{p^2}B_{p^2}`$. Of these, only one has a (necessarily split) cyclic metabolizer: the first if $`p1(mod4)`$ and the second if $`p3(mod4)`$. On the other hand, the subgroup of $`\mathrm{\Gamma }`$ consisting of elements of order at most $`p`$ is a metabolizer in all cases. As we know, each of these metabolizers $`G`$ determines an abelian 3-cocycle $`\omega Z^3(G,^{})_{ab}`$ such that $`\mathrm{\Gamma }^\omega \mathrm{\Gamma }`$, and $`\omega `$ is even a coboundary if $`G`$ is cyclic. More is true: if they determine the same pair $`(\mathrm{\Gamma },q)`$ then the corresponding module categories are braided monoidally equivalent. We thus conclude:
###### Example 14.1
Let $`p`$ be an odd prime, $`G_{p^2}`$, $`G^{}_p\times _p`$. Then there is an abelian 3-cocycle $`\omega ^{}Z^3(G^{},^{})_{ab}`$ such that $`D(G)\text{-}\text{Mod}`$ and $`D^\omega ^{}(G^{})\text{-}\text{Mod}`$ are equivalent as braided tensor categories. The same argument applies if we take $`G^{}_{p^n}\times _{p^n}`$ and $`G_{p^{2n}}`$, and to a host of other situations.
Next we discuss a duality between the module categories corresponding to the twisted double of homogeneous $`p`$-groups.
###### Theorem 14.2
Let $`n,k`$ be a pair of positive integers, and let $`G\left(_{p^n}\right)^k`$ and $`G^{}\left(_{p^k}\right)^n`$ ($`p`$ an odd prime). Then the following hold:
1. There are exactly $`\left(\genfrac{}{}{0pt}{}{n+k}{k}\right)`$ equivalence classes of monoidal categories of the form $`D^\omega (G)\text{-}\text{Mod}`$, for some $`\omega Z^3(G,^{})ab`$.
2. There is a canonical bijection between equivalence classes of monoidal categories of the form $`D^\omega (G)\text{-}\text{Mod}`$ and those of the form $`D^\omega ^{}(G^{})\text{-}\text{Mod}`$.
3. There are natural bijections between equivalence classes of braided monoidal categories of the form $`D^\omega (G)\text{-}\text{Mod}`$ and those of the form $`D^\omega ^{}(G^{})\text{-}\text{Mod}`$.
###### Example 14.3
One can also count the number of braided monoidal categories (up to equivalence) of the form $`D^\omega (G)\text{-}\text{Mod}`$, where $`G=\left(_{p^n}\right)^k`$ and $`\omega Z^3(G,^{})_{ab}`$. Thus we have
| $`k`$ | \# braided monoidal categories |
| --- | --- |
| $`1`$ | $`(2n+1)`$ |
| $`2`$ | $`2n^2+2n+1`$ |
| $`3`$ | $`1/3(4n^3+6n^2+8n+3)`$ |
| $`\mathrm{}`$ | $`\mathrm{}`$ |
To prove these assertions, recall from Theorem 13.7 that if $`G=\left(_{p^n}\right)^k`$ is homogeneous, $`p`$ odd, then the pair $`(G,\omega )`$, $`\omega Z^3(G,^{})_{ab}`$, determines a non-degenerate bilinear form $`(\mathrm{\Gamma },b)`$ of the form
$$(\mathrm{\Gamma },b)=(H_1,b_1)\mathrm{}(H_k,b_k)$$
(71)
where $`(H_i,b_i)`$ has a cyclic metabolizer $`G_i`$, $`G=G_1\times \mathrm{}\times G_k`$, and $`H_i`$ is isomorphic to $`_{p^{n+a}}_{p^{na}}`$ for some integer $`a`$ with $`0an`$. Of course, $`\mathrm{\Gamma }`$ is isomorphic to the group of fusion rules $`\mathrm{\Gamma }^\omega `$. Denote by $`\mathrm{\Gamma }_a`$ the group $`_{p^{n+a}}_{p^{na}}`$ (where now $`n`$ is fixed). If $`\mathrm{\Gamma }_a`$ occurs with multiplicity $`m_a`$ in (71), we can represent (71) symbolically in the form
$$\mathrm{\Gamma }=_{a=0}^nm_a\mathrm{\Gamma }_a$$
(72)
at least as far as the fusion rules are concerned. It is easy to see that we have
$$\underset{a=0}{\overset{n}{}}m_a=k.$$
(73)
Furthermore, any $`n+1`$-tuple $`(m_0,m_1,\mathrm{},m_n)`$ of nonnegative integers satisfying (73) may be chosen in (72) and which correspond to a pair $`(\mathrm{\Gamma },b)`$ with metabolizer $`G`$. Thus the number of non-isomorphic groups of fusion rules $`\mathrm{\Gamma }`$ that correspond to $`G`$ is the number of non-negative $`(n+1)`$-tuples $`(m_0,m_1,\mathrm{},m_n)`$ satisfying (73). This is easily seen to equal $`\left(\genfrac{}{}{0pt}{}{n+k}{k}\right)`$. Now part (i) of the Theorem follows from Theorem 9.4.
Note that if also $`G^{}=\left(_{p^k}\right)^n`$ then by part (i), there are $`\left(\genfrac{}{}{0pt}{}{n+k}{k}\right)`$ equivalence classes of tensor categories both of the type $`D^\omega (G)\text{-}\text{Mod}`$ and $`D^\omega ^{}(G^{})\text{-}\text{Mod}`$. We will establish a canonical bijection between these two sets.
With $`(m_0,m_1,\mathrm{},m_n)`$ as above, let $`\lambda `$ be the partition which contains the integer $`i`$ with multiplicity $`m_i`$, $`1in`$. Thus $`\lambda `$ is a partition of the integer $`r=_{i=1}^nim_i`$. Note that $`rnk`$ by (73). Let $`\lambda ^t`$ be the dual (conjugate) partition of $`\lambda `$, and let $`s_j`$ be the multiplicity of $`j`$ in $`\lambda ^t`$ for $`j1`$. So also $`r=_{j=1}^kjs_j`$, the upper limit $`k`$ in the sum arising from the fact that the maximal part of $`\lambda ^t`$ (i.e., the largest $`j`$ with $`s_i>0`$) is equal to the number of non-zero multiplicities $`m_i`$, and this is at most $`k`$ by (73).
Now observe that $`_{j=1}^ks_j=`$ \# of parts of the partition $`\lambda ^t=\mathrm{max}\{i|m_i>0\}n`$. Thus if we define $`s_0=n_{j=1}^ks_j`$ then $`s_j0`$ for $`j0`$ and
$$\underset{j=0}{\overset{k}{}}s_j=n.$$
(74)
Consider
$$\mathrm{\Gamma }^{}=_{j=0}^ks_j\mathrm{\Gamma }_j^{}$$
(75)
where $`\mathrm{\Gamma }_j^{}=_{p^{k+j}}_{p^{kj}}`$. By previous arguments, $`\mathrm{\Gamma }^{}`$ is the group of fusion rules corresponding to a suitable pair $`(G^{},\omega ^{})`$. Thus (74) and (75) enjoy the same relation to $`(G^{},\omega ^{})`$ that (72) and (73) do to $`(G,\omega )`$.
Thus the canonical bijection (duality) which relates $`\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }^{}`$ is achieved via the duality between the partitions $`\lambda `$ and $`\lambda ^t`$ which are associated to them. This explains part (ii) of the theorem.
###### Example 14.4
In Table 1 below, we illustrate this duality in case the two groups in question are $`G=_{p^2}\times _{p^2}\times _{p^2}`$ and $`G^{}=_{p^3}\times _{p^3}`$. Groups $`\mathrm{\Gamma }`$, $`\mathrm{\Gamma }^{}`$ on the same line are in duality.
| $`G=_{p^2}\times _{p^2}\times _{p^2}`$ | | $`G^{}=_{p^3}\times _{p^3}`$ | | |
| --- | --- | --- | --- | --- |
| $`(m_1,m_2)`$ | $`\mathrm{\Gamma }`$ | $`(s_1,s_2,s_3)`$ | $`\mathrm{\Gamma }^{}`$ | $`f`$ |
| $`(0,3)`$ | $`_{p^4}\times _{p^4}\times _{p^4}`$ | $`(0,0,2)`$ | $`_{p^6}\times _{p^6}`$ | 1 |
| $`(1,2)`$ | $`_{p^4}\times _{p^4}\times _{p^3}\times _p`$ | $`(0,1,1)`$ | $`_{p^6}\times _{p^5}\times _p`$ | 2 |
| $`(0,2)`$ | $`_{p^4}\times _{p^4}\times _{p^2}\times _{p^2}`$ | $`(0,2,0)`$ | $`_{p^5}\times _{p^5}\times _p\times _p`$ | 1 |
| $`(2,1)`$ | $`_{p^4}\times _{p^3}\times _{p^3}\times _p\times _p`$ | $`(1,0,1)`$ | $`_{p^6}\times _{p^4}\times _{p^2}`$ | 2 |
| $`(1,1)`$ | $`_{p^4}\times _{p^3}\times _{p^2}\times _{p^2}\times _p`$ | $`(1,1,0)`$ | $`_{p^5}\times _{p^4}\times _{p^2}\times _p`$ | 2 |
| $`(0,1)`$ | $`_{p^4}\times _{p^2}\times _{p^2}\times _{p^2}\times _{p^2}`$ | $`(0,2,0)`$ | $`_{p^4}\times _{p^4}\times _{p^2}\times _{p^2}`$ | 1 |
| $`(2,0)`$ | $`_{p^3}\times _{p^3}\times _{p^2}\times _{p^2}\times _p\times _p`$ | $`(0,1,0)`$ | $`_{p^5}\times _{p^3}\times _{p^3}\times _p`$ | 1 |
| $`(3,0)`$ | $`_{p^3}\times _{p^3}\times _{p^3}\times _p\times _p\times _p`$ | $`(0,0,1)`$ | $`_{p^6}\times _{p^3}\times _{p^3}`$ | 1 |
| $`(1,0)`$ | $`_{p^3}\times _{p^2}\times _{p^2}\times _{p^2}\times _{p^2}\times _p`$ | $`(1,0,0)`$ | $`_{p^4}\times _{p^3}\times _{p^3}\times _{p^2}`$ | 1 |
| $`(0,0)`$ | $`\left(_{p^2}\right)^6`$ | $`(0,0,0)`$ | $`\left(_{p^3}\right)^4`$ | 0 |
Table: 1
Turning to part (iii), we first establish:
###### Lemma 14.5
With the notation of (72), let $`f`$ be the number of indices $`a1`$ such that $`m_a1`$. Then the number of inequivalent, non-degenerate bilinear forms $`b`$ on $`\mathrm{\Gamma }`$ with metabolizer (isomorphic to) $`G`$ is $`2^f`$.
Proof. From Theorem 13.7, we only need to enumerate the equivalence classes of bilinear forms that arise as in (71) and (72) above.
To show that this number is $`2^f`$, let us assume for convenience that $`p1(mod4)`$ (the case $`p3(mod4)`$ is proved similarly). According to \[Wal63\] and the discussion following equation (66), the only forms $`\mathrm{\Gamma }_a`$ with cyclic metabolizer are of type $`AA`$ or type $`BB`$ if $`1an1`$; type $`A`$ or type $`B`$ if $`a=n`$; type $`AA`$ ($``$ type $`BB`$) if $`a=0`$. Furthermore, if $`a1`$ the only forms on $`m_a\mathrm{\Gamma }_a`$ with metabolizer $`\left(_{p^n}\right)^{m_a}`$ are (up to equivalence) of the shape $`m_a(AA)`$ or $`(m_a1)(AA)(BB)`$.Thus, as long as $`m_a1`$ and $`a1`$, there are in any case just two inequivalent forms on $`m_a\mathrm{\Gamma }_a`$ with homogeneous metabolizer $`\left(_{p^n}\right)^{m_a}`$. Now the lemma follows immediately. $`\mathrm{}`$
Note that the integer $`f`$ of Lemma 14.5 is the number of unequal parts of the partition $`\lambda `$. But it is easy to see (and well-known) that this is equal also to the number of unequal parts of $`\lambda ^t`$. So the canonical bijection of part (ii) of Theorem 14.2 also induces bijections between the inequivalent braided monoidal categories of the form $`D^\omega (G)\text{-}\text{Mod}`$ and $`D^\omega ^{}(G^{})\text{-}\text{Mod}`$ which have fusion rules $`\mathrm{\Gamma }`$, $`\mathrm{\Gamma }^{}`$ respectively. This holds for all such $`\mathrm{\Gamma }`$, so that the number of inequivalent braided tensor categories $`D^\omega (G)\text{-}\text{Mod}`$ is the same as the numbers of type $`D^\omega ^{}(G^{})\text{-}\text{Mod}`$.
This is not quite a natural bijection. However, we saw above that the allowable bilinear forms on $`\mathrm{\Gamma }`$ are naturally indexed by $`f`$-tuples consisting of $`A`$’s and $`B`$’s. More precisely, let us again assume for convenience that $`p1(mod4)`$. If $`a1`$ with $`m_a1`$ then we may take all but one $`\mathrm{\Gamma }_a`$ in (72) to be type $`A`$ or $`AA`$ (according to whether $`a=n`$ or $`a<n`$); the remaining $`\mathrm{\Gamma }_a`$ may be taken to be either type $`A`$ or type $`B`$ (or type $`AA`$ or $`BB`$). All such $`f`$-tuples are allowable, so that there are exactly $`2^f`$ such tuples, as claimed. Since (by the lemma) this analysis applies not only to $`\mathrm{\Gamma }`$ but to the dual group $`\mathrm{\Gamma }^{}`$, we obtain a “natural” bijection between bilinear forms, and hence braided monoidal categories, by associating those forms which correspond to the same $`f`$-tuple of $`A`$’s and $`B`$’s. This completes our discussion of part (iii) of the theorem.
Finally, the formulae in example 14.3 follow from our analysis without difficulty. For example, take the case $`k=2`$, so that $`G_{p^n}_{p^n}`$. The possible $`(n+1)`$-tuples $`(m_0,m_1,\mathrm{},m_n)`$ with $`_{a0}m_a=2`$ are trivially enumerated, and each give rise to $`2^f`$ inequivalent forms as in Lemma 14.5. If $`m_0=2`$ then $`f=0`$; if $`m_0=1`$ or $`m_0=0`$ and $`m_a=2`$ for some $`a1`$ then $`f=1`$; otherwise $`f=2`$. Thus the total number of inequivalent allowable forms (or braided monoidal categories) is $`1+2n+2n+4\left(\genfrac{}{}{0pt}{}{n}{2}\right)=2n^2+2n+1`$.
###### Example 14.6
We can read off further examples of symmetry from table 1. Consider the group $`\mathrm{\Gamma }=_{p^4}\times _{p^4}\times _{p^2}\times _{p^2}`$: we see that the two groups $`G=_{p^2}\times _{p^2}\times _{p^2}`$ and $`G^{}=_{p^3}\times _{p^3}`$ each admit $`\mathrm{\Gamma }`$ as the corrresponding fusion rules for some choice of quadratic form. Indeed, the previous discussion shows that there are exactly two non-degenerate quadratic forms on $`\mathrm{\Gamma }`$ with metabolizers equal to both $`G`$ and $`G^{}`$. By Theorem 10.4 it follows that the following holds: there are cohomology classes $`[\omega _1]`$, $`[\omega _2]`$ in $`H^3(G,^{})_{ab}`$ and $`[\omega _1^{}]`$, $`[\omega _2^{}]`$ in $`H^3(G^{},^{})_{ab}`$ such that $`D^{\omega _i}(G)`$ and $`D^{\omega _i^{}}(G^{})`$, $`i=1,2`$, are gauge equivalent.
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# Bloch’s Conjecture and Chow Motives
Bloch’s conjecture, Chow motive, Albanese kernel, Murre’s decomposition
Let $`X`$ be a connected smooth projective complex surface. J. Murre constructed a decomposition of Chow motives for $`X`$, i.e. there exist mutually orthogonal idempotents $`\pi _i\text{CH}^2(X\times X)_{}`$ as correspondences for $`0i4`$ such that $`_i\pi _i`$ is equal to the diagonal and the action of $`\pi _i`$ on $`H^j(X,)`$ is the identity for $`i=j`$, and vanishes otherwise. The decomposition is not uniquely characterized by the above properties. See also . The theory of Chow motives would be rather complicated if the following condition is not satisfied:
(0.1) An idempotent of $`\text{CH}^2(X\times X)_{}`$ is zero if so is its cohomology class.
Here we can also consider a stronger condition:
(0.2) $`\text{Ker}(\text{CH}^2(X\times X)_{}H^4(X\times X,)(2))`$ is a nilpotent ideal.
This is related to Beilinson’s conjectures . If (0.2) is true, the uniqueness of the projectors modulo inner automorphisms can be proved due to Beilinson. See \[7, 7.3\].
Let $`h^i(X)`$ denote the ‘image’ of the projector $`\pi _i`$ in the motivic sense. Then $`h^2(X)`$ carries the Albanese kernel and the Neron-Severi group as well as the transcendental cycles. This is compatible with a conjecture of S. Bloch that the vanishing of the transcendental cycles (i.e. that of $`p_g)`$ would be equivalent to:
(0.3) The Albanese map is injective.
Note that the noninjectivity of the Albanese map in the case $`p_g0`$ is a theorem of D. Mumford , and Bloch’s conjecture is proved at least if $`X`$ is not of general type . See also , for some more cases.
In the second projector $`\pi _2`$ is actually defined as the difference between the diagonal and the sum of the other projectors $`\pi _i`$. In the case $`p_g=0`$, we can construct explicitly a projector $`\stackrel{~}{\pi }_2`$ which is homologically equivalent to $`\pi _2`$ and is orthogonal to the other projectors $`\pi _i`$ $`(i2)`$. Then we can prove Bloch’s conjecture if $`\stackrel{~}{\pi }_2`$ coincides with $`\pi _2`$ (modulo rational equivalence). So the conjecture is reduced to (0.1). Actually, we can show
###### {\bf0.4.~Theorem}
In the case $`p_g=0,`$ the above three conditions (0.1–3) are all equivalent, and the cube of the ideal in (0.2) is zero if it is nilpotent.
The proof of (0.3) $`(0.2)`$ uses an argument similar to together with the bijectivity of the cycle map in the divisor case. Combined with the above mentioned result of , it implies
###### {\bf0.5.~Theorem}
If $`p_g=0`$ and $`X`$ is not of general type (or, if $`X`$ is as in , ), then the cube of the ideal in (0.2) is zero.
We can show that the square of the ideal in (0.2) does not vanish if the irregularity $`q(=dim\mathrm{\Gamma }(X,\mathrm{\Omega }_X^1))`$ is nonzero.
In Sect. 1, we review Murre’s decomposition of Chow motives, and prove (0.1) $``$ (0.3). In Sect. 2 we show (0.3) $``$ (0.2) using a variant of the construction of .
Part of this work was done during my stay at the university of Leiden. I would like to thank J.P. Murre for useful discussions on Chow motives that have originated this work. I thank also the staff of the institute for the hospitality.
1. Chow motives
1.1. Correspondences. For smooth proper complex algebraic varieties $`X,Y`$ such that $`X`$ is purely $`n`$-dimensional, we define the group of correspondences with rational coefficients by
$$C^i(X,Y)_{}=\text{CH}^{n+i}(X\times Y)_{}.$$
For $`\xi C^i(X,Y)_{}`$ and $`\eta C^j(Y,Z)_{}`$, the composition is denoted by $`\eta \xi C^{i+j}(X,Z)_{}`$. For $`\zeta \text{CH}^i(X)_{}`$, let
$$\mathrm{\Gamma }_\zeta C^i(pt,X)_{}(=\text{CH}^i(X)_{})$$
$`(\mathrm{1.1.1})`$
be the element defined by $`\zeta `$. For a morphism $`f:XY`$, we denote by $`\mathrm{\Gamma }_f`$ the graph of $`f`$ which belongs to $`C^0(Y,X)`$. Sometimes we will use the notation
$$f^{}=\mathrm{\Gamma }_f,f_{}={}_{}{}^{t}\mathrm{\Gamma }_{f}^{}.$$
$`(\mathrm{1.1.2})`$
Assume $`X,Y`$ connected. By Hodge theory together with the Künneth decomposition and the duality, we have a canonical isomorphism
$$\frac{\text{CH}^1(X\times Y)_{}}{pr_1^{}\text{CH}^1(X)_{}+pr_2^{}\text{CH}^1(Y)_{}}=\text{Hom}_{\text{HS}}(H^{2n1}(X,)(n1),H^1(Y,)),$$
$`(\mathrm{1.1.3})`$
where the right-hand side is the group of morphisms of Hodge structures. See also , . Let $`\xi \text{CH}_0(X)_{},\xi ^{}\text{CH}_0(Y)_{}`$ with degree one. Then the left-hand side of (1.1.3) is isomorphic to
$$\{\mathrm{\Gamma }C^{1n}(X,Y)_{}|\mathrm{\Gamma }\mathrm{\Gamma }_\xi =0\text{and }{}_{}{}^{t}\mathrm{\Gamma }_{\xi ^{}}^{}\mathrm{\Gamma }=0\}.$$
$`(\mathrm{1.1.4})`$
1.2. Murre’s construction. Let $`X`$ be a connected smooth projective variety of dimension $`n2`$. We choose and fix an embedding of $`X`$ into a projective space. Let $`l`$ denote the multiplication by the hyperplane section class. Let $`C`$ be the intersection of $`n1`$ generic smooth hyperplane sections. (Note that $`[C]\text{CH}^{n1}(X)`$ is independent of the choice of $`C`$.) By (1.1.3–4) there exists uniquely $`\mathrm{\Gamma }\text{CH}^1(X\times X)_{}`$ such that
$`\mathrm{\Gamma }\mathrm{\Gamma }_\xi `$ $`=0,{}_{}{}^{t}\mathrm{\Gamma }_{\xi ^{}}^{}\mathrm{\Gamma }=0,`$ $`\mathrm{1.2.1}`$$`\mathrm{1.2.2}`$
$`\mathrm{\Gamma }_{}l^{n1}`$ $`=id\text{on }H^1(X,).`$
We have $`{}_{}{}^{t}\mathrm{\Gamma }=\mathrm{\Gamma }`$ if $`\xi =\xi ^{}`$. Note that (1.2.1) implies
$$\mathrm{\Gamma }_{}:H^{i+2n2}(X,)(n1)H^i(X,)\text{vanishes for }i1.$$
$`(\mathrm{1.2.3})`$
Let $`i:CX`$ denote the inclusion morphism. Let
$$\pi ^{}=\mathrm{\Gamma }i_{}i^{},$$
where $`i_{}={}_{}{}^{t}\mathrm{\Gamma }_{i}^{},i^{}=\mathrm{\Gamma }_i`$ as in (1.1.2). Following we define
$$\pi _0=\mathrm{\Gamma }_{[X]}{}_{}{}^{t}\mathrm{\Gamma }_{\xi ^{}}^{},\pi _1=\pi ^{}(1{}_{}{}^{t}\pi _{}^{}/2),\pi _{2n1}={}_{}{}^{t}\pi _{1}^{},\pi _{2n}={}_{}{}^{t}\pi _{0}^{},$$
where $`1`$ denotes the diagonal of $`X`$. Then
$$(\pi _i)_{}|H^j(X,)=\delta _{i,j}id\text{for }i=0,1,2n1,2n.$$
$`(\mathrm{1.2.4})`$
If $`n=2`$, we define $`\pi _2=1_{i2}\pi _i`$.
###### {\bf1.3.~Theorem}
(Murre ). $`\pi _i\pi _j=\delta _{i,j}\pi _i`$ for $`\{i,j\}\{0,1,2n1,2n\}`$.
Outline of proof. We recall here some arguments of the proof which will be needed in the proof of the main theorems. See , for details.
We have $`\mathrm{\Gamma }i_{}i^{}\mathrm{\Gamma }=\mathrm{\Gamma }`$ by (1.1.3–4), and $`{}_{}{}^{t}\mathrm{\Gamma }\mathrm{\Gamma }=0`$ in $`\text{CH}^{2n}(X\times X)_{}`$, because it is cohomologically zero by (1.2.3). So we get
$$\pi ^2=\pi ^{},{}_{}{}^{t}\pi _{}^{}\pi ^{}=0.$$
$`(\mathrm{1.3.1})`$
Then we can verify
$$\pi _1^2=\pi _1,\pi _{2n1}\pi _1=0,\pi _1\pi _{2n1}=0.$$
$`(\mathrm{1.3.2})`$
We have furthermore
$$\pi _0\pi ^{}=\pi _{2n}\pi ^{}=\pi ^{}\pi _0=\pi ^{}\pi _{2n}=0$$
$`(\mathrm{1.3.3})`$
Indeed, $`\pi _0\pi ^{}=0`$ by (1.2.1), and $`\pi ^{}\pi _0=\pi _{2n}\pi ^{}=0`$, because $`\pi ^{}\mathrm{\Gamma }_{[X]}C^0(pt,X)`$ is cohomologically zero (and similarly for $`{}_{}{}^{t}\mathrm{\Gamma }_{[X]}^{}\pi ^{})`$. Finally the vanishing of $`\pi ^{}\pi _{2n}`$ follows from $`i_{}i^{}\mathrm{\Gamma }_\xi ^{}C^{2n1}(pt,X)=0`$.
Then we can verify the remaining assertions.
Remark. The Albanese map $`\text{CH}_0(X)_{}^0\text{Alb}_X()_{}`$ induces an isomorphism
$$(\pi _{2n1})_{}\text{CH}_0(X)_{}^0\stackrel{}{}\text{Alb}_X()_{}.$$
If $`n=dimX=2`$, $`(\pi _2)_{}\text{CH}_0(X)_{}^0`$ coincides with the kernel of the Albanese map with $``$-coefficients. See \[7, 7.1\].
1.4. Construction of $`\stackrel{~}{\pi }_2`$. Assume $`n=dimX=2`$ and $`p_g=0`$. Let $`C_i`$ be irreducible curves on $`X`$ such that the cohomology classes of $`[C_i]`$ form a basis of $`H^2(X,)(1)`$. Let $`\stackrel{~}{C}_i`$ be the normalization of $`C_i`$, and $`\stackrel{~}{C}`$ the disjoint union of $`\stackrel{~}{C}_i`$ with $`\stackrel{~}{i}:\stackrel{~}{C}X`$ the canonical morphism. Let $`A=(A_{i,j})`$ be the intersection matrix of the $`C_i`$ (i.e. $`A_{i,j}=C_iC_j)`$. Let $`B=(B_{i,j})`$ be the inverse of $`A`$.
Let $`\mathrm{\Gamma }_BC^1(\stackrel{~}{C},\stackrel{~}{C})_{}=\text{CH}^0(\stackrel{~}{C}\times \stackrel{~}{C})_{}`$ defined by the matrix $`B`$. Let
$$\stackrel{~}{\mathrm{\Gamma }}=\stackrel{~}{i}_{}\mathrm{\Gamma }_B\stackrel{~}{i}^{}.$$
Since the composition
$$\stackrel{~}{i}^{}\stackrel{~}{i}_{}:H^0(\stackrel{~}{C},)H^2(\stackrel{~}{C},)(1)$$
is given by the matrix $`A`$ (using the projection formula), we see that
$$\stackrel{~}{\mathrm{\Gamma }}_{}|H^i(X,)=\delta _{i,2}id.$$
$`(\mathrm{1.4.1})`$
(The assertion for $`i2`$ follows from the definition of $`\stackrel{~}{\mathrm{\Gamma }}`$ .) Note that $`{}_{}{}^{t}\stackrel{~}{\mathrm{\Gamma }}=\stackrel{~}{\mathrm{\Gamma }}`$ because $`B`$ is symmetric. We define
$$\stackrel{~}{\pi }_2=(1\pi ^{})\stackrel{~}{\mathrm{\Gamma }}(1{}_{}{}^{t}\pi _{}^{})$$
Then the symmetry $`{}_{}{}^{t}\stackrel{~}{\pi }_{2}^{}=\stackrel{~}{\pi }_2`$ is clear.
###### {\bf1.5.~Proposition}
$`\stackrel{~}{\pi }_2`$ is an idempotent, and is orthogonal to $`\pi _i`$ for $`i2`$.
###### Demonstration
Proof.
We have $`\pi ^{}\stackrel{~}{\pi }_2=0`$ by (1.3.1). Since
$$\mathrm{\Gamma }_B\stackrel{~}{i}^{}\mathrm{\Gamma }C^2(X,\stackrel{~}{C})_{}=\text{CH}^0(X\times \stackrel{~}{C})_{}$$
is cohomologically zero by (1.2.3), we get
$$\stackrel{~}{\mathrm{\Gamma }}\pi ^{}=0,\stackrel{~}{\pi }_2\pi ^{}=0,$$
$`(\mathrm{1.5.1})`$
using (1.3.1). Then we have $`{}_{}{}^{t}\pi _{}^{}\stackrel{~}{\pi }_2=\stackrel{~}{\pi }_2{}_{}{}^{t}\pi _{}^{}=0`$ by transpose, and
$$\pi _i\stackrel{~}{\pi }_2=\stackrel{~}{\pi }_2\pi _i=0\text{for }i=1,3.$$
$`(\mathrm{1.5.2})`$
For $`i=0,4`$, we have $`\stackrel{~}{\mathrm{\Gamma }}\pi _0=\stackrel{~}{\mathrm{\Gamma }}\pi _4=0`$ by $`\mathrm{\Gamma }_B\stackrel{~}{i}^{}\mathrm{\Gamma }_{[X]}C^1(pt,\stackrel{~}{C})=0`$ and $`\stackrel{~}{i}^{}\mathrm{\Gamma }_\xi ^{}C^2(pt,\stackrel{~}{C})=0`$. Then (1.5.2) holds also for $`i=0,4`$ using (1.3.3), because $`{}_{}{}^{t}\stackrel{~}{\mathrm{\Gamma }}=\stackrel{~}{\mathrm{\Gamma }}`$.
Finally we have
$$\stackrel{~}{\pi }_1^2=\stackrel{~}{\pi }_2.$$
$`(\mathrm{1.5.3})`$
Indeed, $`\stackrel{~}{\mathrm{\Gamma }}`$ is an idempotent, because $`\mathrm{\Gamma }_B\stackrel{~}{i}^{}\stackrel{~}{i}_{}\mathrm{\Gamma }_B\text{CH}^0(\stackrel{~}{C}\times \stackrel{~}{C})_{}`$ coincides with $`\mathrm{\Gamma }_B`$ (cohomologically). Then (1.5.3) follows from (1.5.1) and (1.3.1).
1.6. Proof of (0.1) $`(0.3)`$. Applying (0.1) to $`1(_{i2}\pi _i+\stackrel{~}{\pi }_2)`$, we get $`\stackrel{~}{\pi }_2=\pi _2`$ by (0.1). So it is enough to show $`(\stackrel{~}{\pi }_2)_{}\text{CH}^2(X)_{}=0`$ by \[7, 7.1\] (because the Albanese kernel is torsion free by ). See Remark after (1.3). By the definition of $`\stackrel{~}{\pi }_2`$, it is enough to show $`\stackrel{~}{i}^{}\text{CH}^2(X)_{}=0`$. But this is clear.
2. Cycle maps and correspondences
2.1. Let $`X`$ be a smooth proper complex variety with the structure morphism $`a_X:Xpt`$. Let $`(j)`$ denote the Tate Hodge structure of type $`(j,j)`$ (see ) which is naturally identified with a mixed Hodge Module on $`pt`$ (see \[10, (4.2.12)\]). Then we have a cycle map
$$cl:\text{CH}^p(X)_{}\text{Ext}^{2p}(a_X^{},a_X^{}(p))=\text{Ext}^{2p}(,(a_X)_{}a_X^{}(p)),$$
$`(\mathrm{2.1.1})`$
where $``$ means $`(0)`$, and Ext is taken in the derived category of mixed Hodge Modules on $`X`$ or $`pt`$. See (4.5.18) of loc. cit. (The last isomorphism of (2.1.1) follows from the adjoint relation.) The target is isomorphic to $``$-Deligne cohomology, and (2.1.1) is an isomorphism for $`p=1`$ as well-known. (This cycle map coincides with Deligne’s cycle map which uses local cohomology, and can also be obtained by using the theory of Bloch and Ogus as was done by Beilinson and Gillet. In particular, its restriction to homologically equivalent to zero cycles coincides with Griffiths’ Abel-Jacobi map.)
Let $`X,Y`$ be smooth proper complex varieties. Assume $`X`$ is purely $`n`$-dimensional. Then the cycle map induces
$`cl:C^i(X,Y)_{}`$ $`\text{Ext}^{2n+2i}(,(a_{X\times Y})_{}a_{X\times Y}^{}(n+i))`$ $`(\mathrm{2.1.2})`$
$`=\text{Ext}^{2i}((a_X)_{}a_X^{},(a_Y)_{}a_Y^{}\text{(i)) }.`$
See \[11, II\]. This is an isomorphism for $`n+i=1`$. By (3.3) of loc. cit, (2.1.2) is compatible with the composition of correspondences.
###### {\bf2.2.~Proposition}
Let $`X,S`$ be connected smooth proper complex varieties, and $`\mathrm{\Gamma }\text{CH}^p(S\times X)_{}`$. Assume $`\mathrm{\Gamma }`$ is homologically equivalent to zero, and the cycle map (2.1.1) for $`X`$ and $`p`$ is injective. Then $`\mathrm{\Gamma }=\mathrm{\Gamma }_1+\mathrm{\Gamma }_2`$ where $`\mathrm{\Gamma }_1`$ is supported on $`D\times X`$ for a divisor $`D`$ on $`S`$, and $`\mathrm{\Gamma }_2=[S]\times \xi `$ for $`\xi \text{CH}^p(X)_{}`$.
###### Demonstration
Proof.
Let $`H^{2p1}(X,)_S`$ denote a constant variation of Hodge structure on $`S`$ such that the fibers are $`H^{2p1}(X,)`$. This is identified with the direct image of $`_{S\times X}`$ by the first projection. Since the restriction of $`\mathrm{\Gamma }`$ to $`\{s\}\times X`$ is homologically equivalent to zero for any $`sS`$, $`\mathrm{\Gamma }`$ determines a normal function
$$e\text{Ext}^1(_S,H^{2p1}(X,)_S(p)),$$
where Ext is taken in the derived category of mixed Hodge Modules (using \[10, 3.27\]). Note that $`e`$ can be identified with a section of $`S\times J^p(X)S`$ (where $`J^p(X)`$ is Griffith’ intermediate Jacobian), if we replace rational coefficients with integral coefficients.
By the adjoint relation for $`a_S:Spt`$, we have a short exact sequence
$`0\text{Ext}^1(,H^{2p1}(X,)`$ $`(p))\text{Ext}^1(_S,H^{2p1}(X,)_S(p))`$
$`\text{Hom}(,H^1(S,)H^{2p1}(X,)(p))0,`$
where the first Ext and the last Hom are taken in the category of mixed Hodge structures. Since $`\mathrm{\Gamma }`$ is homologically equivalent to zero, the image of $`e`$ in the last term is zero, and hence $`e`$ comes from the first term, i.e. it is constant. So there exists $`\xi \text{CH}^p(X)_{}`$ such that $`\xi `$ is homologically equivalent to zero and $`e`$ is the image of $`[S]\times \xi `$. Then replacing $`\mathrm{\Gamma }`$ with $`\mathrm{\Gamma }[S]\times \xi `$, we may assume $`e=0`$.
Let $`k`$ be a finitely generated subfield of $``$ such that $`X,S`$ and $`\mathrm{\Gamma }`$ are defined over $`k`$, i.e. there exist smooth proper $`k`$-varieties $`X_k,S_k`$ and $`\mathrm{\Gamma }_k\text{CH}^p(S_k\times _kX_k)_{}`$ with isomorphisms $`X=X_k_k,`$ etc. Let $`K=k(S_k)`$ the function field of $`S_k`$, and $`X_K`$ the generic fiber of the first projection of $`S_k\times _kX_k`$. Let $`\mathrm{\Gamma }_K\text{CH}^p(X_K)_{}`$ denote the restriction of $`\mathrm{\Gamma }_k`$. Then it is enough to show $`\mathrm{\Gamma }_K=0`$.
We choose an embedding $`K`$ extending $`k`$. Since $`e=0`$, the image of $`\mathrm{\Gamma }_K_K\text{CH}^p(X)_{}`$ by the cycle map is zero, because $`\mathrm{\Gamma }_K_K`$ is identified with the restriction of $`\mathrm{\Gamma }`$ to $`\{s\}\times X`$ where $`sS`$ corresponds to the embedding $`K`$. So $`\mathrm{\Gamma }_K_K`$ is zero by hypothesis. Then the assertion follows from the injectivity of $`\text{CH}^p(X_K)_{}\text{CH}^p(X)_{}`$.
###### {\bf2.3.~Theorem}
Let $`X`$ be a connected smooth proper complex surface such that the Albanese map for $`X`$ is injective. Then the cube of the ideal in (0.2) is zero.
###### Demonstration
Proof.
Let $`\mathrm{\Gamma }\text{CH}^2(X\times X)_{}`$ that is homologically equivalent to zero. Then $`\mathrm{\Gamma }=\mathrm{\Gamma }_1+\mathrm{\Gamma }_2`$ such that $`\mathrm{\Gamma }_1`$ is supported on $`D\times X`$ and $`\mathrm{\Gamma }_2=[X]\times \xi `$ as in (2.2). We have to show
$$\mathrm{\Gamma }^{\prime \prime }\mathrm{\Gamma }^{}\mathrm{\Gamma }_i=0(i=1,2)$$
$`(\mathrm{2.3.1})`$
for any $`\mathrm{\Gamma }^{},\mathrm{\Gamma }^{\prime \prime }C^0(X,X)`$ which are homologically equivalent to zero.
For $`i=1`$, let $`Y`$ denote the normalization of $`D`$ with $`f:YX`$ the canonical morphism. Then there exists $`\mathrm{\Gamma }_1^{}C^0(Y,X)_{}`$ such that $`\mathrm{\Gamma }_1=\mathrm{\Gamma }_1^{}f^{}`$. By the injectivity of the cycle map in the divisor case, it is enough to show that the image of $`\mathrm{\Gamma }^{\prime \prime }\mathrm{\Gamma }^{}\mathrm{\Gamma }_1^{}`$ by (2.1.2) is zero. Since $`\mathrm{\Gamma }^{},\mathrm{\Gamma }^{\prime \prime }`$ are homologically equivalent to zero, and $`\text{Ext}^2`$ vanishes, the assertion follows by using for example a (noncanonical) decomposition
$$(a_X)_{}a_X^{}_iH^i(X,)[i]$$
in the derived category of mixed Hodge Modules on $`pt`$ (or equivalently, of graded-polarizable mixed Hodge structures). See \[10, (4.5.4)\].
The argument is similar for $`i=2`$ by using the injectivity of the Albanese map, because it is enough to show that the image of $`\mathrm{\Gamma }^{\prime \prime }\mathrm{\Gamma }^{}\mathrm{\Gamma }_\xi C^2(pt,X)`$ by the cycle map (2.1.2) is zero. This completes the proof of (2.3).
2.4. Remark. The square of the ideal in (0.2) is nonzero if $`H^1(X,)0`$. Indeed, let $`C`$ be a hyperplane section of $`X`$ with the inclusion morphism $`i:CX`$. By Hodge theory we have a divisor $`D`$ on $`C\times X`$ such that $`D_{}:H^i(C,)H^i(X,)`$ is zero for $`i1`$ and $`D_{}i^{}:H^1(X,)H^1(X,)`$ is an isomorphism. Let $`k`$ be a finitely generated subfield of $``$ such that $`X`$, $`C`$, $`D`$ are defined over $`k`$, i.e. there exist $`X_k`$, $`C_k`$, $`D_k`$ with isomorphisms $`X_k_k=X,`$ etc. Let $`D_i=(pr_i\times id)^{}D_k`$ where $`pr_i:C_k\times _kC_kC_k`$ are the natural projections for $`i=1,2`$. Let $`K=k(C_k\times _kC_k)`$, and $`\mathrm{\Gamma }_{K,i}`$ be the restriction of $`D_i`$ to the generic fiber of $`C_k\times _kC_k\times _kX_kC_k\times _kC_k`$. Then $`\mathrm{\Gamma }_{K,1}\times _K\mathrm{\Gamma }_{K,2}`$ is the restriction of $`D_k\times _kD_k`$. Let $`\mathrm{\Gamma }_i=\mathrm{\Gamma }_{K,i}_KC^1(pt,X)_{}`$. Then the composition of $`\mathrm{\Gamma }_{[C]}{}_{}{}^{t}\mathrm{\Gamma }_{1}^{}`$ and $`\mathrm{\Gamma }_1{}_{}{}^{t}\mathrm{\Gamma }_{[C]}^{}`$ is equal to a nonzero multiple of $`\mathrm{\Gamma }_2{}_{}{}^{t}\mathrm{\Gamma }_{1}^{}`$, and this is nonzero.
References
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A. Beilinson, Height pairing between algebraic cycles, Lect. Notes in Math., vol. 1289, Springer, Berlin, 1987, pp. 1–26.
S. Bloch, Lectures on algebraic cycles, Duke University Mathematical series 4, Durham, 1980.
S. Bloch, A. Kas and D. Lieberman, Zero cycles on surfaces with $`p_g=0,`$ Compos. Math. 33 (1976), 135–145.
P. Deligne, Théorie de Hodge I, Actes Congrès Intern. Math., 1970, vol. 1, 425-430 : II, Publ. Math. IHES, 40 (1971), 5–57; III ibid., 44 (1974), 5–77.
D. Mumford, Rational equivalence of $`0`$-cycles on surfaces, J. Math. Kyoto Univ. 9 (1969), 195–204.
J.P. Murre, On the motive of an algebraic surface, J. Reine Angew. Math. 409 (1990), 190–204.
, On a conjectural filtration on Chow groups of an algebraic variety, Indag. Math. 4 (1993), 177–201.
A. Roitman, The torsion in the group of zero cycles modulo rational equivalence, Ann. Math. 111 (1980), 553–569.
M. Saito, Mixed Hodge Modules, Publ. RIMS, Kyoto Univ., 26 (1990), 221–333.
, Hodge conjecture and mixed motives, I, Proc. Symp. Pure Math. 53 (1991), 283–303; II, in Lect. Notes in Math., vol. 1479, Springer, Berlin, 1991, pp. 196–215.
A. J. Scholl, Classical Motives, Proc. Symp. Pure Math. 55 (1994), Part 1, 163–187.
C. Voisin, Sur les zéro cycles de certaines hypersurfaces munies d’un automorphisme, Ann. Sci. Norm. Sup. Pisa 19 (4) (1992), 473–492.
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# The Globular Cluster Systems in the Coma Ellipticals. III: The Unique Case of IC 4051
## 1 INTRODUCTION
The Coma Cluster ($`d100`$ Mpc), as a rich Abell cluster, is the host environment for a huge range of E/S0 galaxies. Its two central supergiants, NGC 4874 and 4889, are among the very most luminous galaxies known, and the cluster has many other large ellipticals scattered throughout its $`1`$Mpc core region and well beyond. The globular cluster systems (GCSs) around these galaxies are well within reach of the HST cameras, and thus give us an extraordinarily rich range of target galaxies for comparative GCS studies. Coma also has important implications for the determination of the Hubble constant $`H_0`$, since its large recession velocity ($`7100`$ km s<sup>-1</sup>) greatly exceeds any anticipated local peculiar motions, and it is situated at high galactic latitude ($`b=+87\stackrel{}{\mathrm{.}}7`$) nearly unaffected by foreground absorption,
IC 4051 is a giant E2 galaxy located $`14^{}`$ east in projection from the center of Coma and has no luminous neighbors. With an integrated magnitude $`V_T^0=13.20`$, it is the fifth brightest elliptical galaxy within the $`1\mathrm{°}`$ central region of Coma. Thus, IC 4051 presents a good opportunity to study the GCS of a more-or-less “normal” large elliptical. Baum et al. (1997) obtained the first deep WFPC2 exposures of IC 4051 with exactly this purpose in mind, but they limited their published analysis to only the photometry from the PC1 frame on which the galaxy was centered. However, considerably more information on the radial structure of the GCS, its specific frequency, and radial metallicity gradient are potentially available from the outer sections of the WFPC2 field. In this paper, we re-analyze the archival HST/WFPC2 frames from the Baum et al. (1997) program and discuss the global properties of the IC 4051 GCS using the entire body of data available in the raw exposures.
## 2 DATA ANALYSIS
The raw database comprises HST/WFPC2 observations taken on 1995 July 20/21 (Baum et al. 1997). Eight images totalling 20500 s (essentially, a sequence of full-orbit exposures) were taken through the $`F606W`$ (wide $`V`$) filter, and two images totalling 5200 s were taken through the $`F814W`$ ($`I`$) filter. The techniques in our data analysis are similar to those in Papers I and II of this series (Kavelaars et al. (2000); Harris et al. (2000)), in which we analyzed the GCS around NGC 4874, the central cD giant. The first step was to register and combine the images to produce a master $`V`$ and master $`I`$ frame free of cosmic rays. An excellent color composite image constructed from the combined frames is published by Baum et al. For the PC1 frame, we produced a “flattened” image in which the overall elliptical contours of the central galaxy were modelled and then subtracted. In the WF2,3,4 fields, we generated an empirical model of the galaxy light by median filtering each frame and subtracting the smoothed image from the original picture (after a preliminary star-finding and removal).
Our photometry was then carried out on the flattened master frames with the DAOPHOT II and ALLSTAR codes (Stetson (1992)). The instrumental magnitudes returned by ALLSTAR were transformed to the standard $`V`$,$`I`$ system via the equations derived by Holtzman et al. (1995),
$$V=m(F606W)\overline{\mathrm{\Delta }m}+22.093+0.254(VI)+0.012(VI)^2+2.5\mathrm{log}(GR_i),$$
(1)
$$I=m(F814W)\overline{\mathrm{\Delta }m}+20.8390.062(VI)+0.025(VI)^2+2.5\mathrm{log}(GR_i),$$
(2)
where $`m(F606W,F814W)`$ are the instrumental magnitudes returned by ALLSTAR, and the $`\overline{\mathrm{\Delta }m}`$ are constants which shifts these ALLSTAR psf magnitudes to the equivalent magnitude of the light within a $`0\stackrel{}{\mathrm{.}}5`$ aperture. The mean value of $`\mathrm{\Delta }m`$ was determined empirically for each of the four WFPC2 fields from roughly ten moderately bright, isolated stars in each CCD, and applied to all our detected starlike images. The typical internal uncertainties in the mean $`\mathrm{\Delta }m`$ values were $`\pm 0.02`$ mag.
It should be noted that at the distance of Coma, all globular clusters are easily starlike in appearance and thus these transformation equations normalized to the $`0\stackrel{}{\mathrm{.}}5`$ aperture can strictly be applied to them. The calibration equation for $`V`$ quoted above is the synthetic model transformation from Table 10 of Holtzman et al. (1995), while the equation for $`I`$ is the observationally based one from their Table 7. If, instead, we had chosen their Table 10 model transformation for $`I`$, our resulting $`(VI)`$ color indices would have ended up bluer by roughly 0.02 mag at the average color $`(VI)1.2`$ of the IC 4051 globular clusters (see below; this shift would have the effect of reducing the deduced cluster metallicities by $`\mathrm{\Delta }`$\[Fe/H\] = 0.1).<sup>1</sup><sup>1</sup>1Baum et al. (1997) adopted the synthetic transformation for $`I`$ from Table 10 of Holtzman et al., but used the specific curve for the color range $`(VI)<1.0`$. However, the majority of the globular clusters are somewhat redder than this, so that the transformation for $`(VI)>1.0`$ would have been preferable. The net result is to increase the small $`(VI)`$ difference between their scale and ours by a further $`0.020.03`$ mag, although our instrumental color scale $`(F606WF814W)`$ agrees quite closely with theirs. When we add the internal uncertainties in the aperture corrections $`\mathrm{\Delta }m`$, we estimate that the zeropoints of either the magnitude or color scales are uncertain to at least $`\pm 0.03`$ mag.
Finally, as an internal check on the relative zeropoints of the $`V`$ and $`I`$ scales between the PC1 chip and the three WF chips, we inspected the mean $`(VI)`$ color indices of the measured globular clusters in the annulus around the center of IC 4051 ($`R10^{\prime \prime }20^{\prime \prime }`$) that overlapped the PC/WF boundary. We enforced the WF2,3,4 data in this annulus to have the same mean color as those in the PC1 region by adjusting the $`I`$ magnitudes (which are from much shorter exposures than $`V`$, thus internally more uncertain). The final result places our $`(VI)`$ scale in agreement with the Baum et al. (1997) scale to within 0.03 mag, a level entirely consistent with the combined photometric uncertainties.
As was already evident from the Baum et al. (1997) study, IC 4051 has a populous GCS, which appears as a very obvious swarm of faint objects across the whole WFPC2 field. The main source of sample contamination is from very faint, compact background galaxies, with a (nearly negligible) contribution from Galactic foreground stars. A high proportion of the background galaxies can be eliminated through conventional radial-moment image analysis. For this purpose we used the $`r_1`$ radial moment as implemented in Harris et al. (1991),
$$r_1=\left(\frac{rI}{I}\right),$$
(3)
which is an intensity-weighted mean radius for the object calculated over all pixels brighter than the detection threshold (see Harris et al. 1991). A straightforward plot of $`r_1`$ against magnitude then shows a well defined stellar sequence, with nonstellar objects scattering to larger $`r_1`$. These classification graphs for the four CCDs in $`V`$ are shown in Figure 1. The dashed lines indicate the adopted cutoffs applied to the measurements. Similar object classifications were applied to the $`I`$ data (which, however, have shallower limits), with the final culled data lists containing 4058 objects in $`V`$ and 1672 objects in $`(VI)`$.
The faint-end completeness of our photometry was investigated through an extensive series of artificial-star tests on the master images. The procedure performed here was to add 500 artificial stars to each frame over a range of input magnitudes, measure these frames through the normal DAOPHOT sequence, and find out how many were recovered. Fifteen of these trials were carried out, with average resulting completeness fractions as displayed in Figure 2. Convenient fits to the raw points are provided by the Pritchet interpolation function,
$$f=\frac{1}{2}\left[1\frac{\alpha (RR_{lim})}{\sqrt{1+\alpha ^2(RR_{lim})^2}}\right].$$
(4)
Table 1 summarizes the best-fit parameters to Equation 4 for each CCD and bandpass. In the Table, ($`V_{lim},I_{lim}`$) are the magnitudes at which $`f`$ drops to 0.5, and the parameter $`\alpha `$ controls the steepness of the falloff. The curves for all three of the outer chips (WF2,3,4) are nearly identical; for the inner PC1 chip, the limiting magnitudes are brighter, driven by the spread of the PSF over many more pixels and (for $`R{}_{}{}^{<}\mathrm{\hspace{0.17em}10}_{}^{\prime \prime }`$) the brighter background light.
Lastly, Figure 3 shows how the photometric measurement uncertainties (also derived from the ADDSTAR completeness tests) increase with magnitude. At the formal limiting ($`f=0.5`$) magnitude, the rms uncertainty in the photometry reaches 0.15 mag. Wherever possible, we avoid dealing with any features of the data below that limit.
## 3 COLOR AND METALLICITY DISTRIBUTIONS
The distribution in colour of the globular clusters can be used to gain insight into the existence of multiple sub-populations in the GCS. Bimodal color distributions are found about half the time in gE galaxies (e.g., Kundu & Whitmore (1999); Neilsen & Tsvetanov (1999)) and are often interpreted as relics of at least two major phases of star formation in the early history of the galaxy, whether by merger, accretion, or in situ processes. With the conventional “null hypothesis” for giant ellipticals that the clusters are all old ($`{}_{}{}^{>}\mathrm{\hspace{0.17em}10}`$ Gy), the color index is primarily a tracer of cluster metallicity. The colour-magnitude distribution for the 1672 objects measured in both $`V`$ and $`I`$ is shown in Figure 4. At projected galactocentric radii larger than about $`80\mathrm{}`$, we found (see below) that the residual numbers of clusters dropped nearly to zero, so we adopt this outer region as defining a suitable “background” population.
To eliminate a few more contaminating objects, we reject objects bluer than $`(VI)=0.74`$ or redder than $`(VI)=1.46`$ (vertical dashed lines in Fig. 4). These colour limits generously include the range in colours of the known globular clusters in large galaxies (e.g., Harris (1996); Whitmore et al. (1995); Neilsen & Tsvetanov (1999)). We also further limited the color sample to objects brighter than $`V=26.0`$ to ensure high completeness at all colors.
A “clean” color distribution for the GCS was then obtained by subtraction of the background ($`R>80^{\prime \prime }`$) color distribution, normalized to the same total area as the inner population. This procedure left a final total of 479 objects within the magnitude and color limits given above, with a net distribution over $`(VI)`$ as shown in the histogram of Figure 5. The mean color of the sample is $`VI=1.12\pm 0.01`$ (internal uncertainty of the mean), with a dispersion of $`\sigma _{VI}=0.13`$. Subtracting an adopted foreground reddening $`E(VI)=0.014`$ and using the calibration of $`(VI)_0`$ in terms of metallicity given in Paper II,
$$(VI)_0=0.17[\mathrm{Fe}/\mathrm{H}]+\mathrm{\hspace{0.17em}1.15},$$
we then estimate that the IC 4051 GCS as a whole has $``$Fe/H$`0.3`$. The peak position of this color distribution is quite similar to the metal-rich components in other giant ellipticals such as NGC 4472 (Geisler et al. (1996)), M87 (Whitmore et al. (1995); Kundu et al. (1999)), and other Virgo members (Neilsen & Tsvetanov (1999)). However, the metal-poor component which is usually found in these same galaxies at a mean color $`(VI)0.95`$ or \[Fe/H\] $`1.5`$ (and which we found in the Coma cD NGC 4874; see Paper II) is entirely missing in IC 4051, or at very most is a fringe component buried in the wings of the main distribution.
We cannot place firm limits on the intrinsic dispersion $`\sigma `$\[Fe/H\], since the mean observational measurement scatter over the sample is $`\sigma _{VI}0.11`$, comparable with the observed sample dispersion of $`\pm 0.13`$. Nevertheless, subtracting off the observational scatter in quadrature, we estimate roughly $`\sigma _0`$\[Fe/H\] $`0.4`$, which is in close agreement (for example) with the value $`\sigma `$\[Fe/H\] = 0.38 found by Geisler et al. (1996) to fit each of the metal-rich and metal-poor components in NGC 4472. In the Milky Way, the well known bimodal MDF has been found to be fit by Gaussian functions with dispersions near $`0.3`$ dex (Zinn (1985); Armandroff & Zinn (1988); Harris (2000)). For IC 4051, a single Gaussian with the same mean and standard deviation as the sample (Fig. 5) matches the MDF with a $`\chi ^214.6`$ over 14 degrees of freedom, which provides no strong evidence for bimodality (but see below).
Trends of mean color with either galactocentric distance or magnitude were also searched for. Table 2 shows the mean color and dispersion in 0.5-magnitude bins from $`V=22.5`$ to 26.0. These binned means reveal no significant change in color with luminosity. However, slightly more interesting features emerge in the graph of color versus radius (Figure 6). Binned mean colors, listed in Table 3, indicate no systematic change in color for $`R{}_{}{}^{>}\mathrm{\hspace{0.17em}10}_{}^{\prime \prime }`$, but within $`10^{\prime \prime }`$ the clusters are indeed slightly redder than the overall mean. The distribution in its entirety is barely suggestive of two sub-populations: one centered on $`(VI)1.2`$ which is found at all radii; and a second, slightly bluer one centered near $`(VI)1.0`$. The lack of bluer clusters within $`R{}_{}{}^{<}\mathrm{\hspace{0.17em}10}_{}^{\prime \prime }`$ is then largely responsible for the inner color gradient of the whole sample mentioned above. Much stronger versions of this same effect have shown up in some other giant E or cD galaxies with far more obvious bimodal MDFs (e.g., Secker et al. (1995); Geisler et al. (1996); Lee et al. (1998); Ostrov & Forte (1998)). In these, the different central concentrations of the metal-rich and metal-poor subsystems produce a steady outward change in the relative proportions of blue-to-red clusters with radius and thus a mean metallicity gradient.
Using Fig. 6 as a guide, we divided the sample of objects at $`(VI)=1.07`$ and tested the radial distributions of the bluer and redder halves. A standard Kolmogorov-Smirnov two-sample test indicated that their spatial distributions are significantly different (the redder half is more centrally concentrated) at the 99% level, suggesting to us that the inner gradient is indeed a real effect.
We therefore very tentatively suggest that the IC 4051 system may contain a bimodal MDF in which the two modes are rather closely spaced in mean metallicity, thus heavily blurred out by the raw photometric measurement uncertainty. Numerical experiments with various two-component fittings of the entire MDF lead to models of the form shown in Figure 7. Here, a sample twin-Gaussian fit is shown in which the bluer (metal-poor) component is centered at $`(VI)=1.00`$ or \[Fe/H\] $`0.96`$, the redder (metal-rich) one at $`(VI)=1.17`$ or \[Fe/H\] $`+0.04`$, both have dispersions $`\sigma (VI)=0.10`$, and the redder one contains about 55% of the total sample. The combined components now represent the total shape of the MDF better, with its modest skewness toward the red side (the total $`\chi ^2`$ is 12.6). The relative proportions of blue and red components, however, are quite uncertain (the formal uncertainties are $`\pm 0.1`$, but variations of factors of two in the proportions give scarcely different overall fits).
Clearly, this particular two-component model is only illustrative of the range of possibilities: the moderately small difference in color between the two components, and the very significant broadening of the MDF by photometric measurement uncertainty, do not justify more extensive analysis. However, it would clearly be of value to measure the MDF of this populous globular cluster system with a photometric index much more sensitive to metallicity than $`(VI)`$, in which the subpopulations would be far more clearly revealed. A more sensitive color index would also permit establishment of the true mean \[Fe/H\] with much less zeropoint uncertainty.
Lastly, it is worth comparing the mean colors of the GCS components to that of the halo light of the central galaxy. Mehlert et al. (1998) find $`(VI)1.30`$ at a projected radius $`R=10^{\prime \prime }`$, increasing inward to $`(VI)=1.35`$ at the very center. This color range is distinctly redder than the typical levels $`(VI)1.20\pm 0.03`$ for giant E galaxies (Buta & Williams (1995); Prugniel & Héraudeau (1998)). The measured absorption line indices (Mg, Fe, H$`\beta `$) lead Mehlert et al. (1998) to conclude, in line with the integrated color, that the core of IC 4051 is extremely old and very metal-rich, perhaps as high as \[Fe/H\] $`=+0.25`$. However, the deduced metallicity from the line indices becomes lower at larger radii, dropping to an equivalent \[Fe/H\] $`0.5`$ for $`R{}_{}{}^{>}\mathrm{\hspace{0.17em}20}_{}^{\prime \prime }`$ (the effective radius $`r_e`$ of the light profile), similar to the inner GCS. Mehlert et al. find that IC 4051 harbors an old, co-rotating core with an unusually large “break radius” (it is detectable out to $`5^{\prime \prime }`$ or 3.4 kpc) but which contributes $`{}_{}{}^{}{}_{}{}^{<}`$ 1% of the total light of the galaxy. If this inner stellar disk is a signature of a dissipational merging event, it is likely to have occurred at early times.
## 4 THE LUMINOSITY DISTRIBUTION
As Baum et al. (1997) showed, the $`V`$ photometry reaches faint enough to reveal the “turnover point” (peak frequency) in the globular cluster luminosity function (GCLF). By adding in the photometry from the WF chips, we have been able to double the total sample of clusters and thus improve the definition of the GCLF.
The distribution of all the detected objects classified as “starlike” and used to define the GCLF is shown in Figure 8. These are, quite evidently, strongly concentrated to the center of IC 4051 (much more so than in the GCS of the Coma supergiant NGC 4874; see Papers I and II). More or less arbitrarily, we take the region $`R>80^{\prime \prime }`$ marked by the outer dashed line in Fig. 8, as defining the luminosity function of the background population, to be subtracted statistically from the inner ($`10^{\prime \prime }<R<80^{\prime \prime }`$) zone after correction for photometric incompleteness.
The results of this exercise for each of the four CCD chips separately are shown in Figure 9. Aside from the noticeably brighter completeness limit for the PC1 zone, no significant differences in the GCLF shape or turnover from place to place are evident. (The GCLF peak for the PC1 region shows an apparent peak fainter than $`V28`$, but this is fainter than the 50% completeness limit and so cannot be given much weight.) We therefore add all four sectors to form the composite GCLF shown in Figure 10. The numerical results in 0.3-mag bins are listed in Table 4: here, successive columns give (1) the $`V`$ magnitude range of the bin (2) the number of detected starlike objects in the inner ($`10^{\prime \prime }80^{\prime \prime }`$) zone (3) the number in the outer ($`>80^{\prime \prime }`$) background zone (4) the number in the inner zone corrected for completeness, and (5) the net GCLF, after subtraction of the area-normalized background counts.
To estimate the turnover level and shape of the GCLF, we fit a standard Gaussian interpolation function (Harris (1991); Jacoby et al. (1992)) to the data shown in Fig. 10, setting the standard deviation $`\sigma _V`$ of the curve and then solving for the best-fit turnover level $`V^0`$. Trials with different adopted $`\sigma _V`$’s gave the results summarized in Table 5. The reduced $`\stackrel{~}{\chi }^2`$ values favor a solution in the broad range $`\sigma _V1.41.8`$, with little to choose among values in this range in a formal sense. However, it is well known that both $`\sigma _V`$ and $`V^0`$ tend to be overestimated in situations like these where the magnitude limit of the data reaches barely past the actual turnover (e.g., Hanes & Whittaker (1987); Paper I) since the solutions for the two parameters are correlated. For this reason, we favor a choice in the narrower range $`\sigma _V1.41.6`$ and $`V^027.628.0`$. Sample Gaussian curves for the extremes of this range are shown in Fig. 10.
Our final adopted pair of parameters is $`V^0=27.8\pm 0.2`$, $`\sigma _V=1.5\pm 0.1`$. For comparison, Baum et al. (1997) found $`V^0=27.72`$ employing a different and more complex fitting function. As is discussed more extensively in Paper I, this turnover level is also similar to what we found in the central cD NGC 4874. Using both of them combined, along with a calibration of the absolute magnitude of the turnover point based on the Virgo ellipticals, we find $`d100`$ Mpc for Coma along with a Hubble constant $`H_070`$ (see Paper I).
A second and more physically oriented way to display the same material is as the luminosity distribution function (LDF), or number of clusters per unit (linear) luminosity. (The relation between the GCLF and LDF forms is exhaustively discussed by McLaughlin (1994).) The LDF is shown in Figure 11. At levels brighter than the GCLF “turnover” (which in turn is only slightly brighter than the photometric completeness limit), the LDF clearly approximates a power-law falloff toward higher luminosity, $`N(L)dLL^\alpha `$. To second order, however, the slope $`\alpha =d`$log($`N`$)/$`d`$log($`L`$) appears to steepen slightly at the upper end: an unweighted least-squares fit to all bins brighter than the turnover yields $`\alpha =2.05`$, while exclusion of the half-dozen very brightest bins yields $`\alpha =1.75`$.
These power-law forms – as well as logarithmic slope values $`\alpha 2`$ – are entirely similar to what has been found in a wide range of other galaxies from dwarf ellipticals to spirals (Harris & Pudritz (1994); Durrell et al. (1996)). However, in most giant ellipticals studied to date, the slopes tend to be somewhat flatter at $`\alpha 1.5\pm 0.3`$ (Harris & Pudritz (1994)). The total shape for log $`(L/L_{}){}_{}{}^{>}\mathrm{\hspace{0.17em}5}`$, complete with its progressive steepening toward higher luminosity, can be well matched by a formation model in which protocluster clouds build up by collisional agglomeration and in which the more massive clouds have shorter lifetimes against star formation (McLaughlin & Pudritz (1996); Harris (2000)). Our data for IC 4051 add further to the general body of material which indicates a remarkable place-to-place similarity in the luminosities of old globular clusters, and thus a quasi-universal formation process.
## 5 RADIAL DISTRIBUTION AND SPECIFIC FREQUENCY
Because the GCS around IC 4051 is quite centrally concentrated (see Fig. 8), we can use the complete WFPC2 data to define the spatial distribution outward nearly to its limits. The radial profile of the raw counts for all starlike objects brighter than $`V=27.0`$, for which the data are highly complete nearly in to the central core of the galaxy, is shown in Figure 12. The inner core ($`R{}_{}{}^{<}\mathrm{\hspace{0.17em}5}_{}^{\prime \prime }`$), in which the projected density of clusters is nearly flat, continues outward to a steep power-law falloff which covers most of our survey area. Finally, for $`R>80^{\prime \prime }`$, the number density $`\sigma `$ begins to level off towards its eventual far-field background level; more or less arbitrarily, we set this background at $`\sigma _b=(0.02\pm 0.01)`$ arcsec<sup>-2</sup> as representing nearly the average of the outermost two points. (As will be seen below, small differences in the adopted $`\sigma _b`$ will not have major effects on any of our subsequent conclusions.)
The complete profile data broken into circular annuli are listed in Table 6, giving the number of objects in each bin, the surface area of the annulus, and the projected density $`\sigma `$. The residual number density of clusters, $`\sigma _{cl}=\sigma \sigma _b`$, is plotted in Figure 13. Simple King (1966) models can be fitted to the $`\sigma _{cl}`$ data points to give rough estimates of the GCS core radius and central concentration: performing a weighted fit in the manner described in Paper II, and ignoring the very uncertain outermost three points, we find a formal best-fit core radius $`R_c=10\stackrel{}{\mathrm{.}}25`$ (equivalent to 5.1 kpc at the adopted Coma distance) as well as a concentration index $`c=1.45`$ for a dimensionless central potential $`W_0=6.26`$. The core radius is four times smaller in IC 4051 than the $`22`$ kpc value we found in the much more extended NGC 4874 (Paper II).
A second comparison can be made with the halo light of the galaxy. It is conventionally found in giant ellipticals that the GCS is a more spatially extended system as a whole than the halo (Harris (1991), 1999). IC 4051 is no exception, despite its overall compact structure. In Fig. 13, we show the wide-field surface intensity profiles in $`\mu _R`$ measured by Strom & Strom (1978) and Jorgensen et al. (1992). Although the Strom data are photographically measured, their profile agrees tolerably with the more recent CCD measurements of Jorgensen et al. (1992) over their region of overlap.
The bulk of the GCS profile is more extended than the halo light, except possibly for the outer ($`R{}_{}{}^{>}\mathrm{\hspace{0.17em}30}_{}^{\prime \prime }`$) regions where their slopes are more nearly similar. For $`R{}_{}{}^{>}\mathrm{\hspace{0.17em}20}_{}^{\prime \prime }`$, the GCS profile behaves as $`\sigma _{cl}R^2`$, although at the largest radii little weight can be placed on the very uncertain outermost half-dozen points. There is a strong hint from the halo light profile that the galaxy may be truncated past $`R60^{\prime \prime }`$ (about 30 kpc), though here again the profile is very sensitive to slight differences in the adopted background level, so not much meaning can be ascribed to the slope differences between the GCS and the halo there. For giant E galaxies in general, a rough mean relation between galaxy luminosity $`M_V^T`$ and the radial falloff outside the central core is (Kaisler et al. (1996)) $`d`$log$`\sigma `$/$`d`$log$`R`$ $`0.29M_V^T8.00`$. For IC 4051, this relation would predict $`\sigma _{cl}R^{1.65}`$, somewhat flatter than the observed $`R^2`$ trend.
Calculating the total GCS population and specific frequency is now a straightforward matter. From Table 6, we multiply $`\sigma _{cl}(R)`$ by the area of each annulus, then sum the annuli to get the total cluster population out to the limits of our survey. We find $`N=(1845\pm 165)`$ for $`V27.0`$ and $`R{}_{}{}^{<}\mathrm{\hspace{0.17em}130}_{}^{\prime \prime }`$. If the true GCLF turnover magnitude is at $`V^0=27.8\pm 0.2`$ (see above), then we must multiply this raw total by $`(3.35\pm 0.52)`$ to estimate the total cluster population over all magnitudes, giving $`N_{cl}=6180\pm 1100`$. The integrated luminosity of the galaxy is $`V^T=13.20`$ (RC3 catalog value), corresponding to $`M_V^T=21.9`$ for our adopted Coma distance. Thus, the specific frequency is
$$S_N=N_{cl}10^{0.4(M_V^T+15)}=10.8\pm 1.9.$$
In strict terms this is a lower limit to the true global $`S_N`$, since we have not accounted for any cluster population outside the $`120^{\prime \prime }`$ radial limit of our WFPC2 field. However, given that the halo is clearly declining quite steeply in this region (Fig. 13), any such population correction is likely to be small. A generous but reasonable upper limit estimate to the total population can be made if we assume that the GCS profile continues as $`\sigma _{cl}R^2`$ outward to the nominal tidal radius at $`R_t230^{\prime \prime }`$. This assumption gives an additional $`380\pm 300`$ clusters brighter than $`V=27`$, which then translates to $`S_N=12.6\pm 2.6`$.
Placing more weight on the lower limit – which reflects the steep falloff of the system near the radial limit of our data – we adopt a final estimate
$$S_N(\mathrm{final})=11\pm 2.$$
Remarkably, this GCS population ratio is several times higher than the $`S_N{}_{}{}^{<}\mathrm{\hspace{0.17em}2}`$ value found in NGC 4881 (Baum et al. (1995)), a galaxy which is quite comparable with IC 4051 in luminosity, structure, and location on the outskirts of the Coma core. This high $`S_N`$, in fact, places IC 4051 in the range which is conventionally reserved for the central-giant cD galaxies like M87 and many other BCGs (Harris et al. (1998); Blakeslee (1997), 1999). It is, perhaps, particularly noteworthy that IC 4051 has a specific frequency three times higher than the central cD in its own host galaxy cluster, NGC 4874 (see Paper II). No other instance of such a large contrast between a low$`S_N`$ central cD and a higher$`S_N`$ outlying elliptical is known. IC 4051 provides striking evidence that a central location in a rich cluster environment is not required to form a high population of globular clusters.
## 6 DISCUSSION
A brief summary of our findings for IC 4051 is that its GCS is (a) almost entirely metal-rich, albeit possibly with two narrowly separated subcomponents; (b) relatively compact in radial structure; and (c) a “high specific frequency” system despite that fact that its host galaxy is not a central giant elliptical nor one with a cD-type envelope.
Just as in Paper II for NGC 4874, we now attempt to use the integrated characteristics to reconstruct a partial history of the system. Formation scenarios for giant ellipticals tend to fall into three basic camps: (a) “in situ” formation, whereby the galaxy condenses by dissipative collapse of gas clouds in its immediate vicinity, in one or more major bursts; (b) later mergers of pre-existing disk-type galaxies with both gas and stars; or (c) successive mergers or accretions of smaller gas-poor satellites. Various combinations of these extremes are, of course, possible, and even likely.
For IC 4051, the lack of low-metallicity clusters already places fairly strong constraints on the range of possible formation events. For example, the mechanism investigated by Côté et al. (1998) – in which an original metal-rich “seed” gE accretes dozens or hundreds of smaller satellites – is unlikely, since these dwarf satellites would have brought in a population of hundreds or even thousands of low-metallicity clusters, which we do not see.
Similarly, merger-formation models in which gas-rich disk galaxies combine to build a composite elliptical (Ashman & Zepf (1992)) would predict a strong component of metal-poor clusters in the resulting MDF from the globular clusters that were present in the pre-merger galaxies. These merger models also have severe difficulty in generating high specific frequency products, since increasing the cluster population relative to the field stars by a large enough amount to produce high $`S_N`$ requires very large ($`>10^{10}M_{}`$) input gas masses, more than is routinely available in disk galaxies today. The normal merger route does appear to be quite effective as a logical source for low$`S_N`$ field ellipticals (see Harris (2000) or Whitmore & Schweizer (1995) for much more extensive discussion).
However, if either the merger or accretion processes are taken to an extreme form in which the merging objects are almost completely gaseous, then they become closely similar to the in situ route, and the conundrum of the missing low-metallicity clusters can be more easily circumvented. If the gas supply – however it was assembled – underwent most or all of its star formation in the high-pressure, high-density environment of the protoelliptical, then the conversion of gas to stars would have run much further to completion and built up the metallicity to the high levels that we now observe. Later gaseous mergers are, of course, not ruled out: the central corotating disk in the core of IC 4051 (Mehlert et al. 1998) with its very high metallicity is a likely signature of such an event, though at its $`{}_{}{}^{<}\mathrm{\hspace{0.17em}1}`$% contribution to the present-day luminosity, it probably did not form more than a few dozen globular clusters along with it, and even these would have mostly disrupted by now if they resided in the central few kpc of the core.
The relatively compact structure of the galaxy may be the result of tidal trimming (“harrassment”) from the Coma potential well (e.g., Moore et al. (1996)). The radial velocity of IC 4051 (4940 km s<sup>-1</sup>) is almost two standard deviations away from the Coma centroid (6850 km s<sup>-1</sup>; see Colless & Dunn (1996)), indicating that this galaxy oscillates back and forth through the cluster and is now passing through the dense Coma core at high speed.
These elements of an evolutionary scenario for IC 4051 are in strong contrast to NGC 4874, for which we argued (Paper II) that a large fraction of its clusters (which are almost entirely low-metallicity) could have been acquired by accretions of smaller satellites. In IC 4051, we are forced to argue that the bulk of its clusters formed in situ. The globular clusters in these two galaxies provide unique evidence for the view that large E galaxies can form by radically different evolutionary routes.
One of the most challenging elements of IC 4051 to interpret is certainly the high specific frequency of its GCS. In the previous literature (Harris (1991), 2000; Blakeslee (1997), 1999; Harris et al. (1998); McLaughlin (1999)) it has become conventional to associate high $`S_N`$ with giant galaxies at the centers of rich clusters. These central BCG’s or cD’s can have had histories of star and cluster formation through inflowing gas clouds and filaments, mergers, and accretions (e.g., Dubinski (1998)) that were much more extended than for normal outlying ellipticals. Recently, the view has been developed that such high$`S_N`$ galaxies should be regarded not as “cluster-rich” but rather as “star-poor” (Blakeslee (1997), 1999; Harris et al. (1998); McLaughlin (1999)). In this scheme, we postulate that the protogalactic gas started forming globular clusters at early times at a normal efficiency rate, but was then disrupted (perhaps by supernova-driven galactic winds, or by tidal shredding during infall; cf. the papers cited above) before its star formation could run to completion. The leftover gas now remains around these galaxies as their hot X-ray halos. This picture, however, assumes that the globular clusters form earlier than the bulk of the field stars in any given round of star formation – not an implausible requirement given the bulk of the observational evidence for starburst systems (see Harris 2000) and given that globular clusters emerge from the densest, most massive protocluster clouds.
McLaughlin (1999) defines a globular cluster formation efficiency, measured empirically as the mass ratio
$$ϵ=\frac{M_{cl}}{M_{}+M_{gas}}$$
where $`M_{}`$ and $`M_{gas}`$ are the masses within the galaxy in the form of visible stars and in the X-ray gas respectively. He finds that $`ϵ`$ is essentially identical in the well studied Virgo and Fornax systems M87, NGC 4472, and NGC 1399 (despite their very different $`S_N`$), providing evidence for a “universal” globular cluster formation efficiency $`ϵ0.26`$% relative to the initial protogalactic gas supply. The total mass ratio $`ϵ`$ is a more important indicator of cluster formation than $`S_N`$, which is only a measure of the cluster numbers (or equivalently total mass) relative to the galaxy light. In other words, $`S_N`$ is a measure of only the gas mass $`M_{}`$ that got converted to stars. Additional support for the near-universality of $`ϵ`$ in several other BCG’s is found by Blakeslee (1999) .
In this view, any high$`S_N`$ galaxy should then be surrounded by a massive X-ray gaseous component whether or not it is a centrally dominant giant. Notably, IC 4051 is indeed one of the few Coma ellipticals with an individually detected X-ray halo. Dow & White (1995), from ROSAT observations of the Coma core region, find that IC 4051 is detectable at the $`2\sigma `$ level in the soft X-ray range $`0.20.4`$ keV, but not in the higher $`0.42.4`$ keV range. If it were at the $`6.3`$ times closer distance of Virgo, IC 4051 would have a total $`L_X5\times 10^{41}`$ erg s<sup>-1</sup>. This level makes it quite comparable with the Virgo giant NGC 4472, which has $`L_X6\times 10^{41}`$ erg s<sup>-1</sup> in the soft X-ray regime (Fabbiano et al. (1992); Irwin & Sarazin (1996); Matsumoto et al. (1997); Buote & Fabian (1998)). However, this amount of X-ray gas corresponds to only $`5`$% of the stellar mass (McLaughlin (1999)) and NGC 4472, as expected, has only a “normal” specific frequency level $`S_N5`$.
With our adopted distance ratios for Virgo and Coma, we find that IC 4051 is about half as luminous as NGC 4472, so if it has a roughly similar amount of X-ray gas mass, this gas would only make up $`10`$% of its stellar mass. Along with $`S_N11`$, we find that these parameters convert to a present-day value for the mass ratio in IC 4051 of $`ϵ0.005`$, twice as large as McLaughlin’s (1999) fiducial value.
Nominally, it therefore seems that IC 4051 acts against the paradigm of a universal globular cluster formation efficiency. An obvious possibility, however, is that IC 4051 originally did possess much more gas shortly after its main era of globular cluster formation, but that most of this unused material was quickly stripped away as IC 4051 went through its first few passages of the Coma core. This gas would have joined the general reservoir of hot gas spread throughout the Coma potential well. The same mechanism which resulted in this galaxy’s compact structure might then have plausibly left it with the unusual combination of high $`S_N`$ and modest amount of X-ray gas that we now see.
A situation which would act much more strongly to falsify McLaughlin’s case for a universal $`ϵ`$ would be the opposite one: that is, a galaxy with a massive X-ray halo but a “normal” or subnormal $`S_N{}_{}{}^{<}\mathrm{\hspace{0.17em}4}`$. In such a case it would be much harder to avoid the conclusion that the formation efficiency of globular clusters was genuinely different (and low). Does the central Coma giant NGC 4874 present us with such a case? As we found in Paper II, NGC 4874 is not a high-$`S_N`$ system and is embedded within a very massive X-ray envelope. This X-ray gas is, however, so extended that must belong to the general Coma potential well as a whole, with no detectable concentrated component that can be associated with NGC 4874 itself (Dow & White (1995)). Thus there are ambiguities in the interpretation that are hard to circumvent. Better candidates would be E galaxies with massive X-ray halos that are not at the centers of rich clusters.
Finally, we may compare the interesting case of IC 4051 with that of its Coma neighbor NGC 4881 (Baum et al. (1995)), a giant E galaxy of similar location, size, and structure. Curiously, NGC 4881 holds a GCS of low specific frequency ($`S_N{}_{}{}^{<}\mathrm{\hspace{0.17em}2}`$) which appears to be almost entirely metal-poor, just the opposite of IC 4051. It has no significant amounts of X-ray gas (Dow & White (1995)). We speculate that NGC 4881 may have resulted from the merger of smaller galaxies in which these metal-poor globulars had already formed. These mergers should have been rather gas-poor to prevent the formation of newer and more metal-rich clusters. This is, however, an extremely sketchy interpretation, and there is an obvious problem with the much higher metallicity of the host galaxy light (how did the bulk of the giant E galaxy form at higher metallicity without leaving behind some metal-rich globular clusters? See Paper II for additional discussion).
The Coma ellipticals clearly present a wide range of GCS characteristics that strongly challenge the array of current galaxy formation models.
This research was supported through a grant to WEH from the Natural Sciences and Engineering Research Council of Canada. We thank Peter Macdonald of Ichthus Data Systems and the Department of Mathematics and Statistics at McMaster for providing the MIX multicomponent fitting code. Bill Baum provided several constructive comments on the first version of the text.
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# The SCUBA Local Universe Galaxy Survey I. First Measurements of the Submillimetre Luminosity and Dust Mass Functions
## 1 Introduction
With the advent of the SCUBA bolometer array (Holland et al. 1999) on the James Clerk Maxwell Telescope<sup>1</sup><sup>1</sup>1The JCMT is operated by the Joint Astronomy Center on behalf of the UK Particle Physics and Astronomy Research Council, the Netherlands Organization for scientific Research and the Canadian National Research Council. it has, for the first time, become feasible to map the submillimetre (100$`\mu `$m$`<\lambda <`$1 mm) emission of a large ($`>`$100) sample of galaxies. While this may seem a paltry feat in comparison to surveys conducted at other wavelengths, submillimetre astronomy has lagged behind the rest due to the technical difficulties of building sensitive enough instruments. Until recently, all that existed were a handful of submillimetre fluxes of nearby galaxies made with single element bolometers (Eales, Wynn-Williams & Duncan 1989; Stark et al. 1989; Clements, Andreani & Chase 1993). No information about the distribution of the dust was available and even the fluxes themselves were subject to considerable uncertainty, especially when beam corrections had to be made in order to compare with measurements in the far infra-red (FIR : 10$`\mu `$m$`<\lambda <`$100$`\mu `$m) made with the IRAS satellite. Quite surprisingly, even before SCUBA, there had been some successful submillimetre observations of high redshift objects (Dunlop et al. 1994; Isaak et al. 1994; Ivison 1995) but interpretation of the measurements, particularly the high implied dust masses, was hampered by a lack of knowledge of local submillimetre properties, such as luminosity and dust mass functions (Eales & Edmunds 1996, 1997 hereafter EE96,97). We have started to use SCUBA to carry out a survey of a large number of nearby galaxies and hope that this first statistical submillimetre survey will answer a number of questions.
### 1.1 How much dust does a galaxy contain?
The amount of dust in a galaxy is a measure of the quantity of heavy elements in the interstellar medium (ISM), since $``$ 50 per cent of the heavy elements are locked up in dust and there is some evidence that this fraction is constant from galaxy to galaxy (EE96). This is a different and complementary way of investigating the heavy element content from simply measuring the metallicity, which is of course, the mass of heavy elements per unit mass of gas not including that which is locked up in dust. EE96 have shown that the way in which the dust mass of a galaxy evolves depends critically on the mode of evolution:– whether the galaxy can be treated as a closed box or whether gas is inflowing onto or outflowing from it. Dust masses have been estimated from IRAS measurements but the resulting gas-to-dust ratios are a factor of $``$ 5–10 higher than for our own Galaxy (Devereux & Young 1990; Sanders et al. 1991) suggesting a problem with the method. The submillimetre offers a number of advantages for estimating dust masses. At longer submillimetre wavelengths ($`\lambda >350\mu `$m) we are sampling the Rayleigh-Jeans part of the Planck function, where the flux is least sensitive to temperature and most sensitive to the mass of the emitting material. This is a good thing because accurate dust temperatures ($`T_\mathrm{d}`$) are difficult to estimate. Dust radiates as a ‘grey body’, which is a Planck function modified by an emissivity term Q$`{}_{em}{}^{}\nu ^\beta `$, where the emissivity index ($`\beta `$) is believed to lie between 1 and 2 (Hildebrand 1983). To derive the temperature, an assumption must be made about $`\beta `$ or vice-versa. If an incorrect temperature is assumed at say 850$`\mu `$m, then the dust mass will be wrong by approximately the fractional error in $`T_\mathrm{d}`$, but at shorter FIR wavelengths (100$`\mu `$m) the error will be much larger as flux goes as $`T_\mathrm{d}^{(4+\beta )}`$. Generally, using IRAS measurements alone to determine the temperature and mass will always over-estimate the first and so under-estimate the second. The reason for this is that because of the sensitivity to temperature of emission on the Wien side of the grey-body curve, the FIR emission from a galaxy is dominated by warm dust, even when there is relatively little mass in this component. Thus fits of grey-body curves to FIR fluxes are biased towards higher temperatures than is warranted by the relative masses of warm and cold dust. This is, of course, the probable explanation of the high gas-to-dust ratios observed by Devereux & Young (1990), and also explains why it has been so difficult to demonstrate the presence of cold dust in galaxies. Longer wavelength studies (200$`\mu `$m – 850$`\mu `$m) using COBE, ISO and SCUBA (Sodroski et al. 1994; Reach et al. 1995; Alton et al. 1998a,b; Davies et al. 1999; Frayer et al. 1999; Papadopoulos & Seaquist 1999) have now confirmed the existence of cold dust components at 15 $`<T_\mathrm{d}<`$ 25 K in nearby spiral galaxies as predicted for grains heated by the general interstellar radiation field (Cox, Krügel & Mezger 1986), rather than in star forming regions.
There is one further difficulty in estimating dust masses at any wavelength:– our poor knowledge of the dust mass opacity coefficient $`\kappa _d(\nu )`$, which is needed to give absolute values for mass. This varies with frequency in the same way as the emissivity, and so current measurements at 120–200$`\mu `$m (Hildebrand 1983, Draine & Lee 1984) must be extrapolated to the wavelength of interest using the assumed value of $`\beta `$, raising the possibility of large errors since $`\beta `$ itself is poorly known. Hughes, Dunlop & Rawlings (1997) estimate the uncertainty in this coefficient to be a factor $``$8 at 850$`\mu `$m. For many issues, including galaxy evolution, absolute values of mass are not important (EE96) provided the relative masses are correct.
### 1.2 How much optical light is absorbed by dust?
There is an increasing need to understand the effects of dust on our optical view of the universe. Optical astronomers have recently shown that the UV luminosity density in the universe may rise from now to $`z1`$ and then fall off at higher redshifts (Lilly et al. 1996; Madau et al. 1996). There have been claims that this represents the ‘star-formation history of the universe’ but how much is this affected by dust? In general, how important are selection effects caused by dust, and is the high redshift universe significantly attenuated by foreground dusty objects? (Davies et al. 1997). Mapping large numbers of galaxies at submillimetre wavelengths will help determine how the dust is distributed relative to the stars, and how effective it is at absorbing optical light.
### 1.3 How much cosmological evolution is seen in the <br>submillimetre?
With the deep SCUBA surveys now taking place (Smail, Ivison & Blain 1997; Hughes et al. 1998; Barger et al. 1998,1999; Blain et al. 1999a; Eales et al. 1999; Lilly et al. 1999), there is now probably more information about the submillimetre properties of the distant universe than the local universe. The galaxies in the deep submillimetre surveys appear very similar to the ULIRGs found nearby (Smail et al. 1998; Hughes et al. 1998; Lilly et al. 1999). If the dust in these high-redshift objects is heated by young stars, as seems to be largely the case for nearby ULIRGs (Genzel et al. 1998) then, since these sources make up $`>`$20 per cent of the total extra-galactic background emission, this implies that $``$10 per cent of all the stars that have ever formed did so in in this kind of extreme object (Eales et al. 1999). Studies of the cosmological evolution of this population suggests that it is very similar to that seen at optical wavelengths in the Lilly-Madau curve. What is derived for the cosmological evolution, however, depends critically on what is assumed about the submillimetre properties of the local universe (cf. results of Eales et al. 1999 and Blain et al. 1999b). Currently, rather than working from a local submillimetre luminosity function, most investigations have perforce started from a local IRAS 60$`\mu `$m luminosity function and extrapolated to submillimetre wavelengths using some assumptions about the average FIR-submm SED. The lack of a direct submillimetre luminosity function has been a severe limitation when interpreting the results of deep surveys. Even before SCUBA, submillimetre observations of high-redshift radio galaxies and quasars have found unusually high dust masses. Here again, explanation has been restricted by our ignorance of the local universe, in this case the statistics of dust masses in galaxies (EE 96).
### 1.4 Is CO a good tracer of molecular hydrogen?
Carbon monoxide (CO) has long been used as a tracer for molecular gas (H<sub>2</sub>). The conversion factor between CO and H<sub>2</sub> (known as the ‘$`X`$’ factor) is quite uncertain and may be a function of physical environment such as density, temperature and metallicity (Maloney 1990; Wilson 1995). Dust is an alternative tracer of H<sub>2</sub>, and may help support the case for or against CO. Any variations of the ratio of CO to dust would provide insight on the way in which the $`X`$ factor depends on galaxy properties.
### 1.5 The scope of the survey
The ideal way to carry out the survey would be to do a blank field survey of large parts of the sky and measure the redshifts of all objects detected, a method which would be largely free of any selection effects. Unfortunately, this is impractical at the moment because the field of view of SCUBA is only $``$2 arcmins, making the time needed for such a survey prohibitively long. Instead, we decided to use the less ideal but more practical method of observing galaxies drawn from as many different complete samples selected in as many different wavebands as possible. This procedure still allows us to produce unbiased estimates of the submillimetre luminosity function and of the dust mass function using ‘accessible volume’ techniques (Avni & Bahcall 1980). These will not be biased unless there is a class of galaxies which is not represented in any of the samples. Provided at least one member of such a class of objects is present in one of the samples then the estimates will be unbiased although clearly with large random errors.
In this paper we present the results from a sample selected at 60$`\mu `$m. We present first estimates of the luminosity and dust mass functions, and examine the extent to which these may be biased. Subsequent papers on other samples, in particular a sample selected at optical wavelengths, will allow us to refine these estimates. We will also begin to address the other questions raised in this section.
Section 2 describes the observations and the methods used to reduce the data. In Section 3 we present the results from the IRAS-selected sample, in Section 4 we discuss the submillimetre luminosity and dust mass functions while in Section 5 we compare our results with other data (gas masses, optical properties) taken from the literature.
We assume a Hubble constant of 75 km s<sup>-1</sup> Mpc<sup>-1</sup> throughout.
## 2 OBSERVATIONS AND DATA REDUCTION
### 2.1 The Sample
The IRAS-selected sample was taken from the revised Bright Galaxy Sample (BGS), which is complete to a flux limit of $`S_{60}>5.24`$ Jy at all $`b>30^{}`$ and $`\delta >30^{}`$ (Soifer et al. 1989). We observed a subset of this sample with SCUBA, consisting of all galaxies with declination from $`10^{}<\delta <50^{}`$ and with velocity $`>`$ 1900 km s<sup>-1</sup>, a limit imposed to try to ensure the galaxies fitted within the SCUBA field of view. The SCUBA sample covers an area of $``$ 10,400 sq and contains 104 objects. Many of these are interacting pairs and although most were resolved by SCUBA they were not by IRAS, even with subsequent HIRES processing (Surace et al. 1993). Table LABEL:fluxT lists the sample.
### 2.2 Observations
We observed the galaxies between July 1997 and September 1998 using the SCUBA bolometer array at the 15-m James Clerk Maxwell Telescope (JCMT) on Mauna Kea, Hawaii. SCUBA has two arrays of 37 and 91 bolometers for operation at long (850$`\mu `$m) and short (450$`\mu `$m) wavelengths respectively. They operate simultaneously with a field of view of $``$ 2.3 arcmins (slightly smaller at 450$`\mu `$m). The arrays are cooled to 0.1 K and have typical sensitivities (NEFDs) of 90 mJy Hz<sup>-1/2</sup> at 850$`\mu `$m and 700 mJy Hz<sup>-1/2</sup> at 450$`\mu `$m (Holland et al. 1999). Beam sizes were measured to be $`15`$ arcsec and 8 arcsec at 850 and 450$`\mu `$m respectively, depending on chop throw and conditions. We made our observations in ‘jiggle-map’ mode, which is the most efficient mapping mode for sources smaller than the field of view. The arrangement of the bolometers is such that the sky is instantaneously under-sampled and so the secondary mirror is stepped in a 16-point pattern to ‘fill in the gaps’. For observations using both arrays, a 64 offset pattern is required to fully sample the sky as the spacing/size of the feedhorns is larger for the long wavelength array. Generally we used the 64-point jiggle, except in some cases where it was clear that the atmospheric conditions were too poor to obtain useful 450$`\mu `$m data, when we instead used the 16-point jiggle to give better sky cancellation.
The telescope has a chopping secondary mirror operating at 7.8 Hz which provides cancellation of rapid sky variations. To compensate for linear sky gradients (the effect of a gradual increase or decrease in sky brightness), the telescope is nodded to the ‘off’ position when in ‘jiggle’ mode. We used a chop throw of 120 arcsec in azimuth in all observations, except for galaxies with nearby companions, when we chose a chop direction which avoided the second galaxy. Skydips were performed regularly to measure the zenith opacity $`\tau _{850}`$. This varied over the course of the observations resulting in some data being taken in excellent conditions with $`\tau _{850}<0.2`$ while others done in ‘back-up’ mode had $`\tau _{850}`$ 0.5–0.6. Due to the large range of observing conditions, only about one third of the galaxies have useful 450$`\mu `$m data and, in fact, we will not consider the 450$`\mu `$m data further in this paper although they will be presented in another paper at a later date. We checked the pointing regularly and found it generally to be good to $`2`$ arcsec. Due to the uncertainty of the IRAS positions, we centred our observations on the optical coordinates taken from the Digitised Sky Survey (DSS)<sup>2</sup><sup>2</sup>2The Digitised Sky Surveys were produced at the Space Telescope Science Institute under U.S. Government grant NAGW-2166. The images of these surveys are based on photographic data obtained using the Oschin Schmidt Telescope on Palomar Mountain and the U.K. Schmidt Telescope. The plates were processed into the present digital form with the permission of these institutions..
Integration times differed depending on conditions and source strength. For bright sources in good weather we typically used 6 integrations (15 mins), but many sources had to be observed on more than one occasion due to poor S/N or because the galaxy fell on noisy bolometers. In total we integrated for $`44`$ hours, mostly in weather band 3 ($`0.3\tau _{850}0.5`$). We calibrated our data by making jiggle maps of planets (Uranus and Mars) and, when they were unavailable, of the secondary calibrators CRL 618 and HL Tau. Planet fluxes were taken from the JCMT fluxes program and we assumed that CRL 618 and HL Tau had fluxes of 4.56 Jy beam<sup>-1</sup> and 2.32 Jy beam<sup>-1</sup> respectively. The calibration error at $`850\mu `$m was taken to be 10 per cent (except for 7 galaxies observed in very poor conditions where 15 per cent was used instead. These are denoted by a * in Table LABEL:fluxT.). This calibration error is included in the flux error quoted in Table LABEL:fluxT.
### 2.3 Data reduction
We reduced the data using the standard surf package (Jenness & Lightfoot 1998). This consisted of flat-fielding the data after the on-off positions had been subtracted and then correcting for atmospheric extinction using opacities derived from skydip measurements. Any noisy bolometers were flagged as ‘bad’ and large spikes removed.
Notes to Table LABEL:fluxT:
(1) Most commonly used name taken from the IRAS BGS (2) Right ascension J2000 epoch (3) Declination J2000 epoch (4) Recession velocity taken from NED<sup>3</sup><sup>3</sup>3The NASA/IPAC Extragalactic Database (NED) is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration. (5) 60$`\mu `$m flux from Soifer et al. 1989 (7) 850$`\mu `$m flux (this work) (8) error on 850$`\mu `$m flux, calculated in the manner described in Section 2.4 and inclusive of a 10% calibration uncertainty. A * indicates that a 15% calibration uncertainty was used. (9) Dust temperature derived from a single component fit to the 60, 100 and 850$`\mu `$ data points as described in Section 3.1 (10) Statistical uncertainty in the dust temperature as described in Section 3.1 (11) Emissivity index derived fron the single component fit, described in Section 3.1 (12) Statistical uncertainty in the emissivity index as described in Section 3.1 <sup>p</sup> indicates a close or interacting pair which was resolved by SCUBA. Individual fluxes are listed in Table 2.
Except on the most dry and stable nights a residual sky noise was seen which was correlated across all the pixels. In most cases we removed this using the surf task remsky, which uses bolometers chosen by the user to define a sky region, and then subtracts this sky level from the other bolometers. The bolometers used to measure the ‘sky’ were chosen to be those away from the source emission using a rough submillimetre map and optical information as a guide. After sky noise removal, the data was despiked and regridded to form an image on 1 arcsec pixels. If an object had more than one dataset, a co-add was made in which each map was weighted as the inverse square of its measured noise.
The nodding and chopping should ideally leave a background with zero mean and some correlated noise superimposed on top. In reality the background level was not zero and could, in bad conditions, become quite large (either positive or negative). Often this residual sky level would vary linearly across the array giving a ‘tilted sky plane’, the surf task remsky, which was designed to remove the sky noise is rather simplistic and cannot remove these sky planes (Jenness, Lightfoot & Holland 1998). Due to the short integration times of most of our objects and the poor conditions in which many were observed, this problem of sky planes became the limiting factor in obtaining sufficient S/N in many observations. To this end a program was written to remove spatially varying sky backgrounds from the array as well as the temporal ‘sky noise’ dealt with by remsky.
Two fairly typical 850$`\mu `$m images are shown in Figure 1 along with optical images from the DSS. The SCUBA images clearly contain large amounts of structural information, which we will delay considering until a later paper. Here we will simply consider the 850$`\mu `$m flux measurements.
### 2.4 Flux measurement and error analysis
We used the submillimetre images to choose an aperture over which to integrate the flux, trying to ensure that it contained all the flux associated with the galaxy and as little of the surrounding sky as possible.
The noise /error on the flux measurement has two components.
i) The error in subtracting the background sky level which, despite our best attempts at removing it, often had a non-zero value that varied over the image. A number of smaller apertures were placed on regions containing no source emission and used to estimate the mean sky level. The standard deviation between the sky levels in each aperture was then used to find the error in the mean sky as $`\sigma _{ms}=S.D./\sqrt{\mathrm{No}.\mathrm{aps}}`$, where ‘No. aps’ is the number of sky apertures used to determine the mean sky value. The noise in the object aperture due to this sky error is then $`\sigma _{sky}=\sigma _{ms}N_{ap}`$ where $`N_{ap}`$ is the number of pixels in the object aperture.
ii) Shot/Poisson noise from pixel to pixel variations within the sky aperture, which would usually be determined from the standard deviation of the pixels within the sky aperture and contribute to the flux error as $`\sigma _{shot}=\sigma _{pix}\sqrt{N_{ap}}`$. In ordinary CCD astronomy this would be the only error estimated and we will refer to it as the ‘theoretical error’. It is actually quite a large underestimate of the true shot noise because the noise is correlated across pixels. This is because, unlike a CCD image, the pixels of a SCUBA image do not contain independent samples of the sky, i.e. two neighbouring pixels on a SCUBA image are constructed from overlapping sets of bolometer measurements, (also the pixels on the SCUBA maps are ‘artificial’ as they are smaller than the size of the beam). To quantify this effect, simulations of blank SCUBA fields were made by replacing the data for each bolometer in an observation by the output of a Gaussian random number generator. A large number of artificial maps were made in this way, and on each map the noise in an aperture of given size was measured in the traditional way using $`\sigma _{shot}=\sigma _{pix}\sqrt{N_{ap}}`$. This was compared with the standard deviation of all the fluxes (in that size aperture) between the different maps, deemed to be the ‘real shot noise’. For the range of aperture sizes actually used in measuring fluxes, the ratio of real to predicted noise was $`8`$. The total error for each flux was calculated from
$$\sigma _{tot}=(\sigma _{cal}^2+\sigma _{sky}^2+\sigma _{shot}^2)^{1/2}$$
where $`\sigma _{cal}10\%`$
$$\sigma _{sky}=\sigma _{ms}N_{ap}$$
$$\sigma _{shot}=8\sigma _{pix}\sqrt{N_{ap}}$$
To convert the aperture flux in volts to Janskys, a measurement of the calibrator flux for that night was made using the same aperture as was used for the object. The orientation of the aperture relative to the chop throw was also translated to the calibrator map (this did have a significant effect for more elliptical apertures).
Fluxes were found to have total errors $`\sigma _{tot}`$ in the range 10–25 per cent at 850$`\mu `$m and these are listed in Table LABEL:fluxT. All objects were detected at $`>3\sigma `$ except for IR 0857+39, which was only detected at $`2.4\sigma `$, although the flux is in good agreement with that measured by Rigopoulou et al. (1996). Where the system is interacting and not resolved by IRAS, the 850$`\mu `$m fluxes given are for both galaxies combined. For those paired galaxies which were resolved by SCUBA, the individual fluxes are listed in Table 2.
A few objects were detected with good signal to noise on more than one occasion and these could be used to keep check on the consistency of our flux measuring and calibration procedures. These objects and their relative flux errors are detailed in Table 3 and it can be seen that the differences in measured fluxes are within the errors calculated in the manner described above.
## 3 RESULTS
### 3.1 Spectral fits
We have fitted the IRAS 60, 100 and 850$`\mu `$m fluxes with a single-component temperature model by minimising the sum of the chi-squared residuals. The best fit was found for both the emissivity index $`\beta `$ and temperature $`T_\mathrm{d}`$, with the values given in Table LABEL:fluxT. Some examples of fits are shown in Fig. 2 along with the $`1\sigma `$ ranges. The uncertainty of each fitted value of $`T_\mathrm{d}`$ and $`\beta `$ was estimated in the following way. For each galaxy, a Gaussian random number generator was used to create 100 artificial flux sets from the original fluxes and measurement errors. These new data sets were then fitted in the same way and the standard deviation in the new parameters was taken to represent the uncertainty in the parameters found from the real data set. Average values and standard deviations (S.D.) for the whole sample (i.e. using only the best fitting parameters for the real data sets) are $`\overline{T_\mathrm{d}}=35.6\pm \mathrm{\hspace{0.17em}4.9}`$ K and $`\overline{\beta }=1.3\pm \mathrm{\hspace{0.17em}0.2}`$.
The value of $`\beta `$ is remarkably uniform (the overall S.D. for the whole sample is of the same order as the individual uncertainties, derived from the above technique), and lower than the value of 1.5–2.0 obtained from multi-wavelength studies of our Galaxy (Masi et al. 1995; Reach et al. 1995), and also of NGC 891 (Alton et al. 1998b). The number of galaxies at $`\beta 1.5`$ is equivalent to those at $`\beta 1.1`$, and in fact, the distribution of $`\beta `$ about the mean is well represented by a Gaussian, with those galaxies at $`\beta >1.5`$ being the tail of that Gaussian. The distribution of $`\beta `$ values for the sample is shown in Fig. 3. In fact, when the fluxes of the Milky Way and NGC 891 are fitted in the same way as we have fitted the IRAS galaxies (i.e. a single temperature model using only the 60,100 and 850$`\mu `$m points) the $`\beta `$ value we find is 0.7. It is only when the full SED is used (which inlcudes fluxes between 200 and 800$`\mu `$m) that the colder component at $`<20`$ K can be identified, leading to the higher $`\beta `$ of 1.5–2. So rather than implying that $`\beta `$ truly does have a low value, our results probably indicate that there is dust in these galaxies at colder temperatures than is indicated by a single-component fit. The uniformity of the $`\beta `$ values in the sample could be indicative of similar cold component properties in all of the IRAS galaxies.
The lower $`\beta `$ of 0.7 determined in this way for the Milky Way and NGC 891, which are optically selected, may be due to them having larger fractions of colder dust than the IRAS galaxies. We can partially test this using preliminary data from the optically selected sample. Figure 4 shows a colour-colour plot for the IRAS galaxies (stars) plus the galaxies observed so far from the optical sample (circles). There is clear evidence that the optically selected galaxies occupy a different area of the space, implying that they are not merely the less luminous equivalent of IRAS galaxies. We can make a more thorough investigation when the optically selected sample is completed and we have 350 and 450$`\mu `$m data for the samples.
### 3.2 Dust masses
Using the 850$`\mu `$m flux and the temperature derived in the above way the dust mass is calculated from
$$M_\mathrm{d}=\frac{S_{850}D^2}{\kappa _d(\nu )B(\nu ,T_\mathrm{d})}$$
(1)
where $`\kappa _d(\nu )`$ is the dust mass opacity coefficient, $`B(\nu ,T_\mathrm{d})`$ is the value of the Planck function at 850$`\mu `$m for a temperature $`T_\mathrm{d}`$, and $`D`$ for $`\mathrm{\Omega }_0=`$1 is given by
$$D=\frac{2c}{H_0}\left(1\frac{1}{\sqrt{1+z}}\right)$$
(2)
The dust masses are listed in Table LABEL:massT, with values for individual members of pairs in Table 5. We have assumed a value of 0.077 m<sup>2</sup> kg<sup>-1</sup> for $`\kappa _d(\nu )`$, which is intermediate between values for graphite and silicates as given by Draine & Lee (1984) and by Hughes et al. (1993), and also that $`\kappa _d(\nu )`$ has the same value at 850$`\mu `$m for all our galaxies. Even though the value of $`\kappa _d(\nu )`$ is notoriously uncertain, the relative values of our dust masses will be correct as long as the dust has similar properties in all galaxies. The uncertainties in the relative dust masses then depend only on the errors in $`S_{850}`$ and $`T_\mathrm{d}`$. The formal statistical uncertainties in $`T_\mathrm{d}`$ are only $``$ a few K, but there is the possibility that our assumption of a single temperature has biased our estimates of $`T_\mathrm{d}`$ to higher values, and thus our mass estimates to lower values. We can estimate the size of this effect in the following way: A single temperature fit will produce a lower value for $`\beta `$ than is actually true if more than one component is present. In their COBE–FIRAS study of our Galaxy, Reach et al. (1995) concluded that the highest observed value of $`\beta `$ in any region was likely to be closest to the true value, with any excess emission in other regions due to a colder component. If we adopt the same approach and assume that all our galaxies have a true $`\beta `$ of 2 (our highest observed value is 1.9) and a uniform cold dust temperature of 20 K, we can then fit a two-component temperature model and calculate a new dust mass. We find that for the same mass coefficient as before, the masses of dust estimated from this two-component model are between 1.5 and 3 times higher than those derived from a single temperature fit. The discrepancy is larger for galaxies with lower $`\beta `$ and higher temperatures and some galaxies which have steep SEDs and cold temperatures have little extra dust when modelled in this way (e.g. NGC 5371, NGC 772). The dust masses calculated assuming this form for the cold component are listed alongside the single temperature masses in Table LABEL:massT. They will be referred to as ‘cold dust masses’ in future.
### 3.3 Gas masses
Neutral hydrogen (Hi) fluxes were taken from the literature<sup>4</sup><sup>4</sup>4See notes to Table LABEL:massT and converted to masses (in solar units) using
$$M_{\mathrm{HI}}=2.36\times 10^5D^2S_{\mathrm{HI}}$$
(3)
where $`D^2`$ (Eqn. 2) is in Mpc and $`S_{\mathrm{HI}}`$ is in Jy km s<sup>-1</sup>
Molecular gas (H<sub>2</sub>) masses were calculated using CO fluxes taken from the literature<sup>3</sup> and scaled to telescope independent units (Jy km s<sup>-1</sup>). When an object appeared in more than one reference, an average flux was used to determine the mass. Conversion factors used for different telescopes are given in Table 6. Molecular masses in M are given by
$$M_{\mathrm{H}_2}=1.1\times 10^4D^2S_{\mathrm{CO}}$$
(4)
(Kenny & Young 1989), in which the parameters have the same units as in Eqn. 3. This assumes a CO to H<sub>2</sub> conversion factor of $`X=2.8\times 10^{20}`$ H<sub>2</sub> cm<sup>-2</sup>/\[K($`T_R`$)km s$`{}_{}{}^{1}]`$. Atomic and molecular gas masses are also given in Table LABEL:massT.
### 3.4 Optical luminosities
Blue magnitudes were taken from the Lyon-Meudon Extra-galactic Database (LEDA) and converted to luminosities using $`M_B=5.48`$. Blue luminosities are given in Table LABEL:massT and are corrected for Galactic extinction but not for internal extinction or inclination effects. In future discussions involving $`L_\mathrm{B}`$, mention will be made of ‘corrected’ blue luminosities. This will refer to the values in Table LABEL:massT after further correction for the internal effects of dust in the galaxy, and for inclination, following the prescription given in the 3rd Reference Catalogue of Bright Galaxies (RC3) (de Vaucouleurs et al. 1991).
### 3.5 Far Infra-Red Luminosities
The FIR luminosity ($`L_{\mathrm{fir}}`$) is usually calculated using the 60 and 100$`\mu `$m IRAS fluxes, as described in the Appendix of Catalogued Galaxies and Quasars Observed in the IRAS Survey (Version 2, 1989):
$$FIR=1.26\times 10^{14}(2.58S_{60}+S_{100})$$
and
$$L_{\mathrm{fir}}=4\pi D^2\times FIR\times C$$
Notes to Table LABEL:massT
(1) 60$`\mu `$m luminosity. (2) 850$`\mu `$m luminosity. (3) FIR luminosity calculated by integrating the measured SED from 40-1000$`\mu `$m. (4) Dust mass calculated using a single temperature, derived from fitting the 60,100 and 850$`\mu `$m fluxes. (5) Dust mass ($`M_\mathrm{d}^{cold}`$) calculated using a two-component temperature fit to the data, assuming a cold $`T_\mathrm{d}=20`$ K and $`\beta =2`$. (6) Hi refs:- Bottinelli et al. (1990), Huchtmeier & Richter (1989), Theureau et al. (1998) (7) H<sub>2</sub> refs:- Young et al. (1995), Solomon et al. (1997), Sanders et al. (1991), Chini, Krügel & Lemke (1996), Maiolino et al. (1997), Casoli et al. (1996), Lavezzi & Dickey (1998), Sanders et al. (1986), Sanders & Mirabel (1985). (8) Blue luminosity calculated from B<sub>T</sub> taken from the LEDA database and corrected for galactic extinction but not for internal extinction or inclination effects. (9) $`G_\mathrm{d}`$ is the gas-to-dust ratio calculated from Hi \+ H<sub>2</sub> and the single temperature dust mass. (10) The FIR luminosity per unit gas mass (molecular); often used as a measure of the star formation efficiency of a galaxy. <sup>p</sup> NGC 7469 includes masses for IC 5283, except for Hi which is for NGC 7469 only.
where $`D`$ is defined in Eqn. 2 and $`C`$ is a colour correction factor which depends on the ratio of $`S_{60}/S_{100}`$ and the assumed value of the emissivity index. The correction factor, designed to account for emission outside the IRAS bands, varies between 1.3 and 2.4 and is explained in more detail by Helou et al. (1988). Having submillimetre fluxes, we can use our derived temperatures and $`\beta `$ to integrate the total flux under the SED directly out to 1000$`\mu `$m, hopefully leading to more accurate values for $`L_{\mathrm{fir}}`$ since no general assumptions are being made. We integrate from $`401000\mu `$m, since if we were integrating to shorter wavelengths we should really include the IRAS 12 and 25$`\mu `$m data points which would require a multi-component temperature model, beyond the scope of what we are doing here. However, for our objects the ratio of $`S_{60}/S_{25}`$ is greater than 2.4, meaning that the contribution to the integral at $`\lambda <40\mu `$m is not very significant (if we integrate our SED with only the three data points out to 1$`\mu `$m we find increases in $`L_{\mathrm{fir}}`$ of only a few per cent). Our values for $`L_{\mathrm{fir}}`$ are similar to those calculated using the standard IRAS formulation described above, although slightly lower ($`10`$ per cent) than if the appropriate value of $`C`$ were chosen for an emissivity index of one. This is probably because calculating a temperature from the 60/100 flux ratio tends to overestimate the temperature and therefore the correction factor required. It must be noted that the true SED of these galaxies may not be well represented by a single temperature model and the values of $`L_{\mathrm{fir}}`$ would change accordingly. Our integrated values for $`L_{\mathrm{fir}}`$ are listed in Table LABEL:massT.
## 4 LUMINOSITY AND DUST MASS FUNCTIONS
The luminosity function (LF) is given by
$$\mathrm{\Phi }(L)\mathrm{\Delta }L=\underset{i}{}\frac{1}{V_i}$$
(5)
where $`\mathrm{\Phi }(L)\mathrm{\Delta }L`$ is the number density of sources (Mpc<sup>-3</sup>) in the luminosity range $`L`$ to $`L+\mathrm{\Delta }L`$, the sum is over all the sources in the sample in this luminosity range, and $`V_i`$ is the accessible volume of the ith source in the original sample (Avni & Bahcall 1980). For this sample, the accessible volume is the maximum volume in which the object could be seen and still be in the IRAS Bright Galaxy Sample and so to calculate this, the 60$`\mu `$m luminosity of each galaxy is used. In Equation 5, however, the luminosity is the luminosity at 850$`\mu `$m, the wavelength of our survey. The volume from $`cz=0`$ to $`cz=`$ 1900 km s<sup>-1</sup> is not included in the calculation of $`V_i`$ as galaxies with velocities less than this were excluded from our sample. $`\mathrm{\Phi }(L)`$ has been normalised to dex<sup>-1</sup> by dividing by $`\mathrm{\Delta }L`$. The dust mass function is estimated in the same way as the luminosity function but substituting dust mass for luminosity in Eqn. 5. This was done for both values of the dust mass; using a single temperature, and assuming a colder component at $`T_\mathrm{d}=20`$ K. The luminosity function at 850$`\mu `$m and both of the dust mass functions are shown in Fig. 5, and tabular forms for the functions are given in Table 7.
The BGS contains many close pairs which are resolved at 850$`\mu `$m but not at the IRAS wavelengths. A few do have HIRES fluxes (Surace et al. 1993) and, when the resolved 60$`\mu `$m flux is still above the 5.24 Jy limit of the sample, can be treated as separate objects for the purposes of calculating the luminosity and dust mass functions. The effect on the luminosity function of separating sources in this way would be to steepen it, as one luminous source with a large accessible volume becomes two less luminous sources with smaller volumes. We attempted to quantify this in the following way. For galaxies without HIRES fluxes, we estimated 60$`\mu `$m fluxes for the individual galaxies in each pair using the FIR-radio correlation (Helou et al. 1985). We removed any galaxy whose flux fell below the flux limit of the sample and re-calculated the luminosity function. Figure 6 shows this luminosity function compared with our original one (calculated using the sum of the 850$`\mu `$m emission from a pair of galaxies unresolved at 60$`\mu `$m). There is only a significant difference in the highest luminosity bin, which suffers from noise anyway due to the small number of objects. In future, we will refer only to the luminosity function constructed using the combined pair fluxes (Fig. 5a), since this can be most readily compared with the 60$`\mu `$m LF.
The luminosity function estimator (Eqn. 5) is unbiased providing there is no population of submillimetre emitting galaxies which have a high space density, yet are completely absent from the original IRAS sample. The only conceivable type of galaxy to which this could apply would be a hypothetical ‘cold’ population with $`T_\mathrm{d}<25`$ K. To investigate this possibility, let us assume that there are two populations of galaxies with equal space densities and 850$`\mu `$m luminosities, one with dust temperatures at 20 K and the other at 35 K. The relative numbers in an 850$`\mu `$m flux-limited survey would be given by
$$\frac{N_A}{N_B}=\frac{\mathrm{\Phi }_A}{\mathrm{\Phi }_B}\times \frac{V_A}{V_B}\frac{\mathrm{\Phi }_A}{\mathrm{\Phi }_B}\frac{(L_{850,A})^{3/2}}{(L_{850,B})^{3/2}}$$
where A and B refer to the two populations, $`\mathrm{\Phi }`$ is the space density, $`V`$ is the ‘accessible volume’ in which one of the galaxies would have been detected by the survey, and $`L_{850}`$ is the luminosity at 850$`\mu `$m. Given the assumptions we have made, $`N_A/N_B=1`$ and therefore equal numbers of the two populations should be found by the survey.
Now let us consider the relative numbers that would be found by a 60$`\mu `$m flux-limited survey (the BGS). Replacing the 850$`\mu `$m luminosity in the above equation with 60$`\mu `$m luminosity, we now predict that the ratio $`N_{35K}/N_{20K}`$ will be $`724`$. Clearly, in a sample of $`100`$ IRAS galaxies, finding a member of this hypothetical cold population would be unlikely and so our attempt to construct the 850$`\mu `$m LF from our sample would be an underestimate. Making a similar calculation for galaxies at $`T_\mathrm{d}=25`$ K and 35 K we now expect to find $``$ 32 times as many warm galaxies as cold ones, so in our sample we would expect to see around 3 objects at $`T_\mathrm{d}=25`$ K. There are actually two galaxies with $`T_\mathrm{d}25`$ K in our sample, which is consistent with the idea that we could be missing a significant cold population.
Is there truly a possibility that there is a missing population of galaxies? An important point to realise is that our dust temperatures are merely those corresponding to a best-fitting single-component model and are not actual dust temperatures. One way to assess the possibility that we are missing galaxies is to carry out precisely the same fitting procedure for optically-selected galaxies. For example, NGC 891 has been studied at many submillimetre wavelengths (Alton et al. 1998b), but if we throw away all the measurements except those at 60, 100 and 850$`\mu `$m, and carry out our fitting procedure, we obtain $`T_\mathrm{d}=34`$ K and $`\beta =0.7`$. While this temperature is similar to the average temperature for our sample, indicating that there is a warm component in NGC 891, the value of 0.7 for $`\beta `$ is significantly lower than for most of our galaxies suggesting a large amount of colder dust, and when the full submillimetre data-set was used the bulk of the dust in NGC 891 was found to be at 15 K. Consider now our own Galaxy : if the IRAS FIR and ARGO submillimetre data (Sodroski et al. 1989; Masi et al. 1995) for the Milky Way is combined and fitted in our usual way we find a temperature of 28 K with a $`\beta `$ of 0.7, so again the SED of the Milky Way would not necessarily have excluded it from our sample. The question of whether we are missing cold galaxies remains open (primarily how many cold galaxies there are to miss). To address this properly, and to constrain the luminosity function better at low luminosities we will need to complete the survey of an optically-selected sample, which will not be biased by temperature selection as is the BGS.
The 60$`\mu `$m luminosity function of the BGS from Soifer et al. (1987), can be compared to the 850$`\mu `$m LF by making an extrapolation based on an assumption about a galaxy’s SED from the far-infrared to the submillimetre. In fact, all of the 60$`\mu `$m LFs using IRAS data are consistent with each other (with the exception the one from Saunders et al. (1990), see Lawrence et al. 1999) and so it does not really matter which one we use. We have used the Soifer et al. one as it represents the same sample of galaxies which we observed. In the past this extrapolation has been done by assuming a single temperature and $`\beta `$ for all galaxies, which for 60 and 850$`\mu `$m gives the following
$$L_{850}=L_{60}\times \left(\frac{\nu _{850}}{\nu _{60}}\right)^{3+\beta }\times \left(\frac{e^{240.2/T_\mathrm{d}}1}{e^{16.8/T_\mathrm{d}}1}\right)$$
(6)
Extrapolations of the 60$`\mu `$m LF using a range of plausible values of $`\beta `$ and $`T_\mathrm{d}`$ are shown over the measured 850$`\mu `$m LF in Fig 7. As one would expect, an 850$`\mu `$m luminosity function obtained by extrapolating in wavelength from a 60$`\mu `$m LF is highly dependent on the values assumed for $`T_\mathrm{d}`$ and $`\beta `$. In practice, luminosity functions obtained in this way have had much lower amplitudes than the one we actually measure, essentially because dust temperatures deduced from IRAS data alone are always higher than the temperatures we estimate here. This can have significant consequences when trying to model the evolution required to fit the observed submillimetre number counts from deep surveys and the submillimetre background and the implications will be discussed in a future paper (Eales et al. in prep). The other difference in the luminosity functions is the slope: the luminosity function we measure is steeper than the ones extrapolated from the 60$`\mu `$m LF. The difference in slope between the two LFs is due to the correlation of dust temperature with 60$`\mu `$m luminosity (Fig. 8). As the temperature changes with 60$`\mu `$m luminosity, extrapolating every $`L_{60}`$ to 850$`\mu `$m with an average $`T_\mathrm{d}`$ will underestimate the 850$`\mu `$m luminosity at lower $`L_{60}`$ and overestimate it at high $`L_{60}`$. Using the fitted dependences of $`L_{60}`$ on $`T_\mathrm{d}`$ and also $`\beta `$ on $`T_\mathrm{d}`$ (see Table 8) we again extrapolate the 60$`\mu `$m LF and this time the agreement with the 850$`\mu `$m LF is much better (Fig. 9).
The 850$`\mu `$m luminosity and dust mass functions are well fitted by Schechter functions of the form (Press & Schechter 1974; Schechter 1975)
$$\mathrm{\Phi }(L)dL=\varphi (L)\left(\frac{L}{L_{}}\right)^\alpha e^{(L/L)}dL/L_{}$$
The best fitting parameters for the 850$`\mu `$m LF and the dust mass functions along with the reduced chi-squared values ($`\chi _\nu ^2`$) for the fits, are given in Table 7. The $`\chi ^2`$ contours showing the joint confidence intervals on $`\alpha `$ and $`L_{}`$ for all the functions are shown in Figs. 10(a–c).
There are several points to note here. Firstly, the 60$`\mu `$m LF cannot be fitted by a simple Schechter function (Lawrence et al. 1986; Rieke & Lebofsky 1986) as the high luminosity end does not fall off steeply enough. The fact that the 850$`\mu `$m LF and the dust mass function can, suggests that the Schechter function is surprisingly universal and that in the case of the 60$`\mu `$m LF, the Schechter function exponential fall-off is concealed by the sensitivity of the 60$`\mu `$m emission to dust temperature. Secondly, the slope of the 850$`\mu `$m LF is steeper than $`2`$ ($`\alpha =2.18`$ ) at lower submillimetre luminosities. This is significant because if this slope continued to zero luminosity, the submillimetre sky would be infinitely bright (a submillimetre Olbers Paradox). Of course, this means that the slope must flatten out at lower luminosities not yet probed by our survey. The shape of the dust mass function is strongly affected by the temperature distribution assumed, i.e. all of the dust at a single temperature, or most of the dust at some colder temperature with a small warmer component being responsible for the 60$`\mu `$m flux. The shape of the ‘cold dust mass’ function bears more resemblance to the 850$`\mu `$m LF simply because of our assumption of a common universal temperature for the cold component. We do believe however, that the true distribution of the dust masses lies somewhere between the two; the temperature of a cold component may not be the same in all galaxies but will probably vary less than the temperature of the warm components and this degree of variation along with the relative amounts of cold and warm dust will determine how similar the dust mass function is in shape to the 850$`\mu `$m LF.
## 5 DISCUSSION
### 5.1 Correlations
#### 5.1.1 Dust temperature
Luminosity is often strongly correlated with mass in astrophysical situations, and so at first sight the $`L_{60}T_\mathrm{d}`$ correlation (Fig. 8) suggests that the more luminous systems are hotter. We can test this by plotting dust temperature versus dust mass and blue luminosity, often used as a measure of the mass of a galaxy. In the first case (Fig. 11) there is no significant correlation (see Table 8 for coefficients), so there is no tendency for galaxies with lots of dust to be hotter, and in the second case (Fig. 12) there is an inverse correlation. The $`L_{60}T_\mathrm{d}`$ correlation is therefore most naturally explained by the sensitivity of the 60$`\mu `$m flux to the dust temperature. If, for whatever reason (and not because of its mass) the dust is hotter, this will greatly increase the 60$`\mu `$m luminosity. What is rather more interesting (and perhaps surprising) is the anti-correlation of the corrected $`L_\mathrm{B}`$ versus $`T_\mathrm{d}`$ (Fig. 12). The significance is at the 98 per cent level, making it the least secure of the correlations we present here and the probable link is with the star formation efficiency rather than dust temperature itself. Using $`L_{\mathrm{fir}}/M_{\mathrm{H}_2}`$ as a measure of star formation efficiency, Young (1999) has recently found that the star formation efficiency of galaxies in many different types of environment decreases as galaxy size increases (as measured by the optical linear diameter $`D_{25}`$). If we now plot $`L_{\mathrm{fir}}/M_{\mathrm{H}_2}`$ against $`L_\mathrm{B}`$ for our galaxies a much better correlation is seen than with $`T_\mathrm{d}`$, indicating that this is where the true dependence lies (see Fig. 13 and Table 8). The explanation Young gives for this relationship is that larger galaxies have flatter rotation curves over more of their disks and so experience more shear. She suggests that the effect of this would be to reduce the ability of molecular clouds to form stars by increasing the turbulence within them, and to dampen the effects of star-formation triggers.
#### 5.1.2 Gas and dust masses
We compared the dust masses, derived from the submillimetre fluxes using a single temperature, to the H<sub>2</sub>, Hi and H<sub>2</sub> \+ Hi masses as shown in Figure 14a-c. The strongest correlation is with the molecular gas which suggests that the dust is primarily found in molecular clouds. The larger scatter on the Hi plots may be because we are using global Hi values which may include gas at very large radii where there would be much less dust than in the inner disk. Devereux & Young (1990) found that if they only used the Hi in the inner disk, in addition to the molecular gas, then the correlation between gas and dust was improved. We cannot, however, investigate the spatial connection between dust and atomic gas in the region they both occupy as we do not yet have Hi maps. We repeated the process using the ‘cold dust masses’ but found no differences in the correlations or slopes of the plots and so have not included these.
We would like to try to address the reasons behind the tightness of the relationship between dust and molecular gas. The scatter in this plot (Fig. 14a) is very small, the r.m.s dispersion about the line of best fit is only $`60`$ per cent. Some of this is due to observational errors; if we assume 30 per cent errors in the fluxes and subtract this in quadrature from the overall value then we are left with 52 per cent as the intrinsic scatter. There are some observational-type errors which are hard to account for, such as aperture corrections and the probability that our assumption of a single dust temperature is wrong, but the small scatter does suggest that any aperture effects must be similar for all the galaxies and that if our assumption about temperature is wrong then it must be wrong in the same way for all of the galaxies. This is borne out be the similarity of the relationship when the ‘cold dust masses’ are used.
In order to determine the physical meaning of the 50 per cent scatter we will review the current views on the formation of dust and CO and the expected dependences on metallicity.
There are no firm conclusions about how dust grains form. Suggested sites of formation range from post-AGB stars to supernova remnants (Whittet 1992). The grains are thought to grow in the darkest molecular clouds by accretion of icy mantles and are probably destroyed by sputtering in the diffuse ISM, after the dispersion of the molecular cloud by supernova shocks and HII regions. Despite the complications in determining how and where the various stages of dust evolution occur, there is some evidence that the fraction of metals bound up in dust grains is a constant, both from observations of nearby galaxies (Issa et al. 1990) and from observations in our own Galaxy, where the carbon and oxygen abundances (the main constituents of dust) are constant over a wide range ISM densities (Cardelli et al. 1996; Meyer et al. 1998). This suggests that $`M_\mathrm{d}ZM_g`$ where $`M_g`$ is the total mass of gas and $`Z`$ is the metallicity.
The formation of CO is better understood than that of dust and is believed to occur through a network of gas-phase reactions in the giant molecular clouds (GMCs). CO will form readily, the rate increasing with density, but is easily destroyed by UV radiation. In order to remain intact, the CO needs to be shielded by dust grains and/or large column densities of H<sub>2</sub>. Assuming that the dominant mass in the GMCs is in the form of H<sub>2</sub>, then since CO forms from the C and O available in the gas-phase, (and O is always more abundant than C in the ISM), a metallicity dependence is suggested where $`L_{\mathrm{CO}}ZM_{\mathrm{H}_2}`$. However, since most of the shielding required to form the CO is produced by dust rather than H<sub>2</sub>, the true metallicity dependence may be even higher than this. When there is insufficient dust in a molecular cloud to provide adequate shielding from the UV radiation field, the volume of the cloud occupied by CO may not be as large as that occupied by H<sub>2</sub>, because H<sub>2</sub> is more efficient at self-shielding. This is the likely case in low metallicity systems such as the LMC or where the volume and/or column densities are low. The possibilities that CO and H<sub>2</sub> may not always be co-extensive, and that there may be a dependence of CO luminosity on metallicity, can lead to problems with the application of the $`X`$ factor to CO observations in order to determine the mass of H<sub>2</sub>. For systems with significantly lower densities and/or metallicities a larger conversion may be needed relative to galaxies such as the Milky Way. Observations do indicate a dependence of $`X`$ on metallicity for low metallicity systems (Wilson 1995) with $`X/X_G=(Z/Z_{})^{0.7}`$ for systems with less than solar metallicity and staying roughly constant for higher than solar metallicities (Frayer & Brown 1997). This observed dependence is less steep than one would have imagined from the above simplistic arguments. It could be argued that for metallicities above some critical value, where there is sufficient shielding, the $`X`$ factor may show less dependence on metallicity, as the abundance of CO relative to H<sub>2</sub> depends on the availability of C and O and not on dust any longer. While the formation of CO should still be greater for higher metallicities, H<sub>2</sub> formation is also dependent on dust grains as catalysts so it is possible that this could cause the $`X`$ factor to saturate : more metals lead to more CO, but they also lead to more dust and thus more H<sub>2</sub> formed from atomic hydrogen, producing the observation that the $`X`$ factor is constant with metallicity above a threshold value. Theory does predict significant depletion of CO onto dust grains (forming mantles) for favourable conditions in the cloud which would make $`X`$ smaller with increasing metallicity, possibly offsetting any trend for an increase but in the Milky Way this depletion is only observed to be 5 to 40 per cent (Whittet 1992), much less than predicted, indicating that there is still much which is not understood about the interaction of dust and molecules in the GMCs.
We do not have metallicity measurements for our galaxies, but the range of metallicities of spiral galaxies of similar absolute magnitudes is 2–3 (Henry & Worthey 1999). Thus the small scatter in the H<sub>2</sub>–dust mass diagram strongly suggests that the metallicity dependence of dust mass and of the CO $`X`$ factor is very similar. An additional test of this will be extending our study to galaxies likely to have a larger range of metallicity.
#### 5.1.3 Star formation efficiency
The ratio of $`L_{\mathrm{fir}}/M_{\mathrm{H}_2}`$ is often used as a measure of the star formation efficiency of a galaxy since $`L_{\mathrm{fir}}`$ is believed to trace the star formation rate and so this ratio gives the star formation rate per unit gas mass. Figure 15 shows the strong correlation between molecular gas and $`L_{\mathrm{fir}}`$ and gives some idea of the range of $`L_{\mathrm{fir}}/M_{\mathrm{H}_2}`$ values for these galaxies. The L/M = 4 line represents ‘normal’ galaxies such as the Milky Way (Scoville & Good 1989) and most of the galaxies in our sample show elevated star formation efficiencies relative to the Milky Way. Many of the galaxies lying above the general trend (SFE$`20`$) are often interacting or have suspected active nuclei (e.g. IR 0857+39 with SFE $`100`$). Values of $`L_{\mathrm{fir}}/M_{\mathrm{H}_2}`$ are given in Table LABEL:massT.
### 5.2 Gas to dust ratios
The average gas-to-dust ratios (using the single temperature dust masses) for the three components are $`M_{\mathrm{H}_2}/M_\mathrm{d}=293\pm 19`$, $`M_{\mathrm{HI}}/M_\mathrm{d}=304\pm 24`$ and $`M_{\mathrm{H}_2+\mathrm{HI}}/M_\mathrm{d}=581\pm 43`$, where the $`\pm `$ indicates the error in the mean. In future, we will refer to $`M_{\mathrm{H}_2+\mathrm{HI}}/M_\mathrm{d}`$ as the gas-to-dust ratio $`G_\mathrm{d}`$ and the values are listed in Table LABEL:massT. In order to compare our $`G_\mathrm{d}`$ with the Galactic value, we must be careful that we account for the differences in $`\kappa _d(\nu )`$ and $`X`$ used by others and ourselves. The gas-to-dust ratio for the Milky Way was measured by Sodroski et al. (1994), using COBE data at 140 and 240$`\mu `$m, therefore in order to compare our measurement with theirs we must also scale the $`\kappa _d(240)`$ value of the opacity coefficient they used to 850$`\mu `$m, using an assumed value for $`\beta `$. Sodroski et al. use a $`\kappa _d(240)=0.72\mathrm{m}^2\mathrm{kg}^1`$; taking this to 850$`\mu `$m gives $`\kappa _d(850)=0.20.057\mathrm{m}^2\mathrm{kg}^1`$ for $`\beta =12`$. Comparing this to our value of 0.077 m<sup>2</sup> kg<sup>-1</sup> and accounting for the differences in $`X`$ factors (theirs being 0.71 times ours) gives a range of Galactic $`G_\mathrm{d}`$ values of 85–302 for a $`\beta `$ of $`12`$. The middle of the range at $`\beta =1.5`$ is $`G_\mathrm{d}=160`$ and this is the nominal value we will make comparisons with. It can now be seen that our average gas-to-dust ratio of 581 is more than a factor of three larger than the Galactic value of 160 (at the least a factor of 2 larger if $`\beta =2`$). Of course, the measurement made by Sodroski et al. is a very difficult one to make, but it is supported by arguments based on element abundances and depletions which predict a ratio of $`130`$ (Eales et al. in prep). The Galaxy has a FIR luminosity of $`1.1\times 10^{10}\mathrm{L}_{}`$ and a dust mass of $`2.9\times 10^7\mathrm{M}_{}`$ (Sodroski et al. 1994) which assumes the same value for $`\kappa _d(\nu )`$ that gives a $`G_d=160`$. This places the Galaxy at the bottom end of the range of luminosities and dust masses found in this IR bright sample. Previous determinations of the gas-to-dust ratios in IR luminous objects using IRAS data have always been high $`1000`$ (Young et al. 1989; Devereux & Young 1990; Sanders et al. 1991) which has been attributed to the lack of sensitivity of IRAS to cold dust. It would seem that our submillimetre observations have reduced this tendency, giving a lower gas-to-dust value, but when the gas-to-dust ratios are re-calculated using the ‘cold dust masses’ we find a further reduction to $`M_{\mathrm{H}_2}/M_\mathrm{d}^{cold}=143\pm 9`$, $`M_{\mathrm{HI}}/M_\mathrm{d}^{cold}=158\pm 14`$ and $`M_{\mathrm{H}_2+\mathrm{HI}}/M_\mathrm{d}^{cold}=293\pm 24`$. This is now just within the range of values we calculated for the Galaxy, and in good agreement if the $`\beta =2`$ case is taken. We perform the calculations with the ‘cold masses’ simply to illustrate the effect of overlooking this component. A slightly colder temperature (say 15–18 K) would increase the dust masses by a larger amount, further reducing the observed gas-to-dust ratios and so giving better agreement with the Galactic values.
The gas-to-dust ratios calculated using the single temperature dust masses vary from 1700 to 200, the large range being mostly due to the variations in $`M_{\mathrm{HI}}/M_\mathrm{d}`$ as discussed in Section 5.1.2. Some of the larger values of $`M_{\mathrm{H}_2}/M_\mathrm{d}`$ could be due to the large angular sizes of some of the nearer galaxies, meaning that our dust masses are underestimated when compared to the CO fluxes produced from multi-positional single dish measurements (this particularly applies to NGC 772 and NGC 7479); however this is not always the case as there are some galaxies of small angular size which have high $`M_{\mathrm{H}_2}/M_\mathrm{d}`$ ratios but more normal $`M_{\mathrm{HI}}/M_\mathrm{d}`$, (e.g. IR 0335+15, UGC 4881 and Mkn 331). These galaxies are interacting or have Seyfert nuclei and typically have high dust temperatures, which may suggest that the lack of apparent dust mass is simply due to using a high temperature when calculating $`M_\mathrm{d}`$ and not accounting for any colder dust. High temperatures cannot be the sole explanation, however, as other galaxies with high $`T_\mathrm{d}`$ and signs of interaction/activity have rather low values of $`M_{\mathrm{H}_2}/M_\mathrm{d}`$ (e.g. Arp 220). It could be that an enhanced $`M_{\mathrm{H}_2}/M_\mathrm{d}`$ is a stage that some merging/starburst galaxies pass through, either by an enhancement of CO relative to dust, or a decrease in dust emission. A change in the $`X`$ factor, causing the H<sub>2</sub> to be overestimated could also be responsible (this has been suggested for many extreme active/starburst nuclei (Solomon et al. 1997)). It is also possible that CO which was depleted from the gas phase onto grain mantles in dark clouds could, at a certain stage in the starburst evolution, become vaporized back into the gas phase, giving a higher CO flux while keeping the atomic, dust and H<sub>2</sub> content the same. One particular case which deserves a mention here is NGC 7714 which has $`M_{\mathrm{HI}}/M_\mathrm{d}1300`$, causing its abnormally high gas-to-dust ($`G_\mathrm{d}`$) value of 1700 (see column 9 of Table LABEL:massT). NGC 7714 is interacting with NGC 7715 although they are widely spaced enough (1.9 arcmin) that only NGC 7714 was observed with SCUBA and the IRAS flux is due entirely to NGC 7714 (Surace et al. 1993). The Hi maps of Smith et al. (1997) show rings and bridges of gas connecting the two galaxies, so it is probable that some of the Hi listed for NGC 7714 in the single dish measurement (Table LABEL:massT) is really associated with the whole region and cannot be compared with what we observed. The amount of Hi contained in the 50 $`\times `$ 50 arcsec region centred on NGC 7714 (similar to the area of the dust emission) is given as only $`1.7\times 10^9`$ M compared to $`7\times 10^9`$ M for the whole system. Using this lower value would give a more reasonable value of 355 for $`M_{\mathrm{HI}}/M_\mathrm{d}`$, bringing $`G_\mathrm{d}`$ down to 802.
One further caveat to the whole $`M_{\mathrm{H}_2}/M_\mathrm{d}`$ business is the large discrepancies in the CO fluxes between single-dish measurements and interferometer mappings. Maps have been found for four of our galaxies (Sanders et al. 1988; Yun & Hibbard 1999; Laine et al. 1999) and always the flux is considerably less ($`1/2`$) than the single-dish measurement (Young et al. 1995). This is usually explained by resolution effects as these interferometer measurements are insensitive to structures larger than $`30`$ arcsec but even in sources unresolved by the interferometer at $`cz10,000`$ km s<sup>-1</sup> (which are not expected to have large scale angular structure), the difference in flux is still as large. If we used the interferometer measurements for NGC 520, NGC 7469, NGC 7479 and UGC 4881, the corresponding H<sub>2</sub> masses would decrease by factors of 2, 2, 4, and 2 and bringing the $`M_{\mathrm{H}_2}/M_\mathrm{d}`$ ratios down to 229, 195, 163 and 260 respectively.
## 6 CONCLUSIONS
We have undertaken the first statistical survey of the local universe with the SCUBA bolometer array on the JCMT. We present here the initial results from the first of the samples we are surveying: a sample of 104 galaxies selected at 60$`\mu `$m from the IRAS Bright Galaxy Sample.
(i) The 60, 100 and 850$`\mu `$m fluxes are well fitted by single temperature SEDs. The mean and standard deviation (S.D.) in the best-fitting dust temperature is $`\overline{T_\mathrm{d}}=35.6\pm \mathrm{\hspace{0.17em}4.9}`$ K with the mean and S.D. for the emissivity index being $`\overline{\beta }=1.3\pm \mathrm{\hspace{0.17em}0.2}`$. We do not however, rule out the possibility of colder dust and a steeper emissivity. If $`\beta =2`$ and we assume there is a cold dust component with a temperature of 20 K, we then obtain dust masses a factor of 1.5–3 higher. The 450$`\mu `$m data obtained for 30 per cent of the galaxies may eventually constrain the submillimetre emissivity index and therefore the presence of a cold component.
(ii) We have presented the first direct measurements of the submillimetre luminosity and dust mass functions. They are well fitted by Schechter functions, in contrast to the IRAS 60$`\mu `$m LFs (Lawrence et al. 1986; Rieke & Lebofsky 1986) which do not have the exponential cut-off required by the Schechter function. The slope of the 850$`\mu `$m LF at low luminosities is steeper than $`2`$ implying that the LF must flatten at lower luminosities than were probed by this survey. The optically selected sample currently being observed will take the submillimetre LF to lower luminosities and constrain the ‘knee’ of the LF. The shape of the dust mass function is affected by the model assumed for the temperature distribution.
(iii) We have shown that a simple extrapolation of the IRAS 60$`\mu `$m LF to 850$`\mu `$m using a single dust temperature and emissivity index does not reproduce the measured submillimetre LF, both in terms of normalisation and shape.
(iv) A correlation was found between the fitted dust temperature and 60$`\mu `$m luminosity. This is best explained by the sensitivity of the 60$`\mu `$m flux to temperature (due to its position on the Wien side of the grey-body curve) rather than by a dependence on galaxy mass. Accounting for this temperature dependence when extrapolating the 60$`\mu `$m LF to submillimetre wavelengths produces a much better match to the observed 850$`\mu `$m LF.
(v) If there is a population of cold ($`T_\mathrm{d}<25`$ K) galaxies which are also luminous submillimetre sources, then our submillimetre LF is likely to be biased, as these objects would not have appeared in the original IRAS bright galaxy sample which was selected at 60$`\mu `$m. The question of a missing population will be addressed by the optically selected sample.
(vi) We find an average value for the gas-to-dust ratio ($`G_\mathrm{d}`$) of $`581\pm 43`$ where the gas mass is taken to be $`M_{\mathrm{HI}}+M_{\mathrm{H}_2}`$. This is lower than previous values determined using IRAS fluxes alone ($`1000`$) indicating that the submillimetre is a better place to measure the dust mass. It is still, however, a factor of $`3.5`$ higher than the value obtained for the Galaxy (Sodroski et al. 1994). Using ‘cold dust masses’, calculated on the basis of there being a $`20\mathrm{K}`$ component in addition to the warm dust, $`G_\mathrm{d}`$ was reduced to 293. This supports the idea that the previous large discrepancies in $`G_\mathrm{d}`$ for our own Galaxy (using COBE data), and that of other galaxies (which use IRAS measurements), are due to a cold dust component which has gone undetected by IRAS.
(vii) The relationship between the mass of molecular hydrogen (as determined from CO observations) and dust mass shows very little intrinsic scatter ($`50`$ per cent) which implies that the CO to H<sub>2</sub> conversion, ‘$`X`$’, depends on metallicity in the same way as the dust mass.
(vii) The star formation efficiency as traced by $`L_{\mathrm{fir}}/M_{\mathrm{H}_2}`$ was found to be anti-correlated with galaxy size (as measured by blue luminosity), a result also found by Young (1999) using a different technique. The star formation efficiency for these galaxies is higher on average than that for the Galaxy by a factor of $`3`$ (13.3 for the sample compared to 4 for the Milky Way).
#### FUTURE WORK
Observations which would further our understanding of the properties of dust in galaxies and refine our initial estimates of the LF are as follows:
* Shorter wavelength (300 – 600$`\mu `$m) data from SCUBA and from SHARC on the CSO will allow us to ascertain the presence of any cold dust component in these galaxies and whether it is uniform (i.e. is the cold component similar in all the galaxies or does it vary as a function of some other property). Determining the existence of such a cold dust will allow a much more accurate estimate of the dust mass (and hence the dust mass function) to be made.
* In order to resolve some of the questions unanswered by this paper we will need to complete the survey of an optically selected sample. This will help to constrain the behaviour of the low luminosity end of the 850$`\mu `$m LF (i.e. the turn-over) and also establish whether there is a population of ‘cold’ galaxies missing from this IRAS sample. We will then be able to investigate the differences between the IRAS galaxies when compared to ‘normal’ optically selected galaxies in terms of their dust properties.
* The 850$`\mu `$m LF also requires more data points at the high luminosity end. This sample contained only a few very luminous objects because of the high flux limit for the BGS. Observations of a complete sample of ULIRGS at 850$`\mu `$m would prove very useful in this respect.
* More CO and HI measurements, for those galaxies in both the IRAS and optically selected samples which do not have them already, would increase the statistical significance of the gas and dust comparisons. In addition, detailed mapping in both HI and CO would be of great importance in determining the relationship of the dust to the different gas phases of the ISM. Work on obtaining this data is in progress.
* Optical imaging would enable a comparison of the optical structures with those in the submillimetre. Of importance is the radial profile and scale height of the dust compared to the stars. Extinction maps could be created from multi-colour optical data and compared with the submillimetre maps, a technique described by Trewhella (1998).
## ACKNOWLEDGMENTS
We would like to thank all of the support staff at the JCMT as well as any observers who took data on our behalf when the conditions were too poor for their own projects. In particular, we are grateful to Wayne Holland for the ’out of hours’ support after our runs were over. We also thank Paul Alton for useful discussions and information about NGC 891. The work of L. Dunne, S. Eales, M. Edmunds, R. Ivison and D. Clements is supported by PPARC.
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# The Formation of Galaxies, the Formation of Old Globular Clusters and the Link with High-Redshift Objects
## 1. Introduction
Statistics is always a key-point in scientific studies and it is no surprise that, as any scientist, the astronomer is looking for large, statistically significant samples to properly analyze the Universe and its content. Unfortunately, the intrinsic size of the Universe is turning this simple point into a difficult brain teaser due to the faintness of the above objects. Globular clusters are likely to contain some of the oldest known stellar populations of the Universe. As such, they potentially hold a cosmologically significant information on their formation and more generally on the conditions that prevailed more than 10 Gyrs ago.
In brief, we have started a study that is heading at selecting the oldest globular clusters (GCs) from the largest available sample of extragalactic GC systems. Previous studies often assumed the systems as an homogeneous population and used the mean properties (metallicity) of the GC systems as the main parameter. Only in the most recent works GCs have been split up in sub-populations (Forbes et al. 1997; Côté et al. 1998). Going back to the very formation of the galaxies (and maybe before), asks to make sure that only the oldest (i.e. reliable fossils) GCs are picked up. Our choice is to select the metal-poor GCs. Indeed, if we can find old metal-rich globular clusters (Ortolani et al. 1995; Puzia et al. 1999), only in very limited cases could we have a late formation of metal-poor GCs. An important step has been to discover that galaxies other than our own contain metal-poor sub-populations that can be associated to a halo component (Puzia et al. 1999).
A more detailed report of this work will be published elsewhere (Burgarella, Kissler-Patig & Buat, 2000).
## 2. The compilation of Old Globular Cluster Populations
Our goal is to select the oldest GCs around galaxies and to compare their metallicities with the host galaxy properties, as well as to compare the systems with each other. The compilation includes galaxies of all types, however spiral galaxies are under-represented while bright elliptical galaxies dominate the sample. The detection of several peaks in the metallicity distribution function is always a problem and we use the mixture-modeling algorithm (KMM) developed by Ashman et al. (1994) to detect and quantify the bimodality and estimate the mean metallicity of the metal-poor GC populations around the sampled galaxies. This compilation of 38 GCs systems includes galaxies of all types and our sample includes galaxies over 10 magnitudes in absolute brightness (see Burgarella et al. 2000).
## 3. Mean metallicity against galaxy luminosity
Before a clear separation of metal-poor and metal-rich populations could be performed in other galaxies than in the Milky Way, the mean metallicities of the whole GC system was thought to correlate with the galaxy luminosity (van den Bergh 1975; Brodie & Huchra 1991). Actually, this apparent correlation was mainly due to the fact that the the brightest galaxies are ellipticals which have, on average, a higher GC mean metallicity than spirals and dwarfs (e.g. Ashman & Zepf 1998; Gebhardt & Kissler-Patig 1999).
The sample of GC systems presented in this paper is the largest database to-date, and about 3 times more numerous than Forbes et al.’s (1997) initial dataset. The mean \[Fe/H\] lies at \[Fe/H\]$`=1.40\pm 0.06`$ with a dispersion of $`\sigma =0.24\pm 0.05`$, that is slightly more metal-poor on average than, and exhibiting a scatter similar to the Forbes et al. sample. Fig. 1 shows the relative percentage of GC systems within each bin of the metallicity function. Indeed, the immediately apparent result is that the mean-metallicities of metal-poor GCs are not distributed at random: most of them are lying around \[Fe/H\]$`1.4`$, with 64 % within -1.5 $`<`$ \[Fe/H\] $`<`$ -1.3 and 80 % within -1.7 $`<`$ \[Fe/H\] $`<`$ -1.1. We plot in Fig. 1 the metallicities of the GC systems as a function of the absolute magnitude M<sub>V</sub>. The average metallicity of the metal-poor GCs is constant over a very large range in absolute magnitude of the host galaxy ($`23<\mathrm{M}_V<16`$).
## 4. Consequences on the formation scenarios of GCs, globular cluster systems, and galactic halos
Following Fall & Rees (1988), GC formation models “can be classified as primary, secondary or tertiary depending on whether GCs are assumed to form before, during or after the collapse of proto-galaxies.”. It seems, however, that the borderline between the three classes is not always very clear. To better identify the origins of GCs, we prefer to split the GCs on whether they are external to the galaxy, and not associated with the final host galaxy, or whether they formed internally, i.e. are associated in some form with the final host galaxy. This terminology is relatively unambiguous if we specify that pre-galactic fragments are not considered to be galaxies. And since we consider only old, metal-poor GCs assumed to have formed before or early in the galaxy formation process, we do not take into account mergers of already formed galaxies.
Now, if we concentrate our attention on the Milky Way GC system, it seems that halo GCs of the Milky Way host (at least) two populations that Zinn (1993) distinguished from their horizontal branch types. He called them ’old’ and ’younger’ halo GCs. There are hints for a similar differentiation in M33 (Ashman & Bird 1993). A probable internal halo old GC population which would have formed internally in the early galaxy lifetime by a dissipative collapse in a few Myr, and an external halo population which would have formed around other satellite galaxies and accreted afterwards. If such a complex GC formation history is valid for our own Galaxy and its nearest neighbors, it cannot be ruled out for other galaxies either.
We could retain that : i) the mean metallicity of halo GCs is independent of the host galaxy properties (M<sub>V</sub>, type, environment (Fig. 2), metallicity) and ii) halo GC populations have very similar mean metallicities in all galaxies. These two points can be added to the dynamical information available for a number of metal-poor outer globular clusters that tends to show that these clusters are on tangentially biased orbits, as opposed to radially biased orbits expected if they had formed in a collapse (Eggen et al. 1962).
The bottom line from the above facts, is that the early cluster and star formation was remarkably homogeneous in the local universe (within several tens of Mpc). The first collapsing fragments were extremely similar in mass and abundances over large scales and collapse in very similar fashions independently of the potential well (dark halo) in which they were located. Presumably, the distinction between galaxy types only appeared after the first formation of stars and clusters in fragments.
## 5. Time and sites of formation of the metal-poor GCs
### 5.1. The measurement of metallicity at high redshift
In this section, we will look for measurements of metallicity at high redshift in order to compare with the average metallicity of our old GCs. Indeed, if the GC formation is the first stellar formation episode of what will become a new galaxy, the first-formed stars might have kept a memory of the genuine intergalactic medium as it was before the galaxy formation. Among the objects observed at high redshifts for which the metallicity can be estimated, two of them seem of interest: Damped Ly$`\alpha `$ systems (DLAs), and Lyman Break Galaxies (LBGs). Both the typical column density in HI and the observed metallicities for DLAs and LBGs are plotted in Fig. 3, together with the same quantities for GCs (and the Lyman $`\alpha `$ forest for completeness). DLAs (Pettini et al. 1997) give a measurement for the metallicity as a function of redshift of high density neutral gas objects. Lyman Break Galaxies (Steidel et al. 1996) can, in addition, be used to estimate the star formation rate as a function of redshift. Assuming a metal ejection rate (Pettini et al. 1997), we can infer a chemical evolution of the Universe and compare it with other estimates, in particular with the mean metallicity of the metal-poor globular cluster systems.
However, \[Fe/H\] may not be a reliable estimate of the metallicity of DLAs, since some Fe may be locked up in dust and thus the measured \[Fe/H\] too low. Pettini et al. (1997) showed that \[Zn/H\] is a more reliable estimator because it essentially measures the metallicity independently of dust depletion. From a \[Zn/H\] analysis of 34 DLAs, Pettini et al. (1997) showed that z $`>`$ 1 DLAs are generally metal-poor (log ($`<Z>/Z_{}`$) $`<`$ -1.0) with a possible trend for z $`>`$ 3 DLAs towards a lower metallicity. However, the value at z = 3 is an upper limit and we would need better high redshift values. Although the dust depletion problem may make the direct use of DLA \[Fe/H\] questionable, we try here to use its variation with redshift in order to compare it with the information from GC systems. The compilation is given in Burgarella et al. (2000). The \[Zn/H\] and \[Fe/H\] variations as a function of the redshift are plotted in Fig. 3. We use the column-density weighted abundances : $`[<\mathrm{M}/\mathrm{H}_{\mathrm{DLA}}>]=\mathrm{log}<(\mathrm{M}/\mathrm{H})_{\mathrm{DLA}}>\mathrm{log}(\mathrm{M}/\mathrm{H})_{}`$ where $`<(\mathrm{M}/\mathrm{H})_{\mathrm{DLA}}`$ (M = Fe or M = Zn) and the associated standard deviations as defined in Pettini et al. (1997).
The data presented in Fig. 3 can be used to constrain the GC system formation. In the first place, the analysis of the CMDs of old Galactic GCs suggests the age of halo GCs to be more than 10 Gyr which corresponds to a GC formation not later than z $``$ 1.6 (H$`{}_{0}{}^{}=50`$ km.s<sup>-1</sup>.Mpc<sup>-1</sup> and q$`{}_{0}{}^{}=0.5`$). On the other hand, the chemical evolution of DLAs and LBGs is below the lower limit for 80 % of our metal-poor globular clusters at z $``$ 4. The conclusion suggested by these data is that the GC formation occurred in average in the redshift range 1.6 $`<`$ z $`<`$ 4 (i.e approximately in the range 10 $`<`$ age (Gyrs) $`<`$ 12 with the assumed cosmology).
From the above discussion we retain that DLAs (and LBGs) have approximately the same range of metallicities and are observed in the redshift range expected for the formation of metal-poor GCs. Note, however, that DLAs contain neutral gas while LBGs are star-forming objects. As already suggested by Fynbo et al. (1999), we may wonder whether we are not observing the same objects at different location in space or in time. For instance, DLAs would be the source of dense gas out of which old GCs formed while LBGs would be star-forming regions e.g. spheroids as proposed by Giavalisco et al. (1996) and Steidel et al. (1996) but also surrounding fragments in the same potential well which are only directly visible at high redshift when the star formation turns on. Eventually these fragments would be accreted by the large galaxy to produce a MW-like object.
#### Acknowledgments.
We would like to thank K. Gebhardt for his help in handling the data of the metal poor GCs and M. Pettini for helful discussions.
## References
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Ashman K.M., Zepf S.E. 1998, Globular Cluster Systems, Cambridge Univ. Press
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Barger A.J., Cowie L.L., Sanders D.B. 1999, ApJ 518, 5
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Fall S.M., Rees M.J. 1988, in ”Globular Clusters in Galaxies”, IAUS 126, eds. Grindlay J.E. and Davis A.G., p.323
Forbes D.A., Brodie J.P., Grillmair C.J. 1997, AJ 113, 1652
Fynbo J.U., Moller P., Warren S.J. 1999, MNRAS 305, 849
Gebhardt K., Kissler-Patig M. 1999 AJ 118, 1526
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Ortolani S., Renzini A., Gilmozzi R., Marconi G., Barbuy B., Bica E., Rich R.M. 1995, Nat. 377, 701
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Steidel C.C., Adelberger K.L., Giavalisco M., Dickinson M., Pettini M. 1999, ApJ 519,1
van den Bergh S. 1975, IAU Symp. 58, ”The Formation and Dynamics of Galaxies”, J.R. Shakeshaft (ed.), p.157
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# Conditions satisfied by characteristic polynomials in fields and division algebras
## 1. Introduction
Let $`E/F`$ be a field extension of degree $`n`$ and $`det:EF`$ be the norm function. For $`xE`$, we define $`\sigma ^{(i)}(x)`$ by
(1.1)
$$det(\lambda 1_Fx)=\lambda ^n+\sigma ^{(1)}(x)\lambda ^{n1}+\mathrm{}+\sigma ^{(n1)}(x)\lambda +\sigma ^{(n)}(x).$$
In particular, $`\sigma ^{(1)}(x)=\mathrm{tr}(x)`$ and $`\sigma ^{(n)}(x)=(1)^ndet(x)`$. In the sequel, whenever we write $`\sigma ^{(i)}(x)`$, we shall always understand $`i`$ to be an integer between $`1`$ and $`n`$. If the reference to the extension $`E/F`$ is not clear from the context, we will sometimes write $`\sigma _{E/F}^{(i)}(x)`$ in place of $`\sigma ^{(i)}(x)F`$.
If $`A`$ is a central simple algebra of degree $`n`$ with center $`F`$ then we can define $`\sigma ^{(i)}=\sigma _{A/F}^{(i)}`$ in the same way. Here $`det`$ in formula (1.1) should be intepreted as the reduced norm in $`A_FF(\lambda )`$.
A number of interesting results, both in the theory of polynomials and in the theory of central simple algebras, can be stated in terms of the existence (or nonexistence) of nontrivial solutions to systems of equations of the form
(1.2)
$$\sigma ^{(i)}(x)=0\text{for}i=i_1,\mathrm{},i_r.$$
###### Example 1.1.
(Hermite \[H\], Joubert \[J\]; see also Coray \[C\]) If $`E/F`$ is a field extension of degree $`5`$ or $`6`$ and $`\mathrm{char}(F)3`$ then there exists an element $`xE`$ such that $`E=F(x)`$ and $`\sigma ^{(1)}(x)=\sigma ^{(3)}(x)=0`$.
In classical language, this means that for $`n=5`$ or $`6`$ every polynomial $`f(t)=t^n+a_1t^{n1}+\mathrm{}+a_nF[t]`$ can be reduced, via the Tschirnhaus transformation $`tx`$, to the form $`f(t)=t^n+b_1t^{n1}+\mathrm{}+b_nF[t]`$ with $`b_1=b_3=0`$; for details we refer the reader to \[BR\].
###### Example 1.2.
Let $`A`$ be a central simple algebra of degree $`n`$ whose center contains a primitive $`n`$th root of unity. Then $`A`$ is cyclic iff there exists an element $`x`$ such that
$$\sigma ^{(1)}(x)=\mathrm{}=\sigma ^{(n1)}(x)=0.$$
A conjecture of Albert asserts that every $`A`$ of prime (or, equivalently, square-free) degree is cyclic. This conjecture is known to be true for $`n=2`$, $`3`$ and $`6`$ (see \[Ro<sub>1</sub>, Section 3.2\]); the remaining cases are open.
###### Example 1.3.
(Haile \[Ha\]; see also Brauer \[Ro<sub>3</sub>, Proposition 7.1.43\]) Suppose $`A`$ is a central simple algebra of degree $`n`$ with center $`F`$. Then there exists an $`(n1)`$-dimensional $`F`$-subspace $`W`$ of $`A`$ such that $`\sigma ^{(1)}(x)=\sigma ^{(n1)}(x)=0`$ for any $`xW`$.
###### Example 1.4.
(Rowen \[Ro<sub>2</sub>, Corollary 5\]) If $`A`$ is a central simple algebra of odd degree with center $`F`$ then there exists an element $`xA\{0\}`$ such that $`\sigma ^{(1)}(x)=\sigma ^{(2)}(x)=0`$.
Note that if $`\mathrm{char}(F)2`$, this follows easily from a theorem of Springer (see e.g., \[Re<sub>1</sub>, Remark 14.3\]); however, the above result is true even if $`\mathrm{char}(F)=2`$.
In \[Re<sub>1</sub>\] the first author showed that in many cases equations of the form $`\sigma ^{(i)}(x)=0`$ or $`\mathrm{tr}(x^i)=0`$ and systems of the form $`\sigma ^{(1)}(x)=\sigma ^{(i)}(x)=0`$ or $`\mathrm{tr}(x)=\mathrm{tr}(x^i)=0`$ do not have nontrivial solutions. In particular, the theorem of Hermite and Joubert, cited in Example 1.1, fails for field extensions of degree $`n=3^m`$ or $`3^m+3^l`$, with $`m>l0`$. In this paper we revisit this subject from a more geometric point of view.
### Notational conventions
Throughout this paper $`n`$ will denote the degree of the field extension or division algebra we are considering, and $`\mathrm{sqf}(n)`$ will denote the square-free part of $`n`$. We will always work over a fixed ground field $`k`$.
Let $`K`$ be a field containing a primitive $`r`$th root of unity $`\zeta _r`$ (in particular, we assume that $`r`$ is prime to $`\mathrm{char}(K)`$), and let $`z,wK`$. Recall that a symbol algebra $`(z,w)_r`$ is defined as
(1.3)
$$(z,w)_r=K\{x,y\}/(x^r=z,y^r=w,yx=\zeta _rxy);$$
cf. \[Ro<sub>3</sub>, p. 194\]. We now define the algebra $`D_n`$ as follows. Write $`n=p_1\mathrm{}p_s`$ as a product of (not necessarily distinct) primes. Let $`K=k(z_1,w_1,\mathrm{},z_s,w_s)`$, where $`z_1,w_1,\mathrm{},z_s,w_s`$ are independent variables over $`k`$ and let
(1.4)
$$D_n=(z_1,w_1)_{p_1}_K\mathrm{}_K(z_s,w_s)_{p_s}.$$
Note that $`D_n`$ is a division algebra of degree $`n`$ and exponent $`\mathrm{sqf}(n)`$, with center $`K`$.
Finally recall that the universal division algebra $`\mathrm{UD}(n)`$ is the subalgebra of $`\mathrm{M}_n(k(s_{ij},t_{ij}))`$ generated, as a division algebra, by two generic $`n\times n`$-matrices $`(s_{ij})`$ and $`(t_{ij})`$. Here $`s_{ij}`$ and $`t_{ij}`$ are $`2n^2`$ independent variables over $`k`$. For details of this construction, see, e.g., \[Ro<sub>1</sub>, Section 3.2\].
### Main results
###### Theorem 1.5.
Suppose $`\mathrm{char}(k)n!`$ and $`D=D_n`$ or $`\mathrm{UD}(n)`$. Then the system
(1.5)
$$\{\begin{array}{c}\sigma ^{(i)}(x_1)=\mathrm{}=\sigma ^{(i)}(x_m)\hfill \\ \sigma ^{(j)}(x_1\mathrm{}x_m)=0\hfill \end{array}$$
has no nontrivial solutions in $`D`$, provided that $`i`$ and $`m`$ are divisible by $`\mathrm{sqf}(n)`$.
Here, as usual, a solution $`(x_1,\mathrm{},x_s)`$ is trivial if $`x_1=\mathrm{}=x_s=0`$ and nontrivial otherwise. Note that the assertion of the theorem for $`\mathrm{UD}(n)`$ is a formal consequence of the assertion for $`D_n`$, because of the specialization property of $`\mathrm{UD}(n)`$. However, our proof will treat the two cases in parallel, since both are proved by the same argument. Theorem 1.5 can be generalized in several directions; some generalizations are discussed at the end of Section 4.
We now record three consequences of Theorem 1.5, which we feel deserve a special mention.
###### Corollary 1.6.
Suppose $`\mathrm{char}(k)n!`$, $`D=D_n`$ or $`\mathrm{UD}(n)`$, and $`m`$ is divisible by $`\mathrm{sqf}(n)`$.
(a) $`\sigma ^{(m)}(x)0`$ for any $`xD\{0\}`$.
(b) If $`det(x_1)=\mathrm{}=det(x_m)`$ for some $`x_1,\mathrm{},x_mD\{0\}`$ then $`\mathrm{tr}(x_1\mathrm{}x_m)0`$.
(c) $`D`$ does not have an element of (reduced) norm 1 and (reduced) trace 0.
To prove part (a), we assume the contrary and substitute $`i=m`$, $`x_1=x`$ and $`x_2=\mathrm{}=x_m=0`$ into (1.5) to obtain a contradiction. To prove part (b), we apply Theorem 1.5 with $`i=n`$ and $`j=1`$. Finally, if $`det(x)=1`$ then setting $`x_1=x`$ and $`x_2=\mathrm{}=x_m=1`$ in part (b), we obtain $`\mathrm{tr}(x)0`$, thus proving part (c). ∎
The commutative counterpart of the universal division algebra is the general field extension $`L_n/K_n`$ defined as follows:
(1.6) $`K_n=k(a_1,\mathrm{},a_n)\text{and}L_n=K_n[x]/(x^n+a_1x^{n1}+\mathrm{}+a_n),`$
where $`a_1,\mathrm{},a_n`$ are algebraically independent indeterminates over $`k`$.
###### Theorem 1.7.
Let $`n_1`$ and $`n_2`$ be positive integers, and $`L_n/K_n`$ be the general field extension of degree $`n=n_1+n_2`$. Then the system of equations
(1.7)
$$\mathrm{tr}(x^{m_1})=\mathrm{tr}(x^{m_2})=0$$
has no nontrivial solutions $`xL_n^{}`$, provided that
(i) $`n_1n_20`$ and $`(\frac{n_2}{n_1})^{m_2m_1}1`$ in $`k`$.
(ii) each $`\mathrm{sqf}(n_i)`$ ($`i=1,2`$) divides $`m_1`$ or $`m_2`$ (and possibly both).
Note that if $`\mathrm{char}(k)=0`$ then condition (i) holds unless $`m_1=m_2`$ or $`n_1=n_2`$ and $`m_2m_1`$ is even. If we replace (i) by a more complicated condition, we can also show that the system $`\sigma ^{(m_1)}(x)=\sigma ^{(m_2)}(x)=0`$ has no nontrivial solutions; see Section 6.
It is interesting to note that Theorem 1.5 and Corollary 1.6 remain true if $`D`$ is replaced by $`L_n`$; see Remark 4.6. On the other hand, Theorem 1.7 fails if $`L_n`$ is replaced by $`\mathrm{UD}(n)`$; see Remark 5.2.
All of the main results in this paper are proved by the same general method. based on the Going Down Theorem 2.1. This method is outlined in Section 2. In particular, our proofs of Theorems 1.5 and 1.7, given in Sections 4 and 5, are applications of Propositions 2.4 and 2.2 respectively. Proposition 2.4 says a system of equations, such as (1.5), has no nontrivial solutions in a “sufficiently generic” division algebra if a certain projective $`\mathrm{PGL}_n`$-variety, constructed from this system, does not have $`H`$-fixed points for some abelian subgroup $`H`$ of $`\mathrm{PGL}_n`$. Proposition 2.2 gives a similar criterion for nonexistence of solutions in field extensions. Other applications of this approach and some generalizations are presented in Sections 68.
### Acknowledgements
The authors would like to thank A. R. Wadsworth for helpful discussions.
## 2. The Going Down Theorem and its applications
The following result will play a key role in the sequel. A simple proof, due to Kollár and Szabó, can be found in \[RY<sub>1</sub>, Appendix\]. Assume that $`k`$ is an algebraically closed base field, and that all varieties, group actions and maps are defined over $`k`$.
###### Theorem 2.1 (The Going Down Theorem).
Let $`H`$ be a finite abelian group acting on algebraic varieties $`X`$ and $`Y`$ and let $`f:XY`$ be an $`H`$-equivariant rational map. If $`X`$ has a smooth $`H`$-fixed point and $`Y`$ is projective then $`Y`$ has an $`H`$-fixed point. ∎
### $`\mathrm{S}_n`$-varieties
Let $`L/K`$ be a separable field extension of degree $`n`$, let $`L^{}`$ be the normal closure of $`L`$ over $`K`$, and $`\mathrm{Gal}(L^{}/K)=G`$. Note that $`G`$ acts on the set of embeddings $`LL^{}`$ and thus is naturally realized as a transitive subgroup of $`\mathrm{S}_n`$. For each $`i=1,\mathrm{},n`$ choose $`g_i\mathrm{S}_n`$ such that $`g_i(1)=i`$. The embedding of $`G`$ in $`\mathrm{S}_n`$ defines a (permutation) action of $`G`$ on $`𝔸^n`$ and thus a diagonal actions on $`(𝔸^n)^m`$ for every $`m1`$.
Let $`P(x_{11},\mathrm{},x_{1n};\mathrm{};x_{m1},\mathrm{},x_{mn})k[(𝔸^n)^m]`$ be a $`G`$-invariant polynomial and let $`a_1,\mathrm{},a_mL`$. Then we can define $`P(a_1,\mathrm{},a_m)`$ as $`P(a_{11},\mathrm{},a_{1n};\mathrm{};a_{m1},\mathrm{},a_{mn})`$, where $`a_{ij}=g_j(a_i)L^{}`$. A priori, $`P(a_1,\mathrm{},a_m)L^{}`$; however, since $`P`$ is $`G`$-invariant polynomial, $`P(a_1,\mathrm{},a_m)`$ actually lies in $`(L^{})^G=K`$.
In the sequel we shall assume that $`K`$ is finitely generated over $`k`$ (and hence, so are $`L`$ and $`L^{}`$).
###### Proposition 2.2.
Let $`Y`$ be the subvariety of $`((𝔸^n)^m)`$ given by $`G`$-invariant homogeneous polynomial equations $`P_1=\mathrm{}=P_s=0`$. Suppose that $`Y`$ does not have $`H`$-fixed points for some abelian subgroup $`HG`$. Assume that there exists a $`G`$-variety $`X`$ which has a smooth $`H`$-fixed point and such that $`k(X)=L^{}`$ as fields with $`G`$-action. Then the system of equations
(2.1)
$$P_1(a_1,\mathrm{},a_m)=\mathrm{}=P_s(a_1,\mathrm{},a_m)=0$$
has no nontrivial solutions in $`L`$.
We remark that if $`\mathrm{char}(k)=0`$ then a $`G`$-variety $`X`$ such that $`k(X)=L^{}`$ (as $`G`$-fields) always exists; see \[Re<sub>2</sub>, Proposition 8.6 and Example 8.4c\]. Moreover, we can choose $`X`$ to be smooth and projective; see \[RY<sub>2</sub>, Proposition 2.2\]. In view of Theorem 2.1, the presence of an $`H`$-fixed point on such an $`X`$ is a birational invariant, i.e., is independent of the choice of the (smooth projective) model.
###### Proof.
Suppose $`(a_1,\mathrm{},a_m)L^mk(X)^m`$ is a non-trivial solution of (2.1) and let $`a_{i1},\mathrm{},a_{in}`$ be the conjugates of $`a_i`$ in $`L^{}`$. Then
$$f:x[a_{11}(x):a_{12}(x):\mathrm{}:a_{mn}(x)]$$
is a $`G`$-equivariant rational map $`X((𝔸^n)^m)`$. By our choice of $`a_1,\mathrm{},a_m`$, the image of $`f`$ lies in $`Y`$. Applying Theorem 2.1 to the rational map $`f:XY`$, we conclude that $`Y`$ has an $`H`$-fixed point, a contradiction. ∎
In the sequel we shall use use Proposition 2.2 only for $`m=1`$; the statement for general $`m`$ is intended to make it parallel to Proposition 2.4 below.
### $`\mathrm{PGL}_n`$-varieties
Let $`Pk[(\mathrm{M}_n)^m]^{\mathrm{PGL}_n}`$; it is a polynomial in the entries of $`m`$ matrices $`U_1,\mathrm{},U_m`$ invariant under simultaneous conjugation. If $`A`$ is a central simple algebra of degree $`n`$ and $`a_1,\mathrm{},a_mA`$ then we can define $`P(a_1,\mathrm{},a_m)`$ as follows. Split $`A`$ by the algebraic closure $`\overline{K}`$ of $`K`$: $`A_K\overline{K}\mathrm{M}_n(\overline{K})`$. Thus $`A\mathrm{M}_n(\overline{K})`$, and we can evaluate $`P(a_1,\mathrm{},a_m)\overline{K}`$.
###### Lemma 2.3.
$`P(a_1,\mathrm{},a_m)`$ lies in $`K`$ and is independent of the choice of the isomorphism $`A_K\overline{K}\mathrm{M}_n(\overline{K})`$.
###### Proof.
Any two choices of the isomorphism $`A_K\overline{K}\mathrm{M}_n(\overline{K})`$ differ by conjugation by some $`g\mathrm{PGL}_n(\overline{K})`$. Since $`P`$ is $`\mathrm{PGL}_n`$-invariant, conjugation by $`g`$ does not change the value of $`P(a_1,\mathrm{},a_m)`$.
Consider the action of $`\mathrm{Gal}(\overline{K}/K)`$ on $`\mathrm{M}_n(\overline{K})`$; for any $`\sigma \mathrm{Gal}(\overline{K}/K)`$ and $`B_1,\mathrm{},B_n\mathrm{M}_n(\overline{K})`$, $`P(\sigma (B_1),\mathrm{},\sigma (B_n))=\sigma (P(B_1,\mathrm{},B_m))`$. The composition
$$\mathrm{M}_n(\overline{K})\stackrel{}{}A_K\overline{K}\mathrm{@}>\mathrm{Id}\sigma ^1>>A_K\overline{K}\stackrel{}{}\mathrm{M}_n(\overline{K})\mathrm{@}>\sigma >>\mathrm{M}_n(\overline{K})$$
is an automorphism of $`\mathrm{M}_n(\overline{K})`$ whose restriction to the center $`\overline{K}`$ is trivial. Hence, this composition is given by conjugation by some $`g\mathrm{PGL}_n(\overline{K})`$. It follows that for $`a_1,\mathrm{},a_mA`$, $`P(a_1,\mathrm{},a_m)`$ is fixed by $`\mathrm{Gal}(\overline{K}/K)`$ and thus lies in $`K`$. ∎
Note that the Lemma is an immediate consequence of the fact that $`k[(\mathrm{M}_n)^m]^{\mathrm{PGL}_n}`$ is generated by elements of the form $`\sigma ^{(i)}(U)`$, where $`U`$ is a monomial in the $`m`$-matrices $`U_1,\mathrm{},U_m`$. The latter was proved by Sibirskii \[Si\] and Procesi \[P1\] in the case $`\mathrm{char}(k)=0`$ and, more recently, by Donkin \[D\] in prime characteristic. The elementary argument given above allows us to avoid appealing to this more difficult result.
Next we recall that is $`F`$ be a finitely generated field extension of $`k`$ then an element of $`H^1(F,\mathrm{PGL}_n)`$ may be interpreted either as a central simple algebra $`D`$ of degree $`n`$ with center $`F`$ or, alternatively, as a generically free $`\mathrm{PGL}_n`$-variety $`X`$ such that $`k(X)^{\mathrm{PGL}_n}=F`$. It is shown in \[Re<sub>1</sub>\] (under the assumption $`\mathrm{char}(k)=0`$) that $`D\stackrel{}{=}\mathrm{RMaps}_{\mathrm{PGL}_n}(X,\mathrm{M}_n)`$ = the algebra of $`\mathrm{PGL}_n`$-equivariant rational maps from $`X`$ to $`\mathrm{M}_n`$; see also \[RY<sub>2</sub>, Section 3\]. Note that the above isomorphism is an isomorphism of $`F`$-algebras, where we identify $`fF=k(X)^{\mathrm{PGL}_n}`$ with the $`\mathrm{PGL}_n`$-equivariant rational map $`X\mathrm{M}_n(k)`$ given by $`xf(x)I_n`$. (Here $`I_n`$ denotes the $`n\times n`$-identity matrix.)
###### Proposition 2.4.
Let $`Y`$ be the subvariety of $`((\mathrm{M}_n)^m)`$ cut out by $`\mathrm{PGL}_n`$-invariant homogeneous polynomial equations $`P_1=\mathrm{}=P_s=0`$. Suppose $`Y`$ has no fixed points for some finite abelian subgroup $`H`$ of $`\mathrm{PGL}_n`$. Then the system of equations
(2.2)
$$P_1(x_1,\mathrm{},x_m)=\mathrm{}=P_s(x_1,\mathrm{},x_m)=0$$
has no nontrivial solutions in any central simple algebra $`D`$ of the form $`D=\mathrm{RMaps}_{\mathrm{PGL}_n}(X,\mathrm{M}_n)`$, where $`X`$ is a generically free $`\mathrm{PGL}_n`$-variety which has a smooth $`H`$-fixed point.
###### Proof.
Suppose the system (2.2) has a nontrivial solution $`(x_1,\mathrm{},x_m)`$. As $`D=\mathrm{RMaps}_{\mathrm{PGL}_n}(X,\mathrm{M}_n)`$, each $`x_i`$ can be interpreted as a rational $`\mathrm{PGL}_n`$-invariant map $`X\mathrm{M}_n`$; collectively, these elements define a rational $`\mathrm{PGL}_n`$-equivariant map $`f:X((\mathrm{M}_n)^m)`$. By our choice of $`x_1,\mathrm{},x_m`$, the image of this map lies in $`Y`$. By Theorem 2.1, $`Y`$ has a $`H`$-fixed point, a contradiction. ∎
## 3. Abelian subgroups
In order to use Propositions 2.2 and 2.4, we need a description of abelian subgroups $`H`$ of $`\mathrm{S}_n`$ and $`\mathrm{PGL}_n`$. In this section we introduce the abelian subgroups that will be used in subsequent applications.
We shall assume that the base field $`k`$ contains all roots of unity. For a finite abelian group $`A`$ of order prime to $`\mathrm{char}(k)`$, we shall denote its dual group $`\mathrm{Hom}(A,k^{})`$ by $`A^{}`$.
### Abelian subgroups of $`\mathrm{S}_n`$
Let $`A=\{a_1,\mathrm{},a_n\}`$ be an abelian group of order $`n`$. The right multiplication action of $`A`$ on itself gives rise to an embedding
$$\psi _A:A\mathrm{S}_n.$$
(Note that if we relabel the elements of $`A`$, $`\psi _A`$ will change by an inner automorphism of $`\mathrm{S}_n`$.) Given a character $`\chi :Ak^{}`$, let
$$R_\chi =(\chi (a_1),\mathrm{},\chi (a_n)).$$
It is easy to see that $`k^n=_{\chi A^{}}\mathrm{Span}_k(R_\chi )`$ is a decomposition of $`k^n`$ as a direct sum of 1-dimensional character spaces for the permutation action of $`A`$ on $`k^n`$ (via $`\psi _A`$); moreover, the character associated to $`\mathrm{Span}_k(R_\chi )`$ is precisely $`\chi ^1`$.
In the sequel we will be interested in the permutation action of
(3.1)
$$H=\psi _{A_1}(A_1)\times \psi _{A_2}(A_2)\mathrm{S}_{n_1}\times \mathrm{S}_{n_2}\mathrm{S}_n$$
on $`k^n`$. Here $`A_1`$ and $`A_2`$ are abelian groups of order $`n_1`$ and $`n_2`$ respectively and $`n=n_1+n_2`$. For future reference, we decompose this action as a direct sum of character spaces. We shall write elements of $`k^n=k^{n_1+n_2}`$ as $`(R^{},R^{\prime \prime })`$, where $`R^{}k^{n_1}`$ and $`R^{\prime \prime }k^{n_2}`$. Let $`V_0=\{(\underset{n_1\text{ times}}{\underset{}{a,\mathrm{},a}},\underset{n_2\text{ times}}{\underset{}{b,\mathrm{},b}})a,bk\}`$.
###### Lemma 3.1.
$$k^n=V_0\left(_{\chi A_1^{}}\mathrm{Span}_k(R_\chi ,0)\right)\left(_{\eta A_2^{}}\mathrm{Span}_k(0,R_\eta )\right)$$
is a decomposition of $`k^n`$ as a direct sum of character spaces for the $`H`$-action defined above. Here $`V_0`$ is a 2-dimensional subspace with trivial associated character; the remaining $`n2`$ summands are 1-dimensional subspaces with distinct nontrivial characters.
###### Proof.
The proof of this lemma amounts to verifying that the summands of the above decomposition are, indeed, character spaces and finding their characters. We leave the details of the reader. ∎
### Abelian subgroups of $`\mathrm{PGL}_n`$
Let $`A`$ be an abelian subgroup of order $`n`$ and $`V=k[A]`$. The group $`A`$ acts on $`V`$ by the regular representation $`aP_a\mathrm{GL}(V)`$, where
$$P_a(_{bA}c_bb)=_{bA}c_bab$$
for any $`aA`$ and $`c_bk`$. The dual group $`A^{}`$ acts on $`V`$ by the representation $`\chi D_\chi \mathrm{GL}(V)`$, where
$$D_\chi (_{aA}c_aa)=_{aA}c_a\chi (a)a$$
for any $`\chi A^{}`$ and $`c_ak`$. Note that in the basis $`\{a|aA\}`$ of $`V`$, each $`P_a`$ is represented by a permutation matrix and each $`D_\chi `$ is represented by a diagonal matrix; this explains our choice of the letters $`P`$ and $`D`$. It is easy to see that
(3.2)
$$D_\chi P_a=\chi (a)P_aD_\chi ;$$
hence, we have constructed an embedding
(3.3)
$$\varphi _A:A\times A^{}\mathrm{PGL}(V)=\mathrm{PGL}_n$$
given by $`(a,\chi )\overline{P_a}\overline{D_\chi }`$, where $`\overline{P_a}`$ and $`\overline{D_\chi }`$ are the elements of $`\mathrm{PGL}(V)`$, represented, respectively, by $`P_a`$ and $`D_\chi \mathrm{GL}(V)`$.
For future reference we record two simple lemmas.
###### Lemma 3.2.
For each $`aA`$ and $`\chi A^{}`$, $`V_{a,\chi }=\mathrm{Span}_k(P_aD_\chi )`$ is a 1-dimensional $`H`$-invariant subspace of $`\mathrm{M}_n`$, with associated character $`(b,\eta )\chi ^1(b)\eta (a)`$. Moreover, the $`n^2`$ matrices $`P_aD_\chi `$ form a $`k`$-basis of $`\mathrm{M}_n`$.
###### Proof.
The first assertion is immediate from (3.2). Since the $`n^2`$ characters associated to the spaces $`V_{a,\chi }`$ are distinct, the second assertion now follows from linear independence of characters. ∎
###### Lemma 3.3.
Let $`A`$ be an abelian group of order $`n`$ and $`(a,\chi )`$ be an element of order $`c`$ in $`A\times A^{}`$.
(a) $`(P_aD_\chi )^c=ϵI_n`$, where $`ϵ=\chi (a)^{\frac{1}{2}c(c1)}=\pm 1`$ and $`I_n`$ is the $`n\times n`$-identity matrix.
(b) The characteristic polynomial of $`P_aD_\chi `$ is $`r(t)=(t^cϵ)^{\frac{n}{c}}`$.
(c) Assume $`\mathrm{char}(k)(\frac{n}{c})!`$. Then $`\sigma ^{(i)}(P_aD_\chi )0`$ for any $`i`$ divisible by $`c`$.
###### Proof.
(a) The identity $`(P_aD_\chi )^c=ϵI_n`$, where $`ϵ=\chi (a)^{\frac{1}{2}c(c1)}`$, is immediate from (3.2). To see that $`ϵ=1`$ or $`1`$, note that $`ϵ^2=(\chi (a)^c)^{c1}=1^{c1}=1`$.
(b) Let $`C`$ be the cyclic subgroup of $`A\times A^{}`$ generated by $`(a,\chi )`$, so that $`c=|C|`$. For each $`\alpha (A\times A^{})/C`$, let $`V_\alpha `$ be the vector subspace of $`\mathrm{M}_n`$ spanned by $`(b,\eta )\alpha `$. Each $`V_\alpha `$ is a $`c`$-dimensional subspace of $`\mathrm{M}_n`$, which is stable under right multiplication by $`P_aD_\chi `$. Since the matrices $`P_aD_\chi `$ form a basis of $`\mathrm{M}_n`$ as $`(a,\chi )`$ ranges over $`A\times A^{}`$ (see Lemma 3.2), we can write
(3.4)
$$\mathrm{M}_n=_{\alpha (A\times A^{})/C}V_\alpha .$$
By part (a), $`(P_aD_\chi )^c=ϵI_n`$. It is now easy to see that the characteristic polynomial for the action of $`P_aD_\chi `$ on each $`V_\alpha `$ is $`p(t)=t^cϵ`$. Consequenly, the charactersistic polynomial for the left multiplication action of $`P_aD_\chi `$ on $`\mathrm{M}_n`$ is $`q(t)=p(t)^{n^2/c}`$ (one factor of $`p(t)`$ for each subspace $`V_\alpha `$ in (3.4)), and the characteristic polynomial of the $`n\times n`$-matrix $`P_aD_\chi `$ (or, equivalently, of its action on $`n\times 1`$-column vectors) is
$$r(t)=q(t)^{1/n}=p(t)^{n/c}=(t^cϵ)^{n/c},$$
as claimed.
(c) The binomial formula tells us that under our assumption on $`\mathrm{char}(k)`$, every monomial of the form $`t^{ni}`$ with $`i`$ divisible by $`c`$ (and $`in`$), appears in $`r(t)`$ with a nonzero coefficient. In other words, for these values of $`i`$, $`\sigma ^{(i)}(P_aD_\chi )0`$, as claimed. ∎
## 4. Proof of Theorem 1.5
We may (and will, throughout this section) assume without loss of generality that $`k`$ is an algebraically closed field. Otherwise we can simply replace $`D`$ by $`\overline{D}=D_k\overline{k}`$, where $`\overline{k}`$ is the algebraic closure of $`k`$: if the system (1.5) has no nontrivial solutions in $`\overline{D}`$, it cannot have one in $`D`$.
Our goal is to deduce Theorem 1.5 as a special case of Proposition 2.4. We shall now proceed to introduce the finite abelian group $`H`$ and the $`\mathrm{PGL}_n`$-varieties $`X`$ and $`Y`$ and to show that they satisfy the conditions of Proposition 2.4. We will then apply Proposition 2.4 with these $`H`$, $`X`$, and $`Y`$, to conclude that the system (1.5) has no nontrivial solutions in $`D_n`$ or $`\mathrm{UD}(n)`$.
### The group $`H`$
We define $`H`$ to be the finite abelian subgroup of $`\mathrm{PGL}_n`$ given by
(4.1) $`H=A\times A^{}\stackrel{\varphi _A}{}\mathrm{PGL}_n,`$ where $`A=/p_1\times \mathrm{}\times /p_s.`$
Here, as in Section 1, $`n=p_1\mathrm{}p_s`$, where $`p_1,\mathrm{},p_s`$ are not necessarily distinct primes; the inclusion $`\varphi _A`$ is as in (3.3). Note that the assumption $`\mathrm{char}(k)n!`$ of Theorem 1.5 implies that $`|H|=n^2`$ is prime to $`\mathrm{char}(k)`$.
### The variety $`X`$
We shall now write the algebras that come up in the statement of Theorem 1.5, namely $`D=\mathrm{UD}(n)`$ and $`D=D_n`$, in the form $`\mathrm{RMaps}_{\mathrm{PGL}_n}(X,\mathrm{M}_n)`$ for specific $`\mathrm{PGL}_n`$-varieties $`X`$. Note that we do not assume $`\mathrm{char}(k)=0`$.
###### Lemma 4.1.
(Procesi) $`\mathrm{UD}(n)=\mathrm{RMaps}_{\mathrm{PGL}_n}(X,\mathrm{M}_n)`$, where $`X=(\mathrm{M}_n)^2`$ and $`\mathrm{PGL}_n`$ acts on $`X`$ by simultaneous conjugation.
###### Proof.
See \[Sa, Theorem 14.16\], cf. also \[P2, Theorem 2.1\] or \[RY<sub>2</sub>, Example 3.1\]. ∎
Let $`G`$ be an algebraic group, $`S`$ be a closed subgroup of $`G`$, and $`Y`$ be an affine $`S`$-variety. The groups $`S`$ and $`G`$ act on $`G\times Y`$ via respectively, $`s(g,y)=(gs^1,sy)`$ and $`g^{}(g,y)=(g^{}g,y)`$; moreover, the two actions commute. Thus the quotient $`(G\times Y)//S=\mathrm{Spec}(k[G\times Y]^S)`$ is a $`G`$-variety; we will denote it by $`G_SY`$. We will restrict our attention to the case where $`S`$ is a finite group of order prime to $`\mathrm{char}(k)`$. In this case a theorem of Hilbert and Noether (see, e.g., \[Sm, Theorem 1.1\]) tells us that $`k[G\times Y]^S`$ is a finitely generated $`k`$-algebra, i.e., $`G_SY`$ is again an affine variety (of finite type).
###### Lemma 4.2.
There exists a faithful $`2s`$-dimensional linear representation $`V`$ of $`H`$ such that $`D_n\mathrm{RMaps}_{\mathrm{PGL}_n}(X,\mathrm{M}_n)`$, where $`X=\mathrm{PGL}_n_HV`$.
###### Proof.
Choose a set of generators $`a_1,\mathrm{},a_s`$ for $`A`$ and a “dual” set of generators $`\chi _1,\mathrm{},\chi _s`$ for $`A^{}`$ so that
$$\chi _i(a_j)=\{\begin{array}{cc}1\hfill & \text{if }ij\hfill \\ \zeta _{p_i}\hfill & \text{if }i=j,\hfill \end{array}$$
where $`\zeta _{p_i}`$ is the same primitive $`p_i`$th root of unity used in defining $`(z_i,w_i)_{p_i}`$; see (1.3) and (1.4). Consider the faithful action of $`H=A\times A^{}`$ on $`V=k^{2s}`$ given by
$`(a,\chi ):(\alpha _1,\mathrm{},\alpha _s,\beta _1,\mathrm{},\beta _s)`$
$`(\chi ^1(a_1)\alpha _1,\mathrm{},\chi ^1(a_s)\alpha _s,\chi _1(a)\beta _1,\mathrm{},\chi _s(a)\beta _s).`$
Set $`X=\mathrm{PGL}_n_HV`$ and $`R=\mathrm{RMaps}_{\mathrm{PGL}_n}(X,\mathrm{M}_n)`$. Note that
(4.2)
$$k(X)^{\mathrm{PGL}_n}=k(\mathrm{PGL}_n\times V)^{\mathrm{PGL}_n\times H}=k(V)^H=k(\alpha _1^{p_1},\beta _l^{p_1},\mathrm{},\alpha _s^{p_s},\beta _s^{p_s}).$$
Define elements $`\pi _i`$ and $`\eta _i`$ of $`R`$ by
(4.3) $`\pi _i:[g,(\alpha _1,\mathrm{},\alpha _s,\beta _1,\mathrm{},\beta _s)]`$ $`\alpha _igP_{a_i}g^1`$
$`\eta _i:[g,(\alpha _1,\mathrm{},\alpha _s,\beta _1,\mathrm{},\beta _s)]`$ $`\beta _igD_{\chi _i}g^1.`$
These elements are well-defined because $`\pi _i(g,v)=\pi _i(gh^1,hv)`$ and $`\eta _i(g,v)=\eta _i(gh^1,hv)`$ for every $`hH`$ and $`i=1,\mathrm{},s`$; see (3.2). Note that since $`P_{a_i}`$ and $`D_{\chi _i}`$ generate $`\mathrm{M}_n(k)`$ as a $`k`$-algebra, as $`i`$ ranges from $`1`$ to $`n`$ (cf. Lemma 3.2), there exists a dense Zariski dense open subset $`X_0X`$ such that
(4.4) $`\pi _i(x)`$ and $`\eta _i(x)`$ generate $`\mathrm{M}_n(k)`$ for every $`xX_0`$.
In particular, if $`f`$ is a central element of $`R`$ then $`f(x)`$ is a scalar matrix for every $`xX_0`$. Consequently, the center $`Z(R)`$ consists of rational maps $`X\mathrm{M}_n`$ whose image lies in the subspace of scalar matrices. In other words,
(4.5)
$$Z(R)=k(X)^{\mathrm{PGL}_n}$$
where, as before, we identify $`fk(X)^{\mathrm{PGL}_n}`$ with the $`\mathrm{PGL}_n`$-equivariant rational map $`X\mathrm{M}_n(k)`$ given by $`xf(x)I_n`$.
We are now ready to construct an isomorphism between $`D_n`$ and $`R`$. First we identify $`D_n`$ with the skew-polynomial ring
$$D_n=Z(R)\{x_1,y_1,\mathrm{},x_s,y_s\},$$
where $`x_i^{p_i}=\alpha _i^{p_i}`$, $`y_i^{p_i}=\beta _i^{p_i}`$, $`y_ix_i=\zeta _{p_i}x_iy_i`$ and all other pair of variables commute. (Recall that $`Z(R)`$ is the purely transcendental extension of $`k`$ generated by $`\alpha _1^{p_1},\beta _1^{p_1},\mathrm{},\alpha _s^{p_s},\beta _s^{p_s}`$; see (4.2) and (4.5).) Let $`\varphi :D_nR`$ be the $`Z(R)`$-algebra homomorphism given by $`\varphi (x_i)=\pi _i`$ and $`\varphi (y_i)=\eta _i`$. This homomorphism is well-defined because $`\pi _i`$ and $`\eta _i`$ satisfy the same relations as $`x_i`$ and $`y_i`$; see (4.3) and (3.2).
We claim $`\varphi `$ is an isomorphism. Indeed, $`\varphi `$ is injective since $`D_n`$ is a simple algebra. Moreover, since $`dim_k(\mathrm{M}_n)=n^2`$, it is easy to see that $`dim_{Z(R)}Rn^2`$ (see, e.g. , \[Re<sub>2</sub>, Lemma 7.4(a)\] for a characteristic-free proof). This shows that $`\varphi `$ is an isomorphism and thus completes the proof of Lemma 4.2. ∎
### The variety $`Y`$
We now define the $`\mathrm{PGL}_n`$-variety $`Y`$ by
(4.8) $`Y=\{(y_1:\mathrm{}:y_m)((\mathrm{M}_n)^m)|\begin{array}{c}\sigma ^{(i)}(y_1)=\mathrm{}=\sigma ^{(i)}(y_m)\hfill \\ \sigma ^{(j)}(y_1\mathrm{}y_m)=0\hfill \end{array}\},`$
as in Proposition 2.4. Recall that our goal is to use Proposition 2.4 to show that the system (1.5) has no nontrivial solutions.
###### Lemma 4.3.
Under the assumptions of Theorem 1.5 (i.e., $`\mathrm{char}(k)n!`$, $`\mathrm{sqf}(n)m`$ and $`\mathrm{sqf}(n)i`$), $`H`$ acts on $`Y`$ without fixed points.
###### Proof.
The $`H`$-fixed points in $`((\mathrm{M}_n)^m)`$ are of the form $`y=(y_1:\mathrm{}:y_m)`$, where each $`y_i`$ is either 0 or an element of $`\mathrm{M}_n`$ which spans a 1-dimensional character space for $`H`$. Moreover, the associated characters of all non-zero $`y_i`$ have to be the same. Thus, in view of Lemma 3.2, there exists an element $`(a,\chi )A\times A^{}`$ such that $`y_i=t_iP_aD_\chi `$ for some $`t_1,\mathrm{},t_mk`$. Note that at least one $`t_i`$ has to be non-zero, since otherwise $`y=(0:\mathrm{}:0)`$ is not a well-defined point of $`((\mathrm{M}_n)^m)`$.
Now suppose $`y`$ is an $`H`$-fixed point of $`Y`$. Substituting $`y_i=t_iP_aD_\chi `$ into the defining equations for $`Y`$, we obtain
(4.9)
$$\{\begin{array}{c}t_1^i\sigma ^{(i)}(P_aD_\chi )=\mathrm{}=t_m^i\sigma ^{(i)}(P_aD_\chi ),\hfill \\ t_1\mathrm{}t_m\sigma ^{(j)}((P_aD_\chi )^m)=0.\hfill \end{array}$$
Let $`c`$ be the order of $`(a,\chi )`$ in $`A\times A^{}`$. Then $`c\mathrm{exp}(A)`$, $`\mathrm{exp}(A)=\mathrm{sqf}(n)`$, $`\mathrm{sqf}(n)m`$, $`\mathrm{sqf}(n)i`$, and thus, $`cm`$ and $`ci`$. By Lemma 3.3(a), $`(P_aD_\chi )^m=\pm I_n`$, and hence, $`\sigma ^{(j)}((P_aD_\chi )^m)0`$. By Lemma 3.3(c), $`\sigma ^{(i)}(P_aD_\chi )0`$. Therefore, we can rewrite (4.9) as
$$\{\begin{array}{c}t_1^i=\mathrm{}=t_m^i,\hfill \\ t_1\mathrm{}t_m=0.\hfill \end{array}$$
This system has no solutions other than $`t_1=\mathrm{}=t_m=0`$, a contradiction. We conclude that $`Y`$ has no $`H`$-fixed points, as claimed. ∎
### Conclusion of the proof
In order to complete the proof, it remains to show that $`X`$ has a smooth $`H`$-fixed point; the desired conclusion will then follow by applying Proposition 2.4 to the abelian group $`H`$ and $`\mathrm{PGL}_n`$-varieties $`X`$ and $`Y`$ we introduced above.
If $`D=\mathrm{UD}(n)`$ then $`X=(\mathrm{M}_n)^2`$ (see Lemma 4.1), and the origin is a smooth $`H`$-fixed point of $`X`$.
If $`D=D_n`$ then $`X=\mathrm{PGL}_n_HV=(\mathrm{PGL}_n\times V)//H`$; see Lemma 4.2. Since $`\mathrm{PGL}_n\times V`$ is a smooth variely, and $`H`$ acts freely on it, $`X`$ is also smooth. Moreover, the point of $`X`$ represented by $`(1,0)\mathrm{PGL}_n\times V`$, is clearly fixed by $`H`$. Thus $`X`$ has a smooth $`H`$-fixed point, as claimed.
This completes the proof of Theorem 1.5. ∎
### Refinements
A slight modification of the above argument proves the following more general variant of Theorem 1.5.
###### Theorem 4.4.
Let $`P(z_1,\mathrm{},z_v)k\{z_1,\mathrm{},z_v\}`$ be a homogeneous (non-commutative) polynomial of degree $`d`$ in $`v`$ variables. The system of equations
(4.10)
$$\{\begin{array}{c}\sigma ^{(i)}(x_1^u)=\mathrm{}=\sigma ^{(i)}(x_v^u)\hfill \\ \sigma ^{(j)}(P(x_1,\mathrm{},x_v))=0\hfill \end{array}$$
has no nontrivial solutions in $`D_n`$ or $`\mathrm{UD}(n)`$, provided that
(i) $`iu`$ and $`jd`$ are divisible by $`\mathrm{sqf}(n)`$.
(ii) $`P(\zeta _1,\mathrm{},\zeta _v)0`$ for any (not necessarily primitive) $`ij`$-th roots of unity $`\zeta _1,\mathrm{},\zeta _v`$.
Note that if we set $`u=1`$, $`d=v=m`$ and $`P(z_1,\mathrm{},z_v)=z_1\mathrm{}z_v`$, then we recover Theorem 1.5 from Theorem 4.4.
###### Remark 4.5.
Suppose $`K=k(a_1,b_2,\mathrm{},a_l,b_l)`$ and
$$D=(a_1,b_1)_{r_1}_K\mathrm{}_K(a_l,b_l)_{r_l}$$
be a tensor product of generic symbol algebras of degree $`n=r_1\mathrm{}r_l`$. Denote the least common multiple of $`r_1,\mathrm{},r_l`$ by $`e`$. (Equivalently, $`e`$ is the exponent of $`D`$.) Then the system (1.5) has no solutions in $`D`$ as long as $`i`$ and $`m`$ are divisible by $`e`$. The proof is the same as above, except that instead of choosing $`H`$ and $`A`$ as in (4.1), we take $`H=\varphi _A(A\times A^{})`$ with $`A=(/r_1)\times \mathrm{}\times (/r_l)`$. Similarly, the system (4.10) has no solutions in $`D`$, provided that $`iu`$ and $`jd`$ is divisible by $`e`$, and condition (ii) of Theorem 4.4 holds.
###### Remark 4.6.
Theorem 1.5 remains true if $`D`$ is replaced by the general field extension $`L_n/K_n`$. The reason is that there is a natural embedding $`\alpha :L_n\mathrm{UD}(n)`$ such that
$$\alpha :\sigma _{L_n/K_n}^{(i)}(y)\sigma _{\mathrm{UD}(n)/Z(n)}^{(i)}(\alpha (y))$$
for every $`yL_n`$ and every $`i=1,\mathrm{},n`$. Indeed, recall that $`\mathrm{UD}(n)`$ is generated by two generic $`n\times n`$-matrices, $`X=(s_{ij})`$ and $`Y=(t_{ij})`$: we can define $`\alpha (x)=X`$ and $`\alpha (a_i)=\sigma ^{(i)}(X)`$, see, e.g., \[P1, Lemma II.1.4\]. If system (1.5) had a nontrivial solution in $`L_n`$, it would then have a nontrivial solution in $`\mathrm{UD}(n)`$, contradicting Theorem 1.5.
###### Remark 4.7.
Suppose $`\mathrm{char}(k)=0`$, $`n=p^r`$ and $`D^{}`$ as a prime-to-$`p`$ extension of $`D_n`$ or $`\mathrm{UD}(n)`$. Then Theorem 1.5, Corollary 1.6 and Theorem 4.4 remain valid if $`D`$ is replaced by $`D^{}`$. Indeed, let $`X`$ be as in Lemma 4.1 (if $`D=\mathrm{UD}(n)`$) and Lemma 4.2 (if $`D=D_n`$). Then we can write $`D^{}`$ as $`\mathrm{RMaps}_{\mathrm{PGL}_n}(X^{},\mathrm{M}_n)`$, where $`X^{}X`$ is a $`\mathrm{PGL}_n`$-invariant rational cover, of degree prime to $`p`$. We may assume that $`X^{}`$ is smooth and projective. (This follows from canonical resolution of singularities; see \[RY<sub>2</sub>, Proposition 2.2\].) Since $`H`$ is a $`p`$-group, the Going Up Theorem says that $`X^{}`$ has an $`H`$-fixed point; see \[RY<sub>1</sub>, Proposition A.4\]. The desired conclusion now follows from Proposition 2.4.
## 5. Proof of Theorem 1.7
We may assume without loss of generality that $`k`$ is an algebraically closed field; otherwise we may simply replace $`K_n`$ and $`L_n`$ by $`K_n_k\overline{k}`$ and $`L_n_k\overline{k}`$ respectively, where $`\overline{k}`$ is the algebraic closure of $`k`$.
Let $`f(x)=x^n+a_1x^{n1}+\mathrm{}+a_n`$ and $`L_n=K_n[x]/(f(x))`$, as in (1.6). The normal closure of $`L_n`$ over $`K_n`$ is the field $`L^{}=K_n(x_1,\mathrm{},x_n)=k(x_1,\mathrm{},x_n)`$, where $`x_1,\mathrm{},x_n`$ are the roots of $`f`$; they are algebraically independent over $`k`$. We will identify $`L_n`$ with $`K_n(x_1)`$ by identifying $`xL_n`$ with $`x_1k(x_1,\mathrm{},x_n)`$ and $`a_i`$ with $`(1)^is_i(x_1,\mathrm{},x_n)`$, where $`s_i`$ is the $`i`$th elementary symmetric polynomial.
We shall deduce Theorem 1.7 as a particular case of Proposition 2.2, with $`m=1`$, $`K=K_n`$, $`L=L_n`$, $`L^{}`$ as above, and $`G=\mathrm{Gal}(L^{}/L_n)=\mathrm{S}_n`$. We will now define the remaining objects that appear in the statement of Proposition 2.2, namely the abelian subgroup $`H`$ of $`G=\mathrm{S}_n`$ and the $`G`$-varieties $`X`$ and $`Y`$.
We set $`H=H_1\times H_2`$, with $`H_1=\psi _{A_1}(A_1)\mathrm{S}_{n_1}`$, $`H_2=\psi _{A_2}(A_2)\mathrm{S}_{n_2}`$, as in (3.1); here for $`i=1,2`$, $`A_i`$ is an abelian subgroup of order $`n_i`$ and exponent $`\mathrm{sqf}(n_i)`$. More precisely, if $`n_1=p_1\mathrm{}p_s`$ and $`n_2=q_1\mathrm{}q_t`$ are written as products of (not necessarily distinct) primes then
$`H_1A_1=(/p_1)\times \mathrm{}\times (/p_s)`$
(5.1) and
$`H_2A_2=(/q_1)\times \mathrm{}\times (/q_t).`$
We define $`X=𝔸^n`$, with the natural permutation action of $`G=\mathrm{S}_n`$. If we denote the coordinates on $`𝔸^n`$ by $`x_1,\mathrm{},x_n`$ then $`k(X)=k(x_1,\mathrm{},x_n)=L^{}`$ as fields with $`\mathrm{S}_n`$-action. The origin is a smooth point of $`X`$ fixed by $`\mathrm{S}_n`$ and, hence, by $`H`$.
The $`S_n`$-variety $`Y`$ is defined as the subvariety of $`(𝔸^n)=^{n1}`$ given by
(5.2)
$$\{\begin{array}{cc}\hfill x_1^{m_1}+\mathrm{}+x_n^{m_1}& =0\hfill \\ \hfill x_1^{m_2}+\mathrm{}+x_n^{m_2}& =0.\hfill \end{array}$$
In order to apply Proposition 2.2, it is now sufficient to prove the following:
###### Lemma 5.1.
Under the assumptions (i) and (ii) of Theorem 1.7, $`Y`$ has no $`H`$-fixed points.
###### Proof.
By Lemma 3.1, the fixed points $`y`$ for the $`H`$-action on $`^{n1}=^{n_1+n_21}`$ are of one of the following three types:
Type I: $`y=R_{a,b}=(\underset{n_1\text{ times}}{\underset{}{a:\mathrm{}:a}}:\underset{n_2\text{ times}}{\underset{}{b:\mathrm{}:b}})`$, for some $`a,bk`$, not both 0.
Type II: $`y=(R_\chi ,0)=(\chi (\alpha _1):\mathrm{}:\chi (\alpha _{n_1}):0:\mathrm{}:0)`$, where $`H_1=\{\alpha _1,\mathrm{},\alpha _{n_1}\}`$ and $`\chi `$ is a character of $`H_1`$.
Type III: $`y=(0,R_\eta )=(0:\mathrm{}:0:\eta (\beta _1):\mathrm{}:\eta (\beta _{n_2}))`$, where $`H_2=\{\beta _1,\mathrm{},\beta _{n_2}\}`$ and $`\eta `$ is a character of $`H_2`$.
Consider a point of type I. Substituting the coordinates of $`R_{a,b}`$ into (5.2), we see that $`R_{a,b}`$ lies in $`Y`$ if and only if $`(a,b)`$ is a nontrivial solution of the homogeneous system
(5.3)
$$\{\begin{array}{c}n_1a^{m_1}+n_2b^{m_1}=0\hfill \\ n_1a^{m_2}+n_2b^{m_2}=0.\hfill \end{array}$$
An elementary computation shows that under assumption (i) of Theorem 1.7 this system has no nontrivial solutions. Hence we conclude that no point of type I can lie on $`Y`$.
We now turn to points of types II and III. Since $`H_1`$ has exponent $`\mathrm{sqf}(n_1)`$, we see that $`\chi (\alpha _i)^{\mathrm{sqf}(n_1)}=1`$ for every $`\alpha _iH_1`$. It follows from the assumptions of Theorem 1.7 that $`n_10`$ in $`k`$ and either $`m_1`$ or $`m_2`$ is divisible by $`\mathrm{sqf}(n_1)`$; consequently, $`(R_\chi ,0)`$ does not lie on $`Y`$. Similarly, $`(0,R_\eta )`$ does not lie on $`Y`$. Hence, no point of type II or III lies on $`Y`$. This completes the proof of the lemma and thus of Theorem 1.7. ∎
###### Remark 5.2.
Theorem 1.7 fails if the field extension $`L_n/K_n`$ is replaced by the generic division algebra $`\mathrm{UD}(n)`$. Suppose, for simplicity, that $`k`$ is an algebraically closed field of characteristic zero. Then, by a theorem of Wedderburn, $`\mathrm{UD}(3)`$ is cyclic; thus it has an elements $`x`$ and $`y`$ such that $`x=\zeta _3yxy^1`$, where $`\zeta _3`$ is a primitive cube root of 1. It is now easy to see that $`\mathrm{tr}(x)=\mathrm{tr}(x^2)=0`$. On the other hand, Theorem 1.7 with $`n_1=m_1=1`$ and $`n_2=m_2=2`$, says that no such element can exist in $`L_3`$.
Another example of this kind can be constructed for $`n=6`$. The algebra $`D=\mathrm{UD}(6)`$ is known to be cyclic; hence, it has a non-zero element $`z`$ such that $`\mathrm{tr}(z^i)=0`$ for $`i=1,\mathrm{},5`$. On the other hand, Theorem 1.7 says that the systems $`\mathrm{tr}(x)=\mathrm{tr}(x^5)=0`$ or $`\mathrm{tr}(x^2)=\mathrm{tr}(x^4)=0`$ have no solutions in $`L_6^{}`$.
###### Remark 5.3.
Let $`n_1=p_1\mathrm{}p_s`$ and $`n_2=q_1\mathrm{}q_t`$, where $`p_1,\mathrm{},p_s,q_1,\mathrm{},q_t`$ are (not necessarily distinct) primes. Suppose $`z_1,\mathrm{},z_s`$ and $`w_1,\mathrm{},w_t`$ are independent variables over $`k`$. Set $`E_1=k(z_1,\mathrm{},z_s,w_1^{q_1},\mathrm{},w_t^{q_t})`$, $`E_2=k(z_1^{p_1},\mathrm{},z_s^{p_s},w_1,\mathrm{},w_t)`$, and $`F=k(z_1^{p_1},\mathrm{},z_s^{p_s},w_1^{q_1},\mathrm{},w_t^{q_t})`$. Then we can replace $`L_n/K_n`$ by the $`n`$-dimensional etale $`F`$-algebra $`E=E_1E_2`$ (cf. \[Re<sub>1</sub>, Section 4\]) in the statement of Theorem 1.7. In other words,
under assumptions (i) and (ii) of Theorem 1.7 the system of equations $`\mathrm{tr}(x^{m_1})=\mathrm{tr}(x^{m_2})=0`$ has no nontrivial solutions in $`E`$.
The role played by $`E`$ in this setting is analogous to the role played by $`D_n`$ in the setting of Theorem 1.5. In particular, one can show that $`E=\mathrm{RMaps}_{\mathrm{S}_n}(X,𝔸^n)`$, where $`X=\mathrm{S}_n_HV`$, $`V`$ is a faithful $`(s+t)`$-dimensional linear representation of $`H=H_1\times H_2`$, and the algebra structure on $`\mathrm{RMaps}_{\mathrm{S}_n}(X,𝔸^n)`$ is induced from the algebra structure on $`𝔸^n=\underset{n\text{ times}}{\underset{}{k\mathrm{}k}}`$ (compare with Lemma 4.2). Since $`X`$ has a smooth $`H`$-fixed point (namely, the point represented by $`(id,0)\mathrm{S}_n\times V`$), the rest of our argument goes through unchanged.
## 6. Systems of the form $`\sigma ^{(m_1)}(x)=\sigma ^{(m_2)}(x)=0`$
We do not know whether or not the system $`\mathrm{tr}(x^{m_1})=\mathrm{tr}(x^{m_2})=0`$ may be replaced by the system
(6.1)
$$\sigma ^{(m_1)}(x)=\sigma ^{(m_2)}(x)=0.$$
in the statement of Theorem 1.7. (Such a result would be of interest, since it would mean that the general polynomial of degree $`n`$ cannot be transformed, by a Tschirnhaus substitution, into a polynomial $`t^n+b_1t^{n1}+\mathrm{}+b_n`$, with $`b_{m_1}=b_{m_2}=0`$.) Every step of our proof of Theorem 1.7 goes through in this case, except that the system (5.3) is replaced by the system
(6.2)
$$\{\begin{array}{c}s_{m_1}(a,\mathrm{},a,b,\mathrm{},b)=0\hfill \\ s_{m_2}(a,\mathrm{},a,b,\mathrm{},b)=0,\hfill \end{array}$$
where $`(a,\mathrm{},a,b,\mathrm{},b)`$ stands for $`(\underset{n_1\text{ times}}{\underset{}{a,\mathrm{},a}},\underset{n_2\text{ times}}{\underset{}{b,\mathrm{},b}})`$ and $`s_i`$ denotes the $`i`$th elementary symmetric polynomial. Thus:
###### Proposition 6.1.
Let $`n_1`$ and $`n_2`$ be positive integers prime to $`\mathrm{char}(k)`$, and $`L_n/K_n`$ be the general field extension of degree $`n=n_1+n_2`$. Then the system (6.1) has no nontrivial solutions $`xL_n^{}`$, provided that each $`\mathrm{sqf}(n_i)`$ ($`i=1,2`$) divides $`m_1`$ or $`m_2`$ and the system (6.2) has no nontrivial solutions $`(a,b)k^2`$.
Of course, this result is less satisfying than Theorem 1.7 because we do not know for what values of $`n_1`$, $`m_1`$, $`n_2`$ and $`m_2`$ the system (6.2) has no nontrivial solutions. (The analogous question for the system (5.3) is quite easy: the answer is given by condition (i) of Theorem 1.7.) Nevertheless, for low values of $`n`$, Proposition 6.1 gives us a rather complete picture. We shall give two such examples below.
Before preceeding with the examples, we record a simple observation.
###### Remark 6.2.
Let $`E/F`$ be a field extension of degree $`n`$. Multiplying (1.1) by $`det((\lambda x)^1)`$, we easily obtain the identity $`\sigma ^{(ni)}(x^1)=\sigma ^{(i)}(x)/\sigma ^{(n)}(x)`$. In particular, if $`xE`$ satisfies (6.1) then $`\sigma ^{(nm_1)}(x^1)=\sigma ^{(nm_2)}(x^1)=0`$. ∎
###### Example 6.3.
Let $`L_5/K_5`$ be the general field extension of degree $`5`$ and let $`1m_1<m_25`$. Then the system (6.1) has a nontrivial solution $`xL_5^{}`$ if and only if $`(m_1,m_2)=(1,3)`$ or $`(2,4)`$.
###### Proof.
By the theorem of Hermite cited in Example 1.1, the system (6.1) has a solution $`0xL_5`$ for $`(m_1,m_2)=(1,3)`$. Then $`x^1`$ is a solution to (6.1) with $`(m_1,m_2)=(2,4)`$; see Remark 6.2.
It remains to show that there are no solutions for any other values of $`m_1`$ and $`m_2`$. Indeed, we may assume without loss of generality that $`m_25`$, since $`\sigma ^{(5)}(x)=det(x)0`$ for any $`xL_5^{}`$. The remaining possibilities for $`(m_1,m_2)`$ are: $`(1,2)`$, $`(1,4)`$, $`(2,3)`$, and $`(3,4)`$. In view of Remark 6.2, we only need to consider $`(1,2)`$, $`(1,4)`$ and $`(2,3)`$.
$`(m_1,m_2)=(1,2)`$. By Newton’s formulas the system $`\sigma ^{(1)}(x)=\sigma ^{(2)}(x)=0`$ is equivalent to $`\mathrm{tr}(x)=\mathrm{tr}(x^2)=0`$. The latter system has no solutions by Theorem 1.7 with $`n_1=1`$ and $`n_2=4`$. (Alternatively, use Proposition 6.1 with $`n_1=1`$, $`n_2=4`$ or appeal to \[Re<sub>1</sub>, Theorem 1.3(b)\], with $`p=2`$ and $`m=2`$.)
$`(m_1,m_2)=(1,4)`$. Apply Proposition 6.1 with $`n_1=1`$ and $`n_2=4`$. In this case (6.2) reduces to
$$\{\begin{array}{c}s_1(a,b,b,b,b)=a+4b=0\hfill \\ s_4(a,b,b,b,b)=b^4+4ab^3=0.\hfill \end{array}$$
It is easy to see that this system has no nontrivial solutions. (Alternatively, use \[Re<sub>1</sub>, Theorem 6.1b\].)
$`(m_1,m_2)=(2,3)`$. Apply Proposition 6.1 with $`n_1=2`$ and $`n_2=3`$. In this case (6.2) becomes
$$\{\begin{array}{c}s_2(a,a,b,b,b)=a^2+6ab+3b^2=0\hfill \\ s_3(a,a,b,b,b)=3a^2b+6ab^2+b^3=0.\hfill \end{array}$$
This system has no nontrivial solutions. ∎
###### Example 6.4.
Let $`L_6/K_6`$ be the general field extension of degree $`6`$ and let $`1m_1<m_26`$. Then the system (6.1) has a nontrivial solution $`xL_5^{}`$ if and only if $`(m_1,m_2)=(1,3)`$ or $`(3,5)`$.
###### Proof.
The existence of solutions for $`(m_1,m_2)=(1,3)`$ and $`(3,5)`$ follows from Example 1.1 and Remark 6.2.
We may assume $`m_25`$ because $`\sigma ^{(6)}(x)=det(x)0`$ for any $`xL_6^{}`$. It is now enough to show that there are no solutions for $`(m_1,m_2)=(1,2)`$, $`(1,4)`$, $`(1,5)`$, $`(2,3)`$, and $`(2,4)`$; the remaining cases follow from these by Remark 6.2.
$`(m_1,m_2)=(1,2)`$. In this case (6.2) is equivalent to $`\mathrm{tr}(x)=\mathrm{tr}(x^2)=0`$. The latter system has no solutions by Theorem 1.7 with $`n_1=2`$ and $`n_2=4`$. (Alternatively, use Proposition 6.1 with $`n_1=1`$, $`n_2=4`$ or appeal to \[Re<sub>1</sub>, Theorem 1.3(c)\], with $`p=2`$, $`m=2`$ and $`l=1`$.)
$`(m_1,m_2)=(1,4)`$. Apply Proposition 6.1 with $`n_1=2`$, $`n_2=4`$. In this case (6.2) reduces to $`2a+4b=6a^2b^2+8ab^3+b^4=0`$. This system has no nontrivial solutions.
$`(m_1,m_2)=(1,5)`$. Apply Proposition 6.1 with $`n_1=1`$, $`n_2=5`$. In this case (6.2) reduces to $`a+5b=5ab^4+b^5=0`$. There are no nontrivial solutions. (Alternatively, use \[Re<sub>1</sub>, Theorem 1.3(b)\] with $`p=5`$.)
$`(m_1,m_2)=(2,3)`$. Apply Proposition 6.1 with $`n_1=2`$, $`n_2=4`$. In this case (6.2) becomes $`a^2+8ab+6b^2=4a^2b+12ab^2+4b^3=0`$. There are no nontrivial solutions.
$`(m_1,m_2)=(2,4)`$. Use Proposition 6.1 with $`n_1=2`$, $`n_2=4`$. In this case (6.2) becomes $`a^2+8ab+6b^2=6a^2b^2+8ab^3+b^4=0`$. Once again, there are no nontrivial solutions. ∎
## 7. A further generalization
In this section we will show that the assumption that the $`G`$-variety $`Y`$ in Proposition 2.2 has no fixed points can sometimes be weakened. We will present a general result extending Proposition 2.2 and illustrate it with an example. One can generalize Proposition 2.4 in a similar manner; we leave the details to an interested reader.
In this section we assume that $`k`$ is algebraically closed.
###### Proposition 7.1.
Assume
(i) $`L/K`$ is a separable field extension of degree $`n`$, $`L^{}`$ is the normal closure of $`L`$ over $`K`$, $`G=\mathrm{Gal}(L^{},K)`$, and $`H`$ is an abelian subgroup of $`G`$,
(ii) $`YZ`$ are subvarieties of $`(𝔸^n)^m`$ given, respectively, by systems of $`G`$-invariant polynomial equations $`P_1=\mathrm{}=P_s=0`$ and $`Q_1=\mathrm{}=Q_r=0`$,
(iii) there exists a complete $`H`$-variety $`W`$ without $`H`$-fixed points and a regular $`H`$-equivariant map $`h:YZW`$, and
(iv) there exists a $`G`$-variety $`X`$ such that such that $`k(X)=L^{}`$ as fields with $`G`$-action, and $`X`$ has a smooth $`H`$-fixed point.
Then any solution $`(a_1,\mathrm{},a_m)L^m`$ of the system
(7.1)
$$P_1(x_1,\mathrm{},x_m)=\mathrm{}=P_s(x_1,\mathrm{},x_m)=0$$
also satisfies the system
(7.2)
$$Q_1(x_1,\mathrm{},x_m)=\mathrm{}=Q_r(x_1,\mathrm{},x_m)=0.$$
Note that since $`ZY`$, the ideal $`(Q_1,\mathrm{},Q_r)k[(𝔸^n)^m]`$ contains some power of the ideal $`(P_1,\mathrm{},P_s)`$. Hence, any solution of (7.2) in $`L^n`$ is a solution of (7.1). Proposition 7.1 asserts that under assumptions (i)–(iv), the opposite is also true.
###### Proof.
Given a solution $`(a_1,\mathrm{},a_m)`$ of (7.1), we construct a rational map $`f:XY(𝔸^n)^m`$, as in the proof of Proposition 2.2. If $`(a_1,\mathrm{},a_m)`$ does not satisfy (7.2), then $`f(X)Z`$ and hence, the composition $`X\stackrel{𝑓}{}Y\stackrel{}{}W`$ is a well-defined $`H`$-equivariant rational map. As $`X`$ has a smooth $`H`$-fixed point, Theorem 2.1 says that $`W`$ also has one, a contradiction. ∎
###### Remark 7.2.
To see that Proposition 2.2 is a special case of Proposition 7.1, assume that the polynomials $`P_1,\mathrm{},P_s`$ are homogeneous, so that $`Y`$ is a cone in $`(𝔸^n)^m`$, and $`Z`$ is the origin in $`(𝔸^n)^m`$. Note that the origin of $`(𝔸^n)^m`$ can be cut out by $`G`$-invariant homogeneous polynomials (this is true for any finite group representation), thus we can choose $`Q_1,\mathrm{},Q_rk[(𝔸^n)^m]^G`$ to be generators of the ideal of the origin in $`k[(𝔸^n)^m]`$.
Let $`W((𝔸^n)^m)`$ be the projectivisation of the cone $`Y`$, and $`h:YZW`$ the natural projection. If $`W`$ has no $`H`$-fixed points, and $`X`$ has a smooth $`H`$-fixed point then Proposition 7.1 implies that the system (7.1) has no solutions, except for $`x_1=\mathrm{}=x_m=0`$. This is precisely the statement of Proposition 2.2.
###### Remark 7.3.
Proposition 7.1 can be applied in the following situation. Suppose that $`Z`$ is the singular set of $`Y`$. Let $`\stackrel{~}{Y}`$ be the closure of $`Y(𝔸^n)^m=𝔸^{nm}`$ in $`^{nm}𝔸^{nm}`$; note that the $`G`$-action on $`(𝔸^n)^m`$ extends to a regular $`G`$-action on $`^{nm}`$, and $`\stackrel{~}{Y}`$ is $`G`$-invariant. Let $`\pi :W\stackrel{~}{Y}`$ be the canonical resolution of singularities. Such a resolution is known to exist if $`\mathrm{char}(k)=0`$; see the discussion and the references in \[RY<sub>1</sub>, Section 3\]. Note that $`\pi `$ is an isomorphism over $`YZ`$ and thus we can take $`h=\pi ^1:YZW`$. If $`W`$ has no $`H`$-fixed points then Proposition 7.1 applies.
###### Example 7.4.
Suppose $`n`$ is prime and $`n\mathrm{char}(k)`$. Then for any $`ck`$ the equation
(7.3)
$$\underset{i=1}{\overset{n1}{}}\sigma ^{(i)}(x)^n\sigma ^{(n)}(x)^{n1i}+c\sigma ^{(n)}(x)^{2n2}=0,$$
has no nontrivial solutions in the general field extension $`L_n/K_n`$; see (1.6). Here $`\sigma ^{(i)}`$ stands for $`\sigma _{L_n/K_n}^{(i)}`$.
###### Proof.
We may assume without loss of generality that $`k`$ is algebraically closed, and thus, contains the roots of unity.
First consider the case $`c0`$. We apply Proposition 7.1 in the following setting: $`K=K_n`$, $`L=L_n`$, $`G=\mathrm{S}_n`$, $`X=𝔸^n`$ with the natural $`\mathrm{S}_n`$-action, $`H`$ = the cyclic subgroup of $`\mathrm{S}_n`$ generated by the $`n`$-cycle $`h=(\mathrm{1\hspace{0.17em}2}\mathrm{}n)`$, $`s=m=1`$, and $`P_1=s_1^ns_n^{n2}+s_2^ns_n^{n3}+\mathrm{}+s_{n1}^n+cs_n^{2n2}`$, where $`s_i`$ denotes the $`i`$th elementary symmetric polynomial in the coordinates $`x_1,\mathrm{},x_n`$ in $`𝔸^n`$. (To construct $`P_1`$, we replaced $`\sigma ^{(i)}(x)`$ by $`(1)^is_i(x_1,\mathrm{},x_n)`$ in the left hand side of 7.3.) Note that $`P_1`$ is not homogeneous in $`x_1,\mathrm{},x_n`$ as $`c0`$.
We take $`Z`$ to be the origin in $`𝔸^n`$. Similarly to Remark 7.3, let $`\stackrel{~}{Y}`$ the closure of $`YA^n`$ in $`^n`$; then the $`H`$-action on $`\stackrel{~}{Y}Z`$ is free. Let $`\stackrel{~}{^n}^n`$ be the blowup of $`Z`$; we identify its exceptional divisor $`S`$ with $`^{n1}`$. Let $`Y^{}`$ be the strict transform of $`\stackrel{~}{Y}`$; then $`Y^{}\stackrel{~}{Y}`$ is a blowup centered at $`Z`$, and $`SY^{}`$ is the hypersurface in $`^{n1}`$ given by the homogeneous equation $`\overline{P}_1=0`$ where $`\overline{P}_1=s_1^ns_n^{n2}+s_2^ns_n^{n3}+\mathrm{}+s_{n1}^n`$ is the initial form of $`P_1`$.
The intersection $`SY^{}`$ contains $`H`$-fixed points $`q_\zeta =(1:\zeta :\zeta ^2:\mathrm{}:\zeta ^{n1})`$ for each $`n`$-th root of unity $`\zeta 1`$. Let $`WY^{}`$ be the blowup of these $`n1`$ points. We claim that $`W`$ has no $`H`$-fixed points.
To see this, consider the hypersurfaces $`S_i\stackrel{~}{^n}`$ for $`i=1,\mathrm{},n1`$ which are the closures in $`\stackrel{~}{^n}`$ of the hypersurfaces in $`𝔸^nZ`$ given by the equations $`s_i=0`$. For each $`i`$, the intersection $`S_iS`$ is the hypersurface in $`S=^{n1}`$ given by the homogeneous equation $`s_i=0`$; in particular, each $`S_i`$ passes through $`q_\zeta `$. Consider the $`(n1)\times (n1)`$ Jacobian determinants $`D_l(q_\zeta )=det(s_i/x_j)(q_\zeta )`$, where $`i=1,\mathrm{},n1`$ and $`j=1,\mathrm{},\widehat{l},\mathrm{},n`$. By Newton’s formulas $`D_l(q_\zeta )=det(p_i/x_j)(q_\zeta )`$, where $`p_i=x_1^i+\mathrm{}+x_n^i`$. The latter determinant is a Vandermonde determinant, which does not vanish at $`q_\zeta `$. This shows that the hypersurfaces $`S_iS`$ are smooth and intersect transversely (in $`S=^{n1}`$) at each $`q_\zeta `$; hence $`S_1,\mathrm{},S_{n1}`$ and $`S`$ are smooth and intersect transversely (in $`\stackrel{~}{^n}`$) at each $`q_\zeta `$.
Thus the tangent spaces $`T_{q_\zeta }(S_1),\mathrm{},T_{q_\zeta }(S_{n1})`$, together with $`T_{q_\zeta }(S)`$, form a system of coordinate hyperplanes in $`T_{q_\zeta }(\stackrel{~}{^n})`$. Since each $`S_i`$ is $`H`$-invariant, the linear $`H`$-action on $`T_{q_\zeta }(\stackrel{~}{^n})`$ is diagonalized in this coordinate system. The group $`H`$ acts by different characters on each of the coordinate directions; in fact, $`h`$ acts by multiplication by $`\zeta ^i`$ on $`T_{q_\zeta }(\stackrel{~}{^n})/T_{q_\zeta }(S_i)`$, and trivially on $`T_{q_\zeta }(\stackrel{~}{^n})/T_{q_\zeta }(S)`$. Identifying the exceptional divisor $`E_{q_\zeta }`$ of the blowup of $`\stackrel{~}{^n}`$ centered at $`q_\zeta `$, with $`(T_{q_\zeta }(\stackrel{~}{^n}))`$, we see that the $`H`$-fixed points on $`E_{q_\zeta }`$ are the points of $`(T_{q_\zeta }(\stackrel{~}{^n}))`$ that correspond to the directions of the coordinate axes in $`T_{q_\zeta }(\stackrel{~}{^n})`$. The exceptional divisor of $`W`$ over $`q_\zeta `$ is the projectivisation of the tangent cone to $`Y^{}`$ at $`q_\zeta `$, and the latter does not contain the coordinate axes. We conclude that $`W`$ does not have $`H`$-fixed points, as claimed.
Thus, we may apply Proposition 7.1; it shows that the equation (7.3) has no nontrivial solutions, similarly to Remark 7.2.
In case $`c=0`$, we need to make the following changes. Now $`Y`$ is an affine cone; we take $`Z`$ to be the union of $`(n1)!`$ lines that correspond to the points $`(\zeta _1:\mathrm{}:\zeta _n)^{n1}`$ where $`\zeta _1,\mathrm{},\zeta _n`$ are different $`n`$th roots of unity; this includes the lines that correspond to the points $`q_\zeta `$. Now let $`Y^{}`$ be the blowup of $`\stackrel{~}{Y}`$ at the origin as before, and $`W`$ be the blowup of $`Y^{}`$ at the lines that make up the strict transform of $`Z`$ in $`Y^{}`$. (Alternatively, we may take the route similar to Remark 7.2 and set $`W`$ to be the blowup of $`(Y)`$ at the points $`q_\zeta `$.) Then $`W`$ does not have $`H`$-fixed points, and Proposition 7.1 shows that any $`xL_n`$ satisfying (7.3) also satisfies the system (7.2), which in our case is
(7.4)
$$\sigma ^{(1)}(x)=\mathrm{}=\sigma ^{(n1)}(x)=0.$$
One can now show directly that $`L_n`$ does not have a non-zero element $`x`$ satisfying (7.4); otherwise $`L_n/K_n`$ would have to be a cyclic extension, a contradiction. Alternatively, one can show that the system (7.4) has no nontrivial solutions by applying Proposition 7.1 one more time, as follows:
* take the new $`H`$ to be any cyclic subgroup of $`G=\mathrm{S}_n`$ of order different from $`n`$ and $`1`$;
* the new $`Y`$ to be the old $`Z`$, i.e., $`P_i=s_i(x_1,\mathrm{},x_n)`$ for $`i=1,\mathrm{},n1`$.
* the new $`Z`$ to be the origin in $`𝔸^n`$, i.e., $`Q_j=s_j(x_1,\mathrm{},x_n)`$ for $`j=1,\mathrm{},n`$.
* the new $`W`$ to be the normalization of $`Z`$, i.e., the disjoint union of $`(n1)!`$ lines.
Applying Proposition 7.1 we see that the system (7.4) has no nontrivial solutions and, hence, neither does equation (7.3). ∎
## 8. Equations in octonion algebras
### Preliminaries
Let $`F`$ be a field of characteristic $`2`$. Recall that for any $`0a,b,cF`$, the octonion (or Cayley—Dickson) algebra $`𝕆_F(a,b,c)`$ is defined as follows. The quaternion algebra
$$(a,b)_2=F\{i,j\}/(i^2=a,j^2=b,ji=ij)$$
is equipped with an involution $`x\overline{x}`$ given by
(8.1)
$$\overline{x_0+x_1i+x_2j+x_3ij}=x_0x_1ix_2jx_3ij$$
for any $`x_0,\mathrm{},x_3F`$. Now $`𝕆_F(a,b,c)\stackrel{\mathrm{def}}{=}(a,b)_2(a,b)_2l`$ is an 8-dimensional $`F`$-algebra with (non-associative) multiplication given by $`(x+yl)(z+wl)=(xz+c\overline{w}y)+(wx+y\overline{z})l`$. The involution (8.1) extends from $`(a,b)_2`$ to $`𝕆_F(a,b,c)`$ via $`\overline{x+yl}=\overline{x}yl`$. The algebra $`𝕆_F(a,b,c)`$ is also equipped with $`F`$-valued trace and norm functions given by $`\mathrm{tr}(x)=x+\overline{x}`$ and $`n(x)=x\overline{x}=\overline{x}x`$ such that $`x^2\mathrm{tr}(x)x+n(x)=0`$ for any $`x𝕆_F(a,b,c)`$; we can think of $`\mathrm{tr}(x)`$ as $`\sigma ^{(1)}(x)`$ and $`n(x)`$ as $`\sigma ^{(2)}(x)`$. Note that $`\mathrm{tr}(x)`$ is intrinsically defined in $`𝕆_K(a,b,c)`$, i.e., $`\mathrm{tr}(x)=\mathrm{tr}(\sigma (x))`$, where $`\sigma `$ is a $`K`$-algebra automorphism in $`𝕆_K(a,b,c)`$; the same is true of $`n(x)`$. For a more detailed description of octonion algebras we refer the reader to \[Sc\].
Two octonion algebras will be of particular interest to us: the split algebra $`𝕆_F(1,1,1)`$ over $`F`$ and the generic algebra $`𝕆_{gen}=𝕆_K(a,b,c)`$, where $`K=k(a,b,c)`$ and $`a,b,c`$ are algebraically independent over $`k`$.
By a theorem of Zorn \[Sc, III.3.17\], any $`8`$-dimensional $`F`$-algebra $`A`$ such that $`A_FF^{}𝕆_F(1,1,1)`$ for some field extension $`F^{}/F`$, is necessarily isomorphic to $`𝕆_F(a,b,c)`$ for some $`a,b,cF^{}`$. This means that octonion algebras are “forms” of the split octonion algebra $`𝕆_k(1,1,1)`$ in the same way as central simple algebras are “forms” of the matrix algebra $`\mathrm{M}_n(k)`$.
### $`G_2`$-equivariant maps
From now on we shall assume the base field $`k`$ to be algebraically closed and of characteristic $`2`$.
Recall that the automorphism group of the split octonion algebra $`𝕆=𝕆_k(1,1,1)`$ is the exceptional group $`G_2`$. Octonion algebras are related to $`G_2`$-varieties in the same way as central simple algebras are related to $`\mathrm{PGL}_n`$-varieties. In particular, if $`k`$ is of characteristic 0 then any octonion algebra whose center is a finitely generated field extension of $`k`$ can be written in the form $`\mathrm{RMaps}_{G_2}(X,𝕆)`$, where $`𝕆`$ is viewed as an 8-dimensional vector space with the natural $`G_2`$-action and $`X`$ is a generically free $`G_2`$-variety, uniquely determined up to birational isomorphism.
From now on, let $`H(/2)^3`$ be the subgroup of $`G_2`$ generated by $`\tau _1`$, $`\tau _2`$ and $`\tau _3`$, where
(8.2)
$$\begin{array}{ccc}\tau _1(i)=i,\hfill & \tau _1(j)=j,\hfill & \tau _1(l)=l;\hfill \\ \tau _2(i)=i,\hfill & \tau _2(j)=j,\hfill & \tau _2(l)=l;\hfill \\ \tau _3(i)=i,\hfill & \tau _3(j)=j,\hfill & \tau _3(l)=l.\hfill \end{array}$$
###### Lemma 8.1.
The generic octonion algebra $`𝕆_{gen}`$ is isomorphic to $`\mathrm{RMaps}_{G_2}(X,V)`$, where $`X=G_2_HV`$ and $`V=\mathrm{Span}\{i,j,k\}`$ is the 3-dimensional faithful representation of $`H`$ given by (8.2).
###### Proof.
The proof is similar to the proof of Lemma 4.2, so we will only outline it below.
Let $`\alpha ,\beta ,\gamma `$ be the coordinates of $`V`$ relative to the basis $`\{i,j,l\}`$, let $`R=\mathrm{RMaps}_{G_2}(X,𝕆)`$ and let $`\pi _1,\pi _2,\pi _3:X𝕆`$ be the elements of $`R`$ given by
(8.3)
$$\begin{array}{c}\pi _1:[g,(\alpha ,\beta ,\gamma )]\alpha g(i)\hfill \\ \pi _2:[g,(\alpha ,\beta ,\gamma )]\beta g(j)\hfill \\ \pi _3:[g,(\alpha ,\beta ,\gamma )]\gamma g(l).\hfill \end{array}$$
It is easy to see that these maps are well-defined, i.e. $`\pi _a(g,v)=\pi _a(gh^1,hv)`$. Let
$$K=k(X)^{G_2}=k(V)^H=k(\alpha ^2,\beta ^2,\gamma ^2).$$
We now identify $`𝕆_{gen}`$ with $`𝕆_K(\alpha ^2,\beta ^2,\gamma ^2)`$, and define $`\varphi :𝕆_{gen}R`$ by $`\varphi (i)=\pi _1`$, $`\varphi (j)=\pi _2`$ and $`\varphi (l)=\pi _3`$. Then $`\varphi `$ is well-defined; see (8.3). Since $`𝕆`$ is a (non-associative) division algebra, $`\varphi `$ is injective. To see that $`\varphi `$ is an isomorphism, we only need to show that $`dim_K(R)8`$; this follows from \[Re<sub>2</sub>, Lemma 7.4(a)\]. ∎
### $`G_2`$-invariant polynomials
Consider the diagonal $`G_2`$-action on the $`8m`$-dimensional $`k`$-vector space $`W=𝕆^m`$. Let $`Pk[W]^{G_2}`$ be a $`G`$-invariant polynomial and let $`A=𝕆_F(a,b,c)`$ be an octonion algebra. Identifying $`A`$ with an $`F`$-subalgebra of $`A_FF^{}𝕆_F^{}(1,1,1)`$, where $`F^{}=F(\sqrt{a},\sqrt{b},\sqrt{c})`$, we can define $`P(a_1,\mathrm{},a_m)`$ for any $`a_1,\mathrm{},a_mA`$. Arguing as in Lemma 2.3, we see that $`P(a_1,\mathrm{},a_m)`$ is well-defined and lies in $`F`$ for any $`a_1,\mathrm{},a_mA`$. (This also follows from a theorem of Schwarz \[Sw, (3.23)\], which asserts that $`k[W]^{G_2}`$ is generated by elements of the form $`\mathrm{tr}(M)`$, where $`M`$ is a monomial in $`u_1,\mathrm{},u_m𝕆`$.)
###### Proposition 8.2.
Let $`H(/2)^3`$ be the subgroup of $`G_2`$ defined in (8.2). Suppose the subvariety $`Y`$ of $`(𝕆^m)`$, cut out by homogeneous $`G_2`$-invariant polynomials $`P_1=\mathrm{}=P_r=0`$, does not have an $`H`$-fixed point. Then the system
(8.4)
$$P_1(x_1,\mathrm{},x_m)=\mathrm{}=P_r(x_1,\mathrm{},x_m)=0$$
has no non-trivial solutions in any octonion algebra of the form $`\mathrm{RMaps}_{G_2}(X,𝕆)`$, where $`X`$ is a $`G_2`$-variety with a smooth $`H`$-fixed point. In particular, the system (8.4) has no nontrivial solutions in the generic octonion algebra $`O_{gen}`$.
###### Proof.
We argue as in the proof of Proposition 2.4. Assume, to the contrary, that $`(a_1,\mathrm{},a_m)`$ is a nontrivial solution of (8.4). Each $`a_i`$ is a $`G_2`$-equivariant rational map $`X𝕆^m`$; together they define a $`G_2`$-equivariant rational map $`f:XY(𝕆^m)`$. Applying the Going Down Theorem 2.1, we obtain a contradiction.
This proves the first assertion of the proposition. The second assertion follows from Lemma 8.1. Indeed, the variety $`X=G_2_HV`$ defined there has a smooth fixed point, namely $`(1,0)`$. ∎
### A system of equations
We are now ready to state and prove the main result of this section.
###### Theorem 8.3.
Let $`Q(x_1,\mathrm{},x_m)`$ be (a non-commutative and non-associative) homogeneous polynomial of even degree in $`x_1,\mathrm{},x_m`$ such that $`Q(ϵ_1,\mathrm{},ϵ_m)0`$ for any $`(2s)`$-th roots of unity $`ϵ_1,\mathrm{},ϵ_m`$, and let $`m`$ and $`s`$ be positive integers. Then the system
(8.5)
$$\{\begin{array}{c}\mathrm{tr}(x_1^{2s})=\mathrm{}=\mathrm{tr}(x_m^{2s})\\ \mathrm{tr}(Q(x_1,\mathrm{},x_m))=0.\end{array}$$
has no non-zero solutions in any octonion algebra of the form $`\mathrm{RMaps}_{G_2}(X,𝕆)`$, where $`X`$ is a generically free $`G_2`$-variety with a smooth $`H`$-fixed point. In particular, the system (8.5) has no nontrivial solutions in the generic octonion algebra $`𝕆_{gen}`$.
Here $`H=<\tau _1,\tau _2,\tau _3>(/2)^3`$ is the subgroup of $`G_2`$ defined in (8.2).
###### Proof.
According to Proposition 8.2, it is enough to check that the variety
$$Y=\{(U_1:\mathrm{}:U_m)(𝕆^m)|\mathrm{tr}(U_1^{2s})=\mathrm{}=\mathrm{tr}(U_m^{2s}),\mathrm{tr}(Q(U_1,\mathrm{},U_m))=0\}$$
(where $`U_1,\mathrm{},U_m𝕆`$ are taken up to multiplication by an element of $`k`$) has no $`H`$-fixed points.
A point $`(U_1:\mathrm{}:U_m)(𝕆^m)`$ is $`H`$-fixed iff all $`U_r`$ lie in the same character space for the $`H`$-action on $`𝕆`$. In other words, there exists a $`\zeta \{1,i,j,l,ij,il,jl,ijl\}`$ such that every $`U_r`$ is of the form $`U_r=u_r\zeta `$ for some $`u_rk`$. Note that at least one $`u_r`$ is non-zero; otherwise the point $`(U_1:\mathrm{}:U_m)`$ is not well-defined in $`(𝕆^m)`$. The condition that such a fixed point lies in $`Y`$ translates into the system
$$\{\begin{array}{cc}& u_1^{2s}=\mathrm{}=u_m^{2s}\hfill \\ & Q(u_1,\mathrm{},u_m)=0\hfill \end{array}$$
of homogeneous equations in $`u_1,\mathrm{},u_m`$. If $`u_1=0`$ then the remaining $`u_r`$ are also equal to $`0`$, a contradiction. If $`u_10`$ then $`ϵ_r=u_r/u_1`$ is a $`(2s)`$-th root of unity for each $`r=1,\mathrm{},m`$, and $`Q(ϵ_1,\mathrm{},ϵ_m)=0`$, contradicting our assumption on $`Q`$. This shows that $`Y`$ has no $`H`$-fixed points. ∎
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# Supersymmetry Phenomenology
## 1 Motivation
### 1.1 Problems in the Standard Model
The Standard Model of particle physics, albeit extremely successful phenomenologically, has been regarded only as a low-energy effective theory of the yet-more-fundamental theory. One can list many reasons why we think this way, but a few are named below.
First of all, the quantum number assignments of the fermions under the standard $`SU(3)_C\times SU(2)_L\times U(1)_Y`$ gauge group (Table 1) appear utterly bizarre. Probably the hypercharges are the weirdest of all. These assignments, however, are crucial to guarantee the cancellation of anomalies which could jeopardize the gauge invariance at the quantum level, rendering the theory inconsistent. Another related puzzle is why the hypercharges are quantized in the unit of $`1/6`$. In principle, the hypercharges can be any numbers, even irrational. However, the quantized hypercharges are responsible for neutrality of bulk matter $`Q(e)+2Q(u)+Q(d)=Q(u)+2Q(d)=0`$ at a precision of $`10^{21}`$.
The gauge group itself poses a question as well. Why are there seemingly unrelated three independent gauge groups, which somehow conspire together to have anomaly-free particle content in a non-trivial way? Why is “the strong interaction” strong and “the weak interaction” weaker?
The essential ingredient in the Standard Model which appears the ugliest to most people is the electroweak symmetry breaking. In the list of bosons in the Standard Model Table 2, the gauge multiplets are necessary consequences of the gauge theories, and they appear natural. They of course all carry spin 1. However, there is only one spinless multiplet in the Standard Model: the Higgs doublet
$$\left(\begin{array}{c}H^+\\ H^0\end{array}\right)$$
(1)
which condenses in the vacuum due to the Mexican-hat potential. It is introduced just for the purpose of breaking the electroweak symmetry $`SU(2)_L\times U(1)_YU(1)_{\mathrm{QED}}`$. The potential has to be arranged in a way to break the symmetry without any microscopic explanations.
Why is there a seemingly unnecessary three-fold repetition of “generations”? Even the second generation led the Nobel Laureate I.I. Rabi to ask “who ordered muon?” Now we face even more puzzling question of having three generations. And why do the fermions have a mass spectrum which stretches over almost six orders of magnitude between the electron and the top quark? This question becomes even more serious once we consider the recent evidence for neutrino oscillations which suggest the mass of the third-generation neutrino $`\nu _\tau ^{}`$ of about $`0.05`$ eV. This makes the mass spectrum stretching over thirteen orders of magnitude. We have no concrete understanding of the mass spectrum nor the mixing patterns.
### 1.2 Drive to go to Shorter Distances
All the puzzles raised in the previous section (and more) cry out for a more fundamental theory underlying the Standard Model. What history suggests is that the fundamental theory lies always at shorter distances than the distance scale of the problem. For instance, the equation of state of the ideal gas was found to be a simple consequence of the statistical mechanics of free molecules. The van der Waals equation, which describes the deviation from the ideal one, was the consequence of the finite size of molecules and their interactions. Mendeleev’s periodic table of chemical elements was understood in terms of the bound electronic states, Pauli exclusion principle and spin. The existence of varieties of nuclide was due to the composite nature of nuclei made of protons and neutrons. The list would go on and on. Indeed, seeking answers at more and more fundamental level is the heart of the physical science, namely the reductionist approach.
The distance scale of the Standard Model is given by the size of the Higgs boson condensate $`v=250`$ GeV. In natural units, it gives the distance scale of $`d=\mathrm{}c/v=0.8\times 10^{16}`$ cm. We therefore would like to study physics at distance scales shorter than this eventually, and try to answer puzzles whose partial list was given in the previous section.
Then the idea must be that we imagine the Standard Model to be valid down to a distance scale shorter than $`d`$, and then new physics will appear which will take over the Standard Model. But applying the Standard Model to a distance scale shorter than $`d`$ poses a serious theoretical problem. In order to make this point clear, we first describe a related problem in the classical electromagnetism, and then discuss the case of the Standard Model later along the same line.
### 1.3 Positron Analogue
In the classical electromagnetism, the only dynamical degrees of freedom are electrons, electric fields, and magnetic fields. When an electron is present in the vacuum, there is a Coulomb electric field around it, which has the energy of
$$\mathrm{\Delta }E_{\mathrm{Coulomb}}=\frac{1}{4\pi \epsilon _0}\frac{e^2}{r_e}.$$
(2)
Here, $`r_e`$ is the “size” of the electron introduced to cutoff the divergent Coulomb self-energy. Since this Coulomb self-energy is there for every electron, it has to be considered to be a part of the electron rest energy. Therefore, the mass of the electron receives an additional contribution due to the Coulomb self-energy:
$$(m_ec^2)_{\mathrm{𝑜𝑏𝑠}}=(m_ec^2)_{\mathrm{𝑏𝑎𝑟𝑒}}+\mathrm{\Delta }E_{\mathrm{Coulomb}}.$$
(3)
Experimentally, we know that the “size” of the electron is small, re
<1017
<subscript𝑟𝑒superscript1017r_{e}\mathrel{\mathchoice{\lower 3.0pt\vbox{\halign{$\m@th\displaystyle\hfill#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\m@th\textstyle\hfill#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\m@th\scriptstyle\hfill#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\m@th\scriptscriptstyle\hfill#\hfil$\cr<\crcr\sim\crcr}}}}10^{-17} cm. This implies that the self-energy $`\mathrm{\Delta }E`$ is greater than 10 GeV or so, and hence the “bare” electron mass must be negative to obtain the observed mass of the electron, with a fine cancellation like
$$0.511=9999.489+10000.000\mathrm{MeV}.$$
(4)
Even setting a conceptual problem with a negative mass electron aside, such a fine-cancellation between the “bare” mass of the electron and the Coulomb self-energy appears ridiculous. In order for such a cancellation to be absent, we conclude that the classical electromagnetism cannot be applied to distance scales shorter than $`e^2/(4\pi \epsilon _0m_ec^2)=2.8\times 10^{13}`$ cm. This is a long distance in the present-day particle physics’ standard.
The resolution to the problem came from the discovery of the anti-particle of the electron, the positron, or in other words by doubling the degrees of freedom in the theory. The Coulomb self-energy discussed above can be depicted by a diagram where the electron emits the Coulomb field (a virtual photon) which is absorbed later by the electron (the electron “feels” its own Coulomb field). But now that the positron exists (thanks to Anderson back in 1932), and we also know that the world is quantum mechanical, one should think about the fluctuation of the “vacuum” where the vacuum produces a pair of an electron and a positron out of nothing together with a photon, within the time allowed by the energy-time uncertainty principle $`\mathrm{\Delta }t\mathrm{}/\mathrm{\Delta }E\mathrm{}/(2m_ec^2)`$. This is a new phenomenon which didn’t exist in the classical electrodynamics, and modifies physics below the distance scale $`dc\mathrm{\Delta }t\mathrm{}c/(2m_ec^2)=200\times 10^{13}`$ cm. Therefore, the classical electrodynamics actually did have a finite applicability only down to this distance scale, much earlier than $`2.8\times 10^{13}`$ cm as exhibited by the problem of the fine cancellation above. Given this vacuum fluctuation process, one should also consider a process where the electron sitting in the vacuum by chance annihilates with the positron and the photon in the vacuum fluctuation, and the electron which used to be a part of the fluctuation remains instead as a real electron. V. Weisskopf calculated this contribution to the electron self-energy for the first time, and found that it is negative and cancels the leading piece in the Coulomb self-energy exactly:
$$\mathrm{\Delta }E_{\mathrm{pair}}=\frac{1}{4\pi \epsilon _0}\frac{e^2}{r_e}.$$
(5)
After the linearly divergent piece $`1/r_e`$ is canceled, the leading contribution in the $`r_e0`$ limit is given by
$$\mathrm{\Delta }E=\mathrm{\Delta }E_{\mathrm{Coulomb}}+\mathrm{\Delta }E_{\mathrm{pair}}=\frac{3\alpha }{4\pi }m_ec^2\mathrm{log}\frac{\mathrm{}}{m_ecr_e}.$$
(6)
There are two important things to be said about this formula. First, the correction $`\mathrm{\Delta }E`$ is proportional to the electron mass and hence the total mass is proportional to the “bare” mass of the electron,
$$(m_ec^2)_{\mathrm{𝑜𝑏𝑠}}=(m_ec^2)_{\mathrm{𝑏𝑎𝑟𝑒}}\left[1+\frac{3\alpha }{4\pi }\mathrm{log}\frac{\mathrm{}}{m_ecr_e}\right].$$
(7)
Therefore, we are talking about the “percentage” of the correction, rather than a huge additive constant. Second, the correction depends only logarithmically on the “size” of the electron. As a result, the correction is only a 9% increase in the mass even for an electron as small as the Planck distance $`r_e=1/M_{Pl}=1.6\times 10^{33}`$ cm.
The fact that the correction is proportional to the “bare” mass is a consequence of a new symmetry present in the theory with the antiparticle (the positron): the chiral symmetry. In the limit of the exact chiral symmetry, the electron is massless and the symmetry protects the electron from acquiring a mass from self-energy corrections. The finite mass of the electron breaks the chiral symmetry explicitly, and because the self-energy correction should vanish in the chiral symmetric limit (zero mass electron), the correction is proportional to the electron mass. Therefore, the doubling of the degrees of freedom and the cancellation of the power divergences lead to a sensible theory of electron applicable to very short distance scales.
### 1.4 Supersymmetry
In the Standard Model, the Higgs potential is given by
$$V=\mu ^2|H|^2+\lambda |H|^4,$$
(8)
where $`v^2=H^2=\mu ^2/2\lambda =(176\mathrm{GeV})^2`$. Because perturbative unitarity requires that λ
<1
<𝜆1\lambda\mathrel{\mathchoice{\lower 3.0pt\vbox{\halign{$\m@th\displaystyle\hfill#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\m@th\textstyle\hfill#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\m@th\scriptstyle\hfill#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\m@th\scriptscriptstyle\hfill#\hfil$\cr<\crcr\sim\crcr}}}}1, $`\mu ^2`$ is of the order of $`(100\mathrm{GeV})^2`$. However, the mass squared parameter $`\mu ^2`$ of the Higgs doublet receives a quadratically divergent contribution from its self-energy corrections. For instance, the process where the Higgs doublets splits into a pair of top quarks and come back to the Higgs boson gives the self-energy correction
$$\mathrm{\Delta }\mu _{\mathrm{top}}^2=6\frac{h_t^2}{4\pi ^2}\frac{1}{r_H^2},$$
(9)
where $`r_H`$ is the “size” of the Higgs boson, and $`h_t1`$ is the top quark Yukawa coupling. Based on the same argument in the previous section, this makes the Standard Model not applicable below the distance scale of $`10^{17}`$ cm.
The motivation for supersymmetry is to make the Standard Model applicable to much shorter distances so that we can hope that answers to many of the puzzles in the Standard Model can be given by physics at shorter distance scales. In order to do so, supersymmetry repeats what history did with the positron: doubling the degrees of freedom with an explicitly broken new symmetry. Then the top quark would have a superpartner, stop,<sup>1</sup><sup>1</sup>1This is a terrible name, which was originally meant to be “scalar top.” If supersymmetry will be discovered by the next generation collider experiments, we should seriously look for better names for the superparticles. whose loop diagram gives another contribution to the Higgs boson self energy
$$\mathrm{\Delta }\mu _{\mathrm{stop}}^2=+6\frac{h_t^2}{4\pi ^2}\frac{1}{r_H^2}.$$
(10)
The leading pieces in $`1/r_H`$ cancel between the top and stop contributions, and one obtains the correction to be
$$\mathrm{\Delta }\mu _{\mathrm{top}}^2+\mathrm{\Delta }\mu _{\mathrm{top}}^2=6\frac{h_t^2}{4\pi ^2}(m_{\stackrel{~}{t}}^2m_t^2)\mathrm{log}\frac{1}{r_H^2m_{\stackrel{~}{t}}^2}.$$
(11)
One important difference from the positron case, however, is that the mass of the stop, $`m_{\stackrel{~}{t}}`$, is unknown. In order for the $`\mathrm{\Delta }\mu ^2`$ to be of the same order of magnitude as the tree-level value $`\mu ^2=2\lambda v^2`$, we need $`m_{\stackrel{~}{t}}^2`$ to be not too far above the electroweak scale. Similar arguments apply to masses of other superpartners that couple directly to the Higgs doublet. This is the so-called naturalness constraint on the superparticle masses (for more quantitative discussions, see papers).
### 1.5 Other Directions
Of course, supersymmetry is not the only solution discussed in the literature to avoid miraculously fine cancellations in the Higgs boson mass-squared term. Technicolor (see a review) is a beautiful idea which replaces the Higgs doublet by a composite techni-quark condensate. Then $`r_H1`$ TeV is a truly physical size of the Higgs doublet and there is no need for fine cancellations. Despite the beauty of the idea, this direction has had problems with generating fermion masses, especially the top quark mass, in a way consistent with the constraints from the flavor-changing neutral currents. The difficulties in the model building, however, do not necessarily mean that the idea itself is wrong; indeed still efforts are being devoted to construct realistic models.
Another recent idea is to lower the Planck scale down to the TeV scale by employing large extra spatial dimensions. This is a new direction which has just started, and there is an intensive activity to find constraints on the idea as well as on model building. Since the field is still new, there is no “standard” framework one can discuss at this point, but this is no surprise given the fact that supersymmetry is still evolving even after almost two decades of intense research.
One important remark about all these ideas is that they inevitably predict interesting signals at TeV-scale collider experiments. While we only discuss supersymmetry in this lecture, it is likely that nature has a surprise ready for us; maybe none of the ideas discussed so far is right. Still we know that there is something out there to be uncovered at TeV scale energies.
## 2 Supersymmetric Lagrangian
We do not go into full-fledged formalism of supersymmetric Lagrangians in this lecture but rather confine ourselves to a practical introduction of how to write down Lagrangians with explicitly broken supersymmetry which still fulfill the motivation for supersymmetry discussed in the previous section. One can find useful discussions as well as an extensive list of references in a nice review by Steve Martin.
### 2.1 Supermultiplets
Supersymmetry is a symmetry between bosons and fermions, and hence necessarily relates particles with different spins. All particles in supersymmetric theories fall into supermultiplets, which have both bosonic and fermionic components. There are two types of supermultiplets which appear in renormalizable field theories: chiral and vector supermultiplets.
Chiral supermultiplets are often denoted by the symbol $`\varphi `$, which can be (for the purpose of this lecture) regarded as a short-handed notation for the three fields: a complex scalar field $`A`$, a Weyl fermion $`\frac{1\gamma _5}{2}\psi =\psi `$, and a non-dynamical (auxiliary) complex field $`F`$. Lagrangians for chiral supermultiplets consist of two parts, Kähler potential and superpotential. The Kähler potential is nothing but the kinetic terms for the fields, usually written with a short-hand notation $`d^4\theta \varphi ^{}\varphi `$, which can be explicitly written down as
$$d^4\theta \varphi _i^{}\varphi _i=_\mu A_i^{}^\mu A_i+\overline{\psi }_ii\gamma ^\mu _\mu \psi _i+F_i^{}F_i.$$
(12)
Note that the field $`F`$ does not have derivatives in the Lagrangian and hence is not a propagating field. One can solve for $`F_i`$ explicitly and eliminate it from the Lagrangian completely.
The superpotential is defined by a holomorphic function $`W(\varphi )`$ of the chiral supermultiplets $`\varphi _i`$. A short-hand notation $`d^2\theta W(\varphi )`$ gives the following terms in the Lagrangian,
$$d^2\theta W(\varphi )=\frac{1}{2}\frac{^2W}{\varphi _i\varphi _j}|_{\varphi _i=A_i}\psi ^i\psi ^j+\frac{W}{\varphi _i}|_{\varphi _i=A_i}F_i.$$
(13)
The first term describes Yukawa couplings between fermionic and bosonic components of the chiral supermultiplets. Using both Eqs. (12) and (13), we can solve for $`F`$ and find
$$F_i^{}=\frac{W}{\varphi _i}|_{\varphi _i=A_i}.$$
(14)
Substituting it back to the Lagrangian, we eliminate $`F`$ and instead find a potential term
$$V_F=\left|\frac{W}{\varphi _i}\right|_{\varphi _i=A_i}^2.$$
(15)
Vector supermultiplets $`W_\alpha `$ ($`\alpha `$ is a spinor index, but never mind), which are supersymmetric generalization of the gauge fields, consist also of three components, a Weyl fermion (gaugino) $`\lambda `$, a vector (gauge) field $`A_\mu `$, and a non-dynamical (auxiliary) real scalar field $`D`$, all in the adjoint representation of the gauge group with the index $`a`$. A short-hand notation of their kinetic terms is
$$d^2\theta W_\alpha ^aW^{\alpha a}=\frac{1}{4}F_{\mu \nu }+\overline{\lambda }^ai\overline{)}D\lambda ^a+\frac{1}{2}D^aD^a.$$
(16)
Note that the field $`D`$ does not have derivatives in the Lagrangian and hence is not a propagating field. One can solve for $`D^a`$ explicitly and eliminate it from the Lagrangian completely.
Since the vector supermultiplets contain gauge fields, chiral supermultiplets which transform non-trivially under the gauge group should also couple to the vector multiplets to make the Lagrangian gauge invariant. This requires the modification of the Kähler potential $`d^4\theta \varphi ^{}\varphi `$ to $`d^4\theta \varphi ^{}e^{2gV}\varphi `$, where $`V`$ is another short-hand notation of the vector multiplet. Then the kinetic terms in Eq. (12) are then modified to
$`{\displaystyle d^4\theta \varphi _i^{}e^{2gV}\varphi _i}`$
$`=D_\mu A_i^{}D^\mu A_i+\overline{\psi }_ii\gamma ^\mu D_\mu \psi _i+F_i^{}F_i\sqrt{2}g(A^{}T^a\lambda ^a\psi )gA^{}T^aD^aA.`$
(17)
Using Eqs. (16,17), one can solve for $`D^a`$ and eliminate it from the Lagrangian, finding a potential term
$$V_D=\frac{g^2}{2}(A^{}T^aA)^2$$
(18)
General supersymmetric Lagrangians are given by Eqs. (17,15,18).<sup>2</sup><sup>2</sup>2We dropped one possible term called Fayet–Illiopoulos $`D`$-term possible for vector supermultiplets of Abelian gauge groups. They are often not useful in phenomenological models, but there are exceptions.
Even though we do not go into formal discussions of supersymmetric field theories, one important theorem must be quoted: the non-renormalization theorem of the superpotential. Under the renormalization of the theories, the superpotential does not receive renormalization at all orders in perturbation theory.<sup>3</sup><sup>3</sup>3There are non-perturbative corrections to the superpotential, however. See, e.g., a review. We will come back to the virtues of this theorem later on.
Finally, let us study a very simple example of superpotential to gain some intuition. Consider two chiral supermultiplets $`\varphi _1`$ and $`\varphi _2`$, with a superpotential
$$W=m\varphi _1\varphi _2.$$
(19)
Following the above prescription, the fermionic components have the Lagrangian
$$\frac{1}{2}\frac{^2W}{\varphi _i\varphi _j}\psi ^i\psi ^j=m\psi _1\psi _2,$$
(20)
while the scalar potential term Eq. (15) gives
$$\left|\frac{W}{\varphi _i}\right|_{\varphi _i=A_i}^2=m^2|A_1|^2m^2|A_2|^2.$$
(21)
Obviously, the terms Eqs. (20,21) are mass terms for the fermionic (Dirac fermion) and scalar components (two complex scalars) of the chiral supermultiplets, with the same mass $`m`$. In general, fermionic and bosonic components in the same supermultiplets are degenerate in supersymmetric theories.
## 3 Softly Broken Supersymmetry
We’ve discussed supersymmetric Lagrangians in the previous section, which always give degenerate bosons and fermions. In the real world, we do not see such degenerate particles with the opposite statistics. Therefore supersymmetry must be broken. We will come back later to briefly discuss various mechanisms which break supersymmetry spontaneously in manifestly supersymmetric theories. In the low-energy effective theories, however, we can just add terms to supersymmetric Lagrangians which break supersymmetry explicitly. The important constraint is that such explicit breaking terms should not spoil the motivation discussed earlier, namely to keep the Higgs mass-squared only logarithmically divergent. Such explicit breaking terms of supersymmetry are called “soft” breakings.
The possible soft breaking terms have been classified. In a theory with a renormalizable superpotential
$$W=\frac{1}{2}\mu _{ij}\varphi _i\varphi _j+\frac{1}{6}\lambda _{ijk}\varphi _i\varphi _j\varphi _k,$$
(22)
the possible soft supersymmetry breaking terms have the following forms:
$$m_{ij}^2A_i^{}A_j,M\lambda \lambda ,\frac{1}{2}b_{ij}\mu _{ij}A_iA_j,\frac{1}{6}a_{ijk}\lambda _{ijk}A_iA_jA_k.$$
(23)
The first one is the masses for scalar components in the chiral supermultiplets, which remove degeneracy between the scalar and spinor components. The next one is the masses for gauginos which remove degeneracy between gauginos and gauge bosons. Finally the last two ones are usually called bilinear and trilinear soft breaking terms with parameters $`b_{ij}`$ and $`a_{ijk}`$ with mass dimension one.
In principle, any terms with couplings with positive mass dimensions are candidates of soft supersymmetry breaking terms. Possibilities in theories without gauge singlets are
$$\psi _i\psi _j,A_i^{}A_jA_k,\psi _i\lambda ^a$$
(24)
Obviously, the first term is possible only in theories with multiplets with vector-like gauge quantum numbers, and the last term with chiral supermultiplets in the adjoint representation. In the presence of gauge singlet chiral supermultiplets, however, such terms cause power divergences and instabilities, and hence are not soft in general. On the other hand, the Minimal Supersymmetric Standard Model, for instance, does not contain any gauge singlet chiral supermultiplets and hence does admit first two possible terms in Eq. (24). There has been some revived interest in these general soft terms. We will not consider these additional terms in the rest of the discussions. It is also useful to know that terms in Eq. (23) can also induce power divergences in the presence of light gauge singlets and heavy multiplets.
It is instructive to carry out some explicit calculations of Higgs boson self-energy in supersymmetric theories with explicit soft supersymmetry breaking terms. Let us consider the coupling of the Higgs doublet chiral supermultiplet $`H`$ to left-handed $`Q`$ and right-handed $`T`$ chiral supermultiplets,<sup>4</sup><sup>4</sup>4As will be explained in the next section, the right-handed spinors all need to be charged-conjugated to the left-handed ones in order to be part of the chiral supermultiplets. Therefore the chiral supermultiplet $`T`$ actually contains the left-handed Weyl spinor $`(t_R)^c`$. The Higgs multiplet here will be denoted $`H_u`$ in later sections. given by the superpotential term
$$W=h_tQTH_u.$$
(25)
This superpotential term gives rise to terms in the Lagrangian<sup>5</sup><sup>5</sup>5We dropped terms which do not contribute to the Higgs boson self-energy at the one-loop level.
$$h_tQTH_uh_t^2|\stackrel{~}{Q}|^2|H_u|^2h_t^2|\stackrel{~}{T}|^2|H_u|^2m_Q^2|\stackrel{~}{Q}|^2m_T^2|\stackrel{~}{T}|^2h_tA_t\stackrel{~}{Q}\stackrel{~}{T}H_u,$$
(26)
where $`m_Q^2`$, $`m_T^2`$, and $`A_t`$ are soft parameters. Note that the fields $`Q`$, $`T`$ are spinor and $`\stackrel{~}{Q}`$, $`\stackrel{~}{T}`$, $`H_u`$ are scalar components of the chiral supermultiplets (an unfortunate but common notation in the literature). This explicit Lagrangian allows us to easily work out the one-loop self-energy diagrams for the Higgs doublet $`H_u`$, after shifting the field $`H_u`$ around its vacuum expectation value (this also generates mass terms for the top quark and the scalars which have to be consistently included). The diagram with top quark loop from the first term in Eq. (26) is quadratically divergent (negative). The contractions of $`\stackrel{~}{Q}`$ or $`\stackrel{~}{T}`$ in the next two terms also generate (positive) contributions to the Higgs self-energy. In the absence of soft parameters $`m_Q^2=m_T^2=0`$, these two contributions precisely cancel with each other, consistent with the non-renormalization theorem which states that no mass terms (superpotential terms) can be generated by renormalizations. However, the explicit breaking terms $`m_Q^2`$, $`m_T^2`$ make the cancellation inexact. With a simplifying assumption $`m_Q^2=m_T^2=\stackrel{~}{m}^2`$, we find
$$\delta m_H^2=\frac{6h_t^2}{(4\pi )^2}\stackrel{~}{m}^2\mathrm{log}\frac{\mathrm{\Lambda }^2}{\stackrel{~}{m}^2}.$$
(27)
Here, $`\mathrm{\Lambda }`$ is the ultraviolet cutoff of the one-loop diagrams. Therefore, these mass-squared parameters are indeed “soft” in the sense that they do not produce power divergences. Similarly, the diagrams with two $`h_tA_t`$ couplings with scalar top loop produce only a logarithmic divergent contribution.
## 4 The Minimal Supersymmetric Standard Model
Encouraged by the discussion in the previous section that the supersymmetry can be explicitly broken while retaining the absence of power divergences, we now try to promote the Standard Model to a supersymmetric theory. The Minimal Supersymmetric Standard Model (MSSM) is a supersymmetric version of the Standard Model with the minimal particle content.
### 4.1 Particle Content
The first task is to promote all fields in the Standard Model to appropriate supermultiplets. This is obvious for the gauge bosons: they all become vector multiplets. For the quarks and leptons, we normally have left-handed and right-handed fields in the Standard Model. In order to promote them to chiral supermultiplets, however, we need to make all fields left-handed Weyl spinors. This can be done by charge-conjugating all right-handed fields. Therefore, when we refer to supermultiplets of the right-handed down quark, say, we are actually talking about chiral supermultiplets whose left-handed spinor component is the left-handed anti-down quark field. As for the Higgs boson, the field Eq. (1) in the Standard Model can be embedded into a chiral supermultiplet $`H_u`$. It can couple to the up-type quarks and generate their masses upon the symmetry breaking. In order to generate down-type quark masses, however, we normally use
$$i\sigma _2H^{}=\left(\begin{array}{c}H^+\\ H^0\end{array}\right)=\left(\begin{array}{c}H^0\\ H^{}\end{array}\right).$$
(28)
Unfortunately, this trick does not work in a supersymmetric fashion because the superpotential $`W`$ must be a holomorphic function of the chiral supermultiplets and one is not allowed to take a complex conjugation of this sort. Therefore, we need to introduce another chiral supermultiplet $`H_d`$ which has the same gauge quantum numbers of $`i\sigma _2H^{}`$ above.<sup>6</sup><sup>6</sup>6Another reason to need both $`H_u`$ and $`H_d`$ chiral supermultiplets is to cancel the gauge anomalies arising from their spinor components.
In all, the chiral supermultiplets in the Minimal Supersymmetric Standard Model are listed in Table 3.
The particles in the MSSM are referred to as follows.<sup>7</sup><sup>7</sup>7When I first learned supersymmetry, I didn’t believe it at all. Doubling the degrees of freedom looked too much to me, until I came up with my own argument at the beginning of the lecture. The funny names for the particles were yet another reason not to believe in it. It doesn’t sound scientific. Once supersymmetry will be discovered, we definitely need better sounding names! First of all, all quarks, leptons are called just in the same way as in the Standard Model, namely electron, electron-neutrino, muon, muon-neutrino, tau, tau-neutrino, up, down, strange, charm, bottom, top. Their superpartners, which have spin 0, are named with “s” at the beginning, which stand for “scalar.” They are denoted by the same symbols as their fermionic counterpart with the tilde. Therefore, the superpartner of the electron is called “selectron,” and is written as $`\stackrel{~}{e}`$. All these names are funny, but probably the worst one of all is the “sstrange” ($`\stackrel{~}{s}`$), which I cannot pronounce at all. Superpartners of quarks are “squarks,” and those of leptons are “sleptons.” Sometimes all of them are called together as “sfermions,” which does not make sense at all because they are bosons. The Higgs doublets are denoted by capital $`H`$, but as we will see later, their physical degrees of freedom are $`h^0`$, $`H^0`$, $`A^0`$ and $`H^\pm `$. Their superpartners are called “higgsinos,” written as $`\stackrel{~}{H}_u^0`$, $`\stackrel{~}{H}_u^+`$, $`\stackrel{~}{H}_d^{}`$, $`\stackrel{~}{H}_d^0`$. In general, fermionic superpartners of boson in the Standard Model have “ino” at the end of the name. Spin 1/2 superpartners of the gauge bosons are “gauginos” as mentioned in the previous section, and for each gauge groups: gluino for gluon, wino for $`W`$, bino for $`U(1)_Y`$ gauge boson $`B`$. As a result of the electroweak symmetry breaking, all neutral “inos”, namely two neutral higgsinos, the neutral wino $`\stackrel{~}{W}^3`$ and the bino $`\stackrel{~}{B}`$ mix with each other to form four Majorana fermions. They are called “neutralinos” $`\stackrel{~}{\chi }_i^0`$ for $`i=1,2,3,4`$. Similarly, the charged higgsinos $`\stackrel{~}{H}_u^+`$, $`\stackrel{~}{H}_d^{}`$, $`\stackrel{~}{W}^{}`$, $`\stackrel{~}{W}^+`$ mix and form two massive Dirac fermions “charginos” $`\stackrel{~}{\chi }_i^\pm `$ for $`i=1,2`$. All particles with tilde do not exist in the non-supersymmetric Standard Model. Once we introduce $`R`$-parity in a later section, the particles with tilde have odd $`R`$-parity.
### 4.2 Superpotential
The $`SU(3)_C\times SU(2)_L\times U(1)_Y`$ gauge invariance allows the following terms in the superpotential
$`W`$ $`=`$ $`\lambda _u^{ij}Q_iU_jH_u+\lambda _d^{ij}Q_iD_jH_d+\lambda _e^{ij}L_iE_jH_d+\mu H_uH_d`$ (29)
$`+\lambda _u^{ijk}U_iD_jD_k+\lambda _d^{ijk}Q_iD_jL_k+\lambda _e^{ijk}L_iE_jL_k+\mu _i^{}L_iH_u.`$
The first three terms correspond to the Yukawa couplings in the Standard Model (with exactly the same number of parameters). The subscripts $`i,j,k`$ are generation indices. The parameter $`\mu `$ has mass dimension one and gives a supersymmetric mass to both fermionic and bosonic components of the chiral supermultiplets $`H_u`$ and $`H_d`$. The terms in the second line of Eq. (29) are in general problematic as they break the baryon ($`B`$) or lepton ($`L`$) numbers.
If the superpotential contains both $`B`$\- and $`L`$-violating terms, such as $`\lambda _u^{112}U_1D_1D_2`$ and $`\lambda _d^{121}Q_1D_2L_1`$, one can exchange $`\stackrel{~}{D}_2=\stackrel{~}{s}`$ to generate a four-fermion operator
$$\frac{\lambda _u^{112}\lambda _d^{121}}{m_{\stackrel{~}{s}}^2}(u_Rd_R)(Q_1L_1),$$
(30)
where the spinor indices are contracted in each parentheses and the color indices by the epsilon tensor. Such an operator would contribute to the proton decay process $`pe^+\pi ^0`$ at a rate of $`\mathrm{\Gamma }\lambda ^4m_p^5/m_{\stackrel{~}{s}}^4`$, and hence the partial lifetime of the order of
$$\tau _p6\times 10^{13}\mathrm{sec}\left(\frac{m_{\stackrel{~}{s}}}{1\mathrm{TeV}}\right)^4\frac{1}{\lambda ^4}.$$
(31)
Recall that the experimental limit on the proton partial lifetime in this mode is $`\tau _p>1.6\times 10^{33}`$ years. Unless the coupling constants are extremely small, this is clearly a disaster.
### 4.3 $`R`$-parity
To avoid this problem of too-rapid proton decay, a common assumption is a discrete symmetry called $`R`$-parity (or matter parity). The $`Z_2`$ discrete charge is given by
$$R_p=(1)^{2s+3B+L}$$
(32)
where $`s`$ is the spin of the particle. Under $`R_p`$, all standard model particles, namely quarks, leptons, gauge bosons, and Higgs bosons, carry even parity, while their superpartners odd due to the $`(1)^{2s}`$ factor. Once this discrete symmetry is imposed, all terms in the second line of Eq. (29) will be forbidden, and we do not generate a dangerous operator such as that in Eq. (30). Indeed, $`B`$\- and $`L`$-numbers are now accidental symmetries of the MSSM Lagrangian as a consequence of the supersymmetry, gauge invariance, renormalizability and $`R`$-parity conservation.
One immediate consequence of the conserved $`R`$-parity is that the lightest particle with odd $`R`$-parity, i.e., the Lightest Supersymmetric Particle (LSP), is stable. Another consequence is that one can produce (or annihilate) superparticles only pairwise. These two points have important implications on the collider phenomenology and cosmology. Since the LSP is stable, its cosmological relic is a good (and arguably the best) candidate for the Cold Dark Matter particles (see, e.g., a review on this subject). If so, we do not want it to be electrically charged and/or strongly interacting; otherwise we should have detected them already. Then the LSP should be a superpartner of $`Z`$, $`\gamma `$, or neutral Higgs bosons or their linear combination (called neutralino).<sup>8</sup><sup>8</sup>8A sneutrino can in principle be the LSP,, but it cannot be the CDM to avoid constraints from the direct detection experiment for the CDM particles. It becomes a viable candidate again if there is a large lepton number violation. On the other hand, the superparticles can be produced only in pairs and they decay eventually into the LSP, which escapes detection. This is why the typical signature of supersymmetry at collider experiments is the missing energy/momentum.
The phenomenology of $`R`$-parity breaking models has been also studied. If either $`B`$-violating or $`L`$-violating terms exist in Eq. (29), but not both, they would not induce proton decay. However they can still produce $`n`$-$`\overline{n}`$ oscillation and a plethora of flavor-changing phenomena. We refer to a recent compilation of phenomenological constraints for further details.
### 4.4 Soft Supersymmetry Breaking Terms
In addition to the interactions that arise from the superpotential Eq. (29), we should add soft supersymmetry breaking terms to the Lagrangian as we have not seen any of the superpartners of the Standard Model particles. Following the general classifications in Eq. (23), and assuming $`R`$-parity conservation, they are given by
$`_{\mathrm{𝑠𝑜𝑓𝑡}}=_1+_2,`$ (33)
$`_1=m_Q^{2ij}\stackrel{~}{Q}_i^{}\stackrel{~}{Q}_jm_U^{2ij}\stackrel{~}{U}_i^{}\stackrel{~}{U}_jm_D^{2ij}\stackrel{~}{D}_i^{}\stackrel{~}{D}_j`$
$`m_L^{2ij}\stackrel{~}{L}_i^{}\stackrel{~}{L}_jm_E^{2ij}\stackrel{~}{E}_i^{}\stackrel{~}{E}_jm_{H_u}^2|H_u|^2m_{H_d}^2|H_d|^2,`$ (34)
$`_2=A_u^{ij}\lambda _u^{ij}\stackrel{~}{Q}_i\stackrel{~}{U}_jH_uA_d^{ij}\lambda _d^{ij}\stackrel{~}{Q}_i\stackrel{~}{D}_jH_dA_l^{ij}\lambda _e^{ij}\stackrel{~}{Q}_i\stackrel{~}{U}_jH_d+B\mu H_uH_d+c.c.`$
(35)
The mass-squared parameters for scalar quarks (squarks) and scalar leptons (sleptons) are all three-by-three hermitian matrices, while the trilinear couplings $`A^{ij}`$ and the bilinear coupling $`B`$ of mass dimension one are general complex numbers.<sup>9</sup><sup>9</sup>9It is unfortunate that the notation $`A`$ is used both for the scalar components of chiral supermultiplets and the trilinear couplings. Hopefully one can tell them apart from the context.
### 4.5 Higgs Sector
It is of considerable interest to look closely at the Higgs sector of the MSSM. Following the general form of the supersymmetric Lagrangians Eqs. (17,15,18) with the superpotential $`W=\mu H_uH_d`$ in Eq. (29) as well as the soft parameters in Eq. (34), the potential for the Higgs bosons is given as
$`V={\displaystyle \frac{g^2}{2}}\left(H_u^{}{\displaystyle \frac{1}{2}}H_u+H_d^{}{\displaystyle \frac{1}{2}}H_d\right)^2+{\displaystyle \frac{g^2}{2}}\left(H_u^{}{\displaystyle \frac{\stackrel{}{\tau }}{2}}H_u+H_d^{}{\displaystyle \frac{\stackrel{}{\tau }}{2}}H_d\right)^2`$
$`+\mu ^2(|H_u|^2+|H_d|^2)+m_{H_u}^2|H_u|^2+m_{H_d}^2|H_d|^2(B\mu H_uH_d+c.c.)`$ (36)
It turns out that it is always possible to gauge-rotate the Higgs bosons such that
$$H_u=\left(\begin{array}{c}0\\ v_u\end{array}\right),H_d\left(\begin{array}{c}v_d\\ 0\end{array}\right),$$
(37)
in the vacuum. Since only electrically neutral components have vacuum expectation values, the vacuum necessarily conserves $`U(1)_{\mathrm{QED}}`$.<sup>10</sup><sup>10</sup>10This is not necessarily true in general two-doublet Higgs Models. Consult a review. Writing the potential (36) down using the expectation values (37), we find
$$V=\frac{g_Z^2}{8}(v_u^2v_d^2)^2+(v_uv_d)\left(\begin{array}{cc}\mu ^2+m_{H_u}^2& B\mu \\ B\mu & \mu ^2+m_{H_d}^2\end{array}\right)\left(\begin{array}{c}v_u\\ v_d\end{array}\right),$$
(38)
where $`g_Z^2=g^2+g^2`$. In order for the Higgs bosons to acquire the vacuum expectation values, the determinant of the mass matrix at the origin must be negative,
$$\mathrm{det}\left(\begin{array}{cc}\mu ^2+m_{H_u}^2& B\mu \\ B\mu & \mu ^2+m_{H_d}^2\end{array}\right)<0.$$
(39)
However, there is a danger that the direction $`v_u=v_d`$, which makes the quartic term in the potential identically vanish, may be unbounded from below. For this not to occur, we need
$$\mu ^2+m_{H_u}^2+\mu ^2+m_{H_d}^2>2\mu B.$$
(40)
In order to reproduce the mass of the $`Z`$-boson correctly, we need
$$v_u=\frac{v}{\sqrt{2}}\mathrm{sin}\beta ,v_d=\frac{v}{\sqrt{2}}\mathrm{cos}\beta ,v=250\mathrm{GeV}.$$
(41)
The vacuum minimization conditions are given by $`V/v_u=V/v_d=0`$ from the potential Eq. (38). Using Eq. (41), we obtain
$$\mu ^2=\frac{m_Z^2}{2}+\frac{m_{H_d}^2m_{H_u}^2\mathrm{tan}^2\beta }{\mathrm{tan}^2\beta 1},$$
(42)
and
$$B\mu =(2\mu ^2+m_{H_u}^2+m_{H_d}^2)\mathrm{sin}\beta \mathrm{cos}\beta .$$
(43)
Because there are two Higgs doublets, each of which with four real scalar fields, the number of degrees of freedom is eight before the symmetry breaking. However three of them are eaten by $`W^+`$, $`W^{}`$ and $`Z`$ bosons, and we are left with five physics scalar particles. There are two CP-even scalars $`h^0`$, $`H^0`$, one CP-odd scalar $`A^0`$, and two charged scalars $`H^+`$ and $`H^{}`$. Their masses can be worked out from the potential (38):
$$m_A^2=2\mu ^2+m_{H_u}^2+m_{H_d}^2,m_{H^\pm }^2=m_W^2+m_A^2,$$
(44)
and
$$m_{h^0}^2,m_{H^0}^2=\frac{1}{2}\left(m_A^2+m_Z^2\pm \sqrt{(m_A^2+m_Z^2)^24m_Z^2m_A^2\mathrm{cos}^22\beta }\right).$$
(45)
A very interesting consequence of the formula Eq. (45) is that the lighter CP-even Higgs mass $`m_{h^0}^2`$ is maximized when $`\mathrm{cos}^22\beta =1`$: $`m_{h^0}^2=(m_A^2+m_Z^2|m_A^2m_Z^2|)/2`$. When $`m_A<m_Z`$, we obtain $`m_{h^0}^2=m_A^2<m_Z^2`$, while when $`m_A>m_Z`$, $`m_{h^0}^2=m_Z^2`$. Therefore in any case we find
$$m_{h^0}m_Z.$$
(46)
This is an important prediction in the MSSM. The reason why the masses of the Higgs boson are related to the gauge boson masses is that the Higgs quartic couplings in Eq. (36) are all determined by the gauge couplings because they originate from the elimination of the auxiliary $`D`$-fields in Eq. (17).
Unfortunately, the prediction Eq. (46) is modified at the one-loop level, approximately as
$$\mathrm{\Delta }(m_{h^0}^2)=\frac{N_c}{4\pi ^2}h_t^4v^2\mathrm{sin}^4\beta \mathrm{log}\left(\frac{m_{\stackrel{~}{t}_1}m_{\stackrel{~}{t}_2}}{m_t^2}\right).$$
(47)
With the scalar top mass of up to 1 TeV, the lightest Higgs mass is pushed up to about 130 GeV. (See also the latest analysis including resummed two-loop contribution.)
The parameter space of the MSSM Higgs sector can be described by two parameters. This is because the potential Eq. (38) has three independent parameters, $`\mu ^2+m_{H_u}^2`$, $`\mu ^2+m_{H_d}^2`$, and $`B\mu `$, while one combination is fixed by the $`Z`$-mass Eq. (39). It is customary to pick either $`(m_A,\mathrm{tan}\beta )`$, or $`(m_{h^0},\mathrm{tan}\beta )`$ to present experimental constraints. The current experimental constraint on this parameter space is shown in Fig. 1.<sup>11</sup><sup>11</sup>11The large $`\mathrm{tan}\beta `$ region may appear completely excluded in the plot, but this is somewhat misleading; it is due to the parameterization $`(m_{h^0},\mathrm{tan}\beta )`$ which squeezes the $`m_{h^0}`$ region close to the theoretical upper bound to a very thin one. In the $`(m_A,\mathrm{tan}\beta )`$ parameterization, one can see the allowed region much more clearer.
The range of the Higgs mass predicted in the MSSM is not necessarily an easy range for the LHC experiments, but three-years’ running at the high luminosity is supposed to cover the entire MSSM parameter space, by employing many different production/decay modes as seen in Fig. 2.
### 4.6 Neutralinos and Charginos
Once the electroweak symmetry is broken, and since supersymmetry is already explicitly broken in the MSSM, there is no quantum number which can distinguish two neutral higgsino states $`\stackrel{~}{H}_u^0`$, $`\stackrel{~}{H}_d^0`$, and two neutral gaugino states $`\stackrel{~}{W}^3`$ (neutral wino) and $`\stackrel{~}{B}`$ (bino). They have four-by-four Majorana mass matrix
$`{\displaystyle \frac{1}{2}}\times `$ (56)
$`(\stackrel{~}{B}\stackrel{~}{W}^3\stackrel{~}{H}_d^0\stackrel{~}{H}_u^0)\left(\begin{array}{cccc}M_1& 0& m_Zs_Wc_\beta & m_Zs_Ws_\beta \\ 0& M_2& m_Zc_Wc_\beta & m_Zc_Ws_\beta \\ m_Zs_Wc_\beta & m_Zc_Wc_\beta & 0& \mu \\ m_Zs_Ws_\beta & m_Zc_Ws_\beta & \mu & 0\end{array}\right)\left(\begin{array}{c}\stackrel{~}{B}\\ \stackrel{~}{W}^3\\ \stackrel{~}{H}_d^0\\ \stackrel{~}{H}_u^0\end{array}\right).`$
Here, $`s_W=\mathrm{sin}\theta _W`$, $`c_W=\mathrm{cos}\theta _W`$, $`s_\beta =\mathrm{sin}\beta `$, and $`c_\beta =\mathrm{cos}\beta `$. Once $`M_1`$, $`M_2`$, $`\mu `$ exceed $`m_Z`$, which is preferred given the current experimental limits, one can regard components proportional to $`m_Z`$ as small perturbations. Then the neutralinos are close to their weak eigenstates, bino, wino, and higgsinos. But the higgsinos in this limit are mixed to form symmetric and anti-symmetric linear combinations $`\stackrel{~}{H}_S^0=(\stackrel{~}{H}_d^0+\stackrel{~}{H}_u^0)/\sqrt{2}`$ and $`\stackrel{~}{H}_A^0=(\stackrel{~}{H}_d^0\stackrel{~}{H}_u^0)/\sqrt{2}`$.
Similarly two positively charged inos: $`\stackrel{~}{H}_u^+`$ and $`\stackrel{~}{W}^+`$, and two negatively charged inos: $`\stackrel{~}{H}_d^{}`$ and $`\stackrel{~}{W}^{}`$ mix. The mass matrix is given by
$$(\stackrel{~}{W}^{}\stackrel{~}{H}_d^{})\left(\begin{array}{cc}M_2& \sqrt{2}m_Ws_\beta \\ \sqrt{2}m_Wc_\beta & \mu \end{array}\right)\left(\begin{array}{c}\stackrel{~}{W}^+\\ \stackrel{~}{H}_u^+\end{array}\right)+c.c.$$
(58)
Again once M2,μ
>mW
>subscript𝑀2𝜇subscript𝑚𝑊M_{2},\mu\mathrel{\mathchoice{\lower 3.0pt\vbox{\halign{$\m@th\displaystyle\hfill#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\m@th\textstyle\hfill#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\m@th\scriptstyle\hfill#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\m@th\scriptscriptstyle\hfill#\hfil$\cr>\crcr\sim\crcr}}}}m_{W}, the chargino states are close to the weak eigenstates winos and higgsinos.
### 4.7 Squarks, Sleptons
The mass terms of squarks and sleptons are also modified after the electroweak symmetry breaking. There are four different contributions. One is the supersymmetric piece coming from the $`|W/\varphi _i|^2`$ terms in Eq. (15) with $`\varphi _i=Q,U,D,L,E`$. These terms add $`m_f^2`$ where $`m_f`$ is the mass of the quarks and leptons from their Yukawa couplings to the Higgs boson. Next one is combing from the $`|W/\varphi _i|^2`$ terms in Eq. (15) with $`\varphi _i=H_u`$ or $`H_d`$ in the superpotential Eq. (29). Because of the $`\mu `$ term,
$`{\displaystyle \frac{W}{H_u^0}}`$ $`=`$ $`\mu H_d^0+\lambda _u^{ij}\stackrel{~}{Q}_i\stackrel{~}{U}_j,`$ (59)
$`{\displaystyle \frac{W}{H_d^0}}`$ $`=`$ $`\mu H_d^0+\lambda _d^{ij}\stackrel{~}{Q}_i\stackrel{~}{D}_j+\lambda _e^{ij}\stackrel{~}{L}_i\stackrel{~}{E}_j.`$ (60)
Taking the absolute square of these two expressions pick the cross terms together with $`H_d^0=v\mathrm{cos}\beta /\sqrt{2}`$, $`H_u^0=v\mathrm{sin}\beta /\sqrt{2}`$ and we obtain mixing between $`\stackrel{~}{Q}`$ and $`\stackrel{~}{U}`$, $`\stackrel{~}{Q}`$ and $`\stackrel{~}{D}`$, and $`\stackrel{~}{L}`$ and $`\stackrel{~}{E}`$. Similarly, the vacuum expectation values of the Higgs bosons in the trilinear couplings Eq. (35) also generate similar mixing terms. Finally, the $`D`$-term potential after eliminating the auxiliary field $`D`$ Eq. (18) also give contributions to the scalar masses $`m_Z^2(I_3Q\mathrm{sin}^2\theta _W)\mathrm{cos}2\beta `$. Therefore, the mass matrix of stop, for instance, is given as
$`(\stackrel{~}{t}_L^{}\stackrel{~}{t}_R^{})`$ (65)
$`\left(\begin{array}{cc}m_{Q_3}^2+m_t^2+m_Z^2(\frac{1}{2}\frac{2}{3}s_W^2)c_{2\beta }& m_t(A_t\mu \mathrm{cot}\beta )\\ m_t(A_t\mu \mathrm{cot}\beta )& m_{U_3}^2+m_t^2+m_Z^2(\frac{2}{3}s_W^2)c_{2\beta }\end{array}\right)\left(\begin{array}{c}\stackrel{~}{t}_L\\ \stackrel{~}{t}_R\end{array}\right),`$
with $`c_{2\beta }=\mathrm{cos}2\beta `$. Here, $`\stackrel{~}{t}_L`$ is the up component of $`\stackrel{~}{Q}_3`$, and $`\stackrel{~}{t}_R=\stackrel{~}{T}^{}`$. For first and second generation particles, the off-diagonal terms are negligible for most purposes. They may, however, be important when their loops in flavor-changing processes are considered.
### 4.8 What We Gained in the MSSM
It is useful to review here what we have gained in the MSSM over what we had in the Standard Model. The main advantage of the MSSM is of course what motivated the supersymmetry to begin with: the absence of the quadratic divergences as seen in Eq. (27). This fact allows us to apply the MSSM down to distance scales much shorter than the electroweak scale, and hence we can at least hope that many of the puzzles discussed at the beginning of the lecture to be solved by physics at the short distance scales.
There are a few amusing and welcome by-products of supersymmetry beyond this very motivation. First of all, the Higgs doublet in the Standard Model appears so unnatural partly because it is the only scalar field introduced just for the sake of the electroweak symmetry breaking. In the MSSM, however, there are so many scalar fields: 15 complex scalar fields for each generation and two in each Higgs doublet. Therefore, the Higgs bosons are just “one of them.” Then the question about the electroweak symmetry breaking is addressed in a completely different fashion: why is it only the Higgs bosons that condense? In fact, one can even partially answer this question in the renormalization group analysis in the next sections where “typically” (we will explain what we mean by this) it is only the Higgs bosons which acquire negative mass squared (39) while the masses-squared of all the other scalars “naturally” remain positive. Finally, the absolute upper bound on the lightest CP-even Higgs boson is falsifiable by experiments.
However, life is not as good as we wish. We will see that there are very stringent low-energy constraints on the MSSM in the next section.
## 5 Low-Energy Constraints
Despite the fact that we are interested in superparticles in the 100–1000 GeV range, which we are just starting to explore in collider searches, there are many amazingly stringent low-energy constraints on superparticles.
One of the most stringent constraints comes from the $`K^0`$$`\overline{K}^0`$ mixing parameters $`\mathrm{\Delta }m_K`$ and $`\epsilon _K`$. The main reason for the stringent constraints is that the scalar masses-squared in the MSSM Lagrangian Eq. (34) can violate flavor, i.e., the scalar masses-squared matrices are not necessarily diagonal in the basis where the corresponding quark mass matrices are diagonal.
To simplify the discussion, let us concentrate only on the first and the second generations (ignore the third). We also go to the basis where the down-type Yukawa matrix $`\lambda _d^{ij}`$ is diagonal, such that
$$\lambda _d^{ij}v_d=\left(\begin{array}{cc}m_d& 0\\ 0& m_s\end{array}\right).$$
(67)
Therefore the states $`K^0=(d\overline{s})`$, $`\overline{K}^0=(s\overline{d})`$ are well-defined in this basis. In the same basis, however, the squark masses-squared can have off-diagonal elements in general,
$$m_Q^{2ij}=\left(\begin{array}{cc}m_{\stackrel{~}{d}_L}^2& m_{Q,12}^2\\ m_{Q,12}^2& m_{\stackrel{~}{s}_L}^2\end{array}\right),m_D^{2ij}=\left(\begin{array}{cc}m_{\stackrel{~}{d}_R}^2& m_{D,12}^2\\ m_{D,12}^2& m_{\stackrel{~}{s}_R}^2\end{array}\right).$$
(68)
Since their off-diagonal elements will be required to be small (as we will see later), it is convenient to treat them as small perturbation. We insert the off-diagonal elements as two-point Feynman vertices which change the squark flavor $`\stackrel{~}{d}_{L,R}\stackrel{~}{s}_{L,R}`$ in the diagrams. To simplify the discussion further, we assume that all squarks and gluino the are comparable in their masses $`\stackrel{~}{m}`$. Then the relevant quantities are given in terms of the ratio $`(\delta _{12}^d)_{LL}m_{Q,12}^2/\stackrel{~}{m}^2`$ (and similarly $`(\delta _{12}^d)_{RR}=m_{D,12}^2/\stackrel{~}{m}^2`$), as depicted in Fig. 3. The operator from this Feynman diagram is estimated approximately as
$$0.005\alpha _s^2\frac{(\delta _{12}^d)_{LL}^2}{\stackrel{~}{m}^2}(\overline{d}_L\gamma ^\mu s_L)(\overline{d}_L\gamma _\mu s_L).$$
(69)
This operator is further sandwiched between $`K^0`$ and $`\overline{K}^0`$ states, and we find
$`\mathrm{\Delta }m_K^20.005f_K^2m_K^2\alpha _s^2(\delta _{12}^d)_{LL}^2{\displaystyle \frac{1}{\stackrel{~}{m}^2}}`$
$`=1.2\times 10^{12}\mathrm{GeV}^2\left({\displaystyle \frac{f_K}{160\mathrm{MeV}}}\right)^2\left({\displaystyle \frac{\alpha _s}{0.1}}\right)^2(\delta _{12}^d)_{LL}^2<3.5\times 10^{15}\mathrm{GeV}^2,`$
(70)
where the last inequality is the phenomenological constraint in the absence of accidental cancellations. This requires
(δ12d)LL
<0.05(m~500GeV)
<subscriptsuperscriptsubscript𝛿12𝑑𝐿𝐿0.05~𝑚500GeV(\delta_{12}^{d})_{LL}\mathrel{\mathchoice{\lower 3.0pt\vbox{\halign{$\m@th\displaystyle\hfill#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\m@th\textstyle\hfill#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\m@th\scriptstyle\hfill#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\m@th\scriptscriptstyle\hfill#\hfil$\cr<\crcr\sim\crcr}}}}0.05\left(\frac{\tilde{m}}{500~{}{\rm GeV}}\right) (71)
and hence the off-diagonal element $`m_{Q,12}^2`$ must be small. It turns out that the product $`(\delta _{12}^d)_{LL}(\delta _{12}^d)_{RR}`$ is more stringently constrained, especially its imaginary part from $`\epsilon _K`$. Much more careful and detailed analysis than the above order-of-magnitude estimate gives
$$\mathrm{Re}\left[(\delta _{12}^d)_{LL}(\delta _{12}^d)_{RR}\right]<(0.016)^2,\mathrm{Im}\left[(\delta _{12}^d)_{LL}(\delta _{12}^d)_{RR}\right]<(0.0022)^2.$$
(72)
There are many other low-energy observables, such as electron and neutron electric dipole moments (EDM), $`\mu e\gamma `$, which place important constraints on the supersymmetry parameters.
There are various ways to avoid such low-energy constraints on supersymmetry. The first one is called “universality” of soft parameters. It is simply assumed that the scalar masses-squared matrices are proportional to identity matrices, i.e., $`m_Q^2,m_U^2,m_D^2\mathrm{𝟏}`$. Then no matter what rotation is made in order to go to the basis where the quark masses are diagonal, the identity matrices stay the same, and hence the off-diagonal elements are never produced. There has been many proposals to generate universal scalar masses either by the mediation mechanism of the supersymmetry breaking such as the gauge mediated (see reviews), anomaly mediated, or gaugino mediated supersymmetry breaking, or by non-Abelian flavor symmetries. The second possibility is called “alignment,” where certain flavor symmetries should be responsible for “aligning” the quark and squark mass matrices such that the squark masses are almost diagonal in the same basis where the down-quark masses are diagonal. Because of the CKM matrix it is impossible to do this both for down-quark and up-quark masses. Since the phenomenological constraints in the up-quark sector are much weaker than in the down-quark sector, this choice would alleviate many of the low-energy constraints (except for flavor-diagonal CP-violation such as EDMs). Finally there is a possibility called “decoupling,” which assumes first- and second-generation superpartners much heavier than TeV while keeping the third-generation superpartners as well as gauginos in the 100 GeV range to keep the Higgs self-energy small enough. Even though this idea suffers from a fine-tuning problem in general, many models had been constructed to achieve such a split mass spectrum recently.
In short, the low-energy constraints are indeed very stringent, but there are many ideas to avoid such constraints naturally within certain model frameworks. Especially given the fact that we still do not know any of the superparticle masses experimentally, one cannot make the discussions more clear-cut at this stage. On the other hand, important low-energy effects of supersymmetry are still being discovered in the literature, such as muon $`g2`$, and direct CP-violation. They may be even more possible low-energy manifestations of supersymmetry which have been missed so far.
## 6 Renormalization Group Analyses
Once supersymmetry protects the Higgs self-energy against corrections from the short distance scales, or equivalently, the high energy cutoff scales, it becomes important to connect physics at the electroweak scale where we can do measurements to the fundamental parameters defined at high energy scales. This can be done by studying the renormalization-group evolution of parameters. It also becomes a natural expectation that the supersymmetry breaking itself originates at some high energy scale. If this is the case, the soft supersymmetry breaking parameters should also be studied using the renormalization-group equations. We study the renormalization-group evolution of various parameters in the softly-broken supersymmetric Lagrangian at the one-loop level.<sup>12</sup><sup>12</sup>12Recently, there have been developments in obtaining and understanding all-order beta functions for gauge coupling constants and soft parameters. If supersymmetry indeed turns out to be the choice of nature, the renormalization-group analysis will be crucial in probing physics at high energy scales using the observables at the TeV-scale collider experiments.
### 6.1 Gauge Coupling Constants
The first parameters to be studied are naturally the coupling constants in the Standard Model. The running of the gauge couplings constants are described in term of the beta functions, and their one-loop solutions in non-supersymmetric theories are given by
$$\frac{1}{g^2(\mu )}=\frac{1}{g^2(\mu ^{})}+\frac{b_0}{8\pi ^2}\mathrm{log}\frac{\mu }{\mu ^{}},$$
(73)
with
$$b_0=\frac{11}{3}C_2(G)\frac{2}{3}S_f\frac{1}{3}S_b.$$
(74)
This formula is for Weyl fermions $`f`$ and complex scalars $`b`$. The group theory factors are defined by
$`\delta ^{ad}C_2(G)`$ $`=`$ $`f^{abc}f^{dbc}`$ (75)
$`\delta ^{ab}S_{f,b}`$ $`=`$ $`\mathrm{Tr}T^aT^b`$ (76)
and $`C_2(G)=N_c`$ for SU($`N_c`$) groups and $`S_{f,b}=1/2`$ for their fundamental representations.
In supersymmetric theories, there is always the gaugino multiplet in the adjoint representation of the gauge group. They contribute to Eq. (74) with $`S_f=C_2(G)`$, and therefore the total contribution of the vector supermultiplet is $`3C_2(G)`$. On the other hand, the chiral supermultiplets have a Weyl spinor and a complex scalar, and the last two terms in Eq. (74) are always added together to $`S_f=S_b`$. Therefore, the beta function coefficients simplify to
$$b_0=3C_2(G)S_f.$$
(77)
Given the beta functions, it is easy to work out how the gauge coupling constants measured accurately at LEP/SLC evolve to higher energies.
One interesting possibility is that the gauge groups in the Standard Model $`SU(3)_C\times SU(2)_L\times U(1)_Y`$ may be embedded into a simple group, such as $`SU(5)`$ or $`SO(10)`$, at some high energy scale, called “grand unification.” The gauge coupling constants at $`\mu m_Z`$ are approximately $`\alpha ^1=129`$, $`\mathrm{sin}^2\theta _W0.232`$, and $`\alpha _s^1=0.119`$. In the $`SU(5)`$ normalization, the $`U(1)`$ coupling constant is given by $`\alpha _1=\frac{5}{3}\alpha ^{}=\frac{5}{3}\alpha /\mathrm{cos}^2\theta _W`$. It turns out that the gauge coupling constants become equal at $`\mu 2\times 10^{16}`$ GeV given the MSSM particle content (Fig. 4). On the other hand, the three gauge coupling constants miss each other quite badly with the non-supersymmetric Standard Model particle content. This observation suggests the possibility of supersymmetric grand unification.
### 6.2 Yukawa Coupling Constants
Since first- and second-generation Yukawa couplings are so small, let us ignore them and concentrate on the third-generation ones. Their renormalization-group equations are given as
$`\mu {\displaystyle \frac{dh_t}{d\mu }}`$ $`=`$ $`{\displaystyle \frac{h_t}{16\pi ^2}}\left[6h_t^2+h_b^2{\displaystyle \frac{16}{3}}g_3^23g_2^2{\displaystyle \frac{13}{15}}g_1^2\right],`$ (78)
$`\mu {\displaystyle \frac{dh_b}{d\mu }}`$ $`=`$ $`{\displaystyle \frac{h_b}{16\pi ^2}}\left[6h_b^2+h_t^2+h_\tau ^2{\displaystyle \frac{16}{3}}g_3^23g_2^2{\displaystyle \frac{7}{15}}g_1^2\right],`$ (79)
$`\mu {\displaystyle \frac{dh_\tau }{d\mu }}`$ $`=`$ $`{\displaystyle \frac{h_\tau }{16\pi ^2}}\left[4h_\tau ^2+3h_b^23g_2^2{\displaystyle \frac{9}{5}}g_1^2\right].`$ (80)
The important aspect of these equations is that the gauge coupling constants push down the Yukawa coupling constants at higher energies, while the Yukawa couplings push them up. This interplay, together with a large top Yukawa coupling, allows the possibility that the Yukawa couplings may also unify at the same energy scale where the gauge coupling constants appear to unify (Fig. 5). It turned out that the actual situation is much more relaxed than what this plot suggests. This is because there is a significant correction to $`m_b`$ at tanβ
>10
>𝛽10\tan\beta\mathrel{\mathchoice{\lower 3.0pt\vbox{\halign{$\m@th\displaystyle\hfill#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\m@th\textstyle\hfill#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\m@th\scriptstyle\hfill#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\m@th\scriptscriptstyle\hfill#\hfil$\cr>\crcr\sim\crcr}}}}10 when the superparticles are integrated out .
### 6.3 Soft Parameters
Since we do not know any of the soft parameters at this point, we cannot use the renormalization-group equations to probe physics at high energy scales. On the other hand, we can use the renormalization-group equations from boundary conditions at high energy scales suggested by models to obtain useful information on the “typical” superparticle mass spectrum.
First of all, the gaugino mass parameters have very simple behavior that
$$\mu \frac{d}{d\mu }\frac{M_i}{g_i^2}=0.$$
(81)
Therefore, the ratios $`M_i/g_i^2`$ are constants at all energies. If the grand unification is true, both the gauge coupling constants and the gaugino mass parameters must unify at the GUT-scale and hence the ratios are all the same at the GUT-scale. Since the ratios do not run, the ratios are all the same at any energy scales, and hence the low-energy gaugino mass ratios are predicted to be
$$M_1:M_2:M_3=g_1^2:g_2^2:g_3^21:2:7$$
(82)
at the TeV scale. We see the tendency that the colored particle (gluino in this case) is much heavier than uncolored particle (wino and bino in this case). This turns out to be a relatively model-independent conclusion.
The running of scalar masses is given by simple equations when all Yukawa couplings other than that of the top quark are neglected. We find
$`16\pi ^2\mu {\displaystyle \frac{d}{d\mu }}m_{H_u}^2`$ $`=`$ $`3X_t6g_2^2M_2^2{\displaystyle \frac{6}{5}}g_1^2M_1^2,`$ (83)
$`16\pi ^2\mu {\displaystyle \frac{d}{d\mu }}m_{H_d}^2`$ $`=`$ $`6g_2^2M_2^2{\displaystyle \frac{6}{5}}g_1^2M_1^2,`$ (84)
$`16\pi ^2\mu {\displaystyle \frac{d}{d\mu }}m_{Q_3}^2`$ $`=`$ $`X_t{\displaystyle \frac{32}{3}}g_3^2M_3^26g_2^2M_2^2{\displaystyle \frac{2}{15}}g_1^2M_1^2,`$ (85)
$`16\pi ^2\mu {\displaystyle \frac{d}{d\mu }}m_{U_3}^2`$ $`=`$ $`2X_t{\displaystyle \frac{32}{3}}g_3^2M_3^2{\displaystyle \frac{32}{15}}g_1^2M_1^2.`$ (86)
Here, $`X_t=2h_t^2(m_{H_u}^2+m_{Q_3}^2+m_{U_3}^2)`$ and the trilinear couplings are also neglected. Even within this simplified assumptions, one learns interesting lessons. First of all, the gauge interactions push the scalar masses up at lower energies due to the gaugino mass squared contributions. Colored particles are pushed up even more than uncolored ones, and the right-handed sleptons would be the least pushed up. On the other hand, Yukawa couplings push the scalar masses down at lower energies. The coefficients of $`X_t`$ in the Eqs. (83, 85, 86) are simply the multiplicity factors which correspond to 3 of $`SU(3)_C`$, 2 of $`SU(2)_Y`$ and 1 of $`U(1)_Y`$. It is extremely amusing that the $`m_{H_u}^2`$ is pushed down the most because of the factor of three as well as is pushed up the least because of the absence of the gluino mass contribution. Therefore, the fact that the Higgs mass squared is negative at the electroweak scale may well be just a simple consequence of the renormalization-group equations! Since the Higgs boson is just “one of them” in the MSSM, the renormalization-group equations provide a very compelling reason why it is only the Higgs boson whose mass-squared goes negative and condenses. One can view this as an explanation for the electroweak symmetry breaking.
### 6.4 Minimal Supergravity
Of course, nothing quantitative can be said unless one makes some specific assumptions for the boundary conditions of the renormalization-group equations. One common choice called “Minimal Supergravity” is the following set of assumptions:
$`m_Q^{2ij}=m_U^{2ij}=m_D^{2ij}=m_L^{2ij}=m_E^{2ij}=m_0^2\delta ^{ij},`$
$`m_{H_u}^2=m_{H_d}^2=m_0^2,`$
$`A_u^{ij}=A_d^{ij}=A_l^{ij}=A_0`$
$`M_1=M_2=M_3=M_{1/2}`$
at the GUT-scale. The parameter $`m_0`$ is called the universal scalar mass, $`A_0`$ the universal trilinear coupling, and $`M_{1/2}`$ the universal gaugino mass. Once this assumption is made, there are only five parameters at the GUT-scale, $`(m_0,M_{1/2},A_0,B,\mu )`$. This assumption also avoids most of the low-energy constraints easily because the scalar mass-squared matrices are proportional to the identity matrices and hence there is no flavor violation. Of course this is probably an oversimplification of the parameter space, but it still provides useful starting point in discussing phenomenology. Especially most of the search limits from collider experiments have been reported using this assumption. In general, this choice of the boundary conditions, which actually have not much to do with supergravity itself, lead to acceptable and interesting phenomenology including the collider signatures, low-energy constraints as well as cosmology.
## 7 Collider Phenomenology
We do not go into much details of the collider phenomenology of supersymmetry in this lecture notes and we refer to reviews. Here, we give only a very brief summary of collider phenomenology. Supersymmetry is an ideal target for current and new future collider searches. As long as they are within the mass scale expected by the argument given at the beginning of the lecture, we expect supersymmetric particles to be discovered at LEP-II (even though the phase space left is quite limited by now), Tevatron Run-II, or the LHC.
The next two figures Figs. 6, 7 show the discovery reach of supersymmetry at LEP-II, Tevatron Run II, LHC. It is fair to say that the mass range of superparticles relevant to solve the problem of fine cancellation in the Higgs boson self-energy described at the beginning of the lecture is covered by these experiments.
A future $`e^+e^{}`$ linear collider would play a fantastic role in proving that new particles are indeed superpartners of the known Standard Model particles and in determining their parameters. Once such studies will be done, we will exploit renormalization-group analyses trying to connect physics at TeV scale to yet-more-fundamental physics at higher energy scales. Example of such possible studies are shown in Fig. 9. The measurements of gaugino masses were simulated. At the LHC, the measurements are basically on the gluino mass and the LSP mass which is assumed to be the bino state, and their mass difference can be measured quite well. By assuming a value of the LSP mass, one can extract the gluino mass. At the $`e^+e^{}`$ linear colliders, one can even disentangle the mixing in neutralino and chargino states employing expected high beam polarizations and determine $`M_1`$ and $`M_2`$ in a model-independent matter. Combination of both types of experiments determine all three gaugino masses, which would provide a non-trivial test of the grand unification.
## 8 Mediation Mechanisms of Supersymmetry Breaking
One of the most important questions in the supersymmetry phenomenology is how supersymmetry is broken and how the particles in the MSSM learn the effect of supersymmetry breaking. The first one is the issue of dynamical supersymmetry breaking, and the second one is the issue of the “mediation” mechanism.
The problem of the supersymmetry breaking itself has gone through a dramatic progress in the last few years thanks to works on the dynamics of supersymmetric gauge theories by Seiberg. The original idea by Witten was that the dynamical supersymmetry breaking is ideal to explain the hierarchy. Because of the non-renormalization theorem, if supersymmetry is unbroken at the tree-level, it remains unbroken at all orders in perturbation theory. However, they may be non-perturbative effects suppressed by $`e^{8\pi ^2/g^2}`$ that could break supersymmetry. Then the energy scale of the supersymmetry breaking can be naturally suppressed exponentially compared to the energy scale of the fundamental theory (string?). Even though this idea attracted a lot of interest,<sup>13</sup><sup>13</sup>13I didn’t live through this era, so this is just a guess. the model building was hindered by the lack of understanding in dynamics of supersymmetric gauge theories. Only relative few models were convincingly shown to break supersymmetry dynamically, such as the $`SU(5)`$ model with two pairs of $`\mathrm{𝟓}^{}+\mathrm{𝟏𝟎}`$ and the 3-2 model. After Seiberg’s works, however, there has been an explosion in the number of models which break supersymmetry dynamically (see a review and references therein). For instance, some of the models which were claimed to break supersymmetry dynamically, such as $`SU(5)`$ with one pair of $`\mathrm{𝟓}^{}+\mathrm{𝟏𝟎}`$ or $`SO(10)`$ with one spinor $`\mathrm{𝟏𝟔}`$, are actually strongly coupled and could not be analyzed reliably (called “non-calculable”), but new techniques allowed us to analyze these strongly coupled models reliably. Unexpected vector-like models were also found which proved to be useful for model building.
There has also been an explosion in the number of mediation mechanisms proposed in the literature. The oldest mechanism is that in supergravity theories where interactions suppressed by the Planck scale are responsible for communicating the effects of supersymmetry breaking to the particles in the MSSM. For instance, see a review. Even though the gravity itself may not be the only effect for the mediation but there could be many operators suppressed by the Planck-scale responsible for the mediation, this mechanism was sometimes called “gravity-mediation.” The good thing about this mechanism is that this is almost always there. However we basically do not have any control over the Planck-scale physics and the resulting scalar masses-squared are in general highly non-universal. In this situation, the best idea is probably to constrain the scalar masses-squared matrix proportional to the identity matrix by non-Abelian flavor symmetries. Models were constructed where the breaking patterns of the flavor symmetry naturally explain the hierarchical quark and lepton mass matrices, while protecting the squark masses-squared matrices from deviating too far from the identity matrices.
A beautiful idea to guarantee the universal scalar masses is to use the MSSM gauge interactions for the mediation. Then the supersymmetry breaking effects are mediated to the particles in the MSSM in such a way that they do not distinguish particles in different generations (“flavor-blind”) because they only depend on the gauge quantum numbers of the particles. Such a model was regarded difficult to construct in the past. However, a break-through was made by Dine, Nelson and collaborators, who started constructing models where the MSSM gauge interactions could indeed mediate the supersymmetry breaking effects, inducing postive scalar masses-squared and large enough gaugino masses (which used to be one of the most difficult things to achieve). The original models had three independent sectors, one for supersymmetry breaking, one (the messenger sector) for mediation alone, and the last one the MSSM. Later models eliminated the messenger sector entirely (see also reviews).
Difficulty still remained how large enough gaugino masses can be generated in models where the sector of dynamical supersymmetry breaking couples to the MSSM fields only by Planck-scale suppressed interactions. One could go around this problem by a clever choice of the quantum numbers for a gauge singlet field. But it was not realized until recently that the gaugino masses are generated by superconformal anomaly. This observation was confirmed and further generalized by other groups. Randall and Sundrum further realized that one could even have scalar masses entirely from the superconformal anomaly if the sector of dynamical supersymmetry breaking and the MSSM particles are physically separated in the extra dimensions. The consequence was striking: the soft parameters were determined solely by the low-energy theory and did not depend on the physics at high energy scales at all. This makes it attractive as a solution to the problem of flavor-changing neutral currents, as the low-energy interactions of first and second generations are indeed nearly flavor-blind. Even though such models initially suffered from the problem that some of the scalars had negative mass-squared, simple fixes were proposed. One can preserve the virtue of the anomaly mediation, namely ultraviolet insensitivity, and construct realistic models.
Finally a new idea called “gaugino mediation” came out lately. This idea employs an extra dimension where the gauge fields propagate in the bulk. Supersymmetry is broken on a different brane and the MSSM fields learn the supersymmetry breaking effects by the MSSM gauge interactions, and hence solving the flavor-changing problem.
## 9 Conclusion
Supersymmetry is a well-motivated candidate for physics beyond the Standard Model. It would allow us to extrapolate the (supersymmetric version of the) Standard Model down to much shorter distances, giving us hope to connect the observables at TeV-scale experiments to parameters of the much more fundamental theories. Even though it has been extensively studied over two decades, many new aspects of supersymmetry have been uncovered in the last few years. We expect that research along this direction will continue to be fruitful. We, however, really need a clear-cut confirmation (or falsification) experimentally. The good news is that we expect it to be discovered, if nature did choose this direction, at the currently planned experiments.
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# Nonexistence of conformally flat slices of the Kerr spacetime
## I Introduction
Perhaps the most exciting source that might be detected by gravitational wave detectors now in development is radiation from a merger of black holes. This has been one of the motivations for the effort being put into the application of numerical relativity to black hole collisions. In this work supercomputers are used to evolve initial value solutions of Einstein’s field equations. The computation of the initial value solutions is itself a difficult task, and much of the work has taken advantage of the Bowen-York program for initial value solutions, in which the restriction is made that the initial 3-geometry is conformally flat.
Despite the elegance and convenience of the Bowen-York approach, a conformally flat initial solution has a serious shortcoming for work with black holes. Astrophysically realistic black holes will be rapidly rotating. The spatial geometry of the Kerr spacetime of rotating holes is not conformally flat. More specifically, for Kerr spacetime described in standard Boyer-Lindquist coordinates $`t,r,\theta ,\varphi `$, a slice at constant $`t`$ is not conformally flat. Because of this, the use of Bowen-York initial data to study colliding holes entails two difficulties. First, the Bowen-York representation of a rotating hole will be that of a distorted Kerr hole. Numerical evolution of this solution will produce a burst of radiation as each of the colliding holes “relaxes” to an approximately Kerr form. If it were possible to start the collision of the holes at large separation, this initial burst would be easily distinguishable from the radiation arising from the merger itself. But for the present, numerical evolutions must start quite close to the final state of the merger. The second difficulty is that the Bowen-York program cannot give an initial value solution that is a perturbation of the final single stationary hole, since that final solution does not have conformally flat spatial slices. This precludes the application of “close limit” perturbation theory that has proven to be very successful as an approximation scheme for black hole mergers.
It would be extremely useful if the need to deal with Kerr black holes could be reconciled with the convenience of the Bowen-York conformally flat scheme. For that reason, we inquire here whether there might exist conformally flat slices of the Kerr spacetime. We know that the $`t=`$ constant slices are not flat, but it may be that the geometry of a different kind of slicing, of the form $`t=f(r,\theta ,\varphi )`$ is conformally flat. The question of conformal flatness of any 3-geometry can most conveniently be answered with the use of the the Bach (Cotton-York) tensor,
$`B^{ij}`$ $`=`$ $`2ϵ^{ikp}[R_k^j{\displaystyle \frac{1}{4}}\delta _k^jR]_{;p},`$ (1)
where $`R_k^j`$ is the Ricci tensor for the 3-geometry, $`RR_k^k`$ is its Ricci scalar, and the semicolon refers to covariant differentiation with respect to the metric of the 3-geometry. The vanishing of the Bach tensor is a necessary and sufficient condition for a 3-geometry to be conformally flat.
Using the Bach tensor, we investigate here whether a conformally flat slice of the Kerr spacetime might exist, but we make certain restrictions on the search that limit the generality of our conclusion. One restriction is that we consider only axisymmetric slices of the axisymmetric Kerr geometry. One reason for this is practical: Gaining the advantages of conformal flatness while losing axisymmetry would be a Pyrrhic victory. A second reason is that the extension of our conclusion to a nonaxisymmetric slicing turns out to be quite difficult. We restrict ourselves, therefore, to slices of the form $`t=f(a,r,\theta )`$, where $`a`$ is the Kerr spin parameter.
A second restriction we make is to consider only families of slicings that have the property that in the Schwarzschild ($`a0`$) limit, the slicings smoothly go to slices of constant Schwarzschild time. This means that the slicing function $`f(a,r,\theta )`$, would have the limit zero as $`a0`$, since the Boyer-Lindquist coordinates become the Schwarzschild coordinates as ($`a0`$). We assume, furthermore, that $`f`$ can be expanded in $`a`$ so that
$$t=aF(r,\theta )+𝒪(a^2).$$
(2)
Our approach is to compute the Bach tensor for the 3-geometry induced by the slicing in eq. (2) and to expand the Bach tensor in $`a`$. We will show that no slicing can be found that makes the tensor vanish to lowest nontrivial order in $`a`$, and we conclude that no family of slices of this type can be conformally flat.
The assumption in eq. (2) means that our conclusion does not rule out a conformally flat slicing for an isolated value of $`a`$, or a family of conformally flat slicings for a range of $`a`$ that does not include $`a=0`$. These limitations are inherent to our method of expanding about $`a=0`$. A more subtle shortcoming of our method is that it does not rule out a family of slicings of the form
$$t=G(r,\theta )+aF(r,\theta )+𝒪(a^2).$$
(3)
where $`G`$ is not a constant. That is, it does not rule out a family of slicings that, in the Schwarzschild ($`a0`$) limit takes the form $`t=G(r,\theta )`$. But there are slicings of Schwarzschild other than $`t=`$ constant that are conformally flat. In fact any spherically symmetric 3 geometry is conformally flat, so any slicing of the form $`t=G(r)`$ is conformally flat. Such slicings of Schwarzschild, of course, are not orthogonal to the timelike Killing vector for the spacetime geometry. The use of such a slicing would have disadvantages for initial data similar to the disadvantage of a nonaxisymmetric slicing. For the Kerr spacetime of course, the Killing vector is not hypersurface orthogonal, so there is no foliation that is singled out by the symmetry. Still, intuition suggests that a useful slicing for Kerr should have extrinsic curvature that only describes the rotation, and that vanishes as $`a0`$. This is equivalent to a slicing that reduces to one of constant Schwarzschild time.
## II Kerr metric tensor
In Boyer-Lindquist coordinates, the components of the Kerr metric have the explicit form
$`{}_{}{}^{(4)}g_{tt}^{}`$ $`=`$ $`1+2Mr/(r^2+a^2\mathrm{cos}^2\theta )`$ (4)
$`{}_{}{}^{(4)}g_{t\varphi }^{}`$ $`=`$ $`2Mra\mathrm{sin}^2\theta /(r^2+a^2\mathrm{cos}^2\theta )`$ (5)
$`{}_{}{}^{(4)}g_{rr}^{}`$ $`=`$ $`(r^2+a^2\mathrm{cos}^2\theta )/(r^22Mr+a^2)`$ (6)
$`{}_{}{}^{(4)}g_{\theta \theta }^{}`$ $`=`$ $`r^2+a^2\mathrm{cos}^2\theta `$ (7)
$`{}_{}{}^{(4)}g_{\varphi \varphi }^{}`$ $`=`$ $`[(r^2+a^2)^2(r^22Mr+a^2)a^2\mathrm{sin}^2\theta ]\mathrm{sin}^2\theta /(r^2+a^2\mathrm{cos}^2\theta ).`$ (8)
The 3-geometry $`{}_{}{}^{(3)}g_{ij}^{}`$ induced on a spatial slice given by eq. (2) has components
$`{}_{}{}^{(3)}g_{rr}^{}`$ $`=`$ $`{}_{}{}^{(4)}g_{rr}^{}+^{(4)}g_{tt}a^2\left(_rF\right)^2+𝒪(a^3)`$ (9)
$`{}_{}{}^{(3)}g_{\theta \theta }^{}`$ $`=`$ $`{}_{}{}^{(4)}g_{\theta \theta }^{}+^{(4)}g_{tt}a^2\left(_\theta F\right)^2+𝒪(a^3)`$ (10)
$`{}_{}{}^{(3)}g_{r\varphi }^{}`$ $`=`$ $`2^{(4)}g_{t\varphi }a\left(_rF\right)+𝒪(a^3)`$ (11)
$`{}_{}{}^{(3)}g_{\theta \varphi }^{}`$ $`=`$ $`2^{(4)}g_{t\varphi }a\left(_\theta F\right)+𝒪(a^3)`$ (12)
$`{}_{}{}^{(3)}g_{\varphi \varphi }^{}`$ $`=`$ $`{}_{}{}^{(4)}g_{\varphi \varphi }^{}.`$ (13)
From these expressions, and from the fact that $`g_{t\varphi }`$ is proportional to $`a`$, it follows that the metric functions $`{}_{}{}^{(3)}g_{ij}^{}`$ have no terms first order in $`a`$. The deviations of $`{}_{}{}^{(3)}g_{ij}^{}`$ from a $`t=`$ constant slice of the Schwarzschild spacetime are therefore of order $`a^2`$ and higher. Since the $`t=`$ constant Schwarzschild slice is conformally flat, and hence has a vanishing Bach tensor, the components of the Bach tensor for eqs. (9) – (13) vanish to first order in $`a`$.
## III Diagonal components of the Bach tensor
To second order in $`a`$, the diagonal components of the Bach tensor turn out to be given by
$`B^{rr}`$ $`=`$ $`6Ma^2(r2M)\left(r^5\sqrt{r^5\mathrm{sin}^2\theta /(r2M)}\right)^1`$ (15)
$`\times \left[3(\mathrm{sin}^2\theta \mathrm{cos}^2\theta )_\theta F(r,\theta )5\mathrm{cos}\theta \mathrm{sin}\theta _\theta ^2F(r,\theta )\mathrm{sin}^2\theta _\theta ^3F(r,\theta )\right]`$
$`B^{\theta \theta }`$ $`=`$ $`6Ma^2\sqrt{r^5/(r2M)}\left(r^{12}\mathrm{sin}\theta \right)^1`$ (20)
$`\times [(r^4+4M^2r^2\mathrm{cos}^2\theta +r^4\mathrm{cos}^2\theta 4Mr^3\mathrm{cos}^2\theta 4M^2r^2+4Mr^3)_\theta _r^2F(r,\theta )+`$
$`(8r^256M^2+11r^2\mathrm{cos}^2\theta 50Mr\mathrm{cos}^2\theta +56M^2\mathrm{cos}^2\theta +44Mr)_\theta F(r,\theta )+`$
$`(5r^326M^2r\mathrm{cos}^2\theta +23Mr^2\mathrm{cos}^2\theta +26M^2r5r^3\mathrm{cos}^2\theta 23Mr^2)_\theta _rF(r,\theta )+`$
$`r(r2M)\mathrm{cos}\theta \mathrm{sin}\theta _\theta ^2F(r,\theta )]`$
$`B^{\varphi \varphi }`$ $`=`$ $`6Ma^2\left(r^7\mathrm{sin}^2\theta \sqrt{r^5/(r2M)}\right)^1`$ (23)
$`\times [(5r28M)\mathrm{sin}\theta _\theta F(r,\theta )+4r\mathrm{cos}\theta _\theta ^2F(r,\theta )+(13M5r)r\mathrm{sin}\theta _\theta _rF(r,\theta )+`$
$`r^2(r2M)\mathrm{sin}\theta _\theta _r^2F(r,\theta )+r\mathrm{sin}\theta _\theta ^3F(r,\theta )]`$
These three diagonal components are not independent, but are related by the fact that the trace $`B_i^i`$ vanishes, as can easily be checked (to second order in $`a`$) for eqs. (15)–(23).
If $`B^{rr}`$ is to vanish, we have from eq. (15) that
$$0=_\theta \left[3\mathrm{cos}\theta \mathrm{sin}\theta _\theta F(r,\theta )+\mathrm{sin}^2\theta _\theta ^2F(r,\theta )\right].$$
(24)
This equation can be solved by three integrations with respect to $`\theta `$ to give the general solution
$`F(r,\theta )`$ $`=`$ $`{\displaystyle \frac{h(r)}{2\mathrm{sin}^2\theta }}+u(r)\left[{\displaystyle \frac{\mathrm{cos}\theta }{2\mathrm{sin}^2\theta }}+{\displaystyle \frac{1}{2}}\mathrm{ln}(\mathrm{tan}\left[\theta /2\right])\right]+v(r),`$ (25)
in which $`h(r)`$, $`u(r)`$ and $`v(r)`$ are the “constants” introduced in the three integration steps. When this form of $`F(r,\theta )`$ is put into the right hand side of eq. (20), the equation $`B^{\theta \theta }=0`$ takes the form
$$0=\widehat{O}_uu(r)+\mathrm{cos}\theta \widehat{O}_hh(r),$$
(26)
and the vanishing of $`\widehat{O}_hh(r)`$ and of $`\widehat{O}_uu(r)`$ yield the two differential equations
$`0`$ $`=`$ $`(9r^2+56M^246Mr)h(r)+(26M^2r+23Mr^25r^3)_rh(r)+(4Mr^3+r^4+4M^2r^2)_r^2h(r)`$ (27)
$`0`$ $`=`$ $`(8r^2+44Mr56M^2)u(r)+(26M^2r+5r^323Mr^2)_ru(r)+(4M^2r^2r^4+4Mr^3)_r^2u(r).`$ (28)
Finally, the solutions to these equations are
$`h(r)`$ $`=`$ $`A{\displaystyle \frac{r^{7/2}}{\sqrt{r2M}}}+B{\displaystyle \frac{r^{7/2}\mathrm{cosh}^1(\sqrt{r/2M})}{\sqrt{r2M}}}`$ (29)
$`u(r)`$ $`=`$ $`C{\displaystyle \frac{(rM)r^{7/2}}{M\sqrt{r2M}}}+D{\displaystyle \frac{r^4}{M}}.`$ (30)
At this point we have reduced the freedom in $`F(r,\theta )`$ to the constants $`A`$,$`B`$,$`C`$,$`D`$ and the function $`v(r)`$, if the Bach tensor is to vanish.
## IV Offdiagonal components of the Bach tensor
When the form of $`F(r,\theta )`$ required by eqs. (25), (29), and (30) is used, we find that the component $`B^{r\theta }`$ is
$$B^{r\theta }=3a^2BMr^{9/2}\sqrt{r2M}/\mathrm{sin}\theta .$$
(31)
It follows that a conformally flat slicing requires that $`B=0`$.
With $`B`$ set to zero in the form of $`F(r,\theta )`$ required by eqs. (25), (29), and (30), the rather lengthy expression for $`B^{\theta \varphi }`$ can be written as
$$B^{\theta \varphi }=a^2\mathrm{\Xi }\left(\mathrm{sin}^{15}\theta r^{12}(r2M)^3M^2\right)^1,$$
(32)
where $`\mathrm{\Xi }`$ is
$`\mathrm{\Xi }`$ $`=`$ $`c_1(r)\mathrm{cos}\theta +c_2(r)\mathrm{cos}^2\theta +\mathrm{}+c_{16}(r)\mathrm{cos}^{16}\theta `$ (34)
$`+\left[b_1(r)\mathrm{cos}\theta +b_2(r)\mathrm{cos}^2\theta +\mathrm{}+b_{12}(r)\mathrm{cos}^{12}\theta \right]\mathrm{ln}\left(\mathrm{sin}\theta /[1+\mathrm{cos}\theta ]\right)`$
For $`B^{\theta \varphi }`$ to vanish, each of the terms $`c_1(r)\mathrm{}c_{16}(r),b_1(r)\mathrm{}b_{12}(r)`$ must vanish. The sum of $`b_k`$ terms is explicitly
$`{\displaystyle \underset{k}{}}b_k(r)\mathrm{cos}^k\theta `$ $`=`$ $`(2)(1\mathrm{cos}^2\theta )^5P_b(r,\theta ),`$ (35)
with
$`P_b`$ $`=`$ $`[(30Mr^{15}246M^2r^{14}+777M^3r^{13}1158M^4r^{12}+780M^5r^{11}168M^6r^{10})\mathrm{cos}^2\theta `$ (42)
$`+10Mr^{15}82M^2r^{14}+259M^3r^{13}386M^4r^{12}+260M^5r^{11}56M^6r^{10}]DA`$
$`[(30Mr^{29/2}+216M^2r^{27/2}576M^3r^{25/2}+672M^4r^{23/2}288M^5r^{21/2})\mathrm{cos}^2\theta `$
$`10Mr^{29/2}+72M^2r^{27/2}192M^3r^{25/2}+224M^4r^{23/2}96M^5r^{21/2}]\sqrt{r2M}CA`$
$`+[48r^{31/2}+388Mr^{29/2}1208M^2r^{27/2}+1776M^3r^{25/2}1184M^4r^{23/2}+256M^5r^{21/2}]\mathrm{cos}\theta \sqrt{r2M}D^2`$
$`+[48r^{31/2}+388Mr^{29/2}1224M^2r^{27/2}+1872M^3r^{25/2}1376M^4r^{23/2}+384M^5r^{21/2}]\mathrm{cos}\theta \sqrt{r2M}C^2`$
$`[96r^{16}872Mr^{15}+3160M^2r^{14}5740M^3r^{13}+5320M^4r^{12}2192M^5r^{11}+224M^6r^{10}]\mathrm{cos}\theta DC.`$
From the terms in $`P_b`$ that are proportional to $`\mathrm{cos}\theta `$ it follows that $`C`$ and $`D`$ must vanish if $`B^{\theta \varphi }`$ vanishes.
When the simplification $`C=D=0`$ is made in eq. (34) the sum of the $`c_k`$ terms takes the form
$`{\displaystyle \underset{k}{}}c_k(r)`$ $`\mathrm{cos}^k\theta =84M^4(r2M)^{7/2}r^{5/2}(1\mathrm{cos}^2\theta )^8`$ (47)
$`r^7M^2(r2M)^3(3\mathrm{cos}^2\theta +1)(1\mathrm{cos}^2\theta )^5`$
$`\times \left([r^44r^3M+4r^2M^2]_r^3v(r)+[2r^32r^2M4rM^2]_r^2v(r)+[2r^2+8rM5M^2]_rv(r)\right)A`$
$`+2M^2r^{19/2}(r2M)^{7/2}(1\mathrm{cos}^2\theta )^3`$
$`\times [(21M5r)\mathrm{cos}^4\theta 21M\mathrm{cos}^2\theta 3r]A^2.`$
In this expression the terms proportional to $`\mathrm{cos}^{11}\theta \mathrm{}`$$`\mathrm{cos}^{16}\theta `$ do not vanish for any choice of $`A`$ or $`v(r)`$. It follows that there is no function $`F(r,\theta )`$ for which the slicing in eq. (2) gives a 3-geometry that is conformally flat to second order in $`a`$. Under the assumptions stated at the outset this implies that there is no slicing of the Kerr spacetime that is conformally flat.
## V Discussion
We have shown that there can be no spatial slicing of the Kerr spacetime with the following properties: (i) The slicing is conformally flat. (ii) The slicing is axisymmetric. (iii) The slicing, as a function of the Kerr parameter $`a`$, goes smoothly to a slice of constant Schwarzschild time as $`a0`$ and the spacetime approaches the Schwarzschild spacetime. It follows that the Bowen-York method of generating initial value solutions, when applied to configurations with initial Kerr holes, or for the close-limit approximation method, entails the difficulties outlined in Sec. I.
It is natural to ask, aside from the “practical” questions related to black hole collisions, whether any of the restrictions in our conclusion can be modified or removed, so that a more general conclusion can be stated about conformally flat slicings of the Kerr spacetime. The axisymmetric restriction might appear to be particularly simple to remove, since we use what amounts to a perturbation (in $`a`$) expansion in which the “background” (the Kerr spacetime) is axisymmetric. But some of the terms in the Bach tensor are second order in the slicing funtion $`F`$, so a Fourier decomposition of $`F(r,\theta ,\varphi )`$ will result in a mixing of Fourier modes. A solution without a Fourier decomposition would appear to be quite difficult. The key to the relatively simple result we have presented is that the diagonal components of the Bach tensor, to second order in $`a`$, are linear in $`F`$ for axisymmetric slicings. That simplicity disappears in we allow for $`\varphi `$ dependence in $`F`$, and greatly complicates an approach of the type we have used. Since this is a difficult task and has little connection with questions of initial data sets, we have not pursued it.
An attempt to look at slicings with the form in eq. (3) runs into different problems. In this case, for nontrivial $`G`$, the Bach tensor will have terms to first order in $`a`$. These terms will be linear in both $`G`$ and in $`F`$. Choices of $`G`$ and $`F`$ for which the slicing is conformally flat to first order in $`a`$ cannot be ruled out. (If, for example, one chooses both $`F`$ and $`G`$ to be functions only of $`r`$, the slicing to first order in $`a`$ is spherically symmetric and hence conformally flat). It follows that from the first order Bach terms one can only infer restrictions on $`G`$ and on $`F`$. These restrictions must be applied to the equations that arise from the terms to second order in $`a`$, equations that include terms quadratic in $`F`$. Again, we have not pursued a generalization along these lines.
###### Acknowledgements.
We thank Walter Landry, William Krivan and Karel Kuchař for useful discussions. We thank James Bardeen for bringing to our attention the possibility that the Schwarzschild geometry may have slices that are conformally flat other than those of constant Schwarzschild time. Special thanks go to John Whelan for helpful suggestions at the early stages of this work. We gratefully acknowledge the support of the National Science Foundation under grant PHY9734871.
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# Very High-Energy Gamma-Ray Observations of PSR B1509-58 with the CANGAROO 3.8m Telescope
## 1 Introduction
Pulsar nebulae have been suggested as a possible acceleration site of high-energy particles in the galaxy (Harding (1990)). The first order Fermi acceleration mechanism is expected to occur in a shock between the pulsar wind and supernova ejecta, or interstellar matter. Evidence of such energetic phenomena has been obtained through observation of synchrotron emission by accelerated electrons and positrons at radio to gamma-ray ($``$ 10 GeV) energies. However, more direct evidence has become obtainable through Very High-Energy (VHE) gamma-ray ($``$ 300 GeV) observations over the last decade using the Imaging Atmospheric Čerenkov Technique (IACT).
VHE gamma-ray emissions from the directions of three energetic pulsars, the Crab (Weekes et al. (1989); Vacanti et al. (1991); Tanimori et al. (1994)); the Vela pulsar (Yoshikoshi et al. (1997)) and PSR B1706$``$44 (Kifune et al. (1995), Chadwick et al. (1997)), have been detected by ground-based telescopes using the IACT. Although all three pulsars show pulsed emission in the EGRET energy range (100MeV–10GeV), none of the VHE gamma-ray detections have shown any periodicity at the radio pulsar period. This steady VHE gamma-ray emission is usually explained to be a result of the inverse Compton scattering in the pulsar nebula, and not from the pulsar magnetosphere. While the mechanism of the emission from the Crab nebula is well studied (see, for example, de Jager et al. (1996)), information on other pulsars is still sparse. In order to study pulsars and their surrounding environment as possible acceleration sites of the cosmic rays, more examples in the VHE gamma-ray range are required.
PSR B1509$``$58 was discovered as an X-ray pulsar by Seward and Harnden (1982) using the Einstein X-ray Observatory. It is near the center of the supernova remnant MSH15$``$52 (G320.4$``$1.2). Soon after this discovery, pulsed radio emission was found by Manchester, Tuohy and D’Amico (1982). The pulsar has a period of 150 msec and a period derivative of 1.5 $`\times `$10<sup>-12</sup> ss<sup>-1</sup>, the largest known today. The characteristic age of the pulsar is estimated to be $``$1700 years (Manchester et al. (1998)), which makes it the second youngest pulsar after the Crab. <sup>1</sup><sup>1</sup>affiliation: Torii et al. (1997) have reported the discovery of a pulsar 1600 years old. This age is somewhat speculative however as the period derivative of the pulsar has not yet been measured and an association with a historical supernova was assumed to estimate the pulsar age. From the period and the large period derivative, a very strong surface magnetic field of 1.5 $`\times `$ 10<sup>13</sup> G and a large spin down energy loss rate of 1.8 $`\times `$ 10 <sup>37</sup> ergs s<sup>-1</sup> are implied. While the distance to the pulsar is relatively large (4.4 kpc, Taylor et al. (1995)), the expected energy flux received at the Earth is the fifth largest among the known pulsars.
A compact ($`\mathrm{\hspace{0.17em}10}^{}\times \mathrm{\hspace{0.17em}6}^{}`$) synchrotron X-ray nebula has been found to exist around PSR B1509$``$58 (Seward et al. (1984)). The synchrotron emission suggests the existence of non-thermal electrons (positrons) in the nebula, which will also emit VHE gamma-rays via inverse Compton scattering. A detectable VHE gamma-ray flux from this synchrotron nebula was predicted by du Plessis et al. (1995) as a function of the magnetic field strength in the nebula. The expected gamma-ray flux above 1 TeV of 10<sup>-11</sup> to 10$`{}_{}{}^{12}cm^2s^1`$ for nebula magnetic fields 4 to 10 $`\mu `$G is within the sensitivity of the CANGAROO 3.8m telescope. Thus, VHE observations should give a good measurement of the magnetic field strength of this nebula. Du Plessis et al. (1995) also predicted a very hard differential spectral index of $``$1.8 based on the X-ray observations. This prediction provides us with an extreme example of the utility of multiwavelength studies of synchrotron—inverse- Compton emitting objects. Besides the compact nebula, recent X-ray satellite observations suggest various non-thermal phenomena in this remnant. ROSAT observations indicate a non-thermal X-ray component from the central diffuse nebula (CDN) extending to a diameter of 50 ($``$ 60pc) centered on the pulsar (Trussoni et al. (1996)). ASCA observations revealed a non-thermal jet structure between the pulsar and the center of a thermal nebula about 10 north from the pulsar (Tamura et al. (1996)). In order to explain the effective thermalization process of the thermal nebula, Tamura et al. (1996) indicate the existence of accelerated ions as well as electrons in the jet. Furthermore, Gaensler et al. (1998) found synchrotron emission from compact knots in this thermal nebula from 20cm imaging observations with the Australia Telescope Compact Array.
The surface magnetic field strength of the pulsar PSR B1509$``$58 is estimated to be one of the largest among known pulsars. Due to the photon splitting process caused by this strong surface magnetic field, a cut-off in the pulsed emission around MeV energies is predicted by Harding, Baring and Gonthier (1997). In fact, Kuiper et al. (1999) have suggested that a cut-off around 10 MeV exists in the COMPTEL data. EGRET observations have resulted in only an upper limit for the pulsed emission from PSR B1509$``$58 (Thompson et al. (1994)). In contrast, Nel et al. (1992) have reported the detection of transient pulsed VHE gamma-rays from the observations between 1985 and 1988 based on ground–based (non-imaging) Čerenkov telescope observations. However they could not detect any significant pulsed emission in the successive years. They tried to explain their observations with the framework of the outer gap model (Cheng, Ho and Ruderman (1986)). (Bowden et al. (1993) reported a upper limit of the pulsed emission above 0.35 TeV from their observations in 1987 and 1989. Combining with the detection by Nel et al. (1992) in 1987 above 1.5 TeV, power law index of the integral energy spectrum is limited to be harder than $``$ 1.) Interestingly, Kuiper et al. (1999) also indicate a marginal detection of the pulsed emission above 10 MeV, where the origin may differ from that at lower energies. Consequently, we have examined our data for the presence of periodicity as well. Our observations are the first results on this pulsar with using the IACT, which is one order of magnitude more sensitive than non-imaging observations.
For the reasons given above, we believed that PSR B1509$``$58 would be an interesting object to study above 1 TeV energies with the CANGAROO 3.8m IACT telescope in both the steady nebula emission and the pulsed emission. Details of those observations are given in Section-2. The methods of the analysis and results are shown in Section-3. In Section-4, we summarize our results and discuss their implications.
## 2 Observations
The CANGAROO (Collaboration between Australia and Nippon (Japan) for a GAmma-Ray Observatory in the Outback) 3.8m telescope is located at Woomera, South Australia (13647E, 316S and 160m a.s.l.). Čerenkov photons emitted from extensive air showers originated by primary gamma-rays and cosmic rays are collected with a parabolic mirror of 3.8m diameter and detected with an imaging camera at the focal plane. The camera consists of 256 photomultiplier tubes (PMTs) of 10mm$`\times `$10mm size (Hamamatsu R2248). The PMTs are located in a 16$`\times `$16 square grid and the field of view amounts to 3$`{}_{}{}^{}\times `$ 3. When signals from more than 5 tubes exceed 3 photoelectrons each within a gate, a trigger is generated. The amplitude and relative time of each PMT signal, the event time, and the counting rate of each tube are recorded for each event. The absolute time can be obtained with a precision of 200 nsec using a GPS clock. In addition to the GPS clock, the time of a crystal clock with a precision of 100$`\mu `$sec is also recorded. The GPS clock was not available in the 1997 observations due to the installation work of our new data acquisition system. However, because the time indicated by the crystal clock shows a stable drift from that of the GPS clock, we can obtain accurate relative arrival times for events even without the GPS clock. The crystal clock is reset every observation (new moon) period. Therefore, a periodicity analysis based on this clock is valid on a month by month basis. GPS timing was restored in July 1997. Details of the camera and the telescope are described in Hara et al. (1993).
The telescope was pointed in the direction of the pulsar PSR B1509$``$58 (right ascension 15<sup>h</sup>13<sup>m</sup>55<sup>s</sup>.62 and declination –59 08 08<sup>′′</sup>.9 (J2000), Taylor et al. (1995)) in May and June in 1996, from March to May in 1997 and from March to May in 1998. The pulsar (ON source) and an offset region (OFF source), having the same declination as the pulsar but different right ascension, were observed for equal amounts of time each night under moonless and usually clear sky conditions. Typically, the ON source region is observed only once in a night around transit for a few hours. Two OFF source runs are carried out before and after the ON source run. The former one covers the first half of the ON source track and the latter covers the second half. In the off-line analysis, those data obtained when a small patch of cloud was obscuring the source are omitted. At the same time, the corresponding ON (or OFF) source data were also rejected from the analysis. In addition to the weather selection, the data taken when the electronics noise produced an anomalously large trigger rate were not used in the analysis. This happened in the 1996 observations. In the 1998 data, there are many nights which have a large difference of the event rate between the ON and OFF source regions, which is thought to be due to the presence of thin dew on the reflecting mirror. Data taken under these conditions were also omitted. The durations of selected observations after these procedures are $`26^\mathrm{h}30^\mathrm{m}`$ , $`32^\mathrm{h}08^\mathrm{m}`$ and $`21^\mathrm{h}14^\mathrm{m}`$ for the 1996, 1997 and 1998 (both ON and OFF) data, respectively. These data are used for the analysis in this paper.
Observations were carried out under different instrumental conditions in each year. During the 1996 observations, the reflectivity of the mirror was estimated to be $``$ 45%. We recoated the mirror in October 1996 by vacuum evaporation of aluminium at the Anglo Australian Observatory. As a result, the reflectivity of the mirror increased to about 90%. As the reflectivity was improved, the threshold energy of our telescope was lowered. For the 1997 observations, the threshold energy, defined here as the energy at which a differential photon flux with an assumed differential spectral index of 2.5 is maximized in the Monte Carlo calculations, was estimated to be 1.9 TeV, compared to 4.5 TeV before the recoating. By the 1998 observations, the reflectivity had decreased to $``$ 70%, corresponding to a threshold energy of 2.5 TeV. In these estimations, the selection effect of the analysis described in the next section is also taken into account. In the Monte Carlo calculation, we assumed that the observations were made at a zenith angle of 30, which was close to the average value for our observations on PSR B1509$``$58.
The observation times and threshold energies are summarized in Table 1, and as well, the analysis results are shown.
## 3 Analysis and Results
### 3.1 Analysis method
At the beginning of each run, the ADC pedestal and gain for each PMT were measured. To calibrate the gain, a blue LED located at the center of the mirror is used to illuminate the PMTs uniformly. The pedestal value is subtracted from the ADC value and any variations in the PMT gains were normalized using the LED calibration data. PMTs whose TDC value corresponded to a pulse arrival time within $`\pm `$ 30 nsec of the shower plane were regarded as ‘hit ’ tubes and used to calculate image parameters. After omitting some hit tubes which were isolated or which had ADC values less than one standard deviation above the pedestal value, the conventional image parameters (Hillas (1985)) were calculated. (In the 1996 data, a fifth of the PMTs at the bottom in the camera were omitted from analysis to avoid the effect of electronics noise. This makes the threshold energy higher and the effective area smaller. This effect is included in calculating the threshold energy and the flux upper limit.)
The parameter ranges determined from Monte Carlo simulations to optimize the gamma-ray signals are $`:`$ $`0^{}.60<`$ distance $`1^{}.30`$, $`0^{}.04<`$ width $`0^{}.09`$, $`0^{}.10<`$ length $`0^{}.40`$, $`0.35<`$ concentration $`0.70`$ and $`\alpha 10^{}`$. These ranges are slightly narrower than those used in case of the Vela analysis (Yoshikoshi et al. (1997)). The upper limit of $`\alpha `$, $`10^{}`$, is adopted assuming the source is a point-like. Two orientation parameters, $`\alpha `$ and distance, are defined with respect to the assumed source position in the field of view. In this paper, this is fixed at the pulsar position except in the spatial analysis discussed in Section-3.4. To avoid the effect of incomplete images near the edge of the camera, images with centroids located at greater than $`1^{}.05`$ from the center of the camera were also rejected. We also required that the number of hit tubes ($`N_{hit}`$) must be $``$ 5 and the total number of photo-electrons contained in an image ($`N_{p.e.}`$) must be $`40`$ to be able to obtain good image parameters and select only air shower induced events. The upper limit of $`N_{p.e.}`$ is large enough to accept all real events with large numbers of photo-electrons. In Table 1, the numbers of events in the raw data and selected are presented. We can find a large difference between ON and OFF in the raw data. The main reasons of the difference in number are the electronics noise in the 1996 data and the existence of the optically bright stars (M<sub>V</sub> = 4.1 and 4.5) in the field of view in the 1997 data, where the reflectivity of the mirror was the largest. However, the numbers match well after the selection of air shower events. For all the three years’ data analyses we applied the same criteria as described above.
### 3.2 Results of the image analysis
The distributions of the orientation angle ($`\alpha `$) after all other cuts were applied are shown in Figure Very High-Energy Gamma-Ray Observations of PSR B1509$``$58 with the CANGAROO 3.8m Telescope.
Although there was no statistically significant excess of the ON source counts over the OFF source seen in the 1996 data, the 1997 data shows an excess at $`\alpha 10^{}`$ with a statistical significance of 4.1$`\sigma `$. This excess may indicate the presence of a VHE gamma-ray signal from the source. The additional use of the asymmetry parameter showed an excess in the positive (gamma-ray–like) domain, though not at a level which would have increased the overall significance of the excess. More careful study would be necessary in use of this third-moment parameter for the source near the Galactic Center, where the night-sky background level is high. In the 1998 data, we find a small excess in the ON source counts, however, the statistical significance is only 1.4$`\sigma `$ at $`\alpha 10^{}`$. Hereafter, we regard the 1996 and 1998 results as non-detections of the VHE gamma-ray signal and treat the 1997 result as a marginal detection. The corresponding upper limits and flux are calculated as,
$`\mathrm{F}_{99.5\%}(\mathrm{E}\mathrm{\hspace{0.17em}4.5}\mathrm{TeV})\mathrm{\hspace{0.17em}1.9}\times 10^{12}\mathrm{cm}^2\mathrm{s}^1`$
$`\mathrm{F}(\mathrm{E}\mathrm{\hspace{0.17em}1.9}\mathrm{TeV})=(2.9\pm \mathrm{\hspace{0.17em}0.7})\times 10^{12}\mathrm{cm}^2\mathrm{s}^1`$
$`\mathrm{F}_{99.5\%}(\mathrm{E}\mathrm{\hspace{0.17em}2.5}\mathrm{TeV})\mathrm{\hspace{0.17em}2.0}\times 10^{12}\mathrm{cm}^2\mathrm{s}^1`$
for 1996, 1997 and 1998 results, respectively. Here, a differential energy spectral index of 2.5 is assumed. The upper limits and the errors in the flux are estimated based on the numbers of the observed counts. We note that in our calculation of the upper limits the difference of the counts between ON and OFF are also taken into account following the formula introduced by Helene (1983). So the 1998 flux upper limit becomes higher than that from completely null result. If we change the assumption of the differential energy spectral index over the range $`2.5\pm 1.0`$, the corresponding threshold energies are estimated to change by $``$30%. Instrumental uncertainties also affect the estimation of the threshold energies. We estimate the systematic error in determining the absolute threshold energies to be about 40$``$50%. However, because almost all of the systematic errors behave in the same way for the three years’ observations, the uncertainty of the relative threshold energy is smaller than this value.
### 3.3 Consistency and Stability
The positive indication is obtained only from the lowest threshold energy observation. But the derived flux and two flux upper limits require neither variability of the source nor a very soft spectral index, that is, the results from the three years are consistent with each other assuming stable emission with a Crab-like spectral index ($`2.5`$) or the harder index (1.8) expected by du Plessis et al. (1995). We also divided the 1997 data into separate new moon periods to check on consistency. The results are shown in Table 1. Each month’s result has a marginal positive effect on the final result. The excess counting rate is stable during the three observation seasons within the statistical errors.
### 3.4 Spatial analysis
PSR B1509$``$58 and its surrounding environment are complex and there are indications from X-ray data that non-thermal phenomena possibly occur over an extended area of this remnant. So it is possible that the gamma-ray–like signal in the 1997 data is not from a point source at the pulsar position but from some other region near the pulsar. Therefore we have carried out a source search in the $`2^{}\times 2^{}`$ field of view centered on the pulsar position. To do this, we shifted the position of the assumed source over a grid of points around the pulsar and repeated the analysis at each point to obtain the excess counts in the $`\alpha `$ distribution. The resultant map of the significance is shown in Figure Very High-Energy Gamma-Ray Observations of PSR B1509$``$58 with the CANGAROO 3.8m Telescope.
The peak of the excess is found at $`0^{}.1`$ south-west from the pulsar. But when we consider the degrees of freedom of the search, the significance at this maximum should be reduced. And also, from a Monte Carlo calculation, where the observed counts of signal and background are taken into account, we estimated that the precision to determine the source position is $`0^{}.10`$ at the 1$`\sigma `$ level. We conclude, therefore, that the position of the excess is consistent with the pulsar position within the statistics of our observations.
### 3.5 Periodicity analysis
The recorded arrival times of the gamma-ray–like signals ($`\alpha 10^{}`$ after all the image cuts) were converted to the Solar System Barycenter arrival times using the solar system ephemeris based on epoch 2000 (DE200) (Standish, 1982). We then carried out a phase analysis with the phase parameters summarized in Table 2 (Manchester et al. (1998)). Because Nel et al. (1990) pointed out a possibility of a light curve with triple peaks in the TeV energy range, we applied the H-test (de Jager, Swanepoel and Raubenheimer (1989)) to obtain the statistical significance. The virtue of the H-test lies in the fact that it requires no assumptions about bin size and bin location and is also independent of the shape of the light curve. The results are summarized in Table 3.
The results of 1997 are divided into separate observational periods (months), because GPS timing information was not available in 1997 as mentioned in Section-2. The relative arrival time of the events is calculated for the 1997 data from the time of the crystal clock, having a constant drift rate relative to the GPS clock. The H-statistics and the corresponding probabilities against a uniform distribution are shown in Table 3. No evidence for the 150 ms periodicity is found in any of the observation seasons. To calculate the flux upper limit for the pulsed emission, we used the formula given by de Jager (1994). This formula combines the observed counts (N) and pulsed fraction (p) through a parameter, $`\chi `$, as, $`\chi =p\sqrt{N}`$. When the H-statistic is considered as a non-detection of periodicity, $`\chi `$ giving 3 $`\sigma `$ upper limit of p is expressed as,
$$\chi _{3\sigma }=(1.5+10.7\delta )(0.174H)^{0.17+0.14\delta }exp\left[(0.08+0.15\delta )\left\{log_{10}(0.174H)\right\}^2\right]$$
Here, H is the value of the H-test as shown in Table 3. (For H$`<`$0.3 we should take H=0.3 in calculating $`\chi _{3\sigma }`$.) $`\delta `$ is the duty cycle of the pulse profile. In case of PSR B1509$``$58, we assumed $`\delta `$ to be 0.3 using the X-ray observation by Kawai et al. (1991). The 3$`\sigma `$ upper limits for the pulsed VHE gamma-ray emission are also shown in Table 3.
## 4 Discussion
Our observations can be summarized as follows $`:`$ (1) In the observations with the lowest detection threshold energy, a 4.1$`\sigma `$ excess of gamma-ray–like events is found. Null results in the observations of the other years (when the detection threshold energies were higher) are not in conflict with this marginal positive result $`:`$ neither variability of the source nor an especially soft energy spectrum needs to be invoked. (2) From the result in the 1997 observations, there is no evidence of a variability on a monthly time-scale during three observation seasons. (3) In the 1997 data, the peak emission source position is shifted slightly to the south-west direction from the pulsar position. However, considering the statistical error including the real event numbers observed, this is consistent with the pulsar position. (4) The periodicity of the events modulated with the radio pulsar period is studied. We found no evidence of the 150 ms pulsar periodicity using the H-test in any of the observations for three years.
The statistical significance of the 1997 excess, 4.1$`\sigma `$, is too small to claim as the detection of a VHE gamma-ray source, however, it is sufficiently suggestive to allow discussion supposing the excess was due to a VHE gamma-ray signal. With this scheme the simplest and most straightforward explanation can be made assuming that the emission is found from the pulsar nebula surrounding the pulsar. VHE gamma-ray emission from a pulsar nebula is usually considered as a result of inverse Compton scattering by relativistic electrons. From the emission processes of synchrotron and inverse Compton radiations, a simple equation, $`\frac{\dot{\mathrm{E}}_{\mathrm{synch}}}{\dot{\mathrm{E}}_{\mathrm{iC}}}=\frac{ϵ_\mathrm{B}}{ϵ_{\mathrm{ph}}}`$, can be obtained. Here $`\dot{\mathrm{E}}_{\mathrm{synch}}`$ and $`\dot{\mathrm{E}}_{\mathrm{iC}}`$ are the luminosities through synchrotron radiation (mainly resulting in quanta in the X-ray energy range) and inverse Compton scattering (mainly producing VHE gamma-rays), respectively, and $`ϵ_\mathrm{B}`$ and $`ϵ_{\mathrm{ph}}`$ are the energy densities of the magnetic field and the target photons for inverse Compton scattering at the emission region. Assuming isotropic emission of both X-rays and gamma-rays, $`\frac{\dot{\mathrm{E}}_{\mathrm{synch}}}{\dot{\mathrm{E}}_{\mathrm{iC}}}`$ can be equated to $`\frac{\mathrm{F}_{\mathrm{synch}}}{\mathrm{F}_{\mathrm{iC}}}`$. Here $`\mathrm{F}_{\mathrm{synch}}=7.2\times 10^{11}\mathrm{ergs}\mathrm{cm}^2\mathrm{s}^1(0.12.4\mathrm{keV})`$ as given by Trussoni et al. (1996) and $`\mathrm{F}_{\mathrm{iC}}=2.7\times 10^{11}\mathrm{ergs}\mathrm{cm}^2\mathrm{s}^1`$, obtained by integrating the 1997 flux above 1.9 TeV assuming a differential spectral index of 2.5. (The corresponding luminosity at the pulsar, $`L_{iC}`$, is $`6.2\times 10^{34}\mathrm{ergs}\mathrm{s}^1`$ assuming the pulsar distance of 4.4 kpc. That is 0.34% of the pulsar rotating energy loss.) If the 3 K Microwave Background Radiation (MBR) is the only target of the inverse Compton radiation, i.e., $`ϵ_{\mathrm{ph}}=ϵ_{3\mathrm{K}}=3.8\times 10^{13}\mathrm{ergs}\mathrm{cm}^3`$, one obtains $`ϵ_\mathrm{B}=1.0\times 10^{12}\mathrm{ergs}\mathrm{cm}^3`$. This, then, leads to a value for the magnetic field strength $`\mathrm{B}\mathrm{\hspace{0.17em}5}\mu \mathrm{G}`$. Considering the large uncertainties in the arguments above, this value agrees well with the previously estimated value of $`\mathrm{B}\mathrm{\hspace{0.17em}7}\mu \mathrm{G}`$, from the equipartition of energy between the particles and the magnetic field (Seward et al. (1984)). According to the prediction of du Plessis et al. (1995), our result corresponds to a magnetic field strength of $`\mathrm{B}5\mu \mathrm{G}`$. These three estimated values of the magnetic field agree very well with each other.
An alternative source of the target photons is the IR source IRAS 15099$``$5856, known to be positionally coincident with the pulsar (Arendt (1991)). Du Plessis et al. (1995) estimated that the contribution from the IR photons to the VHE gamma-ray flux would be at the same level as that from the 3 K MBR. However the association between IRAS 15099$``$5856 and the pulsar is uncertain. In case that the IRAS source found at 25$`\mu `$m supplies the target photon for the inverse Compton process, the resultant VHE gamma-ray spectrum is expected to be softer than that made from the 3 K MBR. This is because the critical energy of the parent electrons in the Klein-Nishina cross section is $`6\times 10^{12}`$ eV against 25$`\mu `$m IR radiation while it is $`10^{15}`$ eV for the 3 K MBR. Therefore, the VHE gamma-ray spectrum should have a rapid softening over the TeV energy range. To understand the association of this IRAS source, detailed spectral measurements with future observations are required as well as the X-ray observations discussed below.
While our observations do not place any interesting limit on the spectral index, the very hard spectrum predicted by du Plessis et al. (1995) should be discussed. Their prediction was based on the observational results of the X-ray spectrum which showed a hardening of the index in the energy range below a few keV (photon index 1.4$`{}_{0.2}{}^{}{}_{}{}^{+0.4}`$ below 4 keV while 2.15$`\pm `$0.02 between 2 keV and 60 keV). However, recent X-ray observations do not confirm this hardening. The photon indices obtained in the wide X-ray energy band are consistent with a value around 2.2 (Trussoni et al. (1996); Tamura (1997); Marsden et al. (1997)) though the error of the ROSAT result is large. To discuss the synchrotron spectrum in detail, we need information from radio observations. But, even with the recent high resolution observations, a radio pulsar wind nebula has not been discovered (Gaensler et al. (1998)).
The upper limits set to the periodic signal in this paper are one order of magnitude below the previously reported flux in the same energy band (Nel et al. (1992)). Although Nel et al. reported upper limits from observations after 1988, our results should provide a far stricter limit on models. The VHE pulsed emission is in conflict with the observed cut-off around 10 MeV as predicted by the polar-cap model. To explain the VHE pulsed emission, an additional hard component, probably outer-gap emission, is required. Future observations by GLAST may reveal the existence of this component and studies of its flux and spectral variability may hint at large variability in the VHE range. The flux of the transient VHE pulsed emission reported in 1985, $`\mathrm{F}(\mathrm{E}\mathrm{\hspace{0.17em}1.5}\mathrm{TeV})=(3.9\pm \mathrm{\hspace{0.17em}0.9})\times 10^{11}\mathrm{cm}^2\mathrm{s}^1`$, would make this source the brightest known VHE gamma-ray source in the southern hemisphere. We could detect this kind of activity even with short duration monitoring. Semi-simultaneous monitoring of this pulsar with the future large IACT arrays in the southern hemisphere (CANGAROO-III, HESS) and GLAST would be of great interest if the pulsar were to display such an active phase in the future.
Finally, it is notable that, unlike the other pulsar nebulae detected at VHE energies, PSR B1509$``$58 is not firmly detected by EGRET onboard the CGRO satellite. In contrast, this pulsar and its surroundings show a variety of the non-thermal phenomena as introduced in Section-1. A comparison of nonthermal X-ray emission with VHE gamma-ray emission is becoming very useful in the search for VHE gamma-ray sources and study of their environment. Combined with the recent studies of pulsar nebulae (Kawai and Tamura (1996)), the new generation of the Imaging Atmospheric Čerenkov Telescopes (e.g. Matsubara et al. (1997)) will result in an improved understanding of pulsar nebulae and particle acceleration. The CANGAROO II 7m telescope started observations at Woomera in mid-1999. From new observations with a lower energy threshold, we will be able to measure the gamma-ray spectrum precisely and obtain a better estimation of the physical parameters, especially the magnetic field strength, in pulsar nebulae.
This work is supported by a Grant-in-Aid in Scientific Research from the Japan Ministry of Education, Science, Sports and Culture, and also by the Australian Research Council. We would like to thank to Dr. R. N. Manchester who provided us the latest radio ephemeris data on the pulsar. We are grateful to the AAO staffs in the recoating work of the 3.8m mirror. The receipt of JSPS Research Fellowships (JH, AK, TN, GPR, KS, GJT and TY) is also acknowledged. Finally, we thank the anonymous referee whose comments helped us improve the manuscript.
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# 1. Introduction and statement of the main results
## 1. Introduction and statement of the main results
Let $`(M,\omega )`$ be a symplectic manifold of integer class, i.e. $`[\omega /2\pi ]H^2(M;)`$ lifts to an integer cohomology class. Such symplectic manifold has an associated line bundle $`L`$ with first Chern class $`c_1(L)=[\omega /2\pi ]`$, which is equipped with a connection $``$ of curvature $`i\omega `$.
In his outbreaking work \[Do96\] S. Donaldson proved the existence of symplectic submanifolds of $`M`$ that realize the Poincaré dual of a large enough integer multiple of $`[\omega /2\pi ]`$. These are constructed as zero sets of appropriate sections of $`L^k`$. This extends a classical result in Kähler geometry saying that $`L`$ is ample, so $`L^k`$ has holomorphic sections with smooth holomorphic, and so symplectic, zero sets.
Later on, D. Auroux and R. Paoletti have proved independently an extension of Donaldson’s theorem, where now more symplectic submanifolds are constructed as the zero sets of asymptotically holomorphic sections of vector bundles. These bundles are obtained by tensoring an arbitrary complex bundle with large powers of the canonical line bundle $`L`$ \[Au97\], \[Au99\], \[Pa99\]. In his paper, D. Auroux also shows that, asymptotically, all the sequences of submanifolds constructed from a given vector bundle $`E`$ are isotopic. (For a summary of these results see for example the review paper \[Do98\].)
The key idea to understand these works is the concept of ampleness of a complex holomorphic bundle. This concept allows the flexibilization of the bundles in the holomorphic category by means of increasing their curvatures. Donaldson \[Do96\] has translated the definition of ampleness to the symplectic category. For this he studies the asymptotical behaviour of sequences of sections of the bundles $`L^k`$. Similarly, the important point in our work is the definition of the concept of asymptotic holomorphicity for sequences of embeddings constructed from very ample linear systems defined over vector bundles more and more twisted.
The change to the non-integrable setting is controlled by this concept. To define it we need to fix a compatible almost complex structure $`J`$ in $`(M,\omega )`$. So the pair $`(\omega ,J)`$ gives a metric $`g`$ in the tangent bundle. We have a sequence of metrics $`g_k=kg`$ indexed by integers $`k1`$.
###### Definition 1.1.
Let $`X`$ be a Hogde manifold with complex structure $`J_0`$. Let $`\gamma >0`$. A sequence of embeddings $`\varphi _k:MX`$ is $`\gamma `$-asymptotically holomorphic if it verifies the following conditions:
1. $`d\varphi _k:T_xMT_{\varphi _k(x)}X`$ has a left inverse $`\theta _k`$ of norm less than $`\gamma ^1`$ at every point $`xM`$. (The norm is taken with respect to the metric $`g_k`$.)
2. $`|(\varphi _k)_{}JJ_0|_{g_k}=O(k^{1/2})`$ on the subspace $`(\varphi _k)_{}T_xM`$.
3. $`|^p\varphi _k|_{g_k}=O(1)`$ and $`|^{p1}\overline{}\varphi _k|_{g_k}=O(k^{1/2})`$, for all $`p1`$.
A sequence of embeddings is asymptotically holomorphic if there is some $`\gamma >0`$ such that it is $`\gamma `$-asymptotically holomorphic.
The first important result is a generalization to the symplectic category of the classical Kodaira’s embedding Theorem:
###### Theorem 1.2.
Given $`(M,\omega )`$ a closed symplectic $`2n`$-dimensional manifold of integer class endowed with a compatible almost complex structure, there exists an asymptotically holomorphic sequence of embeddings $`\varphi _k:M^{2n+1}`$ with $`\varphi _k^{}[\omega _{FS}]=[k\omega ]`$. Moreover, given two such sequences of embeddings asymptotically holomorphic with respect to two compatible almost complex structures, then they are isotopic for $`k`$ large enough.
A sharper, in a sense, result than this has been obtained independently by Bortwick and Uribe in \[BU99\] using completely different ideas. Their result also obtains control in the symplectic part (equivalently in the metric part) allowing to obtain asymptotically holomorphic embeddings which are also asymptotically symplectic. Their approach is based on ideas coming from Tian \[Ti90\] to solve the problem in the Kähler case.
Our main interest for proving Theorem 1.2 is given by the possibility of studying “projective symplectic geometry”. We mean by this the study of sequences of asymptotically holomorphic submanifolds, namely obtained as images of asymptotically holomorphic embeddings, in the projective space. The strength of this approach is shown in the following
###### Theorem 1.3.
Let $`\varphi _k`$ be an asymptotically holomorphic sequence of embeddings in $`^{2n+1}`$ with $`\varphi _k^{}[\omega _{FS}]=[k\omega ]`$ and let $`ϵ>0`$. Let us fix a holomorphic submanifold $`N`$ in $`^{2n+1}`$. Then there exists an asymptotically holomorphic sequence of embeddings $`\widehat{\varphi }_k`$, at distance at most $`ϵ`$ in $`C^r`$-norm from the initial sequence and verifying that $`\widehat{\varphi }_k(M)N`$ is symplectic for $`k`$ large enough.
With the notations introduced in Section 2 we will precise a little more the precedent result, assuring that $`M\widehat{\varphi }_k^1(N)`$ is a sequence of asymptotically holomorphic submanifolds.
We will see that this result will imply a projective version of the symplectic Bertini’s Theorem proved in \[Do96\]. But the constructive method could allow to find more general types of symplectic submanifolds. This is shown in a more general situation. For this we generalize Theorem 1.2 to the grassmannian case.
###### Theorem 1.4.
Let $`(M,\omega )`$ be a closed symplectic $`2n`$-dimensional manifold of integer class endowed with a compatible almost complex structure. Suppose also that we have a rank $`r`$ hermitian vector bundle with connection, and that $`N>n+r1`$ and $`r(Nr)>2n`$. Then there exist an asymptotically holomorphic sequence of embeddings $`\varphi _k:M\text{Gr}(r,N)`$ with $`\varphi _k^{}𝒰=EL^k`$, where $`𝒰\text{Gr}(r,N)`$ is the universal rank $`r`$ bundle over the grassmannian. Moreover, given two such sequences of embeddings asymptotically holomorphic with respect to two compatible almost complex structures, then they are isotopic for $`k`$ large enough.
In Section 5 we will take profit of this result to extend the construction of determinantal submanifolds to the symplectic category in the following way.
###### Definition 1.5.
Let $`M`$ be a differentiable manifold and let $`E,F`$ be complex vector bundles over $`M`$. Given a morphism of vector bundles $`\phi :EF`$, the $`r`$-determinantal set $`\mathrm{\Sigma }^r(\phi )`$ is defined as
$$\mathrm{\Sigma }^r(\phi )=\{xM|\text{rank}\phi _x=r\}.$$
In the smooth category we can find for any morphism $`\phi :EF`$, another morphism $`\widehat{\phi }:EF`$ arbitrarily close to $`\phi `$ in $`C^p`$-norm, such that $`\mathrm{\Sigma }^r(\widehat{\phi })`$ is a smooth submanifold in $`M`$ of codimension $`2(r_er)(r_fr)`$, where $`r_e`$ and $`r_f`$ are the ranks of $`E`$ and $`F`$, respectively (if this number is greater than the dimension of $`M`$ then the set is empty). There exists a similar result in the algebraic category if the vector bundle $`E^{}F`$ is very ample. Our objective will be to adapt the algebraic discussion to the symplectic category to prove
###### Theorem 1.6.
Let $`(M,\omega )`$ be a closed symplectic manifold of integer class. Let $`E`$ and $`F`$ be hermitian vector bundles of rank $`r_e`$ and $`r_f`$, respectively. Then, for $`k`$ large enough, there exists a morphism $`\phi _k:E(L^{})^kFL^k`$ verifying that
1. $`\mathrm{\Sigma }^r(\phi _k)`$ is an open symplectic submanifold of $`M`$.
2. $`\text{codim }\mathrm{\Sigma }^r(\phi _k)=2(r_er)(r_fr)`$. The set of manifolds $`\{\mathrm{\Sigma }^r(\phi _k)\}_r`$ constitutes a stratified submanifold, called determinantal submanifold.
Moreover, given two stratified determinantal submanifolds constructed following the process described in the proof then there exists an ambient isotopy making the $`r`$-determinantal submanifolds associated to each stratified submanifold isotopic.
Theorem 1.6 was the original motivation of this paper. The idea of studying this kind of submanifolds is inspired in algebraic geometry. Note that in algebraic geometry the manifolds constructed as zeroes of sections of vector bundles have many topological restrictions, namely they satisfy the Lefschetz hyperplane Theorem, their Chern classes are very special, etc. So the set of submanifolds of a given manifold constructed in this way is very special in the set of all the submanifolds. However the determinantal submanifolds are very generic in the set of submanifolds. For instance, every codimension $`2`$ submanifold of an algebraic manifold can be constructed as the determinantal degeneration loci of certain bundle homomorphism \[Vo78\].
An obvious guess is that in symplectic geometry things are similar. Recall that the most general submanifolds constructed using asymptotically holomorphic techniques, prior to Theorem 1.6 are the Auroux’ ones \[Au97\]. These are zeroes of sections of vector bundles, so its topological properties are very special. In fact, Auroux cannot easily assure that these submanifolds are different from the ones constructed by Donaldson in \[Do96\]. In Subsection 5.4 we compute some Chern numbers of determinantal submanifolds showing that they are clearly different from the Chern numbers of Auroux’ and Donaldson’s submanifolds. So the symplectic type, and even the topological type, of the constructed submanifolds is necessarily different. This shows that the class of determinantal submanifolds is far more general.
Remark that, in any case, all the precedents results are obtained by means of twisting vector bundles with large powers of the line bundle $`L`$. So the submanifolds constructed in this way are quite special. It would be desirable to avoid this restriction, but this generalization cannot be made with the Donaldson’s techniques developped in \[Do96\].
From a symplectic point of view determinantal submanifolds are also interesting. They constitute a step in the study of singular symplectic submanifolds following the program sketched by Gromov \[Gr86\]. Donaldson and Auroux have attacked this question in \[Do99\] and \[Au99\]. Donaldson studies the local symplecticity of the fibers of asymptotically holomorpic applications $`f:^n`$ at a neigborhood of a critical point, it is solved by a local perturbation argument. The conclusion of Donaldson’s work is that the topological behaviour of that kind of functions is similar to the holomorphic Morse functions. Auroux studies the local symplecticity of asymptotically holomorphic applications $`f:^2^2`$ at the neighboorhood of a critical point, showing that are topologically equivalent to one of the two generic models of a holomorphic application \[Ar82\]. From this point of view Theorem 1.6 can be considered, in part, an extension of these results to generic singularities.
The organization of the paper is as follows. In Section 2 we will give the basic ideas of the Donaldson-Auroux’ theory needed in our work and prove Theorem 1.2. In Section 3.2 we prove Theorem 1.3. For this we explain some euclidean notions concerning the estimation of angles between subspaces. In Section 4 we generalize all the discussion to the case of the grassmannian embbedings, proving Theorem 1.4. This allows us to prove Theorem 1.6 in Section 5 and to analyze the topological properties of the constructed submanifolds.
Acknowledgements. We want to acknowledge D. Auroux, S. Donaldson and R. Paoletti for his kindness communicating us their results. Also we thank A. Ibort and D. Martínez for a lot of interesting discussions. Second named author was conducting his research financed by the Ph.D. program of the Consejería de Educación de Madrid.
## 2. Asymptotically holomorphic embeddings in projective space
As in the introduction, let $`(M,\omega )`$ be a symplectic manifold of integer class with associated line bundle $`L`$ and a compatible almost complex structure $`J`$. In the Kähler setting this line bundle supports a holomorphic structure and it is ample in the algebraic geometry sense, i.e. $`L^k`$ has a lot of holomorphic sections. This allows to embed $`M`$ in the projective space $`^N`$, for some $`N`$. In this Section we shall extend this classical result to the symplectic case inspired in the ideas of \[Do96\], thus proving Theorem 1.2.
### 2.1. Asymptotically holomorphic sequences
In this Subsection we collect the relevant results of the asymptotically holomorphic theory, as stated by D. Auroux in \[Au99\], that we shall use extensively along this work.
###### Definition 2.1.
A sequence of sections $`s_k`$ of hermitian bundles $`E_k`$ with connections on $`M`$ is called asymptotically $`J`$-holomorphic if there exist constants $`(C_p)_p`$ such that, for all $`k`$, at every point of $`M`$, $`|s_k|C_0`$, $`|^ps_k|C_p`$ for all $`p0`$, and $`|^{p1}\overline{}s_k|C_pk^{1/2}`$ for all $`p1`$. The norms are evaluated with respect to the metrics $`g_k`$.
In Donaldson’s first work on the subject \[Do96\], $`E_k=L^k`$. In that work Donaldson imposed an additional condition of improved transversality to the sequence of sections to assure that its zero sets are symplectic submanifolds for $`k`$ large enough. This condition is stated as follows.
###### Definition 2.2.
A section $`s_k`$ of the line bundle $`L^k`$ is said to be $`\eta `$-transverse to $`0`$ if for every point $`xM`$ such that $`|s_k(x)|<\eta `$ then $`|s_k(x)|>\eta `$.
If we get an asymptotically $`J`$-holomorphic sequence $`s_k`$ of sections of $`L^k`$ verifying that all of them are $`\eta `$-transverse to $`0`$, with $`\eta >0`$ independent of $`k`$ then we can assure that $`|s_k(x)|>|\overline{}s_k(x)|`$ if $`x`$ is a zero of $`s_k`$, for $`k`$ large enough. A simple linear algebra argument assures that the zeroes of $`s_k`$ are symplectic submanifolds for $`k`$ large enough.
In \[Au97\] D. Auroux extended the notion of transversality to the case of higher rank bundles. Let $`E`$ be a rank $`r`$ hermitian bundle with connection.
###### Definition 2.3.
A section $`s_k`$ of the bundle $`EL^k`$ is $`\eta `$-transverse to $`0`$ if for every $`xM`$ such that $`|s_k(x)|<\eta `$ then $`s_k(x)`$ has a right inverse $`\theta _k`$ such that $`|\theta _k|<\eta ^1`$.
We name universal constant to a number which only depends on the manifold geometry and on the constants involved in the data given to start with, i.e. a number independent of $`k`$ and the point $`xM`$. Similarly a universal polynomial is a polynomial only depending on the geometry of the manifold and on the constants provided in the original data. Donaldson uses highly localized asymptotically holomorphic sections, verifying the following definition.
###### Definition 2.4.
A sequence of sections $`s_k`$ of hermitian bundles $`E_k`$ with connections has Gaussian decay in $`C^r`$-norm away from the point $`xM`$ if there exists a universal polynomial $`P`$ and a universal constant $`\lambda >0`$ such that for all $`yM`$, $`|s(y)|`$, $`|s(y)|_{g_k}`$, $`\mathrm{}`$, $`|^rs(y)|_{g_k}`$ are bounded by $`P(d_k(x,y))\mathrm{exp}(\lambda d_k(x,y))`$. Here $`d_k`$ is the distance associated to the metric $`g_k`$.
The starting point for Donaldson’s construction is the following existence Lemma.
###### Lemma 2.5 (\[Do96, Au97\]).
Given any point $`xM`$, for $`k`$ large enough, there exist asymptotically holomorphic sections $`s_{k,x}^{\text{ref}}`$ of $`L^k`$ over $`M`$ satisfying the following bounds: $`|s_{k,x}^{\text{ref}}|>c_s`$ at every point of a ball of $`g_k`$-radius $`1`$ centered at $`x`$, for some universal constant $`c_s>0`$; the sections $`s_{k,x}^{\text{ref}}`$ have Gaussian decay away from $`x`$ in $`C^r`$-norm.
Moreover, given a one-parameter family of compatible almost-complex structures $`(J_t)_{t[0,1]}`$, there exist one-parameter families of sections $`s_{t,k,x}^{\text{ref}}`$ which depend continuously on $`t`$ and satisfy the same precedent properties. $`\mathrm{}`$
The proof of this Lemma uses in particular a refined version of Darboux’ Theorem taking into account the holomorphic structure, which we also enunciate for later use.
###### Lemma 2.6 (Lemma 3 in Chapter 3 of \[Au99\]).
Near any point $`xM`$, for any integer $`k1`$, there exist local complex Darboux coordinates $`(z_k^1,\mathrm{},z_k^n)=\mathrm{\Phi }_k:(M,x)(^n,0)`$ for the symplectic structure $`k\omega `$ such that the followings bounds hold universally: $`|\mathrm{\Phi }_k(y)|^2=O(d_k(x,y)^2)`$ on a ball $`B_{g_k}(x,c)`$ of universal radius $`c`$ around $`x`$; $`|^r\mathrm{\Phi }_k^1|_{g_k}=O(1)`$ for all $`r1`$ on a ball $`B(0,c^{})`$ of universal radius $`c^{}`$ around $`0`$; and, with respect to the almost-complex structure $`J`$ on $`X`$ and the canonical complex structure $`J_0`$ on $`^n`$, $`|\overline{}\mathrm{\Phi }_k^1(z)|_{g_k}=O(k^{1/2}|z|)`$ and $`|^r\overline{}\mathrm{\Phi }_k^1|_{g_k}=O(k^{1/2})`$ for all $`r1`$ on $`B(0,c^{})`$.
Moreover, given a one-parameter continuous family of compatible $`(J_t)_{t[0,1]}`$ and a continuous family of points $`(x_t)_{t[0,1]}`$, there exists a continuous family of Darboux coordinates $`\mathrm{\Phi }_{t,k}`$ satisfying the same estimates and depending continuously on $`t`$.
Proof. In \[Au99\] the result is stated only for the case $`n=2`$ but the proof extends to the case $`n>2`$ trivially. $`\mathrm{}`$
In \[Au99\] D. Auroux used three asymptotically holomorphic sections to set up a projection from a symplectic $`4`$-manifold $`M`$ to $`^2`$. To control the behaviour of this projection he needs to assure global transversality conditions between the sections. He developes a very useful scheme to pass from local transversality conditions to global ones by means of a globalization process inspired in the results of \[Do96\]. Now we explain his idea to formalize Donaldson’s techniques.
###### Definition 2.7.
A family of properties $`P(ϵ,x)_{xM,ϵ>0}`$ of sections of bundles over $`M`$ is local and $`C^r`$-open if, given a section $`s`$ satisfying $`P(ϵ,x)`$, any section $`\sigma `$ such that $`|\sigma (x)s(x)|_{C^r}<\eta `$ satisfies $`P(ϵC\eta ,x)`$, where $`C`$ is universal.
For example, the property $`|s(x)|>ϵ`$ is local and $`C^0`$-open. The property that $`s`$ be $`ϵ`$-transverse to $`0`$ at a point $`x`$ is local and $`C^1`$-open.
###### Proposition 2.8 (Proposition 3 in Chapter 3 of \[Au99\]).
Let $`P(ϵ,x)_{xM,ϵ>0}`$ be a local and $`C^r`$-open family of properties of sections of vector bundles $`E_k`$ over $`M`$. Assume that there exist universal constants $`c`$, $`c^{}`$, $`c^{\prime \prime }`$ and $`p`$ such that given any $`xM`$, any small $`\delta >0`$, and asymptotically holomorphic sections $`s_k`$ of $`E_k`$, there exist, for all large enough $`k`$, asymptotically holomorphic sections $`\tau _{k,x}`$ of $`E_k`$ with the following properties:
1. $`|\tau _{k,x}|_{C^r,g_k}<c^{\prime \prime }\delta `$.
2. The sections $`\frac{1}{\delta }\tau _{k,x}`$ have Gaussian decay away from $`x`$ in $`C^r`$-norm.
3. The sections $`s_k+\tau _{k,x}`$ satisfy the property $`P(\eta ,y)`$ for all $`yB_{g_k}(x,c)`$, with $`\eta =c^{}\delta (\mathrm{log}(\delta ^1))^p`$.
Then, given any $`\alpha >0`$ and asymptotically holomorphic sections $`s_k`$ of $`E_k`$, there exist, for all large enough $`k`$, asymptotically holomorphic sections $`\sigma _k`$ of $`E_k`$ such that $`|s_k\sigma _k|_{C^r,g_k}<\alpha `$ and the sections $`\sigma _k`$ satisfy $`P(ϵ,x)`$ for all $`xM`$ with $`ϵ>0`$ independent of $`k`$.
Moreover, the result holds for one-parameter families of sections, provided the existence of sections $`\tau _{t,k,x}`$ satisfying properties $`1`$, $`2`$ and $`3`$ and depending continuously on $`t`$.
Proof. We only have added the constant $`c^{\prime \prime }`$ to the original statement in Proposition 3 in Chapter 3 of \[Au99\], which can be absorbed into the formula for $`\eta `$ just by enlarging $`p`$ universally. $`\mathrm{}`$
The heart of these techniques is a series of local transversality results which allow to apply Proposition 2.8. These results are based on ideas of complexity of real polynomials coming from the real algebraic geometry. The most powerful result is the following, proved in \[Do99\].
###### Definition 2.9.
A function $`f:^n^r`$ is $`\sigma `$-transverse to $`0`$ at a point $`x^n`$ if it verifies at least one the following properties:
1. $`|f(x)|>\sigma `$.
2. $`df(x)`$ has a right inverse $`\theta `$ such that $`|\theta |<\sigma ^1`$.
###### Proposition 2.10.
(Theorem 12 in \[Do99\]) There exists a universal integer $`p`$ verifying the following property: for $`0<\delta <\frac{1}{2}`$ let $`\sigma =\delta (\mathrm{log}(\delta ^1))^p`$. Let $`f`$ be a function with values in $`^r`$ defined over the ball $`B^+=B(0,\frac{11}{10})^n`$ satisfying the following bounds over $`B^+`$,
$$|f|1,|\overline{}f|\sigma ,|\overline{}f|\sigma .$$
Then there exists $`w`$ with $`|w|<\delta `$ such that $`fw`$ is $`\sigma `$-tranverse to $`0`$ over the unit ball in $`^n`$. The same result holds for one-parameter families of functions $`f_t`$ depending continously on $`t[0,1]`$, where we obtain a continuous path $`w:[0,1]B(0,\delta )`$. $`\mathrm{}`$
This Proposition is a generalization of Theorem 20 of \[Do96\], where the case $`r=1`$ is proved. Later on D. Auroux in \[Au97, Au99\] extended the result to the parametric case with $`r=1`$ and to the case $`r>m`$ respectively. Proposition 2.10 covers all the range of possibilities. We mention also that in \[IMP99\] the result is refined to control the derivatives of the path $`w_t`$ allowing so a generalization to the contact case of the asymptotically holomorphic techniques.
### 2.2. Asymptotically holomorphic embeddings in $`^{2n+1}`$
Through this Section we will study the existence of asymptotically holomorphic embeddings of a closed symplectic manifold $`(M,\omega )`$ of integer class and dimension $`2n`$, endowed with a compatible almost complex structure $`J`$, in the projective space $`^{2n+1}`$. In Section 4 we will develop the techniques to study the more general grassmannian embeddings. We want to prove the following
###### Theorem 2.11.
Given an asymptotically $`J`$-holomorphic sequence of sections $`s_k`$ of the vector bundles $`^{2n+2}L^k`$ and $`\alpha >0`$ then there exists another sequence $`\sigma _k`$ verifying that:
1. $`|s_k\sigma _k|_{C^1,g_k}<\alpha `$.
2. $`(\sigma _k)`$ is an asymptotically holomorphic sequence of embeddings in $`^{2n+1}`$ for $`k`$ large enough.
3. $`[k\omega ]=[\varphi _k^{}\omega _{FS}]`$.
Moreover, let us have two asymptotically holomorphic sequences $`\varphi _k^0`$ and $`\varphi _k^1`$ of embeddings in $`^{2n+1}`$, with respect to two compatible almost complex structures. Then for $`k`$ large enough, there exists an isotopy of asymptotically holomorphic embeddings $`\varphi _k^t`$ connecting $`\varphi _k^0`$ and $`\varphi _k^1`$.
This result gives a proof of Theorem 1.2. We shall proceed by steps to obtain asymptotically holomorphic embeddings of $`M`$ into $`^{2n+1}`$.
###### Definition 2.12.
A sequence of asymptotically $`J`$-holomorphic sections $`s_k`$ of the vector bundles $`^{2n+2}L^k`$ is $`\gamma `$-projectizable if for all $`xM`$, $`|s_k(x)|>\gamma `$.
This is a sufficient condition to get a map to $`^{2n+1}`$ defined as $`\varphi _k=(s_k):M^{2n+1}`$, as the $`\gamma `$-projectizability assures that the sections $`s_k=(s_k^0,\mathrm{},s_k^{2n+1})`$ are not simultaneously zero and so the $``$ operator is well defined. To get local injectivity we need to impose the following.
###### Definition 2.13.
Let $`s_k`$ be a sequence of asymptotically $`J`$-holomorphic $`\gamma `$-projectizable sections of the vector bundles $`^{2n+2}L^k`$ for some $`\gamma >0`$ and let $`0ln`$. Then $`s_k`$ is $`\eta `$-generic of order $`l`$, with $`\eta >0`$, if $`|^l(s_k)(x)|_{g_k}>\eta `$ for all $`xM`$. For $`l=0`$ the condition is vacuus.
We have the following result that will be proved in the following two Subsections.
###### Proposition 2.14.
Let $`s_k`$ be an asymptotically $`J`$-holomorphic sequence of sections of the vector bundles $`^{2n+2}L^k`$ and $`\alpha >0`$. Then there exists another asymptotically holomorphic sequence $`\sigma _k`$ verifying:
1. $`|s_k\sigma _k|_{C^1,g_k}<\alpha `$.
2. $`\sigma _k`$ is $`\gamma `$-projectizable and $`\gamma `$-generic of order $`n`$ for some $`\gamma >0`$.
Moreover, the result holds for one-parameter families of sections where the sections and almost complex structures depend continuously on $`t[0,1]`$.
With this result we can give the proof of Theorem 2.11.
Proof of Theorem 2.11. We first prove the existence result. The last property is obvious since the hyperplane bundle of $`^{2n+1}`$ restricts by construction to $`L^k`$. Let us begin with an asymptotically $`J`$-holomorphic sequence $`\sigma _k`$ of sections of the bundles $`^{2n+2}L^k`$. Now we perturb it using Proposition 2.14 to obtain an asymptotically holomorphic sequence $`s_k`$ with $`|s_k\sigma _k|_{C^1,g_k}<\alpha `$, which is $`\gamma `$-projectizable and $`\gamma `$-generic of order $`n`$, for some $`\gamma >0`$. We have only to check that the sequence $`\varphi _k=(s_k)`$ satisfies the required properties in Definition 1.1. More specifically, we shall check that $`\varphi _k`$ is an immersion of $`M`$ in $`^{2n+1}`$, for $`k`$ large. To get rid of the possible self-intersection we take into account that $`2\text{dim }M<\text{dim }^{2n+1}`$ so we can make a generic $`C^r`$-perturbation of norm less than $`O(k^{1/2})`$ to get an embedding keeping the asymptotic holomorphicity and the genericity of order $`n`$.
Choose a point $`xM`$. By a rotation with an element of $`U(2n+2)`$ acting on $`^{2n+2}`$, we can assure that $`s_k(x)=(s_k^0(x),\mathrm{},s_k^{2n+1}(x))=(s_k^0(x),0,\mathrm{},0)`$. The transformation is constant on $`M`$ and only produces a global isometric transformation of $`\varphi _k(M)`$ in $`^{2n+1}`$. Now using the $`\gamma `$-projectizable property we know that $`|s_k^0(x)|\gamma `$. By the asymptotically holomorphic bounds of $`s_k^0`$ there is a universal $`c`$ such that $`|s_k^0|\gamma /2`$ on $`B_{g_k}(x,c)`$ for all $`k`$. We define the application:
$`f_k:B_{g_k}(x,c)`$ $``$ $`^{2n+1}`$
$`y`$ $``$ $`({\displaystyle \frac{s_k^1(y)}{s_k^0(y)}},\mathrm{},{\displaystyle \frac{s_k^{2n+1}(y)}{s_k^0(y)}}).`$
This application can be written as $`f_k=\mathrm{\Phi }_0\varphi _k`$, where $`\mathrm{\Phi }_0`$ is the standard trivialization application in $`^{2n+1}`$ defined for the chart $`U_0=\{x=[x_0,\mathrm{},x_{2n+1}]|x_00\}`$. It is well known that $`\mathrm{\Phi }_0`$ is an isometry at the point $`[1,0,\mathrm{},0]`$ if we use the standard metric structure of $`^{2n+1}`$. So we can compute the bounds required in Definition 1.1 using $`f_k`$ instead of $`\varphi _k`$. The asymptotic holomorphicity of $`s_k`$ and the bound $`|s_k^0|\gamma /2`$ imply that $`|^pf_k(x)|=O(1)`$ and $`|^p\overline{}f_k(x)|=O(k^{1/2})`$, for $`p0`$. This proves condition 3 in Definition 1.1.
Now we pass to the issue of the existence of a left inverse. We have the decomposition
$$^nd\varphi _k=^n\varphi _k+O(k^{1/2}),$$
where the last term is obtained thanks to $`|\overline{}\varphi _k|_{g_k}=O(k^{1/2})`$. By the $`\gamma `$-genericity of order $`n`$ of $`\varphi _k`$, $`|^n\varphi _k|_{g_k}\gamma `$, so $`|^nd\varphi _k|_{g_k}\gamma /2`$ for $`k`$ large. Let
$$\widehat{\theta }_k=(d\varphi _k)^1:(\varphi _k)_{}T_xMT_xM.$$
By the asymptotic holomorphicity condition, we have $`|d\varphi _k|_{g_k}C_0`$ for a universal constant $`C_0`$, so $`|\widehat{\theta }_k|C\gamma ^1`$ for another universal constant $`C`$. Now define $`\theta _k=\widehat{\theta }_k\mathrm{pr}^{}`$, where $`\mathrm{pr}^{}`$ is the orthogonal projection of $`T_{\varphi _k(x)}^{2n+1}`$ onto $`(\varphi _k)_{}T_xM`$ to get the sought right inverse (reducing $`\gamma `$ conveniently).
Finally we compute the norm of $`(\varphi _k)_{}JJ_0:(\varphi _k)_{}T_xMT_{\varphi _k(x)}^{2n+1}`$. The expression can be written as
$$(\varphi _k)_{}JJ_0=d\varphi _kJ\widehat{\theta }_kJ_0=(d\varphi _k+J_0d\varphi _kJ)J\widehat{\theta }_k=2\overline{}\varphi _kJ\widehat{\theta }_k=O(k^{1/2}),$$
proving condition 2 in Definition 1.1.
For the isotopy result we follow the ideas of \[Au97\]. We need the following auxiliary result, which we prove in Subsection 2.5.
###### Lemma 2.15.
Let $`\varphi _k:M^{2n+1}`$ be a sequence of asymptotically holomorphic embeddings with $`\varphi _k^{}[\omega _{FS}]=[k\omega ]`$. Then there exists a sequence of asymptotically holomorphic sections $`s_k`$ of $`^{2n+2}L^k`$, for $`k`$ large enough, which is $`\gamma `$-projectizable and $`\gamma `$-generic of order $`n`$, for some $`\gamma >0`$, such that $`\varphi _k=(s_k)`$. The same holds for continuous one-parameter families of embeddings and compatible almost complex structures.
Using Lemma 2.15, we can suppose that $`\varphi _k^i=(s_k^i)`$, $`i=0,1`$, where $`s_k^0`$ and $`s_k^1`$ are two asymptotically holomorphic sequences which are $`\gamma `$-projectizable and $`\gamma `$-generic of order $`n`$, $`\gamma >0`$. We construct the following family of sequences of asymptotically holomorphic sections:
$$s_k^t=\{\begin{array}{ccc}(13t)s_k^0,\hfill & \mathrm{with}J_t=J_0,\hfill & t[0,1/3]\hfill \\ 0,\hfill & \mathrm{with}J_t=\text{Path}(J_0,J_1),\hfill & t[1/3,2/3]\hfill \\ (3t2)s_k^1,\hfill & \mathrm{with}J_t=J_1,\hfill & t[2/3,1].\hfill \end{array}$$
Choose $`\alpha >0`$ such that any perturbation of $`s_k^t`$ of $`C^1`$-norm less than $`\alpha `$ is still $`\gamma /2`$-projectizable and $`\gamma /2`$-generic of order $`n`$. Applying Proposition 2.14 to $`s_k^t`$ with this $`\alpha `$, we obtain a family $`\sigma _k^t`$ which is $`\eta `$-projectizable and $`\eta `$-generic of order $`n`$ for some $`\eta >0`$. We define the family of sequences of asymptotically holomorphic sections:
$$\tau _k^t=\{\begin{array}{cc}(13t)s_k^0+3t\sigma _k^0,\hfill & t[0,1/3]\hfill \\ \sigma _k^{3t1},\hfill & t[1/3,2/3]\hfill \\ (3t2)s_k^1+(33t)\sigma _k^1,\hfill & t[2/3,1].\hfill \end{array}$$
These are $`ϵ`$-projectizable and $`ϵ`$-generic of order $`n`$ sequences of sections, with $`ϵ=\mathrm{min}\{\gamma /2,\eta \}`$, so that $`\varphi _k^t=(\tau _k^t)`$ are asymptotically holomorphic embeddings (maybe after a further small perturbation to get rid of self-intersections). This implies that $`\varphi _k^0`$ and $`\varphi _k^1`$ are isotopic for $`k`$ large enough. $`\mathrm{}`$
An important corollary is the existence of symplectic embeddings of $`M`$. The following result is similar to \[Ti77\], but we do not obtain an exact symplectic embedding. On the other hand the dimension of the projective space is controlled in our case.
###### Corollary 2.16.
Let $`(M,\omega )`$ be a closed symplectic manifold of dimension $`2n`$ with symplectic form of integer class. Then there exists a symplectic embedding $`\varphi :M^{2n+1}`$ verifying that $`k\omega =\varphi ^{}\omega _{FS}`$, for $`k`$ large enough.
Proof. Take a $`\gamma `$-asymptotically holomorphic sequence $`\varphi _k`$ of embeddings of $`M`$ in $`^{2n+1}`$. The key idea is that the linear segment of forms $`\omega _t`$ joining two symplectic forms compatible with a fixed $`J`$ is symplectic for every $`t`$. In our case we have this condition asymptotically. Define the family of $`2`$-forms in $`M`$ given by $`\omega _t=(1t)k\omega +t\varphi _k^{}(\omega _{FS})`$, where $`t[0,1]`$. All of them are cohomologous, so to apply Moser’s trick \[MS94\] we only need to prove that they are symplectic. Suppose that there exists $`t`$ such that $`\omega _t`$ is not symplectic. Then there is a unitary tangent vector $`vT_xM`$, for some $`xM`$, such that $`\omega _t(v,w)=0`$, for all $`wT_xM`$. In particular $`\omega _t(v,Jv)=0`$. Now expanding this expression we obtain:
$`\omega _t(v,Jv)`$ $`=`$ $`(1t)k\omega (v,Jv)+t\varphi _k^{}\omega _{FS}(v,Jv)`$
$`=`$ $`(1t)kg(v,v)+t\omega _{FS}(d\varphi _kv,J_0\varphi _kvJ_0\overline{}\varphi _kv)`$
$`=`$ $`(1t)kg(v,v)+tg_{FS}(d\varphi _kv,\varphi _kv)tg_{FS}(d\varphi _kv,\overline{}\varphi _kv)`$
$`=`$ $`(1t)kg(v,v)+tg_{FS}(d\varphi _kv,d\varphi _kv)2tg_{FS}(d\varphi _kv,\overline{}\varphi _kv)`$
$`=`$ $`(1t)kg(v,v)+tg_{FS}(d\varphi _kv,d\varphi _kv)tO(k^{1/2}).`$
Thanks to the $`\gamma `$-asymptotically holomorphic embeddings, we have that $`g_{FS}(d\varphi _kv,d\varphi _kv)\gamma ^2`$. So for $`k`$ large enough we get a contradiction. $`\mathrm{}`$
### 2.3. Construction of $`\gamma `$-projectizable sections.
Our objective is to prove the following perturbation result.
###### Proposition 2.17.
Let $`s_k`$ be an asymptotically $`J`$-holomorphic sequence of sections of vector bundles $`^{2n+2}L^k`$. Then given $`\alpha >0`$, there exists an asymptotically $`J`$-holomorphic sequence of sections $`\sigma _k`$ verifying:
1. $`|s_k\sigma _k|_{C^1,g_k}<\alpha `$.
2. $`\sigma _k`$ is $`\eta `$-projectizable for some $`\eta >0`$.
Moreover, the result can be extended to continuous one-parameter families of asymptotically $`J_t`$-holomorphic sequences $`s_{t,k}`$ obtaining continuous one-parameter families of sections $`\sigma _{t,k}`$ verifying the two precedent conditions.
Proof. The result is a simple generalization of Proposition 1 in \[Au99\] where the result for $`4`$-manifolds is proved. The high dimensional case can be treated with the same techniques.
We will proceed by using the globalization argument described in Proposition 2.8. First we deal with the non-parametric case. For this we define the local and $`C^0`$-open property $`P(ϵ,x)`$ as $`|s_k(x)|>ϵ`$. Let $`\delta >0`$. We only need to find for a point $`xM`$ a section $`\tau _{k,x}`$ with Gaussian decay away from $`x`$, assuring that $`s_k+\tau _{k,x}`$ verifies $`P(\eta ,y)`$ in a ball of universal $`g_k`$-radius $`c`$, with $`\eta =c^{}\delta (\mathrm{log}(\delta ^1))^p`$, $`c^{}`$ and $`p`$ universal constants.
For this choose a section $`s_{k,x}^{\text{ref}}`$ verifying the conditions of Lemma 2.5. Then we select $`c=1`$ (obviously, universal). The lower bound of $`s_{k,x}^{\text{ref}}`$ in the ball $`B_x=B_{g_k}(x,1)`$ let us define the application
$$f_{k,x}=\frac{s_k}{s_{k,x}^{\text{ref}}}:B_x^{2n+2}.$$
Using the lower bound of $`s_{k,x}^{\text{ref}}`$ together with the asymptotic holomorphicity of $`s_k`$ is easy to show that
(1)
$$|f_{k,x}|<C,|\overline{}f_{k,x}|<Ck^{1/2},|\overline{}f_{k,x}|<Ck^{1/2},$$
where $`C`$ is a universal constant. With the aid of Lemma 2.6 we can build $`f_k=f_{k,x}\mathrm{\Phi }_k^1`$ defined on a fixed ball $`B(0,c^{})^n`$. Scaling the coordinates by a universal constant $`\frac{11}{10}(c^{})^1`$ we can suppose that $`f_k`$ is defined on $`B^+`$. In this ball, the bounds (1) yield
(2)
$$|f_k|<C_0,|\overline{}f_k|<C_0k^{1/2},|\overline{}f_k|<C_0k^{1/2},$$
where $`C_0`$ is a universal constant. The application $`g_k=\frac{1}{C_0}f_k`$ is in the hypothesis of Proposition 2.10 and then there exists, for $`k`$ large enough, a number $`w_kB(0,\delta )`$ such that $`|g_kw_k|>\sigma =\delta (\mathrm{log}(\delta ^1))^p`$. Therefore $`|f_kC_0w_k|>C_0\sigma `$ on $`B`$. Now define $`\tau _{k,x}=C_0w_ks_{k,x}^{\text{ref}}`$, so that $`|\tau _{k,x}|_{C^r,g_k}<c^{\prime \prime }\delta `$, for some universal constant $`c^{\prime \prime }`$. Using the lower bound of $`s_{k,x}^{\text{ref}}`$ we obtain that $`|s_k+\tau _{k,x}|c^{}\delta (\mathrm{log}(\delta ^1))^p`$, with $`c^{}`$ and $`p`$ universal constants. Then Proposition 2.8 applies and the proof is concluded in the non-parametric case.
The globalization to the one-parameter case is trivial because all the ingredients in the proof can be easily chosen in a continuous way. $`\mathrm{}`$
### 2.4. Inductive construction of sections $`\gamma `$-generic of order $`l`$
Now we study the problem of perturbing a $`\gamma `$-projectizable sequence of sections to achieve genericity of order $`n`$. We shall do this in steps. The result to be proved is the following
###### Proposition 2.18.
Let $`s_k`$ be an asymptotically $`J`$-holomorphic sequence of sections of the vector bundles $`^{2n+2}L^k`$ which is $`\gamma `$-projectizable and $`\gamma `$-generic of order $`l`$. Then given $`\alpha >0`$, there exists an asymptotically $`J`$-holomorphic sequence of sections $`\sigma _k`$ verifying:
1. $`|s_k\sigma _k|_{C^1,g_k}<\alpha `$.
2. $`\sigma _k`$ is $`\eta `$-generic of order $`l+1`$ for some $`\eta >0`$.
Moreover, this can be extended to continuous one-parameter families of asymptotically $`J_t`$-holomorphic sequences $`s_{t,k}`$ obtaining continuous one-parameter families of sections $`\sigma _{t,k}`$ verifying conditions 1 and 2.
Proof. We construct local $`1`$-forms to control the perturbations. For this at a neighborhood of a point $`xM`$ we fix local complex Darboux coordinates $`(z_k^1,\mathrm{},z_k^n)`$ using Lemma 2.6. As in proof of Theorem 2.11, by applying a unitary transformation to $`^{2n+2}`$, we can suppose that $`s_k(x)=(s_k^0(x),0,\mathrm{},0)`$. Also there exists a ball with center $`x`$ and universal $`g_k`$-radius $`c`$ on which $`|s_k^0|\gamma /2`$. We define, following Auroux’ notations \[Au99\], a local basis of asymptotically holomorphic $`1`$-forms:
$$\mu _k^j=\left(\frac{z_k^js_{k,x}^{\text{ref}}}{s_k^0}\right),$$
where $`s_{k,x}^{\text{ref}}`$ are given by Lemma 2.5. They have Gaussian decay away from $`x`$ thanks to the behaviour of $`s_{k,x}^{\text{ref}}`$. At $`x`$ they form an orthonormal basis of $`T_x^{}M`$. We use the trivialization $`\mathrm{\Phi }_0`$ to define the application
(3) $`f_k:B_{g_k}(x,c)`$ $``$ $`^{2n+1}`$
$`y`$ $``$ $`({\displaystyle \frac{s_k^1(y)}{s_k^0(y)}},\mathrm{},{\displaystyle \frac{s_k^{2n+1}(y)}{s_k^0(y)}}),`$
which is almost an isometry on $`B_{g_k}(x,c)`$.
The case $`l=0`$ without parameters is the easiest. We say that a section $`\gamma /2`$-projectizable verifies $`P(ϵ,x)`$ if $`|\varphi _k(x)|>ϵ`$. This property is local and open in $`C^1`$-sense. We are going to apply Proposition 2.8 to assure the existence of a $`\eta `$-generic of order $`1`$ sequence of sections arbitrarily near the given $`s_k`$ in $`C^1`$-norm, for some $`\eta >0`$. For this let $`0<\delta <\gamma /2c^{\prime \prime }`$, $`c^{\prime \prime }`$ a universal constant whose precise value will appear later. We have to build a local perturbation $`\tau _{k,x}`$ with $`|\tau _{k,x}|<c^{\prime \prime }\delta `$ and Gaussian decay to achieve the property $`P(\eta ,y)`$ in a neighborhood of $`x`$ of universal $`g_k`$ radius $`c`$, with $`\eta =c^{}\delta (\mathrm{log}(\delta ^1))^p`$. (As we only perturb with sections of $`C^0`$-norm less than $`\gamma /2`$ we can assure that all the sections still have the property $`\gamma /2`$-projectizable.)
Fixing $`xM`$, we have the applications $`f_k`$ of (3). It is easy to check that there is a ball of universal radius $`c_0`$ where
$$f_k=(u_k^{11}\mu _k^1+u_k^{12}\mu _k^2+\mathrm{}+u_k^{1n}\mu _k^n,\mathrm{},u_k^{2n+1,1}\mu _k^1+\mathrm{}+u_k^{2n+1,n}\mu _k^n),$$
for some $`u_k^{ij}`$. Then we obtain an application $`u_k:B_{g_k}(x,c_0)^{n\times (2n+1)}`$. Using a complex Darboux chart we can trivialize $`B_{g_k}(x,c_0)`$ to obtain (scaling the coordinates by an appropriate universal constant $`C`$) an application $`\widehat{u}_k:B^+^{n\times (2n+1)}`$ which is asymptotically holomorphic by construction. So we can apply Proposition 2.10 to get $`w_k^{}^{n\times (2n+1)}`$ such that $`|\widehat{u}_kw_k^{}|>\eta =\delta (\mathrm{log}(\delta ^1))^p`$ on $`B`$, where $`|w_k^{}|<\delta `$. Rescaling and passing to the manifold we obtain that $`|u_kCw_k^{}|>C\delta (\mathrm{log}(\delta ^1))^p`$. We denote $`w_k=Cw_k^{}`$ and define the section $`\tau _{k,x}=(0,w_k^{11}z_k^1s_{k,x}^{\text{ref}}+w_k^{12}z_k^2s_{k,x}^{\text{ref}}+\mathrm{}+w_k^{1n}z_k^ns_{k,x}^{\text{ref}},\mathrm{},w_k^{2n+1,1}z_k^1s_{k,x}^{\text{ref}}+\mathrm{}+w_k^{2n+1,n}z_k^ns_{k,x}^{\text{ref}})`$ of $`^{2n+2}L^k`$. This section verifies the properties required in Proposition 2.8.
To check the one-parameter case we have only to get a continuous family of unitary transformations verifying that $`s_{t,k}(x)=(s_{t,k}^0(x),0,\mathrm{},0)`$ for all $`t[0,1]`$. This is clearly possible because of the contractibility of $`[0,1]`$.
Now we pass to the case $`l>0`$. We define the following property for sections $`s_k`$ which are $`\gamma /2`$-projectizable and $`\gamma /2`$-generic of order $`l`$. A section $`s_k`$ has the property $`P(ϵ,x)`$ if $`|^{l+1}s_k(x)|>ϵ`$. This property is local and open in $`C^1`$-sense. For applying Proposition 2.8 we need to build, for $`0<\delta <\gamma /2c^{\prime \prime }C`$, a local perturbation $`\tau _{k,x}`$ with $`|\tau _{k,x}|<c^{\prime \prime }\delta `$ and Gaussian decay with the property $`P(\eta ,y)`$ in a neighborhood of $`x`$ of universal $`g_k`$ radius $`c`$, with $`\eta =c^{}\delta (\mathrm{log}(\delta ^1))^p`$. (Here $`C`$ is the constant of the $`C^1`$-openness of $`P(ϵ,x)`$ in Definition 2.7.) We define $`f_k`$ as in (3). Then it is easy to see that there exists a universal constant $`c`$ such that
$$\frac{|^{l+1}(s_k)|}{|^{l+1}f_k|}>1/2$$
on $`B_{g_k}(x,c)`$. So we can do the computations for the applications $`f_k`$. By a unitary transformation in $`U(2n+1)`$ (on $`^{2n+2}`$ fixing $`(1,0,\mathrm{},0)`$) and other in $`U(n)`$ (on the complex Darboux coordinate chart) we can assure that
(4)
$$f_k(x)=\left(\begin{array}{ccccccc}u_k^{11}(x)& 0& \mathrm{}& & & \mathrm{}& 0\\ 0& u_k^{22}(x)& 0& \mathrm{}& & & 0\\ 0& \mathrm{}& \mathrm{}& 0& \mathrm{}& & 0\\ 0& \mathrm{}& 0& u_k^{nn}(x)& 0& \mathrm{}& 0\end{array}\right),$$
where $`|u_k^{11}(x)\mathrm{}u_k^{ll}(x)|>\gamma /C^{}`$, $`C^{}`$ a universal constant. Shrinking $`c`$ if necessary we can assure that $`|(f_k^1\mathrm{}f_k^l)_{\mu _k^1\mathrm{}\mu _k^l}|>\gamma /2C^{}`$ for all the points of the ball $`B_{g_k}(x,c)`$, where we denote by $`(f_k^1\mathrm{}f_k^l)_{\mu _k^1\mathrm{}\mu _k^l}`$ the component of $`f_k^1\mathrm{}f_k^l`$ in the direction of $`\mu _k^1\mathrm{}\mu _k^l`$. This $`l`$-form is an element of the basis composed by the $`l`$-wedge products of the $`1`$-forms $`\mu _k^1,\mathrm{},\mu _k^n`$. In matrix form we are denoting the order $`l`$ left upper minor of $`f_k`$. Now we construct the $`(l+1)`$-form
$$\theta _k(y)=(f_k^1\mathrm{}f_k^l)_{\mu _k^1\mathrm{}\mu _k^l}\mu _k^{l+1}.$$
We can suppose that $`|\theta _k|>c_s\gamma `$ with $`c_s>0`$ a universal constant. We also consider the following family of $`(l+1)`$-forms
$$M_k^p=(f_k^1\mathrm{}f_k^lf_k^p)_{\mu _k^1\mathrm{}\mu _k^l\mu _k^{l+1}},l+1p2n+1.$$
These forms are components of $`^{l+1}f_k`$. If we perturb so that the norm of $`M_k=(M_k^{l+1},\mathrm{},M_k^{2n+1})`$ is bigger than $`\eta =c^{}\delta (\mathrm{log}(\delta ^1))^p`$ then we have finished because if $`|M_k|>\eta `$ then $`|^{l+1}f_k|>C_0\eta `$ where $`C_0`$ is again a universal constant (using that the basis $`\{\mu _k^{i_1}\mathrm{}\mu _k^{i_{l+1}}\}_{1i_1<\mathrm{}<i_{l+1}n}`$ is almost orthogonal on the ball $`B_{g_k}(x,c)`$, in fact orthogonal at $`x`$).
We define the following sequence of asymptotically holomorphic applications,
$$g_k=(g_k^{l+1},\mathrm{},g_k^{2n+1})=(\frac{M_k^{l+1}}{\theta _k},\mathrm{},\frac{M_k^{2n+1}}{\theta _k}).$$
So we obtain, scaling the coordinates by universal constants if necessary, $`\widehat{g}_k:B^+^{2n+1l}`$ which is asymptotically holomorphic thanks to the lower bound of $`\theta _k`$ and to the asymptotic holomorphicity of $`M_k`$ and $`\theta _k`$. We have that $`n<2n+1l`$ and so we can find $`|w_k|<\delta `$ such that $`|g_kw_k|>\eta =\delta (\mathrm{log}(\delta ^1))^p`$. Thus we obtain that $`|(M_k^{l+1}w_k^{l+1}\theta _k,\mathrm{},M_k^{2n+1}w_k^{2n+1}\theta _k)|>c_s\gamma \eta `$. Recall that all the constants depend on $`\gamma `$ and the asymptotic holomorphicity constants of $`s_k`$, so they are independent of $`x`$ and $`k`$. The perturbation $`(w_k^{l+1}\theta _k,\mathrm{},w_k^{2n+1}\theta _k)`$ is achieved by adding the section $`\tau _{k,x}=(0,\stackrel{(l)}{\mathrm{}},0,w_k^{l+1}z_k^{l+1}s_{k,x}^{\text{ref}},\mathrm{},w_k^{2n+1}z_k^{l+1}s_{k,x}^{\text{ref}})`$ to $`s_k`$. This section verifies the Gaussian decay bounds required in Proposition 2.8 and $`|\tau _{k,x}|_{C^1,g_k}<c^{\prime \prime }\delta `$ for some universal constant $`c^{\prime \prime }`$. This completes the proof in the non-parametric case.
Now we pass to the one-parameter case. With appropriate continuous unitary transformations, we may assume that $`s_{t,k}(x)=(s_{t,k}^0(x),0,\mathrm{},0)`$ and that $`f_{t,k}(x)`$ is written as in (4). The interval $`[0,1]`$ may be split in a finite number of subintervals $`[t_i,t_{i+1}]`$ such that, for every $`xM`$ and for each of the subintervals, there is a fixed order $`l`$ minor of $`f_{t,k}(x)`$ with norm bigger than $`\gamma /C^{}`$, for every $`t`$ in the subinterval. This allows to find global small perturbations of $`s_{t,k}`$ in every $`[t_i,t_{i+1}]`$. Reducing $`\alpha `$ and enlarging $`C^{}`$ we may suppose that the same happens to any perturbation of the original $`s_{t,k}`$ at $`C^1`$-distance at most $`\alpha `$.
Now work as follows. For the first subinterval, consider $`s_{t,k}^1=s_{t,k}`$, $`t[0,t_1]`$. We find a perturbation $`\sigma _{t,k}^1`$, $`t[0,t_1]`$, such that $`|s_{t,k}^1\sigma _{t,k}^1|<\alpha /2`$ and $`\sigma _{t,k}^1`$ is $`\eta _1`$-generic of order $`l+1`$, for some $`\eta _1>0`$. Set $`\sigma _{t,k}=\sigma _{t,k}^1`$ for $`t[0,t_1]`$. In the second subinterval, perturb $`s_{t,k}^2=s_{t,k}^1+(\sigma _{t_1,k}^1s_{t_1,k}^1)`$, $`t[t_1,t_2]`$, to find $`\sigma _{t,k}^2`$ satisfying $`|s_{t,k}^2\sigma _{t,k}^2|<\alpha /4`$ and $`\sigma _{t,k}^2`$ is $`\eta _2`$-generic of order $`l+1`$, for some $`\eta _2>0`$. To glue this perturbation with the previous one puts
$$\sigma _{t,k}=\{\begin{array}{cc}s_{t,k}^2+\frac{tt_1}{ϵ}(\sigma _{t,k}^2s_{t,k}^2),\hfill & t[t_1,t_1+ϵ]\hfill \\ \sigma _{t,k}^2,\hfill & t[t_1+ϵ,t_2].\hfill \end{array}$$
Here $`ϵ>0`$ is chosen so small that $`|s_{t,k}^2\sigma _{t_1,k}^1|_{C^1}<\rho /2`$, for $`t[t_1,t_1+ϵ]`$, and we require also that the perturbation satisfies $`|s_{t,k}^2\sigma _{t,k}^2|<\rho /2`$, where $`\rho >0`$ is a number such that any perturbation of $`\sigma _{t_1,k}^1`$ of $`C^1`$-norm less than $`\rho `$ is $`\eta _1/2`$-generic of order $`l+1`$. This defines $`\sigma _{t,k}`$ for $`t[0,t_2]`$ already.
Proceeding in this way we finally find $`\sigma _{t,k}`$, $`t[0,1]`$, which is $`\eta `$-generic of order $`l+1`$, for some $`\eta >0`$, with $`|\sigma _{t,k}s_{t,k}|<\alpha `$. $`\mathrm{}`$
### 2.5. Lifting asymptotically holomorphic embeddings
In this Subsection we aim to prove that the sequences of asymptotically holomorphic embeddings into $`^{2n+1}`$ that we are considering in Theorem 1.2 come always from asymptotically holomorphic sequences of sections $`s_k`$ of $`^{2n+2}L^k`$ which are $`\gamma `$-projectizable and $`\gamma `$-generic of order $`n`$, for some $`\gamma >0`$ (at least for $`k`$ large).
Proof of Lemma 2.15. Suppose that we have a sequence of $`\gamma `$-asymptotically holomorphic embeddings $`\varphi _k:M^{2n+1}`$, for some $`\gamma >0`$, with $`\varphi _k^{}𝒰=L^k`$. Here $`𝒰`$ is the hyperplane line bundle defined over the projective space. The dual of $`𝒰`$ is the universal line bundle
$$=\{(l,s)|sl\}^{2n+1}\times ^{2n+2}=\underset{¯}{}^{2n+2},$$
interpreted as a sub-bundle of the trivial bundle $`\underset{¯}{}^{2n+2}`$.
Consider the following sequence of line bundles, $`E_k=\varphi _k^{}L^k=\underset{¯}{}^{2n+2}L^k`$, which are topologically trivial. We look for everywhere non-zero sections $`s_k`$ of $`E_k^{2n+2}L^k`$ as they satisfy $`\varphi _k=(s_k)`$.
Let $`P(ϵ,x)`$ be the $`C^1`$-open property for sequences of sections $`s_k`$ of $`E_k`$ of being $`ϵ`$-transverse to $`0`$ at the point $`x`$ (see Definition 2.2). We shall use Proposition 2.10 to find sequences of sections $`s_k`$ which are $`\eta `$-transverse to $`0`$, for some $`\eta >0`$. Fix any asymptotically holomorphic sequence $`s_k`$ of $`E_k`$ (e.g. the zero sections) which will act as the starting point of our perturbation process. Let $`xM`$. Consider the sections $`s_{k,x}^{\text{ref}}`$ of $`L^k`$ given by Lemma 2.5 and define also the local sections of the line bundle $`\varphi _k^{}\underset{¯}{}^{2n+2}`$,
$$\sigma _k:B_{g_k}(x,c)^{2n+2},$$
by setting $`\sigma _k(x)`$ any vector of norm $`1`$ in the direction defined by $`\varphi _k(x)`$ and satisfying the condition $`_r\sigma _k(y)\sigma _k(y)`$, for any $`yB_{g_k}(x,c)`$, where $`r`$ is the radial vector field from $`x`$. This determines $`\sigma _k`$ uniquely. The following estimates hold:
$`|\sigma _k(y)|=1,|\sigma _k(y)|=O(1+d_k(x,y)),`$
(5) $`|\overline{}\sigma _k(y)|=O(k^{1/2}(1+d_k(x,y))),|\overline{}\sigma _k(y)|=O(k^{1/2}(1+d_k(x,y))).`$
The first one follows from $`_r\sigma _k,\sigma _k=_r\sigma _k,\sigma _k+\sigma _k,_r\sigma _k=0`$. For the second one, write $`\sigma _k=\varphi _k+\sigma _k,\sigma _k\sigma _k`$, where we identify $`T_{\varphi _k(y)}^{2n+1}=[\sigma _k(y)]^{}^{2n+2}`$, isometrically. We already know that $`|\varphi _k|=O(1)`$. So
$`_r\sigma _k,\sigma _k`$ $`=`$ $`_r\sigma _k,\sigma _k+\sigma _k,_r\sigma _k=`$
$`=`$ $`_r\sigma _k,\sigma _k+\sigma _k,_r\sigma _k=`$
$`=`$ $`_r\sigma _k,\sigma _k+\sigma _k,_r\sigma _k=`$
$`=`$ $`_r\varphi _k,\varphi _k+\varphi _k,_r\varphi _k=O(1),`$
The first equality uses that $``$ is the standard derivative for functions with values in $`^{2n+2}`$, and hence the second derivatives commute. The second equality follows from $`_r\sigma _k,\sigma _k=0`$. So we have that $`\sigma _k,\sigma _k=O(d_k(x,y))`$ and hence $`|\sigma _k|=O(1+d_k(x,y))`$. The other two cases are worked out analogously.
Now define the application
$$f_k=\frac{s_k}{s_{k,x}^{\text{ref}}\sigma _k}:B_{g_k}(x,c),$$
which is asymptotically holomorphic by construction. Using a complex Darboux chart we trivialize $`B_{g_k}(x,c)`$ to obtain (scaling the coordinates by appropiate universal constants) an application $`\widehat{f}_k:B^+`$ to which we apply Proposition 2.10 to obtain $`w_kB(0,\delta )`$ such that $`\widehat{f}_kw_k`$ is $`\eta `$-transverse to $`0`$ in $`B`$, where $`\eta =\delta (\mathrm{log}(\delta ^1))^p`$. Rescaling and passing to the manifold, we have that $`f_kCw_k`$ is $`C^{}\eta `$-transverse to $`0`$, for some universal constants $`C`$ and $`C^{}`$. Define the sequence of sections $`\tau _{k,x}=w_ks_{k,x}^{\text{ref}}\sigma _k`$ of $`E_k`$, which is asymptotically holomorphic and has Gaussian decay by (5), to get a perturbation satisfying the conditions in Proposition 2.8.
Thus there exists an asymptotically holomorphic sequence $`s_k`$ of sections of $`E_k`$ which is $`\eta `$-transverse to $`0`$, for some $`\eta >0`$. For $`k`$ large enough, the zeroes of $`s_k`$ is a symplectic submanifold representing the trivial homology class, hence the empty set. So $`s_k`$ is nowhere vanishing and hence $`\varphi _k=(s_k)`$.
We have that $`s_k`$ is an asymptotically holomorphic sequence of sections of $`^{2n+2}L^k`$. Let us check that $`s_k`$ is $`\eta `$-projectizable, i.e. that $`|s_k|\eta `$ everywhere. Suppose that this is not the case and take the point $`xM`$ where $`|s_k|`$ attains its minimum. As $`|s_k(x)|<\eta `$, $`\eta `$-transversality implies that $`|s_k(x)|\eta `$. Also $`s_k`$ is asymptotically holomorphic, so for $`k`$ large $`s_k(x):T_xM(E_k)_x`$ is surjective. Take $`vT_xM`$ such that $`_vs_k(x)=s_k(x)`$. Evaluating the equality
$$|s_k|^2=s_k,s_k+s_k,s_k.$$
at the point $`x`$ and along the direction of $`v`$, we obtain $`|s_k(x)|^2=0`$, which is impossible since we have already proved that $`s_k`$ is nowhere vanishing.
Finally the extension to the one-parameter case is trivial. $`\mathrm{}`$
## 3. Estimated intersections of symplectic submanifolds
### 3.1. Notions on estimated euclidean geometry
In order to set up the definitions needed in Subsection 3.2 we state the relevant notions and results on angles between subspaces of euclidean spaces that we shall need. From now on we assume that we are in $`^n`$ equipped with the standard euclidean inner product, but all the proofs apply to a general finite dimensional euclidean space.
The angle between two non-zero vectors $`v,w^n`$ is defined as
$$\mathrm{}(v,w)=\mathrm{arccos}\left(\frac{v,w}{|v||w|}\right)[0,\pi ].$$
The angle is symmetric and satisfies the classical triangular inequality,
$$\mathrm{}(u,w)\mathrm{}(u,v)+\mathrm{}(v,w),$$
for non-zero vectors $`u,v,w^n`$. Also the angle of a vector $`u0`$ respect to a subspace $`V\{0\}`$ is defined as
$$\mathrm{}(u,V)=\underset{vV\{0\}}{\mathrm{min}}\{\mathrm{}(u,v)\}=\mathrm{}(u,v(u))[0,\frac{\pi }{2}],$$
where $`v:^n^n`$ is the orthogonal projection onto $`V`$, well understood that when $`v(u)=0`$ the angle is $`\pi /2`$.
###### Definition 3.1.
The maximum angle of a subspace $`U\{0\}`$ with respect to a subspace $`V\{0\}`$ is defined as
$$\mathrm{}_M(U,V)=\underset{uU\{0\}}{\mathrm{max}}\mathrm{}(u,V).$$
Notice that this angle is not in general symmetric. But in the case $`\text{dim }U=\text{dim }V`$ symmetry does hold. This is easily checked by constructing an orthogonal transformation permuting the two subspaces. Indeed the maximum angle $`\mathrm{}_M(U,V)`$ gives a notion of proximity between $`U`$ and $`V`$ whenever $`\text{dim }U\text{dim }V`$.
###### Lemma 3.2.
Given $`U,V,W`$ non zero-subspaces in $`^n`$ then:
$$\mathrm{}_M(U,W)\mathrm{}_M(U,V)+\mathrm{}_M(V,W).$$
Proof. We will denote by $`v(u)`$ the orthogonal projection of the vector $`u`$ onto the subspace $`V`$. In the following inequalities, if $`v(u)=0`$, we suppose that the angle in which this expression appears is $`\pi /2`$. We have
$`\mathrm{}_M(U,W)`$ $`=`$ $`\underset{uU\{0\}}{\mathrm{max}}\{\underset{wW\{0\}}{\mathrm{min}}\{\mathrm{}(u,w)\}\}`$
$``$ $`\underset{uU\{0\}}{\mathrm{max}}\{\underset{wW\{0\}}{\mathrm{min}}\{\mathrm{}(u,v(u))+\mathrm{}(v(u),w)\}\}=`$
$`=`$ $`\underset{uU\{0\}}{\mathrm{max}}\{\mathrm{}(u,v(u))+\underset{wW\{0\}}{\mathrm{min}}\{\mathrm{}(v(u),w)\}\}`$
$``$ $`\underset{uU\{0\}}{\mathrm{max}}\{\mathrm{}(u,v(u))\}+\underset{uU\{0\}}{\mathrm{max}}\{\underset{wW\{0\}}{\mathrm{min}}\{\mathrm{}(v(u),w)\}\}`$
$``$ $`\mathrm{}_M(U,V)+\underset{vV\{0\}}{\mathrm{max}}\{\underset{wW\{0\}}{\mathrm{min}}\{\mathrm{}(v,w)\}\}`$
$``$ $`\mathrm{}_M(U,V)+\mathrm{}_M(V,W).`$
$`\mathrm{}`$
###### Definition 3.3.
The minimum angle between two non-zero subspaces $`U,V`$ of $`^n`$ is defined as follows:
* If $`\text{dim }U+\text{dim }V<n`$ then $`\mathrm{}_m(U,V)=0`$.
* If their intersection is not transversal then $`\mathrm{}_m(U,V)=0`$.
* If their intersection is transversal then let $`W`$ be their intersection. Define $`U_c`$ as the orthogonal subspace in $`U`$ to $`W`$, and $`V_c`$ in the same way. Then $`\mathrm{}_m(U,V)=\mathrm{min}_{uU_c\{0\}}\{\mathrm{}(u,V_c)\}[0,\pi /2]`$.
The definition is symmetric because (in the transversal case)
$$\mathrm{}_m(U,V)=\underset{uU_c\{0\}}{\mathrm{min}}\{\underset{vV_c\{0\}}{\mathrm{min}}\{\mathrm{}(u,v)\}\}$$
and the two minima commute. Also $`\mathrm{}_m(U,V)=\mathrm{min}_{uU_c\{0\}}\{\mathrm{}(u,V)\}`$.
###### Lemma 3.4.
For non-zero subspaces $`U`$ and $`V`$ of $`^n`$ we have that
$$\mathrm{}_m(U,V)=\underset{uU^{}\{0\}}{\mathrm{min}}\{\underset{vV^{}\{0\}}{\mathrm{min}}\{\mathrm{}(u,v)\}\}$$
Proof. This is trivial in the case $`\text{dim }U+\text{dim }V<n`$ or when $`U`$ and $`V`$ do not intersect transversely. In the transversal case, we can restrict ourselves to the subspace $`(UV)^{}`$ to compute the angles. So without loss of generality we can suppose that $`UV=^n`$, $`U_c=U`$ and $`V_c=V`$. As $`\text{dim }U=\text{dim }V^{}`$, we may construct an orthogonal transformation $`\varphi `$ permuting $`U`$ and $`V^{}`$, i.e. $`\varphi (U)=V^{}`$ and $`\varphi (V^{})=U`$. Therefore also $`\varphi (V)=U^{}`$. So
$$\mathrm{}_m(U,V)=\mathrm{}_m(\varphi (U),\varphi (V))=\mathrm{}_m(V^{},U^{})=\underset{uU^{}\{0\}}{\mathrm{min}}\{\underset{vV^{}\{0\}}{\mathrm{min}}\{\mathrm{}(u,v)\}\},$$
which proves the lemma. $`\mathrm{}`$
###### Proposition 3.5.
For non-zero subspaces $`U,V,W`$ of $`^n`$ we have that
$$\mathrm{}_m(U,V)\mathrm{}_M(U,W)+\mathrm{}_m(W,V).$$
Proof. By Lemma 3.4 we have that
$`\mathrm{}_m(U,V)`$ $`=`$ $`\underset{uU^{}\{0\}}{\mathrm{min}}\{\underset{vV^{}\{0\}}{\mathrm{min}}\{\mathrm{}(u,v)\}\}`$
$``$ $`\underset{uU^{}\{0\}}{\mathrm{min}}\{\mathrm{}(u,w)\}+\underset{vV^{}\{0\}}{\mathrm{min}}\{\mathrm{}(w,v)\},`$
for any $`w^n`$. Choose $`w_0W^{}\{0\}`$ satisfying
$$\mathrm{}_m(W,V)=\underset{wW^{}\{0\}}{\mathrm{min}}\{\underset{vV^{}\{0\}}{\mathrm{min}}\{\mathrm{}(w,v)\}\}=\underset{vV^{}\{0\}}{\mathrm{min}}\{\mathrm{}(w_0,v)\}.$$
Then we have
$$\mathrm{}_m(U,V)\underset{uU^{}\{0\}}{\mathrm{min}}\{\mathrm{}(u,w_0)\}+\mathrm{}_m(W,V)\mathrm{}_M(W^{},U^{})+\mathrm{}_m(W,V).$$
The result follows once we show that $`\mathrm{}_M(W^{},U^{})=\mathrm{}_M(U,W)`$. Put $`\mathrm{}_M(U,W)=\alpha `$. Let $`uU`$ with $`\mathrm{}(u,W)=\alpha `$. Denoting by $`w`$ the projection of $`u`$ onto $`W^{}`$, we have that $`\mathrm{}(u,W^{})=\mathrm{}(u,w)=\frac{\pi }{2}\alpha `$. So $`\mathrm{}(w,U)\frac{\pi }{2}\alpha `$ and hence $`\mathrm{}(w,U^{})\alpha `$. This implies that $`\mathrm{}_M(W^{},U^{})\alpha =\mathrm{}_M(U,W)`$. The opposite inequality follows by symmetry. $`\mathrm{}`$
###### Corollary 3.6.
Given non-zero subspaces $`U,U^{},V`$ of $`^n`$ with $`\mathrm{}_m(U,V)>ϵ`$ and $`\mathrm{}_M(U,U^{})<\delta `$ then $`\mathrm{}_m(U^{},V)>ϵC\delta `$, where $`C`$ is a universal constant ($`C=1`$ in fact). $`\mathrm{}`$
The following result will be very important for our purposes.
###### Proposition 3.7.
Given $`ϵ>0`$ and $`U\text{Gr}(m,n)`$, $`V\text{Gr}(r,n)`$ subspaces verifying that $`\mathrm{}_m(U,V)>ϵ`$. Then there are $`\gamma _0>0`$ and a constant $`C`$, depending only on $`ϵ`$, such that for any $`\gamma <\gamma _0`$, if $`U^{}Gr(m,n)`$ and $`V^{}Gr(r,n)`$ verify that
$$\mathrm{}_M(U,U^{})<\gamma ,\mathrm{}_M(V,V^{})<\gamma ,$$
then $`U^{}`$ and $`V^{}`$ intersect transversally and $`\mathrm{}_M(UV,U^{}V^{})<C\gamma `$.
Proof. By Proposition 3.5 choosing $`\gamma _0>0`$ small enough, only depending on $`ϵ`$, we can assure that the following intersections are transversal $`UV=W`$, $`UV^{}`$, $`U^{}V`$ and $`U^{}V^{}=W^{}`$ and that $`\mathrm{}_m(U^{},V^{})ϵ/2`$. By Lemma 3.2 we have
$$\mathrm{}_M(W,W^{})\mathrm{}_M(W,UV^{})+\mathrm{}_M(W^{},UV^{}).$$
We are going to bound the first term in the right hand side of the inequality, the bounding of the second term being analogous.
Put $`s=\text{dim }W=r+mn`$. Choose an orthonormal basis $`(e_1,\mathrm{},e_s)`$ of $`W`$, extend it to an orthonormal basis $`(e_1,\mathrm{},e_r)`$ of $`V`$ and finally extend it to an orthonormal basis $`(e_1,\mathrm{},e_n)`$ of $`^n`$. Note that $`(e_{s+1},\mathrm{},e_r)`$ is an orthonormal basis of $`V_c`$. As $`\mathrm{}_m(U,V)=\mathrm{}_m(U_c,V)>ϵ`$ and $`\mathrm{}_M(V,V^{})<\gamma _0`$ we have $`\mathrm{}_m(U_c,V^{})>ϵ/2`$ (decreasing $`\gamma _0`$ if necessary). So $`U_cV^{}=\{0\}`$. Recalling that $`VU_c=^n`$, we see that there is a basis $`(e_1+\epsilon _1,\mathrm{},e_r+\epsilon _r)`$ for $`V^{}`$ where $`\epsilon _jU_c`$. Using that $`\mathrm{}_m(U,V)>ϵ`$ and that the decomposition $`^n=WV_cV^{}`$ is orthogonal, we have
$`\mathrm{pr}_W^{}(\epsilon _j)`$ $`=`$ $`0,`$
$`\mathrm{pr}_{V_c}^{}(\epsilon _j)`$ $``$ $`|\mathrm{cos}ϵ||\epsilon _j|,`$
$`\mathrm{pr}_V^{}^{}(\epsilon _j)`$ $``$ $`\sqrt{1|\mathrm{cos}ϵ|^2}|\epsilon _j|=|\mathrm{sin}ϵ||\epsilon _j|.`$
Checking the angle of $`e_j+\epsilon _j`$ with respect to $`V`$, we get that
(6)
$$\mathrm{}_M(V,V^{})\mathrm{arctan}\frac{|\mathrm{sin}ϵ||\epsilon _j|}{1+|\mathrm{cos}ϵ||\epsilon _j|}\mathrm{arctan}\left(\frac{\mathrm{sin}ϵ}{1+|\epsilon _j|}|\epsilon _j|\right).$$
For $`\gamma _0<\mathrm{arctan}\frac{\mathrm{sin}ϵ}{2}`$, (6) implies that $`|\epsilon _j|<1`$ and hence we get that $`\mathrm{}_M(V,V^{})\mathrm{arctan}(\frac{\mathrm{sin}ϵ}{2}|\epsilon _j|)\frac{4}{\pi }\frac{\mathrm{sin}ϵ}{2}|\epsilon _j|`$, or said otherwise $`|\epsilon _j|<C\mathrm{}_M(V,V^{})`$ for a constant $`C`$ depending on $`ϵ`$.
Now let us compute $`\mathrm{}_M(W,UV^{})`$. The intersection $`UV^{}`$ has basis $`(e_1+\epsilon _1,\mathrm{},e_s+\epsilon _s)`$. Take a general vector $`u=_{i=1}^sa_i(e_i+\epsilon _i)`$ in $`UV^{}`$ and compute $`\mathrm{}(u,W)`$. We may suppose that $`a=(a_1,\mathrm{},a_s)`$ has norm one. Write $`\epsilon =_{i=1}^sa_i\epsilon _i`$. Then
$$\mathrm{}(u,W)=\mathrm{arccos}\frac{1}{\sqrt{1+|\epsilon |^2}}=\mathrm{arctan}|\epsilon ||\epsilon |.$$
Finally
$$\mathrm{}_M(W,UV^{})\underset{|a|=1}{\mathrm{max}}|\underset{i=1}{\overset{s}{}}a_i\epsilon _i|=\underset{1is}{\mathrm{max}}|\epsilon _i|C\mathrm{}_M(V,V^{})C\gamma .$$
$`\mathrm{}`$
Now we are going to set up the relationship between the transversality of maps in the Donaldson-Auroux approach and the angles defined above. This is the content of the following
###### Lemma 3.8.
Let $`U,V`$ be two non-zero subspaces of $`^n`$ and let $`g:UV`$ and $`h:UV^{}`$ be the projections from $`U`$ with respect to the decomposition $`^n=VV^{}`$. If $`h`$ has a right inverse $`\theta `$ satisfying $`|\theta |<\gamma ^1`$ for some $`\gamma >0`$ then $`\mathrm{}_m(U,V)>\gamma `$.
Proof. In the first place, as $`h`$ is onto, the intersection between $`U`$ and $`V`$ is transversal. Let $`W=UV`$. Define $`\widehat{\theta }=\mathrm{pr}_{U_c}^{}\theta :V^{}U_c`$, which is an inverse of $`h:U_cV^{}`$ such that $`|\widehat{\theta }|<\gamma ^1`$. Now consider any $`uU_c\{0\}`$ and put $`v=h(u)`$. Then
$$\mathrm{}(u,V)=\mathrm{arcsin}\frac{|h(u)|}{|u|}=\mathrm{arcsin}\frac{|v|}{|\widehat{\theta }(v)|}>\mathrm{arcsin}\frac{1}{\gamma ^1}>\gamma ,$$
and the proof is concluded. $`\mathrm{}`$
### 3.2. Projective symplectic geometry
In this Subsection we will prove Theorem 1.3. This will provide a geometric proof of Bertini’s theorem, the main result of \[Do96\]. Although our proof is more technical and long, it has the advantage of giving us a more general kind of symplectic submanifolds than those in \[Do96, Au97\]. In fact our technique will allow us a simple generalization to solve the problem of constructing determinantal symplectic submanifolds in Section 5. First of all, in order to measure the holomorphicity of submanifolds, let us introduce the complex angle of even dimensional subspaces $`V^n`$ as
$`\beta :\text{Gr}_{}(2r,2n)`$ $``$ $`[0,\pi /2]`$
$`V`$ $``$ $`\mathrm{}_M(V,JV).`$
Clearly $`\beta (V)=0`$ if and only if $`V`$ is complex and $`\beta (V)<\pi /2`$ if and only if $`V`$ is symplectic.
###### Definition 3.9.
Let $`(M,\omega )`$ be a symplectic submanifold endowed with a compatible almost complex structure $`J`$. A sequence of submanifolds $`S_kM`$ is asymptotically holomorphic if $`\beta (TS_k)=O(k^{1/2})`$.
Note that if $`S_k`$ are asymptotically holomorphic submanifolds then they are symplectic for $`k`$ large. If $`\varphi _k:M^N`$ is a sequence of asymptotically holomorphic embeddings then $`\varphi _k(M)`$ is a sequence of asymptotically holomorphic submanifolds.
###### Proposition 3.10.
Let $`\varphi _k^1:(M_1,J_1)^N`$ and $`\varphi _k^2:(M_2,J_2)^N`$ be two sequences of asymptotically holomorphic embeddings. Suppose that there exists $`ϵ>0`$ independent of $`k`$ such that for any $`x\varphi _k^1(M_1)\varphi _k^2(M_2)`$, the minimum angle between $`(\varphi _k^1)_{}TM_1(x)`$ and $`(\varphi _k^2)_{}TM_2(x)`$ is greater than $`ϵ`$. Then $`S_k=\varphi _k^1(M_1)\varphi _k^2(M_2)`$ is a sequence of asymptotically holomorphic submanifolds (hence symplectic for $`k`$ large). Also $`S_k^j=(\varphi _k^j)^1(S_k)`$ is a sequence of asymptotically holomorphic submanifolds of $`M_j`$, $`j=1,2`$. Moreover there exists a sequence of compatible almost complex structures $`J_k^j`$ of $`M_j`$ such that $`S_k^j`$ is pseudoholomorphic for $`J_k^j`$, $`|J_k^jJ_j|=O(k^{1/2})`$ and $`\varphi _k^j`$ restricted to $`(S_k^j,J_k^j)`$ is a sequence of asymptotically holomorphic embeddings in $`^N`$, $`j=1,2`$.
The same statement holds for the case of one-parameter families of embeddings $`(\varphi _{t,k}^1)_{t[0,1]}`$ and $`(\varphi _{t,k}^2)_{t[0,1]}`$.
Remark that $`M_1`$ and $`M_2`$ are not necessarily compact manifolds.
Proof. Let $`J_0`$ be the standard complex structure of $`^{2n+1}`$. Then $`\mathrm{}_M((\varphi _k^j)_{}TM,J_0(\varphi _k^j)_{}TM)=O(k^{1/2})`$ for $`j=1,2`$. By Proposition 3.7, $`\mathrm{}_M(TS_k,J_0TS_k)=O(k^{1/2})`$. As $`|(\varphi _k^j)_{}J_jJ_0|=O(k^{1/2})`$ on $`(\varphi _k^j)_{}TM`$, we have $`\mathrm{}_M(TS_k,(\varphi _k^j)_{}J_jTS_k)=O(k^{1/2})`$ and so $`\mathrm{}_M(TS_k^j,J_jTS_k^j)=O(k^{1/2})`$, i.e. $`S_k^j`$ is a sequence of asymptotically holomorphic submanifolds of $`M_j`$.
Finally we have to build $`J_k^j`$ on $`M_j`$ such that $`|J_k^jJ_j|=O(k^{1/2})`$ and $`S_k^j`$ is $`J_k^j`$-holomorphic. Take the composition $`\stackrel{~}{J}_k^j:TS_k^jTM\stackrel{J_j}{}TM\stackrel{\mathrm{pr}^{}}{}TS_k^j`$ with square close to $`1`$, for $`k`$ large enough. So we can homotop it to an almost complex structure $`J_k^j`$ on $`S_k^j`$. Then we extend this $`J_k^j`$ to a small tubular neighborhood of $`S_k^j`$ by giving a complex structure to the normal bundle of $`S_k^j`$. Finally a homotopy between $`J_k^j`$ and $`J_j`$ allows us to extend $`J_k^j`$ off a little bigger neighborhood of $`S_k^j`$ matching with $`J_j`$ on the border. This gives the required $`J_k^j`$.
The result for continuous one-parameter families is trivial from the non-parametric case. $`\mathrm{}`$
Let us have a smooth submanifold $`N`$ of a manifold $`X`$. If we fix a metric on $`X`$ we can define a geodesic flow $`\phi _t`$. In particular, following the perperdicular directions to $`N`$ we can identify a tubular neighborhood of the zero section of the normal bundle of $`N`$ (defined as $`|n|<t_0`$, $`n\nu (N)`$, for some small $`t_0>0`$) with a tubular neighborhood $`U_NX`$ of $`N`$. So we can define an integrable distribution $`D_N`$ in $`U_N`$ as
$$D_N(\phi _n(x))=(\phi _n)_{}T_xN,xN,n\nu (N),|n|<t_0.$$
where $`(\phi _n)_{}`$ denotes parallel transport along the geodesic tangent to $`n`$.
###### Definition 3.11.
Suppose $`\varphi _k:MX`$ is a sequence of asymptotically holomorphic embeddings into a Hodge manifold $`X`$. Let us fix a complex submanifold $`NX`$. We say that $`\varphi _k`$ is $`\sigma `$-transverse to $`N`$, with $`\sigma <t_0`$, if for all $`xM`$ and all $`k`$,
$$d(\varphi _k(x),N)<\sigma \mathrm{}_m((\varphi _k)_{}(T_xM),D_N(\varphi _k(x)))>\sigma .$$
This property is $`C^1`$-open, i.e. given $`\varphi _k`$ an embedding $`\eta `$-transverse to $`N`$, then a perturbation of $`\widehat{\varphi }_k`$ with $`d_{C^1}(\varphi _k,\widehat{\varphi }_k)<\delta `$ is $`(\eta C\delta )`$-transverse to $`N`$, where $`C`$ is a universal constant.
Obviously a $`\sigma `$-transverse sequence of embeddings $`\varphi _k`$ verifies the conditions of Proposition 3.10 with $`\varphi _k^1=\varphi _k:MX`$ and $`\varphi _k^2=i:NX`$. The following result then completes the proof of Theorem 1.3
###### Theorem 3.12.
Let $`\varphi _k=(s_k)`$, where $`s_k`$ is an asymptotically holomorphic sequence of sections of $`^{2n+2}L^k`$ which is $`\gamma `$-projectizable and $`\gamma `$-generic of order $`n`$, for some $`\gamma >0`$. Let us fix a holomorphic submanifold $`N`$ in $`^{2n+1}`$. Then for any $`\delta >0`$ there exists an asymptotically holomorphic sequence of sections $`\sigma _k`$ of $`^{2n+2}L^k`$ such that
1. $`|\sigma _ks_k|_{g_k,C^1}<\delta `$.
2. $`\widehat{\varphi }_k=(\sigma _k)`$ is a $`\eta `$-asymptotically holomorphic embedding in $`^{2n+1}`$ which is $`ϵ`$-transverse to $`N`$, for some $`\eta >0`$ and $`ϵ>0`$. In the case $`\text{dim }M+\text{dim }N<2n+1`$ we actually have that $`d_{FS}(\widehat{\varphi }_k(M),N))>ϵ`$, for $`k`$ large enough.
Moreover the result can be extended to one-parameter continuous families of complex submanifolds $`(N_t)_{t[0,1]}`$, taking in this case as starting point a continuous family $`\varphi _{t,k}=(s_{t,k})`$ where $`s_{t,k}`$ are asymptotically $`J_t`$-holomorphic sections of $`^{2n+2}L^k`$ which are $`\gamma `$-projectizable and $`\gamma `$-generic of order $`n`$, for some $`\gamma >0`$.
The proof of this result will be the content of Subsection 3.3. Now we shall extract some corollaries from it. The first one is the main Theorem of \[Do96\].
###### Corollary 3.13.
Given a compact symplectic manifold $`(M,\omega )`$, suppose that $`[\omega /2\pi ]H^2(M,)`$ is the reduction of an integral class $`h`$. Then for $`k`$ large enough there exists symplectic submanifolds realizing the Poincaré dual of $`kh`$. Moreover, perhaps by increasing $`k`$, we can assure that all the symplectic submanifolds realizing this Poincaré dual, constructed as transverse intersections with a fixed complex hyperplane of asymptotically holomorphic sequences of embeddings with respect to two compatible almost complex structures, are isotopic. The isotopy can be made by symplectomorphisms.
Recall that we obtain an isotopy result similar to \[Au97\], where the isotopy of the submanifolds obtained as zero sets of a special set of sections of the line bundle $`L^k`$ is obtained. The Auroux’ more general case of vector bundles will be proved in Section 4.
Proof. The existence result is a direct consequence of the previous statements. By Theorem 2.11 we build an asymptotically holomorphic sequence of embeddings to $`^{2n+1}`$. In $`^{2n+1}`$ we choose a complex hyperplane $`H`$. By Theorem 3.12 we perturb the sequence of embeddings to find a new asymptotically holomorphic sequence of embeddings $`\varphi _k`$ such that $`\varphi _k(M)`$ intersects $`H`$ with minimum angle greater than $`ϵ>0`$. Finally using Proposition 3.10 we obtain that $`\varphi _k(M)H=H_M`$ is an asymptotically holomorphic sequence of submanifolds, and these manifolds are symplectic for $`k`$ large enough. Also $`\varphi _k^1(H_M)`$ is a symplectic submanifold of $`M`$ for $`k`$ large enough. A direct topological argument shows us that it is Poincaré dual of $`kh`$.
For the isotopy statement, let us assume that there are two sequences of symplectic submanifolds $`W_k^0`$ and $`W_k^1`$, both Poincaré dual of $`kh`$, obtained as intersections between two $`\eta `$-asymptotically $`J_j`$-holomorphic sequences $`(s_{k,j})`$, $`j=0,1`$, and two fixed complex hyperplanes $`H_0`$ and $`H_1`$ in $`^{2n+1}`$ with angles greater than a fixed $`ϵ>0`$. Then we will prove that in this case they are isotopic. We only have to construct the straight segment $`H_t`$, in the dual space, of hyperplanes connecting $`H_0`$ and $`H_1`$. Also we define the following family of asymptotically holomorphic sequences:
$$s_{t,k}=\{\begin{array}{ccc}(13t)s_{0,k},\hfill & \text{with }J_t=J_0,\hfill & t[0,1/3]\hfill \\ 0,\hfill & \text{with }J_t=\text{Path}(J_0,J_1),\hfill & t[1/3,2/3]\hfill \\ (3t2)s_{1,k},\hfill & \text{with }J_t=J_1,\hfill & t[2/3,1].\hfill \end{array}$$
By means of Theorem 3.12, we obtain a family $`\varphi _{t,k}=(\sigma _{t,k})`$ of asymptotically $`J_t`$-holomorphic embeddings which are $`\eta /2`$-transverse to $`N`$, choosing the perturbation $`\delta >0`$ in the statement of the theorem, in such a way that
(7)
$$\eta C\delta >\eta /2,$$
where $`C`$ is the universal constant of the $`C^1`$-openness of the transversality to $`N`$. This gives us a family of symplectic isotopic submanifolds $`(W_k^t)^{}`$ in $`M`$ for each fixed large $`k`$. The problem is that $`W_k^0`$ does not coincide with $`(W_k^0)^{}`$ (and respectively for $`t=1`$). Using (7) we can assure that they are isotopic, in fact the linear segment $`((1t)\sigma _{0,k}+ts_{0,k})_{t[0,1]}`$ provides a family of asymptotically holomorphic embeddings transverse to $`H_0`$, for $`k`$ large enough giving the desired isotopy. $`\mathrm{}`$
The constructive technique of Theorem 3.12 is more general because we do not have to choose hyperplanes in $`^{2n+1}`$ to make the intersection. However, the difficulty in finding topological information about the constructed submanifolds makes that we cannot assure that they are more general that the ones produced in \[Au97\]. To overcome this problem we are going to construct in Section 5 a special kind of submanifolds where we can compute symplectic invariants using similar results from algebraic geometry.
### 3.3. Estimated intersections in $`^{2n+1}`$.
Now we aim to prove Theorem 3.12. Our objective is to find sequences $`\varphi _k`$ of asymptotically holomorphic embeddings which are $`\sigma `$-transverse to $`N`$.
Proof of Theorem 3.12. As usual we begin with the simplest case, when the complex codimension of $`N`$ is $`1`$. Also we consider the non-parametric case, being the parametric one a simple generalization. We say that a sequence of sections $`s_k`$ which is $`\gamma /2`$-projectizable and $`\gamma /2`$-generic of order $`n`$ verifies $`P(ϵ,x)`$ if $`(s_k)`$ is $`ϵ`$-transverse to $`N`$ at the point $`x`$. This property is local and open in $`C^1`$-sense, for $`ϵ<t_0`$. To make use of Proposition 2.8 we need to find local sections with Gaussian decay obtaining local transversality. To achieve this local transversality we are going to use Proposition 2.10. (We could have used instead the case $`m=1`$ proved in \[Do96, Au97\], by increasing a little the complications of the globalization process, which is the way followed by Auroux in \[Au97, Au99\].)
As $`N`$ is a fixed holomorphic submanifold, we may fix a finite covering of $`^{2n+1}`$ by balls $`U_j`$ such that $`N`$ is defined as the zero set of a holomorphic function $`f_j:U_j`$ in each $`U_j`$ and such that for any $`z_1,z_2U_jU_N`$, $`\mathrm{}_M(D_N(z_1),D_N(z_2))\epsilon `$, and for any $`z_1,z_2U_j`$, $`\mathrm{}_M(\mathrm{ker}df_j(z_1),\mathrm{ker}df_j(z_2))\epsilon `$, with $`\epsilon >0`$ an arbitrarily small number fixed along the proof.
We choose a constant $`C`$ independent of $`k`$ such that $`|\varphi _k|_{g_k}C`$. Therefore $`\varphi _k(B_{g_k}(x,c))B_{g_{FS}}(\varphi _k(x),Cc)`$, for any $`c`$. Now we choose $`c>0`$ small enough satisfying the following premises:
1. Let $`xM`$. With a transformation of $`U(2n+2)`$ in $`^{2n+2}`$, we may suppose that $`s_k(x)=(s_k^0(x),0,\mathrm{},0)`$. As $`s_k`$ is $`\gamma `$-projectizable and asymptotically holomorphic, we can choose a universal $`g_k`$-radius $`c`$ with $`|s_k^0|\gamma /2`$ on $`B_{g_k}(x,20c)`$. Also the sections $`s_{k,x}^{\text{ref}}`$ of Lemma 2.5 satisfy $`|s_{k,x}^{\text{ref}}|c_s`$ on $`B_{g_k}(x,20c)`$. Note that $`\varphi _k(B_{g_k}(x,20c))B_{g_{FS}}(\varphi _k(x),20Cc)`$.
2. We use the standard chart $`\mathrm{\Phi }_0`$ for $`^{2n+1}`$ around $`p=\varphi _k(x)=[1,0,\mathrm{},0]`$ to trivialize the ball $`B_{g_{FS}}(p,20Cc)`$. We may choose $`c`$ small enough so that $`\mathrm{\Phi }_0`$ is near an isometry, in the sense that
$$\frac{2}{3}|\mathrm{\Phi }_0(q)|d_{FS}(p,q)2|\mathrm{\Phi }_0(q)|.$$
for $`qB_{g_{FS}}(p,20Cc)`$. Also we require $`|\mathrm{\Phi }_0|2`$ in such ball. With respect to this trivialization the map $`\varphi _k`$ is given locally as
$`f_k=\mathrm{\Phi }_0\varphi _k:B_{g_k}(x,20c)`$ $``$ $`B(0,40Cc)`$
$`y`$ $``$ $`({\displaystyle \frac{s_k^1(y)}{s_k^0(y)}},\mathrm{},{\displaystyle \frac{s_k^{2n+1}(y)}{s_k^0(y)}}).`$
Clearly $`|f_k|2C`$ uniformly in $`k`$.
3. We can reduce $`c`$ so that, for any $`p`$, $`B_{g_{FS}}(p,20Cc)U_j`$ for some $`U_j`$. Therefore $`N`$ is defined in $`B(0,15Cc)`$ by a function $`f:B(0,15Cc)`$. Call $`Z=Z(f)`$ in such ball. The angle condition means that $`\mathrm{ker}df(z_1)`$, $`\mathrm{ker}df(z_2)`$ are close enough (say less than $`\pi /6`$) for $`z_1,z_2Z`$.
Let $`xM`$. In the case $`d(\varphi _k(x),N)2Cc`$, as we perform a small perturbation, say of norm $`\delta >0`$ such that $`d_{FS}(\varphi _k(x),\widehat{\varphi }_k(x))<\frac{1}{2}Cc`$, for all $`xM`$, there is still $`\frac{1}{2}Cc`$-transversality at a $`c`$-neighbourhood of $`x`$. So we are finished.
Suppose $`d(\varphi _k(x),N)<2Cc`$. Then take a point $`z_0B(0,4Cc)Z`$ which gives the minimum distance from $`0`$ to $`Z`$. If $`0Z`$, take $`v=(v_1,\mathrm{},v_{2n+1})^{2n+1}`$ a unitary vector in the direction of the complex line from $`0`$ to $`z_0`$. This vector is perpendicular to $`T_{z_0}Z`$. If $`0Z`$ then let $`v`$ be a unitary vector orthogonal to $`T_0Z`$. Therefore
(8)
$$df(z),v\frac{1}{2}|df(z)|$$
for any $`zZB(0,15Cc)`$, by the condition on the angle (taking $`ϵ>0`$ small enough).
Let $`r_0`$ with $`r_0v=z_0Z`$. We look for a function $`r_k=r_k(y):B_{g_k}(x,c)`$ such that $`r_k(x)=r_0`$ and
(9)
$$f(f_k^1(y)+r_kv_1,\mathrm{},f_k^{2n+1}(y)+r_kv_{2n+1})=0.$$
This corresponds to tracing a straight line from the image of the point $`yB_{g_k}(x,c)`$ to $`Z`$ with direction $`v`$. Such $`r_k`$ can be found with the use of the implicit function theorem applied to the function $`F:B_{g_k}(x,c)\times B(r_0,4Cc)`$ given as the left hand side of (9). This $`F`$ is well-defined since $`f`$ is defined on $`B(0,10Cc)\mathrm{\Phi }_0(U_j)`$. To guarantee the existence of $`r_k=r_k(y)`$ for all $`yB_{g_k}(x,c)`$ we have to check that
$$\left|\frac{_yF}{F/r_k}\right|=\left|\frac{df,f_k}{df,v}\right|4C,$$
which holds thanks to (8). This gives the existence of $`r_k`$ in the whole of the ball $`B_{g_k}(x,c)`$ as well as the bound $`|r_k|4C`$, and hence $`|r_k|8Cc`$.
Now our task will be to prove that $`r_k`$ is asymptotically holomorphic, so we change a geometrical transversality problem into a local one. For this let us compute $`\overline{}r_k`$. Recall that $`f_k`$ is asymptotically holomorphic and $`f`$ is holomorphic. Differentiate the equality $`f(f_k(y)+r_k(y)v)=0`$ to get
(10) $`0`$ $`=`$ $`\overline{}(f(f_k(y)+r_k(y)v))=f(z)(\overline{}f_k(y)+\overline{}r_k(y)v)=`$
$`=`$ $`O(k^{1/2})+df(z),v\overline{}r_k(y),`$
with $`z=f_k(y)+r_k(y)v`$. Using (8) we get that $`\overline{}r_k=O(k^{1/2})`$. We already know that $`|r_k|=O(1)`$. Differentiating (10) one easily obtains also that $`|\overline{}r_k|=O(k^{1/2})`$. So $`r_k`$ is asymptotically holomorphic. We shall achieve transversality for the function
$$h_k=r_k\frac{s_k^0}{s_{k,x}^{\text{ref}}}:B_{g_k}(x,c),$$
which is also asymptotically holomorphic.
Dividing $`h_k`$ by an appropriate constant, using the chart $`\mathrm{\Phi }_k`$ defined in Lemma 2.6 and scaling the coordinates by a universal constant, we obtain a function $`\stackrel{~}{h}_k`$ defined on $`B^+`$ satisfying the hypothesis of Proposition 2.10, for $`k`$ large enough. So going back to $`h_k`$ through universal constants, we find $`|w_k|<\delta `$ such that $`h_kw_k`$ is $`\eta `$-transverse to $`0`$ with $`\eta =c^{}\delta (\mathrm{log}(\delta ^1))^p`$.
Now we have a direction $`v`$ and a modulus $`w_k`$ for a perturbation. The perturbation we give is
$$\tau _{k,x}=(0,w_kv_1s_{k,x}^{\text{ref}},\mathrm{},w_kv_{2n+1}s_{k,x}^{\text{ref}}).$$
Let us look at the perturbed map $`\widehat{\varphi }_k=(s_k+\tau _{k,x})`$. It is asymptotically holomorphic and $`\gamma ^{}`$-projectizable and $`\gamma ^{}`$-generic of order $`n`$, for some $`\gamma ^{}>0`$, with $`|\tau _{k,x}|<c^{\prime \prime }\delta `$ (for $`\delta >0`$ small enough). Let us check that $`\widehat{\varphi }_k`$ is $`\eta `$-transverse to $`N`$ with $`\eta =c^{}\delta (\mathrm{log}(\delta ^1))^p`$ and $`c^{}`$ a constant depending only on $`c`$ and the asymptotically holomorphic bounds of $`s_k`$. With this, applying Proposition 2.8, the proof in this case is concluded. Only a little problem may appear, that the deformed embedding can become an immersion, but then an arbitrarily small perturbation solves the problem.
The $`h_k`$ associated to $`\widehat{\varphi }_k`$ is $`\widehat{h}_k=h_kw_k`$. The final point is to set up the relationship between the transversality of $`\widehat{h}_k`$ to $`0`$ and the transversality of $`\widehat{\varphi }_k`$ to $`N`$. Note that we have $`\widehat{r}_k=\widehat{h}_k\frac{s_{k,x}^{\text{ref}}}{s_k^0}`$, $`\widehat{f}_k=\mathrm{\Phi }_0\widehat{\varphi }_k`$ and $`\widehat{\pi }_k=\widehat{f}_k+\widehat{r}_kv=\pi _k`$.
Using that $`|s_{k,x}^{\text{ref}}/s_k^0|`$ is bounded above and below uniformly and that $`|(s_{k,x}^{\text{ref}}/s_k^0)|=O(1)`$, it is easy to prove that if $`\widehat{h}_k`$ is $`\eta `$-transverse to $`0`$ then $`\widehat{r}_k`$ is $`c_0\eta `$-transverse to $`0`$, for some universal constant $`c_0`$.
Let $`yB_{g_k}(x,c)`$. If $`|\widehat{r}_k(y)|c_0\eta `$ then $`d(\widehat{\varphi }_k(y),N)c_1\eta `$, for some universal constant $`c_1`$. Otherwise $`|\widehat{r}_k(y)|>c_0\eta `$. We shall use Lemma 3.8 for the subspaces $`U=(d\widehat{f}_k)_{}T_yM`$ and $`V=T_{\pi _k(y)}Z`$ of $`^{2n+1}`$. Let $`V^{}=[v]`$. The projections from $`U`$ to the summands of the decomposition $`^{2n+1}=VV^{}`$ are given respectively by $`g=d\pi _k(d\widehat{f}_k)^1`$ and $`h=vd\widehat{r}_k(d\widehat{f}_k)^1`$. This follows from $`d\pi _k=d\widehat{f}_k+d\widehat{r}_kv`$ which gives $`\text{Id}=d\pi _k(d\widehat{f}_k)^1vd\widehat{r}_k(d\widehat{f}_k)^1`$. The map $`h`$ has a right inverse of norm bounded by $`C^{}\eta ^1`$, for some universal constant $`C^{}`$ (here we use that $`\varphi _k`$ is generic of order $`n`$ and that the perturbations are small). It is easy to check that Lemma 3.8 is still valid when $`V`$ and $`V^{}`$ are almost orthogonal (and not just orthogonal), so we have
$$\mathrm{}_m((d\widehat{f}_k)_{}T_yM,T_{\pi _k(y)}Z)c_2\eta .$$
Push forward the distribution $`D_N`$ through the chart $`\mathrm{\Phi }_0`$ to a distribution $`D_Z`$ in $`B(0,15Cc)`$. Then there exists a constant $`C^{\prime \prime }`$ independent of $`k`$ such that
$$\mathrm{}_M(T_zZ,D_Z(z+\lambda v))<C^{\prime \prime }d(z+\lambda v,Z),$$
for $`zZ`$, $`\lambda `$ with $`|z|<14Cc`$, $`|\lambda |<Cc`$. Now use Proposition 3.5 to get
$$\mathrm{}_m((d\widehat{f}_k)_{}T_yM,D_Z(\widehat{f}_k(y)))>c_2\eta C^{\prime \prime }d(\widehat{f}_k(y),Z).$$
For $`d(\widehat{f}_k(y),Z)<c_2\eta /2C^{\prime \prime }`$ we get $`\mathrm{}_m((d\widehat{f}_k)_{}T_yM,D_Z(\widehat{f}_k(y)))>c_2\eta /2`$. Passing to the manifold we get $`\mathrm{}_m((d\widehat{\varphi }_k)_{}T_yM,D_N(\widehat{\varphi }_k(y))>c_2^{}\eta `$, whenever $`d(\widehat{\varphi }_k(y),N)<c_1^{}\eta `$, for some universal constants $`c_1^{}`$ and $`c_2^{}`$.
To achieve the solution when the codimension of $`N`$ is $`r>1`$, we follow the same ideas than in the precedent case. In this case $`f:B(0,15Cc)^r`$ and one chooses the point $`z_0`$ giving the minimum distance from $`0`$ to $`Z`$ which yields a vector $`v_1`$ orthogonal to $`Z`$ at $`z_0`$. Then one completes to an unitary basis $`(v_1,\mathrm{},v_r)`$ for the orthogonal to $`T_{z_0}Z`$. The function $`r_k:B_{g_k}(x,c)^r`$ is defined by the condition $`f(f_k+r_k^1v_1+\mathrm{}+r_k^rv_r)=0`$. The perturbation will be of the form
$$\tau _{k,x}=(0,w_k^1v_1^1s_{k,x}^{\text{ref}}+\mathrm{}+w_k^rv_r^1s_{k,x}^{\text{ref}},\mathrm{},w_k^1v_1^{2n+1}s_{k,x}^{\text{ref}}+\mathrm{}+w_k^rv_r^{2n+1}s_{k,x}^{\text{ref}}),$$
where $`v_i=(v_i^1,\mathrm{},v_i^{2n+1})`$, $`i=1,\mathrm{},r`$ and $`w_k=(w_k^1,\mathrm{},w_k^r)^r`$. The proof above works out in this case. $`\mathrm{}`$
## 4. Asymptotically holomorphic embeddings to grassmannians
Let $`(M,\omega )`$ be a symplectic manifold of integer class and let $`L`$ stand for the hermitian line bundle with a connection $``$ with curvature $`i\omega `$. Let $`E`$ be a rank $`r`$ hermitian bundle over $`M`$ endowed with an hermitian connection. Fix a compatible almost complex structure $`J`$ on $`M`$. In this Section we shall deal with the issue of constructing sequences of embeddings of $`M`$ into the grassmannian $`\text{Gr}(r,N)`$ which are asymptotically $`J`$-holomorphic in the sense of Definition 1.1. More specifically, we aim to prove the following result from which Theorem 1.4 follows.
###### Theorem 4.1.
Suppose $`N>n+r1`$ and $`r(Nr)>2n`$. Given an asymptotically $`J`$-holomorphic sequence of sections $`s_k`$ of the vector bundles $`^NEL^k`$ and $`\alpha >0`$ then there exists another sequence $`\sigma _k`$ verifying that:
1. $`|s_k\sigma _k|_{C^1,g_k}<\alpha `$.
2. $`\varphi _k=\text{Gr}(\sigma _k)`$ is an asymptotically holomorphic sequence of embeddings in $`\text{Gr}(r,N)`$ for $`k`$ large enough.
3. $`\varphi _k^{}𝒰=EL^k`$, where $`𝒰\text{Gr}(r,N)`$ is the universal rank $`r`$ bundle over the grassmannian.
Moreover given two asymptotically holomorphic sequences $`\varphi _k^0`$ and $`\varphi _k^1`$ of embeddings in $`\text{Gr}(r,N)`$ with respect to two compatible almost complex structures, then for $`k`$ large enough there exists an isotopy of asymptotically holomorphic embeddings $`\varphi _k^t`$ connecting $`\varphi _k^0`$ and $`\varphi _k^1`$.
### 4.1. Proof of main result.
First let us fix some notation. A point $`s\text{Gr}(r,N)`$ corresponds to an $`r`$-dimensional subspace $`V_s^N`$. Choosing a basis $`s_1,\mathrm{},s_r`$ for $`V_s`$, we denote
$$s=\left[\begin{array}{c}s_1\\ \mathrm{}\\ s_r\end{array}\right]=\left[\begin{array}{ccccc}s_{11}& s_{12}& \mathrm{}& s_{1N}& \\ \mathrm{}& & \mathrm{}& \mathrm{}& \\ s_{r1}& s_{r2}& \mathrm{}& s_{rN}& \end{array}\right].$$
This identifies $`s`$ as the equivalence class of $`r\times N`$ matrices of rank $`r`$ under the action of $`\text{GL}(r,)`$ on the left. The standard metric $`g_{Gr}`$ for $`\text{Gr}(r,N)`$ is the metric induced by the Fubini-Study metric $`g_{FS}`$ under the Plücker embedding \[GH78, Chapter 1, Section 5\]
$`\text{Gr}(r,N)`$ $``$ $`({\displaystyle ^r}^N)`$
$`\left[\begin{array}{c}s_1\\ \mathrm{}\\ s_r\end{array}\right]`$ $``$ $`s_1\mathrm{}s_r.`$
We proceed by steps to obtain asymptotically holomorphic embeddings.
###### Definition 4.2.
Let $`\gamma >0`$ and $`0lr`$. A sequence of asymptotically $`J`$-holomorphic sections $`s_k=(s_k^1,\mathrm{},s_k^N)`$ of the vector bundles $`^NEL^k`$ is said to be $`\gamma `$-grassmannizable of order $`l`$ if for all $`xM`$, $`|^ls_k(x)|>\gamma `$. It is $`\gamma `$-grassmannizable when it is $`\gamma `$-grassmannizable of order $`r`$. (Here $`s_k=(s_k^1,\mathrm{},s_k^N)`$ is interpreted as a morphism of bundles $`\underset{¯}{}^NEL^k`$ and $`^ls_k`$ is the corresponding $`l`$-fold wedge product.)
If we have the condition of $`\gamma `$-grassmannizability for a section $`s_k`$ then we obtain a morphism $`\varphi _k=\text{Gr}(s_k):M\text{Gr}(r,N)`$, called the grassmannization of $`s_k`$, as follows. At a point $`x`$ take a basis $`(e_1,\mathrm{},e_r)`$ for the fibre of $`E`$ at $`x`$. Then
$$\varphi _k(x)=[s_k^1(x),\mathrm{},s_k^N(x)]=\left[\begin{array}{ccccc}s_k^{11}& s_k^{12}& \mathrm{}& s_k^{1N}& \\ \mathrm{}& & \mathrm{}& \mathrm{}& \\ s_k^{r1}& s_k^{r2}& \mathrm{}& s_k^{rN}& \end{array}\right].$$
where $`s_k^i(x)=s_k^{1i}e_1+\mathrm{}+s_k^{ri}e_r`$. This is well-defined and independent of the chosen basis.
###### Definition 4.3.
Let $`\eta >0`$ and $`0ln`$. A sequence of asymptotically $`J`$-holomorphic $`\gamma `$-grassmannizable sections $`s_k`$ of vector bundles $`^NEL^k`$ is $`\eta `$-generic of order $`l`$, with $`\eta >0`$, if given $`\text{Gr}(s_k)`$ then for all $`xM`$, $`|^l\text{Gr}(s_k)(x)|_{g_k}>\eta `$.
In order to prove Theorem 4.1 we shall use the following auxiliar Proposition that will be proved in the following Subsections. Also we state the analogue of Lemma 2.15 which will be proved in Subsection 4.4.
###### Proposition 4.4.
Suppose $`N>n+r1`$ and $`r(Nr)>2n`$. Let $`s_k`$ be an asymptotically $`J`$-holomorphic sequence of sections of the vector bundles $`^NEL^k`$ and $`\alpha >0`$. Then there exists another sequence $`\sigma _k`$ verifying:
1. $`|s_k\sigma _k|_{C^1,g_k}<\alpha `$.
2. $`\sigma _k`$ is $`\gamma `$-grassmannizable and $`\gamma `$-generic of order $`n`$ for some $`\gamma >0`$.
Moreover, the result holds for one-parameter families of sections where the sections and the compatible almost complex structures depend continuously on $`t[0,1]`$.
###### Lemma 4.5.
Let $`\varphi _k:M\text{Gr}(r,N)`$ be a sequence of asymptotically holomorphic embeddings with $`\varphi _k^{}𝒰=EL^k`$. Then there exists a sequence of asymptotically holomorphic sections $`s_k`$ of $`^NEL^k`$, for $`k`$ large enough, which is $`\gamma `$-grassmannizable and $`\gamma `$-generic of order $`n`$, for some $`\gamma >0`$, such that $`\varphi _k=\text{Gr}(s_k)`$. The same holds for continuous one-parameter families of embeddings and compatible almost complex structures.
Proof of Theorem 4.1. Note that the last property is obvious by the construction. Let us begin with an asymptotically $`J`$-holomorphic sequence $`\sigma _k`$ of sections of the bundles $`^NEL^k`$ and perturb it using Proposition 4.4 to obtain an asymptotically holomorphic $`\gamma `$-grassmannizable and $`\gamma `$-generic of order $`n`$ sequence of sections $`s_k`$. The first property implies that $`\varphi _k=\text{Gr}(s_k)`$ is well-defined, the second that it is an immersion. To get an embedding we use that $`2\text{dim }M<\text{dim }\text{Gr}(r,N)=2r(Nr)`$ to find a generic $`C^p`$-perturbation of norm less than $`O(k^{1/2})`$ to get rid of the self-intersections and keeping the asymptotic holomorphicity, the grassmannizability and the genericity of order $`n`$. Now we only have to check that the sequence $`\varphi _k=\text{Gr}(s_k)`$ verifies the required conditions in Definition 1.1.
Choose a point $`xM`$ and trivialize $`E`$ in a neighborhood of $`x`$ by fixing an orthonormal basis $`e_1,\mathrm{},e_r`$. Now by a rotation with an element of $`U(N)`$ acting on $`^N`$ and an element of $`U(r)`$ acting on $`E`$, we can assure that
(12)
$$s_k(x)=\left(\begin{array}{ccccccc}s_k^{11}(x)& 0& \mathrm{}& & & \mathrm{}& 0\\ 0& s_k^{22}(x)& 0& \mathrm{}& & & 0\\ 0& \mathrm{}& \mathrm{}& 0& \mathrm{}& & 0\\ 0& \mathrm{}& & s_k^{rr}(x)& 0& \mathrm{}& 0\end{array}\right)$$
where $`s_k^{ij}`$ are sections of $`L^k`$. This corresponds to an isometric transformation of $`\text{Gr}(r,N)`$. The $`\gamma `$-grassmannizable property implies that $`|s_k^{11}\mathrm{}s_k^{rr}|\gamma `$. By the asymptotic holomorphicity bounds it is $`|s_k|=O(1)`$, so that $`|s_k^{ii}|\gamma /C`$, for some universal constant $`C`$. Therefore on a ball $`B_{g_k}(x,c)`$ of fixed universal radius $`c`$, the first $`r\times r`$ minor of $`s_k(y)`$ has an inverse of norm bounded by $`C^{}\gamma ^1`$, for some universal constant $`C^{}`$.
Let $`v_1,\mathrm{},v_N`$ be the canonical basis of $`^N`$. As $`\varphi _k(x)=\mathrm{\Pi }_0=\left[\begin{array}{c}v_1\\ \mathrm{}\\ v_r\end{array}\right]`$, we consider the standard local chart for $`\text{Gr}(r,N)`$ around $`\mathrm{\Pi }_0`$ for the open set $`U_0=\{\mathrm{\Pi }|\mathrm{\Pi }[v_{r+1},\mathrm{},v_N]=\{0\}\}`$, given by
$`\mathrm{\Phi }_0:U_0`$ $``$ $`^{r\times (Nr)}`$
$`\left[\begin{array}{ccc}s_{11}& \mathrm{}& s_{1N}\\ \mathrm{}& \mathrm{}& \mathrm{}\\ s_{r1}& \mathrm{}& s_{rN}\end{array}\right]`$ $``$ $`\left(\begin{array}{ccc}s_{11}& \mathrm{}& s_{1r}\\ \mathrm{}& \mathrm{}& \mathrm{}\\ s_{r1}& \mathrm{}& s_{rr}\end{array}\right)^1\left(\begin{array}{cccc}s_{1,r+1}& s_{1,r+2}& \mathrm{}& s_{1N}\\ \mathrm{}& & \mathrm{}& \mathrm{}\\ s_{r,r+1}& s_{r,r+2}& \mathrm{}& s_{rN}\end{array}\right)`$
It is easy to check that $`\mathrm{\Phi }_0`$ is an isometry at the point $`\mathrm{\Pi }_0`$.
The application $`f_k=\mathrm{\Phi }_0\varphi _k`$ is given by
$`f_k:B_{g_k}(x,c)`$ $``$ $`^{r\times (Nr)}`$
$`y`$ $``$ $`\left(\begin{array}{ccc}s_k^{11}(y)& \mathrm{}& s_k^{1r}(y)\\ \mathrm{}& \mathrm{}& \mathrm{}\\ s_k^{r1}(y)& \mathrm{}& s_k^{rr}(y)\end{array}\right)^1\left(\begin{array}{ccc}s_k^{1,r+1}(y)& \mathrm{}& s_k^{1N}(y)\\ \mathrm{}& \mathrm{}& \mathrm{}\\ s_k^{r,r+1}(y)& \mathrm{}& s_k^{rN}(y)\end{array}\right)`$
We can compute the bounds required in Definition 1.1 using $`f_k`$ instead of $`\varphi _k`$. Now the arguments in the proof of Theorem 2.11 carry over verbatim. For the isotopy result we use Lemma 4.5. $`\mathrm{}`$
### 4.2. Construction of $`\gamma `$-grassmannizable sections.
Our objective is to prove the following perturbation result:
###### Proposition 4.6.
Suppose $`N>n+r1`$. Let $`s_k`$ be an asymptotically $`J`$-holomorphic sequence of sections of the vector bundles $`^NEL^k`$ which is $`\gamma `$-grassmannizable of order $`l`$, for some $`\gamma >0`$. Then given $`\alpha >0`$, there exists an asymptotically $`J`$-holomorphic sequence of sections $`\sigma _k`$ verifying:
1. $`|s_k\sigma _k|_{C^1,g_k}<\alpha `$.
2. $`\sigma _k`$ is $`\eta `$-grassmannizable of order $`l+1`$ for some $`\eta >0`$.
Moreover, the result can be extended to continuous one-parameter families depending continuously of $`t[0,1]`$.
Proof. Again we use the globalization argument described in Proposition 2.8. Let us do the non-parametric case, the other one being a trivial extension by now. Define the local and $`C^0`$-open property $`P(ϵ,x)`$ as $`|^{l+1}s_k(x)|>ϵ`$. We only need to find for a point $`xM`$ a section $`\tau _{k,x}`$ with Gaussian decay away from $`x`$, assuring that $`s_k+\tau _{k,x}`$ verifies $`P(\eta ,y)`$ in a ball of universal $`g_k`$-radius $`c`$.
Choose a point $`xM`$. Fix an orthonormal basis $`e_1,\mathrm{},e_r`$ trivializing $`E`$ in a neighbourhood of $`x`$, so $`s_k`$ may be interpreted as a morphism $`^N^rL^k`$. By a rotation with an element of $`U(N)`$ on $`^N`$ and an element of $`U(r)`$ on $`E`$, we can assure that
$$s_k(x)=\left(\begin{array}{ccccccc}s_k^{11}(x)& 0& \mathrm{}& & & \mathrm{}& 0\\ 0& s_k^{22}(x)& 0& \mathrm{}& & & 0\\ 0& \mathrm{}& \mathrm{}& 0& \mathrm{}& & 0\\ 0& \mathrm{}& & s_k^{rr}(x)& 0& \mathrm{}& 0\end{array}\right)$$
with $`|s_k^{11}(x)\mathrm{}s_k^{ll}(x)|\gamma `$. So $`|s_k^{11}\mathrm{}s_k^{ll}|>\gamma /2`$ on a ball $`B_{g_k}(x,c)`$ of fixed radius $`c`$. Let $`s_{k,x}^{\text{ref}}`$ be the sections given by Lemma 2.5 and define $`\theta _k=s_k^{11}\mathrm{}s_k^{ll}s_{k,x}^{\text{ref}}`$. Clearly $`|\theta _k|>c_s\gamma /2`$ on $`B_{g_k}(x,c)`$. Consider the family of functions
$$M_k^p=s_k^{11}\mathrm{}s_k^{ll}s_k^{l+1,p},l+1pN.$$
These are components of $`^{l+1}s_k`$. If we perturb $`s_k`$ so that the norm of $`M_k=(M_k^{l+1},\mathrm{},M_k^N)`$ is bigger than $`\eta =c^{}\delta (\mathrm{log}(\delta ^1))^p`$ then we have finished. For this we define $`g_k=(g_k^{l+1},\mathrm{},g_k^N)=(\frac{M_k^{l+1}}{\theta _k},\mathrm{},\frac{M_k^N}{\theta _k})=(\frac{s_k^{l+1,l+1}}{s_{k,x}^{\text{ref}}},\mathrm{},\frac{s_k^{l+1,N}}{s_{k,x}^{\text{ref}}})`$. We obtain, scaling the coordinates by universal constants if necessary, $`g_k:B^+^{Nl}`$ which is asymptotically holomorphic. As $`n<Nl`$, we can find $`|w_k|<\delta `$ such that $`|g_kw_k|>\delta (\mathrm{log}(\delta ^1))^p`$. Then we obtain that $`|(M_k^{l+1}w_k^{l+1}\theta _k,\mathrm{},M_k^Nw_k^N\theta _k)|>\eta =c^{}\delta (\mathrm{log}(\delta ^1))^p`$, for some universal $`c^{}`$. This perturbation term is achieved by adding the section $`\tau _{k,x}=(0,\stackrel{(l)}{\mathrm{}},0,w_k^{l+1}e_{l+1}s_{k,x}^{\text{ref}},\mathrm{},w_k^Ne_{l+1}s_{k,x}^{\text{ref}})`$ of the bundles $`^NEL^k`$. This finishes the proof. $`\mathrm{}`$
###### Remark 4.7.
We cannot improve the condition $`N>n+r1`$ in Proposition 4.6. As we shall see in Section 5, we expect the locus of points of $`M`$ where the rank of $`s_k:\underset{¯}{}^NEL^k`$ is not maximum to have codimension $`Nr+1`$.
### 4.3. Inductive construction of sections $`\gamma `$-generic of order $`l`$
Now we study the problem of perturbing the sequence $`s_k`$ to achieve genericity of order $`n`$. The result to be proved is the following.
###### Proposition 4.8.
Suppose $`r(Nr)>2n`$. Let $`s_k`$ be an asymptotically $`J`$-holomorphic sequence of sections of the vector bundles $`^NEL^k`$, which is $`\gamma `$-grassmannizable and $`\gamma `$-generic of order $`l`$. Then given $`\alpha >0`$, there exists an asymptotically $`J`$-holomorphic sequence of sections $`\sigma _k`$ verifying:
1. $`|s_k\sigma _k|_{C^1,g_k}<\alpha `$.
2. $`\sigma _k`$ is $`\eta `$-generic of order $`l+1`$ for some $`\eta >0`$.
Moreover, this result can be extended to continuous one-parameter families of sections and almost complex structures.
Proof. Define the property $`P(ϵ,x)`$ for a section $`s_k`$ which is $`\gamma /2`$-grassmannizable and $`\gamma /2`$-generic of order $`l`$ as $`|^{l+1}\text{Gr}(s_k)(x)|>ϵ`$. A perturbation of our initial section verifies the hypothesis if we perturb by adding sections of $`C^1`$ norm smaller than $`\gamma /2C`$, $`C`$ some universal constant. For applying Proposition 2.8 we need to build, for $`0<\delta <\gamma /2Cc^{\prime \prime }`$, a local perturbation $`\tau _{k,x}`$ with $`|\tau _{k,x}|<c^{\prime \prime }\delta `$ and Gaussian decay with the property $`P(\eta ,y)`$ on $`B_{g_k}(x,c)`$ with $`\eta =c^{}\delta (\mathrm{log}(\delta ^1))^p`$.
Choose a point $`xM`$. By a rotation with an element of $`U(N)`$ acting on $`^N`$ and an element of $`U(r)`$ acting on $`E`$, we can assure that $`s_k(x)`$ is as in (12). By the $`\gamma `$-grassmannizability, $`|s_k^{11}(x)\mathrm{}s_k^{rr}(x)|\gamma `$. The asymptotically holomorphic bounds imply that $`|s_k|=O(1)`$, so that $`|s_k^{ii}(x)|\gamma /C`$ for some universal constant $`C`$. There is a fixed universal radius $`c`$ such that the first $`r\times r`$ minor of $`s_k(y)`$ has an inverse of norm bounded by $`C^{}\gamma ^1`$, for some universal constant $`C^{}`$, on $`B_{g_k}(x,c)`$. Then we can use the trivialization $`\mathrm{\Phi }_0`$ to define the applications
$`f_k:B_{g_k}(x,c)`$ $``$ $`^{r\times (Nr)}`$
$`y`$ $``$ $`\left(\begin{array}{ccc}s_k^{11}(y)& \mathrm{}& s_k^{1r}(y)\\ \mathrm{}& & \mathrm{}\\ s_k^{r1}(y)& \mathrm{}& s_k^{rr}(y)\end{array}\right)^1\left(\begin{array}{cccc}s_k^{1,r+1}(y)& \mathrm{}& s_k^{1N}(y)& \\ \mathrm{}& & \mathrm{}& \\ s_k^{r,r+1}(y)& \mathrm{}& s_k^{rN}(y)& \end{array}\right)`$
Now consider the sections $`s_{k,x}^{\text{ref}}`$ of Lemma 2.5. We define the applications
$`\stackrel{~}{f}_k:B_{g_k}(x,c)`$ $``$ $`^{r\times (Nr)}`$
$`y`$ $``$ $`{\displaystyle \frac{1}{s_{k,x}^{\text{ref}}(y)}}\left(\begin{array}{cccc}s_k^{1,r+1}(y)& s_k^{1,r+2}(y)& \mathrm{}& s_k^{1N}(y)\\ \mathrm{}& & \mathrm{}& \\ s_k^{r,r+1}(y)& s_k^{r,r+2}(y)& \mathrm{}& s_k^{rN}(y)\end{array}\right)`$
Clearly $`f_k=\mathrm{\Psi }\stackrel{~}{f}_k`$ where $`\mathrm{\Psi }:B_{g_k}(x,c)\text{GL}(r,)`$ satisfies $`|\mathrm{\Psi }|=O(1)`$, $`|\mathrm{\Psi }^1|=O(1)`$, $`|\mathrm{\Psi }|=O(1)`$ and $`|\mathrm{\Psi }^1|=O(1)`$. Therefore it is enough to get a perturbation which has $`|^{l+1}\stackrel{~}{f}_k|>\eta `$ on $`B_{g_k}(x,c)`$.
Spreading out the entries of the matrix $`\stackrel{~}{f}_k`$ in one row we can write $`\stackrel{~}{f}_k(y)=(\stackrel{~}{f}_k^{11}(y),\mathrm{},\stackrel{~}{f}_k^{r,Nr}(y))`$. Using the local forms $`dz_k^1,\mathrm{},dz_k^n`$, we may write
$$\stackrel{~}{f}_k=(u_k^{111}dz_k^1+u_k^{112}dz_k^2+\mathrm{}+u_k^{11n}dz_k^n,\mathrm{},u_k^{r,Nr,1}dz_k^1+\mathrm{}+u_k^{r,Nr,n}dz_k^n),$$
for some $`u_k^{ijl}`$. Using a unitary transformation of $`U(n)`$ on the complex Darboux coordinate chart and relabeling horizontally the coordinates, we can suppose that
(17)
$$\stackrel{~}{f}_k(x)=\left(\begin{array}{ccccccc}u_k^{11}(x)& & \mathrm{}& & & \mathrm{}& \\ 0& u_k^{22}(x)& & \mathrm{}& & & \\ 0& \mathrm{}& \mathrm{}& & \mathrm{}& & \\ 0& \mathrm{}& 0& u_k^{nn}(x)& & \mathrm{}& \end{array}\right),$$
where $`|u_k^{11}(x)\mathrm{}u_k^{ll}(x)|>\gamma /C_0`$, $`C_0`$ a universal constant. The relabeling is given by a function $`\alpha \{1,\mathrm{},r(Nr)\}(i(\alpha ),j(\alpha ))\{1,\mathrm{},r\}\times \{1,\mathrm{},Nr\}`$. Shrinking $`c`$ if necessary we can assure that $`|(\stackrel{~}{f}_k^1\mathrm{}\stackrel{~}{f}_k^l)_{dz_k^1\mathrm{}dz_k^l}|>\gamma /2C_0`$ for all the points of the ball $`B_{g_k}(x,c)`$. Now we construct the $`(l+1)`$-form
$$\theta _k(y)=(\stackrel{~}{f}_k^1\mathrm{}\stackrel{~}{f}_k^l)_{dz_k^1\mathrm{}dz_k^l}dz_k^{l+1}.$$
and the family of $`(l+1)`$-forms
$$M_k^p=(\stackrel{~}{f}_k^1\mathrm{}\stackrel{~}{f}_k^l\stackrel{~}{f}_k^p)_{dz_k^1\mathrm{}dz_k^ldz_k^{l+1}},l+1pr(Nr),$$
which are components of $`^{l+1}\stackrel{~}{f}_k`$. If we perturb so that the norm of $`M_k=(M_k^{l+1},\mathrm{},M_k^{r(Nr)})`$ gets bigger than $`\eta =c^{}\delta (\mathrm{log}(\delta ^1))^p`$ then we are done. We define the following sequence of asymptotically holomorphic applications: $`h_k=(\frac{M_k^{l+1}}{\theta _k},\mathrm{},\frac{M_k^{r(Nr)}}{\theta _k})`$. So we obtain, scaling the coordinates by universal constants if necessary, $`h_k:B^+^{r(Nr)l}`$ which is asymptotically holomorphic. As $`n<r(Nr)l`$ we can find $`|w_k|<\delta `$ such that $`|h_kw_k|>\delta (\mathrm{log}(\delta ^1))^p`$. Thus $`|(M_k^{l+1}w_k^{l+1}\theta _k,\mathrm{},M_k^{r(Nr)}w_k^{r(Nr)}\theta _k)|>\eta =c^{}\delta (\mathrm{log}(\delta ^1))^p`$. The perturbation term $`(w_k^{l+1}\theta _k,\mathrm{},w_k^{r(Nr)}\theta _k)`$ is achieved by adding the section
$$\tau _{k,x}=(0,\stackrel{(r)}{\mathrm{}},0,\underset{j(\alpha )=r+1,\alpha >l}{}w_k^\alpha z_{l+1}e_{i(\alpha )}s_{k,x}^{\text{ref}},\mathrm{},\underset{j(\alpha )=N,\alpha >l}{}w_k^\alpha z_{l+1}e_{i(\alpha )}s_{k,x}^{\text{ref}}).$$
This finishes the proof in the non-parametric case. The other case is left to the reader. $`\mathrm{}`$
### 4.4. Lifting asymptotically holomorphic embeddings in grassmannians
This Subsection is devoted to a proof of Lemma 4.5, which states that any asymptotically holomorphic embedding into a grassmannian is of the form provided by Theorem 4.1.
Proof of Lemma 4.5. Suppose that we have a sequence of $`\gamma `$-asymptotically holomorphic embeddings $`\varphi _k:M\text{Gr}(r,N)`$, for some $`\gamma >0`$, with $`\varphi _k^{}𝒰=EL^k`$, where $`𝒰`$ is the universal rank $`r`$ bundle over the grassmannian. The dual of $`𝒰`$ is given by
$$𝒰^{}=\{(\mathrm{\Pi },s)|s\mathrm{\Pi }\}\text{Gr}(r,N)\times ^N=\underset{¯}{}^N,$$
interpreted as a sub-bundle of the trivial bundle $`\underset{¯}{}^N`$. We consider the sequence of hermitian bundles, $`E_k=\varphi _k^{}𝒰^{}EL^k=\mathrm{End}E^NEL^k`$. We look for sequences of sections $`s_k`$ of $`E_k`$ which are $`\sigma `$-grassmannizable of order $`n`$ such that they are asymptotically holomorphic when considered as sections of $`^NEL^k`$. Let $`S_k^l=\mathrm{Tr}(^ls_k)`$, which is an asymptotically holomorphic sequence of sections of the trivial vector bundle $`\underset{¯}{}`$. We want to prove that $`|S_k^r|\sigma `$ for $`k`$ large. We shall prove that we can find sequences $`s_k`$ with $`|S_k^l|\eta _l`$, for some $`\eta _l>0`$, by induction on $`l`$.
Suppose that $`s_k`$ is an asymptotically holomorphic sequence of sections of $`E_k`$ such that $`|S_k^l|\gamma `$. Let $`P(ϵ,x)`$ be the $`C^1`$-open property for sequences of sections $`s_k`$ of $`E_k`$ given as $`S_k^{l+1}=\mathrm{Tr}(^{l+1}s_k)`$ is $`ϵ`$-transverse to $`0`$ at $`x`$.
Let $`xM`$. We want to find a local perturbation with Gaussian decay obtaining the property $`P(\eta ,y)`$ in a ball of universal $`g_k`$-radius $`c`$ around $`x`$. For this, define the local sections $`\sigma _k`$ of $`\varphi _k^{}𝒰^{}E^NE`$ as follows. Locally, $`\sigma _k`$ is a map
$$\sigma _k:B_{g_k}(x,c)^N^r,$$
such that for $`yB_{g_k}(x,c)`$, $`\sigma _k(y)`$ is a $`N\times r`$ matrix, i.e. a linear map $`\sigma _k(y):^N^r`$. The point $`\varphi _k(y)\text{Gr}(r,N)`$ corresponds to the image of the embedding $`\sigma _k(y)^T:^r^N`$. Note that one may identify the tangent space $`T_{\varphi _k(y)}\text{Gr}(r,N)`$ to the set of linear maps $`^N^r`$ which are zero on $`\mathrm{im}(\sigma _k(y)^T)`$, i.e. maps $`\phi `$ such that $`\phi \sigma _k(y)^{}=0`$. With this, $`\sigma _k=\varphi _k+(\sigma _k\sigma _k^{})\sigma _k`$. So it is natural to require $`(_r\sigma _k(y))\sigma _k(y)^{}=0`$, for any $`yB_{g_k}(x,c)`$, where $`r`$ is the radial vector field from $`x`$. We fix $`\sigma _k(x)`$ satisfying $`\sigma _k(x)\sigma _k(x)^{}=\text{I}`$. This determines $`\sigma _k`$ uniquely. The following bounds are proved as in Subsection 2.5,
$`\sigma _k(y)\sigma _k(y)^{}=\text{I},|\sigma _k(y)|=O(1),|\sigma _k(y)|=O(1+d_k(x,y)),`$
$`|\overline{}\sigma _k(y)|=O(k^{1/2}(1+d_k(x,y))),|\overline{}\sigma _k(y)|=O(k^{1/2}(1+d_k(x,y))).`$
Trivialize $`E`$ in a ball $`B_{g_k}(x,c)`$, so that $`s_k/s_{k,x}^{\text{ref}}`$ can be considered as an application $`B_{g_k}(x,c)^{r\times N}`$. Define the application
$$f_k=\frac{s_k\sigma _k^{}}{s_{k,x}^{\text{ref}}}:B_{g_k}(x,c)^{r\times r},$$
so that $`f_k\sigma _k=s_k/s_{k,x}^{\text{ref}}`$. Then $`f_k`$ is asymptotically holomorphic and we may check property $`P(\eta ,y)`$ for $`f_k`$ instead of $`s_k`$. Let $`F_i=\mathrm{Tr}(^if_k)`$, so that $`|F_l|C\gamma `$ for some universal constant $`C`$. For any $`w`$ we have
$$\mathrm{Tr}(^{l+1}(f_kw\text{I}))=F_{l+1}w(nl)F_l+w^2\left(\genfrac{}{}{0pt}{}{nl+1}{2}\right)F_{l1}+\mathrm{}+(w)^{l+1}\left(\genfrac{}{}{0pt}{}{n}{l+1}\right)F_0$$
By the standard argument, we may obtain $`|w|<\delta `$ such that $`\frac{F_{l+1}}{F_l}w`$ is $`\eta `$-transverse to $`0`$, with $`\eta =\delta (\mathrm{log}(\delta ^1))^p`$, in $`B_{g_k}(x,c)`$. Then it is easy to see that $`\mathrm{Tr}(^{l+1}(f_k\frac{w}{nl}\text{I})`$ is $`c^{}\eta `$-transverse to $`0`$, where $`c^{}`$ is a universal constant. The perturbation
$$\tau _{k,x}=\frac{w}{nl}\sigma _ks_{k,x}^{\text{ref}}$$
is a sequence of sections of $`E_k=\varphi _k^{}𝒰^{}EL^k`$, with Gaussian decay such that $`|\tau _{k,x}|<c^{\prime \prime }\delta `$ and $`s_k+\tau _{k,x}`$ satisfies $`P(\eta ,y)`$ for $`yB_{g_k}(x,c)`$, with $`\eta =c^{}\delta (\mathrm{log}(\delta ^1))^p`$. By Proposition 2.8, there exists an asymptotically holomorphic sequence of sections of $`E_k`$, which we denote by $`s_k`$ again, such that $`S_k^{l+1}=\mathrm{Tr}(^{l+1}s_k)`$ is $`\eta `$-transverse to $`0`$, for some $`\eta >0`$. For $`k`$ large enough, the zeroes of $`S_k^{l+1}`$ is a symplectic submanifold representing the trivial homology class, hence the empty set. So $`|S_k^{l+1}|\eta `$. This completes the proof. The extension to the one-parameter case is trivial. $`\mathrm{}`$
### 4.5. Zero sets of vector bundles
Following the ideas of Subsection 3.3 and using Proposition 3.10 we can prove the following two results
###### Theorem 4.9.
Given $`\varphi _k=\text{Gr}(s_k)`$, where $`s_k`$ is a sequence of asymptotically holomorphic sections of $`^NEL^k`$, which are $`\gamma `$-grassmannizable and $`\gamma `$-generic of order $`n`$, for some $`\gamma >0`$. Fix a holomorphic submanifold $`V`$ in $`\text{Gr}(r,N)`$. Then for any $`\alpha >0`$ there exists a sequence of asymptotically holomorphic sections $`\sigma _k`$ of $`^NEL^k`$ such that
1. $`|\sigma _ks_k|_{g_k,C^1}<\alpha `$.
2. $`\text{Gr}(\sigma _k)`$ is an $`\eta `$-asymptotically holomorphic embedding in $`\text{Gr}(r,N)`$ which is $`ϵ`$-transverse to $`V`$, with $`\eta >0`$ and $`ϵ>0`$ independent of $`k`$. In the case $`\text{dim }M+\text{dim }V<2r(Nr)`$ we have that $`d_{Gr}(\varphi _k(M),V))>ϵ`$, for $`k`$ large enough.
Moreover the result can be extended to one-parameter continuous families of complex submanifolds $`(V_t)_{t[0,1]}`$, taking in this case as starting point a continuous family $`\varphi _{t,k}=\text{Gr}(s_{t,k})`$, where $`s_{t,k}`$ is a continuous family of $`\gamma `$-grassmannizable and $`\gamma `$-generic of order $`n`$ sequences asymptotically $`J_t`$-holomorphic.
Proof. The proof is similar to that of Theorem 3.12. We just briefly point out the differences. For simplicity we suppose that the codimension of $`V`$ is $`1`$.
For $`xM`$, we may suppose that $`s_k(x)`$ is as in (12). We use the chart $`\mathrm{\Phi }_0`$ to get the local maps
$`f_k=\mathrm{\Phi }_0\varphi _k:B_{g_k}(x,c)`$ $``$ $`^{r\times (Nr)}`$
$`y`$ $``$ $`\left(\begin{array}{ccc}s_k^{11}(y)& \mathrm{}& s_k^{1r}(y)\\ \mathrm{}& & \mathrm{}\\ s_k^{r1}(y)& \mathrm{}& s_k^{rr}(y)\end{array}\right)^1\left(\begin{array}{cccc}s_k^{1,r+1}(y)& \mathrm{}& s_k^{1N}(y)& \\ \mathrm{}& & \mathrm{}& \\ s_k^{r,r+1}(y)& \mathrm{}& s_k^{rN}(y)& \end{array}\right)`$
This time we have a vector $`v^{r\times (Nr)}`$. We define the functions $`h_k:B_{g_k}(x,c)`$ by the condition
$$f\left(f_k+r_ks_{k,x}^{\text{ref}}\left(\begin{array}{ccc}s_k^{11}& \mathrm{}& s_k^{1r}\\ \mathrm{}& \mathrm{}& \mathrm{}\\ s_k^{r1}& \mathrm{}& s_k^{rr}\end{array}\right)^1\left(\begin{array}{ccc}s_k^{11}(x)& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& s_k^{rr}(x)\end{array}\right)\left(\begin{array}{ccc}v^{11}& \mathrm{}& v^{1,Nr}\\ \mathrm{}& \mathrm{}& \mathrm{}\\ v^{r1}& \mathrm{}& v^{r,Nr}\end{array}\right)\right)=0,$$
and prove that they are asymptotically holomorphic. Then we find $`|w_k|<\delta `$ such that $`h_kw_k`$ is $`\eta `$-transverse to $`0`$ with $`\eta =c^{}\delta (\mathrm{log}(\delta ^1))^p`$. Finally the perturbation will be
$$\tau _{k,x}=\left(\begin{array}{cccccc}0& \mathrm{}& 0& w_kv^{11}s_{k,x}^{\text{ref}}& \mathrm{}& w_kv^{1,Nr}s_{k,x}^{\text{ref}}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& 0& w_kv^{r1}s_{k,x}^{\text{ref}}& \mathrm{}& w_kv^{r,Nr}s_{k,x}^{\text{ref}}\end{array}\right).$$
The arguments run parallel to those in the proof of Theorem 3.12, although the constants have to be arranged suitably, but we leave this task to the careful reader. $`\mathrm{}`$
We call universal planes to the zero sets of algebraic sections transverse to zero of the universal bundle $`𝒰`$ over the grassmannian $`\text{Gr}(r,N)`$. Now we can deduce the main result of \[Au97\].
###### Corollary 4.10.
Let $`(M,\omega )`$ be a compact symplectic manifold of integer class. Let $`E`$ be a hermitian vector bundle over $`M`$. Then for $`k`$ large enough there exist symplectic submanifolds obtained as zero sets of the bundles $`EL^k`$. Moreover, perhaps by increasing $`k`$, we can assure that all the symplectic submanifolds constructed as transverse intersections of asymptotically holomorphic sequences with a fixed universal plane are isotopic. The isotopy can be made through symplectomorphisms.
The proof follows the steps of the proof of Corollary 3.13. Remark also that the result is a corollary of Theorem 5.4 to be proved in Section 5.
## 5. Determinantal submanifolds of closed symplectic manifolds
Let $`(M,\omega )`$ be a symplectic $`4`$-manifold of integer class, endowed with a compatible almost complex structure. Let $`E`$ and $`F`$ be two vector bundles of ranks $`r_e`$ and $`r_f`$, respectively. Recall that for any morphism $`\phi :EF`$ we have defined in Definition 1.5 the $`r`$-determinantal set as
$$\mathrm{\Sigma }^r(\phi )=\{xM|\text{rank}\phi _x=r\}.$$
We want to prove Theorem 1.6, which allows to construct $`\mathrm{\Sigma }^r(\phi )`$ as a symplectic submanifold, after twisting $`E`$ and $`F`$ with large powers of $`L`$. The solution to this problem goes through embedding $`M`$ in a product of two grassmannians and cutting its image with suitable “generalized Schur cycles”. We shall do this in next Section.
###### Remark 5.1.
A direct approach to proving Theorem 1.6 consists on reducing it to Auroux’ case by taking the $`r`$-fold wedge product of $`\phi _k`$,
$`{\displaystyle ^r}\phi _k:{\displaystyle ^r}E(L^{})^k`$ $``$ $`{\displaystyle ^r}FL^k`$
$`s_1\mathrm{}s_r`$ $``$ $`\phi _k(s_1)\mathrm{}\phi _k(s_r).`$
So the zero set of $`^r\phi _k`$ is generically a stratified submanifold $`\mathrm{\Sigma }^0(\phi _k)\mathrm{}\mathrm{\Sigma }^r(\phi _k)`$. If we suppose that $`\phi _k`$ is an asymptotically $`J`$-holomorphic sequence of sections of the bundle $`E^{}FL^{2k}`$, one could try to use Donaldson’s techniques to obtain a new sequence of sections transverse in an adequate sense to assure the symplecticity. The following example shows the main obstacle to this approach. Take a symplectic $`4`$-manifold in the hypothesis of Theorem 1.6 with two hermitian vector bundles $`E`$ and $`F`$ of rank $`2`$. Using Auroux’ techniques we can assure that the zero sets of $`\phi _k`$ are $`\eta `$-transverse to $`0`$, for some $`\eta >0`$. When we go to $`^2\phi _k`$, the condition to be satisfied is
$$|\overline{}^2\phi _k|<|^2\phi _k|.$$
However, at any $`xZ(\phi _k)`$ we obtain $`|\overline{}^2\phi _k(x)|=|^2\phi _k(x)|=0`$, so we cannot impose a global transversality property for the section $`^2\phi _k`$. This case is very similar to that in \[Do99\] and can be treated with an “ad hoc” argument, but more general cases do not admit a treatment based on the use of normal forms of the singularities, because for higher dimensions the problem becomes intractable \[Ar82\].
### 5.1. Bigrassmannian embeddings
The idea to prove Theorem 1.6 is based in the following observations. Choose two sequences of sections $`s_k^e`$ and $`s_k^f`$ of the bundles $`^NE^{}L^k`$ and $`^NFL^k`$ respectively, which are $`\gamma `$-grassmannizable and $`\gamma `$-generic of order $`n`$, for some $`\gamma >0`$, providing by Theorem 4.1, asymptotically holomorphic sequences of embeddings $`\text{Gr}(s_k^e)`$ and $`\text{Gr}(s_k^f)`$ of $`M`$ in $`\text{Gr}(r_e,N)`$ and $`\text{Gr}(r_f,N)`$, respectively, for $`N`$ a large integer number.
Performing the cartesian product we obtain an asymptotically holomorphic sequence of embeddings of $`M`$ into the bigrassmannian $`\text{Bi}(r_e,r_f,N)=\text{Gr}(r_e,N)\times \text{Gr}(r_f,N)`$,
$$\varphi _k=\text{Gr}(s_k^e)\times \text{Gr}(s_k^f):M\text{Gr}(r_e,N)\times \text{Gr}(r_f,N)=\text{Bi}(r_e,r_f,N).$$
Let $`𝒰_e`$ and $`𝒰_f`$ be the universal bundles over $`\text{Gr}(r_e,N)`$ and $`\text{Gr}(r_f,N)`$ respectively, which are very ample. Define $`\pi _e:\text{Bi}(r_e,r_f,N)\text{Gr}(r_e,N)`$ as the projection onto the first factor (and analogously $`\pi _f`$). Therefore $`𝒰_{ef}=\pi _e^{}(𝒰_e)\pi _f^{}(𝒰_f)=𝒰_{ef}`$ is very ample on $`\text{Bi}(r_e,r_f,N)`$. Recall that $`\text{Gr}(s_k^e)^{}(𝒰_e)=E^{}L^k`$ and $`\text{Gr}(s_k^f)^{}(𝒰_f)=FL^k`$. Then $`\varphi _k^{}𝒰_{ef}=E^{}FL^{2k}`$. $`𝒰_{ef}`$ has a holomorphic section $`s`$ verifying that:
1. $`D_r=\mathrm{\Sigma }^r(s)`$ is an open complex submanifold in $`\text{Bi}(r_e,r_f,N)`$.
2. $`\text{codim}_{}D_r=(r_er)(r_fr)`$.
If we assure that, for each $`r`$, $`\varphi _k`$ is transverse to $`D_r`$ with an angle $`ϵ>0`$ independent of $`k`$, we have finished the proof of Theorem 1.6 by Proposition 3.10. This is carried out as follows.
###### Lemma 5.2.
Let $`\varphi _k:M\text{Bi}(r_e,r_f,N)`$ be a $`\gamma `$-asymptotically holomorphic sequence of embeddings. Suppose that $`\varphi _k`$ is $`\sigma `$-transverse to $`D_r`$. Then there exists $`ϵ>0`$, depending only on $`\gamma `$, $`\sigma `$ and the universal bounds of the derivatives of the sequence, such that $`\varphi _k`$ is $`\sigma /2`$-transverse to $`D_r^{}`$, $`r^{}>r`$, when we restrict to an $`ϵ`$-neighborhood of $`D_r`$.
In other words we do not have to care about the behaviour of the angle near the border of the strata.
Proof. Choose a point $`xD_r\varphi _k(M)`$. Recall that by $`\sigma `$-transversality, the minimum angle between $`T_xD_r`$ and $`T_x\varphi _k(M)`$ is greater than $`\sigma `$. We trivialize $`\text{Bi}(r_e,r_f,N)`$ by a chart $`\mathrm{\Phi }_0`$ defined as the cartesian product of two standard charts in the grassmannians, which is an isometry at the origin and verifies that $`\mathrm{\Phi }_0(x)=0`$, namely,
$$\mathrm{\Phi }_0:\text{Bi}(r_e,r_f,N)^{r_e(Nr_e)}\times ^{r_f(Nr_f)}.$$
Since $`D_r`$ is contained in the closure of $`D_r^{}`$, we have
(19)
$$|y|<\delta \mathrm{}_M(T_0\mathrm{\Phi }_0(D_r),T_y\mathrm{\Phi }_0(D_r^{}))<c_D\delta ,yB(0,c_u)\mathrm{\Phi }_0(D_r^{}).$$
The angles are measured with respect to the standard euclidean metric which is close to that induced by the bigrassmannian if we choose $`c_u`$ small enough. Here $`c_D`$ is universal. Also by the asymptotic holomorphicity bounds of $`\varphi _k`$ we know that
$`|y|<\delta \mathrm{}_M(T_0\mathrm{\Phi }_0(\varphi _k(M)),T_y\mathrm{\Phi }_0(\varphi _k(M)))<c_\varphi \delta ,`$
(20) $`yB(0,c_u){\displaystyle \mathrm{\Phi }_0(\varphi _k(M))},`$
where $`c_\varphi `$ is universal. Now Proposition 3.5 says that
$$\mathrm{}_m(T_0\mathrm{\Phi }_0(D_r),T_0\mathrm{\Phi }_0(\varphi _k(M)))\mathrm{}_M(T_0\mathrm{\Phi }_0(D_r),T_y\mathrm{\Phi }_0(D_r^{}))+$$
$$+\mathrm{}_m(T_y\mathrm{\Phi }_0(D_r^{}),T_y\mathrm{\Phi }_0(\varphi _k(M)))+\mathrm{}_M(T_y\mathrm{\Phi }_0(\varphi _k(M)),T_0\mathrm{\Phi }_0(\varphi _k(M))).$$
Using inequalities (19) and (20) and remembering that all the angles have to be measured with respect to the bigrassmannian metric (which is related to the standard metric in the ball $`B(0,c_u)`$ by non zero universal constants), we get the required result. $`\mathrm{}`$
With Lemma 5.2 the proof of Theorem 1.6 reduces to the following result, whose proof is similar to that of Theorem 4.9.
###### Proposition 5.3.
Let $`s_k^e`$ and $`s_k^f`$ be two asymptotically holomorphic sequences of the vector bundles $`E^{}L^k`$ and $`FL^k`$ which are $`\gamma `$-grassmannizable and $`\gamma `$-generic of order $`n`$, defining so an asymptotically holomorphic embedding in $`\text{Bi}(r_e,r_f,N)`$. Fix an algebraic open submanifold $`V`$ in $`\text{Bi}(r_e,r_f,N)`$ with compactification $`\overline{V}=VW`$. Then for any $`ϵ,\alpha >0`$, there exist $`\eta >0`$ and two asymptotically holomorphic sequences $`\sigma _k^e`$ and $`\sigma _k^f`$ of sections of the vector bundles $`E^{}L^k`$ and $`FL^k`$ respectively, verifying:
1. $`|\sigma _k^es_k^e|_{g_k,C^1}<\alpha `$ and $`|\sigma _k^fs_k^f|_{g_k,C^1}<\alpha `$.
2. $`\varphi _k=\text{Gr}(\sigma _k^e)\times \text{Gr}(\sigma _k^f)`$ is a sequence of $`\eta `$-asymptotically holomorphic embeddings in $`\text{Bi}(r_e,r_f,N)`$.
3. Denoting by $`V_ϵ^{}`$ the compact submanifold of $`V`$ obtained by removing an $`ϵ`$-neighborhood of $`W`$, we obtain that $`\varphi _k`$ is $`\eta `$-transverse to $`V_ϵ^{}`$.
Moreover the result can be extended to continuous one-parameter families of sections $`(s_{k,t}^e)_{t[0,1]}`$ and $`(s_{k,t}^f)_{t[0,1]}`$ providing embeddings to the bigrassmannian and to continuous one-parameter families of open submanifolds $`V_t`$. Thus we obtain continuous families of sequences $`\sigma _{k,t}^e`$ and $`\sigma _{k,t}^f`$ verifying the required conditions. $`\mathrm{}`$
### 5.2. Dependence loci of sections of a vector bundle
Suppose that $`E`$ is an hermitian vector bundle of rank $`n`$ and consider $`s_1,\mathrm{},s_m`$ sections of $`E`$. Then we can interpret $`s=(s_1,\mathrm{},s_m)`$ as a morphism of bundles $`s:\underset{¯}{}^mE`$. The $`r`$-determinantal set of $`s`$ is
$$\mathrm{\Sigma }^r(s)=\{xM|\text{dim }[s_1(x),\mathrm{},s_m(x)]=r\},$$
and it is called the $`r`$-dependence locus of the sections $`s_1,\mathrm{},s_m`$.
###### Theorem 5.4.
Let $`(M,\omega )`$ be a closed symplectic manifold of integer class and let $`E`$ be a rank $`n`$ hermitian vector bundle. Then, for $`k`$ large enough, there exist $`s_k=(s_k^1,\mathrm{},s_k^m)`$ sections of $`^mE`$ such that
1. $`\mathrm{\Sigma }^r(s_k)`$ is an open symplectic submanifold of $`M`$.
2. $`\text{codim }\mathrm{\Sigma }^r(s_k)=2(mr)(nr)`$. The set of manifolds $`\{\mathrm{\Sigma }^r(s_k)\}_r`$ constitutes a stratified submanifold.
Moreover, any two stratified submanifolds constructed by the process in the proof below are isotopic.
Proof. The proof is similar to the arguments developed in Subsection 5.1. Let $`𝒰`$ be the universal bundle over $`\text{Gr}(n,N)`$ and consider $`m`$ holomorphic sections $`s_1,\mathrm{},s_m`$ verifying that:
1. $`D_r=\mathrm{\Sigma }^r(s)`$ is an open complex submanifold in $`\text{Gr}(n,N)`$.
2. $`\text{codim}_{}D_r=(mr)(nr)`$.
Now we choose a sequence of asymptotically holomorphic embeddings $`\varphi _k:M\text{Gr}(n,N)`$ such that $`\varphi _k^{}𝒰=EL^k`$. If we assure that, for each $`r`$, $`\varphi _k`$ is transverse to $`D_r`$ with an angle $`ϵ>0`$ independent of $`k`$, we have finished the proof because of Proposition 3.10. But we may perturb $`\varphi _k`$ by using analogues of Lemma 5.2 and Proposition 5.3 for the case of just one grassmannian. $`\mathrm{}`$
### 5.3. Homology and homotopy groups of the determinantal submanifolds
In this Subsection we prove a result concerning the topology of smooth determinantal submanifolds analogous to Proposition 39 in \[Do96\] (symplectic Lefschetz hyperplane theorem) and Proposition 2 in \[Au97\]. The main result is
###### Proposition 5.5.
Let $`E,F`$ be vector bundles of ranks $`r_e`$, $`r_f`$, respectively, over a closed symplectic manifold $`(M,\omega )`$ of integer class and let $`D_r^k`$ be a sequence of determinantal submanifolds constructed, by using the vector bundles $`E(L^{})^k`$ and $`FL^k`$, as a transverse intersection of an asymptotically holomorphic sequence of embeddings in $`\text{Bi}(r_e,r_f,N)`$ with the determinantal varieties of a fixed generic section $`s`$ of the universal bundle $`𝒰_{ef}`$ over $`\text{Bi}(r_e,r_f,N)`$. Assume that the stratified determinantal submanifold has only one stratum $`D_r^k`$. Then the inclusion $`i:D_r^kM`$ induces, for $`k`$ large enough, an isomorphism on homotopy groups $`\pi _p`$ for $`p<\frac{1}{2}\text{dim }D_r^k`$ and a surjection on $`\pi _p`$ for $`p=\frac{1}{2}\text{dim }D_r^k`$. The same property also holds for homology groups.
Remark that the asumption of only one stratum implies that $`r=\mathrm{min}\{r_e,r_f\}1`$ and $`2(r_er+1)(r_fr+1)=4(|r_er_f|+2)>\text{dim }M`$. Along the proofs we will suppose that $`r_er_f`$, leaving the details of the other case to the reader. We proceed in several steps.
#### 5.3.1. Determinant vector spaces
Let $`V,W`$ be vector spaces of dimensions $`m`$ and $`n`$ ($`mn`$) respectively. We need some results about the behaviour of the determinant vector space $`^r(V^{})^rW`$ associated to the vector space of linear morphisms $`V^{}W`$. We define the $`r`$-fold wedge product $`^r`$ of a linear application $`\phi \mathrm{Hom}(V,W)`$ as the linear application
$`{\displaystyle ^r}\phi :{\displaystyle ^r}V`$ $``$ $`{\displaystyle ^r}W`$
$`v_1\mathrm{}v_r`$ $``$ $`\phi (v_1)\mathrm{}\phi (v_r).`$
Thus we obtain a non-linear map $`^r:\mathrm{Hom}(V,W)\mathrm{Hom}(^rV,^rW)`$. The previous definition extends in an obvious way to any pair of vector bundles $`E`$ and $`F`$ providing a non-linear map of vector bundles $`^r:\mathrm{Hom}(E,F)\mathrm{Hom}(^rE,^rF)`$. With this notation a rank $`r1`$ determinantal submanifold $`D_{r1}`$ associated to a morphism $`\phi `$ between vector bundles $`E`$ and $`F`$ is the set
$$D_{r1}=\{xM:^r\phi (x)=0\}.$$
###### Lemma 5.6.
Let $`V,W`$ be vector spaces of dimensions $`m`$ and $`n`$ ($`mn`$) respectively, then the set $`R(V,W)=^n(\mathrm{Hom}(V,W))\{0\}`$ is a smooth open complex submanifold of $`\mathrm{Hom}(^nV,^nW)`$ of dimension $`mn+1`$. Moreover $`R(V,W)`$ is invariant under multiplication by non-zero complex scalars, and so given any point $`dR(V,W)`$ then, using the standard identification between a vector space and its tangent space at a point, $`dT_dR(V,W)`$.
Proof. The last statement is obvious. For the first one, fix basis $`(e_1,\mathrm{},e_m)`$ in $`V`$ and $`(f_1,\mathrm{},f_n)`$ in $`W`$. First notice that $`R(V,W)`$ is invariant under the actions of the groups $`\text{GL}(V)`$ and $`\text{GL}(W)`$. Thus for computing $`T_{^n(\phi )}R(V,W)`$ we can restrict our atention to the point
(21)
$$\phi =\underset{i=1}{\overset{n}{}}f_ie_i^{}.$$
Notice that this is possible since the condition $`^n\phi 0`$ implies that the linear map $`\phi `$ has rank $`n`$ and therefore suitable changes of basis provide the expression (21). Now, we only have to compute the images of the tangent basis $`\phi _{ij}=\frac{d}{dt}|_{t=0}(\phi +tb_{ij})`$, where $`b_{ij}=f_je_i^{}`$. First assume that $`in`$, then we obtain
$$(^n)_{}\phi _{ij}=\{\begin{array}{cc}\phi ,\hfill & i=j,\hfill \\ 0,\hfill & ij\hfill \end{array}$$
However for the cases $`i>n`$ we obtain
$$(^n)_{}\phi _{ij}=(1)^{nj}f_1\mathrm{}f_ne_1^{}\mathrm{}e_{j1}^{}e_{j+1}^{}\mathrm{}e_n^{}e_i^{}.$$
Then the image of this tangent basis has dimension $`mn+1`$. This happens at any point of $`R(V,W)`$. Now, the image of an application of constant rank is locally a submanifold.
Finally we have to check that the counterimages of $`^n`$ are connected, i.e. given two morphisms $`\phi _0`$ and $`\phi _1`$ such that $`^n\phi _0=^n\phi _1`$ then there exists a path $`\{\phi _t\}_{t[0,1]}`$ connecting the two morphisms and satisfying $`^n\phi _t=^n\phi _0`$. For this, note that the kernels of $`\phi _0`$ and $`\phi _1`$ coincide. Therefore there exists an endomorphism $`A`$ in $`GL(W)`$ such that $`A\phi _0=\phi _1`$. Such $`A`$ is forced to be in $`SL(W)`$. Now fix a path $`A_t`$, $`t[0,1]`$, connecting the identity with $`A`$ and put $`\phi _t=A_t\phi _0`$. $`\mathrm{}`$
This Lemma extends trivially to vector bundles to obtain the following
###### Lemma 5.7.
Let $`E,F`$ be vector bundles of ranks $`m`$ and $`n`$ ($`mn`$) respectively, then the fibration $`R(E,F)`$, given at any point $`xM`$ by $`^n(\mathrm{Hom}(E_x,F_x))\{0\}`$, has smooth fibers which are open complex submanifolds of $`\mathrm{Hom}(^nE_x,^nF_x)`$ of dimension $`mn+1`$. Moreover $`R(E,F)`$ is invariant under multiplication by a never null complex-valued function, and so given any point $`dR(E,F)`$ we have, using the standard identification between a vector space and its tangent space at a point, that $`dT_dR(E,F)`$.
#### 5.3.2. Generalized asymptotically holomorphic sequences of sections of vector bundles
Now we recall the process of construction of a sequence of symplectic determinantal submanifolds. Let $`E`$, $`F`$ be vector bundles of ranks $`r_e`$ and $`r_f`$, respectively, and suppose $`r_er_f`$. Write
$$r=\mathrm{min}\{r_e,r_f\}=r_f.$$
Fix a generic section $`s`$ of the universal bundle $`𝒰_{ef}`$ over $`\text{Bi}(r_e,r_f,N)`$. We embed $`M`$ in $`\text{Bi}(r_e,r_f,N)`$ constructing an asymptotically holomorphic sequence $`\varphi _k`$ of embeddings. Using Lemma 5.2 and Proposition 5.3 we assure that the sequence is transverse to the holomorphic determinantal varieties defined by $`s`$ in $`\text{Bi}(r_e,r_f,N)`$. We can define a sequence of sections of the bundles $`E^{}FL^{2k}`$ as
$$s_k=\varphi _k^{}s.$$
We consider now the connection $`\widehat{}_k`$ defined on $`E^{}FL^{2k}`$ as the pull-back of the canonical one defined in $`𝒰_{ef}`$. Also we consider in $`M`$ the sequence of metrics $`\widehat{g}_k`$ defined as the pull-back through $`\varphi _k`$ of the standard metric on the bigrassmanian $`\text{Bi}(r_e,r_f,N)`$. Then using properties 1 and 2 of Definition 1.1, we obtain that the sequence $`s_k`$ is asymptotically holomorphic with respect to the fixed complex structure $`J`$ in $`M`$, computing the derivatives respect to $`\widehat{}_k`$ and the norms respect to $`\widehat{g}_k`$. Analogously taking the pull-back of the connection associated to $`^r\pi _e^{}(𝒰_e)^r\pi _f^{}(𝒰_f)`$, we obtain connections for the bundles $`^r(E^{}L^k)^r(FL^k)`$. Then the sequence $`^rs_k`$ is asymptotically $`J`$-holomorphic with respect to these connections and to the metric $`\widehat{g}_k`$.
Now we look for a condition to express when the sections $`^rs_k`$ are transversal in a certain sense. The key property is
###### Lemma 5.8.
Let $`E`$ and $`F`$ be vector bundles with connections $`^e`$ and $`^f`$ respectively. Suppose $`s`$ is a section of the bundle of morphisms $`E^{}F`$ equipped with the connection $`^{ef}`$ induced by $`^e`$ and $`^f`$. If $`^rs(x)0`$ at a point $`xM`$, then $`^{ef}^rs(x)T_{^rs(x)}R(E_x,F_x)`$.
Proof. To check this we only have to show that the following diagram is commutative
$$\begin{array}{ccc}\mathrm{\Omega }^0(E^{}F)& \stackrel{id^{ef}}{}& \mathrm{\Omega }^0(E^{}F)\mathrm{\Omega }^1(E^{}F)\\ ^r& & T^r\\ \mathrm{\Omega }^0(^r(E^{})^rF)& \stackrel{id^{ef}}{}& \mathrm{\Omega }^0(^r(E^{})^rF)\mathrm{\Omega }^1(^r(E^{})^rF).\end{array}$$
The map $`T^r`$ is defined as
$`T{\displaystyle ^r}:\mathrm{\Omega }^0(E^{}F)\mathrm{\Omega }^1(E^{}F)`$ $``$ $`\mathrm{\Omega }^0({\displaystyle ^r}(E^{}){\displaystyle ^r}F)\mathrm{\Omega }^1({\displaystyle ^r}(E^{}){\displaystyle ^r}F)`$
$`(s_0,s_1)`$ $``$ $`(s_0,\underset{t0}{lim}{\displaystyle \frac{^r(s_0+ts_1)}{t}}).`$
To check this one fixes local frames in $`E`$ and $`F`$ and carries out the computation explicitly. $`\mathrm{}`$
Given a generic section $`s`$ of the bundle of morphisms $`E^{}F`$ then we denote by $`D_{r2}^ϵ`$ the $`ϵ`$-neighborhood of the determinantal set $`D_{r2}`$ associated to $`s`$.
###### Definition 5.9.
Let $`E`$ and $`F`$ be vector bundles over $`M`$ of ranks $`r_e`$ and $`r_f`$ ($`r_er_f`$) respectively. Put $`r=r_f`$. We say that the section $`s`$ is $`\eta `$-$`^r`$-transverse to $`0`$, for some $`\eta >0`$, if for any $`xMD_{r2}^\eta `$ such that $`|^rs(x)|<\eta `$ then the covariant derivative $`\widehat{s}(x)=^rs(x)`$ has rank $`r_er_f+1`$ and also there exists a right inverse $`\theta :T_{^rs(x)}R^r(E,F)T_xM`$ of $`\widehat{s}(x)`$ with norm less that $`\eta ^1`$.
We cannot impose the estimated transversality near the stratum $`D_{r2}`$ because the section $`^rs`$ is always critical in that stratum, so if we want to obtain a notion of estimated transversality we need to remove a neighborhood of $`D_{r2}`$.
Observe that given any small $`\eta >0`$, the section $`s`$ is $`\eta `$-$`^r`$-transverse to $`0`$.
Using that $`\varphi _k(M)`$ is transverse to $`D_{r1}`$ we can check that $`s_k`$ is $`\eta ^{}`$-$`^r`$-transverse to $`0`$ on $`M`$, for some universal $`\eta ^{}>0`$, with the connections and metrics defined in the prededent lines. Observe that to guarantee this property is absolutely necessary that the minimum distance from $`\varphi _k(M)`$ to $`D_{r2}`$ be greater than $`\eta `$, but this is true by construction.
#### 5.3.3. Proof of Proposition 5.5.
We have as starting data a sequence of asymptotically holomorphic sections of the bundles $`E^{}FL^{2k}`$ obtained by pull-back of a fixed section $`s`$ of the universal bundle $`𝒰_{ef}`$. As before, we may suppose that $`r_er_f`$ and write $`r=r_f`$. Therefore the only non-empty stratum is $`D_{r1}^k`$, by assumption. We assume also that $`s_k`$ is $`\eta `$-$`^r`$-transverse to $`0`$, for a universal $`\eta >0`$. The stratum $`D_{r2}`$ is empty and so the $`\eta `$-$`^r`$-transversality is checked all over $`M`$. We can follow the ideas of \[Do96, Au97\] to develop the proof.
We define the function $`f_k=\mathrm{log}|^rs_k|^2`$. Clearly $`f_k(\mathrm{})=D_{r1}^k`$. Denote the complex dimension of $`D_{r1}^k`$ by $`N`$. We are going to show that all the critical points of $`f_k`$ are of index at least $`N+1`$. Therefore a standard Morse-theoretic argument will finish the proof.
Denote $`\sigma _k=^rs_k`$. First notice that if $`x`$ is a critical point of $`|\sigma _k|^2`$ then $`\sigma _k(x)`$ is not in the image of $`\sigma _k`$ and so $`\sigma _k`$ is not surjective to $`T_{\sigma _k(x)}R(E_x,F_x)`$. It follows from the $`\eta `$-$`^r`$-transversality property that $`|\sigma _k(x)|>\eta `$.
Now we differentiate $`f_k`$ to obtain
$$f_k=\frac{1}{|\sigma _k|^2}(\sigma _k,\sigma _k+\sigma _k,\overline{}\sigma _k).$$
At a critical point $`x`$, $`f_k(x)=0`$. Using the asymptotic holomorphic bounds we obtain
(22)
$$|\sigma _k,\sigma _k|=|\overline{}\sigma _k,\sigma _k|Ck^{1/2}|\sigma _k|.$$
Differentiating a second time we obtain, evaluating at a critical point, the expression
$$\overline{}\mathrm{log}|\sigma |^2=\frac{1}{|\sigma |^2}(\overline{}\sigma ,\sigma \sigma ,\sigma +\overline{}\sigma ,\overline{}\sigma +\sigma ,\overline{}\sigma ),$$
where we omit the subindex $`k`$ for simplicity. Recall that $`\overline{}+\overline{}`$ equals the $`(1,1)`$-part of the curvature of the bundle $`^r(E^{}L^k)^r(FL^k)`$. Its $`(1,1)`$-curvature $`R`$ is the pull-back through $`\varphi _k`$ of the $`(1,1)`$-curvature $`\stackrel{~}{R}`$ of $`^r𝒰_e^r𝒰_f`$. So we obtain
$$\overline{}f_k=\frac{1}{|\sigma |^2}(R\sigma ,\sigma \overline{}\sigma ,\sigma +\sigma ,\overline{}\sigma \sigma ,\sigma +\overline{}\sigma ,\overline{}\sigma ).$$
We define the subspace
$$V=\{vT_xM|_v\sigma (x)=\lambda \sigma (x),\text{ for some }\lambda \}.$$
Using the inequality (22) we obtain, for any $`vV`$, that
$$|_v\sigma ,\sigma |=|_v\sigma ||\sigma |Ck^{1/2}|\sigma |.$$
Restricting $`\overline{}f_k`$ to $`V`$, it equals to $`\frac{1}{|\sigma |^2}R\sigma ,\sigma +O(k^{1/2})`$. Denote the Hessian of $`f`$ by $`H_f`$. We know that $`H_f(u)+H_f(Ju)=2i\overline{}f_k(u,Ju)=2i\frac{1}{|\sigma |^2}R(u,Ju)\sigma ,\sigma +O(k^{1/2})`$, for any unit vector $`uV`$. We claim that it is possible to bound above the expression
(23)
$$2i\frac{1}{|\sigma |^2}R(u,Ju)\sigma ,\sigma $$
by a universal strictly negative constant, where $`u`$ is a unitary vector. For this we need to estimate the curvature $`R`$. We start by computing the curvature of the universal bundle $`𝒰`$ over the grassmannian $`\text{Gr}(r,N)`$. We use the local expression of the curvature of $`𝒰^{}`$ from \[We73, page 82\],
$$R_𝒰^{}=h^1\overline{df}^tdfh^1\overline{df}^tfh^1\overline{f}^tdf,$$
where $`f=(f_1,\mathrm{},f_r)`$ is a frame in an open neighborhood of $`\text{Gr}(r,N)`$ and $`h=\overline{f}^tf`$. We may assume that we are at the point $`\mathrm{\Pi }_0=[\text{I}|\mathrm{𝟎}]`$ of the grassmannian, after suitable change of coordinates. Select the following holomorphic local frame,
$$f=((1,0,\stackrel{(r1)}{\mathrm{}},0,z_{11},\mathrm{},z_{1,nr}),\mathrm{},(0,\mathrm{},0,1,z_{r1},\mathrm{},z_{r,nr})),$$
So at the point $`\mathrm{\Pi }_0`$ we obtain $`R_𝒰^{}=\overline{df}^tdf`$ and
$$R_𝒰=df^t\overline{df}.$$
In the trivialization $`(z_{jk})`$ we take the standard basis $`e_{jk}=\frac{}{z_{jk}}`$. We obtain $`R_𝒰(e_{jk},ie_{jk})=ib_{jj}`$, where the endomorphism $`b_{jj}`$ is defined as $`e_je_j^{}`$. So the endomorphism $`iR_𝒰(u,Ju)`$ is semi-definite negative for $`uT_{\mathrm{\Pi }_0}\text{Gr}(r,N)`$ non-zero. This implies also that $`iR_{^k𝒰}(u,Ju)`$ is semi-definite negative, for $`1kr`$. Moreover computing $`iR_{^r𝒰}(u,Ju)`$ in $`\text{Gr}(r,N)`$, or recalling that $`𝒰`$ is very ample, we obtain that this endomorphism is definite negative. Returning to $`\text{Bi}(r_e,r_f,N)`$ with $`r=r_fr_e`$, we have that the curvature of $`^r𝒰_e^r𝒰_f`$ is
$$\stackrel{~}{R}=R_{\pi _e^{}^r𝒰_e}\text{I}_1+\text{I}_\nu R_{\pi _f^{}^r𝒰_f},$$
where $`\nu =\left(\begin{array}{c}r_e\\ r\end{array}\right)`$. So $`\stackrel{~}{R}(u,Ju)`$ is definite negative, for $`uT\text{Bi}(r_e,r_f,N)`$ unitary vector. Using that the sequence of embeddings $`\varphi _k=(\varphi _k^e,\varphi _k^f)`$ satisfies properties $`1`$ and $`2`$ of Definition 1.1, we get that the expression (23) is bounded above by a universal strictly negative number.
Therefore, for any unitary $`uV`$, $`H_f(u)+H_f(Ju)`$ is negative for $`k`$ large enough. Recall that from the definition we obtain that $`\text{dim }V2N+2`$. Suppose that there exists a subspace $`PT_xM`$ of real dimension at least $`2nN`$ such that $`H_f`$ in non-negative. The dimension of $`PJP`$ is at least $`2n2N`$, and there the function $`H_f()+H_f(J)`$ is, obviously, non-negative. Therefore $`PJP`$ has to intersect trivially with $`V`$ but $`\text{dim }PJP+\text{dim }V2n+2`$, and this is clearly impossible. So such space $`P`$ does not exist and then the index of $`f_k`$ at $`x`$ is greater than $`N`$. This finishes the proof. $`\mathrm{}`$
### 5.4. Chern classes of the constructed submanifolds
For computing the Chern classes of determinantal submanifolds, we shall use the results of Harris and Tu in \[HT84\]. All their results are stated for holomorphic determinantal submanifolds in a holomorphic manifold, but they apply without the condition of integrability of the complex structure. We state the formulas that we shall use. Following Subsection 5.1 we denote $`r_e=\text{rank}E`$, $`r_f=\text{rank}F`$, $`2n=\text{dim }M`$ and $`D_r`$ is the $`r`$-determinantal loci of a bundle map $`\phi :EF`$ constructed in Theorem 1.6. First of all, set
(24)
$$\mathrm{\Delta }_{i_1,\mathrm{},i_{r_er}}=\left|\begin{array}{cccc}c_{r_fr+i_1}& c_{r_fr+i_1+1}& \mathrm{}& \\ c_{r_fr+i_21}& c_{r_fr+i_2}& \mathrm{}& \\ & & \mathrm{}& \\ & & \mathrm{}& c_{r_fr+i_{r_er}}\end{array}\right|,$$
where $`c_j=c_j(FE)`$. For instance, $`\mathrm{\Delta }_{0,\mathrm{},0}=\mathrm{\Delta }=\mathrm{PD}([D_r])`$, which is the classical Porteous formula for the homology class of a determinantal locus. We can suppose that the indices $`i_j`$ are decreasing, and so if we have any index $`i_j=0`$ we do not write it, e.g. $`\mathrm{\Delta }_{2,1,0}=\mathrm{\Delta }_{2,1}`$.
In \[HT84\] a complete description of the Chern numbers of the tangent bundle of a determinantal submanifold is performed, supposing that $`D_{r1}=\mathrm{}`$ and so $`D_r`$ is smooth. We concentrate ourselves in the cases $`\text{dim}_{}D_r=1`$ and $`\text{dim}_{}D_r=2`$, where Harris and Tu obtain the following formulas:
1. For $`\text{dim}_{}M=(r_er)(r_fr)+1`$, then $`\text{dim}_{}D_r=1`$. We have
$$n_1(D_r)=c_1(D_r),[D_r]=(c_1(M)+(r_er)c_1(EF))\mathrm{\Delta }+(r_er_f)\mathrm{\Delta }_1.$$
2. For $`\text{dim}_{}M=(r_er)(r_fr)+2`$, then $`\text{dim}_{}D_r=2`$. We have
$`n_{11}(D_r)`$ $`=`$ $`c_1^2(D_r),[D_r]=(c_1(M)+(r_er)c_1(EF))^2\mathrm{\Delta }+`$
$`+`$ $`2(r_er_f)(c_1(M)+(r_er)c_1(EF))\mathrm{\Delta }_1+(r_er_f)^2(\mathrm{\Delta }_2+\mathrm{\Delta }_{11}),`$
$`n_2(D_r)`$ $`=`$ $`c_2(D_r),[D_r]=(c_2(M)+(r_er)c_1(M)c_1(EF)+`$
$`+`$ $`(r_er)(c_2(E)c_2(F))+\left({\displaystyle \genfrac{}{}{0pt}{}{r_er}{2}}\right)c_1^2(E)(r_er)^2c_1(E)c_1(F)+`$
$`+`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{r_er+1}{2}}\right)c_1^2(F))\mathrm{\Delta }+`$
$`+`$ $`((r_er)c_1(M)+((r_er)(r_er_f)1)c_1(EF))\mathrm{\Delta }_1`$
$`+`$ $`{\displaystyle \frac{1}{2}}((r_er_f)^2+(r_er)+(r_fr)2)\mathrm{\Delta }_2+`$
$`+`$ $`{\displaystyle \frac{1}{2}}((r_er_f)^2(r_er)(r_fr)2)\mathrm{\Delta }_{11}.`$
In our case, we are going to apply the above formulas to morphisms $`\phi :E(L^{})^kFL^k`$. We have the following asymptotic expansions for Chern classes (we write $`\omega _k=\frac{k\omega }{2\pi }`$ for simplicity)
$`c_p(FL^k)=\left({\displaystyle \genfrac{}{}{0pt}{}{r_f}{p}}\right)\omega _k^p+O(k^{p1}),`$
$`c_p(E(L^{})^k)=\left({\displaystyle \genfrac{}{}{0pt}{}{r_e}{p}}\right)(\omega _k)^p+O(k^{p1}),`$
$`c_p=c_p(FL^kE(L^{})^k)=\text{Coeff}_{x^p}{\displaystyle \frac{(1+x)^{r_f}}{(1x)^{r_e}}}\omega _k^p+O(k^{p1})=`$
(25) $`={\displaystyle \underset{i=0}{\overset{r_f}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{r_f}{i}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{r_e+pi1}{pi}}\right)\omega _k^p+O(k^{p1}).`$
We are going to give two families of examples to show that the symplectic manifolds obtained here are more general than those in \[Au97\].
#### 5.4.1. Example 1.
Choose $`\text{dim}_{}M=(r_er)(r_fr)+1`$ and so we can apply the formulas for the complex $`1`$-dimensional case. Also suppose that $`r=1`$ and $`r_e=2`$, so $`\text{dim}_{}M=r_f=n>1`$. By Proposition 5.5 the submanifolds $`D_1`$ are connected. Now $`\mathrm{PD}[D_1]=\mathrm{\Delta }=c_{n1}`$ and $`\mathrm{\Delta }_1=c_n`$. Using (25) we get that
$`\text{vol}_{\omega _k}(D_1)`$ $`=`$ $`\mathrm{\Delta }\omega _k=(n2^{n1}+O(k^1))\text{vol}_{\omega _k}(M),`$
$`n_1(D_1)`$ $`=`$ $`(n+2)\omega _k\mathrm{\Delta }+(2n)\mathrm{\Delta }_1+O(k^1)\text{vol}_{\omega _k}(M)=`$
$`=`$ $`((n+2)n2^{n1}+(2n)(n2^{n1}+2^n)+O(k^1))\text{vol}_{\omega _k}(M)`$
$`{\displaystyle \frac{n_1(D_1)}{\text{vol}_{\omega _k}(D_1)}}`$ $`=`$ $`22n+{\displaystyle \frac{4}{n}}+O(k^1).`$
To compare with the Auroux’ case we compute the precedent symplectic invariants for this situation. Denote by $`Z`$ the zero set of a transverse section of a bundle of the form $`EL^k`$, we choose $`\text{rank}E=n1`$ to set up the comparison. Suppose that $`Z`$ is symplectic. Using Proposition 5 in \[Au97\] we obtain
$`\text{vol}_{\omega _k}(Z)`$ $`=`$ $`(1+O(k^1))\text{vol}_{\omega _k}(M),`$
$`n_1(Z)`$ $`=`$ $`(1n+O(k^1))\text{vol}_{\omega _k}(M),`$
$`{\displaystyle \frac{n_1(Z)}{\text{vol}_{\omega _k}(Z)}}`$ $`=`$ $`1n+O(k^1).`$
Therefore there does not exist any $`n2`$ such that the quotients $`\frac{n_1(D_1)}{\text{vol}_{\omega _k}(D_1)}`$ coincide with the quotients $`\frac{n_1(Z)}{\text{vol}_{\omega _k}(Z)}`$, obviously for $`k`$ large enough. So Auroux’ sequences of submanifold are not symplectomorphic to our sequences of determinantal submanifolds.
To check that, for $`k`$ large, our determinantal submanifolds do not coincide with Auroux’ examples we work as follows. Suppose that for integers $`k_1,k_2`$ the submanifold $`D_1=D_1^{k_1}`$ is isotopic to $`Z=Z_{k_2}`$. Then they define the same cohomology class and hence $`n2^{n1}k_1=k_2+O(1)`$. Also $`n_1(D_1)=n_1(Z)`$ implies $`(22n+\frac{4}{n})k_1=(1n)k_2+O(1)`$. So, for large enough $`k`$’s, $`(1n)n2^{n1}=22n+\frac{4}{n}`$ and hence $`n=2`$. Therefore for $`n>2`$ and large $`k`$ we get new examples of symplectic submanifolds.
Note that for $`n=r_e=r_f=2`$, the determinantal set $`D_1`$ for a morphism $`\phi :E(L^{})^kFL^k`$ is the zero set of the section $`^2\phi `$ of $`^2E^{}^2FL^{4k}`$. Since this zero set is smooth of the expected codimension, our example is just one of Auroux’ examples.
#### 5.4.2. Example 2.
Now, choose $`\text{dim}_{}M=(r_er)(r_fr)+2`$ and so we can apply the formulas for the complex $`2`$-dimensional case. Again we suppose that $`r=1`$ and $`r_e=2`$, so $`\text{dim}_{}M=r_f+1=n>2`$. By Proposition 5.5 these submanifolds are connected. In this case we have
$`\text{vol}_{\omega _k}(D_1)`$ $`=`$ $`((n1)2^{n2}+O(k^1))\text{vol}_{\omega _k}(M),`$
$`n_{11}(D_1)`$ $`=`$ $`(4(n1)(n^25)2^{n2}+O(k^1))\text{vol}_{\omega _k}(M)`$
$`n_2(D_1)`$ $`=`$ $`(2(n^2+n4)(n1)2^{n2}+O(k^1)\text{vol}_{\omega _k}(M)`$
$`{\displaystyle \frac{n_2(D_1)}{n_{11}(D_1)}}`$ $`=`$ $`{\displaystyle \frac{n^2+n4}{2(n^25)}}+O(k^1).`$
For the Auroux’ case with $`\text{rank}E=n2`$ we obtain
$`\text{vol}_{\omega _k}(Z)`$ $`=`$ $`(1+O(k^1))\text{vol}_{\omega _k}(M),`$
$`n_{11}(Z)`$ $`=`$ $`((n2)^2+O(k^1))\text{vol}_{\omega _k}(M),`$
$`n_2(Z)`$ $`=`$ $`\left({\displaystyle \frac{(n1)(n2)}{2}}+O(k^1)\right)\text{vol}_{\omega _k}(M)`$
$`{\displaystyle \frac{n_2(Z)}{n_{11}(Z)}}`$ $`=`$ $`{\displaystyle \frac{n1}{2(n2)}}+O(k^1).`$
If we compute the symplectic invariants $`\frac{n_{11}(Z)}{\text{vol}_{\omega _k}(Z)}`$ and $`\frac{n_2(Z)}{\text{vol}_{\omega _k}(Z)}`$, it is easy to verify that Auroux’ submanifolds are not symplectomorphic to the determinantal ones constructed in this example.
Moreover, for $`4`$-manifolds, the numbers $`n_2=\chi `$ and $`n_{11}=(2\chi +3\sigma )/4`$ are topological invariants. Therefore $`\frac{n_2}{n_{11}}`$ is a topological invariant. Comparing the Auroux’ case and the determinantal example we find that these symplectic submanifolds are not even homeomorphic, for $`k`$ large enough (even choosing different $`k`$’s in either case).
In general, it is clear that the determinantal class is quite bigger than the Donaldson-Auroux one. We could compute more examples and more precise invariants using recent results from algebraic geometry about the topology of determinantal submanifolds. As a reference it could be useful \[HT84b, Pr88, PP91\]. Remark that in these references the computations are performed even in the singular case. To adapt them to the symplectic category we would need to define the Segre classes of a singular symplectic manifold. This definition seems quite natural.
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# 1 Introduction
## 1 Introduction
Two-Time Physics - is an approach that provides a new perspective for understanding ordinary one-time dynamics from a higher dimensional, more unified point of view including two timelike dimensions. This is achieved by introducing new gauge symmetries that insure unitarity, causality and absence of ghosts. The new phenomenon in two-time physics is that the gauge symmetry can be used to obtain various one-time dynamical systems from the same simple action of two-time physics, through gauge fixing, thus uncovering a new layer of unification through higher dimensions.
The principle behind two-time physics is the gauge symmetry . The basic observation in its simplest form is that for any theory the Lagrangian has the form $`L=\frac{1}{2}\left(\dot{x}p\dot{p}x\right)H(x,p)`$ up to an inessential total time derivative. The first term has a global Sp$`(2,R)`$ symmetry that transforms $`(x,p)`$ as a doublet. The basic question we pose is: what modification of the Lagrangian can turn this global symmetry into a local symmetry? The reason to be interested in such a local symmetry is that duality symmetries in M-theory and N=2 super Yang-Mills theory have similarities to gauge symplectic transformations, and their origin in the fundamental theories in physics remains a mystery. Understanding them may well be the key to constructing M-theory. Independent of M-theory, the question is a fundamental one in its own right, and its investigation has already led to a reformulation of ordinary one-time dynamical systems in a new language of two-time physics. This has uncovered previously unnoticed higher symmetries in well known one-time dynamical systems, and provided a new level of unification through higher dimensions for systems that previously would have been considered unrelated to each other . The simplest Sp$`(2,R)`$ gauge symmetry has generalizations when spin , supersymmetry , and extended objects (branes) are part of the theory. Recent works have given an indication that the domain of unification of two-time physics can be enlarged in additional directions in field theory including interactions, and in the world of branes .
In the two-time physics approach the familiar one-time is a gauge dependent concept. From the point of view of a two-time observer the true gauge invariants are identical in a variety of one-time dynamical systems that are unified by the same two-time action. Such gauge invariant quantities can be used to test the validity of the underlying unification. An important gauge invariant concept is the global symmetry of the two-time action, which must be shared by all the gauge fixed one-time dynamical systems. In the simplest case the global symmetry is SO$`(d,2)`$, but this can be different in the presence of background fields as we will see in the current paper. In the simple case, the SO$`(d,2)`$ symmetry has been shown to be present in the same irreducible representation in all the one-time dynamical systems derived from the same two-time action. The presence of such symmetries, which remained unknown even in elementary one-time systems until the advent of two-time physics, can be considered as a test of the underlying unification within a two-time theory .
Two-Time Physics has been generalized to include global space-time supersymmetry and local kappa supersymmetry with two-times . This led to a framework which suggests that M-theory could be embedded in a two-time theory in 13 dimensions, with a global OSp(1$`|`$64) symmetry. In this scenario the different corners of M-theory correspond to gauge fixed sectors of the 13D theory, and the dualities in M-theory are regarded as gauge transformations from one fixed gauge to another fixed gauge. Then the well known supersymmetries of various corners of M-theory appear as subsupergroups of OSp$`\left(1|64\right)`$. This mechanism has been illustrated through explicit examples of dynamical particle models which may be regarded as a toy-M-theory. In the 11D-covariant gauge fixed corner, the supergroup OSp(1$`|`$64) is interpreted as the conformal supergroup in 11-dimensions, with 32 supersymmetries and 32 superconformal symmetries. But in other gauge fixed sectors, the same OSp(1$`|`$64) symmetry of two-time physics is realized and interpreted differently, thus revealing various corners of toy-M-theory on which a subsupergroup is linearly realized while the rest is non-linearly realized. Indeed OSp(1$`|`$64) contains various embeddings that reveal 13,12,11 dimensional supersymmetries, as well as the usual 10-dimensional type-IIA, type-IIB, heterotic, type-I, and AdS$`{}_{D}{}^{}`$S<sup>k</sup> type supersymmetries in $`D+k=11,10`$ and lower dimensions. The explicit models provided by illustrate these ideas while beginning to realize dynamically some of the observations that suggested two-time physics in the framework of branes, dualities and extended supersymmetries in M-theory, F-theory, and S-theory -.
In this paper we generalize the worldline formulation of two-time physics by including background gravitational and gauge fields and other potentials. To keep the discussion simple we concentrate mainly on particles without supersymmetry. For spinless particles, as in the case of the free theory, local Sp($`2,R`$) gauge symmetry is imposed as the underlying principle. For the gauge symmetry to be valid, the gravitational and gauge fields and other potentials must obey certain differential equations. We show that the gauge field obeys an equation that generalizes a similar one discovered by Dirac in 1936 in the flat background, while the gravitational field satisfies a closed homothety condition. When all fields are simultaneously present they obey coupled equations. Examples of background fields that solve these equations are provided.
A similar treatment for spinning particles in background fields is given. As in the free theory, local OSp$`\left(n|2\right)`$ gauge symmetry is imposed as the underlying principle. The set of background fields is now richer. The generalizations of Dirac’s equation and the closed homothety conditions in the presence of spin are derived. Instead of OSp$`\left(n|2\right)`$ gauge symmetry it may also be possible to consider other supergroups that contain Sp$`(2,R)SL(2,R).`$
In the presence of the background fields one learns that much larger classes of one-time dynamical systems can now be reformulated as gauge fixed versions of the same two-time theory. This extends the domain of unification of one-time systems through higher dimensions and a sort of duality symmetry (the Sp$`(2,R)`$ gauge symmetry and its generalizations in systems with spin and/or spacetime supersymmetry, and branes). Furthermore, with the results of this paper it becomes evident that all one-time particle dynamics can be reformulated as particle dynamics in two-time physics. This provides a much broader realm of possible applications of the two-time physics formalism.
One possible practical application of the formulation is to provide a tool for solving problems by transforming a complicated one-time dynamical system (one fixed gauge) to a simpler one-time dynamical system (another fixed gauge), as in duality transformations in M-theory. Although this may turn out to be the computationally useful aspect of this formulation, it is not explored in the present paper since our main aim here is the formulation of the concepts.
The two-time formulation also has deeper ramifications. By providing the perspective of two-time physics for ordinary physical phenomena, the familiar “time” dimension appears to play a less fundamental role in the formulation of physics. Since the usual “time” is a gauge dependent concept in the new formulation, naturally one is led to a re-examination of the concept of “time” in this new setting.
## 2 Local and global symmetry
We start with a brief summary of the worldline formulation of two-time physics for the simplest case of spinless particle dynamics without background fields and without a Hamiltonian - (i.e. the “free” case). Just demanding local symmetry for the first term in the Lagrangian gives a surprisingly rich model based on $`Sp(2,R)`$ gauge symmetry described by the action
$`S_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle 𝑑\tau D_\tau X_i^MX_j^N\epsilon ^{ij}\eta _{MN}}`$ (1)
$`=`$ $`{\displaystyle 𝑑\tau (_\tau X_1^MX_2^N\frac{1}{2}A^{ij}X_i^MX_j^N)\eta _{MN}}.`$ (2)
Here $`X_i^M(\tau )`$ is an $`Sp(2,R)`$ doublet, consisting of a coordinate and its conjugate momentum ($`X_1^MX^M`$ and $`X_2^MP^M`$). The indices $`i,j=1,2`$ denote the doublet of $`Sp(2,R)`$; they are raised and lowered by the antisymmetric Levi-Civita symbol $`\epsilon _{ij}`$. The covariant derivative $`D_\tau X_i^M`$ that appears in (1) is defined as
$$D_\tau X_i^M=_\tau X_i^M\epsilon _{ik}A^{kl}X_l^M,$$
(3)
where $`A^{ij}\left(\tau \right)`$ are the three Sp$`(2,R)`$ gauge potentials in the adjoint representation written as a 2$`\times 2`$ symmetric matrix. The local $`Sp(2,R)`$ acts as $`\delta X_i^M=\epsilon _{ik}\omega ^{kl}X_l^m`$ and $`\delta A^{ij}=\omega ^{ik}\epsilon _{kl}A^{kj}+\omega ^{jk}\epsilon _{kl}A^{ik}+_\tau \omega ^{ij}`$, where $`\omega ^{ij}\left(\tau \right)`$ are the Sp$`(2,R)`$ gauge parameters. The second form of the action (2) is obtained after an integration by parts so that only $`X_1^M`$ appears with derivatives. This allows the identification of $`X,P`$ by the canonical procedure ($`X_1^MX^M`$ and $`X_2^MP^M=S_0/\dot{X}_{1M}`$). A third form of the action can be obtained by integrating out $`X_2^M`$ and writing it in terms of $`X^M`$ and $`\dot{X}^M`$ . Then the local Sp$`(2,R)=`$SO$`(1,2)`$ can also be regarded as the local conformal group on the worldline (including $`\tau `$ reparametrization, local scale transformations, and local special conformal transformations) and the theory can be interpreted as conformal gravity on the worldline .
The gauge fields $`A^{11}`$, $`A^{12}`$, and $`A^{22}`$ act as Lagrange multipliers for the following three first class constraints
$$Q_{ij}^0=X_iX_j=0X^2=P^2=XP=0,$$
(4)
as implied by the local $`Sp(2,R)`$ invariance. From the basic quantum rules for $`(X^M,P^M)`$ one can verify that the $`Q_{ij}^0`$ form the Sp$`(2,R)`$ algebra
$`[Q_{ij},Q_{kl}]`$ $`=`$ $`i\epsilon _{jk}Q_{il}+i\epsilon _{ik}Q_{jl}+i\epsilon _{jl}Q_{ik}+i\epsilon _{il}Q_{jk},or`$ (5)
$`[Q_{11},Q_{22}]`$ $`=`$ $`4iQ_{12},[Q_{11},Q_{12}]=2iQ_{11},[Q_{22},Q_{12}]=2iQ_{22}.`$ (6)
The two timelike dimensions are not put in by hand, they are implied by the local $`Sp(2,R)`$ symmetry. It is precisely the solution of the constraints $`Q_{ij}^0=0`$ that require the global metric $`\eta _{MN}`$ in (1) to have a signature with two-time like dimensions: if $`\eta _{MN}`$ were purely Euclidean the only solution would be vanishing vectors $`X_i^M`$, if it had Minkowski signature (one time) the only solution would be two lightlike parallel vectors $`X_i^M`$ without any angular momentum, if it had more than two timelike dimensions there would be ghosts that would render the theory non-unitary. The local Sp$`(2,R)`$ is just enough gauge symmetry to remove the ghosts due to two timelike dimensions. Thus, $`\eta _{MN}`$ stands for the flat metric on a ($`d,2`$) dimensional space-time. It is the only signature consistent with absence of ghosts, unitrarity or causality problems.
We now turn to the global symmetries that are gauge invariant under Sp$`(2,R)`$. The metric $`\eta _{MN}`$ is invariant under SO$`(d,2).`$ Hence the action (1-2) has an explicit global $`SO(d,2)`$ invariance. Like the two times, the $`SO(d,2)`$ symmetry of the action (1) is also implied by the local $`Sp(2,R)`$ symmetry when background fields are absent. The SO$`(d,2)`$ Lorentz generators
$$L^{MN}=X^MP^NX^NP^M=\epsilon ^{ij}X_i^MX_j^N$$
(7)
commute with the Sp$`(2,R)`$ generators, therefore they are gauge invariant. As we mentioned above, different gauge choices lead to different one-time particle dynamics (examples: free massless and massive particles, H-atom, harmonic oscillator, particle in AdS$`{}_{D}{}^{}\times `$S<sup>k</sup> etc.) all of which have $`SO(d,2)`$ invariant actions that are directly obtained from (1-2) by gauge fixing. Since the action (1-2) and the generators $`L^{MN}`$ are gauge invariant, the global symmetry SO$`(d,2)`$ is not lost by gauge fixing. This explains why one should expect a hidden (previously unnoticed, non-linearly realized) global symmetry SO$`(d,2)`$ for each of the one-time systems that result by gauge fixing<sup>2</sup><sup>2</sup>2A well known case is the SO$`(4,2)`$ conformal symmetry of the massless particle. Less well known is the SO$`(4,2)`$ symmety of the H-atom action, which acts as the dynamical symmetry for the quantum H-atom. Previously unknown is the SO$`(4,2)`$ symmetry of the massive non-relativistic particle action $`S=𝑑\tau `$ $`\dot{𝐱}^2/2m`$. Others are the SO$`(10,2)`$ symmetry of a particle in the AdS$`{}_{5}{}^{}\times S^5`$ background, or the SO$`(11,2)`$ symmetry in the AdS$`{}_{7}{}^{}\times S^4`$ and the AdS$`{}_{4}{}^{}\times S^7`$ backgrounds, etc. These and more examples of such non-linearly realized SO$`(d,2)`$ hidden symmetries for familiar systems in any space-time dimension $`d`$ are explicitly given in .. Furthermore all of the resulting one-time dynamical systems are quantum mechanically realized in the same unitary representation of SO$`(d,2)`$ -. This fact can be understood again as a simple consequence of representing the same quantum mechanical two-time system in various fixed gauges. The gauge choices merely distinguish one basis versus another basis within the same unitary representation of SO$`(d,2)`$ without changing the Casimir eigenvalues of the irreducible representation. Such relations among diverse one-time systems provide evidence that there is an underlying unifying principle behind them. The principle is the local Sp$`(2,R)`$ symmetry, and its unavoidable consequence of demanding a spacetime with two timelike dimensions which provides a basis for the global symmetry.
To describe spinning systems, worldline fermions $`\psi _\alpha ^M\left(\tau \right)`$, with $`\alpha =1,2,\mathrm{},n`$ are introduced. Together with $`X^M,P^M`$, they form the fundamental representation $`(\psi _\alpha ^M,X^M,P^M)`$ of the supergroup OSp$`\left(n/2\right)`$. Gauging this supergroup instead of Sp$`(2,R)`$ produces a Lagrangian that has $`n`$ local supersymmetries plus $`n`$ local conformal supersymmetries on the worldline, in addition to local Sp$`(2,R)`$ and local SO$`\left(n\right)`$. The full set of first class constraints that correspond to the generators of these gauge (super)symmetries are, at the classical level,
$$XX=PP=XP=X\psi _\alpha =P\psi _\alpha =\psi _{[\alpha }\psi _{\beta ]}=0.$$
(8)
The classical solution of these constraints, with a flat spacetime metric $`\eta ^{MN},`$ require a signature with two timelike dimensions. Therefore, as in the spinless case the global symmetry of the theory is SO$`(d,2)`$. It is applied to the label $`M`$ on $`(\psi _\alpha ^M,X^M,P^M)`$. The global SO$`(d,2)`$ generators $`J^{MN}`$ that commute with all the OSp$`\left(n/2\right)`$ gauge generators (8) now include the spin
$$J^{MN}=L^{MN}+S^{MN},S^{MN}=\frac{1}{2i}\left(\psi _\alpha ^M\psi _\alpha ^N\psi _\alpha ^N\psi _\alpha ^M\right).$$
(9)
As in the spinless case, by gauge fixing the bosons as well as the fermions, one finds a multitude of spinning one-time dynamical systems that are unified by the same two-time system both at the classical and quantum levels. All of these have SO$`(d,2)`$ hidden symmetry realized in the same representation, where the representation is different for each $`n`$ (number of local supersymmetries on the worldline, which is related also to the spin of the particle).
## 3 Interactions with background fields
The simple action in (2) is written in a flat two-time spacetime with metric $`\eta _{MN}`$ which could be characterized as a “free” theory. Interactions in the one-time systems emerged because of the first class constraints $`X^2=P^2=XP=0`$, not because of explicit interactions in the two time theory. The constraints generate the Sp$`(2,R)`$ gauge symmetry. This symmetry was realized linearly on the doublet $`X_i^M=(X^M,P^M)`$ and its generators were $`Q_{ij}^0=X_iX_j.`$
We now generalize the “free” theory to an “interacting” theory by including background gravitational and gauge fields and other potentials. This will be done by generalizing the worldline Hamiltonian (canonical conjugate to $`\tau `$) $`Q_{22}^0=P_MP_N\eta ^{MN}`$ to a more general form that includes a metric $`G^{MN}\left(X\right),`$ a gauge potential<sup>3</sup><sup>3</sup>3It is possible to generalize this discussion by promoting $`A`$ to a non-Abelian Yang-Mills potential coupled to a non-Abelian charge, which is an additional dynamical degree of freedom. To keep the discussion simple we take an Abelian $`A`$ in the present paper. to gauge-covariantize the momentum $`P_M+A_M\left(X\right),`$ and an additional potential $`U\left(X\right)`$ that is added to the kinetic term. Generalizing $`Q_{22}`$ in this way requires also generalizing all $`Q_{ij}^0`$ to $`Q_{ij}(X,P)`$ whose functional form will be determined. The Lagrangian is formally similar to the “free” case (2)
$$S=𝑑\tau (_\tau X^MP_M\frac{1}{2}A^{ij}Q_{ij}(X,P)).$$
(10)
Whatever the expressions for $`Q_{ij}(X,P)`$ are, by the equations of motion of the gauge potentials $`A^{ij},`$ they are required to form first class constraints that close under the Sp$`(2,R)`$ commutation rules (5), which should follow from the basic commutation rules of $`(X^M,P^M).`$ Furthermore, the local Sp$`(2,R)`$ transformation properties of the dynamical variables should be given by these generators under commutation rules
$`\delta X^M`$ $`=`$ $`{\displaystyle \frac{i}{2}}\omega ^{ij}\left(\tau \right)[Q_{ij}(X,P),X^M]={\displaystyle \frac{1}{2}}\omega ^{ij}\left(\tau \right){\displaystyle \frac{Q_{ij}(X,P)}{P_M}}`$ (11)
$`\delta P^M`$ $`=`$ $`i\omega ^{ij}\left(\tau \right)[Q_{ij}(X,P),P^M]={\displaystyle \frac{1}{2}}\omega ^{ij}\left(\tau \right){\displaystyle \frac{Q_{ij}(X,P)}{X^M}}`$ (12)
$`\delta A^{ij}`$ $`=`$ $`_\tau \omega ^{ij}+\omega ^{ik}\epsilon _{kl}A^{lj}+\omega ^{jk}\epsilon _{kl}A^{li}.`$ (13)
These certainly hold for the free case with $`Q_{ij}^0=X_iX_j,`$ but now we discuss the general case. Substituting these transformation laws into the Lagrangian we have (ignoring orders of operators at the classical level)
$$\delta L=_\tau \left(\delta X^M\right)P_M+_\tau X^M\delta P_M\frac{1}{2}\delta A^{ij}Q_{ij}(X,P)\frac{1}{2}A^{ij}\delta Q_{ij}(X,P)$$
(14)
where $`\delta Q_{ij}(X,P)=\frac{Q_{ij}}{X^M}\delta X^M+\frac{Q_{ij}}{P_M}\delta P_M`$. After an integration by parts of the first term, using (11-13) this becomes
$$\delta L=\frac{1}{2}_\tau \left(\omega ^{ij}Q_{ij}\right)\frac{1}{2}\left(\omega ^{ik}\epsilon _{kl}A^{lj}+\omega ^{jk}\epsilon _{kl}A^{li}\right)Q_{ij}\frac{1}{4}A^{ij}\omega ^{kl}\{Q_{ij},Q_{kl}\},$$
(15)
where $`\{Q_{ij},Q_{kl}\}`$ is the Poisson bracket
$$\{Q_{ij},Q_{kl}\}=\frac{Q_{ij}}{X^M}\frac{Q_{kl}}{P_M}\frac{Q_{ij}}{P_M}\frac{Q_{kl}}{X^M}.$$
(16)
Thus, if the $`Q_{ij}`$ satisfy the Sp$`(2,R)`$ algebra (5), then the Poisson bracket term cancels the second term, and $`\delta L`$ is a total derivative. Hence to insure the gauge invariance of the action $`S`$ we must require the differential constraints
$$\frac{Q_{ij}}{X^M}\frac{Q_{kl}}{P_M}\frac{Q_{ij}}{P_M}\frac{Q_{kl}}{X^M}=\epsilon _{jk}Q_{il}+\epsilon _{ik}Q_{jl}+\epsilon _{jl}Q_{ik}+\epsilon _{il}Q_{jk}.$$
(17)
With these restrictions we look for $`Q_{ij}(X,P)`$ that can be interpreted as dynamics with background fields, as opposed to dynamics in flat spacetime. To be able to integrate out the momenta $`P^M`$ we restrict these expressions to contain at the most two powers of $`P^M`$ (this restriction could be lifted to construct even more general systems <sup>4</sup><sup>4</sup>4The coefficients of higher powers of $`P^M`$ have the interpretation of higher spin fields). Also, keeping the analogy to the flat case, we will take $`Q_{11}`$ to have no powers of $`P^M,`$ $`Q_{12}`$ to have at most one power of $`P^M,`$ and $`Q_{22}`$ to have at the most two powers of $`P^M,`$ as follows
$`Q_{11}`$ $`=`$ $`W\left(X\right),Q_{12}={\displaystyle \frac{1}{2}}V^M\left(P_M+A_M\right)+{\displaystyle \frac{1}{2\sqrt{G}}}\left(P_M+A_M\right)\sqrt{G}V^M,`$ (18)
$`Q_{22}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{G}}}\left(P_M+A_M\right)\sqrt{G}G^{MN}\left(P_N+A_N\right)+U\left(X\right).`$ (19)
The functions $`W\left(X\right),V^M\left(X\right),G^{MN}\left(X\right),A_N\left(X\right),U\left(X\right)`$ will satisfy certain constraints. The expression for $`Q_{22}`$ is a generalization of the free worldline “Hamiltonian” in flat space $`\eta ^{MN}P_MP_N.`$ The factors of $`\sqrt{G}`$ are inserted to insure hermiticity of the operators in a quantum theory as applied on wavefunctions with a norm $`\sqrt{G}\psi ^{}\psi `$. In the classical theory the factors of $`\sqrt{G}`$ in $`Q_{12},Q_{22}`$ cancel since orders of operators are neglected, but in any case a reordering amounts to a redefinition of $`A_M\left(X\right)`$ and $`U\left(X\right)`$.
The combination $`P_M+A_M\left(X\right)`$ is gauge invariant under $`\delta _\mathrm{\Lambda }A_M\left(X\right)=_M\mathrm{\Lambda }\left(X\right)`$ and $`\delta _\mathrm{\Lambda }P_M=_M\mathrm{\Lambda }\left(X\right)`$, where $`\mathrm{\Lambda }\left(X\left(\tau \right)\right)`$ is a gauge function of spacetime. The Lagrangian has this gauge symmetry since it transforms into a total derivative under the gauge transformation $`\delta _\mathrm{\Lambda }L=_\tau X^M_M\mathrm{\Lambda }\left(X\right)=_\tau \mathrm{\Lambda }`$. Furthermore, the Lagrangian is a scalar under spacetime general coordinate transformations, since the $`Q_{ij}`$ are scalars when all the background fields are transformed as tensors, while the term $`_\tau X^MP_M`$ is invariant under $`\delta _\epsilon X^M=\epsilon ^M\left(X\right)`$ and $`\delta _\epsilon P_M=_M\epsilon ^NP_N`$. Of course, if the background fields are fixed, the general covariance and gauge symmetries are not generally valid, and only a subgroup that corresponds to Killing symmetries of the combined gauge and reparametrization transformations survive.
By integrating out $`P_M`$ we can rewrite the Lagrangian purely in terms of $`X^M\left(\tau \right)`$ and its derivatives $`\dot{X}^M\left(\tau \right)`$
$$L=\frac{1}{2A^{22}}\left(\dot{X}^MA^{12}V^M\right)G_{MN}\left(\dot{X}^NA^{12}V^N\right)\frac{A^{22}}{2}U\frac{A^{11}}{2}W\dot{X}^MA_M.$$
(20)
By inspection of (19) or (20) we interpret $`A_M\left(X\right)`$ as a gauge field, $`G_{MN}\left(X\right)`$ as a spacetime metric and $`U\left(X\right)`$ as an additional potential. The function $`W\left(X\right)0`$ is the constraint that replaces $`XX0`$ and the vector $`V^M\left(X\right)`$ can be thought of as a general coordinate transformation since the action of $`Q_{12}`$ on phase space is $`\delta _{12}X^M=V^M\left(X\right)`$ and $`\delta _{12}P_M=_MV^KP_K+_M\left(VA\right)`$ which looks like a general coordinate transformation up to a gauge transformation.
The classical local Sp$`(2,R)`$ transformation laws for $`(X^M,P_M)`$ in phase space follow from $`(\text{11},\text{13})`$
$`\delta X^M`$ $`=`$ $`\omega ^{12}\left(\tau \right)V^M+\omega ^{22}\left(\tau \right)G^{MN}\left(P_N+A_N\right)`$ (21)
$`\delta P_M`$ $`=`$ $`{\displaystyle \frac{1}{2}}\omega ^{11}\left(\tau \right)_MW\omega ^{12}\left(\tau \right)\left[\left(_MV^N\right)P_N+_M\left(VA\right)\right]`$
$`{\displaystyle \frac{1}{2}}\omega ^{22}\left(\tau \right)\left[\left(_MG^{KL}\right)\left(P_K+A_K\right)\left(P_L+A_L\right)+_MU+2G^{KL}_MA_K\left(P_L+A_L\right)\right]`$
This, together with (13), is a local symmetry of the action provided (17) is satisfied. These conditions give the following differential constraints on the functions $`W\left(X\right),`$ $`V^M\left(X\right),`$ $`G^{MN}\left(X\right),`$ $`A_N\left(X\right),`$ $`U\left(X\right)`$. From $`\{Q_{11},Q_{22}\}=4Q_{12}`$ we learn
$$V^M=\frac{1}{2}G^{MN}_NW.$$
(23)
From $`\{Q_{11},Q_{12}\}=2Q_{11}`$ we learn
$$V^M_MW=2W,orG^{MN}\left(_MW\right)\left(_NW\right)=4W.$$
(24)
Finally from $`\{Q_{22},Q_{12}\}=2Q_{22}`$ we learn (from the coefficients of each power of $`P_M)`$ that
$$\mathrm{\pounds }_VG^{MN}=2G^{MN},V^M_MU=2U,V^MF_{MN}=0,$$
(25)
where $`\mathrm{\pounds }_VG^{MN}`$ is the Lie derivative of $`G^{MN}`$ (an infinitesimal general coordinate transformation)
$$\mathrm{\pounds }_VG^{MN}V^K_KG^{MN}_KV^MG^{KN}_KV^NG^{MK},$$
(26)
and $`F_{MN}=_MA_N_NA_M`$ is the gauge field strength. The differential equation $`\mathrm{\pounds }_VG^{MN}=2G^{MN}`$ together with (23) was called a “closed homothety” condition on the geometry <sup>5</sup><sup>5</sup>5I learned this term when I came across ref., after having derived these equations independently sometime ago. The physical problem in the present paper is quite different than where our spacetime index $`M`$ (with (d,2) signature) is replaced by a particle label for multiparticles in (with Euclidean signature); nevertheless the mathematics formally coincide with ref.. After the current paper was submitted for publication, I was informed that similar equations were obtained in in the context of conformally invariant sigma models on a p+1 dimensional worldvolume, using a very different approach than ours. Although the case of p=0 (worldline) relevant for our case was missed by these authors, when their expressions are continued to p=0 they agree with our results. While there are formal similarities, an important difference between our work and those of and is that we have local SO(1,2)=Sp(2) symmetry as opposed to their global symmetry. This requires the constraints $`Q_{ij}(X,P)=0`$ which demand a spacetime with two timelike dimensions, thus leading to conceptually very different physics.. We have the added generalization of the gauge field $`A_M`$ in our case. When all fields are present they are coupled to each other.
The differential equation for the gauge field may also be rewritten in terms of the Lie derivative on the vector $`\mathrm{\pounds }_VA_M=_M\left(VA\right)`$, where the Lie derivative on the vector is $`\mathrm{\pounds }_VA_M=V^K_KA_M+_MV^KA_K`$ (an infinitesimal general coordinate transformation). Using the gauge invariance of the physics, without loss of generality one may choose an axial gauge $`VA=0`$. There still is a remaining gauge symmetry $`\delta _\mathrm{\Lambda }A_M=_M\mathrm{\Lambda }`$, for all $`\mathrm{\Lambda }`$ that satisfy $`V^K_K\mathrm{\Lambda }=0`$. Thus, the gauge field equation may be rewritten in the form
$$\mathrm{\pounds }_VA_M=0,VA=0,$$
(27)
with a remaining gauge symmetry of these equations $`\left\{\mathrm{\Lambda };V^K_K\mathrm{\Lambda }=0\right\}`$ which we will make use of later.
Any solution to the coupled equations (23, 24, 25, 27) gives an action with local Sp$`(2,R)`$ symmetry. Such an action provides a two-time physics theory including interactions with background fields. The global symmetries correspond to Killing symmetries in the presence of backgrounds, which is a subgroup embedded in general coordinate transformations combined with gauge transformations. This is the global symmetry, which in the flat and free case becomes SO$`(d,2)`$.
The Sp$`(2,R)`$ gauge symmetry may be gauge fixed to define a “time” and analyze the system from the point of view of one-time physics. The global symmetry described in the previous paragraph survives after gauge fixing the Sp$`(2,R)`$ local symmetry, since it commutes with it (recall the $`Q_{ij}`$ are invariant under general coordinate and gauge transformations). This global symmetry would then become the non-linearly realized hidden global symmetries in each of the one-time dynamical systems that emerge after gauge fixing (in the “free” case it is SO$`(d,2)`$). The symmetry must be realized in the same representation for each one-time dynamical system that belongs to the same class, where the class is fixed by a given set of background fields.
## 4 Pure gauge field background
When the background metric is flat $`G^{MN}=\eta ^{MN}`$ the only solution of the homothethy condition $`\mathrm{\pounds }_VG^{MN}=2G^{MN}`$ is $`V^M=X^M.`$ This immediately gives $`W=XX,`$ and $`U`$ is any homogeneous function of $`X^M`$ of degree -2. The global symmetry of the metric is SO$`(d,2).`$ If we want to keep the SO$`(d,2)`$ symmetry, $`U`$ could only be $`U=g/XX`$ (however, without the SO$`(d,2)`$ symmetry one can allow some other $`U`$ of degree -2).
The equations for the gauge field (27) simplify in flat space. The remaining gauge symmetry parameter is homogeneous of degree zero $`X\mathrm{\Lambda }=0`$ in $`d+2`$ dimensions. This is sufficient to fix further the gauge $`_MA^M=0`$ since according to the equations $`A_M`$ also is homogeneous of degree $`1`$ in this gauge. The three equations satisfied by the gauge field are now
$$XA\left(X\right)=0,\left(X+1\right)A_M\left(X\right)=0,_MA^M=0.$$
(28)
There still remains gauge symmetry in these equations for $`\mathrm{\Lambda }`$ that satisfy $`X\mathrm{\Lambda }=\mathrm{\Lambda }=0.`$ The content of these equations for $`\mathrm{\Lambda }`$ is still non-trivial.
These equations were proposed by Dirac in 1936 as subsidiary conditions to describe the usual 4-dimensional Maxwell theory of electromagnetism (in the Lorentz gauge), as a theory in 6 dimensions which automatically displays SO$`(4,2)`$ symmetry. Dirac’s aim was to linearize the conformal symmetry of the 4 dimensional Maxwell theory. The subsidiary conditions can be regarded as “kinematics” while dynamics is given by a Klein-Gordon type equation in 6-dimensions that may include interactions with other fields. As Dirac showed, the linear SO$`(4,2)`$ Lorentz symmetry of the 6 dimensional theory is indeed the non-linear conformal symmetry of the Maxwell theory.
Actually, in the framework of two-time physics, conformal symmetry is only one of the possible interpretations of the SO$`(4,2)`$ global symmetry of these equations. In two-time physics this interpretation relies on a particular choice of “time” among the two available timelike dimensions, while with other gauge choices the interpretation of the SO$`(4,2)`$ symmetry is completely different than conformal symmetry. To illustrate this, denote the components of the 6 dimensions as $`X^M=(X^+^{},X^{^{}},X^\mu )`$ with metric $`XX=2X^+^{}X^{^{}}+X_\mu X^\mu .`$ The Sp$`(2,R)`$ gauge choices $`P^+^{}\left(\tau \right)=0`$, $`X^+^{}\left(\tau \right)=1`$ eliminates one timelike and one spacelike dimensions and brings down the two-time formulation in $`d+2`$ dimension to a one time formulation in $`d`$ dimensions. It is convenient to use the electromagnetic gauge choice $`A^+^{}\left(X\right)=0`$ (instead of Dirac’s $`_MA^M=0`$). Then the solution of the gauge choices and constraints (including $`Q_{11}=Q_{12}=0`$), $`XX=XP=XA=0,`$ is given in the following form
$`X^M\left(\tau \right)`$ $`=`$ $`(1,x^2/2,x^\mu \left(\tau \right)),P^M=(0,xp,p^\mu \left(\tau \right)),`$ (29)
$`A^M\left(X\right)`$ $`=`$ $`(\mathrm{\hspace{0.17em}0},xA,A^\mu \left(x\left(\tau \right)\right)).`$ (30)
The dynamics of the remaining degrees of freedom $`(x^\mu \left(\tau \right),p^\mu \left(\tau \right))`$ are obtained by substituting these solutions into the gauge invariant 6-dimensional action (20). The result is the standard 4-dimensional action for the massless relativistic particle coupled to the electromagnetic gauge potential $`A_\mu \left(x\right)`$
$$L=\frac{1}{2A^{22}}\left(\dot{x}^\mu \right)^2\dot{x}^\mu A_\mu \left(x\right).$$
(31)
Thus the original two-time action displays explicitly the hidden SO$`(4,2)`$ symmetry of the one-time action. The general coordinate transformation of the previous section, specialized to $`\epsilon ^M=\epsilon ^{MN}X_N`$ with constant antisymmetric $`\epsilon ^{MN},`$ is the SO$`(4,2)`$ global Lorentz symmetry of the 6-dimensional action, including the gauge field. This 6-dimensional Lorentz symmetry is also the non-linearly realized conformal symmetry of the gauge fixed action above, since the global symmetry commutes with the gauge symmetry, and gauge fixing of the gauge invariant action could not destroy the global symmetry. Indeed the generators of conformal transformations are the gauge invariant $`L^{MN}=X^MP^NX^NP^N`$ now expressed in terms of the gauge fixed coordinates and momenta as shown in . This agrees with Dirac’s interpretation of the conformal SO$`(4,2)`$ symmetry as being the Lorentz symmetry of 6 dimensions.
However, if one chooses another gauge for time instead of $`X^+^{}\left(\tau \right)=1`$, as was done with many illustrations in , other $`d`$-dimensional dynamical systems arise, which now are coupled to a gauge potential. Then the SO$`(d,2)`$ symmetry generated by the same $`L^{MN}`$ has a different interpretation than conformal symmetry, as explained in . The presence of the gauge field background now produces a large class of dynamical systems with hidden SO$`(d,2)`$ symmetries, and Sp$`(2,R)`$ duality relations among them.
The two-time physics approach - was developed without being aware of the field equations invented by Dirac. While Dirac was interested in linearizing conformal symmetry<sup>6</sup><sup>6</sup>6I thank Vasilev for informing me of Dirac’s work and the line of research that followed the same trend of thought in relation to conformal symmetry . A field theoretic formulation of two-time physics has been derived recently and its relation to Dirac’s work has been established. It is shown in that two-time physics in a field theoretic setting, as in the particle dynamics setting, unifies different looking one-time field theories as being the same two-time field theory, while simultaneously revealing previously unnoticed hidden symmetries in field theory, including interactions. Such duality and global symmetry properties of two-time physics go well beyond Dirac’s goal of linearizing conformal symmetry., the motivation for the work in - came independently from duality, and signals for two-timelike dimensions in M-theory and its extended superalgebra including D-branes . Driven by different motivations, and unaware of Dirac’s approach to conformal symmetry, two-time physics produced new insights that include conformal symmetry but go well beyond it. Besides providing a deeper Sp$`(2,R)`$ gauge symmetry as the fundamental basis for Dirac’s approach (see further ), two-time physics unifies classes of one-time physical systems in $`d`$ dimensions that previously would have been thought of as being unrelated to each other. The SO$`(d,2)`$ symmetry is interpreted as conformal symmetry in a certain one-time system, but in other dually related dynamical systems it is a hidden symmetry with a different interpretation, but realized in exactly the same irreducible representation. . The unifying aspect in all the interpretations is that the symmetry is the underlying spacetime symmetry in a spacetime that includes two timelike dimensions.
## 5 Gravitational background
We now seek a solution of (23-27) that includes gravity in $`d`$ dimensions. It is convenient to make a change of variables $`X^M=X^M(\kappa ,w,x^\mu )`$ such that the function $`W\left(X\right)`$ is identified with the product of new coordinates $`2w\kappa `$, while the coordinate $`x^\mu `$ is in $`d`$ dimensions. The inverse of this change of variables is, $`\kappa =K\left(X\right),`$ $`w=W\left(X\right)/2K\left(X\right)`$ and $`x^\mu =x^\mu \left(X\right).`$ Before we look for a solution to (23-27) it is instructive to consider the example of the flat case that has components $`X^M=(X^+^{},X^{^{}},X^\mu )`$ with the constraint $`W\left(X\right)=XX=2X^+^{}X^{^{}}+X_\mu X^\mu .`$ The change of variables and the inverse relations for this case are
$`X^+^{}`$ $`=`$ $`\kappa ,X^{^{}}={\displaystyle \frac{\kappa x^2}{2}}+w,X^\mu =\kappa x^\mu `$ (32)
$`\kappa `$ $`=`$ $`X^+^{},w={\displaystyle \frac{XX}{2X^+^{}}},x^\mu ={\displaystyle \frac{X^\mu }{X^+^{}}}`$ (33)
This change of variables is a special case of a general coordinate transformation. The flat metric in the new variables takes the form
$`ds^2`$ $`=`$ $`dX^MdX^N\eta _{MN}=2dX^+^{}dX^{^{}}+dX^\mu dX^\nu \eta _{\mu \nu }`$ (34)
$`=`$ $`2d\kappa dw+\kappa ^2dx^\mu dx^\nu \eta _{\mu \nu }.`$ (35)
For this choice of basis we have $`V^M=(\kappa ,w,0)`$ and $`W=2\kappa w`$ and the homothety conditions are easily verified. Taking this form as a model we seek a similar solution. With a choice of coordinates we can always take $`V^M=(\kappa ,w,0)`$. In the new coordinate system $`W(\kappa ,w,x^\mu )`$ needs to be determined consistently with the closed homothety conditions. We will make an ansatz which may not be the most general, but is adequate to provide a sufficiently large set of solutions. Thus, we will take $`W(\kappa ,w,x)=2w\kappa `$ to have the same form as the free case, and insert these forms of $`V,W`$ in the closed homothety conditions with a general $`G^{MN}`$. The homothety condition reads
$$\left(\kappa _\kappa +w_w\right)G^{MN}\delta _\kappa ^MG^{\kappa N}\delta _w^MG^{wN}\delta _\kappa ^NG^{\kappa M}\delta _w^NG^{wM}=2G^{MN}.$$
(36)
From $`V^M=\frac{1}{2}G^{MN}_NW\left(X\right)`$ we learn further
$`V^\mu `$ $`=`$ $`0=G^{\mu \kappa }wG^{\nu w}\kappa G^{\mu \kappa }={\displaystyle \frac{1}{\kappa }}W^\mu ,G^{\mu w}={\displaystyle \frac{w}{\kappa ^2}}W^\mu ,`$ (37)
$`V^\kappa `$ $`=`$ $`\kappa =G^{\kappa \kappa }wG^{\kappa w}\kappa ,G^{\kappa \kappa }={\displaystyle \frac{\kappa }{w}}(1+G^{\kappa w})`$ (38)
$`V^w`$ $`=`$ $`w=G^{w\kappa }wG^{ww}\kappa ,G^{ww}={\displaystyle \frac{w}{\kappa }}(1+G^{\kappa w})`$ (39)
Specializing the indices in the homothety condition gives the solutions for all components of $`G^{MN}`$ in the form
$$G^{MN}=\left(\begin{array}{ccc}\frac{\kappa }{w}\left(\gamma 1\right)& \gamma & \frac{1}{\kappa }W^\nu \\ \gamma & \frac{w}{\kappa }\left(\gamma 1\right)& \frac{w}{\kappa ^2}W^\nu \\ \frac{1}{\kappa }W^\mu & \frac{w}{\kappa ^2}W^\mu & \frac{g^{\mu \nu }}{\kappa ^2}\end{array}\right)$$
(40)
where the functions $`\gamma (x,\frac{w}{\kappa }),`$ $`W^\mu (x,\frac{w}{\kappa }),`$ $`g^{\mu \nu }(x,\frac{w}{\kappa })`$ are arbitrary functions of only $`x^\mu `$ and the ratio $`\frac{w}{\kappa }`$.
In this coordinate system we can also solve the kinematic conditions for the gauge field (27), which become
$$\left(w_w+\kappa _\kappa \right)A_M+\delta _M^wA_w+\delta _M^\kappa A_\kappa =0,wA_w+\kappa A_\kappa =0.$$
(41)
The general solution is
$$A_w=\frac{1}{\kappa }B(\frac{w}{\kappa },x),A_\kappa =\frac{w}{\kappa ^2}B(\frac{w}{\kappa },x),A_\mu =A_\mu (\frac{w}{\kappa },x).$$
The remaining gauge symmetry $`V^M_M\mathrm{\Lambda }=0`$ is just sufficient to set $`B=0`$ in this solution, if so desired. Finally the solution for $`U(w,\kappa ,x)`$ that satisfies $`V^M_MU=2U`$ is
$$U=\frac{1}{\kappa ^2}u(\frac{w}{\kappa },x).$$
(42)
For this solution, the generators of Sp$`(2,R)`$ in $`(\text{18},\text{19})`$ become, in the gauge $`B=0`$,
$`Q_{11}`$ $`=`$ $`2\kappa w,Q_{12}=\kappa p_\kappa +wp_w,`$ (43)
$`Q_{22}`$ $`=`$ $`2\gamma p_wp_\kappa +\left(p_\kappa ^2{\displaystyle \frac{\kappa }{w}}+p_w^2{\displaystyle \frac{w}{\kappa }}\right)\left(\gamma 1\right)`$
$`+{\displaystyle \frac{2}{\kappa ^2}}\left(\kappa p_\kappa wp_w\right)W^\mu p_\mu +{\displaystyle \frac{H}{\kappa ^2}},`$
where
$$H=\frac{1}{\sqrt{g}}\left(p_\mu +A_\mu \right)\sqrt{g}g^{\mu \nu }\left(p_\nu +A_\nu \right)+u.$$
(45)
It is easy to verify directly that they close correctly for any background fields $`\gamma ,g_{\mu \nu },W^\mu ,A_\mu ,u`$ that are arbitrary functions of $`(\frac{w}{\kappa },x^\mu )`$.
Imposing the Sp$`(2,R)`$ constraints $`Q_{ij}=0`$ is now easy. It is convenient to choose a Sp$`(2,R)`$ gauge, which we know will produce a one-time theory. A gauge choice that is closely related to the massless relativistic particle is taken by analogy to the flat theory. At the classical level we choose the Sp$`(2,R)`$ gauges $`\kappa \left(\tau \right)=1`$ and $`p_w\left(\tau \right)=0,`$ and solve $`Q_{11}=Q_{12}=0`$ in the form $`w\left(\tau \right)=p_\kappa \left(\tau \right)=0`$. There remains unfixed one gauge subgroup of Sp$`(2,R)`$ which corresponds to $`\tau `$ reparametrization, and the corresponding Hamiltonian constraint $`H0`$, which involves the background fields $`g_{\mu \nu }\left(x\right),`$ $`A_\mu \left(x\right),`$ $`u\left(x\right)`$ that now are functions of only the $`d`$ dimensional coordinates $`x^\mu ,`$ since $`w/\kappa =0.`$ In this gauge, the background fields $`\gamma ,W^\mu `$ decouple from the dynamics that govern the time development of $`x^\mu \left(\tau \right).`$ The two-time theory described by the original Lagrangian (20) reduces to a one-time theory
$$L=\frac{1}{2A^{22}}\dot{x}^\mu \dot{x}^\nu g_{\mu \nu }\left(x\right)\frac{A^{22}}{2}u\left(x\right)\dot{x}^\mu A_\mu \left(x\right).$$
which controls the dynamics of the remaining degrees of freedom $`x^\mu \left(\tau \right).`$ Evidently this Lagrangian describes a particle moving in arbitrary gravitational, electromagnetic gauge fields and other potential $`g_{\mu \nu }\left(x\right),`$ $`A_\mu \left(x\right)`$, $`u\left(x\right)`$ in the remaining $`d`$ dimensional spacetime.
We have therefore demonstrated that all usual interactions experienced by a particle, as described in the one-time formulation of dynamics, can be embedded in two time physics as a natural solution of the two-time equations (23-27), taken in a fixed Sp$`(2,R)`$ gauge.
## 6 Spinning particles in background fields
To describe spinning particles in two time physics we need local superconformal symmetry instead of local conformal symmetry, as demonstrated in flat space in . There the Sp$`(2,R)`$ gauge group was replaced by the supergroup OSp$`\left(n|2\right)`$ as described at the end of section 2 of this paper. To generalize this approach to curved space we need a soldering form $`E_M^a`$ and its inverse $`E_a^M`$ (analog of vierbein) that transforms curved base space indices to flat tangent space indices and vice versa. The metric in tangent space is $`\eta _{ab}`$ while the general metric is given by $`G_{MN}=E_M^aE_N^b\eta _{ab}.`$ Next consider phase space including spin degrees of freedom $`(X^M,P_M,\psi _\alpha ^a)`$ where $`a`$ is a tangent space index and $`\alpha =1,2,\mathrm{},n`$ denote the $`n`$ supersymmetries. The canonical commutation rules are
$$[X^M,P_N]=i\delta _N^M,\{\psi _\alpha ^a,\psi _\beta ^b\}=\eta ^{ab}\delta _{\alpha \beta }.$$
(46)
The $`\psi _\alpha ^a`$ form a Clifford algebra and may be represented by gamma matrices if so desired.
A Lagrangian that has the desired OSp$`\left(n|2\right)`$ local symmetry has the same form as the flat case given in with some modifications
$$L=\dot{X}^MP_M+\frac{i}{2}\psi _\alpha ^a\dot{\psi }_\alpha ^b\eta _{ab}\frac{1}{2}A^{ij}Q_{ij}+iF^{i\alpha }Q_{i\alpha }\frac{1}{2}B^{\alpha \beta }Q_{\alpha \beta },$$
(47)
The OSp$`\left(n|2\right)`$ gauge fields may be arranged into the form of a $`\left(n+2\right)\times \left(n+2\right)`$ supermatrix
$$\left(\begin{array}{cc}B^{\left[\alpha \beta \right]}\hfill & F^{\alpha i}\hfill \\ \epsilon _{ij}F^{j\beta }\hfill & A^{ij}\hfill \end{array}\right),A,B=bose,F=fermi$$
(48)
They obey the standard transformation rules for gauge fields, as given in . The OSp$`\left(n|2\right)`$ generators $`Q_{ij},Q_{i\alpha },Q_{\alpha \beta }`$ are to be taken as non-linear functions in phase space, including background fields. As in the purely bosonic case, our task is to find the forms of the background fields that have an interpretation as gravitational, gauge or other interactions experienced by spinning particles in two-time physics. The gauge field equations of motion require the first class constraints $`Q_{ij}Q_{i\alpha }Q_{\alpha \beta }0,`$ whose solution will require two timelike dimensions, as in the flat theory or as in the curved purely bosonic theory. These are then the generators of infinitesimal transformations that tell us how to transform $`\delta X^M,\delta P_M,\delta \psi _\alpha ^M`$ under the local OSp$`\left(n|2\right).`$ As in the purely bosonic theory treated earlier in this paper, it is easy to show that the Lagrangian has the local symmetry provided these first class constraints close into the algebra of OSp$`\left(n|2\right).`$ This requirement gives the differential equations for the background fields.
In the flat case the OSp$`\left(n|2\right)`$ generators are given by $`Q_{ij}^0=X_iX_j,`$ $`Q_{i\alpha }^0=X_i\psi _\alpha ,`$ and $`Q_{\alpha \beta }^0=\frac{i}{2}\psi _{[\alpha }\psi _{\beta ]}.`$ To include background fields we first generalize the fermionic generators $`P\psi _\alpha `$ ($`n`$ local supersymmetries) and $`X\psi _\alpha `$ ($`n`$ local superconformal symmetries) by introducing a tangent space vector $`V_a\left(X\right)`$, a soldering from $`E_M^a\left(X\right),`$ a spin connection $`\omega _M^{ab}\left(X\right),`$ a gauge field $`A_M\left(X\right),`$ and replacing the momentum by the covariant momentum
$$\mathrm{\Pi }_a(X,P,\psi )=E_a^M\left(P_M+A_M+\frac{1}{2}\omega _M^{ab}S_{ab}\right)$$
(49)
The spin connection, which generally includes torsion, is coupled to the spin operator $`S^{ab}=\frac{1}{2i}\left(\psi _\alpha ^a\psi _\alpha ^b\psi _\alpha ^b\psi _\alpha ^a\right)`$ to form the covariant momentum. The generalized fermionic generators are as follows
$$Q_{1\alpha }=\psi _\alpha ^aV_a\left(X\right),Q_{2\alpha }=\frac{1}{2}\left(\psi _\alpha ^a\mathrm{\Pi }_a+\stackrel{~}{\mathrm{\Pi }}_a\psi _\alpha ^a\right).$$
(50)
The bosonic generators are computed from the closure of the OSp$`\left(n|2\right)`$ commutation relations
$$\{Q_{1\alpha },Q_{1\beta }\}=\delta _{\alpha \beta }Q_{11},\{Q_{2\alpha },Q_{2\beta }\}=\delta _{\alpha \beta }Q_{22},\{Q_{1\alpha },Q_{2\beta }\}=\delta _{\alpha \beta }Q_{12}+Q_{\alpha \beta }.$$
(51)
where $`Q_{\alpha \beta }`$ is the antisymmetric SO$`\left(n\right)`$ generator, and $`Q_{ij}`$ are the symmetric Sp$`\left(2\right)`$ generators. Note that $`Q_{2\alpha }`$ contains up to cubic terms in the fermions. $`\stackrel{~}{\mathrm{\Pi }}_a`$ is given by $`\stackrel{~}{\mathrm{\Pi }}_a=\left(\sqrt{G}\right)^1\mathrm{\Pi }_a\sqrt{G}`$, where the factors of $`\sqrt{G}`$ insure hermiticity in a quantum theory with correct factor ordering, but for the invariance of the classical action, where we only need Poisson brackets instead of the commutators as explained in the spinless case, these factors may be neglected.
For simplicity we will impose the flat $`Q_{\alpha \beta }=Q_{\alpha \beta }^0`$
$$Q_{\alpha \beta }=\frac{i}{2}\psi _{[\alpha }\psi _{\beta ]}$$
(52)
but will compute $`Q_{ij}`$ as a function of the background fields<sup>7</sup><sup>7</sup>7We could have included also $`E_I^MW_M^{ab}Q_{ab}^0`$ as part of $`\mathrm{\Pi }_I,`$ with $`W_M^{ab}`$ a gauge field that acts in the SO$`\left(n\right)`$ space within OSp$`\left(n|2\right).`$ In that case we could also introduce a vielbein $`E_a^A`$ for an internal space. For simplicity we will omit these complications and seek a solution with a “flat” SO$`\left(n\right)`$ space, implying that the metric in SO$`\left(n\right)`$ space is $`\delta _{ab}`$ instead of a curved space metric $`G_{AB}=E_A^aE_B^a.`$ Recall that in the final analysis we are interested in imposing $`Q_{ab}=0`$ as part of the singlet condition. I the presence of non-singlet background fields such as $`E_A^a,W_M^{ab}`$ this condition is harder to satisfy.. This condition requires that $`E_M^a`$ be determined in terms of $`V^a,\omega _M^{ab}`$
$$E_M^a=D_MV^a=_MV^a+\omega _M^{ab}V_b,$$
(53)
while
$$Q_{11}=V^aV^b\eta _{ab},Q_{12}=\frac{1}{2}\left(V^a\mathrm{\Pi }_a+\stackrel{~}{\mathrm{\Pi }}_aV^a\right),Q_{22}=\frac{1}{n}\left[\frac{1}{2}\left(\psi _\alpha ^a\mathrm{\Pi }_a+\stackrel{~}{\mathrm{\Pi }}_a\psi _\alpha ^a\right)\right]^2.$$
(54)
Note that $`Q_{22}`$ contains several powers of the fermions. The closure (51) is possible provided the gauge field strength and the curvature are transverse to $`V`$
$$V^MF_{MN}=0,V^MR_{MN}^{ab}=0,$$
(55)
where
$$V^M=E_a^MV^a,$$
(56)
and
$$F_{MN}=_MA_N_NA_M+[A_M,A_N],R_{MN}^{ab}=_M\omega _N^{ab}_N\omega _M^{ab}+[\omega _M,\omega _N]^{ab}.$$
(57)
Furthermore, since $`E_M^a=D_MV^a`$ the torsion is determined in terms of the curvature and $`V`$ as
$$T_{MN}^a=D_ME_N^aD_NE_M^a=R_{MN}^{ab}V_b,$$
(58)
and is automatically transverse to $`V`$ provided the curvature is.
There remains to check the Sp$`\left(2\right)\times SO\left(n\right)`$ closure of the bosonic generators. The SO$`\left(n\right)`$ part is trivial. The Sp$`\left(2\right)`$ part is similar to the purely bosonic case of the previous section and is subject to the same conditions (23)-(25) discussed there. However now $`W,G^{MN}`$ are given by $`W=V^aV_a`$ and $`G_{MN}=E_M^aE_N^b\eta _{ab}`$ and $`U=0.`$ These forms automatically satisfy (23)-(25) provided $`E_M^a`$ is of the form (53). In particular, (23) is satisfied as follows
$$V^M=\frac{1}{2}G^{MN}_NW=G^{MN}\left(D_NV^a\right)V_a=G^{MN}E_N^aV_a=E_b^MV^b$$
(59)
which agrees with the definition (56). Meanwhile, the homothety condition (25) is equivalent to
$$\mathrm{\pounds }_VE_M^a=E_M^a$$
(60)
where $`\mathrm{\pounds }_VE_M^a=V^ND_NE_M^a+_MV^NE_N^a.`$ This is also satisfied automatically for the geometry constructed above in terms of $`V^a`$ and $`\omega _M^{ab}`$ as follows
$`\mathrm{\pounds }_VE_M^a`$ $`=`$ $`V^ND_NE_M^a+_MV^NE_N^a=V^NT_{NM}^a+V^ND_ME_N^a+_MV^NE_N^a`$ (61)
$`=`$ $`V^NT_{NM}^a+D_M\left(V^NE_N^a\right)=V^NT_{NM}^a+D_MV^a`$ (62)
$`=`$ $`E_M^a`$ (63)
where we have used the orthogonality of $`V`$ to the curvature or torsion. Related equations appear in , but our approach provides a OSp$`\left(n|2\right)`$ gauge symmetry basis for introducing Eq.(53) and the rest of the geometrical equations. Also, a similar problem was discussed in in a less geometrical formalism and in the absence of the gauge field $`A_M.`$ In our case we are interested in solutions of the equations that permit the imposition of the constraints $`Q_{ij}Q_{i\alpha }Q_{\alpha \beta }0.`$
The geometry described by $`E_M^a`$ is fully determined by the functions $`\omega _M^{ab}\left(X\right)`$ and $`V^a\left(X\right)`$ which are constrained only by the transversality condition $`V^MR_{MN}^{ab}=0,`$ but are otherwise arbitrary. The solution space includes the most general gravitational metric in $`d`$ dimensions as already seen in the previous section. The formalism in this section provides a more covariant solution and permits the construction of the general interacting two-time physics for spinning particles.
## 7 Conclusion and discussion
The choice of coordinates $`\kappa ,w,x^\mu `$ and solution of background fields used above emphasizes a basis that is convenient for deriving the free massless relativistic particle from two time physics in the case of zero background fields. In this basis it was easy to eliminate one timelike and one spacelike coordinates through the gauge choices $`\kappa \left(\tau \right)=1,`$ $`p_w=0,`$ leaving the usual timelike coordinate as a component of the $`d`$-dimensional vector $`x^\mu \left(\tau \right).`$ With this choice of time we interpreted the theory and the background fields, as discussed above. However, as we have already seen in the flat case, other choices of the time coordinate produce very different physical interpretations from the point of view of the one-time observer, even though the two time physics theory is the same. In the general theory it is also possible to work in other coordinates that are convenient to solve the Sp$`(2,R)`$ constraints in other Sp$`(2,R)`$ gauges. Then the choice of “time” embedded in the two-time theory is different.
It follows that the same background fields given above would give rise to very different interpretation of the dynamics in one-time physics in different Sp$`(2,R)`$ gauges. For example, in the flat spinless case, with $`\gamma =g_{\mu \nu }=W^\mu =A_\mu =u=0,`$ different Sp$`(2,R)`$ gauges produced a class of related one-time dynamics that included the free massless relativistic particle, the free massive relativistic particle, the free massive non-relativistic particle, the H-atom, the harmonic oscillator in one less dimension, the particle in AdS$`{}_{dk}{}^{}\times `$S<sup>k</sup> backgrounds for any $`k=0,1,\mathrm{},d2,`$ etc. In a similar way, in the general theory all possible choices of time define a class of one-time dynamical theories related to the same two-time dynamics with a fixed set of background fields. Changing the background fields changes the class of related one-time dynamical models.
In the flat case the global symmetry was SO$`(d,2).`$ In the general case the Killing symmetries of the background fields (which is embedded in the general coordinate and gauge transformations) replaces the global SO$`(d,2)`$ symmetry. The global symmetries should be realized in the same representation for all of the different one-time dynamical models in the same class derived from the same two-time physics theory.
## 8 Acknowledgements
I thank E. Witten for a discussion on the homothety conditions.
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# Stability of Oscillating Hexagons in Rotating Convection
## 1 Introduction
Convection has played a key role in the elucidation of the spatio-temporal dynamics arising in non-equilibrium pattern forming systems. The interplay of well-controlled experiments with analytical and numerical theoretical work has contributed to a better understanding of various mechanisms that can lead to complex behavior. From a theoretical point of view the effect of rotation on roll convection has been particularly interesting because it can lead to spatio-temporal chaos immediately above threshold where the small amplitude of the pattern allows a simplified treatment. Early work of Küppers and Lortz showed that for sufficiently large rotation rate the roll pattern becomes unstable to another set of rolls rotated with respect to the initial one. Due to isotropy the new set of rolls is also unstable and persistent dynamics are expected. Later Busse and Heikes confirmed experimentally the existence of this instability and the persistent dynamics arising from it. They proposed an idealized model of three coupled amplitude equations in which the instability leads to a heteroclinic cycle connecting three sets of rolls rotated by 120<sup>o</sup> with respect to each other. Recently the Küppers-Lortz instability and the ensuing dynamics have been subject to intensive research, both experimentally and theoretically . It is found that in sufficiently large systems the switching between rolls of different orientation looses coherence and the pattern breaks up into patches in which the rolls change orientation at different times. The shape and size of the patches changes persistently due to the motion of the fronts separating them.
In this paper we are interested in the effect of rotation on hexagonal rather than roll (stripe) patterns as they arise in systems with broken up-down symmetry (e.g. non-Boussinesq or surface-tension driven convection with rotation). The dynamics of strictly periodic hexagon patterns with broken chiral symmetry have been investigated in detail by Swift and Soward . They found that the heteroclinic orbit of the Busse-Heikes model is replaced by a periodic orbit arising from a secondary Hopf bifurcation off the hexagons. Their results have been confirmed in numerical simulations of a Swift-Hohenberg-type model in which an alternation among the three modes that compose the hexagonal pattern is observed. In the following we will call this state ‘oscillating hexagons’. The oscillations can be homogeneous in space or can take on the form of traveling waves. Starting from coupled Ginzburg-Landau equations, we have previously derived evolution equations for the oscillating hexagons that are valid close to the Hopf bifurcation. Within this framework the oscillating hexagons were found to support a state of spatio-temporal chaos that is characterized by defects, in which the oscillation amplitude vanishes . At the band center the oscillating hexagons and their chaotic state are described by the single complex Ginzburg-Landau equation (CGLE). In general, however, the oscillation amplitude is coupled to the phases of the underlying pattern as it is to be expected for a secondary bifurcation. Here we study how this coupling affects the stability of the oscillations. In particular, it will modify the stability properties of the waves emitted by the spirals in the defect chaotic state. We find that the additional coupling leads to new long-wave instabilities if the rotation is strong enough or if the wavenumber of the hexagons is far away from the band-center.<sup>1</sup><sup>1</sup>1For stripe patterns (e.g. convection rolls) the generic equation for a secondary Hopf bifurcation is well known , and has been studied both theoretically and experimentally . There, the coupling to the phase of the pattern can delay the occurrence of long-wave instabilities of the oscillatory mode.
The paper is organized as follows. In the second section we introduce the appropriate coupled Ginzburg-Landau equations that describe the hexagon patterns and derive the amplitude-phase equations close to the secondary Hopf bifurcation. The stability analysis of these equations is addressed in section III. In section IV we study numerically the behavior resulting from these instabilities. Conclusions are given in section V. Details of the calculations are given in two appendices.
## 2 Amplitude-Phase equations for Oscillating Hexagons
We consider small-amplitude hexagon patterns in systems with broken chiral symmetry. In order to analyze the possibility of modulational instabilities we include spatial derivatives. Due to the strong coupling between modes of different orientation, we take the gradients in both directions to be of the same order and retain only linear gradient terms (a study of the influence of nonlinear gradient terms in the stability of oscillating hexagons is addressed in appendix B). After rescaling the amplitude, time, and space we arrive at the equations,
$`_tA_1`$ $`=`$ $`\mu A_1+(𝐧_1)^2A_1+\overline{A}_2\overline{A}_3A_1|A_1|^2`$
$`(\nu +\gamma )A_1|A_2|^2(\nu \gamma )A_1|A_3|^2,`$
where the equations for the other two amplitudes are obtained by cyclic permutation of the indices and $`\mu `$ is a parameter related to the distance from threshold. The overbar represents complex conjugation. These equations can be obtained from the corresponding physical equations (e.g. Navier-Stokes) using a perturbative technique. The broken chiral symmetry manifests itself by the cross-coupling coefficients not being equal. Hence $`\gamma `$ is a measure of rotation.
For completeness it should be noted that rotation leads in convection not only to a chiral symmetry breaking but also to a (weak) breaking of the translation symmetry due to the centrifugal force. In the following we will consider it to be negligible. We focus on not too small Prandtl numbers, in which case the primary bifurcation is always steady .
Equation (2) admits hexagon solutions $`A_j=Re^{iq\widehat{𝐧}_j𝐱+i\varphi _j}`$ with a slightly off-critical wavenumber ($`\stackrel{~}{𝐪}_j=\stackrel{~}{𝐪}_j^c+𝐪_j`$, $`|𝐪_j||\stackrel{~}{𝐪}_c|\stackrel{~}{q}_c`$), with
$$R=\frac{1\pm \sqrt{1+4(\mu q^2)(1+2\nu )}}{2(1+2\nu )},\mathrm{\Phi }\varphi _1+\varphi _2+\varphi _3=0.$$
(2)
The stability of this solution to perturbations with the same wave vectors has been studied by several authors . Typical results are sketched in the bifurcation diagrams shown in Fig. 1. The hexagons appear through a saddle-node bifurcation at $`\mu =\mu _{sn}`$,
$$\mu _{sn}=\frac{1}{4(1+2\nu )}+q^2,$$
(3)
and become unstable via a Hopf bifurcation at $`\mu =\mu _H`$, with a frequency $`\omega _H`$,
$$\mu _H=\frac{(2+\nu )}{(\nu 1)^2}+q^2,\omega _H=2\sqrt{3}\gamma /(\nu 1)^2.$$
(4)
The Hopf bifurcation is supercritical and for $`\mu >\mu _H`$ stable oscillations in the three amplitudes of the hexagonal pattern arise with a phase shift of $`2\pi /3`$ between them , resulting in what we call oscillating hexagons. As $`\mu `$ is increased further, eventually a point $`\mu =\mu _{het}`$ is reached at which the branch of oscillating hexagons ends on the branch corresponding to a mixed-mode solution in a global bifurcation involving a heteroclinic connection (see Fig. 1a). Above this point the only stable solution is the roll solution whose stability region is bounded below by
$$\mu _R=\frac{1}{(\nu +\gamma 1)(\nu \gamma 1)}+q^2.$$
(5)
From Eqs. (4), (5) it is easy to see that, when $`\gamma ^2>(\nu 1)(\nu +1)/(\nu +2)`$, the transition to oscillating hexagons occurs at a value of $`\mu `$ for which the rolls are still unstable. There is then a parameter regime in which the oscillating hexagons are the only stable solution. Furthermore, when $`|\gamma |>\nu 1\gamma _{KL}`$ rolls are never stable and the limit cycle persists for arbitrarily large values of $`\mu `$ (see Fig. 1b). In the absence of the quadratic terms in Eq. (2) this condition corresponds to the Küppers-Lortz instability of rolls. When the quadratic term in Eq. (2) is small (i.e. small non-Boussinesq effects) it can be considered as a perturbation of the usual three mode model for rotating roll convection. Far above the Hopf bifurcation ($`\mu \mu _H`$) the periodic orbit is expected to become asymmetrical and the resulting state similar to that encountered in the usual rotating Rayleigh-Bénard convection.
The stability of steady hexagons with respect to side-band perturbations in the presence of rotation has been studied previously . Due to the rotation steady and oscillatory, long-wave and short-wave instabilities have been found. It turns out that in the presence of rotation the steady hexagons can be stable up to the Hopf bifurcation over quite a range of wavenumbers. Thus, the oscillating hexagons may be stable near the Hopf bifurcation.
The focus of the present paper are side-band instabilities of the oscillating hexagons. To address this question analytically we focus on the vicinity of the Hopf bifurcation where the oscillation amplitude is small and a weakly nonlinear analysis is possible. Since we are dealing with a secondary bifurcation the two phases of the underlying hexagons have to be taken into account as well. In fact, from a linear stability analysis of Eq. (2) it is easy to see that there are four marginal modes at $`\mu =\mu _H`$. Two correspond to the Hopf bifurcation and can be described by a complex amplitude $``$. The other two correspond to a two-dimensional phase vector $`\stackrel{}{\varphi }`$, related to translations in the x- and y-directions . It can be written as a combination of the phases of the three modes: $`\stackrel{}{\varphi }=(\varphi _x,\varphi _y)(\varphi _2\varphi _3,(\varphi _2\varphi _3)/\sqrt{3})`$, satisfying the locking condition $`\varphi _1+\varphi _2+\varphi _3=\mathrm{\Phi }`$. The global phase $`\mathrm{\Phi }`$ is a fast variable that relaxes rapidly to its stationary values $`\mathrm{\Phi }=0,\pi `$ (up or down hexagons). To study the nonlinear behavior of the oscillating hexagons, the amplitudes $`A_i`$ are expanded as:
$$A_n=(R_H+e^{2\pi ni/3}\sqrt{ϵ}e^{i\omega _Ht}+c.c.+𝒪(ϵ))e^{i𝐪_i𝐱+\sqrt{ϵ}\varphi _i},$$
(6)
where $`ϵ`$ is a small parameter related to the distance from the bifurcation line. The amplitude of the steady hexagons at the bifurcation point, $`R_H=1/(\nu 1)`$, is independent of $`\mu `$ and $`q`$. Eliminating the fast variables, at order $`ϵ^{3/2}`$ (see appendix A) we arrive at an equation for the amplitude of the oscillation $``$, coupled to the phase vector of the underlying hexagonal pattern,
$`_T`$ $`=`$ $`\epsilon \delta _1+\xi ^2\delta _2\stackrel{}{\varphi }\rho ||^2,`$ (7)
$`_T\stackrel{}{\varphi }`$ $`=`$ $`D_{}^2\stackrel{}{\varphi }+D_{}(\stackrel{}{\varphi })+D_{\times _1}(𝐞_z\times ^2\stackrel{}{\varphi })`$
$`+D_{\times _2}(𝐞_z\times )(\stackrel{}{\varphi })+\alpha ||^2+\beta _1(𝐞_z\times )||^2`$
$`i\beta _2(\overline{}\overline{})+i\eta [(𝐞_z\times )\overline{}\overline{}(𝐞_z\times )],`$
with $`\epsilon =\mu \mu _H`$ and
$`v=3R_H(1+2R_H),`$ (9)
$`\delta _1={\displaystyle \frac{2R_H}{v}}{\displaystyle \frac{2i\omega _H}{v}},\delta _2=q\delta _1,`$ (10)
$`\xi ={\displaystyle \frac{1}{2}}{\displaystyle \frac{3q^2R_H}{9R_H^2+\omega _H^2}}{\displaystyle \frac{iq^2}{\omega _H}}{\displaystyle \frac{9R_H^2+2\omega _H^2}{9R_H^2+\omega _H^2}},`$ (11)
$`\rho ={\displaystyle \frac{8(3R_H+1)}{v}}{\displaystyle \frac{4i\omega _H(1+4R_H)}{R_Hv}}{\displaystyle \frac{32i}{3\omega _H}},`$ (12)
$`D_{}={\displaystyle \frac{1}{4}},D_{}={\displaystyle \frac{1}{2}}{\displaystyle \frac{2q^2}{v}},D_{\times _1}={\displaystyle \frac{q^2}{\omega _H}},D_{\times _2}=0,`$ (13)
$`\alpha ={\displaystyle \frac{2\omega _H^2q}{9R_H^2+\omega _H^2}}{\displaystyle \frac{2q(1+6R_H)}{R_Hv}},`$ (14)
$`\beta _1={\displaystyle \frac{6\omega _Hq}{R_H(9R_H^2+\omega _H^2)}},\beta _2=\beta _1,\eta ={\displaystyle \frac{18q}{9R_H^2+\omega _H^2}}.`$ (15)
It is worth pointing out that the phase-amplitude equations (7,2) can be deduced by means of symmetry arguments alone and are, therefore, generic to this order in $`ϵ`$. In fact, they could be derived directly from the fluid equations without the use of the Ginzburg-Landau equations (2). Thus, keeping higher order terms in (2) would change the values of the coefficients, but not their form.
When deriving (7,2) using symmetry arguments, one interesting aspect has to be taken into account. In most secondary bifurcations, the oscillating amplitude is either even or odd under reflection symmetry (see, for instance, ). In our case, however, the field $``$ transforms under reflection $`𝐱𝐱`$ as $`\overline{}`$. The temporal phase of this complex amplitude is therefore a pseudo-scalar, changing sign under reflection. This is because this phase is related to the oscillating frequency which, in turn, depends linearly on the chiral symmetry breaking coefficient ($`\omega _H=2\sqrt{3}R_H^2\gamma `$). This implies that, in Eq. (2), the term $`(\overline{}\overline{})`$ breaks the chiral symmetry, but not the term $`[(𝐞_z\times )\overline{}\overline{}(𝐞_z\times )]`$, as one could have naively expected. Looking at the values for the coefficients we see that, in fact, $`\beta _2`$ changes sign while $`\eta `$ is invariant under $`\omega _H\omega _H`$.
It is interesting to note that at the band-center ($`q=0`$) the system (7,2) decouples. In this case the usual CGLE for the amplitude of the oscillation is recovered, which after rescaling the amplitude, time, and space can be written as :
$$_tH=H+(1+ib_1)^2H(b_3i\mathrm{sign}(\omega _H))H|H|^2,$$
(16)
with
$`b_1={\displaystyle \frac{\xi _i}{\xi _r}}={\displaystyle \frac{2(R_H^2+2\omega _H^2)q^2}{(2q^2R_HR_H^2\omega _H^2)\omega _H}},`$ (17)
$`b_3={\displaystyle \frac{\rho _r}{|\rho _i|}}={\displaystyle \frac{2|\omega _H|R_H(3R_H+1)}{\omega _H^2(1+4R_H)+8R_H^2(1+2R_H)}},`$ (18)
where the sub-indices r and i indicate real and imaginary part, respectively. Furthermore, at the band-center $`b_1=0`$.
The CGLE has been studied extensively. It possesses an extraordinary variety of solutions, including a phase chaotic state, defect chaos, a frozen vortex state and stable plane waves . For the case considered here the values of the parameters $`b_1`$ and $`b_3`$ are always such that the system is in a regime in which stable plane waves coexist with defect chaos. In the present context it has to be emphasized that the complex amplitude of the CGLE represents the amplitude of oscillation of the hexagon pattern. Thus, a solution of Eq. (7) with spatially uniform amplitude $`=He^{i\mathrm{\Omega }t}`$,
$$H=\sqrt{\frac{\epsilon \delta _{1r}}{\rho _r}}=\sqrt{\frac{\epsilon R_H}{4(1+3R_H)}},\mathrm{\Omega }=\delta _{1i}\epsilon \rho _iH^2,$$
(19)
corresponds to a state in which all three amplitudes oscillate in time phase-shifted with respect to each other but the phase of each amplitude is constant in space, as illustrated in Fig. 2. In a traveling wave solution (TW),
$$=He^{i\mathrm{\Omega }t+i𝐤𝐱},$$
(20)
with
$$H=\sqrt{\frac{\epsilon \delta _{1r}\xi _rk^2}{\rho _r}},\mathrm{\Omega }=\epsilon (\delta _{1i}\frac{\rho _i}{\rho _r}\delta _{1r})(\xi _i\frac{\rho _i}{\rho _r}\xi _r)k^2,k=|𝐤|,$$
(21)
the phase of each amplitude is space dependent and in different parts of the system roll-like hexagons with different orientation are dominant at any given time. This is shown in Fig. 3. Note that such a state has two different wavenumbers: that of the underlying regular hexagon pattern and that of the wave modulating the oscillation amplitude. At the band-center, i.e. for hexagons that have the critical wavenumber, the modulation of the oscillation amplitude does not affect the phase of the underlying hexagons.
Away from the band-center ($`q0`$) the complex amplitude $``$ and the phase vector $`\stackrel{}{\varphi }`$ become coupled. For the traveling-wave state of the oscillation amplitude (20) this implies that even the underlying hexagon pattern drifts with a velocity that depends on the wave vector of the modulation. Specifically, the phase is given by $`\stackrel{}{\varphi }=\stackrel{}{\omega }t`$, with
$$\stackrel{}{\omega }=2H^2[\beta _2𝐤\eta (\widehat{𝐞}_z\times 𝐤)],$$
(22)
which implies a drift of the hexagons with a speed $`v=2H^2\sqrt{\beta _2^2+\eta ^2}/\stackrel{~}{q}`$, at an angle $`\theta =\mathrm{arctan}(\eta /\beta _2)`$ with respect to the wave vector of the traveling wave. Substituting the values of the coefficients (9-15) one obtains for the angle $`\theta =\mathrm{arctan}(3R_H/\omega _H)`$. Thus, for $`\omega _H0`$ the angle becomes $`\theta =\pm \pi /2`$, i.e. the drift is perpendicular to the wave vector of the modulation. When $`\omega _H\mathrm{}`$, on the other hand, they drift in the parallel direction ($`\theta \mathrm{\hspace{0.33em}0},\pi `$). The speed $`v`$ is given by
$$v=4\left(\frac{H}{R_H}\right)^2\frac{qk}{\stackrel{~}{q}\sqrt{1+\left(\frac{\omega _H}{3R_H}\right)^2}}.$$
(23)
The traveling waves exist above the curve $`\epsilon \delta _{1r}\xi _rk^2=0`$, up to the global bifurcation at $`\mu =\mu _{het}`$. However they can be unstable to side-band perturbations. In the following we study their stability properties. In particular, we will focus on the effect of the coupling of the oscillation amplitude $``$ to the phase of the underlying hexagons.
## 3 Linear stability analysis
To consider the linear stability of the oscillating hexagons, we perturb the traveling-wave state (20) as
$$=(H+h)e^{i(\mathrm{\Omega }t+𝐤𝐱+\phi )},\stackrel{}{\varphi }=\stackrel{}{\omega }t+\stackrel{}{\varphi }_1.$$
(24)
For simplicity we write $`\stackrel{}{\varphi }_1`$ in the following as $`\stackrel{}{\varphi }`$. Substituting (24) in Eqs. (7,2) we obtain the linear equations for the perturbations:
$`_Th`$ $`=`$ $`2\xi _rH𝐤\phi 2\xi _i𝐤h+\xi _r^2h\xi _iH^2\phi \delta _{2r}H\stackrel{}{\varphi }`$
$`2\rho _rH^2h,`$
$`_T\phi `$ $`=`$ $`2\xi _i𝐤\phi +{\displaystyle \frac{2\xi _r}{H}}𝐤h+\xi _r^2\phi +{\displaystyle \frac{\xi _i}{H}}^2h\delta _{2i}\stackrel{}{\varphi }2\rho _iHh,`$ (26)
$`_T\stackrel{}{\varphi }`$ $`=`$ $`D_{}^2\stackrel{}{\varphi }+D_{}(\stackrel{}{\varphi })+D_{\times _1}^2(𝐞_z\times \stackrel{}{\varphi })+D_{\times _2}(𝐞_z\times )(\stackrel{}{\varphi })`$
$`+2\alpha Hh+2\beta _1H(𝐞_z\times )h2\beta _2H^2\phi +2\eta H^2(𝐞_z\times )\phi `$
$`4\beta _2Hh𝐤+4\eta Hh(\widehat{𝐞}_z\times 𝐤).`$
This leads to a $`4\times 4`$ linear eigenvalue problem, which must be solved numerically. In the long-wave limit the perturbation in the amplitude $`h`$ becomes slaved to the gradients of the phases $`\phi `$ and $`\stackrel{}{\varphi }`$. The resulting $`3\times 3`$ system is, however, still rather involved. A substantial simplification occurs in the case $`𝐤=0`$, i.e. for homogeneously oscillating hexagons. In order to gain physical insight we will consider this case first.
### 3.1 Homogeneous Oscillations
For homogeneously oscillating hexagons $`\stackrel{}{\omega }=0`$ and the perturbation equations (3,26,3) reduce to
$`_Th`$ $`=`$ $`\xi _r^2h\xi H^2\phi \delta _{2r}H\stackrel{}{\varphi }2\rho _rH^2h,`$ (28)
$`_T\phi `$ $`=`$ $`\xi _r^2\phi +{\displaystyle \frac{\xi _i}{H}}^2h\delta _{2i}\stackrel{}{\varphi }2\rho _iHh,`$ (29)
$`_T\stackrel{}{\varphi }`$ $`=`$ $`D_{}^2\stackrel{}{\varphi }+D_{}(\stackrel{}{\varphi })+D_{\times _1}^2(𝐞_z\times \stackrel{}{\varphi })+D_{\times _2}(𝐞_z\times )(\stackrel{}{\varphi })`$ (30)
$`+2\alpha Hh+2\beta _1H(𝐞_z\times )h2\beta _2H^2\phi +2\eta H^2(𝐞_z\times )\phi .`$
In the limit of long-wave perturbations the amplitude $`h`$ can be eliminated adiabatically
$$h\frac{1}{2\rho _rH^2}[\xi _iH^2\phi \delta _{2r}H\stackrel{}{\varphi }].$$
(31)
In this manner we arrive at a system for the three phases, two corresponding to the spatial translations in the plane, the other to a temporal shift,
$`_T\stackrel{}{\varphi }`$ $`=`$ $`D_{}^2\stackrel{}{\varphi }+\left(D_{}{\displaystyle \frac{\alpha \delta _{2r}}{\rho _r}}\right)(\stackrel{}{\varphi })+D_{\times _1}^2(𝐞_z\times \stackrel{}{\varphi })`$ (32)
$`+\left(D_{\times _2}{\displaystyle \frac{\beta _1\delta _r}{\rho _r}}\right)(𝐞_z\times )(\stackrel{}{\varphi }){\displaystyle \frac{\alpha \xi _i}{\rho _r}}^2(\phi )`$
$`{\displaystyle \frac{\beta _1\xi _i}{\rho _r}}^2(𝐞_z\times )\phi 2\beta _2H^2\phi +2\eta H^2(𝐞_z\times )\phi ,`$
$`_T\phi `$ $`=`$ $`\left(\xi _r+{\displaystyle \frac{\rho _i}{\rho _r}}\xi _i\right)^2\phi +\left(\delta _{2r}{\displaystyle \frac{\rho _i}{\rho _r}}\delta _{2i}\right)\stackrel{}{\varphi }.`$ (33)
At the band-center, $`q=0`$, Eqs. (7,2) decouple since $`\beta _1=\beta _2=\delta _2=\alpha =\eta =0`$ (cf. (9-15)). It is easy to show that in that case the eigenvalues are
$`\sigma _{1,2}={\displaystyle \frac{1}{2}}\left[2D_{}+D_{}\pm \sqrt{D_{}^24D_{\times _1}(D_{\times _1}+D_{\times _2})}\right]Q^2,`$ (34)
$`\sigma _3=\left(\xi _r+{\displaystyle \frac{\rho _i}{\rho _r}}\xi _i\right)Q^2.`$ (35)
We obtain therefore the usual expression for the phase instabilities of the underlying hexagons and the Benjamin-Feir instability of the oscillations . The actual values of these eigenvalues, when $`q=0`$, are $`\sigma _1=Q^2/4`$, $`\sigma _2=3Q^2/4`$ and $`\sigma _3=Q^2/2`$. The system is therefore always stable at the band-center.
Away from the band-center ($`q0`$) the system is no longer decoupled. At leading order in the long-wave expansion (32), (33) we have then
$`_T\stackrel{}{\varphi }`$ $`=`$ $`2\beta _2H^2\phi +2\eta H^2(\widehat{e}_z\times )\phi ,`$ (36)
$`_T\phi `$ $`=`$ $`\left(\delta _{2r}{\displaystyle \frac{\rho _i}{\rho _r}}\delta _{2i}\right)\stackrel{}{\varphi }.`$ (37)
These two equations can be combined into a single second-order equation for $`\stackrel{}{\varphi }`$,
$$_T^2\stackrel{}{\varphi }=2\left(\delta _{2r}\frac{\rho _i}{\rho _r}\delta _{2i}\right)H^2(\beta _2\eta (\widehat{e}_z\times ))(\stackrel{}{\varphi }).$$
(38)
Writing in normal modes, $`\varphi _x=\varphi _x^0e^{i𝐐𝐱+\sigma t}`$, $`\varphi _y=\varphi _y^0e^{i𝐐𝐱+\sigma t}`$, $`\phi =\phi ^0e^{i𝐐𝐱+\sigma t}`$, $`Q|𝐐|`$, we arrive at the dispersion relation
$`\sigma _1`$ $`=`$ $`0,`$ (39)
$`\sigma _{2,3}`$ $`=`$ $`\pm \sqrt{2\left(\delta _{2r}\rho _i/\rho _r\delta _{2i}\right)\beta _2}HQ,`$ (40)
indicating the possibility of two different instabilities.
The eigenvalue $`\sigma _1`$ corresponds to the divergence-free part of $`\stackrel{}{\varphi }`$ and does not involve $`\phi `$ (cf. Eq. (37)). It is marginal at this order. The eigenmodes associated with $`\sigma _{2,3}`$ involve both phases $`\stackrel{}{\varphi }`$ and $`\phi `$. For $`\gamma ^2>2/(3R_H^2)`$, $`\sigma _2`$ is always positive (except at the band-center, where $`\beta _2=0`$). Hence there exists a critical value for the rotation $`|\gamma _c|\sqrt{2/3}/R_H`$ above which the system is unstable for $`Q0`$. When $`|\gamma |<|\gamma _c|`$ the eigenvalues $`\sigma _{2,3}`$ are purely imaginary and the stability is determined at next order.
At quadratic order in the perturbation wavenumber $`Q`$ one obtains
$`\sigma _1`$ $`=`$ $`\left(D_{}+{\displaystyle \frac{\eta }{\beta _2}}D_{\times _1}\right)Q^2,`$ (41)
$`\sigma _{2,3}`$ $`=`$ $`\pm \sqrt{2H^2\beta _2\left(\delta _{2r}{\displaystyle \frac{\rho _i}{\rho _r}}\delta _{2i}\right)}Q`$
$`{\displaystyle \frac{1}{2}}\left(D_{}+D_{}\alpha {\displaystyle \frac{\delta _{2r}}{\rho _r}}+\beta _2{\displaystyle \frac{\delta _{2r}\rho _i}{\rho _r^2}}+{\displaystyle \frac{\eta }{\beta _2}}D_{\times _1}+\xi _r+{\displaystyle \frac{\rho _i}{\rho _r}}\xi _i\right)Q^2.`$
Substituting the values of the coefficients from (9-15) in $`\sigma _1`$ we see that it becomes positive provided $`q^2>\gamma ^2R_H^3`$. In Fig. 4a the long-wave stability limits of the oscillating hexagons are shown for $`\nu =2`$ and $`\gamma =0.5`$. These results agree with those obtained solving the full $`4\times 4`$ dispersion relation. Hence, for any given value of the rotation rate there is a value of the wavenumber above which the system becomes unstable. The range of stable wavenumbers decreases as the rotation rate is decreased. In fact, it vanishes as $`\gamma 0`$. In this limit also the range in $`\epsilon `$ over which the oscillating hexagons exist vanishes. Also shown in Fig. 4a are the instabilities of steady hexagons, below the Hopf bifurcation. The solid line represents the long-wave results, while the circles are the results obtained solving the $`6\times 6`$ dispersion relation associated with Eq. (2) . The dash-dotted line represents the line above which the rolls become stable (cf. Eq. (5)).
As $`\gamma `$ is increased the range of wavenumbers that are stable with respect to the diffusive mode ($`\sigma _1<0`$) increases. However, for $`\gamma >\gamma _c`$, $`\sigma _2`$ becomes positive for $`Q0`$. In finite systems the term quadratic in $`Q`$ may not be negligible. In fact, inserting the values from (9-15) into Eq. (3.1) shows that the $`Q^2`$-term in $`\sigma _{2,3}`$ is always stabilizing when $`q1`$ ($`\sigma _2=\sigma _3=5Q^2/8`$, when $`q=0`$), and for $`q𝒪(1)`$ it is only destabilizing when $`\gamma ,\nu 1`$. For $`q=1`$ typical values of $`\gamma ,\nu `$ for which this happens are $`\nu 30,\gamma 70`$. Therefore, in finite systems, in which $`Q_{\mathrm{min}}=2\pi /L`$ cannot be arbitrarily small, there is always a region close to the band-center that is stable, even when $`|\gamma |>|\gamma _c|`$. The stability limit $`\sigma _2=0`$ is given by an expression of the form $`H^2\mu \mu _H=f(q)/(Lq^2)`$, where $`f(q)`$ appears due to the $`q`$-dependence of the second term in (3.1). This situation is shown in Fig. 4b, where the different symbols correspond to several values of the system size. For smaller systems the stable region increases. The shape of the stability limits suggest that $`f(q)`$ does not depend strongly on $`q`$. In Fig. 4b the instability corresponding to $`\sigma _1>0`$ has moved to higher values of $`q`$ ($`q=\pm 0.9`$).<sup>2</sup><sup>2</sup>2The form of the eigenvalues $`\sigma _{2,3}`$ is similar to that encountered in the case of a secondary Hopf bifurcation off a roll pattern . However, in contrast to the one-dimensional case there is a third eigenvalue, $`\sigma _1`$, because the phase has two components. Similar expressions for the eigenvalues are to be expected for other two-dimensional patterns, e.g. squares, undergoing a Hopf bifurcation.
It should be noted that taking the limit $`q0`$ in (41), (3.1) does not give the same results as in (34), (35). In fact, in the limit $`q0`$, the three eigenvalues in (41), (3.1) become $`\sigma _1=Q^2/4`$, $`\sigma _2=\sigma _3=5Q^2/8`$, while Eqs. (34), (35) yield $`\sigma _2=3Q^2/4`$ and $`\sigma _3=Q^2/2`$. This difference is due to the fact that, in order to obtain (41) and (3.1) we assume $`q`$ to be $`𝒪(1)`$, and expand in terms of $`Q`$, while (34) and (35) are obtained by taking the limit $`q0`$ first. If we consider both to be small and of the same order, $`qQ1`$, then:
$`\sigma _1`$ $`=`$ $`D_{}Q^2`$ (43)
$`\sigma _{2,3}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(D_{}+D_{}+\xi _r)Q^2`$
$`\pm \left[2H^2\beta _2\left(\delta _{2r}{\displaystyle \frac{\rho _i}{\rho _r}}\delta _{2i}\right)+{\displaystyle \frac{1}{4}}(D_{}+D_{}\xi _r)^2Q^2\right]^{1/2}Q.`$
When $`q=0`$ we obtain $`\sigma _1=Q^2/4`$, $`\sigma _2=3Q^2/4`$ and $`\sigma _3=Q^2/2`$, while considering $`Qq`$ expressions (39), (40) are recovered at leading order.
### 3.2 Stability of Traveling Waves
For the traveling waves the expressions become quite complicated. We therefore present the analytical results from the long-wave analysis only up to linear order in the gradients. As in the case of homogeneous oscillations, in the long-wave limit $`h`$ becomes slaved (cf. (3)),
$$h\frac{1}{2\rho _rH^2}[2\xi _rH(𝐤)\phi \delta _{2r}H(\stackrel{}{\varphi })].$$
(45)
At the band-center $`\stackrel{}{\varphi }`$ and $`\phi `$ decouple and we recover expression (34) for the eigenvalues $`\sigma _{1,2}`$ associated with $`\stackrel{}{\varphi }`$. The eigenvalue $`\sigma _3`$ describes the long-wave behavior of the single CGLE. It is given by
$$\sigma _3=i𝐯_g𝐐+2(1+\frac{\rho _i^2}{\rho _r^2})\frac{\xi _r^2(𝐤𝐐)^2}{\epsilon \delta _{1r}\xi _rk^2}(\xi _r+\frac{\rho _i}{\rho _r}\xi _i)Q^2,$$
(46)
and yields the usual Eckhaus stability limit for waves. Here we have introduced the group velocity
$$𝐯_g\frac{\mathrm{\Omega }}{𝐤}=2(\xi _i\frac{\rho _i}{\rho _r}\xi _r)𝐤.$$
(47)
The expression $`\epsilon \mu \mu _H=\xi _rk^2/\delta _{1r}`$ gives the neutral surface for the appearance of traveling waves of wavenumber $`𝐤`$.
Away from the band-center Eqs. (3,26,3) reduce at leading order in the long-wave expansion to
$`_T\stackrel{}{\varphi }`$ $`=`$ $`{\displaystyle \frac{4\xi _r}{\rho _r}}[\beta _2𝐤\eta (\widehat{𝐞}_z\times 𝐤)](𝐤)\phi +{\displaystyle \frac{2\delta _{2r}}{\rho _r}}[\beta _2𝐤\eta (\widehat{𝐞}_z\times 𝐤)](\stackrel{}{\varphi })`$
$`2H^2[\beta _2\phi \eta (\widehat{e}_z\times )\phi ],`$
$`_T\phi `$ $`=`$ $`2(\xi _i{\displaystyle \frac{\rho _i}{\rho _r}}\xi _r)𝐤\phi +\left(\delta _{2r}{\displaystyle \frac{\rho _i}{\rho _r}}\delta _{2i}\right)\stackrel{}{\varphi }.`$ (49)
This yields the eigenvalues
$`\sigma _1`$ $`=`$ $`0`$ (50)
$`\sigma _{2,3}`$ $`=`$ $`{\displaystyle \frac{i}{2}}(𝐯_g+{\displaystyle \frac{\delta _{2r}}{\rho _r}}𝐯_h)𝐐\pm \{{\displaystyle \frac{1}{4}}[(𝐯_g+{\displaystyle \frac{\delta _{2r}}{\rho _r}}𝐯_h)𝐐]^2`$
$`+{\displaystyle \frac{4}{\rho _r}}(\delta _{2r}𝐯_g𝐐\xi _r\mathrm{\Gamma }𝐤𝐐)(𝐯_g𝐐)+2H^2\beta _2\mathrm{\Gamma }Q^2\}^{1/2},`$
where we have defined $`\mathrm{\Gamma }\delta _{2r}\rho _i/\rho _r\delta _{2i}`$ and $`𝐯_h\stackrel{}{\omega }/(2H^2)`$ (cf. (22)).
Since the analytical expressions for the eigenvalues are quite complicated we solve Eqs. (3,26,3) directly without the additional long-wave approximation (45) for various values of the wavenumber $`𝐤=(k,0)`$ of the traveling waves. Fig. 5a shows the neutral surface of traveling waves and their stability limit as a function of the hexagon wavenumber $`q`$. Due to the reflection symmetries $`qq`$ and $`kk`$ only one quadrant is shown. For clarity the stability surface has been capped at $`\mu =6`$. For $`k=0`$ the stability limit does not depend on $`\mu `$ and the stability surface is vertical (cf. Fig. 4a). As $`k`$ is increased the stability surface becomes smoother. For $`k0`$, the waves are unstable at onset and become stable above the second surface. Fig. 5b shows cross-sections of the stability surface for $`\gamma >\gamma _c`$, for a system of size $`L=250`$. When $`k=0`$ we recover the results from Fig. 4b. For $`k0`$, but small (cf. $`k=0.1`$ in Fig. 5b) the traveling waves are unstable at onset but they become stable for larger values of $`\mu `$. When $`|q|`$ is large they can become unstable again as $`\mu `$ is increased. For larger values of $`k`$, the latter two stability lines merge and the stability region becomes bounded in $`q`$. At this point the stability limits in Figs. 5a and 5b look similar.
## 4 Numerical simulations
In order to study the nonlinear behavior arising from the instabilities, we have performed numerical simulations of Eqs. (2) and (7), (2). A Runge-Kutta method with an integrating factor that computes the linear derivative terms exactly has been used. Derivatives were computed in Fourier space, using a two-dimensional fast Fourier transform (FFT). The numerical simulations were done in a rectangular box of aspect ratio $`2/\sqrt{3}`$ with periodic boundary conditions. This aspect ratio was used to allow for regular hexagonal patterns.
We investigate the stability of oscillating hexagons simulating both the original amplitude equations (2) and the reduced amplitude-phase equations (7), (2). To that end, we start with homogeneously oscillating hexagons as given by Eq. (19) (or Eq. (6)) and add weak noise. For values of $`\mu `$ close to $`\mu _H`$, the growth rates obtained from (2) and (7), (2) agree with each other and with the results from the linear stability analysis. Both the long-wave instabilities coming from $`\sigma _1>0`$ (41) and $`\sigma _{2,3}>0`$ (3.1) lead to qualitatively similar behavior. The perturbations grow until they reach a saturation suggesting that the bifurcations are supercritical. Since the perturbations involve spatial modulations of the oscillation amplitude the hexagons begin to travel. However, the modulation wavevector $`𝐤`$ varies in space and induces drift velocities that are different in magnitude and direction at different locations in the system, implying a shear of the pattern. This results in deformed hexagon patterns as shown in Fig. 6.
In addition to the long-wave instabilities a short-wave instability appears for larger values of the hexagon wavenumber $`q`$. This instability is induced by the short-wave instability of the steady hexagons and cannot be studied with the amplitude-phase equations (7), (2). As $`\mu `$ is increased above $`\mu _H`$, the oscillating hexagons arise through a Hopf bifurcation off the steady hexagons. Since that bifurcation occurs at zero wavenumber it affects the long-wave properties of the system and we expect that the long-wave stability limits of the steady and of the oscillating hexagons differ qualitatively. The effect of the Hopf bifurcation on short-wave instabilities, on the other hand, should vanish as the amplitude of the oscillations goes to 0. Thus, we expect a continuous transition from the short-wave instability of the steady to that of the oscillating hexagons as the line $`\mu =\mu _H`$ is crossed. In order to check this, we numerically determine the short-wave instability of the oscillating hexagons and compare it with the stability results for the steady hexagons, as obtained from solving the sixth-order dispersion relation associated with Eq. (2) . The transition is indeed continuous, as expected. The stability region of the oscillating hexagons turns out to be reduced as compared to that of the steady hexagons when $`\mu `$ is increased (see Fig. 7).
If the noise added to the oscillating hexagons is sufficiently large, different dynamics may arise even in the parameter range in which the oscillating hexagons are linearly stable. As indicated earlier, within the framework of the three coupled Ginzburg-Landau equations (2) one obtains for the parameters in the CGLE (16) values for which a persistent chaotic state exists while the plane waves are linearly stable. A detailed study of the chaotic state, which is characterized by the creation and annihilation of defects, comparing its description using the coupled Ginzburg-Landau equations (2), the amplitude-phase equations (7,2), and the single CGLE (16) has been presented elsewhere . One of the main results of that study is the observation that this system is one of the few in which the defect chaos of the CGLE should be accessible experimentally.
## 5 Conclusions
In this paper we have studied the effect of a breaking of the chiral symmetry on systems that exhibit hexagon patterns. Classic examples of such systems are non-Boussinesq and surface-tension-driven convection with rotation. We have focused on the dynamics of the oscillating hexagons that arise in a Hopf bifurcation due to the rotation. In the vicinity of this secondary Hopf bifurcation the oscillating hexagons are described by a single complex Ginzburg-Landau equation (CGLE) coupled to the two phases of the underlying hexagons. The resulting amplitude-phase equations have certain similarities with those describing the secondary Hopf bifurcation observed in rectangle patterns in electro-convection in nematics .
Like the CGLE, the amplitude-phase equations support homogeneously oscillating solutions as well as traveling waves. In the latter, the coupling to the phases induces a drift of the hexagons in a direction that is typically oblique to the propagation direction of the traveling waves. The stability analysis of the oscillating and the traveling hexagons reveal two types of long-wave instabilities, one occurring when the hexagons are far away from the band-center, the other for high enough rotation rate. Even in this latter case, in finite systems there is always a stable region close to the band-center, its size depending on the size of the system. In both cases, the instabilities appear to be supercritical, giving rise to a spatially modulated oscillating hexagonal pattern.
Although there is always a region in which the homogeneously oscillating hexagons are linearly stable within the three coupled Ginzburg-Landau equations, they are in fact only meta-stable. Sufficiently large perturbations induce a transition to a state of defect chaos described by the CGLE . There is always bistability between the chaotic state and the ordered oscillations. In that respect the chaotic state resembles spiral-defect chaos as it is observed in Rayleigh-Bénard convection at low Prandtl numbers . As in that system the ordered states can presumably only be obtained by carefully controlled initial conditions. Rotating hexagon convection appears to be the first system in which the defect chaotic regime of the CGLE should be accessible experimentally.
From numerical simulations of the amplitude equations (2) it has been shown that a transition between defect-chaos and a frozen vortex state occurs for wavenumbers far away from the band-center . This transition happens when the asymptotic plane waves emitted by the spirals become absolutely unstable. Therefore, the stability results for the traveling wave solution are relevant in order to determine the range of existence of the chaotic state found in . A complete treatment of this point would involve the determination of the asymptotic wavenumber $`k_{\mathrm{}}`$ for the waves emitted by the spiral solutions of the coupled amplitude-phase equations.
In the present work, we have studied the system close to the bifurcation point, where analytical results can be obtained. However, far away from the Hopf bifurcation, a number of interesting effects are expected. One is related to the strength of the chiral symmetry breaking. To leading order in the perturbation expansion, the usual single CGLE is obtained, which is chirally symmetric. The breaking of the chiral symmetry appears only through the coupling with the phase, and vanishes at the band-center. For larger values of the oscillation amplitude, higher-order terms breaking the chiral symmetry of the CGLE are expected. This leads to an asymmetry between defects with opposite topological charge , with one type of spirals becoming dominant in the frozen vortex state. The asymptotic wavenumber, as well as the onset of absolute instability, becomes different for positively and negatively charged spirals, and the transition to the defect chaotic regime is different than in the chirally symmetric case . We expect that this effect of the chiral symmetry breaking can be observed with Eqs. (2), as the control parameter is increased above the Hopf bifurcation.
Another open question is the relation of the limit cycle corresponding to the oscillating hexagons with the heteroclinic orbit arising in the Küppers-Lortz instability of rolls. As the limit cycle approaches the mixed-mode solutions, the harmonic oscillations are transformed into a switching between the three roll modes making up the hexagons similar to the dynamics arising in the Küppers-Lortz regime. This suggest a connection between the domain chaos found in these systems and the regular and disordered states discussed in this paper. However, it is worth emphasizing that in the defect chaos regime described in , the orientation of the hexagons is well defined and what is spatially chaotic is the modulation of the amplitudes that compose the hexagon pattern. In the Küppers-Lortz regime, on the other hand, patches of rolls with arbitrary orientation are possible, resulting in a state with an isotropic Fourier spectrum. Thus, a quantitative comparison between both states is not possible with Ginzburg-Landau equations such as (2) and generalized Swift-Hohenberg models must be considered .
Acknowledgments
We gratefully acknowledge interesting discussions with F. Sain and M. Silber. The numerical simulations were performed with a modification of a code by G.D. Granzow. This work was supported by D.O.E. Grant DE-FG02-G2ER14303 and NASA Grant NAG3-2113.
## Appendix A Derivation of the Amplitude-Phase Equations
At the Hopf bifurcation, the hexagon solution (2) becomes $`R_H=1/(\nu 1)`$. We will consider perturbations around this solution, both in amplitude and phase:
$$A_i=(R_H+r_i)e^{i𝐪_i𝐱+\varphi _i+\mathrm{\Phi }},$$
(52)
where $`\mathrm{\Phi }=\varphi _1+\varphi _2+\varphi _3`$ is the global phase of the hexagons and the three phases $`\varphi _i`$ can be written as
$`\varphi _1`$ $`=`$ $`\varphi _x,`$
$`\varphi _2`$ $`=`$ $`\varphi _x/2+\sqrt{3}\varphi _y/2,`$
$`\varphi _3`$ $`=`$ $`\varphi _x/2\sqrt{3}\varphi _y/2,`$
where we have defined the phase vector $`\stackrel{}{\varphi }=(\varphi _x,\varphi _y)`$, with $`\varphi _x`$ and $`\varphi _y`$ related to translations of the pattern in the x- and y-directions.
For the perturbations of the modulus $`r_i`$ there are three eigenvalues, one real, $`\sigma _1=2R^2(1+2\nu )+R`$, corresponding to an eigenvector with $`r_1=r_2=r_3`$ and a complex conjugate pair, $`\sigma _{2,3}=2R^2(1\nu )2R\pm 2\sqrt{3}R^2\gamma i`$, whose real part vanishes at $`\mu =\mu _H`$. The corresponding eigenvector satisfies $`r_3=e^{2\pi i/3}r_2=e^{4\pi i/3}r_1`$.
Taking this into account, we consider the expansion:
$`r_1`$ $`=`$ $`ϵr+[({\displaystyle \frac{1}{2}}+i{\displaystyle \frac{\sqrt{3}}{2}})(\sqrt{ϵ}e^{i\omega t}+ϵ(_{10}+\overline{}_{12}e^{2i\omega t}))+c.c.],`$
$`r_2`$ $`=`$ $`ϵr+[({\displaystyle \frac{1}{2}}i{\displaystyle \frac{\sqrt{3}}{2}})(\sqrt{ϵ}e^{i\omega t}+ϵ(_{10}+\overline{}_{12}e^{2i\omega t}))+c.c.],`$
$`r_3`$ $`=`$ $`ϵr+[(\sqrt{ϵ}e^{i\omega t}+ϵ(_{10}+\overline{}_{12}e^{2i\omega t}))+c.c.],`$
$`\varphi _x`$ $`=`$ $`\sqrt{ϵ}\varphi _x+ϵ(\varphi _{x1}e^{i\omega t}+c.c.),`$
$`\varphi _y`$ $`=`$ $`\sqrt{ϵ}\varphi _y+ϵ(\varphi _{y1}e^{i\omega t}+c.c.),`$
$`\mathrm{\Phi }`$ $`=`$ $`ϵ(\mathrm{\Phi }_1e^{i\omega t}+c.c.).`$
We also assume that the resulting state evolves on long time and space scales, specifically: $`_t𝒪(ϵ)`$, $`𝒪(ϵ^{1/2})`$. Substituting the former expressions into Eq. (2), at $`𝒪(ϵ^{1/2})`$ the linear problem is recovered, giving the value for the critical frequency: $`\omega _H=2\sqrt{3}R_H^2\gamma `$.
At $`𝒪(ϵ)`$ an algebraic relation between the slaved and the marginal modes is obtained:
$`\varphi _{x1}`$ $`=`$ $`{\displaystyle \frac{q}{2R_H}}(\sqrt{3}+i)(_xi_y),`$ (53)
$`\varphi _{y1}`$ $`=`$ $`{\displaystyle \frac{qi}{2R_H}}(\sqrt{3}+i)(_xi_y),`$ (54)
$`\mathrm{\Phi }_1`$ $`=`$ $`{\displaystyle \frac{q}{2R_H(3R_H+i\omega _H)}}(\sqrt{3}i1)(_x+i_y),`$ (55)
$`r`$ $`=`$ $`{\displaystyle \frac{1}{3R_H(1+2R_H)}}\left[\mu R_H(1+6R_H)||^2qR_H\stackrel{}{\varphi }\right],`$ (56)
$`_{10}`$ $`=`$ $`{\displaystyle \frac{qR_H}{4\omega _H}}(\sqrt{3}i+1)\left[_x\varphi _y+_y\varphi _x+i(_x\varphi _x_y\varphi _y)\right],`$ (57)
$`_{12}`$ $`=`$ $`{\displaystyle \frac{1}{6}}\left(1+{\displaystyle \frac{8i}{\omega _H}}\right)^2.`$ (58)
At $`𝒪(ϵ^{3/2})`$ we obtain a solvability condition for $``$ and $`\stackrel{}{\varphi }`$:
$`_t`$ $`=`$ $`\mu +{\displaystyle \frac{1}{2}}^2q\stackrel{}{\varphi }\left(1+6R_H{\displaystyle \frac{2\omega _Hi}{R_H}}\right)r`$ (59)
$`{\displaystyle \frac{1}{R_H^2}}\left[3R_H(1+2R_H){\displaystyle \frac{\omega i}{2}}\right]||^2+\left(8{\displaystyle \frac{\omega i}{R_H}}\right)\overline{}_{12}`$
$`+{\displaystyle \frac{R_Hq}{4}}(1+\sqrt{3}i)[_x\varphi _{x1}_y\varphi _{y1}+i(_x\varphi _{y1}+_y\varphi _{x1})]`$
$`+{\displaystyle \frac{R_Hq}{2}}(\sqrt{3}i+1)(_x\mathrm{\Phi }_1i_y\mathrm{\Phi }_1),`$
$`_t\varphi _x`$ $`=`$ $`{\displaystyle \frac{1}{4}}^2\varphi _x+{\displaystyle \frac{1}{2}}_x(\stackrel{}{\varphi })+{\displaystyle \frac{2q}{R_H}}_xr`$ (60)
$`+{\displaystyle \frac{q}{2R_H}}[(\sqrt{3}i1)(_x_{10}i_y_{10})+c.c.]`$
$`+{\displaystyle \frac{\omega }{4R_H}}[(i\sqrt{3})\overline{}(\varphi _{x1}+i\varphi _{y1})+c.c.]`$
$`{\displaystyle \frac{1}{2R_H}}[(1+\sqrt{3}i)(3R_H\omega i)\overline{}\mathrm{\Phi }_1+c.c.],`$
$`_t\varphi _y`$ $`=`$ $`{\displaystyle \frac{1}{4}}^2\varphi _y+{\displaystyle \frac{1}{2}}_y(\stackrel{}{\varphi })+{\displaystyle \frac{2q}{R_H}}_yr`$ (61)
$`+{\displaystyle \frac{q}{2R_H}}[(\sqrt{3}+i)(_x_{10}i_y_{10})+c.c.]`$
$`{\displaystyle \frac{\omega }{4R_H}}[(1+\sqrt{3}i)\overline{}(\varphi _{x1}+i\varphi _{y1})+c.c.]`$
$`{\displaystyle \frac{1}{2R_H}}[(\sqrt{3}i)(3R_H\omega i)\overline{}\mathrm{\Phi }_1+c.c.].`$
Substituting the expressions of the slaved modes into the former equations we obtain the amplitude-phase equations (7), (2).
## Appendix B Nonlinear Gradient Terms
If we retain in the Ginzburg-Landau equations (2) the nonlinear gradient terms that express the dependence of the quadratic coupling term on the hexagon wavenumber we obtain
$`_tA_1`$ $`=`$ $`\mu A_1+(𝐧_1)^2A_1+\overline{A}_2\overline{A}_3A_1|A_1|^2`$ (62)
$`(\nu +\gamma )A_1|A_2|^2(\nu \gamma )A_1|A_3|^2`$
$`+i(\alpha _1+\stackrel{~}{\alpha })\overline{A}_2(𝐧_3)\overline{A}_3+i(\alpha _1\stackrel{~}{\alpha })\overline{A}_3(𝐧_2)\overline{A}_2`$
$`+i\alpha _2\left(\overline{A}_2(𝝉_3)\overline{A}_3\overline{A}_3(𝝉_2)\overline{A}_2\right).`$
For the amplitude-phase equations (7), (2) we obtain then the coefficients:
$`\alpha =1+2q\alpha _1,`$
$`R_H={\displaystyle \frac{\alpha }{\nu 1}},`$
$`ϵ_c={\displaystyle \frac{\alpha ^2(2+\nu )}{(\nu 1)^2}}+q^2=R_H(3R_H+\alpha )+q^2,`$
$`v=3R_H(\alpha +2R_H),`$
$`\omega _H=2\sqrt{3}\gamma R_H^2,`$
$`\delta _1={\displaystyle \frac{2\alpha R_H}{v}}{\displaystyle \frac{2i\omega _H}{v}},`$
$`\delta _2={\displaystyle \frac{2\alpha R_Hq}{v}}+{\displaystyle \frac{4R_H^2\alpha _1(3R_H+\alpha )}{v}}{\displaystyle \frac{2i\omega _H(qR_H\alpha _1)}{v}},`$
$`\xi ={\displaystyle \frac{1}{2}}{\displaystyle \frac{R_H(R_H\alpha _1+q)}{9R_H^2\alpha ^2+\omega _H^2}}[{\displaystyle \frac{\sqrt{3}}{2}}\stackrel{~}{\alpha }\omega _H+3\alpha (q+{\displaystyle \frac{R_H}{2}}(\alpha _1\sqrt{3}\alpha _2))]{\displaystyle \frac{\sqrt{3}R_Hq\stackrel{~}{\alpha }}{\omega _H}}`$
$`{\displaystyle \frac{i(R_H\alpha _1+q)}{9R_H^2\alpha ^2+\omega _H^2}}\left[\omega _Hq+{\displaystyle \frac{R_H}{2}}[\omega _H(\alpha _1\sqrt{3}\alpha _2)3\sqrt{3}R_H\alpha \stackrel{~}{\alpha }]\right]{\displaystyle \frac{iq^2}{\omega _H}}`$
$`+{\displaystyle \frac{iR_H^2}{4\omega _H}}\left[(\alpha _1+\sqrt{3}\alpha _2)^2+3\stackrel{~}{\alpha }^2\right]`$
$`\rho ={\displaystyle \frac{8\alpha (3R_H+\alpha )}{v}}{\displaystyle \frac{4i\omega _H(\alpha +4R_H)}{R_Hv}}{\displaystyle \frac{32i\alpha ^2}{3\omega _H}},`$
$`D_{}={\displaystyle \frac{1}{4}}{\displaystyle \frac{\sqrt{3}R_Hq\stackrel{~}{\alpha }}{\omega _H}},`$
$`D_{}={\displaystyle \frac{1}{2}}{\displaystyle \frac{R_H\alpha _1q}{v}}\left[R_H(\alpha _1\sqrt{3}\alpha _2)2q\right],`$
$`D_{\times _1}={\displaystyle \frac{q^2}{\omega _H}}{\displaystyle \frac{R_H^2}{4\omega _H}}\left[(\alpha _1+\sqrt{3}\alpha _2)^2+3\stackrel{~}{\alpha }^2\right],`$
$`D_{\times _2}={\displaystyle \frac{\sqrt{3}R_H\stackrel{~}{\alpha }}{v}}(R_H\alpha _1q),`$
$`\alpha ={\displaystyle \frac{\alpha +6R_H}{R_Hv}}[2q+R_H(\sqrt{3}\alpha _2\alpha _1)]+{\displaystyle \frac{18\alpha ^2R_H^3\alpha _12\omega _H^2q}{R_H^2(9R_H^2\alpha ^2+\omega _H^2)}},`$
$`\beta _1={\displaystyle \frac{6\omega _H\alpha (q+R_H\alpha _1)}{R_H(9R_H^2\alpha ^2+\omega _H^2)}}+{\displaystyle \frac{\sqrt{3}\stackrel{~}{\alpha }}{v}}(\alpha +6R_H),`$
$`\beta _2={\displaystyle \frac{6\omega _H\alpha (q+R_H\alpha _1)}{R_H(9R_H^2\alpha ^2+\omega _H^2)}}+{\displaystyle \frac{\sqrt{3}\stackrel{~}{\alpha }}{R_H}},`$
$`\eta ={\displaystyle \frac{[9R_H^2\alpha ^2(2qR_H(\alpha _1+\sqrt{3}\alpha _2))R_H\omega _H^2(3\alpha _1+\sqrt{3}\alpha _2)]}{R_H^2(9R_H^2\alpha ^2+\omega _H^2)}}.`$
### B.1 Comparison with the CGLE
After rescaling we obtain the values for the coefficients of the CGLE (16). The coefficient $`b_3`$ now becomes:
$$b_3=\frac{2|\omega _H|R_H\alpha (\alpha +3R_H)}{\omega _H^2(\alpha +4R_H)+8R_H^2\alpha ^2(\alpha +2R_H)}.$$
(63)
For $`q0`$ the maximum value of $`b_3`$ is:
$$b_3^{max}=\frac{\alpha +3R_H}{\sqrt{8(\alpha +2R_H)(\alpha +4R_H)}},$$
(64)
and the limits for small and large $`R_H`$ are the same as in (18), even if now $`b_3`$ depends on $`q`$ (through $`\alpha =1+2q\alpha _1`$). At the band-center $`b_3`$ becomes the same as (18).
The expression for $`b_1`$ is quite involved. At the band-center it reduces to:
$$b_1=\frac{(9R_H^2+\omega _H^2)[(\alpha _1+\sqrt{3}\alpha _2)^2+3\stackrel{~}{\alpha }^2]+2\omega _H\alpha _1[\omega _H(\sqrt{3}\alpha _2\alpha _1)+3\sqrt{3}R_H\stackrel{~}{\alpha }]}{2\omega _H(9R_H^2+\omega _H^2\alpha _1R_H^2[3R_H(\alpha _1\sqrt{3}\alpha _2)+\sqrt{3}\stackrel{~}{\alpha }\omega _H])}.$$
(65)
An important change occurs with respect to the decoupling of the amplitude-phase equation. For $`\alpha _10`$ the amplitude and the phases do not decouple anymore. It would require $`\delta _{2r}=0`$, $`\delta _{2i}=0`$. From the latter we obtain $`q=R_H\alpha _1`$, while the former implies
$$q=\frac{1}{4\alpha _1}\left[(1+4R_H\alpha _1^2)\pm \sqrt{(14R_H\alpha _1)^248R_H^2\alpha _1^2}\right].$$
(66)
The only real solution for both conditions is $`\alpha _1=0`$.
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# Non-equilibrium H2 ortho-to-para ratio in two molecular clouds of the Galactic CenterISO is an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) and with the participation of ISAS and NASA.
## 1 Introduction
The central $`6^{}`$ of our galaxy exhibit a large accumulation of molecular material which is forming big molecular clouds whose masses and sizes are so large as 10<sup>6</sup> M and 15 pc, respectively. These clouds are denser (average densities of 10<sup>4</sup> cm<sup>-3</sup> ), more turbulent (line widths of $``$ 20 km s<sup>-1</sup> ), and hotter (with a warm component with temperatures, $`T`$, up to 200-300 K) than the clouds of the disk of the galaxy (see e.g. Morris & Serabyn morris (1996)). The high temperatures in Galactic Center (GC) clouds were known basically by observations of NH<sub>3</sub> inversion lines over limited regions (Güsten et al. gusten81 (1981); Mauersberger et al. mauersberger (1986)). Hüttemeister et al. (huttemeister93 (1993)) analyzed 36 molecular clouds distributed all along the Central Molecular Zone and the “Clump 2” complex; they showed that high kinetic temperatures are a general characteristic of the GC clouds and not only of those located close to Sgr A and Sgr B2. In the disk of the galaxy, kinetic temperatures higher than 100 K are associated with infrared sources, that is, embedded stars which heat the dust and subsequently the gas by collisions with the dust grains. The typical sizes of such regions are less than 1 pc. The high kinetic temperatures in the GC clouds are found in regions of $``$ 10 pc, where one measures large column densities of cold dust ($`T<`$ 30 K, Odenwald & Fazio odenwald (1984); Cox & Laureijs cox (1989)). This rules out gas-dust collisions as a possible heating mechanism of the warm component. Dissipation of turbulence due to shocks induced by the rotation of the galaxy could be the main heating mechanism in the GC clouds (Wilson et al. wilson82 (1982)).
Unfortunately the NH<sub>3</sub> abundance in the warm component was unknown since one could not estimate the warm H<sub>2</sub> column densities. The Infrared Space Observatory (ISO; Kessler et al. kessler (1996)), has allowed us, for the first time, to measure directly the total column density of warm gas by observing pure-rotational lines of H<sub>2</sub>. These trace gas with temperatures of a few hundreds Kelvin. ISO has also allowed us to study the H<sub>2</sub> ortho-to-para ratio (OTPR), which can help determine the possible heating mechanism and the origin of this molecule. Before ISO, the H<sub>2</sub> OTPR had been studied in regions with temperatures of $``$ 2000 K, using the vibrational lines. In such shock-excited sources, one measures an OTPR of $``$ 3 (Smith et al. smith (1997)), which is the local thermodynamical equilibrium (LTE) value for $`T\stackrel{>}{}`$ 200 K. In contrast, for regions heated mainly by ultraviolet (UV) radiation (Photodissociation regions \[PDRs\]), the vibrational lines give OTPRs in the range of 1.2-2 (see e.g. Chrysostomou et al. chrysostomou (1993)). However, these low OTPRs might not be a consequence of an actual non-LTE ortho-to-para abundances ratio but a result of optical depth effects in the fluorescence-pumping of the ortho-H<sub>2</sub> (Sternberg & Neufeld sternberg (1999)). Using these considerations, one can explain why the PDR in S140 exhibits an OTPR $``$ 2 in the vibrational states but 3 in the lowest rotational levels.
There are two cases of non-equilibrium OTPR measured from the pure-rotational lines: the shock excited source HH54 (Neufeld et al. neufeld (1998)) and the PDR associated with the reflection nebula NGC 7023 (Fuente et al. fuente99 (1999)). The first case has been explained using the shocks model of Timmermann (timmermann (1998)), which involves transient heating by low velocity shocks. To explain the non-LTE OTPR in NGC 7023, it was necessary to invoke a dynamic dissociation front.
To investigate the thermal balance of the GC clouds we have selected 18 clouds from the samples of Hüttemeister et al. (huttemeister93 (1993)) and Martín-Pintado et al. (mp97 (1997)) and we have observed them with the ISO satellite. In this paper we present H<sub>2</sub> observations toward two sources which show similar characteristics (also shown in the NH<sub>3</sub> studies of Hüttemeister et al. huttemeister93 (1993)), indicating that their heating mechanisms are also very similar. In particular, they show a non-LTE OTPR. The detection of OTPRs out of equilibrium in the GC clouds gives us new insights into the heating mechanism, since the gas must be heated to several hundreds K almost without changing the OTPR of cold gas.
In Sects. 2 and 3, we present observations and results, respectively, and in Sect. 4 we discuss the possible heating mechanism and the origin of the non-equilibrium OTPR.
## 2 Observations and data reduction
### 2.1 ISO observations
We observed the pure-rotational H<sub>2</sub> lines S(0), S(1), S(2), and S(3) with the Short Wavelength Spectrometer (SWS; de Graauw et al. graauw (1996)) on board ISO toward two molecular clouds. The galactic coordinates and ISO beam sizes are given in Table 1. These sources are among the farthest from the dynamical center of the galaxy in our sample. M+3.06+0.34 \[$`\alpha `$(2000)= 17<sup>h</sup> 51<sup>m</sup> 26$`\stackrel{s}{.}`$4, $`\delta `$(2000) = -26°08′29$`\stackrel{}{.}`$4\] is located in the “Clump 2” complex (Stark & Bania stark (1986)), while M+1.56$``$0.30 \[$`\alpha `$(2000)= 17<sup>h</sup> 50<sup>m</sup> 26$`\stackrel{s}{.}`$6, $`\delta `$(2000) = -27°45′29$`\stackrel{}{.}`$5\] belongs to the “$`l`$=1$`\stackrel{}{.}`$5-complex” (Bally et al. bally (1988)). The observations were made during orbits 313 (S(0) and S(3) lines), and 467 (S(1) and S(2) lines). The wavelength bands were scanned in the SWS02 mode with a typical on-target time of 100 s. The spectral resolution ($`\lambda /\mathrm{\Delta }\lambda `$) of this mode is $``$ 1000-2000 corresponding to a velocity resolution of $``$ 150-300 km s<sup>-1</sup> . All the lines have broader profiles than those expected for a point source by a factor 1.3-1.5, indicating that the sources are extended in the direction perpendicular to the slit (see Valentijn & Van der Werf valentijn99 (1999)). The flux calibration is believed to be accurate to 30$`\%`$, 20$`\%`$, 25$`\%`$, and 25$`\%`$ for the S(0), S(1), S(2), and S(3) lines, respectively (Salama et al. 1997). Data reduction was carried out with version 6 of the SWS Interactive Analysis at the ISO Spectrometer Data Center at MPE. Further analysis has been made using the ISAP <sup>1</sup><sup>1</sup>1The ISO Spectral Analysis Package (ISAP) is a joint development by the LWS and SWS Instrument Teams and Data Centers. Contributing institutes are CESR, IAS, IPAC, MPE, RAL and SRON. software package. All lines have been rebinned to one fifth of the spectral resolution of the instrument. Fig. 1a-b shows the spectra, and the observed parameters are given in Table 1. The errors in the radial velocities of the H<sub>2</sub> lines listed in this table have been estimated from the Gaussian fits. The wavelength calibration uncertainties, expressed in velocities, are typically of 20-40 km s<sup>-1</sup> for $`\lambda >12\mu `$m and $``$ 25-60 km s<sup>-1</sup> for $`\lambda <12\mu `$m (Valentijn et al. valentijn96 (1996)). Thus, the calibration uncertainties usually dominate the global error in the radial velocities. When one takes into account the errors from the Gaussian fits and the wavelength calibration uncertainties, the central velocities of the H<sub>2</sub> lines are in agreement with those measured from the CO lines (section 2.2). It is noteworthy that, the higher the signal-to-noise ratio of the H<sub>2</sub> lines (S(1) lines), the better the agreement of the H<sub>2</sub> radial velocities with those of CO.
We also present Long Wavelength Spectrometer (LWS; Clegg et al. clegg (1996); Swinyard et al. swinyard (1996)) observations of these sources in grating mode (43-196.7 $`\mu `$m, $`\lambda `$/$`\mathrm{\Delta }\lambda `$ $``$ 200). Fig. 1 e-f shows the LWS spectra. The spectral resolution was 0.29 $`\mu `$m for the 43-93 $`\mu `$m range and 0.6 $`\mu `$m for the 80-196 $`\mu `$m range. The LWS aperture was $`80^\mathrm{"}\times 80^\mathrm{"}`$. The roll angle, which gives the orientation of the apertures, was 90°$`\pm `$ 2°for both the SWS and the LWS observations. Data were taken during orbits 315 and 318 and processed through the LWS Pipeline Version 7. The individual detector scans were calibrated to within 10$`\%`$ of each other, based on overlapping detectors. Post-pipeline analysis (including shifting the different detectors using dark currents and defringing) was performed with ISAP.
### 2.2 IRAM 30-m observations
The J=1–0 line of <sup>13</sup>CO and C<sup>18</sup>O and the J=2–1 line of C<sup>18</sup>O were observed simultaneously with the IRAM 30-m telescope (Pico Veleta, Spain) in May 1997. We used two SIS receivers at 3 and 1.3 mm connected to two $`512\times 1`$ MHz channel filter banks. This configuration provided a velocity resolution of 2.7 and 1.4 km s<sup>-1</sup> for the J=1–0 and J=2–1 lines respectively. Typical system temperatures were $``$ 250 K for the J=1–0 line and $``$ 500 K for the J=2–1 line. The receivers were tuned to single side band with rejections always larger than 10 dB that were checked against standard calibration sources. The beam size of the 30-m telescope was 22<sup>"</sup> and 11<sup>"</sup> at 3 and 1.3 mm respectively. Pointing and focus were monitored regularly. Pointing corrections were always found to be smaller than 3<sup>"</sup>. Calibration of the data was made by observing a hot and cold loads with known temperatures, and the line intensities were converted to main beam brightness temperature, $`T_{\mathrm{MB}}`$, using main beam efficiencies of 0.74 and 0.48 at 3 and 1.3 mm respectively. The spectra are shown in Fig 1c-d and the observed parameters as derived from Gaussian fits are listed in Table 2.
## 3 Analysis
In Fig. 2 we show the H<sub>2</sub> rotational diagrams for the two sources. The open squares correspond to the column densities as measured with ISO, without any correction for the different apertures in the different lines and for the dust extinction. The rotational diagrams for the two sources show a zig-zag distribution since the column densities in the ortho-H<sub>2</sub> levels J=3 and 5 are lower than those expected from the para-H<sub>2</sub> levels for the LTE OTPR. For the typical temperatures involved in these transitions ($``$ 200 K), the LTE OTPR is $``$ 3. Similar rotational diagrams derived from the H<sub>2</sub> pure-rotational lines have been previously found in HH54 by Neufeld et al. (neufeld (1998)) and in NGC7023 by Fuente et al. (fuente99 (1999)). For these sources where extinction is known to be low, the immediate conclusion was that the OTPR was not in LTE.
The H<sub>2</sub> emission has been detected in all sources of our sample indicating that the H<sub>2</sub> emission in the GC must be relatively widespread and extended (Martín-Pintado et al. 1999b ). This is also suggested from the measured linewidths of the H<sub>2</sub> lines (see Section 2). Anyhow, even in the extreme case that the H<sub>2</sub> emission were point-like, the corrections for the different apertures would be small and would not affect substantially the conclusions about the OTPR. For a point-like source the S(0) line will be more diluted than the S(1) and S(2) lines because of the larger beam ($`20^\mathrm{"}\times 27^\mathrm{"}`$ instead of $`14^\mathrm{"}\times 27^\mathrm{"}`$). The situation for the S(3) line will be the opposite since the aperture at this wavelength is $`14^\mathrm{"}\times 20^\mathrm{"}`$. Therefore, in this limit case, the column densities in the level J=2 (derived from the S(0) line) averaged in a beam of $`14^\mathrm{"}\times 27^\mathrm{"}`$ would be larger by a factor of 1.4, while on the opposite, the beam-averaged column density in the J=5 level would be smaller by a factor of 1.4. Hence, the correction for different apertures, cannot explain the zig-zag distribution in the rotational diagram.
A more critical correction is that for the extinction produced by the foreground material. As described by Martin-Pintado et al. (1999b ) the weakness of the S(3) line in the GC clouds should be due to the extinction produced by the silicate feature at 9.7 $`\mu `$m in the foreground dust clouds. In clouds with a LTE H<sub>2</sub> OTPR, one can use the intensity of the S(3) line to estimate the visual extinction once the relative value for the opacity at 9.7 $`\mu `$m to the 0.55 $`\mu `$m opacity is known. One could, in principle, apply corrections for increasing extinctions until the column density in the J=5 level is consistent with the column densities derived for other levels, i.e., until the rotational plot is a straight line (in the case of a Boltzmann distribution with one source temperature) or a smooth curve (in the case of a temperature gradient). In clouds with a non-equilibrium OTPR one could use an equivalent method using only ortho-H<sub>2</sub> levels, but obviously more than two levels are needed. The effect of foreground extinction on the rotational diagram is illustrated in Fig. 2, where the observed fluxes have been corrected for 30 (filled triangles) and 60 mag (filled circles) of visual extinction, using the extinction law of Draine & Lee (draine (1984)). Visual extinctions larger than 60 mag are needed for consistency between the S(1) and S(3) line intensities and a LTE OTPR. In this case, the curvature of the rotational plots suggests the presence of a large temperature gradient in the H<sub>2</sub> emitting region. To constrain the visual extinction toward these sources, in the following sections we will estimate the total column densities of dust and gas from measurements of the continuum dust emission, <sup>13</sup>CO, and C<sup>18</sup>O with a similar resolution to that of the SWS aperture.
### 3.1 H<sub>2</sub> column densities from C<sup>18</sup>O and <sup>13</sup>CO observations
We applied the Large Velocity Gradient (LVG) approximation to our data, to derive the physical conditions and the column densities of molecular gas from the emission of the J=2-1 line of C<sup>18</sup>O and the J=1-0 lines of C<sup>18</sup>O and <sup>13</sup>CO. The lines toward the two sources show complex profiles with two velocity components. From the line intensity ratios one can see that these components have slightly different physical conditions. The C<sup>18</sup>O J=2-1 to J=1-0 line ratio is 1.0-1.4 in M+3.06+0.34 and cannot be determined for the other source. The J=1-0 <sup>13</sup>CO to C<sup>18</sup>O ratio ranges between 5 and 14 in M+3.06+0.34 and is $`>`$ 7 in M+1.56$``$0.30. To within a factor of 2, these values are in agreement with the typical isotopic abundances found in the GC for carbon and oxygen (see Wilson & Matteucci 1994) indicating that the <sup>13</sup>CO lines are optically thin. From the C<sup>18</sup>O J=2-1 to J=1-0 ratio we derive for M+3.06+0.34 the H<sub>2</sub> densities given in Table 3 for two cases: high kinetic temperature ($`T_\mathrm{K}`$=100 K) and low kinetic temperature ($`T_\mathrm{K}`$=20 K). For those H<sub>2</sub> densities we have constrained the total column densities using the <sup>13</sup>CO line intensities. When the C<sup>18</sup>O lines were not detected the range of possible <sup>13</sup>CO(1-0) column densities was obtained by changing the H<sub>2</sub> density between 10<sup>3</sup> and 10<sup>4</sup> cm<sup>-3</sup> for $`T_\mathrm{K}`$=100 K and between 10<sup>3.5</sup> and 10<sup>4.5</sup> cm<sup>-3</sup> for $`T_\mathrm{K}`$= 20 K (see Hüttemeister et al. huttemeister98 (1998)). In the case of cold gas and even higher H<sub>2</sub> densities, the <sup>13</sup>CO column densities will increase only in a factor of 1.3 since for low temperatures and densities $``$ 10<sup>4</sup> the J=1–0 transition of <sup>13</sup>CO is thermalized. The H<sub>2</sub> column densities, $`N_{\mathrm{H}_2}`$, in Table 3 have been derived from the <sup>13</sup>CO column density and a fractional abundance with respect to H<sub>2</sub> of 5 10<sup>-6</sup>. They are typically of a few 10<sup>22</sup> cm<sup>-2</sup>, in good agreement with the values given by Hüttemeister et al. (huttemeister98 (1998)). With these column densities, we have derived the total visual extinction, $`A_\mathrm{v}`$, using the standard conversion factor: $`N_{\mathrm{H}_2}`$( cm<sup>-2</sup> )=$`A_\mathrm{v}(\mathrm{mag})\times 10^{21}`$. Thus the extinctions toward the two GC sources studied in this paper are typically of 15-20 magnitudes.
### 3.2 Dust column densities and temperatures
From the LWS data we can make a direct estimate of the dust temperature and the dust column densities toward both sources. Though the aperture of the LWS is larger than that of the SWS, the dust emission in the GC is relatively smooth (Odenwald & Fazio odenwald (1984)) and one does not expect large variations within the LWS aperture. The spectra for the two sources have very similar shapes with the maximum of the emission at $``$ 100 $`\mu `$m, indicating that the bulk of the dust is relatively cold with temperatures below 30 K, in agreement with previous estimates (Odenwald & Fazio odenwald (1984); Gautier et al. gautier (1984)).
The data cannot be fitted with only one gray body. For simplicity, we have considered a model with two gray bodies of temperatures $`T_1`$ and $`T_2`$. The total flux, $`S_\lambda `$, is given by:
$$S_\lambda =\mathrm{\Omega }[B(T_1,\lambda )(1e^{(1f)\tau (\lambda )})+B(T_2,\lambda )(1e^{f\tau (\lambda )})]$$
(1)
where $`\mathrm{\Omega }`$ is the solid angle of the continuum source, $`B(T)`$ is the Planck function, $`f`$ is the fraction of the opacity due to the warmer component ($`T_2`$), and $`\tau (\lambda )`$ is the total opacity at wavelength $`\lambda `$. In this model, the ratio of the visual extinction, $`A_\mathrm{v}`$, to the total optical depth at 30 $`\mu `$m is taken from the Draine & Lee (draine (1984)) extinction law and the opacity for $`\lambda `$ $`>`$ 30 is given by:
$$\tau (\lambda )=0.014A_\mathrm{v}(30\mu \mathrm{m}/\lambda )^\alpha $$
(2)
where $`\alpha `$ is the spectral index of the dust emission. In accordance with previous estimates for the envelope of Sgr B2 (Martín-Pintado et al. mp90 (1990)) and for the GC background of the cold core GCM 0.25+0.11 (Lis & Menten lis (1998)), we have taken $`\alpha 1`$. We have assumed extended emission ($`\mathrm{\Omega }=\mathrm{\Omega }_{\mathrm{LWS}}`$) and then we have fitted the continuum spectrum with $`f`$, $`A_\mathrm{v}`$, $`T_1`$, and $`T_2`$ as free parameters. As an example, we show in Fig. 3 the best fit to the LWS spectra towards M+1.56$``$0.30 obtained with $`A_\mathrm{v}`$=40, $`T_1`$=15 K, $`T_2`$=27 K and $`f`$=0.1.
Table 4 lists the results of the parameters for the best fits for the two sources. The visual extinctions derived for the two sources are 30 and 40 mag. These values are in agreement, to within a factor of 2, with those derived from the CO data.
The dust emission is dominated by the cool ($`T`$ 15 K) component ($`\tau _{\mathrm{v}_1}(1f)\tau _\mathrm{v}`$), while the slightly warmer component ($`T`$ 30 K) contributes only 10$`\%`$-20$`\%`$ to the total optical depth ($`\tau _{\mathrm{v}_2}f\tau _\mathrm{v}`$). We can also fit the spectra with larger spectral indexes by increasing the dust column densities. For instance, an spectral index of 1.5 will increase the visual extinction to 50-100 mag. These high values of $`A_\mathrm{v}`$ are very unlikely since they are almost one order of magnitude higher than the estimates made from CO (see Table 3).
Since the extinction derived from CO and the continuum accounts for the total gas and dust along the line of sight, they must represent an upper limit to the extinction to the H<sub>2</sub> emitting region. Considering the uncertainties introduced by the unknown spectral index and the many free parameters in the dust column density determination, in the following discussion we will assume that upper limits to the visual extinction of the H<sub>2</sub> emitting region are those derived from the CO emission, namely, 16 magnitudes for M+3.06+0.34 and 20 mag for M+1.56$``$0.30. These values are within a factor of two of estimates obtained from the total dust column density.
### 3.3 Warm H<sub>2</sub>: ortho-to-para ratio and column densities.
As discussed at the beginning of Sect. 3 the H<sub>2</sub> OTPR depends on the correction for extinction. In the previous sections we have estimated the extinction for the two clouds and Fig 4 shows the H<sub>2</sub> rotational diagrams for M+3.06+0.34 and M+1.56$``$0.30 corrected for the estimated extinctions. The error bars take into account the errors in the Gaussian fits of the lines and the calibration uncertainties. From these data, we derive an ortho rotational temperature, $`T_\mathrm{o}`$, from the ortho-H<sub>2</sub> levels J=3 and J=5. In the same way, one can define a para rotational temperature, $`T_\mathrm{p}`$, derived from the para-H<sub>2</sub> levels J=2 and J=4, and an ortho-para temperature, $`T_{\mathrm{op}}`$, derived from the ortho level J=3 and the para level J=2. These temperatures are listed in Table 5. As we see, $`T_\mathrm{p}`$ is $``$ 250 K for both sources while $`T_\mathrm{o}`$ is slightly higher ($``$ 270 K) indicating the presence of a moderate temperature gradient. This effect is more definite in other sources of our sample, where the S(4) and S(5) lines, which trace clearly higher temperatures, have also been observed (Martín-Pintado et al. 1999b ). For the present sample, $`T_{\mathrm{op}}`$ is $`160`$ K, much smaller than $`T_\mathrm{p}`$ and $`T_\mathrm{o}`$ indicating a non-LTE OTPR. In terms of these temperatures, the OTPR measured from our data will be given by:
$$\mathrm{OTPR}=\mathrm{OTPR}_{\mathrm{LTE}}(T_\mathrm{p})\mathrm{exp}(\frac{1}{T_\mathrm{p}}\frac{1}{T_{\mathrm{op}}})$$
(3)
where OTPR<sub>LTE</sub>(T) is the LTE OTPR at temperature $`T`$. As mentioned before, OTPR<sub>LTE</sub> is $``$ 3 for $`T`$ 200 K. Using Eq. (3), one finds an OTPR of $``$ 1 for both sources (see Table 5). Increasing the extinction will make the H<sub>2</sub> OTPR closer to the equilibrium value, however extinctions $`>`$ 70 mag will be required to give an LTE OTPR ratio. Such large visual extinctions are very unlikely from the molecular line and continuum data discussed in the previous sections. We therefore conclude that for the two sources the H<sub>2</sub> OTPR is not in equilibrium. Since the estimated error is $``$ 0.4, we can take $``$ 1.4 as a conservative upper limit for the OTPR in these two sources.
Extrapolating the populations in the J=2 and J=3 levels to the J=0 and J=1 levels , respectively, as two different species at temperature $`T_\mathrm{p}`$, one finds that the total column densities of warm H<sub>2</sub> are $``$ 2 10<sup>21</sup> cm<sup>-2</sup> . This must be considered as a lower limit to the actual warm H<sub>2</sub> column density since the populations of the lowest levels (J=0 and J=1) can be increased significantly by colder, though still warm ($``$ 100 K) gas. Of course, if extinction is higher column densities will also increase. This implies that the measured ratio of warm H<sub>2</sub> to cold gas traced by CO is at least 15$`\%`$.
High gas kinetic temperatures in these two clouds are known to be present from the NH<sub>3</sub> observations of Hüttemeister et al. (huttemeister93 (1993)). The rotational temperatures derived from the (4,4) and the (5,5) metastable inversion lines of NH<sub>3</sub> are in good agreement with the temperatures derived in this paper using the lowest H<sub>2</sub> pure-rotational lines. Extrapolating the populations in the (4,4) and the (5,5) NH<sub>3</sub> levels to lower levels with the rotational temperature derived for each source by Hüttemeister et al. (huttemeister93 (1993)) one finds a column density of warm NH<sub>3</sub> of $`\mathrm{7\; 10}^{14}`$ cm<sup>-2</sup> in both sources. Taking into account the warm H<sub>2</sub> column densities given above, we find a NH<sub>3</sub> abundance of (2-4) 10<sup>-7</sup>, similar to the value obtained by Martín-Pintado et al. (1999a ) in the expanding shells of the envelope of Sgr B2. A similar abundance is also obtained when we compare the column densities of cold ($``$ 20 K) NH<sub>3</sub> (Hüttemeister et al. huttemeister93 (1993)) and the H<sub>2</sub> column densities derived by our <sup>13</sup>CO and C<sup>18</sup>O data.
### 3.4 Warm dust column densities
If the gas and dust are coupled, one expects that the dust associated with the warm H<sub>2</sub> component would be an intense continuum emitter in the mid- and far-IR. There is no hint of such dust component in our data, as shown in Fig. 3, where we represent (as a dotted line) the emission of a gray body with a temperature of 250 K and the size of the SWS aperture attenuated by the total column density of the cold component. The equivalent H<sub>2</sub> column density of warm dust used to simulate the emission in Fig. 3 is only 5 10<sup>18</sup> cm<sup>-2</sup> . Even this small column density should have been detected. Hence, we can rule out a dust component coupled to the warm gas with a column density larger than 2 10<sup>-3</sup> times that of the warm H<sub>2</sub>. On the other hand, the comparison of CO emission with the cold dust emission shows agreement with the standard gas-to-dust ratio within a factor of two.
## 4 Discussion
### 4.1 Heating of the warm component
The large column densities of warm H<sub>2</sub> and the low column densities of associated warm dust require a heating mechanism that heats selectively the gas maintaining the dust at much lower temperatures. A PDR with an incident far-ultraviolet (FUV) flux $`G_0`$ of $``$ 100 (measured in units of 1.6 10<sup>-3</sup> ergs cm<sup>-2</sup> s<sup>-1</sup>) can heat the gas via photoelectric effect in the grains to temperatures of 100-200 K in the external layers of the cloud without heating the dust to temperatures above 30 K (see Hollenbach et al. hollenbach (1991)). However, the large gas phase NH<sub>3</sub> abundance, as derived in Sect. 3.3, is not possible in such a PDR scenario. The evaporation temperature of NH<sub>3</sub> is $``$ 75 K, therefore it cannot be evaporated from grain mantles at only 30 K. Even in the case that evaporation occurs, the UV radiation that heats the dust would destroy the fragile NH<sub>3</sub> molecule. This is the behavior found in NGC 7023 where the NH<sub>3</sub> abundance is $`10^8`$ in the well shielded region and decreases by more than a factor of 30 towards the region where the UV radiation increases and the dust temperature is $``$ 70 K (Fuente et al. fuente90 (1990)).
Shocks have been invoked as an important heating mechanism for the GC clouds (Wilson et al. wilson82 (1982); Martín-Pintado et al. mp97 (1997); Hüttemeister et al. huttemeister98 (1998)). In fact, M+1.56$``$0.30 belongs to the “$`l=`$1$`\stackrel{}{.}`$5-complex”, where Hüttemeister et al. (huttemeister98 (1998)) derived the highest SiO abundance within their sample, while the CS abundance (which traces all dense gas, not just the part that has been subjected to shocks) is not enhanced. They interpreted the SiO enhancement to be produced by large scale dynamic effects, proposing that in this complex, gas sprayed from the intersection of the $`x_1`$ and $`x_2`$ orbits is crashing into material that is still on $`x_1`$ orbits in the context of a bar morphology. In our sample, M+1.56$``$0.30 is also the source with the highest SiO to CS ratio. Furthermore, Dahmen et al. (dahmen (1997)), studying the HNCO emission in this region, found evidence for collisional excitation by shocks.
On the other hand, M+3.06+0.34 is located close to one of the CS cores detected by Stark & Bania (stark (1986)) in the “Clump 2” complex. These dense cores are gravitationally bound but most of the CO is emitted from the lower density gas, not bound to the cores. Stark & Bania (stark (1986)) suggested that this material is the result of tidal stripping of the cores. It is definite that shocks can play a role to explain the large column densities of warm H<sub>2</sub> and the relatively large abundances of NH<sub>3</sub>, as well as the high kinetic temperatures in these clouds. In addition, transient heating by shock waves provides a natural explanation for H<sub>2</sub> OTPRs out of equilibrium.
We have compared the results of the model calculations for slow shocks by Timmermann (timmermann (1998)) with our H<sub>2</sub> data. Interpolating the H<sub>2</sub> line strengths predicted by the model for a preshock OTPR=1 as a function of density, we found that a shock with velocity of 10 km s<sup>-1</sup> and a preshock H<sub>2</sub> density of $`\mathrm{2\; 10}^5`$ cm<sup>-3</sup> reproduces the observed line intensities. The results are displayed in Fig. 5 in the form of rotational plots. Open squares are the predicted column densities, while filled circles are the values derived from observations after correcting for extinction. Though the observed flux in the S(0) seems to be slightly larger than in the model, the agreement is excellent, and calibration errors can account for the discrepancies. The preshock density seems somewhat high but it is plausible since the S(3) line is apparently thermalized, which implies a lower limit to the H<sub>2</sub> density of $`10^4`$ cm<sup>-3</sup> . In any event, the widespread distribution of the HCN emission (Jackson et al. jackson (1996)) shows that densities of $`10^5`$ cm<sup>-3</sup> are common in the GC.
### 4.2 The ortho-to-para ratio
The main processes that affect the OTPR of H<sub>2</sub> are proton exchange collisions with H<sup>+</sup> and reactive collisions with H atoms. Ortho-para conversion in grain surfaces is thought to be less efficient. The rate coefficient for the proton exchange reaction
$$\mathrm{H}_2(\mathrm{ortho})+\mathrm{H}^+\mathrm{H}_2(\mathrm{para})+\mathrm{H}^++170.5\mathrm{K}$$
(4)
is $`\mathrm{3\; 10}^{10}`$ cm<sup>3</sup> s<sup>-1</sup> (Gerlich gerlich (1990)). The analogous reactions with H$`{}_{}{}^{+}{}_{3}{}^{}`$ and H<sub>3</sub>O<sup>+</sup> may also occur at a similar rate (see e.g. Le Bourlot et al. lebourlot (1999)). This rate gives an ortho-para conversion timescale, $`\tau _{\mathrm{conv}}`$, of $`100/n(^+)`$ yr, where $`n(^+)`$ represent the density of H<sup>+</sup>, H$`{}_{}{}^{+}{}_{3}{}^{}`$ or H<sub>3</sub>O<sup>+</sup> in cm<sup>-3</sup> . One should note that the actual conversion time can be a factor of 10 larger than $`\tau _{\mathrm{conv}}`$ (see Flower and Watt 1984).
The rate coefficient for the reactive collisions with H atoms
$$\mathrm{H}_2(\mathrm{ortho})+\mathrm{H}\mathrm{H}_2(\mathrm{para})+\mathrm{H}+170.5\mathrm{K}$$
(5)
is $`\mathrm{8\; 10}^{11}\mathrm{e}^{(3900/T)}`$ cm<sup>3</sup> s<sup>-1</sup> (see e.g. Le Bourlot et al. lebourlot (1999)). Due to the high activation barrier of this reaction (3900 K), in cold and dense molecular clouds the dominant process will be proton exchange collisions. This is also true for low velocity shocks of $``$ 10 km s<sup>-1</sup> since the maximum temperature achieved in the post-shock region is only $`300`$ K. If the H<sup>+</sup> and H$`{}_{}{}^{+}{}_{3}{}^{}`$ densities ($`n`$(H<sup>+</sup>), $`n`$(H$`{}_{}{}^{+}{}_{3}{}^{}`$)) in the postshock region of a 10 km s<sup>-1</sup> shock were as high as $`10^3`$ cm<sup>-3</sup> (see Timmermann timmermann (1998)), $`\tau _{\mathrm{conv}}`$ would be $`10^5`$ yr. It is worth-noting that recent models for ortho-para conversion in shocks by Wilgenbus et al. (wilgenbus (2000)) find much lower H<sup>+</sup> and H$`{}_{}{}^{+}{}_{3}{}^{}`$ densities in the postshock region. In this case, the timescale for ortho-para conversion would be $`>10^5`$ yr. On the other hand, the time needed for the passage of the proposed 10 km s<sup>-1</sup> velocity shock, from the point where the neutral gas starts to heat up to the point where the gas has reached interstellar temperatures again, is $`\stackrel{<}{}10^4`$ yr (see Timmermann timmermann (1998)). However the timescales in which the neutral gas is at high temperatures are much shorter. Hence, if the initial OTPR was lower than 3, the heating-cooling of the gas is too fast for the OTPR to reach the equilibrium at the temperatures of the shocked material.
Shocks with velocities $`>`$ 20 km s<sup>-1</sup> heat the gas to temperatures $`>`$ 700 K. Then, collisions with H would be the main ortho-to-para conversion mechanism, and indeed, the ortho-para conversion timescale would be low enough to obtain at least some conversion in the shock timescale as in the source HH54 (Neufeld et al. neufeld (1998)). However, the lines ratios in M+1.56$``$0.30 and M+3.06+0.34 cannot be explained with a preshock OTPR of $`<`$ 1 and a shock with velocity $`>`$10 km s<sup>-1</sup>. Therefore, the observed OTPR in these clouds must be approximately the preshock OTPR. This conclusion is independent of any shock model since the low temperatures involved by a 10 km s<sup>-1</sup> shock are not sufficient for the H<sub>2</sub>-H reactive collisions to be effective and, even for the largest predicted H<sup>+</sup> and H$`{}_{}{}^{+}{}_{3}{}^{}`$ abundances, the proton exchange reactions are not fast enough to give ortho-para conversion in the shock timescale.
If the OTPR of the preshock gas was in equilibrium at the gas temperature, the temperature should be $``$ 80 K. In this case, the preshock gas should have been already heated before the shock front compresses and heats the gas to 250 K. However, there is no strong reason to believe that the preshock OTPR should be in equilibrium at the preshock temperature. The H<sub>2</sub> molecule is formed mainly on the grain surfaces by a highly exothermic reaction. Thus, if it is rapidly ejected to gas phase the OTPR will be the typical OTPR at high temperature, i.e., 3. On the other hand, if it is not evaporated immediately from the grain there will be ortho-to-para conversion by collisions with radicals, impurities or defects and the OTPR could reach the equilibrium value at $``$ 30 K (dust temperature) of $``$0.01. In our case, the preshock OTPR of $`1`$ suggests that the H<sub>2</sub> molecules were ejected from the grains with OTPR $`>`$1. Afterwards, this ratio could decrease due to proton exchange processes.
The equilibrium proton abundance in dense ($`n`$(H<sub>2</sub>) $`10^5`$ cm<sup>-3</sup> ) clouds, where photoprocesses are not important, depends mainly on the ionization by cosmic rays and on charge exchange reactions with neutral molecules. Modeling the chemistry of dense PDRs, Sternberg and Dalgarno (sternberg2 (1995)) found $`n`$(H<sup>+</sup>) of $`10^5`$ cm<sup>-3</sup> in the well UV-shielded region for a cosmic ray ionization rate ($`\zeta `$) of 5 $`10^{17}\mathrm{s}^1`$, implying $`\tau _{\mathrm{conv}}10^7`$ yr. A similar timescale is obtained for proton exchange collisions with H$`{}_{}{}^{+}{}_{3}{}^{}`$. The density of H<sub>3</sub>O<sup>+</sup> could reach $`10^4`$ cm<sup>-3</sup> and thus $`\tau _{\mathrm{conv}}`$ could decrease by a factor of 10. Nevertheless, the actual time to reach the LTE OTPR would be longer. Flower & Watt (1984) have studied the temporal evolution of the OTPR in molecular clouds. Using the same rate coefficient as above for the proton exchange process, they have shown that for H<sup>+</sup> densities <sup>2</sup><sup>2</sup>2In their model these proton densities were obtained with $`\zeta =10^{17}10^{18}\mathrm{s}^1`$ using a simplified chemical network. of $`10^410^5`$ cm<sup>-3</sup> the actual time needed for an OTPR=3 to be in equilibrium at 30 K (similar to the observed cold component in the GC clouds) is $`10^710^8`$ yr. In particular, if $`n(\mathrm{H}^+)`$ (or $`n(\mathrm{H}_3\mathrm{O}^+)`$) is $`10^4`$ cm<sup>-3</sup> , then $``$ 5 10<sup>6</sup> yr will be needed to have an OTPR=1. Assuming that the H<sub>2</sub> was ejected to gas phase after formation with an OTPR $`\stackrel{<}{}`$3, the shock front reached the cloud approximately 10<sup>6</sup> yr after the formation of the H<sub>2</sub> molecules, since this is the time needed for an OTPR$``$3 to descend to $``$1 in a dense molecular cloud.
## 5 Conclusions
We have presented ISO SWS observations of the S(0), S(1), S(2), and S(3) pure-rotational lines of H<sub>2</sub> and LWS observations of the dust continuum and IRAM-30m <sup>13</sup>CO and C<sup>18</sup>O observations toward the GC molecular clouds M+1.56$``$0.30 and M+3.06+0.34 . Using the CO data and dust column densities from the LWS spectra we estimate $``$ 20 mag of visual extinction toward these sources. The two estimates of the extinction agree within a factor of 2 for the standard gas-to-dust conversion factor. According to the two components scenario proposed by Hüttemeister et al. (huttemeister93 (1993)), the low-J CO emission arises from the cold gas component and is coupled to the dust at a temperature of $`<`$ 30 K. The warm component ($`T`$ 250 K) column density observed directly in H<sub>2</sub> is, at least, $``$ 15 $`\%`$ of the cold one and would have very little warm dust associated with it. From the LWS spectra we set a conservative upper limit to the warm gas column density associated with the warm dust of 2 10<sup>-3</sup> times that of the warm H<sub>2</sub> column density.
After correcting for the dust extinction, we derive an OTPR of 1.0 $`\pm `$ 0.4, which is far from the LTE value expected for the gas temperatures of 250 K. We have also compared the warm H<sub>2</sub> column densities to the NH<sub>3</sub> observations by Hüttemeister et al. (huttemeister93 (1993)), and derived NH<sub>3</sub> abundances of $``$ 2 10<sup>-7</sup>, similar to those in the cold component (Hüttemeister et al. huttemeister93 (1993)).
The low dust temperatures, the high NH<sub>3</sub> abundances, the large CO linewidths, the non-LTE H<sub>2</sub> OTPR, in addition to the high gas temperatures suggest that the warm gas component is heated by low velocity shocks with speeds of $``$ 10 km s<sup>-1</sup> . To explain the OTPR we propose the following scenario. H<sub>2</sub> is formed in the grain surfaces and ejected to gas phase with OTPR$`\stackrel{<}{}`$3. After $`10^6`$ years, the time required to reach the preshock OTPR=1, a low velocity shock heated the gas to the observed temperatures of 250 K, but the OTPR was almost unaltered because the timescale for the passage of such a shock is much shorter than the ortho-to-para conversion timescale. Taking into account the shock timescale this occurred less than 10<sup>4</sup> yr ago. It is interesting to note that the timescale of the cloud’s galactic rotation period is also $`10^6`$ years. This fact suggests that the origin of the shocks can be related to large scale dynamics of the GC region.
###### Acknowledgements.
We thank S. Cabrit for her helpful comments on the ortho-to-para conversion mechanisms. We acknowledge support from the ISO Spectrometer Data Center at MPE, funded by DARA under grant 50 QI 9402 3. NJR-F, JM-P, PdV, and AF have been partially supported by the CYCIT and the PNIE under grants PB96-104, 1FD97-1442 and ESP97-1490-E. NJR-F acknowledges Conserjería de Educación y Cultura de la Comunidad de Madrid for a pre-doctoral fellowship.
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# Supersymmetric Pair Correlation Function of Wilson Loops
## I Introduction
Quantitative universal predictions for the low energy limit of String/M Theory that are independent of specific backgrounds or compactifications are hard to come by. In this paper, we take a step in the direction of quantitative universality, extracting a universal numerical result for the pair potential at short distances between heavy nonrelativistic sources in a supersymmetric gauge theory on some generic D-manifold background of the type I or type II string theories. The result applies irrespective of whether the manifold, $`𝒟`$, is the worldvolume of a higher dimensional Dbrane, the intersection of multiple Dbranes, or the bulk transverse space orthogonal to some configuration of branes.
Our result is obtained from a path integral prescription for the pair correlation function of Wilson loops living in some D-manifold in a weakly coupled background of the type II theory, based on the earlier works . We give a boundary reparameterization invariant computation of the supersymmetric pair correlation function of Wilson loops in the open and closed fermionic string theory with Dbranes . The normalization of the one-loop string vacuum amplitude in such a background can be determined from first principles following a classic method due to Polchinski . However, it will be apparent from our result that the prediction of a short distance potential originating in fluctuations in the vacuum energy density is likely to hold in the broader context of the generic background of String/M theory. Our prescription for determining the phase ambiguities in the fermionic string path integral is derived from the imposition of infrared consistency conditions which follow from matching to an appropriate supergravity theory, the low energy theory at long distances. It is motivated in part by ideas taken from and by the more unified description of self-duality and background fields that appears in the recent papers .
In an earlier paper with Chen we obtained the short distance potential between heavy sources in a gauge theory in some generic background of the bosonic string. The potential is extracted from a covariant string path integral representation of the pair correlation function of Wilson loops. Our results apply both in the background with $`d`$$`=`$$`26`$ spacetime dimensions, or in the presence of a generic background for the Liouville theory with fewer matter fields, $`c_m`$$`<`$$`d`$. We find an attractive, and scale invariant, $`1/r`$ short distance interaction between the heavy gauge theory sources. We note that the bosonic string has a tachyon, which must either be stabilized (see the recent discussions in ), or eliminated— as is possible in the type I and type II string theories. The bosonic results are a useful warm-up for the computation of the annulus in stable backgrounds of the superstring. They also capture the correct qualitative features of the short distance potential in a background of the superstring with a tachyon instability. The calculation in the bosonic string proceeds as follows .
We consider heavy sources in the gauge theory in relative collinear motion with $`r^2`$$`=`$$`R^2`$$`+`$$`v^2\tau ^2`$, $`v`$$`<<`$$`1`$, thus giving a simple realization of coplanar loops while mimicking straight-line trajectories in the Euclideanized $`X^0`$, $`X^1`$ plane. Here $`r`$ is their relative position, and $`\tau `$ is the zero mode of the Euclideanized time coordinate, $`X^0`$. The scattering plane is wrapped into a spacetime cylinder by periodically identifying the coordinate $`X^0`$. Then the closed world-lines of the heavy sources are loops singly wound about this cylinder. We identify these closed world lines with Wilson loops. Following the earlier works of Alvarez , and of Cohen et al , we give a path integral prescription for the pair correlation function of Wilson loops. The loops can be taken to lie in the world-volume of a higher dimensional Dpbrane. Taking the large loop length limit of the correlation function, $`L_i`$$``$$`L_f`$$``$$`T`$$``$$`\mathrm{}`$, with $`R`$ held fixed, we define an effective potential as follows:
$$<M(C_i)M(C_f)>=i\underset{T\mathrm{}}{lim}_T^{+T}𝑑\tau V_{\mathrm{eff}.}[r(\tau ),u].$$
(1)
The dominant contribution to the potential between the sources at short distances is from the massless modes in the open string spectrum. Suppressing the tachyon, and restricting to the massless modes of the bosonic string, the potential can be expressed as a systematic double expansion in small velocities and short distances with the result :
$$V_{\mathrm{bos}.}(r,u)=(d2)\frac{1}{r}+O(z^2,uz/\pi ,u^2),$$
(2)
where $`z`$ is the dimensionless scaling variable, $`z`$$`=`$$`r_{\mathrm{min}.}^2/r^2`$, and $`r_{\mathrm{min}.}^2`$$`=`$$`2\pi \alpha ^{}u`$ is the minimum distance scale probed in the collinear scattering of the heavy point sources. The $`1/r`$ static term receives velocity dependent corrections which are parameterized by the dimensionless variables, $`z^2`$, $`uz/\pi `$, and $`u^2`$. We show that the small velocity short distance approximation is self-consistent for distances in the regime, $`2\pi \alpha ^{}u`$$`<<`$$`r^2`$$`<<`$$`2\pi \alpha ^{}`$, and velocities in the range, $`u`$$`<<`$$`u_+`$, where the upper bound, $`u_+`$, on permissible velocities can be estimated as described in . Thus, String/M theory predicts velocity dependent corrections to the potential between two heavy sources in relative slow motion in a gauge theory, the numerical coefficients of which are given by a systematic expansion.
Evidence of a distance scale shorter than the string scale was originally found in the nonrelativistic scattering of D0branes , which gives a linear potential in the bosonic string theory . The D0branes are assumed to have fixed spatial separation in a direction $`X^{d1}`$, and in relative motion with nonrelativistic velocity $`v`$ in an orthogonal direction $`X^d`$. The static linear potential between a pair of bosonic D0branes corresponds to a shift in the vacuum energy density relative to the background with no Dbrane sources due to a constant background electromagnetic vector potential with vanishing electric field strength : $`A^\mu `$$`=`$$`\overline{A}^\mu `$, $`\mu `$$`=`$$`d`$$``$$`1`$, $`_0\overline{A}^\mu `$$`=`$$`0`$. The systematics of the small velocity short distance double expansion, and the value $`r_{\mathrm{min}.}^2`$$`=`$$`2\pi \alpha ^{}u`$ for the minimum distance probed in the scattering of D0branes, are in precise agreement with our results for the short distance potential between heavy gauge theory sources.
In this paper, we extend these results to the generic stable background of the supersymmetric type I and type II theories. We give a path integral prescription for the supersymmetric annulus amplitude, determining both its phase and its normalization from first principles. Boundary reparameterization invariance is imposed following the analysis of the bosonic string amplitude given in . Our prediction for the short distance potential between heavy gauge theory sources in supersymmetric string theory is:
$$V_{\mathrm{super}.}(r,u)=\frac{u^4}{r^9}2^4\pi ^{7/2}\alpha _{}^{}{}_{}{}^{4}\mathrm{\Gamma }(\frac{9}{2})+O(u^6),$$
(3)
which can be compared with the bosonic string result given in Eq. (2). The systematics of the velocity dependent corrections are much simpler in the fermionic string. This is a consequence of the BPS conditions or, equivalently, as in this example, of spacetime supersymmetry. The key ingredient which enables a prediction of the numerical coefficient in the short distance potential is its relationship to the vacuum energy computation in string theory: unlike in quantum field theories, the one-loop cosmological constant in critical string theory can be unambiguously normalized, an observation due to Polchinski .
We begin in section II with a brief discussion of supersymmetric boundary conditions and spin structure. Spinor conventions, and a recapitulation of the local symmetries of the world-sheet action for the fermionic string, are given in appendix A. Section III contains an evaluation of the supersymmetric annulus from first principles, beginning with the covariant path integral over world-sheets of specified spin structure and specified boundary condition on all of the world-sheet fields. The gauge fixing of the local world-sheet symmetries is carried out in section IIIA where we give a derivation of the supersymmetric annulus with boundaries on parallel and static Dbranes, up to undetermined phases. The precise global structure of the gauge orbit of the supersymmetry and super-Weyl transformations on the world-sheet is not known . We will take the point of view that global ambiguities in the superstring path integral can be eliminated by the imposition of infrared consistency conditions that require matching to an appropriate supergravity theory describing the low energy physics. This prescription for the phases of the fermionic path integrals summed in the vacuum amplitude is given in section IIIB, using simple, and universal, infrared consistency conditions on the long distance physics. The result is an unambiguous determination of both the normalization and the phase of the supersymmetric annulus. Implementing boundary reparameterization invariance in the path integral as in , we derive an expression for the supersymmetric pair correlation function of Wilson loops in section IIIC. Finally, we extend our results in section IIID to generic boundary conditions corresponding to Dbranes in relative motion, imposing appropriate infrared consistency conditions as before. As a consistency check, we compute the short distance potential probed in the nonrelativistic small angle scattering of Dpbranes, recovering the numerical coefficient previously obtained in .
In section IV we adapt these results to the supersymmetric pair correlation function of Wilson loops corresponding to worldlines of heavy sources in relative slow motion. We show that the short distance potential between heavy sources in a supersymmetric gauge theory takes the form of a scale invariant $`1/r`$ fall-off contributed by the bosonic degrees of freedom, multiplicatively corrected by a convergent power series expansion in the dimensionless variable $`z`$$`=`$$`r_{\mathrm{min}.}^2/r^2`$. The leading term in the potential is given by the expression in Eq. (3). Some implications of our results are sketched in the conclusions.
## II Supersymmetric Boundary Conditions and Spin Structure
We begin with the Brink-Di Vecchia-Howe-Deser-Zumino world-sheet action used in Polyakov’s path integral quantization of the fermionic string :
$$S_{SP}=\frac{1}{2\pi \alpha ^{}}d^2\sigma \sqrt{g}[\frac{1}{2}g^{mn}_mX^\mu _nX_\mu +\frac{1}{2}\alpha ^{}\overline{\psi }^\mu \gamma ^m_m\psi _\mu +\sqrt{\frac{\alpha ^{}}{2}}(\overline{\chi }_a\gamma ^m\gamma ^a\psi ^\mu )(_mX_\mu )+\frac{\alpha ^{}}{4}(\overline{\chi }_a\gamma ^b\gamma ^a\psi ^\mu )(\overline{\chi }_b\psi _\mu )],$$
(4)
invariant under both reparameterizations and local $`N`$$`=`$$`1`$ world-sheet supersymmetry transformations. The indices $`m`$,$`n`$$`=`$$`1`$,$`2`$ label the world-sheet coordinates for the string with metric $`g`$, and $`a`$,$`b`$$`=`$$`1`$, $`2`$ label the flat local tangent space to the world-sheet. Spinor conventions and the local symmetries underlying the action are reviewed in the appendix. We use the label $`\mu `$$`=`$$`0`$, $`\mathrm{}`$, $`p`$ for the Neumann directions parallel to the worldvolume of the Dpbranes. The branes are spatially separated by $`R`$ in the $`X^9`$ direction, with $`\mu `$$`=`$$`p`$$`+`$$`1`$, $`\mathrm{}`$, $`9`$ labeling the Dirichlet directions orthogonal to the worldvolume of the Dpbranes. The boundary conditions on the embedding coordinates, $`X`$, and their fermionic superpartners, $`\psi `$, are obtained from the kinetic term in Eq. (A14). The corresponding free field action in a flat embedding spacetime with component fermions, $`\psi _\mu ^\pm `$, is :
$$S[X,\psi ]=\frac{1}{4\pi \alpha ^{}}d^2\sigma [^mX^\mu _mX_\mu +\alpha ^{}(\psi ^\mu (_1i_2)\psi _\mu ^{}+\psi ^{+\mu }(_1+i_2)\psi _\mu ^+)],$$
(5)
which extends to the locally supersymmetric action given above. A variation of the classical action with respect to the embedding coordinate $`X`$ gives the surface term:
$$\delta X^\mu (n^a_aX_\mu )=0.$$
(6)
As possible boundary conditions we list:
$`\mathrm{N}`$ $`:n^a_aX^\mu =0`$ (7)
$`\mathrm{D}`$ $`:X^\mu =y^\mu `$ (8)
$`\mathrm{W}`$ $`:\delta X^\mu t^a_aX^\mu ,`$ (9)
where $`y^\mu `$, $`\mu `$$`=`$$`p+`$$`1`$, $`\mathrm{}`$, $`9`$, gives the spacetime location of the Dbrane. The $`W`$, or modified Dirichlet (MD), boundary condition is motivated by the Wilson loop problem . It permits fluctuations in the world-sheet fields tangential to the boundary. The boundaries of the world-sheet have been identified with the closed world-lines of a heavy quark–antiquark pair in the gauge theory.
A point source undergoing straight line motion with nonrelativistic velocity $`v`$ in the $`X^0`$, $`X^1`$ plane with respect to the origin, $`X^0`$$`=`$$`X^1`$$`=`$$`0`$, and in zero external field, is described by the boundary conditions:
$$\mathrm{V}:n^a_a(X^0vX^1)=0,X^1=vX^0,$$
(10)
with N (D) boundary conditions imposed on the $`X^0`$ $`(X^1)`$ coordinates of the source fixed at the origin. The boundary conditions on two point sources in relative motion in a $`𝒟`$-manifold are rather simple, irrespective of whether the motion occurs within the worldvolume of a higher dimensional Dbrane, the intersection of two or more Dbranes, or in the bulk transverse space orthogonal to some configuration of Dbranes. The point sources are the end-points of open strings. Then we distinguish the $`d`$$`=`$$`10`$ embedding coordinates of the world-sheet as NN, ND, or DD, directions, depending on whether both, one, or neither, point source has nonvanishing spacetime momentum in the direction of the coordinate . In the discussion that follows, we restrict ourselves to NN and DD coordinates alone. Note that identical boundary conditions must be imposed on all of the NN, and all of the DD, fermions in order to preserve the global $`SO(1,p)`$$`\times `$$`SO(9p)`$ symmetry of the vacuum amplitude. The world-sheet gravitino satisfies the same boundary conditions as the NN fermions— this is dictated by supersymmetry.
Consider the variation of the world-sheet action with respect to the fermion field. The vanishing of the surface term dictates the following condition on fermion bilinears on the boundary:
$$\psi ^{+\mu }(\delta \psi _{+\mu })=\psi ^\mu (\delta \psi _\mu ),$$
(11)
with solutions, $`\psi ^{+\mu }`$$`=`$$`\pm \psi ^\mu `$, at any boundary. As a check, we perform a supersymmetry transformation on the surface term. This gives the condition, $`n_a[(_bX_\mu )(\overline{\xi }\gamma ^a\gamma ^b\psi ^\mu )]`$$`=`$$`0`$. Align the world-sheet with coordinate $`\sigma ^2`$ normal, and $`\sigma ^1`$ tangential, to the boundary. Assume Neumann boundary conditions on the embedding coordinate $`X_\mu `$: $`_2X_\mu `$$`=`$$`0`$, with $`_1X_\mu `$$``$$`0`$. Then the requirement that the boundary conditions on the $`\psi ^\mu `$ preserve world-sheet supersymmetry implies the restriction, $`\xi ^+`$$`=`$$`\xi ^{}`$, on permissible supersymmetry transformations on the boundary. With the Dirichlet boundary condition, $`_1X_\mu `$$`=`$$`0`$, we obtain the constraint, $`\overline{\xi }\psi ^\mu `$$`=`$$`0`$. Thus, choosing one or other sign for the NN fermions simultaneously determines the choice of phase for the DD fermions. In component form, the choice $`\xi ^+`$$`=`$$`\xi ^{}`$ implies that $`\chi ^\mu `$$`=`$$`\chi ^{+\mu }`$. The constraints on the fermionic fields for the W and V boundary conditions on the world-sheet are a straightforward generalization of this reasoning.
The annulus amplitude in a flat spacetime background of the supersymmetric string can be represented as a path integral summing over fluctuations of world-sheets with cylindrical topology weighted by the action in Eq. (A14):
$$𝒜=\frac{1}{2}\underset{\beta ,\alpha 0,1}{}C_\alpha ^\beta _{[\beta ,\alpha ]}\frac{[dX][d\psi ][dg][d\chi ]}{\mathrm{Vol}(gauge)}e^{S_{SP}[X,\psi ,g,\chi ]\mu _0_{}d^2\sigma \sqrt{g}S_{\mathrm{ren}.}}.$$
(12)
Thus, the world-sheets of the fermionic string are endowed with additional degrees of freedom, and the quantum fluctuations about some minimum action configuration summed in the path integral must include a consideration of these modes. We consider the simplest background configuration of static parallel Dpbranes separated by a distance $`R`$. The stretched string between the Dbranes contributes a term, $`R^2l/4\pi \alpha ^{}`$, to the action in Eq. (12), further corrections being suppressed at weak coupling. Since the one-loop vacuum amplitude is a sum over surfaces of cylindrical topology with Euler characteristic $`\chi `$$`=`$$`0`$, the amplitude is free of any dependence on the string coupling constant. Also, we can drop boundary cosmological constant terms in favor of the bulk cosmological constant, $`\mu _0`$, since these are not independent Lagrange parameters on the cylinder. We will gauge all of the local symmetries reviewed in appendix A. $`S_{\mathrm{ren}.}`$ contains any additional counterterms that may be necessary in order to obtain amplitudes invariant under both world-sheet diffeomorphisms and Weyl transformations of the metric, as well as local supersymmetry transformations and super-Weyl rescalings of the world-sheet gravitino. Divergent contributions to the path integral arising from local gauge anomalies in the measure will be absorbed in a renormalization of the bare couplings introduced in $`S_{\mathrm{ren}.}`$, including the bulk cosmological constant, the renormalized values being set to zero at the end of the calculation .
We turn next to global aspects of the path integral. We are summing over spin structures and averaging over the $`\pm `$ ambiguity in the boundary condition on the fermions at the boundary of the world-sheet. The label $`\alpha `$ on the path integral refers to a choice of spin structure on the world-sheet: the change in phase in a Weyl fermion upon traversal of a closed path homotopic to either boundary of the cylinder. Under $`\sigma ^1`$$``$$`\sigma ^1+1`$, the left and right-moving component fermions transform as:
$$\psi ^{\pm \mu }(\sigma ^1+1,\sigma ^2)=e^{\pi i\alpha }\psi ^{\pm \mu }(\sigma ^1,\sigma ^2).$$
(13)
The parameter $`\alpha `$$`=`$$`0`$ ($`1`$) labels the string path integral computed with world-sheet spinors that are, respectively, anti-periodic (periodic) around the single closed cycle of the cylinder. The parameter $`\beta `$ denotes the ambiguity, $`\psi ^{+\mu }`$$`=`$$`\pm \psi ^\mu `$, at any boundary, described above. Specifically, for either NN or DD fermions, we will define $`\beta `$ as follows:
$$\psi ^{+\mu }(\sigma ^1,0)=e^{\pi i\beta }\psi ^\mu (\sigma ^1,0),\mathrm{with}\psi ^{+\mu }(\sigma ^1,1)=\psi ^\mu (\sigma ^1,1).$$
(14)
Choosing the phase at the $`\sigma ^2`$$`=`$$`1`$ end-point to correspond to periodic fermions is pure convention. Thus, $`\beta `$$`=`$$`0`$ $`(1)`$ corresponds to reflection with (without) a phase change of $`\pi `$ in the fermionic wavefunction at the $`\sigma ^2`$$`=`$$`0`$ end-point. We remark that this convention corresponds to that in the text . In the critical dimension, and with an unambiguous prescription for the phases, $`C_\alpha ^\beta `$, of the path integrals, there will be no global gravitational or Lorentz anomalies. Our prescription for determining the absence of global phase ambiguities in the one-loop vacuum amplitude will be physical, rather than constructive. It is motivated by infrared consistency conditions on the long distance physics, as will be clarified in section IIIB. For the discussion in section IIIA, we encourage the reader to think of the $`C_\alpha ^\beta `$ as unspecified, and therefore ambiguous, phases that weight the different contributions to the annulus amplitude.
## III Path Integral Evaluation of Supersymmetric Annulus
We will now give a path integral evaluation of the supersymmetric annulus— a sum over orientable world-sheets with boundaries on parallel Dpbranes, paying special attention to the imposition of world-sheet supersymmetry both in the bulk, and on the boundary. We begin with the simplest configuration of parallel and static Dpbranes, generalizing our results in section IIID for the $`V`$ boundary conditions describing Dpbranes in relative motion. The gauge fixing of the local symmetries on the world-sheet, and a derivation from first principles of the annulus amplitude up to unknown phases, is described in section IIIA. In section IIIB, we show how infrared consistency conditions can be used to determine all of the phase ambiguities in the path integral. We extend this derivation in IIIC to an analysis of the supersymmetric pair correlation function of Wilson loops following the treatment in . The extension to generic boundary conditions describing branes in relative motion is given in section IIID. This result will be adapted in section IV to give an expression for the supersymmetric pair correlation function of Wilson loops corresponding to the worldlines of heavy gauge theory sources in slow relative motion. The normalization and the phase of the pair correlation function is therefore precisely, and unambiguously, determined, allowing a derivation of the short distance potential between the sources.
### A Gauge Fixing of the Local World-sheet Symmetries
We begin by gauge fixing the Lorentz transformations in the local tangent space at any point in the world-sheet, thereby eliminating one of the four bosonic gauge parameters. This implies that, although it is convenient to write classically covariant expressions in terms of zweibeins $`e_m^a`$, the number of physical degrees of freedom in the path integral is the same as with metrics: we use local Lorentz rotations to eliminate one of the independent degrees of freedom in the zweibein. Likewise, we eliminate two of four independent modes of the gravitino by choosing super-conformal gauge, $`\gamma ^m\chi _m`$$`=`$$`0`$ , thereby gauge fixing super-Weyl transformations. This immediately creates an apparent problem with supersymmetry since we have only two fermionic but three bosonic gauge parameters, but this is not so. Recall that we work in the critical spacetime dimension gauging Weyl transformations of the metric. As in the path integral quantization of the bosonic string , although it is convenient to keep the Weyl mode in the discussion of the measure, the principle of ultralocality requires that any explicit dependence on the Weyl mode only contribute local, renormalizable, terms to the effective action . The unique choice for such terms is the Liouville action. Thus, in the critical dimension, the Weyl mode entirely decouples.
We will employ similar reasoning in gauging local Lorentz and super-Weyl transformations: ultralocality of the measure in the path integral requires that any explicit dependence on the Weyl and super-Weyl modes can only contribute terms proportional to the supersymmetric Liouville action . Thus, although we find it convenient to keep all four bosonic and fermionic gauge parameters in the discussion below, there are really half as many physical gauge parameters in the critical superstring corresponding, respectively, to diffeomorphisms and local supersymmetry transformations. Following gauge fixing, the path integrals reduce to ordinary integrals over constant modes and moduli. The counting of superconformal Killing spinors and supermoduli on a Riemann surface is given by the supersymmetric analog of the Riemann-Roch theorem. For cylindrical topology, the answer is rather simple, since both the superconformal Killing spinors and the supermoduli are simply constant spinors. These can only exist on a cylinder with both periodic spin structure and the periodic boundary condition at the endpoints of the open string.
The zweibein and metric are related by $`g_{mn}`$$`=`$$`e_m^ae_n^b\delta _{ab}`$. We make the same choice of fiducial world-sheet metric as in the analysis of the bosonic string:
$$ds^2=l^2(d\sigma ^1)^2+(d\sigma ^2)^2,0\sigma ^11,0\sigma ^21.$$
(15)
Thus, the fiducial einbein on the boundary is $`\widehat{e}`$$`=`$$`\sqrt{g}`$, and the modulus $`l`$ can be identified as the boundary length, $`l`$$`=`$$`_0^1𝑑\sigma ^1\widehat{e}`$ . With this choice, the area of the fiducial world-sheet, $`d^2\sigma \sqrt{\widehat{g}}`$, equals $`l`$.
Begin with the integration over the bosonic embedding coordinates, $`X^\mu `$. Integration over $`p`$$`+`$$`1`$ constant modes in the Neumann directions, and upon normalizing the measure for infinitesimal variations as in we obtain the result:
$$[dX]e^{S_{SP}[X,\psi ,g,\chi ]}=2iV_{p+1}(4\pi ^2\alpha ^{})^{(p+1)/2}l^{(p+1)/2}e^{R^2l/4\pi \alpha ^{}}(\mathrm{det}^{}\mathrm{\Delta })_g^5e^{S_{\mathrm{eff}.}[\psi ,\chi ,\widehat{g}]},$$
(16)
where the determinant for the scalar Laplacian is computed with the NN boundary condition on $`p`$$`+`$$`1`$, and DD boundary condition on $`9`$$``$$`p`$, coordinates. The overall factor of two accounts for the two possible orientations of the open string. The effective action for the fermionic fields takes the form:
$`S_{\mathrm{eff}.}=`$ $`{\displaystyle \frac{1}{2\pi \alpha ^{}}}{\displaystyle d^2\sigma \sqrt{g}[\frac{1}{2}\alpha ^{}\overline{\psi }^\mu \gamma ^m_m\psi _\mu +\frac{\alpha ^{}}{2}(\overline{\chi }_a\gamma ^b\gamma ^a\psi ^\mu )(\overline{\chi }_b\psi _\mu )]}`$ (18)
$`{\displaystyle d^2\sigma \sqrt{g}(\overline{\chi }_a\gamma ^m\gamma ^a\psi ^\mu )(\sigma )d^2\sigma ^{}\sqrt{g}(\overline{\chi }_a\gamma _m\gamma ^a\psi _\mu )(\sigma ^{})_{\sigma \sigma ^{}}^2G(\sigma ,\sigma ^{})}.`$
$`G(\sigma ,\sigma ^{})`$ is the Greens function of the scalar Laplacian.
Consider now the gauge fixing of the local symmetries of the effective action. A bosonic deformation of the metric is decomposed into a Weyl transformation, a diffeomorphism continuously connected to the identity, and a change in the length of the boundary of the cylinder. We will gauge both diffeomorphisms and Weyl transformations of the metric. Defining implicitly the measure for infinitesimal variations in the tangent space to the space of metrics as in gives the result:
$$\frac{1}{\mathrm{Order}(\stackrel{~}{D})}\frac{[dg](\mathrm{det}^{}\mathrm{\Delta })_g^5}{\mathrm{Vol}(\mathrm{Diff}_0\times \mathrm{Weyl})}=\frac{1}{2}\frac{[d\delta V]}{\mathrm{Vol}(\mathrm{Diff})_0}\frac{[d\delta \varphi ]}{\mathrm{Vol}(\mathrm{Weyl})}_0^{\mathrm{}}𝑑lJ_b(l;\widehat{g})e^{\frac{10d}{48\pi }S_L[\varphi ,g]}(\mathrm{det}^{}\mathrm{\Delta })_{\widehat{g}}^5,$$
(19)
where $`J_b`$ is the Jacobian from the change of variables computed in :
$$J_b=\frac{(l/2\pi )^{1/2}(\frac{2}{l^2})^{1/2}(\frac{1}{2}l^2\eta ^4(\frac{il}{2}))^{1/2}}{(l^3/2\pi )^{1/2}},$$
(20)
in the cylinder metric specified above. We have divided by the order of the subgroup of the disconnected component of the diffeomorphism group, $`\stackrel{~}{D}`$: discrete diffeomorphisms of the world-sheet left invariant under the choice of superconformal gauge . This gives a factor of two in the denominator of Eq. (19), correcting for the two-fold invariance of the measure under the diffeomorphism: $`\sigma ^1`$$``$$`\sigma ^1`$. The Weyl anomaly of the measure exponentiates to a term proportional to the Liouville action, whose coefficient vanishes in the critical spacetime dimension, $`d`$. Consider the result of performing a local supersymmetry and super-Weyl transformation on this expression. This will induce a super-Weyl anomaly. We must simultaneously include the contributions from world-sheet fermions to the Weyl anomaly.
An arbitrary fermionic deformation of the world-sheet gravitino can be decomposed:
$$\delta \chi _m=D_m\delta \xi +(\delta \zeta )\gamma _m+(_\alpha \chi _m)(\delta \nu ^\alpha ),$$
(21)
where $`\delta \xi `$ is an infinitesimal supersymmetry transformation, $`\delta \zeta `$ is a rescaling of the gravitino, and $`\delta \nu `$ is a change in a possible supermodulus— a constant two component spinor, $`\nu `$, on the cylinder. We work in superconformal gauge, invoking super-Weyl transformations in setting $`\gamma ^m\chi _m`$$`=`$$`0`$ . We make an orthogonal decomposition into infinitesimal deformations parallel (perpendicular) to the gauge slice, respectively, preserving (violating) the restriction to gamma traceless $`\chi _m`$. In addition, we must separate the contribution from superconformal Killing spinors, $`\delta \xi _0`$, which leave the world-sheet gravitino unchanged:
$$D_m(\delta \xi _0)+(\delta \zeta _0)\gamma _m=(D_m+\frac{1}{2}\gamma _m\gamma ^nD_n)(\delta \xi _0)=0,$$
(22)
since the supermodulus, $`\delta \zeta _0`$$`=`$$`\frac{1}{2}\gamma ^mD_m(\delta \xi _0)`$, is in the kernel of the operator $`D_m`$. The superconformal Killing spinor, $`\delta \xi _0`$, is in the kernel of the operator, $`P_{1/2}`$$`=`$$`2D_m`$$`+`$$`\gamma _m\gamma ^nD_n`$. In the fiducial cylinder metric, $`D_m`$$`=`$$`_m`$, and the Killing spinor is simply a constant spinor. Likewise, the supermodulus is also a constant spinor. Constant spinors can only exist on cylindrical world-sheets with both periodic spin structure and periodic boundary condition at the endpoints of the open string. Thus, for $`\alpha `$$`=`$$`\beta `$$`=`$$`1`$, we have one superconformal constant spinor and one supermodulus to be accounted for in the measure. Performing the change of variables in Eq. (21) we can write:
$$_{[\beta ,\alpha ]}\frac{[d\delta \chi _m]}{\mathrm{Vol}(\mathrm{sWeyl}\times \mathrm{sDiff})}=\frac{[d\delta \xi ]^{}}{\mathrm{Vol}(\mathrm{sDiff})_0}\frac{[d\delta \zeta ]}{\mathrm{Vol}(\mathrm{sWeyl})}_{[\beta ,\alpha ]}[d\delta \nu ]J_f^{(\beta ,\alpha )},$$
(23)
where the integration over a supermodulus is absent unless $`\alpha `$$`=`$$`\beta `$$`=`$$`1`$. A norm in the tangent space to the space of the spin $`3/2`$ field $`\chi _m`$ invariant under both reparameterizations and local Lorentz transformations can be written as follows:
$$|\delta \chi _m|^2=id^2\sigma \sqrt{g}(\delta \overline{\chi }^m)(\delta \chi _m)=2d^2\sigma \sqrt{g}(\delta \chi _1^+\delta \chi _1^{}+\delta \chi _2^+\delta \chi _2^{}),$$
(24)
where the measure in the path integral is normalized as in :
$$1=[d\delta \chi _m]e^{|\delta \chi _m|^2/2}[(d\delta \chi _1^+)(d\delta \chi _1^{})(d\delta \chi _2^+)(d\delta \chi _2^{})]e^{|\delta \chi _m|^2/2}.$$
(25)
We can likewise define a norm in the tangent space to the space of $`\xi `$:
$$|\delta \xi |^2=id^2\sigma \sqrt{g}(\delta \overline{\xi })(\delta \xi ),1[(d\delta \xi ^+)(d\delta \xi ^{})]e^{{\scriptscriptstyle d^2\sigma \sqrt{g}(\delta \xi ^+)(\delta \xi ^{})}}.$$
(26)
Separating the ordinary Grassmann integral over the superconformal Killing spinor, $`\xi _0`$, gives:
$$1=[d\delta \xi _0]e^{{\scriptscriptstyle \frac{i}{2}}(\delta \overline{\xi }_0)({\scriptscriptstyle d^2\sigma \sqrt{g}})(\delta \xi _0)}[(d\delta \xi ^+)^{}(d\delta \xi ^{})^{}]e^{|\delta \xi |^2/2}=l[(d\delta \xi ^+)^{}(d\delta \xi ^{})^{}]e^{|\delta \xi |^2/2}.$$
(27)
Likewise, the ordinary Grassmann integration over a supermodulus gives:
$$[d\delta \nu ]e^{{\scriptscriptstyle \frac{i}{2}}(\delta \overline{\nu })({\scriptscriptstyle d^2\sigma \sqrt{g}})(\delta \nu )}=l.$$
(28)
This last term is only necessary when computing the path integral with $`\alpha `$$`=`$$`\beta `$$`=`$$`1`$. The measure for the path integral with constant spinors on the worldsheet is discussed in appendix B. For a configuration of parallel and static Dbranes, the presence of the constant mode on the world-sheet gives a vanishing result for the path integral. Henceforth, we restrict our discussion to the cases $`(\beta ,\alpha )`$$``$$`(1,1)`$, leaving a discussion of the $`\alpha `$$`=`$$`\beta `$$`=`$$`1`$ results to appendix B. Note that the norms in Eqs. (24)–(27), are neither Weyl, nor super-Weyl, invariant. Weyl and super-Weyl transformations generate variations orthogonal to the gauge slice defined by gamma traceless $`\chi `$. Consequently, both symmetries are anomalous, giving contributions to the path integral measure which will exponentiate to terms proportional to the supersymmetric Liouville action .
The result of a local supersymmetry transformation plus super-Weyl scaling on the expression in Eq. (19) is a variation outside the superconformal gauge slice: $`S_L`$ supersymmetrizes to the super-Liouville action and the super-Weyl anomaly exponentiates to a local renormalizable term proportional to the induced $`\gamma ^m(\delta \chi _m)`$. The integration over diffeomorphisms continuously connected to the identity gives the volume of the Weyl group, cancelled by the term in the denominator. Thus, under a local supersymmetry and super-Weyl transformation, the variation of Eq. (19) gives:
$$\frac{[d\delta \varphi ]}{\mathrm{Vol}(\mathrm{Weyl})}_0^{\mathrm{}}𝑑lJ_b(l;\widehat{g})e^{\frac{10d}{48\pi }S_{SL}[\varphi ,\zeta ,g]}(\mathrm{det}^{}\mathrm{\Delta })_{\widehat{g}}^5,$$
(29)
where $`S_{SL}[\varphi ,\zeta ,g]`$ is the super-Liouville action. In the fiducial cylinder metric,
$$S_{SL}[\varphi ,\zeta ,g]=d^2\sigma \sqrt{\widehat{g}}[\frac{1}{2}g^{mn}_m\varphi _n\varphi +\overline{\zeta }\gamma ^m_m\zeta +\mu _1e^\varphi +\mu _2\overline{\zeta }\zeta e^{\varphi /2}],$$
(30)
where $`\mu _1`$ and $`\mu _2`$ are induced violations of, respectively, the Weyl and super-Weyl invariance of the measure.
Likewise, consider the measure for the fermionic fields in the path integral. In the absence of a supermodulus, the effective action for the fermi fields given in Eq. (18) reduces to the free field action for the matter fermions alone: the world-sheet gravitino can be entirely gauged away. As with the gauge fermions, there are no constant modes for the matter fermion on a world-sheet of cylindrical topology when $`\alpha `$,$`\beta `$$``$$`1`$. Integrating over the $`\psi ^\mu `$ with norm,
$$|\delta \psi ^\mu |^2=id^2\sigma \sqrt{g}(\delta \overline{\psi }^\mu )(\delta \psi _\mu ),1=[(d\delta \psi ^{\mu +})(d\delta \psi ^\mu )]e^{{\scriptscriptstyle d^2\sigma \sqrt{g}(\delta \psi ^{\mu +})(\delta \psi _\mu ^{})}},$$
(31)
gives $`(\mathrm{det}\gamma ^m_m)_g^d`$, in the fiducial cylinder metric. Thus, for $`(\beta ,\alpha )(1,1)`$, we obtain:
$$\frac{[d\chi _m](\mathrm{det}\gamma ^m_m)_g^{10}}{\mathrm{Vol}(\mathrm{sDiff}\times \mathrm{sWeyl})}=\frac{[d\delta \xi ]}{\mathrm{Vol}(\mathrm{sDiff})}\frac{[d\delta \zeta ]}{\mathrm{Vol}(\mathrm{sWeyl})}J_f(l;\widehat{g})[\mathrm{det}(\gamma ^m_m)]_{\widehat{g}}^{10}e^{\frac{10d}{96\pi }S_{SL}[\varphi ,\zeta ,g]}.$$
(32)
Combining the expressions in Eqs. (30) and (32) we see that, in the critical spacetime dimension, the Liouville, and super-Liouville, modes entirely decouple. The Jacobian $`J_f`$ describes the change of variables from $`\delta \chi _m`$ to $`(\delta \xi ,\delta \zeta )`$ and was first computed in . In the absence of a supermodulus and superconformal Killing spinor on the world-sheet, the fermionic Jacobian is simply:
$$J_f=(\mathrm{det}P_{1/2}^{}P_{1/2})^{1/2},$$
(33)
where $`P_{1/2}^{}`$$`=`$$`_m`$ on the cylinder, is the adjoint of the operator $`P_{1/2}`$. The extension for $`\alpha `$$`=`$$`\beta `$$`=`$$`1`$ is discussed in appendix B. Combining with results from Eqs. (16), (20), we obtain the following expression for the amplitude:
$$𝒜=\frac{i}{2}V_{p+1}(4\pi ^2\alpha ^{})^{(p+1)/2}_0^{\mathrm{}}\frac{dl}{l}l^{(p+1)/2}e^{R^2l/4\pi \alpha ^{}}\eta (\frac{il}{2})^8\underset{(\beta ,\alpha )(1,1)}{}C_\alpha ^\beta [\mathrm{det}(\gamma ^m_m)]_{(\beta ,\alpha )}^8,$$
(34)
where the contribution from two left-moving and two right-moving component fermions— one each of which is a timelike fermion , has been cancelled against the fermionic Jacobian, $`J_f`$. We are assuming identical boundary condition, $`\beta `$, for all NN, and DD, fermions on a world-sheet with fixed spin structure $`\alpha `$.
The fermionic determinants in Eq. (34) can be computed using the method of zeta function regularization. The component world-sheet fermions are complexified into Weyl fermions, equivalently, using bosonization, into chiral bosons . An advantage is that the result can be readily generalized to the modified boundary conditions describing a pair of Dpbranes rotated relative to each other. Upon analytic continuation of a Euclidean embedding coordinate to Minkowskian time, this is equivalent to imposing boundary conditions describing parallel Dpbranes in relative motion. We begin by grouping the eight left-moving component fermions, $`\psi ^{+i}`$, $`i`$$`=`$$`1`$, $`\mathrm{}`$, $`8`$, into four left-moving Weyl fermions:
$$\psi ^{+1}+i\psi ^{+2},\psi ^{+3}+i\psi ^{+4},\psi ^{+5}+i\psi ^{+6},\psi ^{+7}+i\psi ^{+8}.$$
(35)
Likewise, complexifying eight right-moving component fermions, $`\psi ^i`$, gives four right-moving Weyl fermions. This is the world-sheet fermion content of the fermionic string with both left and right moving $`N`$$`=`$$`1`$ world sheet supersymmetries. The open string boundary condition in Eq. (14) reduces the number of independent fermionic degrees of freedom by half, since it relates corresponding left and right moving Weyl fermions at the end-points. Thus, while the $`4`$$`+`$$`4`$ Weyl fermions are not all independent world-sheet fields, this is a convenient basis in which to express the result. We will work with the type I theory in its T-dual formulation with a type IA, or type IB, supersymmetry, as determined by Dpbranes with even, or odd, $`p`$$``$$`9`$. For convenience, we will keep the full $`SO(8)`$ global symmetry of the transverse coordinates, imposing identical boundary conditions and spin structure on all four independent Weyl fermions on the world-sheet. The required chiral determinants can all be obtained from the single functional determinant:
$$\mathrm{det}^{}\mathrm{\Delta }^{(\beta ,\alpha )}=\underset{n_1,n_2}{\overset{}{}}(\frac{4\pi ^2}{l^2})[(n_1+(\alpha 1)/2)^2+\frac{l^2}{4}(n_2+(\beta 1)/2)^2)].$$
(36)
computed by the method given in . The chiral determinants are formally defined by taking a square root. The result is therefore ambiguous up to a phase. We retain the phase ambiguity, absorbing it in the unknown $`C_\alpha ^\beta `$. For anti-periodic (periodic) Weyl fermions with $`\beta `$$`=`$$`0(1)`$, and world-sheets with spin structure, $`\alpha `$, we obtain the result :
$$\mathrm{det}^{}\mathrm{\Delta }^{(\beta ,\alpha )}=\left[q^{{\scriptscriptstyle \frac{1}{24}}+{\scriptscriptstyle \frac{\beta ^2}{8}}}\underset{m=1}{\overset{\mathrm{}}{}}(1+e^{\pi i\alpha }q^{m{\scriptscriptstyle \frac{1}{2}}(1\beta )})(1+e^{\pi i\alpha }q^{m{\scriptscriptstyle \frac{1}{2}}(1+\beta )})\right]\frac{\mathrm{\Theta }_{(\beta ,\alpha )}(0,\frac{il}{2})}{\eta (\frac{il}{2})}.$$
(37)
As is shown in appendix B, in the case $`\alpha `$$`=`$$`\beta `$$`=`$$`1`$, the left hand side of Eq. (37) is corrected by a contribution from the measure which vanishes for boundaries on static Dbranes. Inserting the expressions for $`(\beta ,\alpha )`$$``$$`(1,1)`$ in Eq. (34) gives :
$$𝒜=\frac{i}{2}V_{p+1}(8\pi ^2\alpha ^{})^{(p+1)/2}_0^{\mathrm{}}\frac{dl}{l}e^{R^2l/2\pi \alpha ^{}}l^{(p+1)/2}\eta (il)^{12}\left[C_0^0\mathrm{\Theta }_{(0,0)}^4(0,il)+C_1^0\mathrm{\Theta }_{(0,1)}^4(0,il)+C_0^1\mathrm{\Theta }_{(1,0)}^4(0,il)\right],$$
(38)
where the $`C_\alpha ^\beta `$ are undetermined phases. Note that we have rescaled by a factor $`l`$$``$$`2l`$ in writing Eq. (38).
### B Infrared Consistency of the Supersymmetric Annulus
We now come to the interesting issue of determining the phase of the path integral. The discussion that follows is based on ideas taken from and also . We will show that the following infrared consistency conditions:
* the elimination of the tachyon
* the absence of a static force between the Dbranes
determine two of the three unknown phases in the expression for the supersymmetric annulus given in Eq. (38). It should be emphasized at the outset that these requirements are insufficient to ensure infrared finiteness of the perturbative string theory . The reason is that the oriented open and closed supersymmetric string has tadpoles for both the dilaton and the Ramond-Ramond scalar fields which are cancelled by contributions to the vacuum amplitude from non-orientable world-sheets in the full unoriented type $`\mathrm{I}^{}`$ string theory . More generally, the presence of classical sources in the generic String/M theory background akin to the orientifold planes of the unoriented string may provide a means to cancel the troublesome tadpoles so we will leave this option open. Tadpole cancellation is an essential requirement for an infrared finite theory . This is sometimes phrased as the requirement of BPS charge conservation in a compact space. We note that we are taking the point of view that an acceptable theory should have a sensible definition both in a noncompact, and a compactified, spacetime. In the compact space, the flux lines of the associated background field are required to close on a configuration of classical Ramond-Ramond sources. While we believe it likely that the infrared consistency conditions on the supersymmetric annulus amplitude listed above are necessary conditions which must be met by the generic stable background of String/M theory— irrespective of whether this is a background described by orientable or non-orientable world-sheets, we should note here the recent work on unstable brane configurations , and on the stabilization of the tachyon in open string field theory . Note also that we have distinguished the BPS condition—the absence of a static force between BPS sources in this example, from the specification of the spacetime supersymmetry. Each of these criteria impacts distinct renormalizability properties of the string theory .
The third phase in the annulus amplitude, which we choose to be $`C_0^0`$, will be determined by computing the static limit of the long distance potential between Dbranes. Define the Minkowskian potential:
$$𝒜(r,u)=i\underset{T\mathrm{}}{lim}_T^{+T}𝑑\tau V[r(\tau ),u],$$
(39)
where $`r^2`$$`=`$$`R^2`$$`+u^2\tau ^2`$, and take the static limit $`u`$$`=`$$`0`$ . This is a special case of the calculation that follows in section IIID. Due to the no-force condition for the BPS configuration of static and parallel Dbranes the amplitude vanishes. However, if we isolate the contribution from the $`(\beta ,\alpha )`$$`=`$$`(0,0)`$ sector alone, we can extract the static long-range Newtonian gravitational potential between the Dbranes. This is required to be attractive, which determines the phase $`C_0^0`$. Thus, simple, and universal, infrared consistency conditions determine all of the unknown phases in the supersymmetric annulus.
Let us carry out this procedure explicitly. The long distance physics of the vacuum amplitude is dominated by the lowest lying closed string modes. This can be exposed more clearly by taking the $`l`$$``$$`0`$ limit of the expression in Eq. (38). Using standard identities for the Jacobi theta functions gives:
$`𝒜(r)=`$ $`{\displaystyle \frac{i}{2}}V_{p+1}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dl}{l}}(8\pi ^2\alpha ^{})^{(p+1)/2}e^{r^2l/2\pi \alpha ^{}}l^{(7p)/2}\eta ({\displaystyle \frac{i}{l}})^{12}{\displaystyle \underset{(\beta ,\alpha )(1,1)}{}}C_\alpha ^\beta \mathrm{\Theta }_{(\alpha ,\beta )}^4(0,{\displaystyle \frac{i}{l}})`$ (40)
$`=`$ $`{\displaystyle \frac{i}{2}}V_{p+1}(8\pi ^2\alpha ^{})^{(p+1)/2}{\displaystyle _0^{\mathrm{}}}𝑑le^{r^2l/2\pi \alpha ^{}}l^{(5p)/2}q^{1/2}`$ (42)
$`\times \left\{q^0(C_0^0+C_0^1)+q^{1/2}[8(C_0^0C_0^1)+16C_1^0]+O(q)\right\},`$
where $`q`$$`=`$$`e^{2\pi /l}`$. The absence of the closed string tachyon appearing at $`O(q^{1/2})`$ requires $`C_0^1`$$`=`$$`C_0^0`$. The absence of a static force between the Dbranes requires that we set $`C_1^0`$$`=`$$`C_0^0`$. The 8+8 massless states in the $`(0,0)`$, $`(0,1)`$ sector contribute to the vacuum amplitude with the opposite spacetime statistics of the 16 massless states in the $`(1,0)`$ sector. Thus, we discover an underlying spacetime supersymmetry in the vacuum amplitude although we have not required it. While the absence of the static force and the requirement of spacetime supersymmetry are equivalent conditions in this example— the T-dualized type I string in a background of parallel and static Dpbranes, the distinction may be of consequence for non-BPS brane configurations.
It remains to determine the overall phase of the annulus. The massless states in the Neveu-Schwarz sector with spin structure, $`(\beta ,\alpha )`$$`=`$$`(0,0)`$, contribute a static long range Newtonian interaction between the Dbranes, analogous to that in the bosonic string . It should be emphasized that this term is universally present in the one-loop vacuum amplitude of any fermionic string theory irrespective of background, prior to possible cancellation by additional contributions to the amplitude. The sign of the potential is determined by the phase $`C_0^0`$. Defining the Minkowskian potential as in Eq. (39) and substituting for the phases in Eq. (42) gives:
$`V_{\mathrm{static}}^{(0,0)}(r)`$ $`=C_0^02V_p(8\pi ^2\alpha ^{})^{(p+1)/2}\frac{1}{2}{\displaystyle _0^{\mathrm{}}}𝑑le^{r^2l/2\pi \alpha ^{}}l^{(5p)/2}[8+O(q^{1/2})]`$ (44)
$`=8C_0^0V_p(8\pi ^2\alpha ^{})^{(p+1)/2}{\displaystyle \frac{1}{r^{7p}}}\mathrm{\Gamma }({\displaystyle \frac{7p}{2}})(2\pi \alpha ^{})^{(7p)/2}+\mathrm{}.`$
The factor of eight counts the transverse polarizations of the $`\mathrm{𝟖}_𝐯`$ multiplet under the $`SO(9,1)`$ Lorentz group. The potential has a static remnant upon setting $`u`$$`=`$$`0`$. The Newtonian potential is required to be attractive, and we can therefore set $`C_0^0`$$`=`$$`1`$. Thus,
$$V_{\mathrm{static}}^{(0,0)}(R)=(d2)\frac{1}{R^{7p}}V_p2^{22p}\pi ^{(53p)/2}\alpha _{}^{}{}_{}{}^{3p}\mathrm{\Gamma }(\frac{7p}{2}),$$
(45)
where $`R`$ is the static separation of the Dpbranes. We emphasize once again that the static interaction in Eq. (45) will be cancelled by contributions to the vacuum amplitude from states in the Ramond sector and, in the full string theory, from the unoriented world-sheets. Nevertheless, it has a simple physical interpretation which allows us to use it to determine an unknown phase in the string path integral.
We have shown that infrared consistency conditions determine all three phases in Eq. (38), $`C_0^0`$$`=`$$`1`$, $`C_1^0`$$`=`$$`C_0^1`$$`=`$$`1`$, giving an unambiguous expression for the supersymmetric annulus. We emphasize that the ambiguity in phase has been determined by requiring consistency with known qualitative features of the long distance physics. On the other hand, the normalization of the string path integral is unambiguously determined leading to the prediction of the numerical coefficient in Eq. (45). We will use these expressions in section IV to make predictions about the short distance physics.
### C Pair Correlation Function of Wilson Loops
In , we gave a path integral prescription for the pair correlation function of macroscopic loop observables, $`M(C_i)`$, $`M(C_f)`$ in the weakly coupled bosonic string theory, following the earlier work of Cohen et al . The loops, $`C_i`$, $`C_f`$, are taken to lie in a flat D-manifold, $`𝒟`$, and are identified with the closed world-lines of heavy point sources in the gauge theory. The result for the short distance potential is independent of whether $`𝒟`$ is the worldvolume of some higher dimensional Dpbrane, the intersection of Dpbranes, or the bulk transverse space orthogonal to some configuration of Dbranes. The key issue is the implementation of boundary reparameterization invariance in the covariant one-loop string path integral. Our interest is in the large loop length limit where the dynamics should be universal, independent of the detailed geometrical parameters of the loops. In , we point out that the large loop length dynamics for generic loops is captured rather simply by summing over reparameterizations of loops with one or more marked points. For such maps the sum over reparameterizations of the boundary, $`M`$, can be easily implemented in closed form even prior to taking the large loop length limit . This gives a well-defined framework for computing the boundary reparameterization invariant pair correlation function which also preserves its normalization : this is the key ingredient that enables a numerical prediction for the short distance potential between heavy sources in the gauge theory arising in fluctuations of the vacuum energy density.
Following Cohen et al , the tree correlation function for a pair of macroscopic string loops is represented as a path integral over embeddings and metrics on world-sheets of cylindrical topology terminating on fixed curves, $`C_i`$, $`C_f`$, within the spacetime $`𝒟`$, weighted by the locally supersymmetric action given in Eq. (A14):
$$<M(C_i)M(C_f)>=\frac{1}{2}\underset{(\beta ,\alpha )(1,1)}{}C_\alpha ^\beta _{[C_i,C_f][\beta ,\alpha ]}\frac{[de][d\chi ]}{\mathrm{Vol}(gauge)_M}\frac{[dX][d\psi ][dg][d\chi ]}{\mathrm{Vol}(gauge)_M}e^{S_{SP}[X,\psi ,g,\chi ]\mu _0_{}d^2\sigma \sqrt{g}S_{\mathrm{ren}.}}.$$
(46)
As has been emphasized in section IIIA-B, we have taken all of the NN, DD, world-sheet fermions to satisfy identical boundary conditions at the end-points $`\sigma ^2`$$`=`$$`0`$, $`1`$, for world-sheets of specified spin structure. We decouple bulk and boundary deformations of the world-sheet fields as in , imposing Dirichlet boundary conditions on all of the spatial embedding coordinates, $`X^\mu `$, $`\mu `$$`=`$$`1`$, $`\mathrm{}`$, $`9`$. Then the boundary conditions on the matter fermions, $`\psi ^\mu `$, are given by Eq. (14). The distinction between the Wilson loop correlation function and the ordinary annulus amplitude computed in section III comes from the inclusion of fluctuations in the world-sheet metric (einbein) on the boundaries of the annulus . We gauge both boundary diffeomorphisms and local supersymmetry variations on the boundary. Anomalies of the measure under Weyl, and super-Weyl, transformations will, as before, be exponentiated as terms in the effective action proportional to the super-Liouville theory with boundary terms included . The quantization of the super-Liouville fields could be performed along the lines of , and citations thereof, but we will be interested in the vacuum amplitude in the critical spacetime dimension where the super-Liouville theory decouples.
Consider the measure for einbeins. As was shown in , a reparameterization $`\delta f(\sigma ^1)`$, tangential to the boundary induces a non-trivial boundary Jacobian computed in . However, a Weyl rescaling of the einbein can always be absorbed in a shift in the Liouville field on the boundary. Consider the variation in the measure for einbeins under a local supersymmetry transformation. From Eqs. (A18)–(A28) of appendix A, we see that the variation in the einbein under a supersymmetry transformation can always be absorbed in a rescaling of the super-Liouville fields, $`(\varphi ,\zeta )`$, on the boundary:
$$\delta _Se=2e_a^m(\delta _Se_m^a)=2(\delta \overline{\xi })(\gamma ^m\chi _m),\delta _We=2\delta \lambda ,$$
(47)
where $`\delta _S`$, $`\delta _W`$, respectively denote the variation under local supersymmetry and Weyl transformations. Likewise, consider the variations in the gravitino on the boundary. In superconformal gauge, setting $`\gamma ^m\chi _m`$$`=`$$`0`$, there are no independent variations of the gravitino on the boundary that have not already been accounted for in the analysis in section IIIB: a variation in $`\chi `$ on the boundary is a departure from superconformal gauge. The resulting super-Weyl anomaly is absorbed in the super-Liouville dynamics. Thus, the sum over boundary deformations of the gravitino in Eq. (46) is pure gauge. We eliminate the contributions to the measure from the zero modes of the Neumann embedding coordinates in the expression for the supersymmetric annulus derived in Eq. (38). The result for the pair correlation function of Wilson loops in the T-dualized type I theory is a remarkably simple extension of the bosonic analysis given in :
$$<M(C_i)M(C_f)>=_0^{\mathrm{}}𝑑le^{S_{\mathrm{saddle}}[\overline{x},\widehat{g}]}\eta (il)^{12}\underset{(\beta ,\alpha )(1,1)}{}C_\alpha ^\beta \mathrm{\Theta }_{(\beta ,\alpha )}^4(0,il).$$
(48)
The annulus terminates on fixed curves, $`C_i`$, $`C_f`$, with fixed spatial separation $`R`$ in some generic D-manifold, i.e., within the world-volume of a Dpbrane, in the intersection of the world-volumes of two or more Dbranes, or in the bulk space transverse to some configuration of branes. The path integral computes quantum fluctuations about a saddle world-sheet configuration stretched between the loops, $`C_i`$, $`C_f`$, and the saddle-point action can be computed as in . In we focussed on the simplest configurations of coplanar loops which can capture the universal features of the large loop length limit which determines the short distance potential between heavy point sources in the gauge theory. For such a configuration, we showed in that the dominant contribution to the saddle action in the small $`R`$, large $`l`$ limit takes the form, $`S_{\mathrm{saddle}}`$$``$$`R^2l/2\pi \alpha ^{}`$, independent of the shape or other geometrical characteristics of the loops.
The reader may wonder if generalizations of this result with new non-trivial fermionic degrees of freedom on the boundary are possible. The answer lies in our understanding of the global structure of supermoduli space. The analysis given above is appropriate for the superconformal gauge fixed perturbative superstring in a Dbrane background. However, since brane dynamics is as yet a poorly understood subject in the wider context of String/M theory, it may be that the global structure of supermoduli space can play an interesting role in the full nonperturbative theory .
### D Generic Boundary Conditions on the Annulus
Consider the supersymmetric annulus derived in Eq. (38) with boundaries on a pair of static and parallel Dpbranes. In this section, we sketch the modifications to Eqs. (38) for a configuration of rotated Dpbranes or, by an analytic continuation of a Neumann coordinate, for a pair of Dpbranes in relative motion, previously derived in . It is convenient to complexify the coordinates, $`X^\mu `$, $`\mu `$$`=`$$`1`$, $`\mathrm{}`$, $`8`$, in pairs, decomposing a generic rotation into independent rotations in the four planes, $`(1,2)`$, $`(3,4)`$, $`(5,6)`$, $`(7,8)`$. The $`(1,1)`$ fermionic path integral no longer vanishes for generic boundary conditions on the fermions. For the generic rotation in all four planes, all four sectors of the fermionic path integral, $`\beta ,\alpha `$$``$$`0,1`$, contribute to the one-loop vacuum amplitude . This case is discussed in appendix B. We will consider here the simpler case of rotation in a single plane. The unknown phases in the annulus amplitude are determined by an extension of the infrared consistency conditions described in section IIIB for configurations of moving Dbranes. As a consistency check, we verify that we recover the long distance velocity dependent potential between Dpbranes in relative motion, including the numerical coefficient previously computed in . This result will be adapted in section IV to obtain the supersymmetric pair correlation function of Wilson loops corresponding to closed world-lines of heavy sources in a gauge theory in relative slow motion.
Consider motion in the plane $`X^0`$$`=`$$`iX^2`$, $`X^1`$. The functional determinant for a complex scalar satisfying the V boundary conditions given in Eq. (10) can be obtained using zeta function regularization . Likewise for the corresponding Weyl fermion; the fermionic functional determinant with spin structure $`\alpha `$ and boundary condition $`\beta `$ is simply the expression given in Eq. (37) with nonvanishing argument for the Jacobi theta functions: $`\mathrm{\Theta }_{(\beta ,\alpha )}(\nu ,\frac{il}{2})`$, with $`\nu `$$`=ul/2\pi `$ . Thus,
$$q^{E_0(\beta ,u)}\underset{m=1}{\overset{\mathrm{}}{}}(1+zq^{m{\scriptscriptstyle \frac{1}{2}}(1\beta )})(1+z^1q^{m{\scriptscriptstyle \frac{1}{2}}(1+\beta )})=\left[\frac{e^{u^2l/2\pi }\mathrm{\Theta }_{(\beta ,\alpha )}(ul/2\pi ,\frac{il}{2})}{\eta (\frac{il}{2})}\right],$$
(49)
where the parameter $`z`$$`=`$$`e^{i\pi (\alpha +ul/\pi )}`$. The vacuum energy of the Weyl fermion with the V boundary condition is $`E_0`$$`=`$$`\frac{1}{24}`$$`+`$$`\frac{1}{2}(\frac{iu}{\pi }`$$`+`$$`\frac{\beta }{2})^2`$. Thus, the one-loop vacuum amplitude with boundaries on Dpbranes in relative motion takes the form :
$$𝒜(r,u)=\frac{1}{2}V_p_0^{\mathrm{}}\frac{dl}{l}(4\pi ^2\alpha ^{}l)^{p/2}\frac{e^{r^2l/4\pi \alpha ^{}}\eta (\frac{il}{2})^9}{i\mathrm{\Theta }_{11}(ul/2\pi ,\frac{il}{2})}\underset{(\beta ,\alpha )(1,1)}{}C_\alpha ^\beta \mathrm{\Theta }_{(\beta ,\alpha )}^3(0,\frac{il}{2})\mathrm{\Theta }_{(\beta ,\alpha )}(ul/2\pi ,\frac{il}{2}).$$
(50)
It is convenient to rescale $`l`$$``$$`2l`$ in the final result.
The phases in the vacuum amplitude can be determined by infrared consistency conditions. As in section IIIB, we require the absence of a tachyon and the vanishing of the static force between the branes. At long distances, we must recover at order $`v^4`$ the attractive $`1/r^{7p}`$ potential between BPS sources required from matching to the low energy type II supergravity theory. The long distance physics of the vacuum amplitude is dominated by the lowest lying closed string modes, exposed by taking the $`l`$$``$$`0`$ limit of the expression in Eq. (50). Using standard identities for the Jacobi theta functions we can write:
$`𝒜(r,u)=`$ $`{\displaystyle \frac{1}{2}}V_p{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dl}{l}}(8\pi ^2\alpha ^{})^{p/2}{\displaystyle \frac{e^{r^2l/2\pi \alpha ^{}}l^{(6p)/2}}{\eta (\frac{i}{l})^9\mathrm{\Theta }_{11}(\frac{iu}{\pi },\frac{i}{l})}}{\displaystyle \underset{(\beta ,\alpha )(1,1)}{}}C_\alpha ^\beta \mathrm{\Theta }_{(\alpha ,\beta )}^3(0,{\displaystyle \frac{i}{l}})\mathrm{\Theta }_{(\alpha ,\beta )}({\displaystyle \frac{iu}{\pi }},{\displaystyle \frac{i}{l}})`$ (51)
$`=`$ $`{\displaystyle \frac{1}{2}}V_p(8\pi ^2\alpha ^{})^{p/2}{\displaystyle _0^{\mathrm{}}}𝑑le^{r^2l/2\pi \alpha ^{}}{\displaystyle \frac{l^{(4p)/2}q^{1/2}}{2\mathrm{S}\mathrm{i}\mathrm{n}(iu)}}`$ (53)
$`\times \left\{q^0(C_0^0+C_0^1)+q^{1/2}[(2\mathrm{C}\mathrm{o}\mathrm{s}(2iu)+6)(C_0^0C_0^1)+16C_1^0\mathrm{Cos}(iu)]+O(q)\right\}.`$
The absence of the closed string tachyon appearing at $`O(q^{1/2})`$ requires $`C_0^1`$$`=`$$`C_0^0`$. With this choice, a small $`u`$ expansion of the $`O(q^0)`$ terms gives:
$$C_0^0(16+8u^2+\frac{8}{3}u^4+O(u^6))+C_1^0(16+8u^2+\frac{2}{3}u^4+O(u^6)).$$
(54)
The absence of a static force between the Dbranes requires that we set $`C_1^0`$$`=`$$`C_0^0`$. As a consequence the leading contribution to the vacuum amplitude is $`O(u^4)`$— at which order, spacetime supersymmetry is broken. The non-vanishing coefficient implies the existence of a long range velocity dependent potential between the Dbranes . The sign of the potential is determined once we fix the phase $`C_0^0`$. Defining the Minkowskian potential as in Eq. (39) and substituting for the phases in Eq. (53) gives:
$`V_{\mathrm{long}}(r,u)`$ $`=C_0^02V_p(8\pi ^2\alpha ^{})^{(p+1)/2}{\displaystyle \frac{1}{2}}{\displaystyle _0^{\mathrm{}}}𝑑le^{r^2l/2\pi \alpha ^{}}l^{(5p)/2}{\displaystyle \frac{\mathrm{tanh}(u)}{2i\mathrm{Sin}(iu)}}[2u^4+O(u^6)]`$ (56)
$`=u^4C_0^0V_p(8\pi ^2\alpha ^{})^{(p+1)/2}{\displaystyle \frac{1}{r^{7p}}}\mathrm{\Gamma }({\displaystyle \frac{7p}{2}})(2\pi \alpha ^{})^{(7p)/2}+O(u^6).`$
The potential is required to be attractive which implies $`C_0^0`$$`=`$$`1`$. Thus, we recover the potential between Dpbranes including the numerical coefficient previously computed in :
$$V_{\mathrm{long}}(r,u)=\frac{u^4}{r^{7p}}V_p2^{22p}\pi ^{(53p)/2}\alpha _{}^{}{}_{}{}^{3p}\mathrm{\Gamma }(\frac{7p}{2})+O(u^6).$$
(57)
Thus, all three phases in Eq. (50) are determined, $`C_0^0`$$`=`$$`1`$, $`C_0^1`$$`=`$$`C_1^0`$$`=`$$`1`$, and we have an unambiguous expression for the amplitude.
An extension of these arguments can be applied to more complicated non-BPS brane configurations. Brane configurations which break one-half (one-quarter) of the spacetime supersymmetries can be distinguished by requiring that an order $`v^2`$ velocity dependent force is respectively absent (present). Similar arguments apply to configurations of mixed, intersecting, or rotated Dbranes: from the low energy correspondence to supergravity, we infer the qualitative form of the long distance potential. This is then applied as an infrared consistency condition on the unknown phases in the supersymmetric annulus.
## IV Short Distance Potential Between Heavy Sources
In an earlier work , we showed that there exists a short distance interaction between heavy sources in a gauge theory traversing fixed spacetime paths in some generic background of the bosonic string. The potential arises in fluctuations in the vacuum energy density. The same phenomenon can be observed in the supersymmetric T-dualized type I string theory. We will obtain in this section an expression for the short distance potential between heavy sources in a supersymmetric gauge theory. Our results are derived in a systematic small velocity short distance double expansion, following an analogous treatment of the short distance potential in bosonic string theory in . The potential is extracted from the supersymmetric pair correlation function of Wilson loops, in the limit of large loop lengths and small spatial separations.
We begin with the pair correlation function of Wilson loops corresponding to world-lines of heavy sources in relative collinear motion with nonrelativistic velocity $`v`$ in the direction $`X^1`$. For coplanar loops, this mimics straight line trajectories in the Euclideanized $`X^0`$$`=iX^2`$, $`X^1`$ plane: $`r^2`$$`=`$$`R^2`$$`+`$$`v^2\tau ^2`$, for small separations $`r`$. The modifications to the expression given in Eq. (48) for these boundary conditions are straightforward, following the results of section IIID. In the large loop length limit, we have:
$$<M(C_i)M(C_f)>=_0^{\mathrm{}}𝑑l\frac{e^{r^2l/2\pi \alpha ^{}}\eta (il)^9}{i\mathrm{\Theta }_{11}(ul/\pi ,il)}\underset{(\beta ,\alpha )(1,1)}{}C_\alpha ^\beta \mathrm{\Theta }_{(\beta ,\alpha )}^3(0,il)\mathrm{\Theta }_{(\beta ,\alpha )}(ul/\pi ,il).$$
(58)
The short distance dynamics is dominated by the lowest lying modes in the open string spectrum. We define the Minkowskian potential :
$$<M(C_i)M(C_f)>=i\underset{T\mathrm{}}{lim}_T^{+T}𝑑\tau V_{\mathrm{eff}.}[r(\tau ),u].$$
(59)
The short distance regime is exposed by expanding the integrand in Eq. (58) in powers of $`q`$$`=`$$`e^{2\pi l}`$. Substituting for the phases, we can express the Minkowskian potential at short distance in a small $`q`$ expansion as in :
$`V_{\mathrm{eff}.}(r,u)=`$ $`2(8\pi ^2\alpha ^{})^{1/2}{\displaystyle _0^{\mathrm{}}}𝑑l\mathrm{tanh}(u)l^{1/2}{\displaystyle \frac{e^{r^2l/2\pi \alpha ^{}}\eta (il)^9}{\mathrm{\Theta }_{11}(ul/\pi ,il)}}{\displaystyle \underset{(\beta ,\alpha )(1,1)}{}}C_\alpha ^\beta \mathrm{\Theta }_{(\beta ,\alpha )}^3(0,il)\mathrm{\Theta }_{(\beta ,\alpha )}(ul/\pi ,il)`$ (60)
$`=`$ $`(8\pi ^2\alpha ^{})^{1/2}{\displaystyle _0^{\mathrm{}}}𝑑le^{r^2l/2\pi \alpha ^{}}{\displaystyle \frac{l^{1/2}\mathrm{tanh}(u)}{q^{1/2}\mathrm{Sin}(ul)}}`$ (62)
$`\times \left[(1+2q^{1/2})^3(1+2q^{1/2}\mathrm{Cos}(2ul))(12q^{1/2})^3(12q^{1/2}\mathrm{Cos}(2ul))16q^{1/2}\mathrm{Cos}(ul)+O(q)\right].`$
The leading non-vanishing terms in this expression, due to massless exchange, contribute at order $`q^0`$:
$$V(r,u)=(8\pi ^2\alpha ^{})^{1/2}_0^{\mathrm{}}𝑑le^{r^2l/2\pi \alpha ^{}}\frac{l^{1/2}\mathrm{tanh}(u)}{\mathrm{Sin}(ul)}\left[12+4\mathrm{C}\mathrm{o}\mathrm{s}(2ul)16\mathrm{C}\mathrm{o}\mathrm{s}(ul)\right].$$
(63)
We have assumed small velocities $`v`$$`=`$$`\mathrm{tanh}(u)`$$``$$`u`$. We now perform a resummation of the integrand in the variables $`r`$,$`u`$, as in . The regime of validity is determined by the behavior of the cosecant function. We perform a Taylor expansion in the first half-period of its argument, $`0`$$``$$`ul`$$`<`$$`\pi `$. For sufficiently small $`u`$ values the oscillations in the integrand are increasingly rapid, smearing out the integral . As in the analysis of the vacuum amplitude for the bosonic string in , we note that the contribution from the integration domain $`ul`$$``$$`\pi `$ can always be bounded, or evaluated by numerical integration, as a self-consistency check on the approximation. This check provides an upper limit, $`u_+`$, on the permissible velocities. With this restriction, the contribution from the domain $`l`$$`>`$$`\pi /u_+`$ can be dropped and we suppress it in what follows. The potential takes the form:
$$V(r,u)=(8\pi ^2\alpha ^{})^{1/2}_0^{\pi /u_+}𝑑le^{r^2l/2\pi \alpha ^{}}l^{1/2}\mathrm{tanh}(u)/u\left[\underset{k=1}{\overset{\mathrm{}}{}}C_k(ul)^{2k}+\underset{k=1}{\overset{\mathrm{}}{}}\underset{m=1}{\overset{\mathrm{}}{}}C_{k,m}(ul)^{2(k+m)}\right],$$
(64)
where the coefficients in the summation are given by:
$`C_k=`$ $`{\displaystyle \frac{4(1)^k(2^{2k}4)}{(2k)!}}`$ (65)
$`C_{k,m}=`$ $`{\displaystyle \frac{8(1)^k(2^{2m1}1)}{(2k)!(2m)!}}|B_{2m}|(2^{2k}4).`$ (66)
The $`B_{2m}`$ are the Bernoulli numbers. Note that the $`k`$$`=`$$`1`$ term vanishes in both sums and the leading velocity dependence of the amplitude is $`O(u^4)`$. Integrating over $`l`$ gives a systematic expansion for the potential in powers of $`u^2/r^4`$. As in , we identify a dimensionless scaling variable, $`z`$$`=`$$`r_{\mathrm{min}.}^2/r^2`$, where $`r_{\mathrm{min}.}^2`$$`=`$$`2\pi \alpha ^{}u`$. The velocity dependent corrections to the potential between heavy sources in the supersymmetric gauge theory are succinctly expressed as a convergent power series in the single dimensionless variable $`z`$:
$`V(r,u)=`$ $`(8\pi ^2\alpha ^{})^{1/2}\mathrm{tanh}(u)/ur^1(2\pi \alpha ^{})^{1/2}[{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}C_kz^{2k}\gamma (2k+1/2,\pi /z)`$ (68)
$`+{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}C_{k,m}z^{2(k+m)}\gamma (2(k+m)+1/2,\pi /z)].`$
Note that the potential takes the form of a scale invariant $`1/r`$ fall-off, contributed by the bosonic degrees of freedom , multiplicatively corrected by a convergent power series in $`z`$:
$$V(r,u)=\frac{\mathrm{tanh}(u)/u}{\mathrm{\Gamma }(\frac{1}{2})}\frac{1}{r}\left[z^4\gamma (\frac{9}{2},\pi /z)+O(z^6)\right],$$
(69)
indicative of its origin in fluctuations of the vacuum energy density. Recall that the regime of validity for the double expansion in small velocities and short distances is $`z`$$`<`$$`1`$, $`u`$$`<`$$`u_+`$. The scale factor $`z`$ determines the magnitude of the velocity dependent corrections and, therefore, the accuracy of the expansion. For a given accuracy, with fixed $`z`$ value, we can probe arbitrarily short distances $`r`$ by simultaneously adjusting the velocity. The power series corrections in the superstring are qualitatively similar, but much simpler than the analogous series in the bosonic string theory : the analogous result in a nonsupersymmetric gauge theory receives corrections in the variables $`z^2`$, $`uz/\pi `$, and $`u^2`$. The leading term in the potential between heavy sources in a supersymmetric gauge theory is therefore $`O(u^4/r^9)`$:
$$V(r,u)=\frac{u^4}{r^9}2^4\pi ^{7/2}\alpha _{}^{}{}_{}{}^{4}\mathrm{\Gamma }(\frac{9}{2})+O(u^6).$$
(70)
## V Conclusions
We have given a derivation from first principles of both the normalization and the phase of the supersymmetric annulus in the generic flat D-manifold background of the type I and type II string theories. The normalization of the string path integral is determined by its symmetries . As a consequence, one can extract numerical predictions from one-loop string amplitudes, free from any dependence on the string coupling constant. We have shown in this paper that phase ambiguities in the fermionic string path integral can be eliminated by the imposition of simple, and universal, infrared consistency conditions on qualitative features of the long distance physics, by matching to an appropriate supergravity theory. This prescription gives the same result as the usual GSO projection in the superstring, but has the hope of wider applicability to generic backgrounds of String/M theory. We note that we have emphasized the BPS conditions over supersymmetry. This is in keeping with the broader goal of understanding the self-Duality of String/M theory in generic backgrounds for the Ramond-Ramond antisymmetric tensor fields associated with Dbrane charges . The preliminary results given here need to be explored in a wider context, extended to an understanding of the consistency conditions that suffice to ensure tadpole cancellation in generic backgrounds of String/M theory.
Extending our earlier results for the bosonic string , we have shown that heavy point sources in a supersymmetric gauge theory in slow relative motion have an attractive, and velocity dependent, interaction at short distances. The potential can be expressed as a convergent power series in the single dimensionless variable $`z`$$`=`$$`r_{\mathrm{min}.}^2/r^2`$, where $`r_{\mathrm{min}.}^2`$$`=`$$`2\pi \alpha ^{}u`$ is the minimum distance probed in this approximation, valid for small velocities and short distances in the regime $`2\pi \alpha ^{}u`$$`<<`$$`r^2`$$`<<`$$`2\pi \alpha ^{}`$. It would be gratifying if this result could be exploited as a window into the short distance physics of String/M theory.
Acknowledgments
This work is supported in part by the National Science Foundation grant NSF-PHY-97-22394.
## A Spinor Conventions and Local Symmetries
In this appendix we establish our conventions, simultaneously deriving the Brink-Di Vecchia-Howe-Deser-Zumino action for the fermionic string with Minkowskian signature metric . The classical action is invariant under both reparameterizations and $`N`$$`=`$$`1`$ world-sheet supersymmetry transformations. The Euclidean action used in the string path integral is obtained by an analytic continuation. We consider free massless spinor fields in a two dimensional Minkowskian space parameterized $`(t,x)`$ with metric:
$$\{\gamma ^\mu ,\gamma ^\nu \}=2\eta ^{\mu \nu },𝒥^{10}=\frac{i}{4}[\gamma ^0,\gamma ^1]=\frac{i}{2}\gamma .$$
(A1)
$`𝒥^{10}`$ is the sole generator of Lorentz transformations in two dimensions and the matrix $`\gamma `$ projects onto spinors of definite chirality. We choose a representation with real gamma matrices. Thus,
$$\gamma ^0=\left(\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 1& \hfill 0\end{array}\right),\gamma ^1=\left(\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 1& \hfill 0\end{array}\right),\gamma =\gamma ^0\gamma ^1=\left(\begin{array}{cc}\hfill 1& \hfill 0\\ \hfill 0& \hfill 1\end{array}\right).$$
(A2)
Fermion bilinears that transform as scalars under Lorentz transformations are obtained by identifying a matrix $`\beta `$ such that:
$$\beta (D(\mathrm{\Lambda }))^{}\beta =D(\mathrm{\Lambda })^1,\beta (𝒥^{\mu \nu })^{}\beta =𝒥^{\mu \nu },\mathrm{where}D(\mathrm{\Lambda })e^{{\scriptscriptstyle \frac{i}{2}}(\omega _{\mu \nu }𝒥^{\mu \nu })}.$$
(A3)
Note that with this choice of gamma matrices the Lorentz generator acts non-unitarily in the spinor representation: $`(D(\mathrm{\Lambda }))^{}`$$``$$`D(\mathrm{\Lambda })^1`$. Thus, $`\beta `$ must be chosen to satisfy the conditions:
$$\beta (\gamma ^\mu )^{}\beta =\gamma ^\mu ,\beta (𝒥^{\mu \nu })^{}\beta =𝒥^{\mu \nu },$$
(A4)
with solution $`\beta `$$`=`$$`i\gamma ^0`$. It is easy to verify that the fermion bilinear:
$$\psi ^+\beta \psi \psi ^+(D(\mathrm{\Lambda }))^{}\beta D(\mathrm{\Lambda })\psi =\psi ^+\beta (D(\mathrm{\Lambda })^1D(\mathrm{\Lambda }))\psi =\psi ^+\beta \psi \overline{\psi }\psi ,$$
(A5)
is Lorentz invariant. Defining components, $`(\psi _\mu )^T`$$`=(\psi _\mu ^+,\psi _\mu ^{})`$, the free fermion Lagrangian on flat world-sheets can be written in component form:
$$=\overline{\psi }^\mu \gamma ^m\pm \psi _\mu ^{}=i(\psi ^\mu )^{}(_0+_1)\psi _\mu ^{}+i(\psi ^{+\mu })^{}(_0_1)\psi _\mu ^+.$$
(A6)
Notice that, with the conventions above, charge conjugation is defined by the Majorana condition, $`\zeta ^{}`$$`=`$$`\gamma \zeta `$, with spinors $`\zeta `$ and $`\gamma (\zeta )^{}`$ transforming identically under Lorentz transformations:
$$\gamma (\gamma ^\mu )\gamma ^1=\gamma (\gamma ^\mu )\gamma =(\gamma ^\mu )^{},\gamma (𝒥^{\mu \nu })\gamma ^1=(𝒥^{\mu \nu })^{}.$$
(A7)
We can therefore choose the component fermions, $`\zeta ^+`$, $`\zeta ^{}`$, to be real.
Analytically continuing to world-sheets with Euclidean signature, we can set $`\sigma ^1`$$`=`$$`it`$, $`\sigma ^2`$$`=`$$`x`$. Thus, we replace, $`_0`$$``$$`i_1`$, $`_1`$$``$$`_2`$, in Eq. (A6) to obtain the Euclidean Lagrangian in component form:
$$_{}=\psi ^\mu (_1i_2)\psi _\mu ^{}+\psi ^{+\mu }(_1+i_2)\psi _\mu ^+.$$
(A8)
It is easy to verify that Eq. (A8) can be recovered from the Lagrangian, $`_{}`$$`=`$$`\overline{\psi }^\mu \gamma ^a_a\psi _\mu `$, $`\overline{\psi }`$$``$$`(\psi )^T\gamma ^1`$ with the following choice of gamma matrices with Euclidean metric:
$$\gamma ^1=\left(\begin{array}{cc}\hfill 0& \hfill i\\ \hfill i& \hfill 0\end{array}\right),\gamma ^2=\left(\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 1& \hfill 0\end{array}\right),\gamma ^5=\gamma =i\gamma ^1\gamma ^2=\left(\begin{array}{cc}\hfill 1& \hfill 0\\ \hfill 0& \hfill 1\end{array}\right).$$
(A9)
From the equation of motion, it is clear that the component fermions, $`\psi ^+`$, $`\psi ^{}`$, transform, respectively, as left-handed, and right-handed, spinors on the world-sheet. We remark that our spinor conventions correspond to those used in the text . The reader can verify the identities:
$$\overline{\xi }\chi =(\xi )^T\gamma ^1\chi =i(\chi ^+\xi ^{}\xi ^{}\chi ^+)=\overline{\chi }\xi ,\overline{\xi }\gamma ^a\chi =\overline{\chi }\gamma ^a\xi ,\overline{\xi }\gamma ^a\gamma ^b\chi =\overline{\chi }\gamma ^b\gamma ^a\xi .$$
(A10)
The free fermion Lagrangian can be extended to a two-dimensional action invariant under a local $`N`$$`=`$$`1`$ supersymmetry following . We must be careful to retain any boundary terms resulting from an integration by parts since our interest is in the classical action for world-sheets with boundary. Appending $`d`$ free fermions to the bosonic string world sheet gives the free field action:
$$S_0=\frac{1}{4\pi \alpha ^{}}d^2\sigma [^aX^\mu _bX_\mu +\alpha ^{}(\psi ^\mu (_1i_2)\psi _\mu ^{}+\psi ^{+\mu }(_1+i_2)\psi _\mu ^+)].$$
(A11)
A variation of the free field action under the global supersymmetry transformation:
$$\delta X^\mu =\sqrt{\frac{\alpha ^{}}{2}}(\overline{\xi }\psi ^\mu ),\sqrt{\frac{\alpha ^{}}{2}}(\delta \psi ^\mu )=\frac{1}{2}(_aX^\mu )\gamma ^a\xi ,\sqrt{\frac{\alpha ^{}}{2}}(\delta \overline{\psi }^\mu )=\frac{1}{2}\overline{\xi }\gamma ^a(_aX^\mu ),$$
(A12)
results in the variation,
$$2\pi \alpha ^{}(\delta _\mathcal{0})=\sqrt{\frac{\alpha ^{}}{2}}\left\{_a[\frac{1}{2}(_bX_\mu )(\overline{\xi }\gamma ^a\gamma ^b\psi ^\mu )]+(_a\overline{\xi })(_bX_\mu )(\overline{\xi }\gamma ^a\gamma ^b\psi ^\mu )\right\}.$$
(A13)
Note that the second term vanishes for a covariantly constant spinor supersymmetry parameter, $`\xi `$. We can identify the Noether current, $`J^a`$$`=`$$`(_bX^\mu )(\gamma ^b\gamma ^a\psi _\mu )`$, and introduce a fermionic source term in the Lagrangian: $`\overline{\chi }_aJ^a`$, where $`_a\xi `$$``$$`(\delta \chi )`$ . The resulting variations close upon inclusion of a new term quartic in the fermions. The result is the Deser-Zumino-Brink-De Vecchia-Howe action for the fermionic string . Introducing the world-sheet metric, $`g_{mn}`$$`=`$$`e_m^ae_{na}`$, we can write the action in covariant form:
$$S=\frac{1}{2\pi \alpha ^{}}d^2\sigma \sqrt{g}[\frac{1}{2}g^{mn}_mX^\mu _nX_\mu +\frac{1}{2}\alpha ^{}\overline{\psi }^\mu \gamma ^m_m\psi _\mu +\sqrt{\frac{\alpha ^{}}{2}}(\overline{\chi }_a\gamma ^m\gamma ^a\psi ^\mu )(_mX_\mu )+\frac{\alpha ^{}}{4}(\overline{\chi }_a\gamma ^b\gamma ^a\psi ^\mu )(\overline{\chi }_b\psi _\mu )],$$
(A14)
invariant under the local supersymmetry transformations:
$`\delta X^\mu `$ $`=\sqrt{\frac{\alpha ^{}}{2}}(\delta \overline{\xi })\psi ^\mu `$ (A15)
$`\sqrt{\frac{\alpha ^{}}{2}}(\delta \psi ^\mu )`$ $`=\frac{1}{2}(_aX^\mu )\gamma ^a\delta \xi +\frac{1}{2}\sqrt{\frac{\alpha ^{}}{2}}(\psi ^\mu \chi _m)\gamma ^m\delta \xi `$ (A16)
$`\delta \chi _m`$ $`=D_m(\delta \xi )`$ (A17)
$`\delta e_m^a`$ $`=(\delta \overline{\xi })\gamma ^a\chi _m.`$ (A18)
The action is invariant under both reparameterizations and supersymmetry transformations on the world-sheet. The constraints that arise from making these compatible on the boundary will be discussed in Section II. Let us briefly recall the local invariances of the world-sheet action in the bulk. Under a diffeomorphism of the world-sheet coordinates parameterized by $`\delta V^m`$, we have:
$`\delta e_m^a`$ $`=\delta V^n(_ne_m^a)+e_n^a(_m\delta V^n)`$ (A19)
$`\delta \chi _m`$ $`=\delta V^n(_n\chi _m)+\chi _n(_m\delta V^n)`$ (A20)
$`\delta X^\mu `$ $`=\delta V^n(_nX^\mu )`$ (A21)
$`\delta \psi ^\mu `$ $`=\delta V^n(_n\psi ^\mu ).`$ (A22)
In the tangent space at any point of the manifold, we can perform independent Lorentz rotations, $`\delta \omega `$:
$`\delta e_m^a`$ $`=(\delta \omega )ϵ^{ab}e_m^b`$ (A23)
$`\delta \chi _m`$ $`=\frac{1}{2}(\delta \omega )\gamma ^5\chi _m`$ (A24)
$`\delta \psi ^\mu `$ $`=\frac{1}{2}(\delta \omega )\gamma ^5\psi ^\mu ,`$ (A25)
which leave world-sheet scalars invariant. The scalars are also invariant under rescalings of the world-sheet metric and gravitino. A Weyl rescaling of the metric, $`\delta \mathrm{\Lambda }`$, induces the variations:
$`\delta e_m^a`$ $`=(\delta \mathrm{\Lambda })e_m^a`$ (A26)
$`\delta \chi _m`$ $`=\frac{1}{2}(\delta \mathrm{\Lambda })\chi _m`$ (A27)
$`\delta \psi ^\mu `$ $`=\frac{1}{2}(\delta \mathrm{\Lambda })\psi ^\mu .`$ (A28)
The action is also invariant under the fermionic rescaling, $`\delta \zeta `$, of the world-sheet gravitino, $`\delta \chi _m`$$`=`$$`(\delta \zeta )\gamma _m`$, known as a super-Weyl transformation, which leaves all other world-sheet fields fixed. Note that the total number of bosonic and fermionic gauge parameters are equal: $`(\delta \lambda ,\delta V^n,\delta \omega )`$ and $`(\delta \xi ,\delta \zeta )`$ contain four parameters each. They can be used to locally fix the world-sheet metric, $`e_m^a`$, and gravitino, $`\chi _m`$, to their values in the superconformal gauge .
## B Constant Spinors and the Path Integral in the $`(\beta ,\alpha )`$$`=`$$`(1,1)`$ Sector
For the Dbrane backgrounds studied in this paper, with either parallel and static branes, or a relative rotation in a single plane, the contribution to the vacuum amplitude from the path integral with $`(\beta ,\alpha )`$$`=`$$`(1,1)`$ vanishes. As mentioned in the text, this is always true unless one considers rotations in all four transverse planes: $`(1,2)`$, $`(3,4)`$, $`(5,6)`$, and $`(7,8)`$. By an analytic continuation, this implies that the contribution from the $`(1,1)`$ sector to the potential between point sources derived in section IV vanishes unless we consider a general motion with velocity components in at least four spatial directions. This can be compared with the discussion of a pair of D4branes in relative motion in four transverse Dirichlet directions given in . From the point of view of the path integral, the vanishing contribution is due to a Grassmann integration over a constant mode which is absent in the action. For motions in four transverse planes, the constant mode is absent for four of the Weyl fermions on the world-sheet. We will verify in this appendix that the Grassmann integration over the constant mode of the remaining Weyl fermion, $`(0,9)`$, can be saturated by an insertion in the path integral that comes from the supermodulus in the $`(1,1)`$ sector.
We complexify the component fermions as in Eq. (35) giving a total of five independent Weyl fermions on the world-sheet upon imposing the open string boundary conditions. Consider separating the constant mode, $`\psi _0^{\pm i}`$, $`i`$$`=`$$`1`$, $`\mathrm{}`$, $`5`$, with $`\psi ^{\pm i}`$$`=`$$`\psi _0^{\pm i}`$$`+`$$`(\psi ^{\pm i})^{}`$. We have,
$$[d\psi ]e^{S(\psi )}d\psi _0[d\psi ^{}]e^{S(\psi ^{})}.$$
(B1)
Since the integrand is independent of the constant mode, the Grassmann integral will vanish. We will now show that the fermionic Jacobian, $`J_f`$, defined in Eq. (23) has a nontrivial dependence on the supermodulus in the $`(1,1)`$ sector. This term can be exponentiated as a correction to the effective action for the fermions with the consequence that the Grassmann integration over one pair of constant fermion modes no longer vanishes. If these are the only fermionic zero-modes present, one obtains a non-zero contribution to the vacuum amplitude from the $`(1,1)`$ sector.
As discussed earlier, in the $`(1,1)`$ sector the cylinder has both a supermodulus and a superconformal Killing spinor, both of which are simply constant spinors on the world-sheet. We consider variations of the gravitino that preserve the gamma tracelessness condition, parallel to the gauge slice. Denoting the supermodulus as $`\nu `$, a constant two component spinor, we can write,
$$\chi _1=\nu ,\chi _2=i\gamma ^5\nu .$$
(B2)
A traceless variation of the gravitino can be decomposed as $`\delta \chi _m=D_m\xi ^{}+\chi _{m,\alpha }d\nu ^\alpha `$, where $`\alpha `$ labels the two spin-components of $`\nu `$. Substituting in Eq. (26) gives:
$`|\delta \chi _m|^2=i{\displaystyle d^2\sigma \sqrt{g}(\delta \overline{\xi }^{}d\nu ^\alpha )\left(\begin{array}{cc}D_mD^m& D^m\chi _{m,\beta }\\ \overline{\chi }_{m,\alpha }D^m& \overline{\chi }_{m,\alpha }\chi _{,\beta }^m\end{array}\right)\left(\begin{array}{c}\delta \xi ^{}\\ d\nu ^\beta \end{array}\right)}.`$ (B7)
Substituting in Eq. (25),
$$1=[d\delta \chi _m]e^{|\delta \chi _m|^2/2}=J_f(\widehat{g})[d\delta \xi _0][d\delta \xi ^{}]𝑑\nu e^{|\delta \chi _m|^2/2}.$$
(B8)
Including the contributions from the Grassmann integrations over the super conformal Killing spinor and the supermodulus given in Eqs. (27), (28), and substituting for the Jacobian matrix in Eq. (B7), gives the result:
$$J_f(\widehat{g})=[det(P_{1/2}^{}P_{1/2})]^{1/2}(1/l^2+1)^1.$$
(B9)
Thus, the fermionic path integral takes the form:
$`{\displaystyle _{[1,1]}}{\displaystyle \frac{[d\delta \chi _m]}{\mathrm{Vol}(\mathrm{sWeyl}\times \mathrm{sDiff})}}`$ $`e^{\frac{1}{4\pi }{\scriptscriptstyle d^2\sigma \sqrt{g}(\overline{\chi }^m\psi ^\mu )(\overline{\chi }_m\psi _\mu )}}`$ (B11)
$`={\displaystyle _{[1,1]}}𝑑\nu [\mathrm{det}^{}(P_{1/2}^{}P_{1/2})]^{1/2}(1/l^2+1)^1e^{(1/l^2+1)\frac{1}{2\pi }{\scriptscriptstyle d^2\sigma \sqrt{g}(\nu ^+\nu ^{}\psi ^{+\mu }\psi _\mu ^{})}},`$
where we use component form for the fermions in the action. Integrating over $`\nu `$, we are left with the following insertion in the path integral for the matter fermions:
$$[\mathrm{det}^{}(P_{1/2}^{}P_{1/2})]^{1/2}\frac{1}{2\pi }d^2\sigma \sqrt{g}(\psi ^{+\mu }\psi _\mu ^{})$$
(B12)
As mentioned in the text, the functional determinant, $`[det^{}(P_{1/2}^{}P_{1/2})]^{1/2}`$, precisely cancels the contribution to the amplitude from the non-constant modes of one pair of component fermions. We choose these to be the $`(0,9)`$ pair. Complexify as described above. Summing on $`i`$$`=`$$`1`$, $`\mathrm{}`$, $`5`$, the insertion takes the form,
$$\frac{1}{2\pi }d^2\sigma \sqrt{g}(\psi ^{+i}\psi ^i)=\frac{l}{2\pi }\psi _0^{+5}\psi _0^5+\mathrm{},$$
(B13)
where the $`\mathrm{}`$ denote the dependence on the remaining constant modes, if present. Thus, precisely one pair of constant modes is saturated by the insertion. If no additional constant modes are present, the resulting path integral gives a non-vanishing result. The fermionic oscillator contributions in the $`(1,1)`$ sector are computed by zeta function regularization as in Eq. (49). As an example, consider the relative motion of a pair of Dpbranes in the directions $`(X^1,X^3,X^5,X^7)`$, with $`v_i`$$`=`$$`\mathrm{tanh}(u_i)`$, $`i`$$`=`$$`1`$, $`\mathrm{}`$, $`4`$. Then the contribution to the annulus from the $`(1,1)`$ sector takes the form:
$$𝒜_{[1,1]}(r,u)=\frac{1}{2}V_p_0^{\mathrm{}}\frac{dl}{l}(4\pi ^2\alpha ^{}l)^{p/2}e^{r^2l/4\pi \alpha ^{}}\underset{i=1}{\overset{4}{}}\frac{\mathrm{\Theta }_{(1,1)}(u_il/2\pi ,\frac{il}{2})}{i\mathrm{\Theta }_{11}(u_il/2\pi ,\frac{il}{2})},$$
(B14)
previously derived in .
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# NEUTRON STAR STRUCTURE AND THE EQUATION OF STATE
## 1 INTRODUCTION
The theoretical study of the structure of neutron stars is crucial if new observations of masses and radii are to lead to effective constraints on the equation of state (EOS) of dense matter. This study becomes ever more important as laboratory studies may be on the verge of yielding evidence about the composition and stiffness of matter beyond the nuclear equilibrium density $`\rho _s2.710^{14}`$ g cm<sup>-3</sup>. Rhoades & Ruffini (1974) demonstrated that the assumption of causality beyond a fiducial density $`\rho _f`$ sets an upper limit to the maximum mass of a neutron star: $`4.2\sqrt{\rho _s/\rho _f}\mathrm{M}_{}`$. However, theoretical studies of dense matter have considerable uncertainty in the high-density behavior of the EOS largely because of the poorly constrained many-body interactions. These uncertainties have prevented a firm prediction of the maximum mass of a beta-stable neutron star.
To date, several accurate mass determinations of neutron stars are available from radio binary pulsars (Thorsett & Chakrabarty 1998), and they all lie in a narrow range ($`1.251.44`$ M). One neutron star in an X-ray binary, Cyg X-2, has an estimated mass of $`1.8\pm 0.2`$ M (Orosz & Kuulkers 1999), but this determination is not as clean as for a radio binary. Another X-ray binary, Vela X-1, has a reported mass around 1.9 M (van Kerkwijk et al. 1995a), although Stickland et al. (1997) argue it to be about 1.4 M. A third object, the eclipsing X-ray binary 4U 1700-37, apparently contains an object with a mass of $`1.8\pm 0.4`$ M (Heap & Corcoran 1992), but Brown, Weingartner & Wijers (1996) have argued that since this source does not pulse and has a relatively hard X-ray spectrum, it may contain a low-mass black hole instead. It would not be surprising if neutron stars in X-ray binaries had larger masses than those in radio binaries, since the latter have presumably accreted relatively little mass since their formation. Alternatively, Cyg X-2 could be the first of a new and rarer population of neutron stars formed with high masses which could originate from more massive, rarer, supernovae. If the high masses for Cyg X-2 or Vela X-1 are confirmed, significant constraints on the equation of state would be realized.
On the other hand, there is a practical, albeit theoretical, lower mass limit for neutron stars, about $`1.11.2`$ M, which follows from the minimum mass of a protoneutron star. This is estimated by examining a lepton-rich configuration with a low-entropy inner core of $`0.6`$ M and a high-entropy envelope (Goussard, Haensel & Zdunik 1998). This argument is in general agreement with the theoretical result of supernova calculations, in which the inner homologous collapsing core material comprises at least 1 M.
Although accurate masses of several neutron stars are available, a precise measurement of the radius does not yet exist. Lattimer et al. (1990) (see also Glendenning 1992) have shown that the causality constraint can be used to set a lower limit to the radius: $`R>3.04GM/Rc^2`$. For a 1.4 M star, this is about 4.5 km.
Estimates of neutron star radii from observations have given a wide range of results. Perhaps the most reliable estimates stem from observations of thermal emissions from neutron star surfaces, which yield values of the so-called “radiation radius”
$$R_{\mathrm{}}=R/\sqrt{12GM/Rc^2},$$
(1)
a quantity resulting from redshifting the stars luminosity and temperature. A given value of $`R_{\mathrm{}}`$ implies that $`R<R_{\mathrm{}}`$ and $`M<0.13(R_{\mathrm{}}/\mathrm{km})`$ M. Thus, a 1.4 M neutron star requires $`R_{\mathrm{}}>10.75`$ km. Those pulsars with at least some suspected thermal radiation generically yield effective values of $`R_{\mathrm{}}`$ so small that it is believed that a significant part of the radiation originates from polar hot spots rather than from the surface as a whole. For example, Golden & Shearer (1999) found that upper limits to the unpulsed emission from Geminga, coupled with a parallactic distance of 160 pc, yielded values of $`R_{\mathrm{}}<9.5`$ km for a blackbody source and $`R_{\mathrm{}}<10`$ km for a magnetized H atmosphere. Similarly, Schulz (1999) estimated emission radii of less than 5 km, assuming a blackbody for eight low mass X-ray binaries.
Other attempts to deduce a radius include analyses (Titarchuk 1994) of X-ray bursts from sources 4U 1705-44 and 4U 1820-30 which implied rather small values, $`9.5<R_{\mathrm{}}<14`$ km. Recently, Rutledge et al. (1999) found that thermal emission from neutron stars of a canonical 10 km radius was indicated by the interburst emission. However, the modeling of the photospheric expansion and touchdown on the neutron star surface requires a model dependent relationship between the color and effective temperatures.
Absorption lines in X-ray spectra have also been investigated (Inoue 1992) with a view to deducing the neutron star radius. Candidates for the matter producing the absorption lines are either the accreted matter from the companion star or the products of nuclear burning in the bursts. In the former case, the most plausible element is thought to be Fe, in which case the relation $`R3.2GM/c^2`$, only slightly larger than the minimum possible value based upon causality (Lattimer et al. 1990; Glendenning 1992) is inferred. In the latter case, plausible candidates are Ti and Cr, and larger values of the radius would be obtained. In both cases, serious difficulties remain in interpreting the large line widths, of order 100–500 eV, in the $`4.1\pm 0.1`$ keV line observed from many sources.
A first attempt at using light curves and pulse fractions from pulsars to explore the $`MR`$ relation suggested relatively large radii, of order 15 km (Page 1995). However, this method, which assumed dipolar magnetic fields, was unable to satisfactorily reconcile the calculated magnitudes of the pulse fractions and the shapes of the light curves with observations.
The discovery of Quasi-Periodic Oscillations (QPOs) from X-ray emitting neutron stars in binaries provides a possible way of limiting neutron star masses and radii. These oscillations are manifested as quasi-periodic X-ray emissions, with frequencies ranging from tens to over 1200 Hz. Some QPOs show multiple frequencies, in particular, two frequencies $`\nu _1`$ and $`\nu _2`$ at several hundred Hz. These frequencies are not constant, but tend to both increase with time until the signal ultimately weakens and disappears. In the beat frequency model (Alpar & Shaham 1985, Psaltis et al. 1998), the highest frequency $`\nu _2`$ is associated with the Keplerian frequency $`\nu _K`$ of inhomogeneities or blobs in an accretion disc. The largest such frequency measured to date is $`\nu _{max}=1230`$ Hz. However, general relativity predicts the existence of a maximum orbital frequency, since the inner edge of an accretion disc must remain outside of the innermost stable circular orbit, located at a radius of $`r_s=6GM/c^2`$ in the absence of rotation. This corresponds to a Keplerian orbital frequency of $`\nu _s=\sqrt{GM/r_s^3}/2\pi `$ if the star is non-rotating. Equating $`\nu _{max}`$ with $`\nu _s`$, and since $`R<r_s`$, one deduces
$$M<1.78\mathrm{M}_{};R<8.86(M/\mathrm{M}_{})\mathrm{km}.$$
(2)
Corrections due to stellar rotation are straightforward to deduce and produce small changes in these limits (Psaltis et al. 1998). The lower frequency $`\nu _1`$ is associated with a beat frequency between $`\nu _2`$ and the spin frequency of the star. This spin frequency is large enough, of order 250-350 Hz, to alter the metric from a Schwarzschild geometry, and increases the theoretical mass limit in equation (2) to about 2.1 M (Psaltis et al. 1998). This remains strictly an upper limit, however, unless further observations support the interpretation that $`\nu _{max}`$ is associated with orbits at precisely the innermost stable orbit.
If the frequency $`\nu _2\nu _1`$ is to be associated with the spin of the neutron star, it should remain constant in time. However, recent observations reveal that it changes with time in a given source. Osherovich & Titarchuk (1999) proposed a model in which $`\nu _1`$ is the Keplerian frequency and $`\nu _2`$ is a hybrid frequency of the Keplerian oscillator under the influence of a magnetospheric Coriolis force. In this model, the frequencies are related to the neutron star spin frequency $`\nu `$ by
$$\nu _2=\sqrt{\nu _1^2+(\nu /2\pi )^2}.$$
(3)
Osherovich & Titarchuk argue that this relation fits the observed variations of $`\nu _2`$ and $`\nu _1`$ in several QPOs. The Keplerian frequency in Osherovich & Titarchuk’s model, being associated with the lower frequency $`\nu _1`$, however, is at most 800 Hz, leading to an upper mass limit that is nearly 3 M and is therefore of little practical use to limit either the star’s mass or radius.
An alternative model, proposed by Stella & Vietri (1999), associates $`\nu _2`$ with the Keplerian frequency of the inner edge of the disc, $`\nu _K`$, and $`\nu _2\nu _1`$ with the precession frequency of the periastron of slightly eccentric orbiting blobs at radius $`r`$ in the accretion disc. In a Schwarzschild geometry, $`\nu _1=\nu _K\sqrt{16GM/rc^2}`$. Note that $`(\nu _K\nu _1)^1`$ is the timescale that an orbiting blob recovers its original orientation relative to the neutron star and the Earth, so that variations in flux are expected to be observed at both frequencies $`\nu _K`$ and $`\nu _K\nu _1`$. Presumably, even eccentricities of order $`10^4`$ lead to observable effects. This model predicts that
$$\nu _1/\nu _2=1\sqrt{16(GM\nu _2)^{2/3}/c^2},$$
(4)
a relation that depends only upon $`M`$. Equation (4) can also fit observations of QPOs, but only if $`1.9<M/\mathrm{M}_{}<2.1`$. This result is not very sensitive to complicating effects due to stellar rotation: the Lense-Thirring effect and oblateness. This mechanism only depends on gravitometric effects, and may apply also to accreting black hole systems (Stella, Vietri & Morsink 1999).
Prospects for a radius determination have improved in recent years with the discovery of a class of isolated, non-pulsing, neutron stars. The first of these is the nearby object RX J185635-3754, initially discovered in X-rays (Walter, Wolk & Neuhaüser 1996) and confirmed with the Hubble Space Telescope (Walter & Matthews 1997). The observed X-rays, from the ROSAT satellite, are consistent with blackbody emission with an effective temperature of about 57 eV and relatively little extinction. The fortuitous location of the star, in the foreground of the R CrA molecular cloud, coupled with the small levels of extinction, limits the distance to $`D<120`$ pc. The fact that the source is not observable in radio and its lack of variability in X-rays implies that it is not a pulsar, unlike other identified radio-silent isolated neutron stars. This gives the hope that the observed radiation is not contaminated with non-thermal emission as in the case for pulsars.
The X-ray flux of RX J185635-3754, combined with a best-fit blackbody $`T_{eff}=57`$ eV, yields $`R_{\mathrm{}}7.3(D/120\mathrm{pc})\mathrm{km}`$. Such a value for $`R_{\mathrm{}}`$, even coupled with the maximum distance of 120 pc, is too small to be consistent with any neutron star with more than 1 M. But the optical flux is about a factor of 4 brighter than what is predicted by the X-ray blackbody. The reconciliation the X-ray and optical fluxes through atmosphere modeling naively implies an increase in $`R_{\mathrm{}}`$ of approximately $`4^{2/3}2.5`$. (This results since the X-ray flux is proportional to $`R_{\mathrm{}}^2T_{eff}^4`$, while the optical flux is on the Rayleigh-Jeans tail of the spectrum and is hence proportional to $`R_{\mathrm{}}^2T_{eff}`$. One seeks to enhance the predicted optical flux by 4 while keeping the X-ray flux fixed, as this is approximately equal to the total flux.) An et al. (2000) determined for non-magnetized heavy element atmospheres that $`R_{\mathrm{}}/D0.18\pm 0.05`$ km pc<sup>-1</sup>, which is rough agreement with the above naive expectations. However, uncertainties in the atmospheric composition and the quality of the existing data precluded obtaining a more precise estimate of $`R_{\mathrm{}}/D`$. An et al. concluded, in agreement with expectations based upon the general results of Romani (1987) and Rajagopal, Romani & Miller (1997), that the predicted spectrum of a heavy element atmosphere, but not a light element atmosphere, was consistent with all the observations. This is in contrast to the conclusions of Pavlov et al. (1996), whose results for RX J185635-3754 implied that the observations in the optical and X-ray bands were incompatible with atmospheric modelling for both heavy element and light element non-magnetized atmospheres, unless the distance to this star is greater than the presumed maximum of 120 pc based upon the star’s location in front of the R Cor Aus molecular cloud. Future prospects for determining the radius of this neutron star are discussed in §7.
Our objectives in this paper are 1) to demonstrate specifically how the accurate measurement of a neutron star radius would constrain the dense matter EOS, and 2) to provide general relationships for other structural quantities, such as the moment of inertia and the binding energy, that are relatively EOS-independent, and which could be used to constrain the neutron star mass and/or radius. We will examine a wide class of equations of state, including those that have extreme softening at high densities. In addition, we will examine analytic solutions to Einstein’s equations which shed light on the results we deduce empirically. In all cases, we will focus on non-rotating, non-magnetized neutron stars at zero temperature.
Lindblom (1992) had suggested that a series of mass and radius measurements would be necessary to accurately constrain the dense matter equation of state. His technique utilizes a numerical inversion of the neutron star structure equation. Our results instead suggest that important constraints on the EOS can be achieved with even a single radius measurement, if it is accurate enough, and that the quality of the constraint is not very sensitive to the mass. The fact that the range of accurately determined neutron star masses is so small, only about 0.2 M to date, further implies that important constraints can be deduced without simulaneous mass-radius measurements. Of course, several measurements of neutron star masses and radii would greatly enhance the constraint on the equation of state.
In § 2, the equations of state selected in this paper are discussed. In § 3, the mass-radius relation for a sample of these equations of state are discussed. A quantitative relationship between the radii of normal neutron stars and the pressure of matter in the vicinity of $`n_s`$ is empirically established and theoretically justified. In turn, how the matter’s pressure at these densities depends upon fundamental nuclear parameters is developed. In § 4, analytic solutions to the general relativistic equations of hydrostatic equilibrium are explored. These lead to useful approximations for neutron star structure and which directly correlate other observables such as moments of inertia and binding energy to the mass and radius. It is believed that the distribution of the moment of inertia inside the star is crucial in the interpretation of glitches observed in the spin down of pulsars, so that measurements of the sizes and frequencies of glitches can constrain neutron star masses and radii (Link, Epstein & Lattimer 1999). In § 5, expressions for the fraction of moment of inertia contained within the stellar crust, as a function of mass, radius, and equation of state, are derived. In § 6, expressions for the binding energy are derived. § 7 contains a summary and outlook.
## 2 EQUATIONS OF STATE
The composition of a neutron star chiefly depends on the nature of strong interactions, which are not well understood in dense matter. Most models that have been investigated can be conveniently grouped into three broad categories: nonrelativistic potential models, relativistic field theoretical models, and relativistic Dirac-Brueckner-Hartree-Fock models. In each of these approaches, the presence of additional softening components such as hyperons, Bose condensates or quark matter, can be incorporated. Details of these approaches have been further considered in Lattimer et al. (1990) and Prakash et al. (1997). A representative sample, and some general attributes, including references and typical compositions, of equations of state employed in this paper are summarized in Table 1.
We have used four equations of state taken from Akmal & Pandharipande (1998). These are: AP1 (the AV18 potential), AP2 (the AV18 potential plus $`\delta v_b`$ relativistic boost corrections), AP3 (the AV18 potential plus a three-body UIX potential ), and AP4 (the AV18 potential plus the UIX potential plus the $`\delta v_b`$ boost). Three equations of state from Müller & Serot (1996), labelled MS1–3, correspond to different choices of the parameters $`\xi `$ and $`\zeta `$ which determine the strength of the nonlinear vector and isovector interactions at high densities. The numerical values used are $`\xi =\zeta =0;\xi =1.5,\zeta =0.06`$; and $`\xi =1.5,\zeta =0.02`$, respectively. Six EOSs from the phenomenological non-relativistic potential model of Prakash, Ainsworth & Lattimer (1988), labelled PAL1–6, were chosen, which have different choices of the symmetry energy parameter at the saturation density, its density dependence, and the bulk nuclear matter incompressibility parameter $`K_s`$. The incompressibilities of PAL1–5 were chosen to be $`K_s=180`$ or 240 MeV, but PAL6 has $`K_s=120`$ MeV. Three interactions from the field-theoretical model of Glendenning & Moszkowski (1991) are taken from their Table II; in order, they are denoted GM1–3. Two interactions from the field-theoretical model of Glendenning & Schaffner-Bielich (1999) correspond, in their notation, to GL78 with $`U_K(\rho _0)=140`$ MeV and TM1 with $`U_K=185`$ MeV. The labels denoting the other EOSs in Table I are identical to those in the original references.
The rationale for exploring a wide variety of EOSs, even some that are relatively outdated or in which systematic improvements are performed, is two-fold. First, it provides contrasts among widely different theoretical paradigms. Second, it illuminates general relationships that exist between the pressure-density relation and the macroscopic properties of the star such as the radius. For example, AP4 represents the most complete study to date of Akmal & Pandharipande (1998), in which many-body and special relativistic corrections are progressively incorporated into prior models, AP1–3. AP1–3 are included here because they represent different pressure-energy density-baryon density relations, and serve to reinforce correlations between neutron star structure and microscopic physics observed using alternative theoretical paradigms. Similarly, several different parameter sets for other EOSs are chosen.
In all cases, except for PS (Pandharipande & Smith 1975), the pressure is evaluated assuming zero temperature and beta equilibrium without trapped neutrinos. PS only contains neutrons among the baryons, there being no charged components. We chose to include this EOS, despite the fact that it has been superceded by more sophisticated calculations by Pandharipande and coworkers, because it represents an extreme case producing large radii neutron stars.
The pressure-density relations for some of the selected EOSs are shown in Figure 1. There are two general classes of equations of state. First, normal equations of state have a pressure which vanishes as the density tends to zero. Second, self-bound equations of state have a pressure which vanishes at a significant finite density.
The best-known example of self-bound stars results from Witten’s (1984) conjecture (also see Fahri & Jaffe 1984, Haensel, Zdunik & Schaeffer 1986, Alcock & Olinto 1988, and Prakash et al. 1990) that strange quark matter is the ultimate ground state of matter. In this paper, the self-bound EOSs are represented by strange-quark matter models SQM1–3, using perturbative QCD and an MIT-type bag model, with parameter values given in Table 2. The existence of an energy ceiling equal to the baryon mass, 939 MeV, for zero pressure matter requires that the bag constant $`B94.92`$ MeV fm<sup>-3</sup>. This limiting value is chosen, together with zero strange quark mass and no interactions ($`\alpha _c=0`$), for the model SQM1. The other two models chosen, SQM2 and SQM3, have bag constants adjusted so that their energy ceilings are also 939 MeV.
For normal matter, the EOS is that of an interacting nucleon gas above a transition density of 1/3 to 1/2 $`n_s`$. Below this density, the ground state of matter consists of heavy nuclei in equilibrium with a neutron-rich, low-density gas of nucleons. In general, a self-consistent evaluation of the equilibrium that exists below the transition density, and the evaluation of the transition density itself, has been carried out for only a few equations of state (e.g., Bethe, Pethick & Sutherland 1972, Negele & Vautherin 1974, Lattimer et al. 1985; Lattimer and Swesty 1990). We have therefore not plotted the pressure below about 0.1 MeV fm<sup>-3</sup> in Figure 1. For densities $`0.001<n<0.08`$ fm<sup>-3</sup> we employ the EOS of Negele & Vautherin (1974), while for densities $`n<0.001`$ fm<sup>-3</sup> we employ the EOS of Bethe, Pethick & Sutherland (1972). However, for most of the purposes of this paper, the pressure in the region $`n<0.1`$ fm<sup>-3</sup> is not relevant as it does not significantly affect the mass-radius relation or other global aspects of the star’s structure. Nevertheless, the value of the transition density, and the pressure there, are important ingredients for the determination of the size of the superfluid crust of a neutron star that is believed to be involved in the phenomenon of pulsar glitches (Link, Epstein & Lattimer 1999).
There are three significant features to note in Figure 1 for normal EOSs. First, there is a fairly wide range of predicted pressures for beta-stable matter in the density domain $`n_s/2<n<2n_s`$. For the EOSs displayed, the range of pressures covers about a factor of five, but this survey is by no means exhaustive. That such a wide range in pressures is found is somewhat surprising, given that each of the EOSs provides acceptable fits to experimentally-determined nuclear matter properties. Clearly, the extrapolation of the pressure from symmetric matter to nearly pure neutron matter is poorly constrained. Second, the slopes of the pressure curves are rather similar. A polytropic index of $`n1`$, where $`P=Kn^{1+1/n}`$, is implied. Third, in the density domain below $`2n_s`$, the pressure-density relations seem to fall into two groups. The higer pressure group is primarily composed of relativistic field-theoretical models, while the lower pressure group is primarily composed of non-relativistic potential models. As we show in § 3, the pressure in the vicinity of $`n_s`$ is mostly determined by the symmetry energy properties of the EOS, and it is significant that relativistic field-theoretical models generally have symmetry energies that increase proportionately to the density while potential models have much less steeply rising symmetry energies.
A few of the plotted normal EOSs have considerable softening at high densities, especially PAL6, GS1, GS2, GM3, PS and PCL2. PAL6 has an abnormally small value of incompressibility ($`K_s=120`$ MeV). GS1 and GS2 have phase transitions to matter containing a kaon condensate, GM3 has a large population of hyperons appearing at high density, PS has a phase transition to a neutral pion condensate and a neutron solid, and PCL2 has a phase transition to a mixed phase containing strange quark matter. These EOSs can be regarded as representative of the many suggestions of the kinds of softening that could occur at high densities.
## 3 NEUTRON STAR RADII
Figure 2 displays the mass-radius relation for cold, catalyzed matter using these EOSs. The causality constraint described earlier and contours of $`R_{\mathrm{}}`$ are also indicated in Figure 2. With the exception of model GS1, the EOSs used to generate Figure 2 result in maximum masses greater than 1.442 M, the limit obtained from PSR 1913+16. From a theoretical perspective, it appears that values of $`R_{\mathrm{}}`$ in the range of 12–20 km are possible for normal neutron stars whose masses are greater than 1 M.
Corresponding to the two general types of EOSs, there are two general classes of neutron stars. Normal neutron stars are configurations with zero density at the stellar surface and which have minimum masses, of about 0.1 M, that are primarily determined by the EOS below $`n_s`$. At the minimum mass, the radii are generally in excess of 100 km. The second class of stars are the so-called self-bound stars, which have finite density, but zero pressure, at their surfaces. They are represented in Figure 2 by strange quark matter stars (SQM1–3).
Self-bound stars have no minimum mass, unlike the case of normal neutron stars for which pure neutron matter is unbound. Unlike normal neutron stars, the maximum mass self-bound stars have nearly the largest radii possible for a given EOS. If the strange quark mass $`m_s=0`$ and interactions are neglected ($`\alpha _c=0`$), the maximum mass is related to the bag constant $`B`$ in the MIT-type bag model by $`M_{max}=2.033(56\mathrm{MeV}\mathrm{fm}^3/B)^{1/2}\mathrm{M}_{}`$. Prakash et al. (1990) and Lattimer et al. (1990) showed that the addition of a finite strange quark mass and/or interactions produces larger maximum masses. The constraint that $`M_{max}>1.44`$ M is thus automatically satisfied for all cases by the condition that the energy ceiling is 939 MeV. In addition, models satisfying the energy ceiling constraint, with any values of $`m_s`$ and $`\alpha _c`$, have larger radii for every mass than the case SQM1. For the MIT model, the locus of maximum masses of self-bound stars is given simply by $`R1.85R_s`$ (Lattimer et al. 1990), where $`R_s=2GM/Rc^2`$ is the Schwarzschild radius, which is shown in the right-hand panel of Figure 2. Strange quark stars with electrostatically supported normal-matter crusts (Glendenning & Weber 1992) have larger radii than those with bare surfaces. Coupled with the additional constraint $`M>1\mathrm{M}_{}`$ from protoneutron star models, MIT-model strange quark stars cannot have $`R<8.5`$ km or $`R_{\mathrm{}}<10.5`$ km. These values are comparable to the possible lower limits for a Bose (pion or kaon) condensate EOS.
Although the $`MR`$ trajectories for normal stars can be strikingly different, in the mass range from 1 to 1.5 M or more it is usually the case that the radius has relatively little dependence upon the stellar mass. The major exceptions illustrated are the model GS1, in which a mixed phase containing a kaon condensate appears at a relatively low density and the model PAL6 which has an extremely small nuclear incompressibility (120 MeV). Both of these have considerable softening and a large increase in central density for $`M>1`$ M. Pronounced softening, while not as dramatic, also occurs in models GS2 and PCL2, which contain mixed phases containing a kaon condensate and strange quark matter, respectively. All other normal EOSs in this figure, except PS, contain only baryons among the hadrons.
While it is generally assumed that a stiff EOS implies both a large maximum mass and a large radius, many counter examples exist. For example, GM3, MS1 and PS have relatively small maximum masses but have large radii compared to most other EOSs with larger maximum masses. Also, not all EOSs with extreme softening have small radii for $`M>1`$ M (e.g., GS2, PS). Nonetheless, for stars with masses greater than 1 M, only models with a large degree of softening (including strange quark matter configurations) can have $`R_{\mathrm{}}<12`$ km. Should the radius of a neutron star ever be accurately determined to satisfy $`R_{\mathrm{}}<12`$ km, a strong case could be made for the existence of extreme softening.
To understand the relative insensitivity of the radius to the mass for normal neutron stars, it is relevant that a Newtonian polytrope with $`n=1`$ has the property that the stellar radius is independent of both the mass and central density. Recall that most EOSs, in the density range of $`n_s2n_s`$, have an effective polytropic index of about one (see Figure 1). An $`n=1`$ polytrope also has the property that the radius is proportional to the square root of the constant $`K`$ in the polytropic pressure law $`P=K\rho ^{1+1/n}`$. This suggests that there might be a quantitative relation between the radius and the pressure that does not depend upon the EOS at the highest densities, which determines the overall softness or stiffness (and hence, the maximum mass).
In fact, this conjecture may be verified. Figure 3 shows the remarkable empirical correlation which exists between the radii of 1 and 1.4 M normal stars and the matter’s pressure evaluated at fiducial densities of 1, 1.5 and 2 $`n_s`$. Table 1 explains the EOS symbols used in Figure 3. Despite the relative insensitivity of radius to mass for a particular EOS in this mass range, the nominal radius $`R_M`$, which is defined as the radius at a particular mass $`M`$ in solar units, still varies widely with the EOS employed. Up to $`5`$ km differences are seen in $`R_{1.4}`$, for example. Of the EOSs in Table 1, the only severe violations of this correlation occurs for PCL2 and PAL6 at 1.4 M for $`n_s`$, and for PS at both 1 and 1.4 M for $`2n_s`$. In the case of PCL2, this is relatively close to the maximum mass, and the matter has extreme softening due to the existence of a mixed phase with quark matter. (A GS model intermediate between GS1 and GS2, with a maximum mass of 1.44 M, would give similar results.) In the case of PS, it is clear from Figure 1 that extensive softening occurs already by $`1.5n_s`$. We emphasize that this correlation is valid only for cold, catalyzed neutron stars, i.e., not for protoneutron stars which have finite entropies and might contain trapped neutrinos.
Numerically, the correlation has the form of a power law:
$$R_MC(n,M)[P(n)]^{0.230.26},$$
(5)
where $`P(n)`$ is the total pressure inclusive of leptonic contributions evaluated at the density $`n`$, and $`C(n,M)`$ is a number that depends on the density $`n`$ at which the pressure was evaluated and the stellar mass $`M`$. An exponent of 1/4 was chosen for display in Figure 3, but the correlation holds for a small range of exponents about this value. Using an exponent of 1/4, and ignoring points associated with EOSs with phase transitions in the density ranges of interest, we find values for $`C(n,M)`$, in units of km fm<sup>3/4</sup> MeV<sup>-1/4</sup>, which are listed in Table 3. The error bars are taken from the standard deviations. The correlation is seen to be somewhat tighter for the baryon density $`n=1.5n_s`$ and $`2n_s`$ cases.
The fact that the exponent is considerably less than the Newtonian value of 1/2 can be quantitatively understood by considering a relativistic generalization of the $`n=1`$ polytrope due to Buchdahl (1967). He found that the EOS
$$\rho =12\sqrt{p_{}P}5P,$$
(6)
where $`p_{}`$ is a constant fiducial pressure independent of density, has an analytic solution of Einstein’s equations. This solution is characterized by the quantities $`p_{}`$ and $`\beta GM/Rc^2`$, and the stellar radius is found to be
$$R=(1\beta )c^2\sqrt{\frac{\pi }{288p_{}G(12\beta )}}.$$
(7)
For completeness, we summarize below the metric functions, the pressure and the mass-energy density as functions of coordinate radius $`r`$:
$`e^\nu `$ $``$ $`g_{tt}=(12\beta )(1\beta u)(1\beta +u)^1;`$ (8)
$`e^\lambda `$ $``$ $`g_{rr}=(12\beta )(1\beta +u)(1\beta u)^1(1\beta +\beta \mathrm{cos}Ar^{})^2;`$ (9)
$`8\pi PG/c^4`$ $`=`$ $`A^2u^2(12\beta )(1\beta +u)^2;`$ (10)
$`8\pi \rho G/c^2`$ $`=`$ $`2A^2u(12\beta )(1\beta 3u/2)(1\beta +u)^2.`$ (11)
where
$`r`$ $`=`$ $`r^{}(1\beta +u)(12\beta )^1;`$ (12)
$`u`$ $`=`$ $`\beta (Ar^{})^1\mathrm{sin}Ar^{};`$ (13)
$`A^2`$ $`=`$ $`288\pi p_{}Gc^4(12\beta )^1.`$ (14)
Note that $`Rp_{}^{1/2}(1+\beta ^2/2+\mathrm{})`$, so for a given value of $`p_{}`$, the radius increases very slowly with mass.
To estimate the exponent, it is instructive to analyze the response of $`R`$ to a change of pressure at some fiducial density $`\rho `$, for a fixed mass $`M`$. (At the relatively low densities of interest, the difference between using $`n`$ or $`\rho `$ in the following analysis is not significant.) We find the exponent to be
$`{\displaystyle \frac{d\mathrm{ln}R}{d\mathrm{ln}P}}|_{\rho ,M}`$ $`=`$ $`{\displaystyle \frac{d\mathrm{ln}R}{d\mathrm{ln}p_{}}}|_\beta {\displaystyle \frac{d\mathrm{ln}p_{}}{d\mathrm{ln}P}}|_\rho \left[1+{\displaystyle \frac{d\mathrm{ln}R}{d\mathrm{ln}\beta }}|_p_{}\right]^1`$ (15)
$`=`$ $`{\displaystyle \frac{1}{2}}\left(1{\displaystyle \frac{5}{6}}\sqrt{{\displaystyle \frac{P}{p_{}}}}\right){\displaystyle \frac{(1\beta )(12\beta )}{(13\beta +3\beta ^2)}}.`$ (16)
In the limit $`\beta 0`$, one has $`P0`$ and the exponent $`d\mathrm{ln}R/d\mathrm{ln}P|\rho ,M1/2`$, the value characteristic of an $`n=1`$ Newtonian polytrope. Finite values of $`\beta `$ and $`P`$ render the exponent smaller than 1/2. If the stellar mass and radius are about 1.4 M and 15 km, respectively, for example, equation (7) gives $`p_{}=\pi /(288R^2)4.8510^5`$ km<sup>-2</sup> (in geometrized units). Furthermore, if the fiducial density is $`\rho 1.5m_bn_s2.0210^4`$ km<sup>-2</sup> (also in geometrized units, with $`m_b`$ the baryon mass), equation (6) implies that in geometrized units $`P8.510^6`$ km<sup>-2</sup>. Since the value of $`\beta `$ in this case is 0.14, one then obtains $`d\mathrm{ln}R/d\mathrm{ln}P0.31`$. This result, while mildly sensitive to the choices for $`\rho `$ and $`R`$, provides a reasonable explanation of the correlation, equation (5). The fact that the exponent is smaller than 1/2 is clearly an effect due to general relativity.
The existence of this correlation is significant because the pressure of degenerate neutron-star matter near the nuclear saturation density $`n_s`$ is, in large part, determined by the symmetry properties of the EOS, as we now discuss. Thus, the measurement of a neutron star radius, if not so small as to indicate extreme softening, could provide an important clue to the symmetry properties of matter. In either case, valuable information will be obtained.
Studies of pure neutron matter strongly suggest that the specific energy of nuclear matter near the saturation density may be expressed as an expansion quadratic in the asymmetry $`(12x)`$, where $`x`$ is the proton fraction, which can be terminated after only one term (Prakash, Ainsworth & Lattimer 1988). In this case, the energy per particle and pressure of cold, beta stable nucleonic matter is
$`E(n,x)`$ $``$ $`E(n,1/2)+S_v(n)(12x)^2,`$
$`P(n,x)`$ $``$ $`n^2[E^{}(n,1/2)+S_v^{}(12x)^2],`$ (17)
where $`E(n,1/2)`$ is the energy per particle of symmetric matter and $`S_v(n)`$ is the bulk symmetry energy (which is density dependent). Primes denote derivatives with respect to density. At $`n_s`$, the symmetry energy can be estimated from nuclear mass systematics and has the value $`S_vS_v(n_s)2736\mathrm{MeV}`$. Attempts to further restrict this range from consideration of fission barriers and the energies of giant resonances have led to ambiguous results. Both the magnitude of $`S_v`$ and its density dependence $`S_v(n)`$ are currently uncertain. Part of the symmetry energy is due to the kinetic energy for noninteracting matter, which for degenerate nucleonic matter is proportional to $`n^{2/3}`$, but the remainder of the symmetry energy, due to interactions, is also expected to contribute significantly to the overall density dependence.
Leptonic contributions must to be added to equation (17) to obtain the total energy and pressure; the electron energy per baryon is $`(3/4)\mathrm{}cx(3\pi ^2nx)^{1/3}`$. Matter in neutron stars is in beta equilibrium, i.e., $`\mu _e=\mu _n\mu _p=E/x`$, which permits the evaluation of the equilibrium proton fraction and the total pressure $`P`$ may be written at a particular density in terms of fundamental nuclear parameters (Prakash 1996). For example, the pressure at the saturation density is simply
$`P_s=n_s(12x_s)[n_sS_v^{}(12x_s)+S_vx_s],`$ (18)
where $`S_v^{}(S_v(n)/n)_{n=n_s}`$ and the equilibrium proton fraction at $`n_s`$ is
$`x_s(3\pi ^2n_s)^1(4S_v/\mathrm{}c)^30.04,`$ (19)
for $`S_v=30`$ MeV. Due to the small value of $`x_s`$, we find that $`P_sn_sS_v^{}`$. The inclusion of muons, which generally begin to appear around $`n_s`$, does not qualitatively affect these results.
Were we to evaluate the pressure at a larger density, contributions featuring other nuclear parameters, including the nuclear incompressibility $`K_s=9(dP/dn)|n_s`$ and the skewness $`K_s^{}=27n_s^3(d^3E/dn^3)|_{n_s}`$, also become significant. For analytical purposes, the nuclear matter energy per baryon, in MeV, may be expanded in the vicinity of $`n_s`$ as
$`E(n,1/2)=16+{\displaystyle \frac{K_s}{18}}\left({\displaystyle \frac{n}{n_s}}1\right)^2{\displaystyle \frac{K_s^{}}{27}}\left({\displaystyle \frac{n}{n_s}}1\right)^3.`$ (20)
Experimental constraints to the compression modulus $`K_s`$, most importantly from analyses of giant monopole resonances (Blaizot et al. 1995; Youngblood et al. 1999), give $`K_s220`$ MeV. The skewness parameter $`K_s^{}`$ has been estimated to lie in the range 1780–2380 MeV (Pearson 1991, Rudaz et al. 1992), but in these calculations contributions from the surface symmetry energy were neglected. For values of $`K_s^{}`$ this large, equation (20) cannot be used beyond about 1.5$`n_s`$. Evaluating the pressure for $`n=1.5n_s`$, we find
$`P(1.5n_s)=2.25n_s[K_s/18K_s^{}/216+n_s(12x)^2S_v^{}].`$ (21)
Assuming that $`S_v(n)`$ is approximately proportional to the density, as it is in most relativistic field theoretical models, $`S_v^{}(n)S_v/n_s`$. Since the $`K_s`$ and $`K_s^{}`$ terms largely cancel, the symmetry term comprises most of the total. Once again, the result that the pressure is mostly sensitive to the density dependence of the symmetry energy is found.
The sensitivity of the radius to the symmetry energy can further demonstrated by the parametrized EOS of PAL (Prakash, Ainsworth & Lattimer 1988). The symmetry energy function $`S_v(n)`$ is a direct input in this parametrization and can be chosen to reproduce the results of more microscopic calculations. Figure 4 shows the dependence of mass-radius trajectories as the quantities $`S_v`$ and $`S_v(n)`$ are alternately varied. Clearly, of the two variations, the density dependence of $`S_v(n)`$ is more important in determining the neutron star radius. Note also the weak sensitivity of the maximum neutron star mass to $`S_v`$, and that the maximum mass depends more strongly upon the function $`S_v(n)`$.
At present, experimental guidance concerning the density dependence of the symmetry energy is limited and mostly based upon the division of the nuclear symmetry energy between volume and surface contributions. Upcoming experiments involving heavy-ion collisions which might sample densities up to $`(34)n_s`$, will be limited to analyzing properties of the nearly symmetric nuclear matter EOS through a study of matter, momentum, and energy flow of nucleons. Thus, studies of heavy nuclei far off the neutron drip lines using radioactive ion beams will be necessary in order to pin down the properties of the neutron-rich regimes encountered in neutron stars.
## 4 MOMENTS OF INERTIA
Besides the stellar radius, other global attributes of neutron stars are potentially observable, including the moment of inertia and the binding energy. These quantities depend primarily upon the ratio $`M/R`$ as opposed to details of the EOS, as can be readily seen by evaluating them using analytic solutions to Einstein’s equations. Although over 100 analytic solutions to Einstein’s equations are known (Delgaty & Lake 1998), nearly all of them are physically unrealistic. However, three analytic solutions are of particular interest in normal neutron star structure.
The first is the well-known Schwarzschild interior solution for an incompressible fluid, $`\rho =\rho _c`$, where $`\rho `$ is the mass-energy density. This case, hereafter referred to as “Inc”, is mostly of interest because it determines the minimum compactness $`\beta =GM/Rc^2`$ for a neutron star, namely 4/9, based upon the central pressure being finite. Two aspects of the incompressible fluid that are physically unrealistic, however, include the fact that the sound speed is everywhere infinite, and that the density does not vanish on the star’s surface.
The second analytic solution, due to Buchdahl (1967), is described in equation (11). We will refer to this solution as “Buch”.
The third analytic solution (which we will refer to as “T VII”) was discovered by Tolman (1939) and corresponds to the case when the mass-energy density $`\rho `$ varies quadratically, that is,
$$\rho =\rho _c[1(r/R)^2].$$
(22)
Of course, this behavior is to be expected at both extremes $`r0`$ and $`rR`$. However, this is also an eminently reasonable representation for intermediate regions, as displayed in Figure 5, which contains results for neutron stars more massive than 1.2 M. A wide variety of EOSs are sampled in this figure, and they are listed in Table 1.
Because the T VII solution is often overlooked in the literature (for exceptions, see, for example, Durgapal & Pande 1980 and Delgaty & Lake 1998), it is summarized here. It is useful in establishing interesting and simple relations that are insensitive to the EOS. In terms of the variable $`x=r^2/R^2`$ and the compactness parameter $`\beta =GM/Rc^2`$, the assumption $`\rho =\rho _c(1x)`$ results in $`\rho _c=15\beta c^2/(8\pi GR^2)`$. The solution of Einstein’s equations for this density distribution is:
$`e^\lambda `$ $`=`$ $`1\beta x(53x),e^\nu =(15\beta /3)\mathrm{cos}^2\varphi ,`$ (23)
$`P`$ $`=`$ $`{\displaystyle \frac{c^4}{4\pi R^2G}}\left[\sqrt{3\beta e^\lambda }\mathrm{tan}\varphi {\displaystyle \frac{\beta }{2}}(53x)\right],n={\displaystyle \frac{(\rho c^2+P)}{m_bc^2}}{\displaystyle \frac{\mathrm{cos}\varphi }{\mathrm{cos}\varphi _1}},`$ (24)
$`\varphi `$ $`=`$ $`(w_1w)/2+\varphi _1,w=\mathrm{log}\left[x5/6+\sqrt{e^\lambda /(3\beta )}\right],`$ (25)
$`\varphi _c`$ $`=`$ $`\varphi (x=0),\varphi _1=\varphi (x=1)=\mathrm{tan}^1\sqrt{\beta /[3(12\beta )]},w_1=w(x=1).`$ (26)
The central values of $`P/\rho c^2`$ and the square of the sound speed $`c_s^2`$ are
$$\frac{P}{\rho c^2}|_c=\frac{2}{15}\sqrt{\frac{3}{\beta }}(\frac{c_s}{c})^2,(\frac{c_s}{c})^2=\mathrm{tan}\varphi _c(\mathrm{tan}\varphi _c+\sqrt{\frac{\beta }{3}}).$$
(27)
This solution, like that of Buchdahl’s, is scale-free, with the parameters $`\beta `$ and $`\rho _c`$ (or $`M`$ and $`R`$). There are obvious limitations to the range of parameters for realistic models: when $`\varphi _c=\pi /2`$, or $`\beta 0.3862`$, $`P_c`$ becomes infinite, and when $`\beta 0.2698`$, $`c_s`$ becomes causal (i.e., $`c`$). Recall that for an incompressible fluid, $`P_c`$ becomes infinite when $`\beta =4/9`$, and this EOS is acausal for all values of $`\beta `$. For the Buchdahl solution, $`P_c`$ becomes infinite when $`\beta =2/5`$ and the causal limit is reached when $`\beta =1/6`$. For comparison, the causal limit for realistic EOSs is $`\beta 0.33`$ (Lattimer et al. 1990, Glendenning 1992), as previously discussed.
The general applicability of these exact solutions can be gauged by analyzing the moment of inertia, which, for a star uniformly rotating with angular velocity $`\mathrm{\Omega }`$, is
$$I=(8\pi /3)_0^Rr^4(\rho +P/c^2)e^{(\lambda \nu )/2}(\omega /\mathrm{\Omega })𝑑r.$$
(28)
The metric function $`\omega `$ is a solution of the equation
$$d[r^4e^{(\lambda +\nu )/2}\omega ^{}]/dr+4r^3\omega de^{(\lambda +\nu )/2}/dr=0$$
(29)
with the surface boundary condition
$$\omega _R=\mathrm{\Omega }\frac{R}{3}\omega _R^{}=\mathrm{\Omega }\left[1\frac{2GI}{R^3c^2}\right].$$
(30)
The second equality in the above follows from the definition of $`I`$ and the TOV equation. Writing $`j=\mathrm{exp}[(\nu +\lambda )/2]`$, the TOV equation becomes
$$j^{}=4\pi Gr(P/c^2+\rho )je^\lambda /c^2.$$
(31)
Then, one has
$$I=\frac{2c^2}{3G}\frac{\omega }{\mathrm{\Omega }}r^3𝑑j=\frac{c^2R^4\omega _R^{}}{6G\mathrm{\Omega }}.$$
(32)
Unfortunately, an analytic representation of $`\omega `$ or the moment of inertia for any of the three exact solutions is not available. However, approximations which are valid in the causal regime to within 0.5% are
$`I_{Inc}/MR^2`$ $``$ $`2(10.87\beta 0.3\beta ^2)^1/5,`$ (33)
$`I_{Buch}/MR^2`$ $``$ $`(2/34/\pi ^2)(11.81\beta +0.47\beta ^2)^1,`$ (34)
$`I_{TVII}/MR^2`$ $``$ $`2(11.1\beta 0.6\beta ^2)^1/7.`$ (35)
In each case, the small $`\beta `$ limit gives the corresponding Newtonian result. Figure 6 indicates that the T VII approximation is a rather good approximation to most EOSs without extreme softening at high densities, for $`M/R0.1`$ M/km. The EOSs with softening fall below this trajectory. Ravenhall & Pethick (1994) had suggested the expression
$$I_{RP}/MR^20.21/(12\beta )$$
(36)
as an approximation for the moment of inertia; however, we find that this expression is not a good overall fit, as shown in Figure 6.
For low-mass stars, none of the analytic approximations are suitable, and the moment of inertia deviates substantially from the behavior of an incompressible fluid. Although neutron stars of such small mass are unlikely to exist, it is interesting to examine the behavior of $`I`$ in the limit of small compactness, especially the suprising result that $`I/MR^20`$ as $`\beta 0`$. It is well known from the work of Baym, Bethe & Pethick (1971) that the adiabatic index of matter below nuclear density is near to, but less than 4/3. As the compactness parameter $`\beta `$ decreases, a greater fraction of the star’s mass lies below $`n_s`$. To the extent that these stars can be approximated as polytropes (i.e., having a constant polytropic index $`n`$), Table 4 shows how the quantity $`I/MR^2`$ varies with $`n`$. For a polytropic index of 3, corresponding to an adiabatic exponent of 4/3, $`I/MR^20.075`$, considerably lower than the value of 0.4 for an incompressible fluid. Calculations of matter at subnuclear density agree on the fact that the adiabatic exponent of matter further decreases with decreasing density, until the neutron drip point (near $`4.3\times 10^{11}`$ g cm<sup>-3</sup>) is approached and the exponent is near zero. Although the central densities of minimum mass neutron stars are about $`2\times 10^{14}`$ g cm<sup>-3</sup>, much of the mass of the star is at considerably lower density, unlike the situation for solar mass-sized neutron stars which are relatively centrally condensed. Thus, as $`\beta `$ decreases, the quantity $`I/MR^2`$ rapidly decreases, approaching the limiting value of zero as an effective polytropic index of nearly 5 is achieved.
Another interesting result from Figure 6 concerns the moments of inertia of strange quark matter stars. Such stars are relatively closely approximated by incompressible fluids, this behavior becoming exact in the limit of $`\beta 0`$. This could have been anticipated from the $`MR^3`$ behavior of the $`MR`$ trajectories for small $`\beta `$ strange quark matter stars as observed in Figure 2.
## 5 CRUSTAL FRACTION OF THE MOMENT OF INERTIA
A new observational constraint involving $`I`$ concerns pulsar glitches. Occasionally, the spin rate of a pulsar will suddenly increase (by about a part in $`10^6`$) without warning after years of almost perfectly predictable behavior. However, Link, Epstein & Lattimer (1999) argue that these glitches are not completely random: the Vela pulsar experiences a sudden spinup about every three years, before returning to its normal rate of slowing. Also, the size of a glitch seems correlated with the interval since the previous glitch, indicating that they represent self-regulating instabilities for which the star prepares over a waiting time. The angular momentum requirements of glitches in Vela imply that $`1.4`$% of the star’s moment of inertia drives these events.
Glitches are thought to represent angular momentum transfer between the crust and another component of the star. In this picture, as a neutron star’s crust spins down under magnetic torque, differential rotation develops between the stellar crust and this component. The more rapidly rotating component then acts as an angular momentum reservoir which occasionally exerts a spin-up torque on the crust as a consequence of an instability. A popular notion at present is that the freely spinning component is a superfluid flowing through a rigid matrix in the thin crust, the region in which dripped neutrons coexist with nuclei, of the star. As the solid portion is slowed by electromagnetic forces, the liquid continues to rotate at a constant speed, just as superfluid He continues to spin long after its container has stopped. This superfluid is usually assumed to locate in the star’s crust, which thus must contain at least 1.4% of the star’s moment of inertia.
The high-density boundary of the crust is naturally set by the phase boundary between nuclei and uniform matter, where the pressure is $`P_t`$ and the density $`n_t`$. The low-density boundary is the neutron drip density, or for all practical purposes, simply the star’s surface since the amount of mass between the neutron drip point and the surface is negligible. One can utilize equation (28) to determine the moment of inertia of the crust alone with the assumptions that $`P/c^2<<\rho `$, $`m(r)M`$, and $`\omega j\omega _R`$ in the crust. Defining $`\mathrm{\Delta }R`$ to be the crust thickness, that is, the distance between the surface and the point where $`P=P_t`$,
$$\mathrm{\Delta }I\frac{8\pi }{3}\frac{\omega _R}{\mathrm{\Omega }}_{R\mathrm{\Delta }R}^R\rho r^4e^\lambda 𝑑r\frac{8\pi }{3GM}\frac{\omega _R}{\mathrm{\Omega }}_0^{P_t}r^6𝑑P,$$
(37)
where $`M`$ is the star’s total mass and the TOV equation was used in the last step. In the crust, the fact that the EOS is of the approximate polytropic form $`PK\rho ^{4/3}`$ can be used to find an approximation for the integral $`r^6𝑑P`$, viz.
$$_0^{P_t}r^6𝑑PP_tR^6\left[1+\frac{2P_t}{n_tm_nc^2}\frac{(1+7\beta )(12\beta )}{\beta ^2}\right]^1.$$
(38)
For most neutron stars, the approximation equation (35) gives $`I`$ in terms of $`M`$ and $`R`$, and equation (30) gives $`\omega _R/\mathrm{\Omega }`$ in terms of $`I`$ and $`R`$, the quantity $`\mathrm{\Delta }I/I`$ can thus be cast as a function of $`M`$ and $`R`$ with the only dependences upon the EOS arising from the values of $`P_t`$ and $`n_t`$; there is no explicit dependence upon the EOS at any other density. However, the major dependence is mostly upon the value of $`P_t`$, since $`n_t`$ enters only as a correction. We then find
$$\frac{\mathrm{\Delta }I}{I}\frac{28\pi P_tR^3}{3Mc^2}\frac{(11.67\beta 0.6\beta ^2)}{\beta }\left[1+\frac{2P_t(1+5\beta 14\beta ^2)}{n_tm_bc^2\beta ^2}\right]^1.$$
(39)
In general, the EOS parameter $`P_t`$, in the units of MeV fm<sup>-3</sup>, varies over the range $`0.25<P_t<0.65`$ for realistic EOSs. The determination of this parameter requires a calculation of the structure of matter containing nuclei just below nuclear matter density that is consistent with the assumed nuclear matter EOS. Unfortunately, few such calculations have been performed. Like the fiducial pressure at and above nuclear density which appears in equation (5), $`P_t`$ should depend sensitively upon the behavior of the symmetry energy near nuclear density.
Since the calculation of the pressure below nuclear density has not been consistently done for most realistic EOSs, we arbitrarily choose $`n_t=0.07`$ fm<sup>-3</sup> and compare the approximation equation (39) with the results of full structural calculations in Figure 7. Two extreme values of $`P_t`$ were assumed in the full structural calculations to identify the core-crust boundary. Irrespective of this choice, the agreement between the analytical estimate equation (39) and the full calculations appears to be good for all EOSs, including ones with extreme softening. We also note that Ravenhall & Pethick (1994) developed a different, but nearly equivalent, analytic formula for the quantity $`\mathrm{\Delta }I/I`$ as a function of $`M,R,P_t`$ and $`\mu _t`$, where $`\mu _t`$ is the neutron chemical potential at the core-crust phase boundary. This prediction is also displayed in Figure 7.
Link, Epstein & Lattimer (1999) established a lower limit to the radius of the Vela pulsar by using equation (39) with $`P_t`$ at its maximum value and the glitch constraint $`\mathrm{\Delta }I/I0.014`$. A minimum radius can be found by combining this constraint with the largest realistic value of $`P_t`$ from any equation of state, namely about 0.65 MeV fm<sup>-3</sup>. Stellar models that are compatible with this constraint must fall to the right of the $`P_t=0.65`$ MeV fm<sup>-3</sup> contour in Figure 7. This imposes a constraint upon the radius, which is approximately equivalent to
$$R>3.9+3.5M/\mathrm{M}_{}0.08(M/\mathrm{M}_{})^2\mathrm{km}.$$
(40)
As shown in the figure, this constraint is somewhat more stringent than one based upon causality. Better estimates of the maximum value of $`P_t`$ should make this constraint more stringent.
## 6 BINDING ENERGIES
The binding energy formally represents the energy gained by assembling $`N`$ baryons. If the baryon mass is $`m_b`$, the binding energy is simply $`BE=Nm_bM`$ in mass units. However, the quantity $`m_b`$ has various interpretations in the literature. Some authors take it to be 939 MeV/$`c^2`$, the same as the neutron or proton mass. Others take it to be about 930 MeV/$`c^2`$, corresponding to the mass of C<sup>12</sup>/12 or Fe<sup>56</sup>/56. The latter choice would be more appropriate if $`BE`$ was to represent the energy released in by the collapse of a white-dwarf-like iron core in a supernova explosion. The difference in these definitions, 10 MeV per baryon, corresponds to a shift of $`10/9390.01`$ in the value of $`BE/M`$. This energy, $`BE`$, can be deduced from neutrinos detected from a supernova event; indeed, it might be the most precisely determined aspect of the neutrino signal.
Lattimer & Yahil (1989) suggested that the binding energy could be approximated as
$$BE1.510^{51}(M/\mathrm{M}_{})^2\mathrm{ergs}=0.084(M/\mathrm{M}_{})^2\mathrm{M}_{}.$$
(41)
Prakash et al. (1997) also concluded that such a formula was a reasonable approximation, based upon a comparison of selected non-relativistic potential and field-theoretical models. In Figure 8, this formula is compared to exact results, which shows that it is accurate at best to about $`\pm 20`$%. The largest deviations are for stars with extreme softening or large mass.
Here, we propose a more accurate representation of the binding energy:
$$BE/M0.6\beta /(10.5\beta ),$$
(42)
which incorporates some radius dependence. Thus, the observation of supernova neutrinos, and the estimate of the total radiated neutrino energy, will yield more accurate information about $`M/R`$ than about $`M`$ alone.
In the cases of the incompressible fluid and the Buchdahl solution, analytic results for the binding energy can be found:
$`BE_{Inc}/M`$ $`=`$ $`{\displaystyle \frac{3}{4\beta }}\left({\displaystyle \frac{\mathrm{sin}^1\sqrt{2\beta }}{\sqrt{2\beta }}}\sqrt{12\beta }\right)1{\displaystyle \frac{3\beta }{5}}+{\displaystyle \frac{9\beta ^2}{14}}+{\displaystyle \frac{5\beta ^3}{6}}+\mathrm{};`$ (43)
$`BE_{Buch}/M`$ $`=`$ $`(11.5\beta )(12\beta )^{1/2}(1\beta )^11{\displaystyle \frac{\beta }{2}}+{\displaystyle \frac{\beta ^2}{2}}+{\displaystyle \frac{3\beta ^3}{4}}+\mathrm{}.`$ (44)
In addition, an expansion for the T VII solution can be found:
$`BE_{TVII}/M{\displaystyle \frac{11\beta }{21}}+{\displaystyle \frac{7187\beta ^2}{18018}}+{\displaystyle \frac{68371\beta ^3}{306306}}+\mathrm{}.`$ (45)
The exact results for the three analytic solutions of Einstein’s equations, as well as the fit of equation (42), are compared to some EOSs in Figure 9. It can be seen that for stars without extreme softening both the T VII and Buch solutions are rather realistic. However, for EOSs with softening, the deviations from this can be substantial. Thus, until information about the existence of softening in neutron stars is available, the binding energy alone provides only limited information about the star’s structure or mass.
## 7 SUMMARY AND OUTLOOK
We have demonstrated the existence of a strong correlation between the pressure near nuclear saturation density inside a neutron star and the radius which is relatively insensitive to the neutron star’s mass and equation of state for normal neutron stars. In turn, the pressure near the saturation density is primarily determined by the isospin properties of the nucleon-nucleon interaction, specifically, as reflected in the density dependence of the symmetry energy, $`S_v(n)`$. This result is not sensitive to the other nuclear parameters such as $`K_s`$, the nuclear incompressibility parameter, or $`K_s^{}`$, the skewness parameter. This is important, because the value of the symmetry energy at nuclear saturation density and the density dependence of the symmetry energy are both difficult to determine in the laboratory. Thus, a measurement of a neutron star’s radius would yield important information about these quantities.
Any measurement of a radius will have some intrinsic uncertainty. In addition, the empirical relation we have determined between the pressure and radius has a small uncertainty. It is useful to display how accurately the equation of state might be established from an eventual radius measurement. This can be done by inverting equation (5), which yields
$$P(n)[R_M/C(n,M)]^4.$$
(46)
The inferred ranges of pressures, as a function of density and for three possible values of $`R_{1.4}`$, are shown in Figure 10. It is assumed that the mass is 1.4 M, but the results are relatively insensitive to the actual mass. Note from Table 3 that the differences between $`C`$ for 1 and 1.4 M are typically less than the errors in $`C`$ itself. The light shaded areas show the pressures including only errors associated with $`C`$. The dark shaded areas show the pressures when a hypothetical observational error of 0.5 km is also taken into account. These results suggest that a useful restriction to the equation of state is possible if the radius of a neutron star can be measured to an accuracy better than about 1 km.
The reason useful constraints might be obtained from just a single measurement of a neutron star radius, rather than requiring a series of simultaneous mass-radii measurements as Lindblom (1992) proposed, stems from the fact that we have been able to establish the empirical correlation, equation (5). In turn, it appears that this correlation exists because most equations of state have slopes $`d\mathrm{ln}P/d\mathrm{ln}n2`$ near $`n_s`$.
The best prospect for measuring a neutron star’s radius may be the nearby object RX J185635-3754. It is anticipated that parallax information for this object will be soon available (Walter, private communication). In addition, it may be possible to identify spectral lines with the Chandra and XMM X-ray facilities that would not only yield the gravitational redshift, but would identify the atmospheric composition. Not only would this additional information reduce the uncertainty in value of $`R_{\mathrm{}}`$, but, both the mass and radius for this object might thereby be determined. It is also possible that an estimate of the surface gravity of the star can be found from further comparisons of observations with atmospheric modelling, and this would provide a further check on the mass and radius.
We have presented simple expressions for the moment of inertia, the binding energy, and the crustal fraction of the moment of inertia for normal neutron stars which are largely independent of the EOS. If the magnitudes of observed glitches from Vela are connected with the crustal fraction of moment of inertia, the formula we derived establishes a more stringent limit on the radius than causality.
We thank A. Akmal, V. R. Pandharipande and J. Schaffner-Bielich for making the results of their equation of state calculations available to us in tabular form. This work was supported in part by the USDOE grants DOE/DE-FG02-87ER-40317 & DOE/DE-FG02-88ER-40388.
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TABLE 1
EQUATIONS OF STATE
| Symbol | Reference | Approach | Composition |
| --- | --- | --- | --- |
| FP | Friedman & Pandharipande (1981) | Variational | np |
| PS | Pandharipande & Smith (1975) | Potential | n$`\pi ^0`$ |
| WFF(1-3) | Wiringa, Fiks & Fabrocine (1988) | Variational | np |
| AP(1-4) | Akmal & Pandharipande (1998) | Variational | np |
| MS(1-3) | Müller & Serot (1996) | Field Theoretical | np |
| MPA(1-2) | Muẗher, Prakash & Ainsworth (1987) | Dirac-Brueckner HF | np |
| ENG | Engvik et al. (1996) | Dirac-Brueckner HF | np |
| PAL(1-6) | Prakash, Ainsworth & Lattimer (1988) | Schematic Potential | np |
| GM(1-3) | Glendenning & Moszkowski (1991) | Field Theoretical | npH |
| GS(1-2) | Glendenning & Schaffner-Bielich (1999) | Field Theoretical | npK |
| PCL(1-2) | Prakash, Cooke & Lattimer (1995) | Field Theoretical | npHQ |
| SQM(1-3) | Prakash, Cooke & Lattimer (1995) | Quark Matter | Q $`(u,d,s)`$ |
NOTES.—- Approach refers to the underlying theoretical technique. Composition refers to strongly interacting components (n=neutron, p=proton, H=hyperon, K=kaon, Q=quark); all models include leptonic contributions.
TABLE 2
PARAMETERS FOR SELF-BOUND STRANGE QUARK STARS
| Model | $`B`$ (MeV fm$`{}_{}{}^{3})`$ | $`m_s`$ (MeV) | $`\alpha _c`$ |
| --- | --- | --- | --- |
| SQM1 | 94.92 | 0 | 0 |
| SQM2 | 64.21 | 150 | 0.3 |
| SQM3 | 57.39 | 50 | 0.6 |
NOTES.—- Numerical values employed in the MIT bag model as described in Fahri & Jaffe (1984).
TABLE 3
THE QUANTITY $`C(n,M)`$ OF EQUATION 5
| $`n`$ | 1 M | 1.4 M |
| --- | --- | --- |
| $`n_s`$ | $`9.53\pm 0.32`$ | $`9.30\pm 0.60`$ |
| $`1.5n_s`$ | $`7.14\pm 0.15`$ | $`7.00\pm 0.31`$ |
| $`2n_s`$ | $`5.82\pm 0.21`$ | $`5.72\pm 0.25`$ |
NOTES.—- The quantity $`C(n,M)`$, in units of km fm<sup>3/4</sup> MeV<sup>-1/4</sup>, which relates the pressure (evaluated at density $`n`$) to the radius of neutron stars of mass $`M`$. The errors are standard deviations.
TABLE 4
MOMENTS OF INERTIA FOR POLYTROPES
| Index $`n`$ | $`I/MR^2`$ | Index $`n`$ | $`I/MR^2`$ |
| --- | --- | --- | --- |
| 0 | 0.4 | 3.5 | 0.045548 |
| 0.5 | 0.32593 | 4.0 | 0.022573 |
| 1.0 | 0.26138 | 4.5 | 0.0068949 |
| 1.5 | 0.20460 | 4.8 | 0.0014536 |
| 2.0 | 0.15485 | 4.85 | 0.00089178 |
| 2.5 | 0.11180 | 4.9 | 0.0004536 |
| 3.0 | 0.075356 | 5.0 | 0 |
NOTES.—- The quantity $`I/MR^2`$ for polytropes, which satisfy the relation $`P=K\rho ^{1+1/n}`$ ($`\rho `$ is the mass-energy density), as a function of the polytropic index $`n`$.
## FIGURE CAPTIONS
FIG. 1.—-The pressure-density relation for a selected set of EOSs contained in Table 1. The pressure is in units of MeV fm<sup>-3</sup> and the density is in units of baryons per cubic fermi. The nuclear saturation density is approximately $`0.16`$ fm<sup>-3</sup>.
FIG. 2.—-Mass-radius curves for several EOSs listed in Table 1. The left panel is for stars containing nucleons and, in some cases, hyperons. The right panel is for stars containing more exotic components, such as mixed phases with kaon condensates or strange quark matter, or pure strange quark matter stars. In both panels, the lower limit causality places on $`R`$ is shown as a dashed line, a constraint derived from glitches in the Vela pulsar is shown as the solid line labelled $`\mathrm{\Delta }I/I=0.014`$, and contours of constant $`R_{\mathrm{}}=R/\sqrt{12GM/Rc^2}`$, are shown as dotted curves. In the right panel, the theoretical trajectory of maximum masses and radii for pure strange quark matter stars is marked by the dot-dash curve labelled $`R=1.85R_s`$.
FIG. 3.—- Empirical relation between pressure, in units of MeV fm<sup>-3</sup>, and radius, in km, for EOSs listed in Table 1. The upper panel shows results for 1 M (gravitational mass) stars; the lower panel is for 1.4 M stars. The different symbols show values of $`RP^{1/4}`$ evaluated at three fiducial densities.
FIG. 4.—- Left panel: Mass-radius curves for selected PAL (Prakash, Ainsworth & Lattimer 1988) forces showing the sensitivity to symmetry energy. The left panel shows variations arising from different choices of $`S_v`$, the symmetry energy evaluated at $`n_s`$; the right panel shows variations arising from different choices of $`S_v(n)`$, the density dependent symmetry energy. In this figure, the shorthand $`u=n/n_s`$ is used.
FIG. 5.—- Profiles of mass-energy density ($`\rho `$), relative to central values ($`\rho _c`$), in neutron stars for several EOSs listed in Table 1. For reference, the thick black lines show the simple quadratic approximation $`1(r/R)^2`$.
FIG. 6.—- The moment of inertia $`I`$, in units of $`MR^2`$, for several EOSs listed in Table 1. The curves labelled “Inc”, “T VII”, “Buch” and “RP” are for an incompressible fluid, the Tolman (1939) VII solution, the Buchdahl (1967) solution, and an approximation of Ravenhall & Pethick (1994), respectively. The inset shows details of $`I/MR^2`$ for $`M/R0`$.
FIG. 7.—- Mass-radius curves for selected EOSs from Table 1, comparing theoretical contours of $`\mathrm{\Delta }I/I=0.014`$ from approximations developed in this paper, labelled “LP”, and from Ravenhall & Pethick (1994), labelled “RP”, to numerical results (solid dots). Two values of $`P_t`$, the transition pressure demarking the crust’s inner boundary, which bracket estimates in the literature, are employed. The region to the left of the $`P_t=0.65`$ MeV fm<sup>-3</sup> curve is forbidden if Vela glitches are due to angular momentum transfers between the crust and core, as discussed in Link, Epstein & Lattimer (1999). For comparison, the region excluded by causality alone lies to the left of the dashed curve labelled “causality” as determined by Lattimer et al. (1990) and Glendenning (1992).
FIG. 8.—- The binding energy of neutron stars as a function of stellar gravitational mass for several EOSs listed in Table 1. The predictions of equation (41), due to Lattimer & Yahil (1989), are shown by the line labelled “LY” and the shaded region.
FIG. 9.—- The binding energy per unit gravitational mass as a function of compactness for the EOSs listed in Table 1. Solid lines labelled “Inc”, “Buch” and “T VII” show predictions for an incompressible fluid, the solution of Buchdahl (1967), and the Tolman (1939) VII solution, respectively. The shaded region shows the prediction of equation (42).
FIG. 10.—The pressures inferred from the empirical correlation equation (5), for three hypothetical radius values (10, 12.5 and 15 km) overlaid on the pressure-density relations shown in Figure 1. The light shaded region takes into account only the uncertainty associated with $`C(n,M)`$; the dark shaded region also includes a hypothetical uncertainty of 0.5 km in the radius measurement. The neutron star mass was assumed to be 1.4 M.
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# Radiative Corrections to Electron-Proton Scattering
## I INTRODUCTION
Electron scattering at intermediate and high energies has been one of the most useful means of investigating nuclear structure for over forty years. With the advent of CW accelerators and high resolution detectors such as MAMI and TJNAF it has become clear that one must have an accurate estimate of the radiative corrections if meaningful cross sections are to be obtained from the experimental measurements. Depending on the experimental conditions – initial beam energy, momentum transfer, and detector resolution or missing mass for the observed particles – the radiative corrections can be as large as $`30\%`$ of the uncorrected cross section. To obtain cross sections which are accurate to 1%, one must then know the radiative correction to 3%.
The theoretical expression for the radiative correction which has been used in the analysis of almost all single arm elastic electron scattering experiments with beam energies below approximately $`25`$ GeV (for which $`W`$ and $`Z`$ exchange are in general not significant) is that given originally by Tsai , in connection with experiments at Stanford, SLAC and CEA. That work involved approximations that were both purely mathematical (made in performing the integrations needed to evaluate the inelastic cross section) and approximations denoted here as “soft-photon approximations” that are directly related to the physics in that the effect of proton structure was neglected; in considering the proton legs, only the soft virtual (infrared) photon contribution is calculated exactly - approximations are made in the hard virtual photon (non-infrared) contribution. In particular, the proton structure is neglected by setting the photon momentum square $`k^2=0`$ in the proton form factor, $`F(k^2)`$, thus simplifying the calculation considerably.
The purpose of the present paper is to study the radiative correction to elastic electron-proton scattering including the internal structure of the nucleon. For this we have considered a simple model in which the proton current is taken to have the usual on-shell form. The model dependence of the radiative correction is clearly an important question for the analysis of electron scattering experiments at the 1% level. This work is an initial study to examine the size of internal structure effects.
The present calculation differs from that of Tsai , in three substantive aspects. First, we evaluate the inelastic cross section (emission of soft real photons) without any approximation; the relevant integrals have been given in closed form by t’Hooft and Veltman . In fact, these exact expressions are simpler in form than the approximate ones given in and . We note in particular that in the limit of the target mass $`M\mathrm{},`$ corresponding to a static Coulomb potential, we obtain exactly the result first given by Schwinger . Second, in the evaluation of the contribution of the box and crossed box diagrams to the elastic cross section we make a less drastic approximation than that made in . Specifically, in the integrands corresponding to the relevant matrix elements, $`M_2`$ and $`M_3`$ (Eqs. (A8) and (A11)), we make a soft photon approximation (setting $`k=0`$ or $`k=q`$) in the numerator (as in ), but not in the denominators. Again, the required integrals (scalar four-point functions) have been given in ; the resulting expressions are again considerably simpler than those obtained in , where the soft-photon approximation is also made in the denominators of $`M_2`$ and $`M_3`$. Finally, in evaluating the proton vertex correction, we have made no soft photon approximation for the virtual photon (as was done in ) and have included form factors for the proton, taking the proton current to be that indicated below in (5).
The organization of the paper is as follows: In Sec. II we discuss questions concerning the electromagnetic nuclear current operator used in this analysis. In Sec. III we give details of the calculation of the matrix elements and cross section for elastic scattering, retaining terms of order $`\alpha `$ relative to the Rosenbluth (one photon exchange) cross section for elastic scattering. Integrals needed for the evaluation of the various matrix elements are written explicitly and expressed in closed form in terms of Spence functions (dilogarithms). Details are given in the Appendices.
In Sec. IV we consider the inelastic cross section in detail; as with the elastic cross section given in Sec. III, the result is expressed in closed form in terms of Spence functions. In Sec. V we add the elastic and inelastic cross sections, giving both an analytic expression and a numerical evaluation of the radiative correction for various values of the pertinent parameters, (initial beam energy, final electron detector resolution, and target nucleus). We compare the values of the radiative correction calculated here with those given in and .
## II Electromagnetic nucleon current operator
We essentially follow in this paper the convention of Björken and Drell . The metric used is defined by
$$p_ip_j=ϵ_iϵ_j𝐩_i𝐩_j$$
(1)
Further, $`\alpha =e^2/4\pi =1/137.036`$; $`m`$ is the electron rest mass; $`M`$ is the target nucleus rest mass; $`Z`$ the charge of the target nucleus; $`\kappa `$ the anomalous magnetic moment of the proton; $`p_1`$ and $`p_3`$ the initial and final electron four-momenta, respectively; $`p_2`$ and $`p_4`$ the initial and final target nucleus four-momenta, respectively; $`q=p_1p_3=p_4p_2`$ is the four-momentum transfer to the target nucleus for elastic scattering. It will prove useful to define, in addition,
$$\rho =p_4+p_2,\rho _m=p_1+p_3,$$
(2)
from which $`\rho ^2=q^2+4M^2`$ and $`\rho _m^2=q^2+4m^2`$. Further, we define
$`x`$ $`=`$ $`(\rho +\rho _1)\mathrm{}(\rho \rho _1)=(\rho +\rho _1)^2\mathrm{}4M^2`$ (3)
$`x_m`$ $`=`$ $`(\rho _m+\rho _1)\mathrm{}(\rho _m\rho _1)=(\rho _m+\rho _1)^2\mathrm{}4m^2`$ (4)
with $`\rho _1^2=q^2`$. In the lab system we have: $`p_1=(ϵ_1,𝐩_1);p_3=(ϵ_3,𝐩_3);p_2=(M,0);p_4=(M+\omega ,𝐪);\omega =q^2/2M`$.
With the aim of presenting expressions which correspond to the experimental conditions of high energy electron scattering, we neglect, in the final expressions given in this paper, terms of relative orders $`m^2/ϵ^2,`$ $`m^2/(q^2)`$, and $`m^2/M^2`$. Neglect of these terms defines our high energy approximation. No assumption is made, however, with regard to the magnitudes of $`M/ϵ_1,`$ $`M/ϵ_3,`$ or $`M^2/(q^2).`$
At low momentum transfer the internal structure of the nucleon can safely be neglected in the determination of the radiative corrections in electron-nucleus scattering. However, with increasing energies and momenta this is in general no longer the case. One of the objectives of this paper is to investigate this in a model for the e.m. interaction of a non-pointlike nucleon. The most general e.m. off-shell nucleon vertex can be characterized by 6 invariant functions . As the most simple model we may consider a vector dominance-like model for the nucleon current, characterized by only two form factors which depend only on the four-momentum square of the photon. It is given by
$$\mathrm{\Gamma }_\mu =F_1(q^2)\gamma _\mu +\kappa F_2(q^2)\frac{i\sigma _{\mu \nu }q^\nu }{2M},$$
(5)
where the form factors $`F_1(q^2)`$ and $`F_2(q^2)`$ are taken to have a monopole or dipole form:
$$F_1(q^2)=F_2(q^2)=\left(\frac{\mathrm{\Lambda }^2}{q^2\mathrm{\Lambda }^2}\right)^n,n=1\text{ or\hspace{0.17em}2 }$$
(6)
with $`\mathrm{\Lambda }`$ being a constant of the order of $`1`$GeV/$`c.`$ Furthermore, $`q=p^{}p`$, $`p`$ and $`p^{}`$ being the momentum of the initial and final nucleon. Although the quantitative predictions of the radiative corrections are expected in general to be dependent on the details of the nucleon model assumed, one should already be able to see most of the salient features in the present model study. In particular, identifying regions in phase space where the finite size of the nucleon may play an important role in the size of radiative corrections can be important. In this way one may hope to get some feeling on the reliability of neglecting the internal structure of the nucleon as is usually done. The present study is intended as a first exploration of the sensitivity on the non-pointlike nature of the e.m. hadronic current. As in , although we are primarily interested in electron-proton scattering, the radiative corrections studied here can also be applied to electron-nucleus scattering, with appropriate changes in $`F_1,`$ $`F_2,\kappa ,`$ and $`M`$. However, even in the case of electron-proton scattering, the factor $`Z`$ is convenient for identifying the contributions from the various diagrams.
It should be noted that the dressed vertex function, $`\stackrel{~}{\mathrm{\Lambda }}_\mu ,`$ with Eq. (5) as e.m. current operator containing the form factors $`F_n`$, satisfies a Ward-Takahashi identity
$$q^\mu \stackrel{~}{\mathrm{\Lambda }}_\mu =F_1(q^2)\left[S^1(p^{})S^1(p)\right]$$
(7)
where $`S`$ is the dressed nucleon propagator. As a direct consequence of (7), one gets for on-mass-shell nucleons, the current conservation
$$q^\mu <p^{}|\stackrel{~}{\mathrm{\Lambda }}_\mu |p>=0.$$
(8)
Obviously, the radiative corrections will in general be sensitive to the choice of the e.m. current. Although interesting in its own right we will not address in this paper the issue of the dependence of the predictions on this ambiguity.
In the study of radiative corrections we may distinguish between the elastic and inelastic contributions, the latter being the real soft photon emission processes from both the electron and hadron. The elastic electron cross section can be determined immediately from the total scattering amplitude $``$ through the well-known expression
$`d\sigma `$ $`=`$ $`{\displaystyle \frac{mM}{\sqrt{(p_1p_2)^2m^2M^2}}}{\displaystyle \underset{spins}{}}{\displaystyle \left|\right|^2(2\pi )^4\delta ^4(p_4+p_3p_2p_1)}`$ (10)
$`\times {\displaystyle \frac{md^3p_3}{(2\pi )^3ϵ_3}}{\displaystyle \frac{Md^3p_4}{(2\pi )^3ϵ_4}},`$
For single-arm experiments with unpolarized electrons in which the final proton is not observed, $`d\sigma `$ must be averaged over initial spins, summed over final spins, and integrated over the final proton four-momentum. Up to order $`\alpha ^2`$ we have for the total scattering amplitude
$$=\underset{n=1}{\overset{6}{}}M_n,$$
(11)
where the various terms correspond to the Feynman graph contributions shown in Fig. 1. $`M_1`$ is the matrix element for the one-photon exchange diagram
$$M_1=Ze^2\overline{u}(p_3)\gamma _\mu u(p_1)\frac{(i)}{q^2+iϵ}\overline{u}(p_4)\mathrm{\Gamma }^\mu (q^2)u(p_2).$$
(12)
Its square gives, for high energy electrons, the Rosenbluth cross section:
$$\frac{d\sigma _0}{d\mathrm{\Omega }}=\frac{\alpha ^2\mathrm{cos}^2\frac{\theta }{2}\left[\left(F_1^2\frac{\kappa ^2q^2}{4M^2}F_2^2\right)\frac{q^2}{2M^2}\left(F_1+\kappa F_2\right)^2\mathrm{tan}^2\frac{\theta }{2}\right]}{4ϵ_1^2\eta \mathrm{sin}^4\left(\theta /2\right)},$$
(13)
where $`\eta `$ is the lab system recoil factor: For $`ϵ_1>>m,ϵ_3>>m,\eta ϵ_1/ϵ_31+(ϵ_1/M)(1\mathrm{cos}\theta )`$ with $`\theta `$ being the electron scattering angle. We note, in particular, that $`1\eta x`$. Furthermore, $`M_2`$ and $`M_3`$ are the matrix elements for the box and crossed box (two-photon exchange) diagrams. $`M_4`$ is the vacuum polarization diagram (only an electron-positron loop is indicated in the figure, but the contribution from higher mass lepton loops can be included without difficulty- see (A13) -(A17)). $`M_5`$ is the electron vertex correction, and $`M_6`$ is the proton vertex correction. For completeness, we list in Appendix A the explicit expressions for the various Feynman diagrams shown in Figs. 1 and 2.
## III Elastic Cross Section
To evaluate the various one-loop corrections to Eq. (11) some tedious algebra has to be carried out. We outline the procedure used to evaluate the matrix elements needed for the radiative correction to the elastic cross section, $`M_2`$ through $`M_6`$, (Eqs. (A1)-(A15)).
### A Proton vertex correction
We begin with the matrix element for the proton vertex correction, $`M_6`$ , given by Eqs. (A2) and (A5); the much simpler matrix element for the electron vertex correction, $`M_5`$, (Eq. (A1)), can be deduced quite easily from that. In (A5), each of the three $`\mathrm{\Gamma }`$’s, given by Eq. (5), contains a term with $`\gamma _\mu `$ (which we denote by $`g`$) and a term with $`\sigma _{\mu \nu }`$ (which we denote by $`s`$). The proton vertex correction $`\mathrm{\Lambda }^\mu (p_4,p_2)`$ then consists of eight terms, which we represent symbolically by $`ggg,gsg,gss,`$ etc. As may be seen after rationalizing the propagators, the $`k`$ dependence of the numerators for $`ggg,gsg,\mathrm{}`$ is such that there are at most four factors of the form $`k/`$. Moreover, the terms with three or four factors $`k/`$ may, with only a minimum of algebra, be written so that two of these factors are adjacent, giving $`k/k/=k^2`$. Although the calculation can equally well be carried out with $`F_1`$ and $`F_2`$ distinct functions, we assume $`F_1=`$ $`F_2=F`$, which simplifies the algebra. The terms $`ggg,gsg,\mathrm{}`$ can then be expressed in terms of the integrals
$$\{I_0;I_\mu ;I_{\mu \nu };J_0;J_\mu ;J_{\mu \nu };K_0\}=\frac{d^4k}{(2\pi )^4}F^2(k^2)\{1;k_\mu ;k_\mu k_\nu ;k^2;k_\mu k^2;k_\mu k_\nu k^2;(k^2)^2\}/D(\lambda ^2)$$
(14)
where
$$D(\lambda ^2)=(k^2\lambda ^2+iϵ)(k^22kp_2+iϵ)(k^22kp_4+iϵ)$$
(15)
For form factors having the form given in (6), the integrals in (14) could all be evaluated as indicated for three-point functions in , Sec.5, and , Appendix E. However, in the interest of obtaining a relatively compact analytic expression in closed form, we have used an alternative procedure. As given here in Appendix B, the integrals may be expressed in terms of their moments, defined by Eqs. (B4)-(B6) and (B14). After straightforward though somewhat tedious algebra, the terms $`ggg,gsg,\mathrm{}`$ are then expressed in terms of these moments. Next, for form factors of the form given in (6), we show that all of the moments may be expressed in terms of three functions, $`\varphi _k`$, which obey a three-term inhomogenous recursion, and this is used for their evaluation. Finally, we note from (B67) - (B98) that the terms $`ggg,gsg,\mathrm{}`$ may be usefully grouped by writing them in the form
$$(g+s)g(g+s)=F(q^2)\left[G_1(q^2)\gamma _\mu +G_2(q^2)\frac{i\sigma _{\mu \nu }q^\nu }{2M}\right]$$
(16)
and
$$(g+s)s(g+s)=\kappa F(q^2)\left[X_1(q^2)\gamma _\mu +X_2(q^2)\frac{i\sigma _{\mu \nu }q^\nu }{2M}\right]$$
(17)
We note in the expressions for $`ggg`$, $`gsg`$,… in Appendix B that the infrared divergent terms are all contained solely within $`ggg`$ and $`gsg`$. These are the terms with a factor $`\varphi _1(\lambda ^2)`$ in (B67) and (B71). Since these are precisely the terms which are retained in the proton vertex correction in , we separate them for the purpose of comparison with that work, writing $`M_6`$ in the form
$$M_6=M_6^{(0)}+M_6^{(1)}$$
(18)
where
$$M_6^{(0)}=\frac{\alpha Z^2}{2\pi }(2M^2q^2)\varphi _1(\lambda ^2)M_1$$
(19)
The function $`\varphi _1(\lambda ^2)`$, defined by (B21), (B30) and (B31) and evaluated in (B40) and (B44), is simply related to the function $`K(p_2,p_4)`$ defined in by
$`K(p_i,p_j)={\displaystyle \frac{2p_ip_j}{i\pi ^2}}{\displaystyle \frac{d^4k}{(k^2\lambda ^2+iϵ)(k^22kp_i+iϵ)(k^22kp_j+iϵ)}}`$viz.,
$$K(p_2,p_4)=2p_2p_4\varphi _1(\lambda ^2)$$
(20)
### B Proton self energy correction
We next consider the contribution of the proton self energy diagrams. It is given by $`\mathrm{\Sigma }^{^{}}`$ where
$$\mathrm{\Sigma }^{^{}}=\frac{1}{4}Tr\left[\frac{\mathrm{\Sigma }}{p/}\right]$$
(21)
in which $`\mathrm{\Sigma }`$ is the lowest order self-energy contribution
$$\mathrm{\Sigma }=ie^2\frac{d^4k}{(2\pi )^4}\frac{1}{k^2\lambda ^2+iϵ}\mathrm{\Gamma }^\nu (k)\frac{1}{(p/k/M+iϵ)}\mathrm{\Gamma }_\nu (k)$$
(22)
Using the Ward-Takahashi identity Eq. (7) we find that
$`\mathrm{\Sigma }^{^{}}=G_1(0)`$where $`G_1(q^2)`$ is the coefficient of $`\gamma _\mu `$ in $`ggg`$ (B67). Explicitly, $`G_1(0)`$ is given by Eqs.(B36) - (B39) and (B45). The addition of this contribution to the lowest order vertex correction modifies the expressions for $`(g+s)g(g+s)`$ and $`(g+s)s(g+s)`$ given above in (16) and (17) so that we now have
$$\overline{(g+s)g(g+s)}=F(q^2)\left[\left(G_1(q^2)G_1(0)\right)\gamma _\mu +G_2(q^2)\frac{i\sigma _{\mu \nu }q^\nu }{2M}\right]$$
(23)
and
$$\overline{(g+s)s(g+s)}=\kappa F(q^2)\left[X_1(q^2)\gamma _\mu +\left(X_2(q^2)G_1(0)\right)\frac{i\sigma _{\mu \nu }q^\nu }{2M}\right]$$
(24)
We then have, for the matrix element including self energy diagrams,
$$\overline{M}_6=\overline{M}_6^{(0)}+\overline{M}_6^{(1)}$$
(25)
where
$$\overline{M}_6^{(0)}=\frac{\alpha Z^2}{2\pi }\left[K(p_2,p_4)+K(p_2,p_2)\right]M_1$$
(26)
which is the expression given in , eq. (II.12). The infrared divergent part of these terms is cancelled exactly by the infrared divergent terms in the inelastic cross section. The contribution of the matrix element $`M_6^{(1)}`$, which depends on the proton form factor, will be considered after we write the electron vertex and box diagram corrections.
### C Electron vertex correction
The electron vertex correction, $`M_5`$, may be obtained directly from the proton vertex correction, $`M_6`$. The expression $`\mathrm{\Lambda }_\mu (p_3,p_1)`$, Eq. (A3), follows from $`\mathrm{\Lambda }^\mu (p_4,p_2)`$ if we retain only the term $`ggg`$, set $`F=1`$, replace $`p_2,`$ $`p_4`$ and $`M`$ by $`p_1,`$ $`p_3`$ and $`m`$, and, after performing the integrations, take the limit $`\mathrm{\Lambda }\mathrm{}`$ (note Eq.(6)). Note that $`\rho =p_4+p_2`$ is then replaced by $`\rho _m=p_3+p_1`$, and $`x,`$ defined in section II, is replaced by $`x_m`$. Details are given in Appendix C. We find
$$ggg=G_1^{(e)}(q^2)\gamma _\mu +G_2^{(e)}(q^2)\frac{i\sigma _{\mu \nu }q^\nu }{2m}$$
(27)
where
$$G_1^{(e)}(q^2)=\frac{\alpha }{2\pi }\left\{K(p_1,p_3)+\left(\frac{3\rho _1^2+8m^2}{2\rho _m\rho _1}\right)\mathrm{ln}x_m+\frac{1}{4}+\frac{1}{2}\mathrm{ln}\left(\frac{\mathrm{\Lambda }^2}{m^2}\right)\right\}$$
(28)
$$G_2^{(e)}(q^2)=\frac{\alpha }{2\pi }\left\{\frac{2m^2}{\rho _m\rho _1}\mathrm{ln}x_m\right\}$$
(29)
Adding the contribution of the electron self energy diagrams gives
$$\overline{ggg}=\left(G_1^{(e)}(q^2)G_1^{(e)}(0)\right)\gamma _\mu +G_2^{(e)}(q^2)\frac{i\sigma _{\mu \nu }q^\nu }{2m}$$
(30)
from which
$$\overline{ggg}=\frac{\alpha }{2\pi }\left\{\left(K(p_1,p_3)+K(p_1,p_1)+\frac{3\rho _1^2+8m^2}{2\rho _m\rho _1}\mathrm{ln}x_m2\right)\gamma _\mu +\left(\frac{2m^2}{\rho _m\rho _1}\mathrm{ln}x\right)\frac{i\sigma _{\mu \nu }q^\nu }{2m}\right\}$$
(31)
For large momentum transfers, $`q^2m^2`$, this reduces to
$$\overline{ggg}=\frac{\alpha }{2\pi }\left\{K(p_1,p_3)+K(p_1,p_1)+\frac{3}{2}\mathrm{ln}\left(\frac{q^2}{m^2}\right)2\right\}\gamma _\mu $$
(32)
Comparing (12) and (A1), we then have, for $`q^2m^2`$,
$$\overline{M}_5=\frac{\alpha }{2\pi }\left\{K(p_1,p_3)+K(p_1,p_1)+\frac{3}{2}\mathrm{ln}\left(\frac{q^2}{m^2}\right)2\right\}M_1$$
(33)
which is the expression given in , eq.(II.5). We note that the infrared divergence is contained entirely within the terms $`K(p_1,p_3)+K(p_1,p_1).`$ The infrared divergent part of these terms is cancelled exactly by the infrared divergent terms in the inelastic cross section.
### D Box and crossed-box diagrams
The matrix elements for the box and crossed-box diagrams, $`M_2`$ and $`M_3`$, are given in (A8) and (A11). After rationalizing the propagators, the required integrals can, for form factors of the form (6), all be written in terms of four-point functions; in principle they can be evaluated using , Sec. 6, and , Appendix E. For the present work, however, we have chosen to evaluate these matrix elements in an approximate manner, but one which is less drastic than that employed in . We note first in $`M_2`$ and $`M_3`$ that the integrands in $`M_2`$ and $`M_3`$ have two infrared divergent factors, $`[(k^2\lambda ^2+iϵ)((kq)^2\lambda ^2+iϵ)]^1`$. The integrands are thus peaked when either of the two exchanged photons is soft, and become divergent when $`k0`$ or when $`kq`$. We therefore evaluate the numerators in $`M_2`$ and $`M_3`$ at these two points but make no changes to the denominators. A simple calculation shows that in fact that each of the numerators has the same value for $`k=0`$ as for $`k=q`$, viz., $`4ip_1p_2q^2M_1`$ in the case of $`M_2`$ and $`4ip_3p_2q^2M_1`$ in the case of $`M_3`$. We then take this factor outside of the integral and are left with a scalar four-point function to evaluate. The result has been given in , Sec. 6 and Appendix E (b) and is expressed simply in terms of logarithms:
$$M_2=\frac{\alpha Z}{\pi }\frac{ϵ_1}{|\text{p}_\mathrm{𝟏}|}\mathrm{ln}\left(\frac{ϵ_1+|\text{p}_\mathrm{𝟏}|}{m}\right)\mathrm{ln}\left(\frac{q^2}{\lambda ^2}\right)M_1$$
(34)
and
$$M_3=\frac{\alpha Z}{\pi }\frac{ϵ_3}{|\text{p}_\mathrm{𝟑}|}\mathrm{ln}\left(\frac{ϵ_3+|\text{p}_\mathrm{𝟑}|}{m}\right)\mathrm{ln}\left(\frac{q^2}{\lambda ^2}\right)M_1$$
(35)
By contrast, in , in addition to the approximation just described, a soft-photon approximation is made in the infrared denominators: Specifically, when $`k=0`$ the factor $`(kq)^2\lambda ^2`$ is set equal to $`q^2\lambda ^2`$ and when $`k=q`$ the factor $`k^2\lambda ^2`$ is set equal to $`q^2\lambda ^2`$, thus giving two terms and reducing the four-point function to three-point functions:
$$M_2=\frac{\alpha Z}{2\pi }\left[K(p_2,p_1)+K(p_4,p_3)\right]M_1$$
(36)
and
$$M_3=\frac{\alpha Z}{2\pi }\left[K(p_2,p_3)+K(p_4,p_1)\right]M_1$$
(37)
(Note , eqs.(II.9) and (II.11)). The infrared divergent terms (those with a factor $`\mathrm{ln}`$ $`\lambda ^2`$) are, for $`M_2`$, the same in (34) and (36), and, for $`M_3`$, the same in (35) and (37) . However, (36) and (37) differ significantly from (34) and (35). These latter expressions are functions of the momentum transfer, $`q^2`$. The integrals $`K(p_i,p_j)`$, on the other hand, are functions only of the scalar invariants $`p_i^2,`$ $`p_j^2`$ and $`p_ip_j`$. In (3.21) and (3.22), $`M_2`$ and $`M_3`$ therefore depend only on the initial and final electron energies, and not on the momentum transfer ( $`p_2p_1=`$ $`p_4p_3=ϵ_1M;`$ $`p_2p_3=`$ $`p_4p_1=ϵ_3M`$).
From (A13)-(A15), (25), (26), (33), (34) and (35), we can now write the square of the matrix element for elastic scattering, including the radiative correction to order $`\alpha `$. Assuming $`q^2m^2,`$ and including only electron-positron pairs in the vacuum polarization matrix element, we have
$`||^2`$ $`=`$ $`|M_1|^2\left\{\begin{array}{c}1+\frac{\alpha }{\pi }\left[\frac{13}{6}\mathrm{ln}\left(\frac{q^2}{m^2}\right)\frac{28}{9}K(p_1,p_3)+K(p_1,p_1)\right]\\ \frac{2\alpha Z}{\pi }\mathrm{ln}\eta \mathrm{ln}\left(\frac{q^2}{\lambda ^2}\right)\\ +\frac{\alpha Z^2}{\pi }\left[K(p_2,p_4)+K(p_2,p_2)\right]\end{array}\right\}`$ (42)
$`+2\text{Re}\left\{M_1^{}\overline{M}_6^{(1)}\right\}`$
### E Contribution of proton form factor
Finally, we consider the contribution of the term $`2`$Re$`\left\{M_1^{}\overline{M}_6^{(1)}\right\}`$, coming from the inclusion of form factors for the proton and integration over the entire range of four-momenta of the virtual photon in the proton vertex correction. Equations (16) and (17) define the functions $`G_1(q^2),G_2(q^2),X_1(q^2),`$ and $`X_2(q^2).`$ From (18) we may write $`M_6^{(1)}=M_6M_6^{(0)}`$, i.e., the term $`M_6^{(1)}`$ is obtained from the full proton vertex correction by subtracting the infrared divergent matrix element $`M_6^{(0)}`$ which is independent of the proton form factor. We therefore define $`G_1^{^{}}(q^2)`$ and $`X_2^{^{}}(q^2)`$ to be the expressions $`G_1(q^2)`$ and $`X_2(q^2)`$ from which we have omitted the terms with factor $`\varphi _1(\lambda ^2)`$ We then write, from (16) and (17),
$$\overline{(g+s)g(g+s)}^{}=F(q^2)\left[\left(G_1^{^{}}(q^2)G_1^{^{}}(0)\right)\gamma _\mu +G_2(q^2)\frac{i\sigma _{\mu \nu }q^\nu }{2M}\right]$$
(43)
and
$$\overline{(g+s)s(g+s)}^{}=\kappa F(q^2)\left[X_1(q^2)\gamma _\mu +\left(X_2^{^{}}(q^2)G_1^{^{}}(0)\right)\frac{i\sigma _{\mu \nu }q^\nu }{2M}\right]$$
(44)
We now define the functions $`F_{1g}(q^2),F_{2g}(q^2),F_{1s}(q^2),`$ and $`F_{2s}(q^2)`$ by
$$F(q^2)\left[\left(G_1^{^{}}(q^2)G_1^{^{}}(0)\right)\gamma _\mu +G_2(q^2)\frac{i\sigma _{\mu \nu }q^\nu }{2M}\right]F_{1g}(q^2)\gamma _\mu +\kappa F_{2g}(q^2)\frac{i\sigma _{\mu \nu }q^\nu }{2M}$$
(45)
$$\kappa F(q^2)[X_1(q^2)\gamma _\mu +(X_2^{^{}}(q^2\}G_1^{^{}}(0))\frac{i\sigma _{\mu \nu }q^\nu }{2M}]F_{1s}(q^2)\gamma _\mu +\kappa F_{2s}(q^2)\frac{i\sigma _{\mu \nu }q^\nu }{2M}$$
(46)
Then with the further definitions
$$\stackrel{~}{F}_1(q^2)F_{1g}(q^2)+F_{1s}(q^2)$$
(47)
$$\stackrel{~}{F}_2(q^2)F_{2g}(q^2)+F_{2s}(q^2)$$
(48)
$$\stackrel{~}{\mathrm{\Gamma }}_\mu \stackrel{~}{F}_1(q^2)\gamma _\mu +\kappa \stackrel{~}{F}_2(q^2)\frac{i\sigma _{\mu \nu }q^\nu }{2M}$$
(49)
we have (apart from factors)
$$\overline{M}_6^{(1)}=\frac{\alpha Z^2}{2\pi }p_3|\gamma ^\mu |p_1p_4|\stackrel{~}{\mathrm{\Gamma }}_\mu |p_2$$
(50)
and
$$2\text{Re}\left\{M_1^{}\overline{M}_6^{(1)}\right\}=\frac{\alpha Z^2}{\pi }\left(p_3|\gamma ^\nu |p_1p_4|\mathrm{\Gamma }_\nu |p_2\right)^{}\left(p_3|\gamma ^\mu |p_1p_4|\stackrel{~}{\mathrm{\Gamma }}_\mu |p_2\right)$$
(51)
This has the same form as
$$M_1^{}M_1=\left(p_3|\gamma ^\nu |p_1p_4|\mathrm{\Gamma }_\nu |p_2\right)^{}\left(p_3|\gamma ^\mu |p_1p_4|\mathrm{\Gamma }_\mu |p_2\right)$$
(52)
with the exception of the replacement $`\mathrm{\Gamma }_\mu \stackrel{~}{\mathrm{\Gamma }}_\mu `$ in the right-hand term. Thus, in place of the Rosenbluth cross section, (13), obtained from $`_{spins}M_1^{}M_1`$, we have
$$\underset{spins}{}2\text{Re}\left\{M_1^{}\overline{M}_6^{(1)}\right\}=\frac{\alpha ^2\mathrm{cos}^2\frac{\theta }{2}}{4ϵ_1^2\eta \mathrm{sin}^4\left(\theta /2\right)}\left(\frac{\alpha Z^2}{\pi }\right)\left\{\right\}$$
(53)
where
$$\left\{\right\}=\left(F_1\stackrel{~}{F}_1\frac{\kappa ^2q^2}{4M^2}F_2\stackrel{~}{F}_2\right)\frac{q^2}{2M^2}\left(F_1+\kappa F_2\right)\left(\stackrel{~}{F}_1+\kappa \stackrel{~}{F}_2\right)\mathrm{tan}^2\frac{\theta }{2}$$
(54)
The purely elastic cross section, including radiative corrections to order $`\alpha ,`$ can thus be written as
$$\left(\frac{d\sigma _0}{d\mathrm{\Omega }}\right)\left(1+\delta _{el}^{(0)}+\delta _{el}^{(1)}\right)$$
(55)
where
$`\delta _{el}^{(0)}`$ $`=`$ $`{\displaystyle \frac{\alpha }{\pi }}\left[\left[\mathrm{ln}\left({\displaystyle \frac{q^2}{m^2}}\right)1\right]\mathrm{ln}\left({\displaystyle \frac{m^2}{\lambda ^2}}\right)+{\displaystyle \frac{13}{6}}\mathrm{ln}\left({\displaystyle \frac{q^2}{m^2}}\right){\displaystyle \frac{28}{9}}{\displaystyle \frac{1}{2}}\mathrm{ln}^2\left({\displaystyle \frac{q^2}{m^2}}\right)+{\displaystyle \frac{\pi ^2}{6}}\right]`$ (56)
$``$ $`{\displaystyle \frac{2\alpha Z}{\pi }}\mathrm{ln}\eta \mathrm{ln}\left({\displaystyle \frac{q^2}{\lambda ^2}}\right)`$ (57)
$`+`$ $`{\displaystyle \frac{\alpha Z^2}{\pi }}\left[\left({\displaystyle \frac{ϵ_4}{|\text{p}_\mathrm{𝟒}|}}\mathrm{ln}x1\right)\mathrm{ln}\left({\displaystyle \frac{M^2}{\lambda ^2}}\right)+{\displaystyle \frac{ϵ_4}{|\text{p}_\mathrm{𝟒}|}}\left[\mathrm{ln}x\mathrm{ln}\left({\displaystyle \frac{\rho ^2}{M^2}}\right)+{\displaystyle \frac{1}{2}}\mathrm{ln}^2x+2L({\displaystyle \frac{1}{x}})+{\displaystyle \frac{\pi ^2}{6}}\right]\right]`$ (58)
and
$$\delta _{el}^{(1)}=\frac{\alpha Z^2}{\pi }\left\{\frac{\left(F_1\stackrel{~}{F}_1\frac{\kappa ^2q^2}{4M^2}F_2\stackrel{~}{F}_2\right)\frac{q^2}{2M^2}\left(F_1+\kappa F_2\right)\left(\stackrel{~}{F}_1+\kappa \stackrel{~}{F}_2\right)\mathrm{tan}^2\frac{\theta }{2}}{\left[\left(F_1^2\frac{\kappa ^2q^2}{4M^2}F_2^2\right)\frac{q^2}{2M^2}\left(F_1+\kappa F_2\right)^2\mathrm{tan}^2\frac{\theta }{2}\right]}\right\}$$
(59)
## IV Inelastic Cross Section
In this section we calculate the inelastic cross section, i.e., the contribution of soft photon emission by the initial and final electron and proton to the radiative correction. The relevant diagrams, with corresponding matrix elements $`M_{b1}`$ and $`M_{b2}`$, are shown in Fig. 2. These matrix elements are given by
$`M_{b1}`$ $`=`$ $`iZe^3(2\pi )^4\delta ^4(p_3+p_4+kp_1p_2){\displaystyle \frac{mM}{\sqrt{2\omega ϵ_1ϵ_3ϵ_2ϵ_4}}}`$ (62)
$`\times \overline{u}(p_3)[ϵ/{\displaystyle \frac{1}{p/_3+k/m+iϵ}}\gamma _\mu +\gamma _\mu {\displaystyle \frac{1}{p/_1k/m+iϵ}}ϵ/]u(p_1)`$
$`\times \overline{u}(p_4)\mathrm{\Gamma }^\mu u(p_2){\displaystyle \frac{1}{(p_1p_3k)^2+iϵ}}`$
$`M_{b2}`$ $`=`$ $`iZ^2e^3(2\pi )^4\delta ^4(p_3+p_4+kp_1p_2){\displaystyle \frac{mM}{\sqrt{2\omega ϵ_1ϵ_3ϵ_2ϵ_4}}}\overline{u}(p_3)\gamma _\mu u(p_1)`$ (64)
$`\times \overline{u}(p_4)[ϵ/{\displaystyle \frac{1}{p/_4+k/m+iϵ}}\mathrm{\Gamma }^\mu +\mathrm{\Gamma }^\mu {\displaystyle \frac{1}{p/_2k/m+iϵ}}ϵ/]u(p_2){\displaystyle \frac{1}{(p_1p_3)^2+iϵ}}`$
Making the soft photon approximation, we rationalize the denominators and drop terms of relative order $`k`$ in the numerator and denominator (but not in the delta function), giving
$`M_{b1}+M_{b2}`$ $`=`$ $`iZe^3(2\pi )^4\delta ^4(p_3+p_4+kp_1p_2){\displaystyle \frac{mM}{\sqrt{2\omega ϵ_1ϵ_3ϵ_2ϵ_4}}}{\displaystyle \frac{1}{q^2}}`$ (66)
$`\times \overline{u}(p_3)\gamma _\mu u(p_1)\overline{u}(p_4)\mathrm{\Gamma }^\mu u(p_2)\left({\displaystyle \frac{p_3ϵ}{p_3k}}{\displaystyle \frac{p_1ϵ}{p_1k}}Z{\displaystyle \frac{p_4ϵ}{p_4k}}+Z{\displaystyle \frac{p_2ϵ}{p_2k}}\right)`$
The cross section for soft bremsstrahlung then follows by squaring the matrix element $`M_{b1}+M_{b2}`$ , dividing by the incident flux and the transition rate and multiplying by the number of final states. Summing over photon polarizations, we then have
$`d\sigma _b`$ $`=`$ $`{\displaystyle \frac{Z^2e^6}{(2\pi )^9}}{\displaystyle \frac{m^2M^2}{\sqrt{(p_1p_2)^2m^2M^2}}}{\displaystyle \underset{spins}{}}{\displaystyle \frac{d^3p_3}{ϵ_3}\frac{d^3p_4}{ϵ_4}\frac{d^3k}{2\omega }(2\pi )^4\delta ^4(p_3+p_4+kp_1p_2)}`$ (68)
$`\times \left|\overline{u}(p_3)\gamma _\mu u(p_1)\overline{u}(p_4)\mathrm{\Gamma }^\mu u(p_2)\right|^2{\displaystyle \frac{1}{(q^2)^2}}\left({\displaystyle \frac{p_3}{p_3k}}{\displaystyle \frac{p_1}{p_1k}}Z{\displaystyle \frac{p_4}{p_4k}}+Z{\displaystyle \frac{p_2}{p_2k}}\right)^2`$
The range of integration in the above expression is determined by the experimental conditions. We assume, as in , that the final proton and emitted photon are undetected; the range of integration in energy and angle of the final electron is determined by the entrance slit and spectrometer. We integrate first over $`d^3p_4`$ and are then left with a single delta function relating the variables of $`𝐤`$ and $`𝐩_3`$: Writing
$$\frac{d^3p_4}{2ϵ_4}=_0^{\mathrm{}}𝑑ϵ_4\delta (p_4^2M^2)d^3p_4=d^4p_4\delta (p_4^2M^2)\theta (p_4^0)$$
(69)
with
$$tp_1+p_2p_3=p_4+k$$
(70)
we then have
$`d\sigma _b`$ $`=`$ $`{\displaystyle \frac{Z^2e^6}{(2\pi )^5}}{\displaystyle \frac{m^2M^2}{\sqrt{(p_1p_2)^2m^2M^2}}}{\displaystyle \underset{spins}{}}{\displaystyle \frac{d^3p_3}{ϵ_3}\frac{d^3k}{\omega }\delta \left((tk)^2M^2\right)\theta (ϵ_4)}`$ (72)
$`\times \left|\overline{u}(p_3)\gamma _\mu u(p_1)\overline{u}(p_4)\mathrm{\Gamma }^\mu u(p_2)\right|^2{\displaystyle \frac{1}{(q^2)^2}}\left({\displaystyle \frac{p_3}{p_3k}}{\displaystyle \frac{p_1}{p_1k}}Z{\displaystyle \frac{p_4}{p_4k}}+Z{\displaystyle \frac{p_2}{p_2k}}\right)^2`$
in which $`p_4=p_1+p_2p_3k`$. We may then transform to the special frame S<sup>0</sup> (defined by t $`=0`$), in which the delta function in (72) is independent of the angle at which the photon is emitted. There
$$(tk)^2M^2=t_0^22t_0\omega +\lambda ^2M^2=0;t_0=ϵ_1+ϵ_2ϵ_3$$
(73)
The photon energy is then given solely by the final electron energy. The procedure used in is to integrate next over the photon energy and angle in S<sup>0</sup> and then transform back to the lab frame to integrate over the energy and angle of the final electron. Instead, we remain in the special frame and integrate first over $`ϵ_3`$, the delta function giving $`ϵ_3`$ in terms of $`\omega `$. However, the range of photon energies is assumed to be sufficiently small compared to all other energies that we can set $`ϵ_3`$ equal to its value for elastic scattering throughout the integrand. In addition, we take the angular range of the final electron to be sufficiently small that we can take some average value for these angles and neglect any variation of the integrand over this angular range. Similarly, we neglect $`k`$ in the above expression for $`p_4`$. We may then take $`\left|\overline{u}(p_3)\gamma _\mu u(p_1)\overline{u}(p_4)\mathrm{\Gamma }^\mu u(p_2)\right|^2\mathrm{}(q^2)^2`$ outside of the integration, giving
$`d\sigma _b`$ $`=`$ $`{\displaystyle \frac{Z^2e^6}{(2\pi )^5}}{\displaystyle \frac{m^2M^2}{\sqrt{(p_1p_2)^2m^2M^2}}}{\displaystyle \underset{spins}{}}\left|\overline{u}(p_3)\gamma _\mu u(p_1)\overline{u}(p_4)\mathrm{\Gamma }^\mu u(p_2)\right|^2{\displaystyle \frac{1}{(q^2)^2}}`$ (75)
$`\times {\displaystyle \frac{\left|\text{p}_3\right|}{2M}}{\displaystyle ^{^{}}}{\displaystyle \frac{d^3k}{\omega }}\left({\displaystyle \frac{p_3}{p_3k}}{\displaystyle \frac{p_1}{p_1k}}Z{\displaystyle \frac{p_4}{p_4k}}+Z{\displaystyle \frac{p_2}{p_2k}}\right)^2`$
The term $`2M`$ in the denominator in (75) comes from the delta function in (72), which contributes the factor $`|d\{\delta \left((tk)^2M^2\right)\}/dϵ_3|=2|t_0\omega |`$ from (73). Again neglecting terms of order $`k,`$ we note from (70) that in S<sup>0</sup>, $`t_0\omega =ϵ_4=M`$. Comparing (68), (10) and (12) and noting that in arriving at (75) we have neglected terms of order $`k`$ in $`p_3`$ and $`p_4`$, we have
$$d\sigma _b=\frac{\alpha }{4\pi ^2}d\sigma _0^{^{}}\frac{d^3k}{\omega }\left(\frac{p_3}{p_3k}\frac{p_1}{p_1k}Z\frac{p_4}{p_4k}+Z\frac{p_2}{p_2k}\right)^2$$
(76)
where $`\omega =\sqrt{𝐤^2+\lambda ^2}`$. There then remains the integration over photon energy (restricted to $`\left|𝐤\right|\mathrm{\Delta }ϵ`$) and angles. The relevant integrals have been evaluated by ’t Hooft and Veltman , Sec. 7. We give here only their final result, rewritten using our metric; the essential steps in the derivation are given in their work. They define
$$L_{ij}=^{^{}}\frac{d^3k}{\omega }\frac{1}{(p_ik)(p_jk)}$$
(77)
in terms of which
$$d\sigma _b=\frac{\alpha }{4\pi ^2}d\sigma _0\left\{\begin{array}{c}m^2L_{11}+m^2L_{33}2p_1p_3L_{13}\\ +Z\left(2p_1p_2L_{12}+2p_3p_2L_{32}+2p_1p_4L_{14}2p_3p_4L_{34}\right)\\ +Z^2\left(M^2L_{22}+M^2L_{44}2p_2p_4L_{24}\right)\end{array}\right\}$$
(78)
As shown in , Sec.7, for the case in which the momenta $`p_i`$ and $`p_j`$ are all on the mass shell, the integrals $`L_{ij}`$ can, provided $`p_i`$ is not a multiple $`p_j`$, be written in the form
$$L_{ij}=\frac{2\pi }{\sqrt{(p_ip_j)^2m_i^2m_j^2}}\left\{S_{ij}^{(1)}+S_{ij}^{(2)}\right\}$$
(79)
where
$$S_{ij}^{(1)}=2\mathrm{ln}\left(\frac{p_ip_j+\sqrt{(p_ip_j)^2m_i^2m_j^2}}{m_im_j}\right)\mathrm{ln}\left(\frac{2\mathrm{\Delta }ϵ}{\lambda }\right)$$
(80)
and
$`S_{ij}^{(2)}`$ $`=`$ $`\mathrm{ln}^2\left({\displaystyle \frac{\beta _i}{m_i\sqrt{t^2}}}\right)\mathrm{ln}^2\left({\displaystyle \frac{\beta _j}{m_j\sqrt{t^2}}}\right)`$ (83)
$`+L\left(1{\displaystyle \frac{\beta _ilt}{t^2\gamma _{ij}}}\right)+L\left(1{\displaystyle \frac{m_i^2lt}{\beta _i\gamma _{ij}}}\right)`$
$`L\left(1{\displaystyle \frac{\beta _jlt}{\alpha t^2\gamma _{ij}}}\right)L\left(1{\displaystyle \frac{m_j^2lt}{\alpha \beta _j\gamma _{ij}}}\right)`$
in which
$$\alpha =\frac{p_ip_j+\sqrt{(p_ip_j)^2m_i^2m_j^2}}{m_i^2}\text{}l=\alpha p_ip_j$$
(84)
$$\beta _{i,j}p_{i,j}t+\sqrt{(p_{i,j}t)^2m_{i,j}^2t^2}\text{}\gamma _{ij}\sqrt{(p_ip_j)^2m_i^2m_j^2}$$
(85)
The evaluation of (77) for $`p_i=p_j`$ is straightforward. The result written in terms of relativistic invariants is
$$L_{ii}=\frac{4\pi }{m_i^2}\left[\mathrm{ln}\left(\frac{2\mathrm{\Delta }ϵ}{\lambda }\right)\frac{p_it}{\sqrt{(p_it)^2m_i^2t^2}}\mathrm{ln}\left(\frac{\beta _i}{m_i\sqrt{t^2}}\right)\right]$$
(86)
In , Sec. 7, the expression for $`S_{ij}^{(2)}`$ is evaluated in the frame S<sup>0</sup>, defined by t$`=0`$. Since we want finally to express the cross section in terms of lab frame energies and momenta, we have, in (83), written $`S_{ij}^{(2)}`$ in terms of relativistic invariants.
The terms of leading order in $`\mathrm{ln}\lambda `$ are apparent in (80) and (86). Substituting these in (78) gives the infrared divergent terms in $`d\sigma _b`$. They are cancelled exactly by the $`\mathrm{ln}\lambda `$ terms in the elastic cross section.
We next express $`d\sigma _b`$ in terms of lab frame energies. To that end, we assume that $`\mathrm{\Delta }ϵ`$ is less than any of the other energies and therefore now neglect $`k`$ in (70), taking $`p_4`$ to be given by its value for elastic (non-radiative) scattering:
$$t=p_1+p_2p_3=p_4$$
(87)
(Note that for $`\mathrm{\Delta }ϵ0,`$ $`S_{ij}^{(2)}`$ remains finite; the only singularity is confined to the term $`\mathrm{ln}(2\mathrm{\Delta }ϵ/\lambda )`$, evident in $`S_{ij}^{(1)}`$.) The relativistic invariants in $`L_{ij}`$ can then be written simply in terms of lab frame energies:
$`p_1p_4`$ $`=`$ $`p_2p_3=Mϵ_3`$ (88)
$`p_3p_4`$ $`=`$ $`p_2p_1=Mϵ_1`$ (89)
$`p_1p_3`$ $`=`$ $`{\displaystyle \frac{1}{2}}q^2+m^2`$ (90)
$`p_2p_4`$ $`=`$ $`{\displaystyle \frac{1}{2}}q^2+M^2`$ (91)
Further, we express $`\mathrm{\Delta }ϵ`$, the maximum momentum of the photon in the frame $`S^0`$, in terms of the final electron detector acceptance in the lab frame, $`\mathrm{\Delta }E`$:
$$\mathrm{\Delta }ϵ=\eta \mathrm{\Delta }E$$
(92)
Details of the derivation are given in Appendix E. Substituting (87) and (89) in (80), (83) and (86) then gives $`L_{ij}`$ in terms of lab frame energies and momenta. We note that although $`L_{ij}`$ as defined in (77) is clearly symmetric in $`i`$ and $`j`$, the expression for $`S_{ij}^{(2)}`$ in (83) is not manifestly symmetric; the form of the expression for $`L_{ij}`$ is rather different from that for $`L_{ji}`$. In writing the explicit expressions for the terms in $`L_{ij}`$, we choose $`i`$ and $`j`$ such that $`L_{ij}`$ simplifies readily for lab frame electron energies and momentum transfers which are very large compared to the electron rest mass. When $`m_im_j`$, this is achieved by choosing $`i`$ and $`j`$ such that $`m_i=m`$ and $`m_j=M`$ (note below that we have written $`S_{32}^{(2)}`$ ).
At this point we make the high energy approximation as defined in section II, in which case the above expressions for $`S_{ij}^{(1)}`$ and $`S_{ij}^{(2)}`$ simplify considerably. For $`S_{ij}^{(1)}`$ this is straightforward. We give the results for $`S_{ij}^{(2)}`$ in Appendix F. Substituting these in $`L_{ij}`$, the high energy approximation for the inelastic cross section given in (78) then becomes
$$d\sigma _b=\frac{\alpha }{\pi }d\sigma _0\left\{\begin{array}{c}\left[\mathrm{ln}\left(\frac{q^2}{m^2}\right)1\right]\mathrm{ln}\left(\frac{(2\eta \mathrm{\Delta }E)^2}{\lambda ^2}\right)\\ \left[\mathrm{ln}\left(\frac{q^2}{m^2}\right)1\right]\mathrm{ln}\left(\frac{4ϵ_1ϵ_3}{m^2}\right)\\ +\frac{1}{2}\mathrm{ln}^2\left(\frac{q^2}{m^2}\right)\frac{1}{2}\mathrm{ln}^2\eta +L(\mathrm{cos}^2\frac{1}{2}\theta )\frac{1}{3}\pi ^2\\ +2Z\left[\begin{array}{c}\mathrm{ln}\eta \mathrm{ln}\left(\frac{(2\eta \mathrm{\Delta }E)^2}{\lambda ^2}\right)\mathrm{ln}\eta \mathrm{ln}x\\ +L(1\frac{\eta }{x})L(1\frac{1}{\eta x})\end{array}\right]\\ +Z^2\left[\begin{array}{c}\left(\frac{ϵ_4}{|\text{p}_4|}\mathrm{ln}x1\right)\mathrm{ln}\left(\frac{(2\eta \mathrm{\Delta }E)^2}{\lambda ^2}\right)\\ \frac{ϵ_4}{|\text{p}_4|}\left[\mathrm{ln}^2x\mathrm{ln}x+L(1\frac{1}{x^2})\right]1\end{array}\right]\end{array}\right\}$$
(93)
## V Radiative Corrections to elastic electron-proton scattering
The results given in (57), (59), and (93) may be added to give the radiative correction, $`\delta `$. The analytic expression is given below in (94) and (100). Numerical evaluation of the radiative correction for various values of the pertinent parameters (initial beam energy, momentum transfer, final electron detector resolution, and target nucleus) are given in Tables I, II, and III. We note that the infrared ($`\mathrm{ln}\lambda `$) terms, which appear in both the purely elastic ( 57) and inelastic (93) contributions to the radiative correction, cancel exactly when added to give the cross section for elastic electron-proton scattering with radiative corrections to first order in $`\alpha `$:
$$d\sigma =d\sigma _0(1+\delta )$$
(94)
where
$`\delta `$ $`=`$ $`{\displaystyle \frac{\alpha }{\pi }}\left[{\displaystyle \frac{13}{6}}\mathrm{ln}\left({\displaystyle \frac{q^2}{m^2}}\right){\displaystyle \frac{28}{9}}\left[\mathrm{ln}\left({\displaystyle \frac{q^2}{m^2}}\right)1\right]\mathrm{ln}\left({\displaystyle \frac{4ϵ_1ϵ_3}{(2\eta \mathrm{\Delta }E)^2}}\right){\displaystyle \frac{1}{2}}\mathrm{ln}^2\eta +L(\mathrm{cos}^2{\displaystyle \frac{1}{2}}\theta ){\displaystyle \frac{\pi ^2}{6}}\right]`$ (100)
$`+{\displaystyle \frac{2\alpha Z}{\pi }}\left[\mathrm{ln}\eta \mathrm{ln}\left({\displaystyle \frac{q^2x}{(2\eta \mathrm{\Delta }E)^2}}\right)+L\left(1{\displaystyle \frac{\eta }{x}}\right)L\left(1{\displaystyle \frac{1}{\eta x}}\right)\right]`$
$`+{\displaystyle \frac{\alpha Z^2}{\pi }}\left[\begin{array}{c}\frac{ϵ_4}{|\text{p}_4|}\left(\frac{1}{2}\mathrm{ln}^2x\mathrm{ln}x\mathrm{ln}\left(\frac{\rho ^2}{M^2}\right)+\mathrm{ln}x\right)\left(\frac{ϵ_4}{|\text{p}_4|}\mathrm{ln}x1\right)\mathrm{ln}\left(\frac{M^2}{(2\eta \mathrm{\Delta }E)^2}\right)+1\\ +\frac{ϵ_4}{|\text{p}_4|}\left(L\left(1\frac{1}{x^2}\right)+2L\left(\frac{1}{x}\right)+\frac{\pi ^2}{6}\right)\end{array}\right]`$
$`+\delta _{el}^{(1)}`$
Here, $`\delta _{el}^{(1)}`$ is the contribution coming from the inclusion of form factors for the proton and integration over the entire range of four-momenta of the virtual photon in the proton vertex correction (see (18), (25), (26), (42)); it is thus not included in the analysis given in and , denoted here as the soft photon approximation. Moreover, $`\delta _{el}^{(1)}`$ has no infrared divergent terms; these are all included in the soft photon approximation.
In Tables I, II, and III we compare the values of the radiative correction, $`\delta `$, calculated in this paper (denoted by MTj) with those given by Mo and Tsai in for various kinematics. The initial beam energies and momentum transfers have been chosen to correspond to experiments proposed or already performed at Jefferson Lab and SLAC . The final electron detector acceptance, $`\mathrm{\Delta }E`$, has been taken throughout to be one percent of the final electron energy, $`ϵ_3`$. In the form factors (see(6)), the parameter $`\mathrm{\Lambda }`$ has been chosen to be 700 MeV/$`c`$ throughout. The contribution of the terms in (100) are grouped according to the power of $`Z`$ which appears there as a factor. The numerical value of each of these groups of terms is given in the rows denoted by $`Z^0,Z^1,Z^2`$. Values given in the column MTj in the row $`Z^2`$ do not include the contribution of the proton form factor (which are contained in $`\delta _{el}^{(1)})`$; they are given for comparison with the values in . In the range of energies and momentum transfers considered here, the correction $`\delta _{el}^{(1)}`$, due to the finite size of the nucleon (and integration over the entire range of four-momenta of the virtual photon in the proton vertex correction), is found in general to be much smaller than the other contributions with factor $`Z^2`$, labeled explicitly in (100) and in Tables I, II, and III. The values given in these tables include only the contribution of electron-positron pairs in the vacuum polarization; the contribution of muon and tau pairs is given by (A13) and (A15).
The curves in Figs. 3 and 4 illustrate the two aspects of the present work: (1) the contribution of nucleonic size effects to the radiative correction, and (2) the improvement of the mathematical treatment of the integrations given in the work of Mo and Tsai , . The nucleonic size effects are all contained in the term $`\delta _{el}^{(1)}`$, (Eq.(59)); its contribution relative to the overall radiative correction factor, $`(1+\delta _{\text{MTj}})`$, is given by the dashed curves marked VTX, where VTX = 100$`\times \delta _{el}^{(1)}/(1+\delta _{\text{MTj}})`$. The dotted curve shows D0 = 100$`\times (`$ $`\delta _{\text{MTj}}^{(0)}\delta _{\text{Tsai}})/(1+\delta _{\text{MTj}})`$, which is that part of the difference between the radiative correction given by Mo and Tsai and the one given in this paper due solely to the improvement of the mathematical treatment of the integrations. Here, $`\delta _{\text{MTj}}^{(0)}`$ is the radiative correction given in (100), excluding the term $`\delta _{el}^{(1)}`$. It will be noticed that VTX is always positive, and that for most of the range of allowed momentum transfers, D0 is negative. Thus their sum, which is the difference between the radiative correction $`\delta _{\text{MTj}}`$ given in this paper in (100), and $`\delta _{\text{Tsai}}`$, given in , is rather small except for the region corresponding to large scattering angles. This sum is given by the solid curves marked D, where D = D0 + VTX = 100$`\times (`$ $`\delta _{\text{MTj}}\delta _{\text{Tsai}})/(1+\delta _{\text{MTj}})`$.
## VI Conclusion
We have calculated the radiative correction to elastic electron-proton scattering to lowest order in $`\alpha `$ using a hadronic model which includes the finite size of the nucleon. The contribution from the emission of real soft photons by the electron and the proton is calculated exactly. The contributions of the box and crossed-box (two-photon exchange) diagrams are calculated in a soft photon approximation which is less drastic than that employed in . A number of observations may be made from the values given in Tables I, II, and III. First, the contributions of the electron vertex correction, vacuum polarization, and real soft photon emission by the electron (the terms in (100) with factor $`\alpha /\pi `$) dominate the radiative correction $`\delta `$. Since our expression for these terms differs from that given by Mo and Tsai solely in that they have omitted the term $`\left(\alpha /\pi \right)\left[L(\mathrm{cos}^2\frac{1}{2}\theta )\frac{\pi ^2}{6}\right]`$ in (100), we find values for $`\delta `$ which differ from theirs by at most 2% for the initial energies and momentum transfers considered here (note that $`\frac{\pi ^2}{6}L(\mathrm{cos}^2\frac{1}{2}\theta )\frac{\pi ^2}{6}0`$). Further, we note that, except for the proton and at the higher energies considered here, the contribution of $`\delta _{el}^{(1)}`$ is negligible. However, for the two highest energies, $`\delta _{el}^{(1)}`$ is between 2% and 3% of the factor ($`1+\delta `$) by which the uncorrected cross section must be multiplied, and hence should be considered in precision measurements for electron-proton scattering at energies above 8 GeV. As an empirical guide, we find that $`\delta _{el}^{(1)}=0.02(1+\delta )`$ for initial energies and scattering angles satisfying $`ϵ_1\mathrm{sin}\theta 8`$ for beam energies between 8 and 16 GeV. Finally, we note that a considerable simplification of the expression in (100) occurs if, in addition to the last two terms multiplying $`\alpha /\pi `$, we neglect the last two terms multiplying $`2\alpha Z/\pi `$ as well as the last three terms multiplying $`\alpha Z^2/\pi `$, each of these sets of terms being always less than $`\pi ^2/6`$ in magnitude. From this study we see that at the energies and momentum transfers considered here, the nucleonic finite size effects are rather small but are expected to become more important at higher energies. The corrections due to the improvement of the high energy behavior of the radiative corrections as described in this paper are not negligible and need to be taken into account at the energies and momentum transfers we have considered.
###### Acknowledgements.
It is a pleasure to acknowledge the valued assistance of W.C. Parke in the realization of this paper.
## A Elastic scattering amplitudes to order $`\alpha ^2`$
Using the notation given in section II, the matrix elements corresponding to the various one-loop diagrams shown in Fig. 1 are
$$M_5=Ze^2\overline{u}(p_3)\mathrm{\Lambda }_\mu (p_3,p_1)u(p_1)\frac{(i)}{q^2+iϵ}\overline{u}(p_4)\mathrm{\Gamma }^\mu (q^2)u(p_2)$$
(A1)
$$M_6=Z^3e^2\overline{u}(p_3)\gamma _\mu u(p_1)\frac{(i)}{q^2+iϵ}\overline{u}(p_4)\mathrm{\Lambda }^\mu (p_4,p_2)u(p_2)$$
(A2)
where
$$\mathrm{\Lambda }_\mu (p_3,p_1)=ie^2\frac{d^4k}{(2\pi )^4}\frac{1}{k^2\lambda ^2+iϵ}\gamma ^\nu \frac{1}{(p/_3k/m+iϵ)}\gamma _\mu \frac{1}{(p/_1k/m+iϵ)}\gamma _\nu $$
(A3)
$`\mathrm{\Lambda }^\mu (p_4,p_2)`$ $`=`$ $`ie^2{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{1}{k^2\lambda ^2+iϵ}\mathrm{\Gamma }^\nu (k^2)\frac{1}{(p/_4k/M+iϵ)}\mathrm{\Gamma }^\mu (q^2)}`$ (A5)
$`\times {\displaystyle \frac{1}{(p/_2k/M+iϵ)}}\mathrm{\Gamma }_\nu (k^2)`$
$`M_2`$ $`=`$ $`(Ze^2)^2{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{1}{k^2\lambda ^2+iϵ}\frac{1}{(k+q)^2\lambda ^2+iϵ}}`$ (A8)
$`\times \left[\overline{u}(p_3)\gamma _\nu {\displaystyle \frac{1}{p/_1k/m+iϵ}}\gamma _\mu u(p_1)\right]`$
$`\times \left[\overline{u}(p_4)\mathrm{\Gamma }^\nu ((k+q)^2){\displaystyle \frac{1}{p/_2+k/M+iϵ}}\mathrm{\Gamma }^\mu (k^2)u(p_2)\right]`$
$`M_3`$ $`=`$ $`(Ze^2)^2{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{1}{k^2\lambda ^2+iϵ}\frac{1}{(k+q)^2\lambda ^2+iϵ}}`$ (A11)
$`\times \left[\overline{u}(p_3)\gamma _\nu {\displaystyle \frac{1}{p/_1k/m+iϵ}}\gamma _\mu u(p_1)\right]`$
$`\times \left[\overline{u}(p_4)\mathrm{\Gamma }^\mu (k^2){\displaystyle \frac{1}{p/_4k/M+iϵ}}\mathrm{\Gamma }^\nu ((k+q)^2)u(p_2)\right]`$
The matrix element, $`M_4`$, for vacuum polarization is, after charge renormalization, related simply to the matrix element $`M_1`$ given above:
$$M_4=\mathrm{\Pi }(q^2)M_1$$
(A12)
For a fermion loop in the photon propagator, $`\mathrm{\Pi }(q^2)`$ is given in in terms of an integral which can be evaluated in closed form , giving
$$\mathrm{\Pi }(q^2)=\mathrm{\Pi }^f(q^2/m_i^2)=\frac{\alpha }{3\pi }\left\{\left(1\frac{u}{2}\right)\sqrt{1+u}\mathrm{log}\left(\frac{\sqrt{1+u}+1}{\sqrt{1+u}1}\right)+u\frac{5}{3}\right\}$$
(A13)
in which $`m_i`$ is the mass of the fermion and $`u=4m_i^2/(q^2).`$
For $`q^2/m_i^21`$ this gives
$$\mathrm{\Pi }^f(q^2/m_i^2)=\frac{\alpha }{\pi }\left\{\frac{1}{3}\mathrm{ln}\left(\frac{q^2}{m_i^2}\right)\frac{5}{9}\right\}$$
(A14)
If we include the vacuum polarization amplitudes from particle-antiparticle loops of different masses, as has been done in several experimental analyses , then we have
$$M_4=M_1\underset{i}{}\mathrm{\Pi }^f(q^2/m_i^2)$$
(A15)
In principle, once one includes particle-antiparticle pairs of mass greater than the electron mass in the vacuum polarization amplitudes, one should consider bosons as well as fermions. The matrix elements for vacuum polarization for a pair of structureless spin zero bosons in the closed loop, first given by Feynman , may be found in a more accessible form in a paper of Tsai . The result, corresponding to the equation above for fermions, is
$$M_4^{boson}=M_1\underset{i}{}\mathrm{\Pi }^b(q^2/m_i^2)$$
(A16)
where
$$\mathrm{\Pi }^b(q^2/m_i^2)=\frac{\alpha }{3\pi }\left\{\frac{1}{2}\sqrt{1+u}\mathrm{log}\left(\frac{\sqrt{1+u}+1}{\sqrt{1+u}1}\right)u\frac{4}{3}\right\}$$
(A17)
A more complete discussion of vacuum polarization should include a consideration of pion structure as well as the contribution of spin-one bosons, in particular the $`\rho `$ meson. A detailed discussion of the hadronic contribution to vacuum polarization may be found in connection with calculations of the anomalous magnetic moment of the muon and in connection with radiative corrections to high energy electron-positron collider experiments .
## B Proton vertex correction
As noted in Sec. IIIA, the terms $`ggg,gsg,\mathrm{}`$ can be expressed in terms of the integrals given in (3.1). There, the integrals $`I_0,J_0,`$ and $`K_0`$ are scalars and hence are functions of the scalars $`p_2^2,p_4^2,`$and $`p_2p_4`$ (and, of course, $`\lambda ^2`$ and $`\mathrm{\Lambda }^2)`$. Since we have on-shell particles in the initial and final states $`(p_2^2=p_4^2=M^2)`$, these integrals are functions of $`M^2`$ and $`q^2`$. The integrals $`I_\mu `$ and $`J_\mu `$ are vectors and hence in principal can be written in the form
$$I_\mu =ap_{2\mu }+bp_{4\mu }$$
(B1)
with a similar equation for $`J_\mu `$ , where $`a`$ and $`b`$ are functions of $`M^2`$ and $`q^2`$. However, the calculation is simplified greatly if we express $`I_\mu `$ in terms of the four-vectors $`\rho =p_4+p_2`$ (which is symmetric in $`p_4`$ and $`p_2`$) and $`q=p_4p_2`$ (which is antisymmetric in $`p_4`$ and $`p_2`$), i.e. $`I_\mu =A\rho _\mu +Bq_\mu `$. Here $`A`$ and $`B`$ are functions of $`M^2`$ and $`q^2`$ and hence are symmetric in $`p_4`$ and $`p_2`$. Further, since the integrands for the vectors $`I_\mu `$ and $`J_\mu `$ are symmetric in $`p_4`$ and $`p_2`$, it follows that $`B=0.`$ We thus have
$$I_\mu =A\rho _\mu $$
(B2)
and a similar equation for $`J_\mu `$. These same considerations of symmetry allow for the simplification of the tensors $`I_{\mu \nu }`$ and $`J_{\mu \nu }`$, which are also symmetric functions of $`p_4`$ and $`p_2`$. We can therefore write
$$I_{\mu \nu }=a_1\rho _\mu \rho _\nu +a_2q_\mu q_\nu +a_3g_{\mu \nu }$$
(B3)
and a similar equation for $`J_{\mu \nu }`$. That the terms $`\rho _\mu q_\nu `$ and $`q_\mu \rho _\nu `$ are absent follows directly by multiplying $`I_{\mu \nu }`$ successively by $`\rho ^\mu q^\nu `$ and $`q^\mu \rho ^\nu `$, using $`\rho q=0`$ and the fact that $`I_{\mu \nu }\rho ^\mu q^\nu `$ and $`I_{\mu \nu }q^\mu \rho ^\nu `$ are antisymmetric in $`p_2`$ and $`p_4`$. Multiplying $`I_\mu `$ ($`J_\mu `$) by $`\rho ^\mu `$, and $`I_{\mu \nu }`$ ($`J_{\mu \nu }`$) successively by $`\rho ^\mu \rho ^\nu ,q^\mu q^\nu ,`$ and $`g^{\mu \nu }`$, the coefficients in the expressions for $`I_\mu ,J_\mu ,I_{\mu \nu }`$ and $`J_{\mu \nu }`$ may be expressed in terms of their moments, defined by
$$g_1=\frac{1}{\rho ^2}I_\mu \rho ^\mu ,h_1=\frac{1}{\rho ^2}J_\mu \rho ^\mu $$
(B4)
$$g_{11}=\frac{1}{\rho ^4}I_{\mu \nu }\rho ^\mu \rho ^\nu ,h_{11}=\frac{1}{\rho ^4}J_{\mu \nu }\rho ^\mu \rho ^\nu $$
(B5)
$$g_{22}=\frac{1}{\rho ^4}I_{\mu \nu }q^\mu q^\nu ,h_{22}=\frac{1}{\rho ^4}J_{\mu \nu }q^\mu q^\nu $$
(B6)
Conversely, the integrals in (14) may be written in terms of the moments:
$$I_\mu =\rho _\mu g_1,J_\mu =\rho _\mu h_1$$
(B7)
$`I_{\mu \nu }`$ $`=`$ $`\rho _\mu \rho _\nu \left[g_{11}{\displaystyle \frac{1}{2\rho ^2}}(h_0\rho ^2g_{11}q^2g_{22})\right]`$ (B10)
$`+q_\mu q_\nu \left[g_{22}{\displaystyle \frac{1}{2q^2}}(h_0\rho ^2g_{11}q^2g_{22})\right]`$
$`+g_{\mu \nu }\left[{\displaystyle \frac{1}{2}}(h_0\rho ^2g_{11}q^2g_{22})\right]`$
$`J_{\mu \nu }`$ $`=`$ $`\rho _\mu \rho _\nu \left[h_{11}{\displaystyle \frac{1}{2\rho ^2}}(k_0\rho ^2h_{11}q^2h_{22})\right]`$ (B13)
$`+q_\mu q_\nu \left[h_{22}{\displaystyle \frac{1}{2q^2}}(k_0\rho ^2h_{11}q^2h_{22})\right]`$
$`+g_{\mu \nu }\left[{\displaystyle \frac{1}{2}}(k_0\rho ^2h_{11}q^2h_{22})\right]`$
where, for convenience of notation, we have defined
$$g_0=I_0,h_0=J_0,k_0=K_0$$
(B14)
The terms $`ggg,gsg,\mathrm{}`$ can be expressed in terms of these moments. Substituting (5) and (6) in the expression for the proton vertex correction, Eq. (A5), and substituting this in turn in the matrix element $`M_6`$, Eq. (A2), the integrals are all of the form given in Eq. (14), which are in turn expressed in terms of the moments by Eqs. (B7)-(B14). We make use of the fact that in $`M_6`$ the vertex correction is taken between free spinors, and finally express $`p_2`$ and $`p_4`$ in terms of $`\rho `$ and $`q`$. We find
$$ggg=ie^2F(q^2)\left\{\begin{array}{c}\left[\begin{array}{c}2(2M^2q^2)g_04(2M^2q^2)g_1\\ 2(8M^2+q^2)g_{11}2(q^4/\rho ^2)g_{22}+8(M^2/\rho ^2)h_0\end{array}\right]\gamma _\mu \\ +\left[8M^2g_1+24M^2g_{11}+8M^2(q^2/\rho ^2)g_{22}8(M^2/\rho ^2)h_0\right]\frac{i\sigma _{\mu \nu }q^\nu }{2M}\end{array}\right\}$$
(B15)
$$gsg=ie^2\kappa F(q^2)\left\{\begin{array}{c}\left[2q^2g_1\right]\gamma _\mu \\ +\left[2(2M^2q^2)g_04(2M^2q^2)g_1\right]\frac{i\sigma _{\mu \nu }q^\nu }{2M}\end{array}\right\}$$
(B16)
$$ggs+sgg=2ie^2\kappa F(q^2)\left\{\begin{array}{c}\left[12M^2g_{11}(q^4/\rho ^2)g_{22}+2(1+2M^2/\rho ^2)h_03h_1\right]\gamma _\mu \\ +\left[\begin{array}{c}(16M^2q^2)g_{11}+(q^2/\rho ^2)(8M^2q^2)g_{22}\\ 4(1+M^2/\rho ^2)h_0+3h_1\end{array}\right]\frac{i\sigma _{\mu \nu }q^\nu }{2M}\end{array}\right\}$$
(B17)
$$gss+ssg=2ie^2\kappa ^2F(q^2)\left\{\begin{array}{c}(q^2/4M^2)\left[\begin{array}{c}(8M^2+q^2)g_{11}(q^4/\rho ^2)g_{22}\\ +(2+4M^2/\rho ^2)h_0+h_1\end{array}\right]\gamma _\mu \\ +\left[3q^2g_{11}+4M^2(q^2/\rho ^2)g_{22}(q^2/\rho ^2)h_0h_1\right]\frac{i\sigma _{\mu \nu }q^\nu }{2M}\end{array}\right\}$$
(B18)
$$sgs=ie^2\kappa ^2F(q^2)\left\{\begin{array}{c}\left[\begin{array}{c}(8M^2+q^2)g_{11}(q^4/\rho ^2)g_{22}+(2+q^2/\rho ^2)h_0\\ (8M^2+q^2)h_{11}/4M^2(q^4/\rho ^2)h_{22}/4M^2+(1+q^2/\rho ^2)k_0/4M^2\end{array}\right]\gamma _\mu \\ +\left[\begin{array}{c}12M^2g_{11}+4M^2(q^2/\rho ^2)g_{22}2(1+2M^2/\rho ^2)h_0\\ +3h_{11}+4(q^2/\rho ^2)h_{22}k_0/\rho ^2\end{array}\right]\frac{i\sigma _{\mu \nu }q^\nu }{2M}\end{array}\right\}$$
(B19)
$$sss=ie^2\kappa ^3F(q^2)\left\{\begin{array}{c}(q^2/4M^2)\left[\begin{array}{c}12M^2g_{11}4M^2(q^2/\rho ^2)g_{22}+2(1+2M^2/\rho ^2)h_0\\ 4h_1+3h_{11}+(q^2/\rho ^2)h_{22}k_0/\rho ^2\end{array}\right]\gamma _\mu \\ +\left[\begin{array}{c}2(2M^2+q^2)g_{11}+4M^2(q^2/\rho ^2)g_{22}\\ (4M^2/\rho ^2+q^2/2M^2)h_0+(q^2/M^2)h_1\\ 2(2M^2+q^2)h_{11}/4M^2+2(q^2/\rho ^2)(2M^2q^2)h_{22}/4M^2\\ +(q^2/\rho ^2)k_0/4M^2\end{array}\right]\frac{i\sigma _{\mu \nu }q^\nu }{2M}\end{array}\right\}$$
(B20)
The expressions do not depend on the particular form of the form factors; we have assumed only that $`F_1=F_2=F`$. However, for form factors of the form given in Eq. (6), the moments may all be expressed more simply in terms of the functions $`C(\mathrm{\Lambda }^2):`$
$$\{C_0(\mathrm{\Lambda }^2);C_\mu (\mathrm{\Lambda }^2);C_{\mu \nu }(\mathrm{\Lambda }^2)\}=d^4k\{1;k_\mu ;k_{\mu \nu }\}/D(\mathrm{\Lambda }^2)$$
(B21)
Using the identity
$$\frac{1}{k^2\lambda ^2}\left(\frac{\mathrm{\Lambda }^2}{k^2\mathrm{\Lambda }^2}\right)^m=\frac{(\mathrm{\Lambda }^2)^m}{(m1)!}T^{m1}\left\{\frac{1}{\mathrm{\Lambda }^2\lambda ^2}\left[\frac{1}{k^2\mathrm{\Lambda }^2}\frac{1}{k^2\lambda ^2}\right]\right\}$$
(B22)
with $`T\frac{}{(\mathrm{\Lambda }^2)}`$, we can write
$$\left\{I\right\}=N_m^{^{}}(\mathrm{\Lambda }^2)^mT^{m1}\left\{\frac{C(\mathrm{\Lambda }^2)C(\lambda ^2)}{\mathrm{\Lambda }^2\lambda ^2}\right\},$$
(B23)
where
$$N_m^{^{}}=\frac{(1)^m}{(m1)!(2\pi )^4}.$$
(B24)
In Eq. (B23) $`I`$ and $`C`$ denote any one of $`I_0,I_\mu ,I_{\mu \nu }`$ and $`C_0,C_\mu ,C_{\mu \nu }`$ respectively. We see from Eq. (B23) that terms in $`C(\mathrm{\Lambda }^2)`$ which are independent of $`\mathrm{\Lambda }^2`$ do not appear in the expression for $`I`$ . In particular, this applies to $`C_{\mu \nu }(\mathrm{\Lambda }^2)`$, which may be evaluated using either dimensional regularization or a convergence factor. The infinities in $`C_{\mu \nu }(\mathrm{\Lambda }^2)`$ are indeed independent of $`\mathrm{\Lambda }^2`$, thus giving a finite result for $`I_{\mu \nu }`$ as it should. In similar fashion, we have
$$\left\{J\right\}=N_m^{^{}}(\mathrm{\Lambda }^2)^mT^{m1}\left\{\frac{\mathrm{\Lambda }^2C(\mathrm{\Lambda }^2)\lambda ^2C(\lambda ^2)}{\mathrm{\Lambda }^2\lambda ^2}\right\}$$
(B25)
in which $`J`$ and $`C`$ denote any one of $`J_0,J_\mu ,J_{\mu \nu }`$ and $`C_0,C_\mu ,C_{\mu \nu }`$ respectively. We see from Eq. (B25) that any terms in $`C(\mathrm{\Lambda }^2)`$ which are independent of $`\mathrm{\Lambda }^2`$ do not appear in the expression for $`J`$ provided $`m>1`$. Finally, for $`K_0`$ we get
$$K_0=N_m^{^{}}(\mathrm{\Lambda }^2)^mT^{m1}\left\{\frac{\mathrm{\Lambda }^4C_0(\mathrm{\Lambda }^2)\lambda ^4C_0(\lambda ^2)}{\mathrm{\Lambda }^2\lambda ^2}\right\}$$
(B26)
We note that, apart from trivial factors, the integrals in (B21) are the three-point functions defined in , Eq.(5.1), and , Eq.(E.1); $`C_0`$ has been evaluated in terms of Spence functions in . The details of the algebra in and being rather lengthy, we choose instead to evaluate the integrals in (B21) using Feynman parameters, writing
$$\frac{1}{D(\mathrm{\Lambda }^2)}=2_0^1_0^1\frac{xdxdy}{[k^22xkp_y\mathrm{\Lambda }^2(1x)+iϵ]^3},$$
(B27)
where $`p_y=p_2y+p_4(1y)`$. Substituting Eq. (B27) in Eq. (B21) and shifting the integration variable $`k(kxp_yk)`$ we then have
$`\{C_0(\mathrm{\Lambda }^2);C_\mu (\mathrm{\Lambda }^2);C_{\mu \nu }(\mathrm{\Lambda }^2)\}`$ (B28)
$`=`$ $`i\pi ^2{\displaystyle _0^1}{\displaystyle _0^1}𝑑x𝑑y{\displaystyle \frac{\{x;x^2p_{y\mu };x^3p_{y\mu }p_{y\nu }+\frac{1}{2}g_{\mu \nu }\lambda ^2(x\frac{1}{2}x^2)\}}{x^2p_y^2+\mathrm{\Lambda }^2(1x)}}+\chi _{\mu \nu },`$ (B29)
where, throughout this section, we denote by $`\chi _{\mu \nu }`$ any terms which are independent of $`\mathrm{\Lambda }`$ . Rewriting $`p_y=\frac{1}{2}\rho +\frac{1}{2}q(12y)`$ in Eqs. (B27) and (B29) we may, neglecting terms which are independent of $`\mathrm{\Lambda }`$, express $`C_0,C_\mu `$ and $`C_{\mu \nu }`$ in terms of the functions
$$\varphi _k(\mathrm{\Lambda }^2)_0^1_0^1\frac{x^kdxdy}{p_y^2x^2+\mathrm{\Lambda }^2(1x)}.$$
(B30)
We get
$$C_0=i\pi ^2\varphi _1(\lambda ^2);C_\mu =i\pi ^2\frac{1}{2}\rho _\mu \varphi _2(\lambda ^2)$$
(B31)
$`C_{\mu \nu }`$ $`=`$ $`i\pi ^2[{\displaystyle \frac{1}{4}}\rho _\mu \rho _\nu \varphi _3(\lambda ^2){\displaystyle \frac{1}{4}}{\displaystyle \frac{\rho ^2}{q^2}}q_\mu q_\nu \varphi _3(\lambda ^2)+`$ (B33)
$`{\displaystyle \frac{q_\mu q_\nu }{q^2}}\lambda ^2[\varphi _1(\lambda ^2)\varphi _2(\lambda ^2)]+{\displaystyle \frac{1}{2}}g_{\mu \nu }\lambda ^2[\varphi _1(\lambda ^2){\displaystyle \frac{1}{2}}\varphi _2(\lambda ^2)]].`$
As shown in Appendix D, the functions $`\varphi _k`$ obey a three-term inhomogeneous recursion, which is used to calculate $`\varphi _k`$ for $`k>1`$:
$`(k+1)\rho ^2\varphi _{k+2}(\mathrm{\Lambda }^2)2(2k+1)\mathrm{\Lambda }^2\varphi _{k+1}(\mathrm{\Lambda }^2)+`$ $`4k\mathrm{\Lambda }^2\varphi _k(\mathrm{\Lambda }^2)`$ (B34)
$`={\displaystyle \frac{2\rho }{\rho _1}}\mathrm{ln}\left({\displaystyle \frac{\rho +\rho _1}{\rho \rho _1}}\right)`$ $`+2\mathrm{\Lambda }^2\left[\varphi _{k+1}^{(0)}(\mathrm{\Lambda }^2)2\varphi _k^{(0)}(\mathrm{\Lambda }^2)\right]`$ (B35)
Here
$$\varphi _k^{(0)}(\mathrm{\Lambda }^2)\varphi _k(\mathrm{\Lambda }^2)|_{q^2=0}=_0^1\frac{x^kdx}{M^2x^2+(1x)\mathrm{\Lambda }^2}$$
(B36)
The functions $`\varphi _k^{(0)}(\mathrm{\Lambda }^2)`$ may in turn be calculated from the recursion
$$M^2\varphi _{k+2}^{(0)}(\mathrm{\Lambda }^2)\mathrm{\Lambda }^2\varphi _{k+1}^{(0)}(\mathrm{\Lambda }^2)+\mathrm{\Lambda }^2\varphi _k^{(0)}(\mathrm{\Lambda }^2)=\frac{1}{k+1}$$
(B37)
To implement the recursions (B34) and (B37) we need
$$\varphi _0^{(0)}(\mathrm{\Lambda }^2)=\frac{1}{\mathrm{\Lambda }\mathrm{\Lambda }_1}\mathrm{ln}\left(\frac{\mathrm{\Lambda }+\mathrm{\Lambda }_1}{\mathrm{\Lambda }\mathrm{\Lambda }_1}\right)$$
(B38)
and
$$\varphi _1^{(0)}(\mathrm{\Lambda }^2)=\frac{1}{2M^2}\left[\mathrm{ln}\frac{M^2}{\mathrm{\Lambda }^2}+\frac{\mathrm{\Lambda }}{\mathrm{\Lambda }_1}\mathrm{ln}\left(\frac{\mathrm{\Lambda }+\mathrm{\Lambda }_1}{\mathrm{\Lambda }\mathrm{\Lambda }_1}\right)\right]$$
(B39)
which follow from Eq. (B36), and $`\varphi _1(\mathrm{\Lambda }^2)`$, which can be expressed in terms of dilogarithms (Spence functions)(see Appendix D):
$$\varphi _1(\mathrm{\Lambda }^2)=\frac{1}{\rho \rho _1}\left\{L\left(1\frac{1}{xy}\right)L\left(1\frac{x}{y}\right)2\mathrm{ln}(x)\mathrm{ln}(1+\frac{1}{y})\right\}$$
(B40)
where
$$L(z)=_0^z\frac{\mathrm{ln}(1t)}{t}𝑑t$$
(B41)
$$x=\frac{\rho +\rho _1}{\rho \rho _1}=\frac{(\rho +\rho _1)^2}{4M^2}$$
(B42)
$$y=\frac{\mathrm{\Lambda }+\mathrm{\Lambda }_1}{\mathrm{\Lambda }\mathrm{\Lambda }_1}=\frac{(\mathrm{\Lambda }+\mathrm{\Lambda }_1)^2}{4M^2}$$
(B43)
We will also want to take the limit $`\lambda 0`$. Neglecting all terms which vanish in this limit, we find,
$$\varphi _1(\lambda ^2)_{\stackrel{}{\lambda 0}}\frac{1}{\rho \rho _1}\left\{2L\left(\frac{1}{x}\right)\frac{\pi ^2}{6}\frac{1}{2}\mathrm{ln}^2x+\mathrm{ln}x\mathrm{ln}\left(\frac{\rho ^2}{\lambda ^2}\right)\right\}$$
(B44)
$$\varphi _1^{(0)}(\lambda ^2)_{\stackrel{}{\lambda 0}}\frac{1}{M^2}\mathrm{ln}\left(\frac{M}{\lambda }\right)$$
(B45)
and for $`k>1,`$
$$\varphi _k(0)=\frac{2}{(k1)\rho \rho _1}\mathrm{ln}x$$
(B46)
$$\varphi _k^{(0)}(0)=\frac{1}{(k1)M^2}$$
(B47)
We then find
$$g_0=N_1\varphi _1(\lambda ^2)+N_m(\mathrm{\Lambda }^2)^mT^{m1}\left\{\frac{1}{\mathrm{\Lambda }^2}\varphi _1(\mathrm{\Lambda }^2)\right\}$$
(B48)
$$g_1=\frac{1}{2}N_m(\mathrm{\Lambda }^2)^mT^{m1}\left\{\frac{1}{\mathrm{\Lambda }^2}\left[\varphi _2(\mathrm{\Lambda }^2)\varphi _2(0)\right]\right\}$$
(B49)
$`g_{11}`$ $`=`$ $`{\displaystyle \frac{1}{4}}N_m(\mathrm{\Lambda }^2)^mT^{m1}\left\{{\displaystyle \frac{1}{\mathrm{\Lambda }^2}}\left[\varphi _3(\mathrm{\Lambda }^2)\varphi _3(0)\right]\right\}`$ (B51)
$`+{\displaystyle \frac{1}{2\rho ^2}}N_m(\mathrm{\Lambda }^2)^mT^{m1}\left\{\varphi _1(\mathrm{\Lambda }^2){\displaystyle \frac{1}{2}}\varphi _2(\mathrm{\Lambda }^2)\right\}`$
$`g_{22}`$ $`=`$ $`{\displaystyle \frac{\rho ^2}{4q^2}}N_m(\mathrm{\Lambda }^2)^mT^{m1}\left\{{\displaystyle \frac{1}{\mathrm{\Lambda }^2}}\left[\varphi _3(\mathrm{\Lambda }^2)\varphi _3(0)\right]\right\}`$ (B53)
$`{\displaystyle \frac{1}{2q^2}}N_m(\mathrm{\Lambda }^2)^mT^{m1}\left\{\varphi _1(\mathrm{\Lambda }^2){\displaystyle \frac{3}{2}}\varphi _2(\mathrm{\Lambda }^2)\right\}`$
$$h_0=N_m(\mathrm{\Lambda }^2)^mT^{m1}\left\{\varphi _1(\mathrm{\Lambda }^2)\right\}$$
(B54)
$$h_1=\frac{1}{2}N_m(\mathrm{\Lambda }^2)^mT^{m1}\left\{\varphi _2(\mathrm{\Lambda }^2)\right\}$$
(B55)
$`h_{11}`$ $`=`$ $`{\displaystyle \frac{1}{4}}N_m(\mathrm{\Lambda }^2)^mT^{m1}\left\{\varphi _3(\mathrm{\Lambda }^2)\right\}`$ (B57)
$`+{\displaystyle \frac{1}{2\rho ^2}}N_m(\mathrm{\Lambda }^2)^mT^{m1}\left\{\mathrm{\Lambda }^2\left[\varphi _1(\mathrm{\Lambda }^2){\displaystyle \frac{1}{2}}\varphi _2(\mathrm{\Lambda }^2)\right]\right\}`$
$`h_{22}`$ $`=`$ $`{\displaystyle \frac{\rho ^2}{4q^2}}N_m(\mathrm{\Lambda }^2)^mT^{m1}\left\{\varphi _3(\mathrm{\Lambda }^2)\right\}`$ (B59)
$`{\displaystyle \frac{1}{2q^2}}N_m(\mathrm{\Lambda }^2)^mT^{m1}\left\{\mathrm{\Lambda }^2\left[\varphi _1(\mathrm{\Lambda }^2){\displaystyle \frac{3}{2}}\varphi _2(\mathrm{\Lambda }^2)\right]\right\}`$
$$k_0=N_m(\mathrm{\Lambda }^2)^mT^{m1}\left\{\mathrm{\Lambda }^2\varphi _1(\mathrm{\Lambda }^2)\right\}$$
(B60)
where
$$N_m=i\pi ^2N_m^{^{}}$$
(B61)
The terms $`ggg,gsg,\mathrm{}sss`$ can now be expressed more simply in terms of the functions $`\varphi _k`$. We get
$`ggg`$ $`=`$ $`ie^2F(q^2)\left\{\begin{array}{c}\begin{array}{c}2(2M^2q^2)N_1\varphi _1(\lambda ^2)\\ +2(2M^2q^2)\left[S^{m1}\left\{\frac{1}{\mathrm{\Lambda }^2}\varphi _1(\mathrm{\Lambda }^2)\right\}S^{m1}\left\{\frac{1}{\mathrm{\Lambda }^2}\left[\varphi _2(\mathrm{\Lambda }^2)\varphi _2(0)\right]\right\}\right]\end{array}\\ +S^{m1}\left\{\varphi _2(\mathrm{\Lambda }^2)\right\}4M^2S^{m1}\left\{\frac{1}{\mathrm{\Lambda }^2}\left[\varphi _3(\mathrm{\Lambda }^2)\varphi _3(0)\right]\right\}\end{array}\right\}\gamma _\mu `$ (B67)
$`ie^2F(q^2)\left\{4M^2S^{m1}\left\{{\displaystyle \frac{1}{\mathrm{\Lambda }^2}}\left[\varphi _2(\mathrm{\Lambda }^2)\varphi _2(0)\right]\right\}+4M^2S^{m1}\left\{{\displaystyle \frac{1}{\mathrm{\Lambda }^2}}\left[\varphi _3(\mathrm{\Lambda }^2)\varphi _3(0)\right]\right\}\right\}{\displaystyle \frac{i\sigma _{\mu \nu }q^\nu }{2M}}`$
$`gsg`$ $`=`$ $`ie^2\kappa F(q^2)\left\{q^2S^{m1}\left\{{\displaystyle \frac{1}{\mathrm{\Lambda }^2}}\left[\varphi _2(\mathrm{\Lambda }^2)\varphi _2(0)\right]\right\}\right\}\gamma _\mu `$ (B71)
$`ie^2\kappa F(q^2)\left\{2(2M^2q^2)\left[\begin{array}{c}N_1\varphi _1(\lambda ^2)+S^{m1}\left\{\frac{1}{\mathrm{\Lambda }^2}\varphi _1(\mathrm{\Lambda }^2)\right\}\\ S^{m1}\left\{\frac{1}{\mathrm{\Lambda }^2}\left[\varphi _2(\mathrm{\Lambda }^2)\varphi _2(0)\right]\right\}\end{array}\right]\right\}{\displaystyle \frac{i\sigma _{\mu \nu }q^\nu }{2M}}`$
$`ggs+sgg`$ $`=`$ $`2ie^2\kappa F(q^2)\left\{\begin{array}{c}(\frac{q^2}{4}3M^2)S^{m1}\left\{\frac{1}{\mathrm{\Lambda }^2}\left[\varphi _3(\mathrm{\Lambda }^2)\varphi _3(0)\right]\right\}\\ +\frac{3}{2}S^{m1}\left\{\varphi _1(\mathrm{\Lambda }^2)\frac{1}{2}\varphi _2(\mathrm{\Lambda }^2)\right\}\end{array}\right\}\gamma _\mu `$ (B77)
$`+2ie^2\kappa F(q^2)\left\{\begin{array}{c}2M^2S^{m1}\left\{\frac{1}{\mathrm{\Lambda }^2}\left[\varphi _3(\mathrm{\Lambda }^2)\varphi _3(0)\right]\right\}\\ 4S^{m1}\left\{\varphi _1(\mathrm{\Lambda }^2)\frac{1}{2}\varphi _2(\mathrm{\Lambda }^2)\right\}\end{array}\right\}{\displaystyle \frac{i\sigma _{\mu \nu }q^\nu }{2M}}`$
$`gss+ssg`$ $`=`$ $`2ie^2\kappa ^2F(q^2){\displaystyle \frac{q^2}{4M^2}}\left\{\begin{array}{c}2M^2S^{m1}\left\{\frac{1}{\mathrm{\Lambda }^2}\left[\varphi _3(\mathrm{\Lambda }^2)\varphi _3(0)\right]\right\}\\ 2S^{m1}\left\{\varphi _1(\mathrm{\Lambda }^2)\frac{1}{2}\varphi _2(\mathrm{\Lambda }^2)\right\}\end{array}\right\}\gamma _\mu `$ (B83)
$`+2ie^2\kappa ^2F(q^2)\left\{\begin{array}{c}(\frac{3}{4}q^2M^2)S^{m1}\left\{\frac{1}{\mathrm{\Lambda }^2}\left[\varphi _3(\mathrm{\Lambda }^2)\varphi _3(0)\right]\right\}\\ \frac{1}{2}S^{m1}\left\{\varphi _1(\mathrm{\Lambda }^2)\frac{1}{2}\varphi _2(\mathrm{\Lambda }^2)\right\}\end{array}\right\}{\displaystyle \frac{i\sigma _{\mu \nu }q^\nu }{2M}}`$
$`sgs`$ $`=`$ $`ie^2\kappa ^2F(q^2)\left\{\begin{array}{c}2M^2S^{m1}\left\{\frac{1}{\mathrm{\Lambda }^2}\left[\varphi _3(\mathrm{\Lambda }^2)\varphi _3(0)\right]\right\}\frac{1}{2}S^{m1}\left\{\varphi _3(\mathrm{\Lambda }^2)\right\}\\ +S^{m1}\left\{\varphi _1(\mathrm{\Lambda }^2)\right\}+\frac{1}{2}S^{m1}\left\{\varphi _2(\mathrm{\Lambda }^2)\right\}\\ \frac{1}{2M^2}S^{m1}\left\{\mathrm{\Lambda }^2\left[\varphi _1(\mathrm{\Lambda }^2)\frac{1}{4}\varphi _2(\mathrm{\Lambda }^2)\right]\right\}\end{array}\right\}\gamma _\mu `$ (B90)
$`ie^2\kappa ^2F(q^2)\left\{\begin{array}{c}2M^2S^{m1}\left\{\frac{1}{\mathrm{\Lambda }^2}\left[\varphi _3(\mathrm{\Lambda }^2)\varphi _3(0)\right]\right\}\\ 2S^{m1}\left\{\varphi _1(\mathrm{\Lambda }^2)\right\}+\frac{1}{2}S^{m1}\left\{\varphi _3(\mathrm{\Lambda }^2)\right\}\end{array}\right\}{\displaystyle \frac{i\sigma _{\mu \nu }q^\nu }{2M}}`$
$`sss`$ $`=`$ $`ie^2\kappa ^3F(q^2){\displaystyle \frac{q^2}{4M^2}}\left\{\begin{array}{c}2M^2S^{m1}\left\{\frac{1}{\mathrm{\Lambda }^2}\left[\varphi _3(\mathrm{\Lambda }^2)\varphi _3(0)\right]\right\}+\frac{1}{2}S^{m1}\left\{\varphi _3(\mathrm{\Lambda }^2)\right\}\\ +2S^{m1}\left\{\varphi _1(\mathrm{\Lambda }^2)\varphi _2(\mathrm{\Lambda }^2)\right\}\end{array}\right\}\gamma _\mu `$ (B97)
$`ie^2\kappa ^3F(q^2)\left\{\begin{array}{c}\frac{1}{2}q^2S^{m1}\left\{\frac{1}{\mathrm{\Lambda }^2}\left[\varphi _3(\mathrm{\Lambda }^2)\varphi _3(0)\right]\right\}(1+\frac{q^2}{2M^2})S^{m1}\left\{\varphi _1(\mathrm{\Lambda }^2)\right\}\\ \frac{1}{2}S^{m1}\left\{\varphi _3(\mathrm{\Lambda }^2)\right\}\\ +\frac{1}{2}(1+\frac{q^2}{M^2})S^{m1}\left\{\varphi _2(\mathrm{\Lambda }^2)\right\}\frac{1}{4M^2}S^{m1}\left\{\mathrm{\Lambda }^2\left[\varphi _1(\mathrm{\Lambda }^2)\varphi _2(\mathrm{\Lambda }^2)\right]\right\}\end{array}\right\}{\displaystyle \frac{i\sigma _{\mu \nu }q^\nu }{2M}}`$
where
$$S^{m1}=N_m(\mathrm{\Lambda }^2)^mT^{m1}$$
(B98)
It should be noted that the terms with $`\varphi _1(\lambda ^2)`$, which appear only in $`ggg`$ and $`gsg`$, constitute the well-known infrared divergence. They are, apart from the hard-photon proton interaction, (5), independent of the proton form factor (in this case independent of $`\mathrm{\Lambda }`$ and $`M`$). This is to be expected, since this term is cancelled by a similar infrared divergent term coming from the cross section for the emission of a real soft photon, which is given by the elastic cross section multiplied by a factor independent of the proton form factor.
## C Electron vertex correction
As noted in Sec. IIIC, the electron vertex correction, $`M_5`$, may be obtained directly from the proton vertex correction, $`M_6`$, if we retain only the term $`ggg`$, set $`F=1`$, replace $`p_2,p_4`$ and $`M`$ by $`p_1,p_3`$ and $`m`$ and take the limit $`\mathrm{\Lambda }\mathrm{}`$. After making these replacements in $`ggg`$ as given in (B67), we need $`\varphi _k(\mathrm{\Lambda }^2)`$ to order $`m^2/\mathrm{\Lambda }^2`$ and $`\varphi _k(0)`$ (see (B46)). From (B40), writing (, p. 389 (B.2))
$`L(1z)=L(z)+{\displaystyle \frac{1}{6}}\pi ^2\mathrm{ln}z\mathrm{ln}(1z)`$and neglecting terms of relative order $`m^2/\mathrm{\Lambda }^2`$, we have
$$\varphi _1(\mathrm{\Lambda }^2)_\stackrel{}{\mathrm{\Lambda }\mathrm{}}\frac{1}{\mathrm{\Lambda }^2}\left\{1+\mathrm{ln}\left(\frac{\mathrm{\Lambda }^2}{m^2}\right)\frac{\rho _m}{\rho _1}\mathrm{ln}x_m\right\}$$
(C1)
and from (B30)
$$\varphi _k(\mathrm{\Lambda }^2)\varphi _{k+1}(\mathrm{\Lambda }^2)_\stackrel{}{\mathrm{\Lambda }\mathrm{}}\frac{1}{\mathrm{\Lambda }^2}\left\{\frac{1}{k+1}\right\}$$
(C2)
Choosing $`m=1`$ in (B67) we then have
$$ggg=\left[G_1^{(e)}(q^2)\gamma _\mu +G_2^{(e)}(q^2)\frac{i\sigma _{\mu \nu }q^\nu }{2m}\right]$$
(C3)
where
$$G_1^{(e)}(q^2)=\frac{\alpha }{4\pi }\left\{2(2m^2q^2)\varphi _1(\lambda ^2)+\left(\frac{3\rho _m^24m^2}{\rho _m\rho _1}\right)\mathrm{ln}x_m+\frac{1}{2}+\mathrm{ln}\left(\frac{\mathrm{\Lambda }^2}{m^2}\right)\right\}$$
(C4)
$$G_2^{(e)}(q^2)=\frac{\alpha }{4\pi }\left\{\frac{4m^2}{\rho _m\rho _1}\mathrm{ln}x_m\right\}$$
(C5)
Adding the contribution of the electron self energy diagrams gives
$$\overline{ggg}=\left[\left(G_1^{(e)}(q^2)G_1^{(e)}(0)\right)\gamma _\mu +G_2^{(e)}(q^2)\frac{i\sigma _{\mu \nu }q^\nu }{2m}\right]$$
(C6)
To facilitate comparison with , we write this in terms of the functions $`K(p_i,p_j)`$. Similarly to (20), we now have
$$(2m^2q^2)\varphi _1(\lambda ^2)=K(p_1,p_3)$$
(C7)
which then gives (31).
## D The functions $`\mathit{\varphi }_𝒌\mathbf{\left(}𝚲^\mathrm{𝟐}\mathbf{\right)}`$
In this Appendix we derive the three-term recurrence relation for the function $`\varphi _k(\mathrm{\Lambda }^2)`$ given in (B34) as well as the expression for the function $`\varphi _1(\mathrm{\Lambda }^2)`$, defined in (B30) and given in terms of Spence functions in (B40). Integrating first over $`y`$ in Eq. (B30) we have
$$\varphi _k(\mathrm{\Lambda }^2)=\frac{1}{\rho _1}_0^1\frac{x^{k1}}{R}\mathrm{log}\left\{\frac{R+x\rho _1}{Rx\rho _1}\right\}𝑑x,$$
(D1)
where $`R^2=\rho ^2x^2+4(1x)\mathrm{\Lambda }^2`$ and $`\rho _1^2=q^2>0`$. Noting that
$$\frac{d}{dx}\left\{x^kR\right\}=x^{k1}\left\{kR^2+\rho ^2x^22x\mathrm{\Lambda }^2\right\}R^1,$$
(D2)
we get
$$(k+1)\rho ^2\varphi _{k+2}2(2k+1)\mathrm{\Lambda }^2\varphi _{k+1}+4k\mathrm{\Lambda }^2\varphi _k=\frac{2}{\rho _1}_0^1\mathrm{log}\left\{\frac{R+x\rho _1}{Rx\rho _1}\right\}d\left(x^kR\right).$$
(D3)
Integration by parts then gives
$$(k+1)\rho ^2\varphi _{k+2}2(2k+1)\mathrm{\Lambda }^2\varphi _{k+1}+4k\mathrm{\Lambda }^2\varphi _k=\frac{2\rho }{\rho _1}\mathrm{log}\left(\frac{\rho +\rho _1}{\rho \rho _1}\right)2\mathrm{\Lambda }^2_0^1\frac{x^k(2x)}{M^2x^2+(1x)\mathrm{\Lambda }^2}$$
(D4)
from which the recursion (B34) follows at once, using Eq. (B36)
From Eq. (D4) it is clear that $`\varphi _2`$ and $`\varphi _3`$ follow once we have evaluated $`\varphi _1`$, which follows. Setting $`k=1`$ in (B30) and integrating first over $`x`$, we get
$$_0^1\frac{xdx}{p_y^2x^2+\mathrm{\Lambda }^2(1x)}=\frac{1}{2p_y^2}\left[\mathrm{ln}\left(\frac{p_y^2}{\mathrm{\Lambda }^2}\right)+\frac{\mathrm{\Lambda }}{\mathrm{\Delta }}\mathrm{ln}\left(\frac{\mathrm{\Lambda }+\mathrm{\Delta }}{\mathrm{\Lambda }\mathrm{\Delta }}\right)\right],$$
(D5)
where $`\mathrm{\Delta }^2=\mathrm{\Lambda }^24p_y^2`$. We next make the change of variable $`y=(1+\omega )/2`$, which gives $`\mathrm{\Delta }^2=\rho _1^2\omega ^2+\mathrm{\Lambda }^2\rho ^2`$, and then make the further change of variable $`\mathrm{\Delta }=\rho _1\omega +s`$, from which
$$\omega =\frac{\mathrm{\Lambda }^2\rho ^2s^2}{2\rho _1s};\mathrm{\Delta }=\frac{\mathrm{\Lambda }^2\rho ^2+s^2}{2s}$$
(D6)
Then integrating (B30) over $`y`$ gives
$`\varphi _1(\mathrm{\Lambda }^2)`$ $`=`$ $`{\displaystyle \frac{2}{\rho _1}}{\displaystyle _s_{}^{s_+}}{\displaystyle \frac{ds}{\rho ^2(\mathrm{\Lambda }s)^2}}\mathrm{ln}\left[{\displaystyle \frac{(\mathrm{\Lambda }+s)^2\rho ^2}{4s\mathrm{\Lambda }}}\right]`$ (D8)
$`{\displaystyle \frac{2}{\rho _1}}{\displaystyle _s_{}^{s_+}}{\displaystyle \frac{ds}{(\mathrm{\Lambda }+s)^2\rho ^2}}\mathrm{ln}\left[{\displaystyle \frac{\rho ^2(\mathrm{\Lambda }s)^2}{4s\mathrm{\Lambda }}}\right]`$
where
$$s_\pm =\mathrm{\Lambda }_1\pm \rho _1$$
(D9)
Factoring the expressions which appear as factors to the logarithms as well as in their arguments, we can write
$$\varphi _1(\mathrm{\Lambda }^2)=\frac{1}{\rho \rho _1}\underset{i=1}{\overset{9}{}}I_i$$
(D10)
where
$`I_1={\displaystyle _s_{}^{s_+}}{\displaystyle \frac{ds}{s\sigma _{}}}\mathrm{ln}(s+\sigma _+);\text{ }I_2={\displaystyle _s_{}^{s_+}}{\displaystyle \frac{ds}{s+\sigma _+}}\mathrm{ln}(s\sigma _{})`$
$`I_3={\displaystyle _s_{}^{s_+}}{\displaystyle \frac{ds}{\sigma _+s}}\mathrm{ln}(s+\sigma _{});\text{ }I_4={\displaystyle _s_{}^{s_+}}{\displaystyle \frac{ds}{s+\sigma _{}}}\mathrm{ln}(\sigma _+s)`$
$$I_5=_s_{}^{s_+}\frac{ds}{s\sigma _{}}\mathrm{ln}\left(\frac{s+\sigma _{}}{2s}\right);\text{ }I_8=_s_{}^{s_+}\frac{ds}{s+\sigma _{}}\mathrm{ln}\left(\frac{s\sigma _{}}{2s}\right)$$
(D11)
$`I_6={\displaystyle _s_{}^{s_+}}{\displaystyle \frac{ds}{s+\sigma _+}}\mathrm{ln}\left({\displaystyle \frac{\sigma _+s}{2s}}\right);\text{ }I_7={\displaystyle _s_{}^{s_+}}{\displaystyle \frac{ds}{\sigma _+s}}\mathrm{ln}\left({\displaystyle \frac{s+\sigma _+}{2s}}\right)`$
$`I_9={\displaystyle _s_{}^{s_+}}𝑑s\left[{\displaystyle \frac{1}{s\sigma _{}}}+{\displaystyle \frac{1}{\sigma _+s}}{\displaystyle \frac{1}{\overline{s}+\sigma _{}}}+{\displaystyle \frac{1}{s+\sigma _+}}\right]\mathrm{ln}(2\mathrm{\Lambda })`$where
$$\sigma _\pm =\mathrm{\Lambda }\pm \rho $$
(D12)
Some of the integrals $`I_i`$ can be integrated out directly. For these terms we get:
$$\underset{n=1}{\overset{4}{}}I_n=\mathrm{ln}\left(\frac{\alpha _{}}{\alpha _+}\right)ln\left(\frac{4\sigma _+\sigma _{}}{(1\alpha _+^2)(1\alpha _{}^2)}\right)\mathrm{ln}\alpha _+\mathrm{ln}(1\alpha _+^2)+\mathrm{ln}\alpha _{}\mathrm{ln}(1\alpha _{}^2)$$
(D13)
and
$$I_9=2\mathrm{ln}\left(\frac{\alpha _{}}{\alpha _+}\right)\mathrm{ln}(2\mathrm{\Lambda })$$
(D14)
where
$`\alpha _\pm ={\displaystyle \frac{\rho \rho _1}{\mathrm{\Lambda }+\mathrm{\Lambda }_1}}`$The remaining terms in Eq. (D10) we may rewrite as
$$\underset{n=5}{\overset{8}{}}I_n=\mathrm{ln}\alpha _+\mathrm{ln}(1\alpha _+^2)\mathrm{ln}\alpha _{}\mathrm{ln}(1\alpha _{}^2)+L(1\alpha _+^2)L(1\alpha _{}^2)$$
(D15)
where $`L`$ is the dilogarithm (Spence) function. We then get the result given in (B40) in Appendix B.
## E Final electron detector acceptance
In this Appendix, we express $`\mathrm{\Delta }ϵ`$, the maximum momentum of the photon in the frame $`S^0`$, in terms of the final electron detector acceptance in the lab frame, $`\mathrm{\Delta }E`$. In $`S^0`$ (p$`{}_{4}{}^{}+`$ k $`=0`$), if $`|`$k$`|`$ $`=\mathrm{\Delta }ϵM`$, we have, from $`(p_1+p_2p_3)^2=(p_4+k)^2,`$ neglecting terms of order $`(\mathrm{\Delta }ϵ/M)^2`$ and $`(m/M)^2`$,
$$p_2(p_1p_3)p_1p_3=M\mathrm{\Delta }ϵ$$
(E1)
Writing this in terms of lab frame energies, we have, for high energies,
$$M(ϵ_1ϵ_3)ϵ_1ϵ_3(1\mathrm{cos}\theta )=M\mathrm{\Delta }ϵ$$
(E2)
For elastic scattering in the lab frame, we have
$$M(ϵ_1ϵ_3^{el})ϵ_1ϵ_3^{el}(1\mathrm{cos}\theta )=0$$
(E3)
Subtracting gives
$$\mathrm{\Delta }E\left(1+\frac{ϵ_1}{M}(1\mathrm{cos}\theta )\right)=\mathrm{\Delta }ϵ$$
(E4)
where
$$\mathrm{\Delta }E=ϵ_3^{el}ϵ_3$$
(E5)
Thus, in terms of lab frame quantities we have
$$\mathrm{\Delta }ϵ=\eta \mathrm{\Delta }E$$
(E6)
## F High energy approximation for $`S_{ij}^{(2)}`$
In this appendix we give the high energy approximation of the terms $`S_{ij}^{(2)}`$ defined in (83), in which we note in particular that for $`i=1`$ or $`3`$ we have $`lt=(\alpha p_ip_j)t\alpha p_it`$. Using transformations of the dilog (Spence) functions , p. 389 (B.3),
$$L(z)=L\left(\frac{1}{z}\right)\frac{1}{6}\pi ^2\frac{1}{2}\mathrm{ln}^2(z)$$
(F1)
$$L(z)=L\left(\frac{z}{z1}\right)\frac{1}{2}\mathrm{ln}^2\left(1z\right)$$
(F2)
the terms in $`S_{ij}^{(2)}`$ simplify considerably. We then obtain
$`S_{12}^{(2)}`$ $`=`$ $`\mathrm{ln}^2\left({\displaystyle \frac{2ϵ_3}{m}}\right)\mathrm{ln}^2x+{\displaystyle \frac{1}{2}}\mathrm{ln}^2\left({\displaystyle \frac{x}{\eta }}\right){\displaystyle \frac{1}{6}}\pi ^2`$ (F4)
$`L\left(1{\displaystyle \frac{1}{x\eta }}\right)+L\left(1{\displaystyle \frac{\eta }{x}}\right)`$
$`S_{32}^{(2)}`$ $`=`$ $`\mathrm{ln}^2\left({\displaystyle \frac{2ϵ_1}{m}}\right)\mathrm{ln}^2\left(x\right)+{\displaystyle \frac{1}{2}}\mathrm{ln}^2\left(x\eta \right){\displaystyle \frac{1}{6}}\pi ^2`$ (F6)
$`+L\left(1{\displaystyle \frac{1}{x\eta }}\right)L\left(1{\displaystyle \frac{\eta }{x}}\right)`$
$$S_{14}^{(2)}=\mathrm{ln}^2\left(\frac{2ϵ_3}{m}\right)\frac{1}{6}\pi ^2$$
(F7)
$$S_{34}^{(2)}=\mathrm{ln}^2\left(\frac{2ϵ_1}{m}\right)\frac{1}{6}\pi ^2$$
(F8)
$`S_{13}^{(2)}`$ $`=`$ $`\mathrm{ln}^2\left({\displaystyle \frac{2ϵ_1}{m}}\right)\mathrm{ln}^2\left({\displaystyle \frac{2ϵ_3}{m}}\right){\displaystyle \frac{1}{3}}\pi ^2+{\displaystyle \frac{1}{2}}\mathrm{ln}^2\left(\mathrm{cos}^2{\displaystyle \frac{1}{2}}\theta \right)`$ (F10)
$`+L\left(\mathrm{cos}^2{\displaystyle \frac{1}{2}}\theta \right)`$
$$S_{24}^{(2)}=\frac{1}{2}\mathrm{ln}^2\left(x\right)+\frac{1}{2}L\left(1\frac{1}{x^2}\right)$$
(F11)
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# Structure of the Sagittarius dwarf galaxy at low Galactic latitudes Based on observations obtained at the European Southern Observatory, La Silla, Chile
## 1 Introduction
The Sagittarius dwarf galaxy is the closest known member of the Local Group orbiting around the Milky Way ($``$25 kpc from the sun, $``$16 kpc from the Galactic Centre), but as a consequence of its location behind the Galactic Centre, it has been discovered only recently (Ibata, Gilmore, Irwin 1994, 1995). Since this discovery it turned out that Sgr presents typical features of a dwarf spheroidal: domination of an old ($``$10 Gyr) metal poor stellar population (Mateo et al. muskkk (1995); Fahlman et al. fahlman (1996); Marconi et al. marconi (1998); Bellazzini et al bfb1 (1999)) and absence of gas (Burton & Lockman bl (1999)). Its highest surface density region is centred on the Globular Cluster M54 (l=5.6, b=-14.0) and it is oriented roughly perpendicular to the Galactic plane so that its Northern extension (in Galactic coordinates) is completely hidden by the MW.
The mapping of Sgr is difficult to achieve because of the combination of its low surface brightness ($`\mu _V25.5`$ mag.arcsec<sup>-2</sup>), contamination by foreground Galactic stars and its large spatial extent (at least 22$`{}_{}{}^{}\times `$8) (Ibata et al. iwgis (1997), hereafter IWGIS). Evidence for the presence of Sgr has been established over 45from b$`3^{}`$ (Alard a96 (1996), hereafter A96; Alcock et al. 1997, hereafter Alc97) down to b$`48^{}`$ (Mateo et al. mom (1998), hereafter MOM), but it is difficult to assess whether these regions still correspond to the main body of Sgr or if we are merely encountering tidal debris (as suggested by Johnston et al. johnston99 (1999)). IWGIS proposed a map of the Southern part of Sgr based on the spatial distribution of the bright main sequence stars in Sgr and covering an area of $`150`$ deg<sup>2</sup> from $`b11^{}`$ down to $`b26^{}`$. However, their method based on statistical decontamination fails at low Galactic latitudes ($`|`$b$`|`$12) where differential reddening and high density of foreground stars (only $``$1 star in 1 000 is in Sgr in these regions) preclude any reliable decontamination, leaving the structure of the Northern extension of Sgr almost unknown. To this point, the detection of RR Lyrae constitutes an essential tool to trace the structure of Sgr in these regions as they can be clearly separated from the RR Lyrae of the MW. This method has already proven successful and $`350`$ RRab were detected between b=-10and b=-4(A96; Alc97). However, a connection between these stars and the centre of Sgr was necessary in order to offer a clear vision of this important region strongly interacting with the MW.
In this paper we report the detection of $``$1 500 RRab members of Sgr and located in its Northern extension. We present a surface density map of Sgr covering $``$50 deg<sup>2</sup> between b=-14and b=-4, based on the spatial distribution of these variables.
The paper is organized as follows : in section 2 we present our data (observations and reduction). Section 3 is devoted to the description of the selection process of RR Lyrae stars as well as a study of its completeness. We then describe the structure of Sgr (section 4). Finally we summarize our results and conclude in section 5.
## 2 Data
### 2.1 Observations
The data discussed in this paper consist of two sets of photographic plates and films taken with the ESO 1m Schmidt telescope at La Silla Observatory (see Table 1), each of them covering an area of $`25`$ deg<sup>2</sup> on the sky.
The first set of plates was part of the DUO project aimed at detecting microlensing events towards the Galactic Bulge (Alard & Guibert a97 (1997), hereafter AG97). This field, centred on Galactic coordinates (l=3.1, b=-7.1), has already been processed and presented in A96. The second set is new and includes 69 films centred on a field shifted towards the centre of density of the Sagittarius dwarf galaxy and slightly overlapping with the former (l=6.6, b=-10.8). Throughout the remainder of this paper, we will call the first field DUO field while the new field will be referred to as the SAG field.
### 2.2 Data reduction
The plates were scanned at CAI/Paris Observatory with the high speed microdensitometer MAMA<sup>1</sup><sup>1</sup>1MAMA (http://dsmama.obspm.fr) is operated by INSU (Institut National des Sciences de l’Univers) and Observatoire de Paris. (Machine Automatique à Mesurer pour l’Astronomie), yielding images with a pixel size of 10$`\mu `$m ($`0.6\mathrm{}`$.).
The photometric reduction has been performed with the software Extractor written by Alard. The process is as follows: first a reference catalogue is extracted from a plate of good quality (seeing$`<1\mathrm{}`$). For all the other plates, a new extraction is performed (implying a new detection of each object) and the new catalogue associated to the reference catalogue. The light curves were built in this way plate by plate and stored in a database. For more details on the photometric reduction process see AG97. The final sample contains light curves for $`14.10^6`$ stars in the DUO field and $`6.10^6`$ stars in the SAG field.
### 2.3 Photometry
#### 2.3.1 Calibration
The DUO field has been calibrated with a CCD sequence taken at the ESO/Danish 1.5 m telescope at La Silla. The photometric system for this field is $`\mathrm{B}_\mathrm{J}=\mathrm{B}0.28(\mathrm{B}\mathrm{V})`$ (Blair & Gilmore bg (1982)) . The Emulsion/Filter combination was different for the SAG field and consisted of a Kodak Tech-Pan 4415 emulsion together with a BG12 Filter. The Tech-Pan 4415 emulsion is an extremely fine-grained, high resolution film with a sensitivity extending to 0.69 $`\mu `$m. For more informations about the 4415 emulsion, see Phillipps & Parker (pp (1993)) and references therein. We were not aware of any photometric relation published for the band used in SAG .
The resulting bandpasses for both fields are shown on Fig. 1. The calibration in SAG has been performed with a sequence provided by Sarajedini & Layden (sl (1995)), located $`12\mathrm{}`$ north of the globular cluster M54 and consisting of 1638 stars calibrated in V and I bands. A polynomial fit to the instrumental magnitude of these stars yielded the photometric system $`\mathrm{B}_\mathrm{i}=\mathrm{V}+1.47(\mathrm{V}\mathrm{I})`$, where B<sub>i</sub> stands for the magnitude in the color band used in SAG . The scatter about this relation is 0.17 mag.
#### 2.3.2 Correction for extinction
The reddening has been estimated separately for each field. For the DUO field we used the well known property that the color of RR Lyrae stars at minimum magnitude is approximately constant and depends only slightly on the period and the metallicity (Sturch, sturch (1966)). A reddening map has been estimated for this field by computing the mean colors (B<sub>J</sub>-R) of RRab in small regions of 10$`\mathrm{}\times `$10$`\mathrm{}`$. The corrected magnitude is $`\mathrm{B}_{J_0}=\mathrm{B}_\mathrm{J}2.84\mathrm{E}(\mathrm{B}\mathrm{R})`$ (Wesselink wes (1987)). For the SAG field, where no color information was available, we used the extinction map of Schlegel et al. (sfd (1998), hereafter SFD) which provides reddening estimates with a precision of 16$`\%`$ for $`|`$b$`|>`$10. However, the SAG field extends to b$``$ -8where, according to SFD, the reddening map might become inaccurate. From the relation $`\mathrm{E}(\mathrm{V}\mathrm{R})=0.74\mathrm{E}(\mathrm{B}\mathrm{V})`$ (Cardelli et al. car (1989); hereafter CCM) we derive $`\mathrm{E}(\mathrm{B}_\mathrm{J}\mathrm{R})/\mathrm{E}(\mathrm{B}\mathrm{V})=1.46`$. This ratio in the overlap between DUO and SAG yields 1.32 $`\pm `$ 0.24, in reasonable agreement with the theoretical expectation, showing that even at the western edge of the SAG field the SFD map provides a satisfactory estimation for the extinction. Assuming E(V-I)=1.55 E(B-V) from CCM and a normal extinction law $`\mathrm{A}_\mathrm{V}`$=3.10 E(B-V) we obtain $`\mathrm{B}_{\mathrm{i}_0}=\mathrm{B}_i`$ \- 5.38 E(B-V) for the de-reddened magnitude in the SAG field.
## 3 Detection of RRab
### 3.1 The selection process
We will describe here the selection process of RRab stars. For the sake of homogeneity, we reprocessed the stars of the DUO field, using the same selection criteria as for the SAG field. The search for RRab in DUO has been performed through the B<sub>J</sub> band.
A first selection was performed by calculating the $`\chi ^2`$ about the mean magnitude ($`\chi _{mean}^2`$) for each light curve. Stars with $`\chi _{mean}^2>`$ 8 were then searched for periodicity. This cut should select all variables with an amplitude $``$ 0.3 mag. A first estimate of the period was done with the string minimization method of Renson (renson (1978)). A more accurate period was then searched in a small window spanning 0.1 day around the first estimate, using a multi-harmonic periodogram method (Schwarzenberg-Czerny czerny (1996)). The next step was to fit a Fourier series (with up to five harmonics) to the folded light curve:
$$B_\mathrm{i}=A_0+\underset{n=1}{\overset{n5}{}}A_n\mathrm{cos}(n\omega t+\varphi _n)$$
The $`\chi _{fit}^2`$ about the fitted light curve was then calculated and all the stars for which $`\chi _{ratio}=\sqrt{\chi _{mean}^2/\chi _{fit}^2}>2`$ have been selected as variable stars. At this step of the process the sample contained $``$ 7 000 variables.
The selection for RR Lyrae stars has been performed through the Fourier coefficients: for each variable we calculated the ratio of the amplitude of the first harmonic relative to the amplitude of the fundamental harmonic $`R_{21}=A_2/A_1`$, and the phase difference $`\varphi _{21}=\varphi _22\varphi _1`$. Fig.2 shows a plot of $`R_{21}`$ versus $`\varphi _{21}`$ for all stars satisfying $`\chi _{mean}^2>8`$ and $`\chi _{ratio}>2`$.
Several clumps lie in this figure. The most obvious one is located at $`R_{21}`$ 0.45 , $`\varphi _{21}`$ 0.7. This clump corresponds to RR Lyrae stars of Bailey type ab (hereafter RRab). For lower values of $`R_{21}`$ (i.e. for more symmetric light curves) we can distinguish two other clumps: one centred on ($`R_{21}`$0.2, $`\varphi _{21}`$3.2) and a shallower one at ($`R_{21}`$0.15, $`\varphi _{21}`$ 1.75), corresponding respectively to contact binaries and RR Lyrae of Bailey type c (RRc). A faint strip across the plot at $`\varphi _{21}3.1`$ is also visible and represents eclipsing binaries of Algol type. The selection of the RRab has been made with an ellipse centred on the clump (see Fig. 2) and finally a cut on periods ($`0.40^d>P>0.85^d`$) has been applied. The final sample contains $``$ 3 000 RRab.
The selected RRab may belong either to the MW or to the Sagittarius dwarf galaxy and we separated them through their distance modulus, assuming absolute magnitudes M$`_{B_\mathrm{J}}`$=0.79 (Wesselink wes (1987)) and M<sub>V</sub>=0.6 (Mateo et al muskkk (1995)). Furthermore, we take the mean color (V-I)<sub>0</sub>=0.46$`\pm `$0.06 after averaging over 27 RRab covering a wide range of metallicities from Table 1 of McNamara (1997). The apparent magnitude of each RRab has been estimated with the constant term of the Fourier series. Taking into account errors on the absolute magnitudes, apparent magnitudes, extinction and colors of RRab, the error on a single distance modulus is $``$0.3 mag in both field, the main source of uncertainty coming from extinction. Fig.3 shows the histogram of distance modulus for both fields before and after correction for extinction.
The histograms were smoothed by estimating the mean magnitude every 0.1 mag in a 0.3 mag bin. Both histograms exhibit similar features: a broad bump centred on (m-M)$`{}_{0}{}^{}`$14.5 (8 kpc) corresponding to RRab of the MW, and a sharp bump centred on (m-M)$`{}_{0}{}^{}`$16.9 (24 kpc) representing RRab members of Sgr. According to current models of RRab densities in the Halo the Galactic contribution to the histograms for (m-M)$`{}_{0}{}^{}>`$16.3 should be no more than 5-10$`\%`$ (Wetterer & McGraw wet (1996)).
The 2D spatial distribution of all the RRab with a distance modulus greater than 16.3 is displayed on Fig.4. This map includes $``$ 1 500 RRab.
The eastern and western box represents respectively the SAG field and the DUO field. The total area covered is about $`50`$ deg<sup>2</sup>, and comprises the globular cluster M54 at (l=5.5,b=-14.0) which is associated to Sgr and located in its highest density region. The image of M54 is completely saturated until $``$1.5 half mass radius on our plates, thus we do not expect this globular cluster to contribute significantly to our RRab sample. The spatial distribution of RRab reveals a density gradient in the SE-NW direction. We also show in Fig. 4 the RRab discovered by the MACHO team (Alc97), confirming that these stars are the continuation of Sgr
### 3.2 Completeness
There are two steps where the completeness of the RRab sample might be affected: first, the detection of stars becomes difficult towards the Galactic Centre because of the increasing stellar density, and some RRab blended by a neighbouring stars are missed. Second, we might miss some RRab during the selection process.
#### 3.2.1 Completeness of the extraction process
To quantify the loss induced by the first effect we simulated a set of 250 000 artificial stars with the same apparent magnitude than the detected RRab. These stars were then injected in small regions of 10$`\mathrm{}\times `$ 10$`\mathrm{}`$ uniformly spread over the fields and we tried to retrieve them with the same detection process as for the real stars. The lower panel of Fig.5 displays the fraction of stars re-detected as a function of Galactic latitude.
Filled circles are stars injected onto the SAG field whereas crosses are stars simulated in the DUO field. The dispersion reflects mainly the dependence on Galactic longitude. One can see that while the fraction of stars re-detected in SAG stays at a high level (above 90$`\%`$) and varies slowly, this is not the case in DUO where this fraction drops abruptly to reach 40$`\%`$ at b$``$-5. The reason for the higher variation rate in DUO is that the stellar density gradient increases at a higher rate towards the Galactic Centre (the mean density gradient is $``$15 stars.arcmin<sup>-2</sup>.deg<sup>-1</sup> in SAG and $``$25 stars.arcmin<sup>-2</sup>.deg<sup>-1</sup> in DUO ). Another feature visible on Fig.5b is that the loss induced by crowding is intrinsically higher in DUO than in SAG as can be seen in the range -10$`{}_{}{}^{}<`$b$`<`$-8($``$10$`\%`$ offset). This can be explained by the lower resolution of the III<sub>aJ</sub> emulsion in DUO relative to the finer grained 4415 emulsion in SAG (Parker & Malin pm (1999)). Furthermore, the lower extinction in DUO (A$`_{B_\mathrm{J}}`$/A$`{}_{B_\mathrm{i}}{}^{}`$0.7) results in a higher number of stars detected (N<sub>stars</sub>(DUO)/N<sub>stars</sub>(SAG)$``$1.25 in the overlap), increasing by this way the crowding.
#### 3.2.2 Completeness of the selection process
##### Amplitudes:
Our selection process might not be able to detect variable of low amplitude. To check the dependency of completeness on amplitude we simulated a set of 1 000 RRab light curves with the same time sampling as the real ones, the Fourier coefficients have been taken from Simon & Teays (1982). The distributions in amplitude, magnitude and period (excluding integer fractions of a day) of the simulated light curves were chosen in a way to match the actual distributions of the detected RRab, and the phasing was uniformly distributed between 0 and $`2\pi `$. This set of simulated light curves was then injected in 100 regions of $`10\mathrm{}\times 10\mathrm{}`$ (uniformly distributed over the fields) from which we took the errors to deteriorate the light curves. These light curves were then reduced in the same way as the real RRab. Fig.5a shows the completeness levels of our selection process as a function of amplitude for the two fields, averaged over 50 000 simulated light curves. The shapes of the completeness curves are nearly identical for both fields and the difference is not significant. Fig.5a shows that the detection rate stays above 95$`\%`$ for amplitude $`>`$0.8 mag, and then drops abruptly down to $``$ 20$`\%`$ at amplitude=0.5 mag. These results signify that our selection process would detect almost all the RRab with an amplitude above 0.8 mag if these were point sources.
However, the completeness levels will differ between SAG and DUO because the amplitudes of the RRab measured in each filter are different, being more important on average for SAG than for DUO . This difference occurs because the color band of SAG peaks at shorter wavelength than the color band used for DUO whereas the amplitude of RRab decreases with increasing wavelength (Smith smith (1995)). A least square fit between the amplitudes of 30 RRab in common in the overlap yielded the relation
$$\mathrm{A}_{DUO}=0.98(\pm 0.10)\mathrm{A}_{SAG}0.05(\pm 0.11)$$
(1)
where A<sub>DUO</sub> and A<sub>SAG</sub> represent respectively the amplitudes measured in DUO and in SAG . In order to construct a consistent density map, we will consider in the remainder of this paper only those RRab satisfying amplitude $`>0.60`$ mag in SAG and amplitude $`>0.54`$ mag in DUO, where 0.54 has been derived from the above relation (1). These cuts have been chosen both to ensure the largest sample as possible and to keep the completeness corrections at a manageable level. The corresponding corrections are 3.7% in SAG and 12.3% in DUO .
##### Periods:
Some RRab are missed because their periods are close to an integer fraction of a day, this causes points of the folded light curve to accumulate in a narrow phase range. The fitted Fourier series is then poorly constrained over a large fraction of the light curve and some of these stars might lie outside the ellipse of our selection process (see Fig.2). Monte-carlo simulations shows that we miss about $`30\%`$ of the RRab within the range 0.49<sup>d</sup> to 0.51<sup>d</sup> for both fields, corresponding to a total loss of $`3\%`$.
### 3.3 Homogeneity between the fields
An important point to inspect for checking the consistency between the two fields is the overlap (Fig.6). Open circles correspond to RRab detected in the SAG field with an amplitude(B<sub>i</sub>) $`>`$ 0.60 mag. while crosses represent RRab detected in the DUO field and having an amplitude(B<sub>J</sub>) $`>`$ 0.54 mag. The dashed lines indicate the limit of each plate, their inclination is due to a slight tilt between the two plates. Note that this overlap applies to the reference frames and is not necessarily constant from plate to plate. Fig.6 reveals that most of the RRab are detected independently in the SAG field and in the DUO field. However, some RRab are not detected twice and it is important to understand the reasons why these stars are missed by one of the fields:
1. RRab not detected in the SAG field (crosses)
* Eight RRab are located very close to the edge of the SAG plate. Such stars usually have fewer points in their light curves because the centre of the plate is not exactly the same at each exposure. For example the three westernmost stars in the SAG field have respectively 47, 55 and 47 points in their light curve (instead of 69 for most of the other light curves).
* Two RRab did not pass through the selection process (one because of its estimated B<sub>J</sub> amplitude and the other one because of its $`\chi _{ratio}`$).
2. RRab not detected in the DUO field (open circles)
* One RRab has not been detected probably because of its low amplitude (0.61 mag in SAG ).
* Two RRab have a period of nearly $``$ 0.50<sup>d</sup> and were detected in SAG only by chance.
* Three RRab were blended by a nearby star. As stated above, DUO is more sensitive to crowding than SAG and the relative loss of three RRab in the overlap is fully consistent with the $``$10$`\%`$ offset observed in Fig.5b in the range -10$`{}_{}{}^{}<`$b$`<`$-8.
Most of the missed RRab will therefore have no statistical incidence and should not bias the density map. The only concern is for the greater sensitivity of the DUO field to crowding. However this effect should be lowered by the crowding correction. Turning now to the western edge of DUO we re-detect 7 RRab out of the 8 detected by the MACHO team in our field (disregarding two RRab located close to the edge). This is a satisfactory result.
## 4 Structure of the Sagittarius dwarf galaxy
### 4.1 Surface density of Sgr
A surface density map is constructed from RRab with the amplitude cuts stated above. The spatial distribution of these RRab has been convolved with a Gaussian on a grid with a step of $`0.1^{}`$ and a variable filter size adapted to the local surface density $`\sigma \rho ^{1/2}`$, constrained between 0.2and 0.5. This map was then corrected for the different completeness in amplitude and crowding (see section 3.2).
The resulting map is shown on Fig 7 where the elongated shape of Sgr is clearly visible. This is the first map of Sgr in these regions, showing that Sgr extends far beyond the outer limit of the map previously published by IWGIS. One of the most striking features of this map is the slow decrease (if any ?) of the density along the main axis of Sgr for $`|`$b$`|`$9.
The main source of uncertainty in the surface density is the Poissonian noise in the star counts, which is variable over the field and tends to increase towards lower $`|`$b$`|`$. To estimate this noise we simulated 1 000 maps by injecting 1 400 stars (corresponding to the number of RRab actually used to construct the final map) onto the field with a probability density matching the surface density of the real map. These spatial distributions were then processed exactly in the same manner as the real one and a 1$`\sigma `$ “noise map” has been deduced.
This map is shown on Fig.8 where the contours are labelled in percent. The typical (relative) uncertainty is constrained between 10$`\%`$ and 15$`\%`$ over the main part of the field, but increases up to $``$40$`\%`$ towards the edges where the number of RRab drops.
### 4.2 Surface density profile of the main axis
The position angle of Sgr has been determined by fitting an exponential to the surface density along various directions. The highest scale length was reached for an angle of 108.4, which we choosed as the direction of the major axis. Fig.9 displays the density profile of Sgr along that axis. This figure is based on a map smoothed on a constant scale of 0.5. The thick line is the density after correction for completeness whereas the dotted line is the density before that correction. The shaded region represents the 1$`\sigma `$ uncertainty issued from simulated maps. A discontinuity in the slope is clearly visible at $``$6from the centre. After this point the surface density seems to be almost constant. It is however disconcerting that this discontinuity occurs near the limit between DUO and SAG and we may wonder if this is not an experimental effect. We have shown in section 3 that the crowding correction in SAG and DUO were consistent with the completeness of each field in the overlap. Furthermore the break is also perceptible in the uncorrected density so it cannot be an effect of the crowding correction. Another possibility is that our amplitude cut in DUO is to low to be consistent with SAG . This is difficult to check and we can only rely on those 30 RRab in common in the overlap, from which we derived the relation in amplitude between DUO and SAG (Eq. 1). However, if we make the assumption that the RRab population is homogeneous over the field, it is possible to search for the relation $`\mathrm{A}_{DUO}=a\mathrm{A}_{SAG}+b`$ for which the amplitude distributions are the most similar (through Kolmogorov-Smirnov test). The resulting coefficients were a=0.95 and b=-0.08 and are within the error bars stated in Eq. 1. The corresponding cuts are $`\mathrm{A}_{SAG}=0.6\mathrm{A}_{DUO}=0.49`$. These cuts would have reinforced the discontinuity, showing that our adopted amplitude cuts are not responsible for the break observed in Fig. 9. We conclude that the discontinuity in the slope of the density profile is probably real and not a consequence of the change of field.
It is also possible to derive an upper limit for the extension of Sgr along the line of sight: the distance modulus histogram can be roughly fitted by a Gaussian with a width of 0.2 mag, corresponding to a depth of $``$4.5 kpc for an assumed distance of 24 kpc.
### 4.3 Model fitting
We define the following analytical functions to fit to the density profile:
$$\rho _K=k\left\{\frac{1}{[1+(r/r_c)^2]^{1/2}}\frac{1}{[1+(r_t/r_c)^2]^{1/2}}\right\}^2$$
(2)
$$\rho _E=\rho _ee^{\frac{r}{r_e}}$$
(3)
$$\rho _G=\rho _ge^{r^2/2\sigma _g^2}$$
(4)
$$\rho _L=a_0+a_1r$$
(5)
Where $`r`$ represents the distance from the Centre of Sgr, and all other parameters are variables to be fitted. Eq. 2 refers to the empirical King model (K) with a core radius r<sub>c</sub> and tidal radius r<sub>t</sub> (King king (1962)). Eq. 3 refers to an exponential model (E) with a radius r<sub>e</sub>. Eq. 4 refers to a Gaussian (G) with a width of $`\sigma _g`$. Finally, Eq. 5 refers to a linear model (L) where the density profile is modeled by a straight line. These models have been fitted along three segments. Two of these segments are located on the main axis: one corresponding to the main axis over its entire length ($``$10), referred to as Maj10; and another one corresponding to the portion of the main axis contained within SAG ($``$6from the center), referred to as Maj6. The latter segment has been chosen in order to avoid fitting the stars past the break and also because it is more consistent since it is entirely contained within SAG . Finally the minor axis is not present within our field and instead we fitted an axis making a large angle relative to the main axis (50 ), referred to as Min . The fit of model L on Maj10has been performed by fitting the density before 6and after 6separately. In order to get uncorrelated points for the fit, we took the densities (after correction for completeness) of RRab inside boxes with a size of 0.5$`\times `$1.0located along each axis (see Fig.10). The results of the fit are shown in Table 2 and in Fig.11.
The single function model that best fits Maj10is model E ($`\chi _{fit}^2/N_{DOF}`$=3.08). However, the fit is significantly improved if we consider the model L (($`\chi _{fit}^2/N_{DOF}`$=1.77) which reproduces the break already observed in Fig. 9. A Fisher test shows that the probability for the ratio of the $`\chi _{fit}^2/N_{DOF}`$ of these two fits to be lower than the observed value by chance is only $``$13$`\%`$. Note that we were unable to fit any convergent two-component model to Maj10: this is due to the almost constant density of the external region which causes one of the component to increase as we move away from the centre in order to compensate the decrease of the other component.
Concerning the core of Sgr (Maj6), the density profile is equally well fitted by model E and L ($`\chi _{fit}^2/N_{DOF}`$2.1). The scale length derived from model E is $``$4.1$`\pm `$0.5(1.7$`\pm `$0.2 kpc). This value is slightly lower to the one derived by MOM who find an inner scale length of 4.7in the Southern part of Sgr. Model K and G also give an acceptable fit to the core of Sgr but they fail to reproduce the high density in the first bin. Furthermore, the uncertainties on the parameters of the empirical King model are quite large and the infinite tidal radius is rather unrealistic.
Finally, the best fit on Min is achieved by model G ($`\chi _{fit}^2/N_{DOF}`$=1.75), but again it fails to reproduce the high density of the first bin. The only model that reproduces the high central density is model E but the $`\chi _{fit}^2/N_{DOF}`$ of this model is worsened by the poor fit on the three last bins. However these bins contain only very few points (between 1 and 5), introducing uncertainties induced by small-numbers statistics.
## 5 Conclusion
To summarize, we presented the detection of $``$1 500 RRab stars located in the Sgr dwarf galaxy. A surface density map based on the spatial distribution of these variables unveiled the structure of this dwarf galaxy in a region that was still almost unexplored so far between b=-14and b=-4. The core of Sgr is best fitted by an exponential with a scale length of 4.1along the major axis. A cross section of this density map revealed a break in the slope occurring at $``$6from the highest density region of Sgr and an almost flat density past the break.
Although the break coincided with the change of field we have shown that this is unlikely to be an experimental effect since it is also perceptible in the uncorrected density, whereas the DUO field is intrinsically more sensitive to crowding than SAG (lower resolution, lower extinction). Also, as shown in Section 4.2, the amplitude cuts used in this study cannot be considered as responsible for the break. Finally, could this break be a consequence of an overestimation of the completeness correction in DUO relative to SAG ? Though not excluded, this would be in conflict with what is observed in the overlap where 3 RRab blended by a neighbouring star were detected in SAG and missed in DUO , a result that is quite consistent with the corrections actually applied. We argue therefore that the break is real. The significance of the break relative to an exponential with a scale length of 4.1is $``$2$`\sigma `$. MOM also observed a break in their density profile in the Southern extension of Sgr. However neither the location (20from the centre) nor the density at the break location ($`\mathrm{\Sigma }_V`$29.0 mag.arcsec<sup>-2</sup>) are consistent with our values (6and $``$26.7 mag.arcsec<sup>-2</sup>) implying that either the main body of Sgr is not symmetric or these “post-break” stars are not directly related to it.
Another striking feature revealed by the surface density profile is its flatness past the break. This feature relies on the accuracy of the completeness correction over the field, a correction that becomes quite important at low Galactic latitudes (up to 60$`\%`$). Yet, the difficulty of modeling point spread functions on photographic plates (due to non-linear response of the emulsion) and potential systematic errors caused by differential sensitivity over the plate makes the crowding correction rather uncertain. Therefore, although our completeness corrections are fairly consistent within the overlap, we cannot exclude that the flatness of the density profile in the outer regions is a consequence of an overcorrection. Wide-field high resolution imaging would be necessary in these extremely crowded regions (up to $``$10<sup>6</sup> stars per square degree at our magnitude limit) to confirm or to rule out this issue. Nevertheless, even if we consider that our completeness corrections are overestimated by a factor of 2 (a quite conservative estimate), it remains that the density profile decreases slowly in the outer regions and Sgr may well be extending even further out towards (beyond ?) the Galactic plane.
Johnston et al. (johnston99 (1999)) recently modeled the Sgr stream as a superposition of a main body and tidal streams of stars stripped on previous peri-centric passages. This scenario has been worked out to explain both the break observed by MOM and the possible detection of stars in the outer region of Sgr with different radial velocities relative to those of the main body (Majewski et al. maj (1999)). Similarly, spectroscopic observations on our RR Lyrae catalogue could allow to determine the nature of the stars in the outer region: if these stars are linked to the main body of Sgr, then they should share almost the same radial velocities as the main body (apart of a gradient along the main axis due to the rapidly varying Galactic potential). On the other hand, if the break we observe corresponds to a transition between the main body and an unbound tidal stream from a previous orbit, it is likely that the two objects will have different radial velocities. This new catalogue of RR Lyrae is an interesting opportunity to study further a region of Sgr that has been poorly investigated so far.
###### Acknowledgements.
We thank René Chesnel for scanning most of the plates used in this paper. We would also like to thank Rodrigo Ibata and Stéphane Léon for interesting discussions. Finally, we thank the anonymous referee for valuable comments which helped to improve this paper.
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# Almost ring theory
## 1. Introduction
Broadly speaking, the aim of this work is to describe “how to do ring theory” within monoidal categories that arise as localisations of categories of modules over certain rings. A reader looking for forerunners of our themes would be drawn inevitably to Gabriel’s “Des categories abeliennes” , and might even conclude that Gabriel’s memoir must have been the main instigation for the present article. In truth, the initial motivation is to be found elsewhere, namely in the want of adequately documented foundations for the method of almost étale extensions that underpins Faltings’ approach to $`p`$-adic Hodge theory as presented in . However, as is often the case with healthy offsprings, our subject matter has eventually resolved to venture beyond its original boundaries and pursue an autonomous existence.
In any case, we are glad to report that our paper remains true to its first vocation, which is to serve as a comprehensive reference, paving the way to deeper aspects of almost étale theory, especially to the difficult purity theorem of . The notions of almost unramified and almost étale morphism are defined and their main properties are established, including the analogues of the classical lifting theorems over nilpotent extensions, and invariance under Frobenius. Also, to any almost finitely presented almost flat morphism we attach an almost trace form, and we characterize almost étale morphisms in terms of this form. Finally, we study some cases of non-flat descent for almost étale maps.
Actually, our terminology is slightly different, in that we replace usual modules and algebras by their “almost” counterparts, which live in the category of almost modules, a localisation of the category of modules. So for instance, instead of almost étale morphisms of algebras, we have étale morphisms of almost algebras.
The categories of almost modules (or almost algebras) and of usual modules (resp. algebras) are linked in manifold ways. First of all we have of course the localisation functor. Then, as it had already been observed by Gabriel, there is a right adjoint to localisation. Furthermore, we show that there is a left adjoint as well, that, to our knowledge, has not been exploited before, in spite of its several useful qualities which will establish it quickly as one of our main tools. The ensemble of localisation and right, left adjoints exhibits some remarkable exactness properties, that are typically associated to open imbeddings of topoi, all of which seems to suggest the existence of some deeper geometric structure, still to be unearthed. We may have encountered here an instance of a general principle, apparently evoked first by Deligne, according to which one should try to do algebraic geometry on arbitrary abelian tensor categories (notice though, that our categories are more general than the tannakian categories of ).
A large part of the paper is devoted to the construction and study of the almost version of Illusie’s cotangent complex, on which we base our deformation theory for almost algebras. Faltings’ original method was based on Hochschild cohomology rather than the cotangent complex. While Faltings’ approach has the advantage of being more explicit and elementary, it also has the drawback of involving a very large number of long and tedious manipulations with cocycles, and requires a painstaking tracking of the “epsilon book-keeping”. The method presented here avoids (or at least removes from view!) these problems, and also leads to more general results (especially, we can drop all finiteness assumptions from the statements of the lifting theorems).
Though we have strived throughout for the widest generality, in a few places one could have gone even further : for instance it would have been possible to globalise all definitions and most results to arbitrary schemes. However, the extension to schemes is completely straightforward, and in practice seems to be scarcely useful. Similarly, there is currently not much incentive to study a notion of “almost smooth morphism”.
## 2. Homological theory
### 2.1. Some ring-theoretic preliminaries
Unless otherwise stated, every ring is commutative with unit. Our basic setup consists of a fixed base ring $`V`$ containing an ideal $`𝔪`$ such that $`𝔪^2=𝔪`$. Starting from section 2.4, we will also assume that $`\stackrel{~}{𝔪}=𝔪_V𝔪`$ is a flat $`V`$-module.
###### Example 2.1.1.
i) The main example is given by a non-discrete valuation ring $`V`$ with valuation $`\nu :V\{0\}\mathrm{\Gamma }`$ of rank one (where $`\mathrm{\Gamma }`$ is the totally ordered abelian group of values of $`\nu `$). Then we can take $`𝔪=\{0\}\{xV\{0\}|\nu (x)>\nu (1)\}`$.
ii) Suppose that $`𝔪=V`$. This is the “classical limit”. In this case almost ring theory reduces to usual ring theory. Thus, all the discussion that follows specialises to, and sometimes gives alternative proofs for, statements about rings and their modules.
We define a uniform structure on the set I of ideals of $`V`$ as follows. For every finitely generated ideal $`𝔪_0𝔪`$ the subset of $`\text{I}\times \text{I}`$ given by $`\{(I,J)|𝔪_0JI\text{and}𝔪_0IJ\}`$ is an entourage for the uniform structure, and the subsets of this kind form a fundamental system of entourages (cp. Ch.II §1). The uniform structure induces a topology on I and moreover the notion of convergent (resp. Cauchy) sequence of ideals is well defined. We will also need a notion of “Cauchy product” : let $`_{n=0}^{\mathrm{}}I_n`$ be a formal infinite product of ideals. We say that the formal product satisfies the Cauchy condition (or briefly : is a Cauchy product) if, for every neighborhood U of $`V\text{I}`$ there exists $`n_00`$ such that $`_{m=n}^{n+p}I_m\text{U}`$ for all $`nn_0`$ and all $`p0`$.
###### Remark 2.1.2.
Suppose that $`J_0J_1\mathrm{}`$ is an increasing infinite sequence of ideals of $`I`$ such that $`\underset{k\mathrm{}}{\text{lim}}J_k=V`$ (convergence for the above uniform structure on I). Then one checks easily that $`_{k=0}^{\mathrm{}}𝔪J_k=𝔪`$.
Let $`M`$ be a given $`V`$-module. We say that $`M`$ is almost zero if $`𝔪M=0`$. A map $`\varphi `$ of $`V`$-modules is an almost isomorphism if both $`\text{Ker}(\varphi )`$ and $`\text{Coker}(\varphi )`$ are almost zero $`V`$-modules.
###### Remark 2.1.3.
(i) It is easy to check that a $`V`$-module $`M`$ is almost zero if and only if $`𝔪_VM=0`$. Similarly, a map $`MN`$ of $`V`$-modules is an almost isomorphism if and only if the induced map $`\stackrel{~}{𝔪}_VM\stackrel{~}{𝔪}_VN`$ is an isomorphism. Notice also that, if $`𝔪`$ is flat, then $`𝔪\stackrel{~}{𝔪}`$.
(ii) Let $`VW`$ be a ring homomorphism. For a $`V`$-module $`M`$ set $`M_W=W_VM`$. We have an exact sequence
(2.1.4)
$$0K𝔪_W𝔪W0$$
where $`K=\text{Tor}_1^V(V/𝔪,W)`$ is an almost zero $`W`$-module. By (i) it follows that $`𝔪_VK(𝔪W)_WK0`$. Then, applying $`𝔪_W_W`$ and $`_W(𝔪W)`$ to (2.1.4) we derive
$$𝔪_W_W𝔪_W𝔪_W_W(𝔪W)(𝔪W)_W(𝔪W)$$
i.e. $`\stackrel{~}{𝔪}_W(𝔪W)\stackrel{~}{}`$. In particular, if $`\stackrel{~}{𝔪}`$ is a flat $`V`$-module, then $`\stackrel{~}{𝔪}_W`$ is a flat $`W`$-module. This means that our basic assumptions on the pair $`(V,𝔪)`$ are stable under arbitrary base extension. Notice that the flatness of $`𝔪`$ does not imply the flatness of $`𝔪W`$. This partly explains why we insist that $`\stackrel{~}{𝔪}`$, rather than $`𝔪`$, be flat.
Before moving on, we want to analyze in some detail how our basic assumptions relate to certain other natural conditions that can be postulated on the pair $`(V,𝔪)`$. Indeed, let us consider the following two hypotheses :
(A) $`𝔪=𝔪^2`$ and $`𝔪`$ is a filtered colimit of principal ideals.
(B) $`𝔪=𝔪^2`$ and, for all integers $`k>1`$, the $`k`$-th powers of elements of $`𝔪`$ generate $`𝔪`$.
Clearly (A) implies (B). Less obvious is the following result.
###### Proposition 2.1.5.
(i) (A) implies that $`\stackrel{~}{𝔪}`$ is flat.
(ii) If $`\stackrel{~}{𝔪}`$ is flat then (B) holds.
###### Proof.
Suppose that (A) holds, so that $`𝔪=\underset{\alpha I}{\text{colim}}Vx_\alpha `$, where $`I`$ is a directed set parametrizing elements $`x_\alpha 𝔪`$ (and $`\alpha \beta Vx_\alpha Vx_\beta `$). For any $`\alpha I`$ we have natural isomorphisms
(2.1.6)
$$Vx_\alpha V/\text{Ann}_V(x_\alpha )(Vx_\alpha )_V(Vx_\alpha ).$$
For $`\alpha \beta `$, let $`j_{\alpha \beta }:Vx_\alpha Vx_\beta `$ be the imbedding; we have a commutative diagram
where $`zV`$ is such that $`x_\alpha =zx_\beta `$, $`\mu _{z^2}`$ is multiplication by $`z^2`$ and $`\pi _\alpha `$ is the projection induced by (2.1.6) (and similarly for $`\pi _\beta `$). Since $`𝔪=𝔪^2`$, for all $`\alpha I`$ we can find $`\beta `$ such that $`x_\alpha `$ is a multiple of $`x_\beta ^2`$. Say $`x_\alpha =yx_\beta ^2`$; then we can take $`z=yx_\beta `$, so $`z^2`$ is a multiple of $`x_\alpha `$ and in the above diagram $`\text{Ker}(\pi _\alpha )\text{Ker}(\mu _{z^2})`$. Hence one can define a map $`\lambda _{\alpha \beta }:(Vx_\alpha )_V(Vx_\alpha )V`$ such that $`\pi _\beta \lambda _{\alpha \beta }=j_{\alpha \beta }j_{\alpha \beta }`$ and $`\lambda _{\alpha \beta }\pi _\alpha =\mu _{z^2}`$. It now follows that for every $`V`$-module $`N`$, the induced morphism $`\text{Tor}_1^V(N,(Vx_\alpha )_V(Vx_\alpha ))\text{Tor}_1^V(N,(Vx_\beta )_V(Vx_\beta ))`$ is the zero map. Taking the colimit we derive that $`\stackrel{~}{𝔪}`$ is flat. This shows (i). In order to show (ii) we consider, for any prime number $`p`$, the following condition
($`_p`$) $`𝔪/p𝔪`$ is generated (as a $`V`$-module) by the $`p`$-th powers of its elements.
Clearly (B) implies ($`_p`$) for all $`p`$. In fact we have :
###### Claim 2.1.7.
(B) holds if and only if ($`_p`$) holds for every prime $`p`$.
Proof of the claim: Suppose that ($`_p`$) holds for every prime $`p`$. The polarization identity
$$k!x_1x_2\mathrm{}x_k=\underset{I\{1,2,\mathrm{},k\}}{}(1)^{k|I|}\left(\underset{iI}{}x_i\right)^k$$
shows that if $`N=_{x𝔪}Vx^k`$ then $`k!𝔪N`$. To prove that $`N=𝔪`$ it then suffices to show that for every prime $`p`$ dividing $`k!`$ we have $`𝔪=p𝔪+N`$. Let $`\varphi :V/pVV/pV`$ be the Frobenius ($`xx^p`$); we can denote by $`(V/pV)^\varphi `$ the ring $`V/pV`$ seen as a $`V/pV`$-algebra via the homomorphism $`\varphi `$. Also set $`\varphi ^{}M=M_{V/pV}(V/pV)^\varphi `$ for a $`V/pV`$-module $`M`$. Then the map $`\varphi ^{}(𝔪/p𝔪)(𝔪/p𝔪)`$ (defined by raising to $`p`$-th power) is surjective by ($`_p`$). Hence so is $`(\varphi ^r)^{}(𝔪/p𝔪)(𝔪/p𝔪)`$ for every $`r>0`$, which says that $`𝔪=p𝔪+N`$ when $`k=p^r`$, hence for every $`k`$.
Next recall (see Exp. XVII 5.5.2) that, if $`M`$ is a $`V`$-module, the module of symmetric tensors $`\text{TS}^k(M)`$ is defined as $`(_V^kM)^{S_k}`$, the invariants under the natural action of the symmetric group $`S_k`$ on $`_V^kM`$. We have a natural map $`\mathrm{\Gamma }^k(M)\text{TS}^k(M)`$ that is an isomorphism when $`M`$ is flat (see loc. cit. 5.5.2.5; here $`\mathrm{\Gamma }^k`$ denotes the $`k`$-th graded piece of the divided power algebra).
###### Claim 2.1.8.
The group $`S_k`$ acts trivially on $`_V^k𝔪`$ and the map $`\stackrel{~}{𝔪}_V𝔪\stackrel{~}{𝔪}`$ ($`xyzxyz`$) is an isomorphism.
Proof of the claim: The first statement is reduced to the case of transpositions and to $`k=2`$. There we can compute : $`xyz=xyz=yxz=yzx`$. For the second statement note that the imbedding $`𝔪V`$ is an almost isomorphism, and apply remark 2.1.3(i).
Suppose now that $`\stackrel{~}{𝔪}`$ is flat and pick a prime $`p`$. Then $`S_p`$ acts trivially on $`_V^p\stackrel{~}{𝔪}`$. Hence
(2.1.9)
$$\mathrm{\Gamma }^p(\stackrel{~}{𝔪})_V^p\stackrel{~}{𝔪}\stackrel{~}{𝔪}.$$
But $`\mathrm{\Gamma }^p(\stackrel{~}{𝔪})`$ is spanned as a $`V`$-module by the products $`\gamma _{i_1}(x_1)\mathrm{}\gamma _{i_k}(x_k)`$ (where $`x_i\stackrel{~}{𝔪}`$ and $`_ji_j=p`$). Under the isomorphism (2.1.9) these elements map to $`\left(\genfrac{}{}{0pt}{}{p}{i_1,\mathrm{},i_k}\right)x_1^{i_1}\mathrm{}x_k^{i_k}`$; but such an element vanishes in $`\stackrel{~}{𝔪}/p\stackrel{~}{𝔪}`$ unless $`i_k=p`$ for some $`k`$. Therefore $`\stackrel{~}{𝔪}/p\stackrel{~}{𝔪}`$ is generated by $`p`$-th powers, so the same is true for $`𝔪/p𝔪`$, and by the above, (B) holds, which shows (ii). ∎
###### Proposition 2.1.10.
Suppose that $`𝔪`$ is countably generated as a $`V`$-module. Then we have :
i) $`\stackrel{~}{𝔪}`$ is countably presented as a $`V`$-module;
ii) if $`\stackrel{~}{𝔪}`$ is a flat $`V`$-module, then it is of homological dimension $`1`$.
###### Proof.
Let $`(\epsilon _i)_{iI}`$ be a countable generating family of $`𝔪`$. Then $`\epsilon _i\epsilon _j`$ generate $`\stackrel{~}{𝔪}`$ and $`\epsilon _i\epsilon _j(\epsilon _k\epsilon _l)=\epsilon _k\epsilon _l(\epsilon _i\epsilon _j)`$ for all $`i,j,k,lI`$. For every $`iI`$, we can write $`\epsilon _i=_jx_{ij}\epsilon _j`$, for certain $`x_{ij}𝔪`$. Let $`F`$ be the $`V`$-module defined by generators $`(e_{ij})_{i,jI}`$, subject to the relations:
$$\epsilon _i\epsilon _je_{kl}=\epsilon _k\epsilon _le_{ij}e_{ik}=\underset{j}{}x_{ij}e_{jk}\text{for all }i,j,k,lI\text{.}$$
We get an epimorphism $`\pi :F\stackrel{~}{𝔪}`$ by $`e_{ij}\epsilon _i\epsilon _j`$. The relations imply that, if $`x=_{k,l}y_{kl}\text{Ker}(\pi )`$, then $`\epsilon _i\epsilon _jx=0`$, so $`𝔪\text{Ker}(\pi )=0`$. Whence $`𝔪_V\text{Ker}(\pi )=0`$ and $`\text{1}_𝔪_V\pi `$ is an isomorphism. We consider the diagram
where $`\varphi `$ and $`\psi `$ are induced by scalar multiplication. We already know that $`\psi `$ is an isomorphism, and since $`F=𝔪F`$, we see that $`\varphi `$ is an epimorphism, so $`\pi `$ is an isomorphism, which shows (i). Now (ii) follows from (i), since it is well-known that a flat countably presented module is of homological dimension $`1`$ (see (Ch.I, Th.3.2) and the discussion in pp.49-50). ∎
### 2.2. Almost categories
If C is a category, and $`X,Y`$ two objects of C, we will usually denote by $`\text{Hom}_\text{C}(X,Y)`$ the set of morphisms in C from $`X`$ to $`Y`$ and by $`\text{1}_X`$ the identity morphism of $`X`$. Moreover we denote by $`\text{C}^o`$ the opposite category of C and by $`s.\text{C}`$ the category of simplicial objects over C, that is, functors $`\mathrm{\Delta }^o\text{C}`$, where $`\mathrm{\Delta }`$ is the category whose objects are the ordered sets $`[n]=\{0,\mathrm{},n\}`$ for each integer $`n0`$ and where a morphism $`\varphi :[p][q]`$ is a non-decreasing map. A morphism $`f:XY`$ in $`s.\text{C}`$ is a sequence of morphisms $`f_{[n]}:X[n]Y[n]`$, $`n0`$ such that the obvious diagrams commute. We can imbed C in $`s.\text{C}`$ by sending each object $`X`$ to the “constant” object $`s.X`$ such that $`s.X[n]=X`$ for all $`n0`$ and $`s.X[\varphi ]=\text{1}_X`$ for all morphisms $`\varphi `$ in $`\mathrm{\Delta }`$.
If C is an abelian category, $`\text{D}(\text{C})`$ will denote the derived category of C. As usual we have also the full subcategories $`\text{D}^+(\text{C}),\text{D}^{}(\text{C})`$ of complexes of objects of C that are exact for sufficiently large negative (resp. positive) degree. If $`R`$ is a ring, the category of $`R`$-modules (resp. $`R`$-algebras) will be denoted by $`R\text{-}\mathrm{𝐌𝐨𝐝}`$ (resp. $`R\text{-}\mathrm{𝐀𝐥𝐠}`$). Most of the times we will write $`\text{Hom}_R(M,N)`$ instead of $`\text{Hom}_{R\text{-}\mathrm{𝐌𝐨𝐝}}(M,N)`$.
We denote by $`\mathrm{𝐒𝐞𝐭}`$ the category of sets. The symbol $``$ denotes the set of non-negative integers; in particular $`0`$.
The full subcategory $`\mathrm{\Sigma }`$ of $`V\text{-}\mathrm{𝐌𝐨𝐝}`$ consisting of all $`V`$-modules that are almost isomorphic to $`0`$ is clearly a Serre subcategory and hence we can form the quotient category $`V\text{-}\mathrm{𝐌𝐨𝐝}/\mathrm{\Sigma }`$. There is a localization functor
$$V\text{-}\mathrm{𝐌𝐨𝐝}V\text{-}\mathrm{𝐌𝐨𝐝}/\mathrm{\Sigma }MM^a$$
that takes a $`V`$-module $`M`$ to the same module, seen as an object of $`V\text{-}\mathrm{𝐌𝐨𝐝}/\mathrm{\Sigma }`$. In particular, we have the object $`V^a`$ associated to $`V`$; it seems therefore natural to use the notation $`V^a\text{-}\mathrm{𝐌𝐨𝐝}`$ for the category $`V\text{-}\mathrm{𝐌𝐨𝐝}/\mathrm{\Sigma }`$, and an object of $`V^a\text{-}\mathrm{𝐌𝐨𝐝}`$ will be indifferently referred to as “a $`V^a`$-module” or “an almost $`V`$-module”. In case we need to stress the dependance on the ideal $`𝔪`$, we can write $`(V,𝔪)^a\text{-}\mathrm{𝐌𝐨𝐝}`$.
Since the almost isomorphisms form a multiplicative system (see e.g. Exerc.10.3.2), it is possible to describe the morphisms in $`V^a\text{-}\mathrm{𝐌𝐨𝐝}`$ via a calculus of fractions, as follows. Let $`V\text{-}\mathrm{𝐚𝐥}.\mathrm{𝐈𝐬𝐨}`$ be the category that has the same objects as $`V\text{-}\mathrm{𝐌𝐨𝐝}`$, but such that $`\text{Hom}_{V\text{-}\mathrm{𝐚𝐥}.\mathrm{𝐈𝐬𝐨}}(M,N)`$ consists of all almost isomorphisms $`MN`$. If $`M`$ is any object of $`V\text{-}\mathrm{𝐚𝐥}.\mathrm{𝐈𝐬𝐨}`$ we write $`(V\text{-}\mathrm{𝐚𝐥}.\mathrm{𝐈𝐬𝐨}/M)`$ for the category of objects of $`V\text{-}\mathrm{𝐚𝐥}.\mathrm{𝐈𝐬𝐨}`$ over $`M`$ (i.e. morphisms $`\varphi :XM`$). If $`\varphi _i:X_iM`$ $`(i=1,2)`$ are two objects of $`(V\text{-}\mathrm{𝐚𝐥}.\mathrm{𝐈𝐬𝐨}/M)`$ then $`\text{Hom}_{(V\text{-}\mathrm{𝐚𝐥}.\mathrm{𝐈𝐬𝐨}/M)}(\varphi _1,\varphi _2)`$ consists of all morphisms $`\psi :X_1X_2`$ in $`V\text{-}\mathrm{𝐚𝐥}.\mathrm{𝐈𝐬𝐨}`$ such that $`\varphi _1=\varphi _2\psi `$. For any two $`V`$-modules $`M,N`$ we define a functor $`\text{F}_N:(V\text{-}\mathrm{𝐚𝐥}.\mathrm{𝐈𝐬𝐨}/M)^oV\text{-}\mathrm{𝐌𝐨𝐝}`$ by associating to an object $`\varphi :PM`$ the $`V`$-module $`\text{Hom}_V(P,N)`$ and to a morphism $`\alpha :PQ`$ the map $`\text{Hom}_V(Q,N)\text{Hom}_V(P,N):\beta \beta \alpha `$. Then we have
(2.2.1)
$$\text{Hom}_{V^a\text{-}\mathrm{𝐌𝐨𝐝}}(M^a,N^a)=\underset{(V\text{-}\mathrm{𝐚𝐥}.\mathrm{𝐈𝐬𝐨}/M)^o}{\text{colim}}\text{F}_N.$$
However, formula (2.2.1) can be simplified considerably, by remarking that, for any $`V`$-module $`M`$, the natural morphism $`\stackrel{~}{𝔪}_VMM`$ is an initial object of $`(V\text{-}\mathrm{𝐚𝐥}.\mathrm{𝐈𝐬𝐨}/M)`$. Indeed, let $`\varphi :NM`$ be an almost isomorphism; the diagram
(cp. remark 2.1.3(i)) allows one to define a morphism $`\psi :\stackrel{~}{𝔪}_VMN`$ over $`M`$. We need to show that $`\psi `$ is unique. But if $`\psi _1,\psi _2:\stackrel{~}{𝔪}_VMN`$ are two maps over $`M`$, then $`\text{Im}(\psi _1\psi _2)\text{Ker}(\varphi )`$ is almost zero, hence $`\text{Im}(\psi _1\psi _2)=0`$, since $`\stackrel{~}{𝔪}_VM=𝔪(\stackrel{~}{𝔪}_VM)`$. Consequently, (2.2.1) boils down to
(2.2.2)
$$\text{Hom}_{V^a\text{-}\mathrm{𝐌𝐨𝐝}}(M^a,N^a)=\text{Hom}_V(\stackrel{~}{𝔪}_VM,N).$$
In particular $`\text{Hom}_{V^a\text{-}\mathrm{𝐌𝐨𝐝}}(M,N)`$ has a natural structure of $`V`$-module for any two $`V^a`$-modules $`M,N`$, i.e. $`\text{Hom}_{V^a\text{-}\mathrm{𝐌𝐨𝐝}}(,)`$ is a bifunctor that takes values in the category $`V\text{-}\mathrm{𝐌𝐨𝐝}`$.
One checks easily (for instance using (2.2.2)) that the usual tensor product induces a bifunctor $`_V`$ on almost $`V`$-modules, which, in the jargon of makes of $`V^a\text{-}\mathrm{𝐌𝐨𝐝}`$ an abelian tensor category. Then an almost $`V`$-algebra is just a commutative unitary monoid in the tensor category $`V^a\text{-}\mathrm{𝐌𝐨𝐝}`$. Let us recall what this means. Quite generally, let $`(\text{C},,U)`$ be any abelian tensor category, so that $`:\text{C}\times \text{C}\text{C}`$ is a biadditive functor, $`U`$ is the identity object of C (see p.105) and for any two objects $`M`$ and $`N`$ in C we have a “commutativity constraint” (i.e. a functorial isomorphism $`\eta _{M|N}:MNNM`$ that “switches the two factors”) and a functorial isomorphism $`\nu _M:UMM`$. Then a C-monoid $`A`$ is an object of C endowed with a morphism $`\mu _A:AAA`$ (the “multiplication” of $`A`$) satisfying the associativity condition
$$\mu _A(\text{1}_A\mu _A)=\mu _A(\mu _A\text{1}_A).$$
We say that $`A`$ is unitary if additionally $`A`$ is endowed with a “unit morphism” $`\underset{¯}{1}_A:UA`$ satisfying the (left and right) unit property :
$$\mu _A(\underset{¯}{1}_A\text{1}_A)=\nu _A\mu _A(\underset{¯}{1}_A\text{1}_A)\eta _{A|U}=\mu _A(\text{1}_A\underset{¯}{1}_A).$$
Finally $`A`$ is commutative if $`\mu _A=\mu _A\eta _{A|A}`$ (to be rigorous, in all of the above one should indicate the associativity constraints, which we have omitted : see ). A commutative unitary monoid will also be simply called an algebra. With the morphisms defined in the obvious way, the C-monoids form a category; furthermore, given a C-monoid $`A`$, a left $`A`$-module is an object $`M`$ of C endowed with a morphism $`\sigma _M:AMM`$ such that $`\sigma _M(\text{1}_A\sigma _M)=\sigma _M(\mu _A\text{1}_M)`$. Similarly one defines right $`A`$-modules and $`A`$-bimodules. In the case of bimodules we have left and right morphisms $`\sigma _{M,l}:AMM`$, $`\sigma _{M,r}:MAM`$ and one imposes that they “commute”, i.e. that
$$\sigma _{M,r}(\sigma _{M,l}\text{1}_A)=\sigma _{M,l}(\text{1}_A\sigma _{M,r}).$$
Clearly the (left resp. right) $`A`$-modules (and the $`A`$-bimodules) form an additive category with $`A`$-linear morphisms defined as one expects. One defines the notion of a submodule as an equivalence class of monomorphisms $`NM`$ such that the composition $`ANAMM`$ factors through $`N`$. Now, if $`f:MN`$ is a morphism of left $`A`$-modules, then $`\text{Ker}(f)`$ exists in the underlying abelian category C and one checks easily that it has a unique structure of left $`A`$-module which makes it a submodule of $`M`$. If moreover $``$ is right exact when either argument is fixed, then also $`\text{Coker}(f)`$ has a unique $`A`$-module structure for which $`N\text{Coker}(f)`$ is $`A`$-linear. In this case the category of left $`A`$-modules is abelian. Similarly, if $`A`$ is a unitary C-monoid, then one defines the notion of unitary left $`A`$-module by requiring that $`\sigma _M(\underset{¯}{1}_A\text{1}_M)=\nu _M`$ and these form an abelian category when $``$ is right exact.
Specialising to our case we obtain the category $`V^a\text{-}\mathrm{𝐀𝐥𝐠}`$ of almost $`V`$-algebras and, for every almost $`V`$-algebra $`A`$, the category $`A\text{-}\mathrm{𝐌𝐨𝐝}`$ of unitary left $`A`$-modules. Clearly the localization functor restricts to a functor $`V\text{-}\mathrm{𝐀𝐥𝐠}V^a\text{-}\mathrm{𝐀𝐥𝐠}`$ and for any $`V`$-algebra $`R`$ we have a localization functor $`R\text{-}\mathrm{𝐌𝐨𝐝}R^a\text{-}\mathrm{𝐌𝐨𝐝}`$.
Next, if $`A`$ is an almost $`V`$-algebra, we can define the category $`A\text{-}\mathrm{𝐀𝐥𝐠}`$ of $`A`$-algebras. It consists of all the morphisms $`AB`$ of almost $`V`$-algebras.
Let again $`(\text{C},,U)`$ be any abelian tensor category. By p.119, the endomorphism ring $`\text{End}_\text{C}(U)`$ of $`U`$ is commutative. For any object $`M`$ of C, denote $`M_{}=\text{Hom}_\text{C}(U,M)`$; then $`MM_{}`$ defines a functor $`\text{C}\text{End}_\text{C}(U)\text{-}\mathrm{𝐌𝐨𝐝}`$. Moreover, if $`A`$ is a C-monoid, $`A_{}`$ is an associative $`\text{End}_\text{C}(U)`$-algebra, with multiplication given as follows. For $`a,bA_{}`$ let $`ab=\mu _A(ab)\nu _U^1`$. Similarly, if $`M`$ is an $`A`$-module, $`M_{}`$ is an $`A_{}`$-module in a natural way, and in this way we obtain a functor from $`A`$-modules and $`A`$-linear morphisms to $`A_{}`$-modules and $`A_{}`$-linear maps. Using (Prop. 1.3), one can also check that $`\text{End}_\text{C}(U)=U_{}`$ as $`\text{End}_\text{C}(U)`$-algebras, where $`U`$ is viewed as a C-monoid using $`\nu _U`$.
All this applies especially to our categories of almost modules and almost algebras : in this case we call $`MM_{}`$ the functor of almost elements. So, if $`M`$ is an almost module, an almost element of $`M`$ is just an honest element of $`M_{}`$. Using (2.2.2) one can show easily that for every $`V`$-module $`M`$ the natural map $`M(M^a)_{}`$ is an almost isomorphism.
Let $`A`$ be an almost $`V`$-algebra. For any two $`A`$-modules $`M,N`$, the set $`\text{Hom}_{A\text{-}\mathrm{𝐌𝐨𝐝}}(M,N)`$ has a natural structure of $`A_{}`$-module and we obtain an internal Hom functor by letting
$$\text{alHom}_A(M,N)=\text{Hom}_{A\text{-}\mathrm{𝐌𝐨𝐝}}(M,N)^a.$$
This is the functor of almost homomorphisms from $`M`$ to $`N`$.
For any $`A`$-module $`M`$ we have also a functor of tensor product $`M_A`$ on $`A`$-modules which, in view of the following proposition 2.2.4 can be shown to be a left adjoint to the functor $`\text{alHom}_A(M,)`$. It can be defined as $`M_AN=(M_{}_A_{}N_{})^a`$ but an appropriate almost version of the usual construction would also work.
With this tensor product, $`A\text{-}\mathrm{𝐌𝐨𝐝}`$ is an abelian tensor category as well, and $`A\text{-}\mathrm{𝐀𝐥𝐠}`$ could also be described as the category of ($`A\text{-}\mathrm{𝐌𝐨𝐝}`$)-algebras. Under this equivalence, a morphism $`\varphi :AB`$ of almost $`V`$-algebras becomes the unit morphism $`\underset{¯}{1}_B:AB`$ of the corresponding monoid. We will sometimes drop the subscript and write simply $`\underset{¯}{1}`$.
###### Remark 2.2.3.
Let $`VW`$ be a map of base rings, $`W`$ taken with the extended ideal $`𝔪W`$. Then $`W^a`$ is an almost $`V`$-algebra so we have defined the category $`W^a\text{-}\mathrm{𝐌𝐨𝐝}`$ using base ring $`V`$ and the category $`(W,𝔪W)^a\text{-}\mathrm{𝐌𝐨𝐝}`$ using base $`W`$. One shows easily that they are equivalent: we have an obvious functor $`(W,𝔪W)^a\text{-}\mathrm{𝐌𝐨𝐝}W^a\text{-}\mathrm{𝐌𝐨𝐝}`$ and an essential inverse is provided by $`MM_{}`$. Similar base comparison statements hold for the categories of almost algebras.
###### Proposition 2.2.4.
i) There is a natural isomorphism $`AA_{}^a`$ of almost $`V`$-algebras.
ii) Let $`R`$ be any $`V`$-algebra. Then the functor $`MM_{}`$ from $`R^a\text{-}\mathrm{𝐌𝐨𝐝}`$ to $`R\text{-}\mathrm{𝐌𝐨𝐝}`$ (resp. from $`R^a\text{-}\mathrm{𝐀𝐥𝐠}`$ to $`R\text{-}\mathrm{𝐀𝐥𝐠}`$) is right adjoint to the localization functor $`R\text{-}\mathrm{𝐌𝐨𝐝}R^a\text{-}\mathrm{𝐌𝐨𝐝}`$ (resp. $`R\text{-}\mathrm{𝐀𝐥𝐠}R^a\text{-}\mathrm{𝐀𝐥𝐠}`$).
iii) The counit of the adjunction $`M_{}^aM`$ is a natural isomorphism from the composition of the two functors to the identity functor $`\text{1}_{A\text{-}\mathrm{𝐌𝐨𝐝}}`$ (resp. $`\text{1}_{A\text{-}\mathrm{𝐀𝐥𝐠}}`$).
###### Proof.
(i) has already been remarked. We show (ii). In light of remark 2.2.3 (applied with $`W=R`$) we can assume that $`V=R`$. Let $`M`$ be a $`V`$-module and $`N`$ an almost $`V`$-module; we have natural bijections
$$\begin{array}{cc}\hfill \text{Hom}_{V^a\text{-}\mathrm{𝐌𝐨𝐝}}(M^a,N)& \text{Hom}_{V^a\text{-}\mathrm{𝐌𝐨𝐝}}(M^a,(N_{})^a)\text{Hom}_V(\stackrel{~}{𝔪}_VM,N_{})\hfill \\ \hfill & \text{Hom}_V(M,\text{Hom}_V(\stackrel{~}{𝔪},N_{}))\text{Hom}_V(M,\text{Hom}_{V^a\text{-}\mathrm{𝐌𝐨𝐝}}(V,(N_{})^a))\hfill \\ \hfill & \text{Hom}_V(M,N_{})\hfill \end{array}$$
which proves (ii). Now (iii) follows by inspecting the proof of (ii), or by (ch.III Prop.3). ∎
###### Remark 2.2.5.
The existence of the right adjoint follows also directly from (chap.III §3 Cor.1 or chap.V §2).
###### Corollary 2.2.6.
The functor $`MM_{}`$ from $`R^a\text{-}\mathrm{𝐌𝐨𝐝}`$ to $`R\text{-}\mathrm{𝐌𝐨𝐝}`$ sends injectives to injectives and injective envelopes to injective envelopes.
###### Proof.
The functor $`MM_{}`$ is right adjoint to an exact functor, hence it preserves injectives. Now, let $`J`$ be an injective envelope of $`M`$; to show that $`J_{}`$ is an injective envelope of $`M_{}`$, it suffices to show that $`J_{}`$ is an essential extension of $`M_{}`$. However, if $`NJ_{}`$ and $`NM_{}=0`$, then $`N^aM=0`$, hence $`𝔪N=0`$, but $`J_{}`$ does not contain $`𝔪`$-torsion, thus $`N=0`$. ∎
###### Corollary 2.2.7.
The categories $`A\text{-}\mathrm{𝐌𝐨𝐝}`$ and $`A\text{-}\mathrm{𝐀𝐥𝐠}`$ are both complete and cocomplete.
###### Proof.
We recall that the categories $`A_{}\text{-}\mathrm{𝐌𝐨𝐝}`$ and $`A_{}\text{-}\mathrm{𝐀𝐥𝐠}`$ are both complete and cocomplete. Now let $`I`$ be any small indexing category and $`M:IA\text{-}\mathrm{𝐌𝐨𝐝}`$ be any functor. Denote by $`M_{}:IA_{}\text{-}\mathrm{𝐌𝐨𝐝}`$ the composed functor $`iM(i)_{}`$. We claim that $`\underset{I}{\text{colim}}M=(\underset{I}{\text{colim}}M_{})^a`$. The proof is an easy application of proposition 2.2.4(iii). A similar argument also works for limits and for the category $`A\text{-}\mathrm{𝐀𝐥𝐠}`$. ∎
Note that the essential image of $`MM_{}`$ is closed under limits. Next recall that the forgetful functor $`A_{}\text{-}\mathrm{𝐀𝐥𝐠}\mathrm{𝐒𝐞𝐭}`$ (resp. $`A_{}\text{-}\mathrm{𝐌𝐨𝐝}\mathrm{𝐒𝐞𝐭}`$) has a left adjoint $`A_{}[]:\mathrm{𝐒𝐞𝐭}A_{}\text{-}\mathrm{𝐀𝐥𝐠}`$ (resp. $`A^{()}:\mathrm{𝐒𝐞𝐭}A_{}\text{-}\mathrm{𝐌𝐨𝐝}`$) that assigns to a set $`S`$ the free $`A_{}`$-algebra $`A_{}[S]`$ (resp. the free $`A_{}`$-module $`A_{}^{(S)}`$) generated by $`S`$. If $`S`$ is any set, it is natural to write $`A[S]`$ (resp. $`A^{(S)}`$) for the $`A`$-algebra $`(A_{}[S])^a`$ (resp. for the $`A`$-module $`(A_{}^{(S)})^a`$. This yields a left adjoint, called the free $`A`$-algebra functor $`\mathrm{𝐒𝐞𝐭}A\text{-}\mathrm{𝐀𝐥𝐠}`$ (resp. the free $`A`$-module functor $`\mathrm{𝐒𝐞𝐭}A\text{-}\mathrm{𝐌𝐨𝐝}`$) to the “forgetful” functor $`A\text{-}\mathrm{𝐀𝐥𝐠}\mathrm{𝐒𝐞𝐭}`$ (resp. $`A\text{-}\mathrm{𝐌𝐨𝐝}\mathrm{𝐒𝐞𝐭}`$) $`BB_{}`$.
Now let $`R`$ be any $`V`$-algebra; we want to construct a left adjoint to the localisation functor $`R\text{-}\mathrm{𝐌𝐨𝐝}R^a\text{-}\mathrm{𝐌𝐨𝐝}`$. For a given $`R^a`$-module $`M`$, let
(2.2.8)
$$M_!=\stackrel{~}{𝔪}_V(M_{}).$$
We have the natural map (unit of adjunction) $`RR_{}^a`$, so that we can view $`M_!`$ as an $`R`$-module.
###### Proposition 2.2.9.
i) The functor $`R^a\text{-}\mathrm{𝐌𝐨𝐝}R\text{-}\mathrm{𝐌𝐨𝐝}`$ defined by (2.2.8) is left adjoint to localisation.
ii) The unit of the adjunction $`MM_!^a`$ is a natural isomorphism from the identity functor $`\text{1}_{R^a\text{-}\mathrm{𝐌𝐨𝐝}}`$ to the composition of the two functors.
###### Proof.
(i) follows easily from (2.2.2) and (ii) follows easily from (i). ∎
###### Corollary 2.2.10.
Suppose that $`\stackrel{~}{𝔪}`$ is a flat $`V`$-module. Then we have :
i) the functor $`MM_!`$ is exact;
ii) the localisation functor $`R\text{-}\mathrm{𝐌𝐨𝐝}R^a\text{-}\mathrm{𝐌𝐨𝐝}`$ sends injectives to injectives.
###### Proof.
By proposition 2.2.9 it follows that $`MM_!`$ is right exact. To show that it is also left exact when $`\stackrel{~}{𝔪}`$ is a flat $`V`$-module, it suffices to remark that $`MM_{}`$ is left exact. Now, by (i), the functor $`MM^a`$ is right adjoint to an exact functor, so (ii) is clear. ∎
Next, let $`B`$ be any $`A`$-algebra. The multiplication on $`B_{}`$ is inherited by $`B_!`$, which is therefore a non-unital ring in a natural way. We endow the $`V`$-module $`VB_!`$ with the ring structure determined by the rule: $`(v,b)(v^{},b^{})=(vv^{},vb^{}+v^{}b+bb^{})`$ for all $`v,v^{}V`$ and $`b,b^{}B_!`$. Then $`VB_!`$ is a (unital) ring. Let $`\mu _𝔪:\stackrel{~}{𝔪}𝔪`$ be the map defined by $`xyxy`$ for all $`x,y𝔪`$; we notice that the subset of all elements of the form $`(\mu (s),s\underset{¯}{1})`$ (for arbitrary $`s\stackrel{~}{𝔪}`$) forms an ideal $`I`$ of $`VB_!`$. Set $`B_{!!}=(VB_!)/I`$. Thus we have a sequence of $`V`$-modules
(2.2.11)
$$0\stackrel{~}{𝔪}VB_!B_{!!}0$$
which in general is only right exact.
###### Definition 2.2.12.
We say that $`B`$ is an exact $`A`$-algebra if the sequence (2.2.11) is exact.
###### Remark 2.2.13.
Notice that if $`\stackrel{~}{𝔪}\stackrel{}{}𝔪`$ (e.g. when $`𝔪`$ is flat), then all $`A`$-algebras are exact. In the general case, if $`B`$ is any $`A`$-algebra, then $`V^a\times B`$ is always exact. Indeed, we have $`(V^a\times B)_{}V_{}^a\times B_{}`$ and, by remark 2.1.3(i), $`\stackrel{~}{𝔪}_VV_{}^a\stackrel{~}{𝔪}`$.
Clearly we have a natural isomorphism $`BB_{!!}^a`$.
###### Proposition 2.2.14.
The functor $`BB_{!!}`$ is left adjoint to the localisation functor $`A_{!!}\text{-}\mathrm{𝐀𝐥𝐠}A\text{-}\mathrm{𝐀𝐥𝐠}`$.
###### Proof.
Let $`B`$ be an $`A`$-algebra, $`C`$ an $`A_{!!}`$-algebra and $`\varphi :BC^a`$ a morphism of $`A`$-algebras. By proposition 2.2.9 we obtain a natural $`A_{}`$-linear morphism $`B_!C`$. Together with the structure morphism $`VC`$ this yields a map $`\stackrel{~}{\varphi }:VB_!C`$ which is easily seen to be a ring homomorphism. It is equally clear that the ideal $`I`$ defined above is mapped to zero by $`\stackrel{~}{\varphi }`$, hence the latter factors through a map of $`A_{!!}`$-algebras $`B_{!!}C`$. Conversely, such a map induces a morhism of $`A`$-algebras $`BC^a`$ just by taking localisation. It is easy to check that the two procedures are inverse to each other, which shows the assertion. ∎
###### Remark 2.2.15.
The functor of almost elements commutes with arbitrary limits, because all right adjoints do. It does not in general commute with colimits, not even with arbitrary infinite direct sums. Dually, the functors $`MM_!`$ and $`BB_{!!}`$ commute with all colimits. In particular, the latter commutes with tensor products.
### 2.3. Almost homological algebra
In this section we fix an almost $`V`$-algebra $`A`$ and we consider various constructions in the category of $`A`$-modules.
###### Remark 2.3.1.
i) Let $`M_1,M_2`$ be two $`A`$-modules. By proposition 2.2.4 it is clear that a morphism $`\varphi :M_1M_2`$ of $`A`$-modules is uniquely determined by the induced morphism $`M_{1}^{}{}_{}{}^{}M_{2}^{}{}_{}{}^{}`$.
ii) It is a bit tricky to deal with preimages of almost elements under morphisms: for instance, if $`\varphi :M_1M_2`$ is an epimorphism (by which we mean that $`\text{Coker}(\varphi )0`$) and $`m_2M_{2}^{}{}_{}{}^{}`$, then it is not true in general that we can find an almost element $`m_1M_{1}^{}{}_{}{}^{}`$ such that $`\varphi _{}(m_1)=m_2`$. What remains true is that for arbitrary $`\epsilon 𝔪`$ we can find $`m_1`$ such that $`\varphi _{}(m_1)=\epsilon m_2`$.
The abelian category $`A\text{-}\mathrm{𝐌𝐨𝐝}`$ satisfies axiom (AB5) (see e.g. (§A.4)) and it has a generator, namely the object $`A`$ itself. It then follows by a general result that $`A\text{-}\mathrm{𝐌𝐨𝐝}`$ has enough injectives. By corollary 2.2.7 any inverse system $`\{M_n|n\}`$ of $`A`$-modules has an (inverse) limit $`\underset{n}{\text{lim}}M`$. As usual, we denote by $`\stackrel{1}{\text{lim}}`$ the right derived functor of the inverse limit functor. Notice that (Cor. 3.5.4) holds in the almost case since axiom (AB4\*) holds in $`A\text{-}\mathrm{𝐌𝐨𝐝}`$ (on the other hand, it is not clear whether (Lemma 3.5.3) holds under (AB4\*), since the proof uses elements).
###### Lemma 2.3.2.
Let $`\{M_n;\varphi _n:M_nM_{n+1}|n\}`$ (resp. $`\{N_n;\psi _n:N_{n+1}N_n|n\}`$) be a direct (resp. inverse) system of $`A`$-modules and morphisms and $`\{\epsilon _n|n\}`$ a sequence of ideals of $`V`$ converging to $`V`$ (for the uniform structure introduced in section 2.1).
i) If $`\epsilon _nM_n=0`$ for all $`n`$ then $`\underset{n}{\text{colim}}M_n0`$.
ii) If $`\epsilon _nN_n=0`$ for all $`n`$ then $`\underset{n}{\text{lim}}N_n0\underset{n}{\text{lim}}^1N_n`$.
iii) If $`\epsilon _n\text{Coker}(\psi _n)=0`$ for all $`n`$ and $`_{j=0}^{\mathrm{}}\epsilon _j`$ is a Cauchy product, then $`\underset{n}{\text{lim}}^1N_n0`$.
###### Proof.
(i) and (ii) : we remark only that $`\underset{n}{\text{lim}}^1N_n\underset{n}{\text{lim}}^1N_{n+p}`$ for all $`p`$ and leave the details to the reader. We prove (iii). From (Cor. 3.5.4) it follows easily that $`(\underset{n}{\text{lim}}^1N_n)^a\underset{n}{\text{lim}}^1N_n`$. It then suffices to show that $`\underset{n}{\text{lim}}^1N_n`$ is almost zero. We have $`\epsilon _n^2\text{Coker}(\psi _n)=0`$ and the product $`_{j=0}^{\mathrm{}}(\epsilon _j^2)`$ is again a Cauchy product. Next let $`N_n^{}=_{p0}\text{Im}(N_{n+p}N_n)`$. If $`J_n=_{p0}(\epsilon _n\epsilon _{n+1}\mathrm{}\epsilon _{n+p})^2`$ then $`J_nN_nN_n^{}`$ and $`\underset{n\mathrm{}}{\text{lim}}J_n=V`$. In view of (ii), $`\underset{n}{\text{lim}}^1N_n/N_n^{}`$ is almost zero, hence we reduce to showing that $`\underset{n}{\text{lim}}^1N_n^{}`$ is almost zero. But
$$J_{n+p+q}N_n^{}\text{Im}(N_{n+p+q}^{}N_n^{})\text{Im}(N_{n+p}^{}N_n^{})$$
for all $`n,p,q`$. On the other hand, by remark 2.1.2 we get $`_{q=0}^{\mathrm{}}𝔪J_{n+p+q}=𝔪`$, hence $`𝔪N_n^{}\text{Im}(N_{n+p}^{}N_n^{})`$ and finally $`𝔪N_n^{}=𝔪^2N_n^{}\text{Im}(𝔪N_{n+p}^{}𝔪N_n^{})`$ which means that $`\{𝔪N_n^{}\}`$ is a surjective inverse system, so its $`\stackrel{1}{\text{lim}}`$ vanishes and the result follows. ∎
###### Example 2.3.3.
Let $`(V,𝔪)`$ be as in example 2.1.1. Then every ideal in $`V`$ is principal, so in the situation of the lemma we can write $`\epsilon _j=(x_j)`$ for some $`x_jV`$. Then the hypothesis in (iii) can be stated by saying that there exists $`c`$ such that $`x_j0`$ for all $`jc`$ and the sequence $`n_{j=c}^n\nu (x_j)`$ is Cauchy in $`\mathrm{\Gamma }`$.
###### Definition 2.3.4.
Let $`M`$ be an $`A`$-module.
i) We say that $`M`$ is flat (resp. faithfully flat) if the functor $`NM_AN`$, from the category of $`A`$-modules to itself is exact (resp. exact and faithful). $`M`$ is almost projective if the functor $`N\text{alHom}_A(M,N)`$ is exact.
ii) We say that $`M`$ is finitely generated if there exists a positive integer $`n`$ and an epimorphism $`A^nM`$. We say that $`M`$ is almost finitely generated if, for arbitrary $`\epsilon 𝔪`$, there exists a finitely generated submodule $`M_\epsilon M`$ such that $`\epsilon MM_\epsilon `$.
iii) We say that $`M`$ is almost finitely presented if, for arbitrary $`\epsilon ,\delta 𝔪`$ there exist positive integers $`n=n(\epsilon )`$, $`m=m(\epsilon )`$ and a three term complex $`A^m\stackrel{\psi _\epsilon }{}A^n\stackrel{\varphi _\epsilon }{}M`$ with $`\epsilon \text{Coker}(\varphi _\epsilon )=0`$ and $`\delta \text{Ker}(\varphi _\epsilon )\text{Im}(\psi _\epsilon )`$.
###### Proposition 2.3.5.
(i) An $`A`$-module $`M`$ is almost finitely generated if and only if for every finitely generated ideal $`𝔪_0𝔪`$ there exists a finitely generated submodule $`M_0M`$ such that $`𝔪_0MM_0`$.
(ii) An $`A`$-module is almost finitely presented if and only if, for every finitely generated ideal $`𝔪_0𝔪`$ there is a complex $`A^m\stackrel{\psi }{}A^n\stackrel{\varphi }{}M`$ with $`𝔪_0\text{Coker}(\varphi )=0`$ and $`𝔪_0\text{Ker}(\varphi )\text{Im}(\psi )`$.
###### Proof.
(i) is easy and we leave it to the reader. To prove (ii), take a finitely generated ideal $`𝔪_1𝔪`$ such that $`𝔪_0𝔪𝔪_1`$, pick a morphism $`\varphi :A^nM`$ whose cokernel is annihilated by $`𝔪_1`$, and apply the following lemma 2.3.6. ∎
###### Lemma 2.3.6.
If $`M`$ is almost finitely presented and $`\varphi :FM`$ is a morphism with $`FA^n`$, then for every finitely generated ideal $`𝔪_1𝔪\text{Ann}_V(\text{Coker}(\varphi ))`$ there is a finitely generated submodule of $`\text{Ker}(\varphi )`$ containing $`𝔪_1\text{Ker}(\varphi )`$.
###### Proof.
We need the following
###### Claim 2.3.7.
Let $`F_1`$ be a finitely generated $`A`$-module and suppose that we are given $`a,bV`$ and a (not necessarily commutative) diagram
such that $`q\varphi =ap`$, $`p\psi =bq`$. Let $`IV`$ be an ideal such that $`\text{Ker}(q)`$ has a finitely generated submodule containing $`I\text{Ker}(q)`$. Then $`\text{Ker}(p)`$ has a finitely generated submodule containing $`abI\text{Ker}(p)`$.
Proof of the claim: Let $`R`$ be the submodule of $`\text{Ker}(q)`$ given by the assumption. We have $`\text{Im}(\psi \varphi ab\text{1}_{F_1})\text{Ker}(p)`$ and $`\psi (R)\text{Ker}(p)`$. We take $`R_1=\text{Im}(\psi \varphi ab\text{1}_{F_1})+\psi (R)`$. Clearly $`\varphi (\text{Ker}(p))\text{Ker}(q)`$, so $`I\varphi (\text{Ker}(p))R`$, hence $`I\psi \varphi (\text{Ker}(p))\psi (R)`$ and finally $`abI\text{Ker}(p)R_1`$.
Now, let $`\delta \text{Ann}_V(\text{Coker}(\varphi ))`$ and $`\epsilon _1,\epsilon _2,\epsilon _3,\epsilon _4𝔪`$. By assumption there is a complex $`A^r\stackrel{t}{}A^s\stackrel{q}{}M`$ with $`\epsilon _1\text{Coker}(q)=0`$, $`\epsilon _2\text{Ker}(q)\text{Im}(t)`$. Letting $`F_1=F`$, $`F_2=A^s`$, $`a=\epsilon _1\epsilon _3`$, $`b=\epsilon _4\delta `$, one checks easily that $`\psi `$ and $`\varphi `$ can be given such that all the assumptions of the above claim are fulfilled. So, with $`I=\epsilon _2V`$ we get that $`\epsilon _1\epsilon _2\epsilon _3\epsilon _4\delta \text{Ker}(\varphi )`$ lies in a finitely generated submodule of $`\text{Ker}(\varphi )`$. But $`𝔪_1`$ is contained in an ideal generated by finitely many such products $`\epsilon _1\epsilon _2\epsilon _3\epsilon _4\delta `$. ∎
The following proposition generalises a well-known characterization of finitely presented modules over usual rings.
###### Proposition 2.3.8.
Let $`M`$ be an $`A`$-module.
i) $`M`$ is almost finitely generated if and only if, for every filtered inductive system $`(N_\lambda ,\varphi _{\lambda \mu })`$ (indexed by a directed set $`\mathrm{\Lambda }`$) the natural morphism
$$\nu :\underset{\mathrm{\Lambda }}{\text{colim}}\text{alHom}_A(M,N_\lambda )\text{alHom}_A(M,\underset{\mathrm{\Lambda }}{\text{colim}}N_\lambda )$$
is a monomorphism.
ii) $`M`$ is almost finitely presented if and only if for every filtered inductive sytem as above, $`\nu `$ is an isomorphism.
###### Proof.
The “only if” part in (i) (resp. (ii)) is first checked when $`M`$ is finitely generated (resp. finitely presented) and then extended to the general case. We leave the details to the reader and we proceed to verify the “if” part. For (i), choose a set $`I`$ and an epimorphism $`p:A^{(I)}M`$. Let $`\mathrm{\Lambda }`$ be the directed set of finite subsets of $`I`$, ordered by inclusion. For $`S\mathrm{\Lambda }`$, let $`M_S=p(A^S)`$. Then $`\underset{\mathrm{\Lambda }}{\text{colim}}(M/M_S)=0`$, so the assumption gives $`\underset{\mathrm{\Lambda }}{\text{colim}}\text{alHom}_A(M,M/M_S)=0`$, i.e. $`\underset{\mathrm{\Lambda }}{\text{colim}}\text{Hom}_A(M,M/M_S)=0`$ is almost zero, so, for every $`\epsilon 𝔪`$, the image of $`\epsilon \text{1}_M`$ in the above colimit is $`0`$, i.e. there exists $`S\mathrm{\Lambda }`$ such that $`\epsilon MM_S`$, which proves the contention. For (ii), we present $`M`$ as a filtered colimit $`\underset{\mathrm{\Lambda }}{\text{colim}}M_\lambda `$, where each $`M_\lambda `$ is finitely presented (this can be done e.g. by taking such a presentation of the $`A_{}`$-module $`M_{}`$ and applying $`NN^a`$). The assumption of (ii) gives that $`\underset{\mathrm{\Lambda }}{\text{colim}}\text{Hom}_A(M,M_\lambda )\text{Hom}_A(M,M)`$ is an almost isomorphism, hence, for every $`\epsilon 𝔪`$ there is $`\lambda \mathrm{\Lambda }`$ and $`\varphi _\epsilon :MM_\lambda `$ such that $`p_\lambda \varphi _\epsilon =\epsilon \text{1}_M`$, where $`p_\lambda :M_\lambda M`$ is the natural morphism to the colimit. If such a $`\varphi _\epsilon `$ exists for $`\lambda `$, then it exists for every $`\mu \lambda `$. Hence, if $`𝔪_0𝔪`$ is a finitely generated subideal, say $`𝔪_0=_j^kV\epsilon _j`$, then there exist $`\lambda \mathrm{\Lambda }`$ and $`\varphi _i:MM_\lambda `$ such that $`p_\lambda \varphi _i=\epsilon _i\text{1}_M`$ for $`i=1,\mathrm{},k`$. Hence $`\text{Im}(\varphi _ip_\lambda \epsilon _i\text{1}_{M_\lambda })`$ is contained in $`\text{Ker}(p_\lambda )`$ and contains $`\epsilon _i\text{Ker}(p_\lambda )`$. Hence $`\text{Ker}(p_\lambda )`$ has a finitely generated submodule $`L`$ containing $`𝔪_0\text{Ker}(p_\lambda )`$. Choose a presentation $`A^mA^n\stackrel{\pi }{}M_\lambda `$. Then one can lift $`𝔪_0L`$ to a finitely generated submodule $`L^{}`$ of $`A^n`$. Then $`\text{Ker}(\pi )+L^{}`$ is a finitely generated submodule of $`\text{Ker}(p_\lambda \pi )`$ containing $`𝔪_0^2\text{Ker}(p_\lambda \pi )`$. Since we also have $`𝔪_0\text{Coker}(p_\lambda \pi )=0`$ and $`𝔪_0`$ is arbitrary, the conclusion follows from proposition 2.3.5. ∎
###### Lemma 2.3.9.
Let $`0M^{}MM^{\prime \prime }0`$ be an exact sequence of $`A`$-modules. Then:
i) If $`M^{}`$, $`M^{\prime \prime }`$ are almost finitely generated (resp. presented) then so is $`M`$.
ii) If $`M`$ is almost finitely presented, then $`M^{\prime \prime }`$ is almost finitely presented if and only if $`M^{}`$ is almost finitely generated.
###### Proof.
These facts can be deduced from proposition 2.3.8 and remark 2.3.11(iii), or proved directly. ∎
###### Lemma 2.3.10.
Let $`𝐏`$ be one of the properties : “flat”, “almost projective”, “almost finitely generated”, “almost finitely presented”. If $`B`$ is a $`𝐏`$ $`A`$-algebra, and $`M`$ is a $`𝐏`$ $`B`$-module, then $`M`$ is $`𝐏`$ as an $`A`$-module.
###### Proof.
Left to the reader. ∎
Let $`R`$ be a $`V`$-algebra and $`M`$ a flat (resp. faithfully flat) $`R`$-module (in the usual sense, see p.45). Then $`M^a`$ is a flat (resp. faithfully flat) $`R^a`$-module. Indeed, the functor $`M_R`$ preserves the Serre subcategory of almost zero modules, so by general facts it induces an exact functor on the localized categories (cp. p.369). For the faithfullness we have to show that an $`R`$-module $`N`$ is almost zero whenever $`M_RN`$ is almost zero. However, $`M_RN`$ is almost zero $``$ $`M_R(𝔪_VN)=0`$ $``$ $`𝔪_VN=0`$ $``$ $`N`$ is almost zero. It is clear that $`A\text{-}\mathrm{𝐌𝐨𝐝}`$ has enough almost projective (resp. flat) objects. Let $`R`$ be a $`V`$-algebra. The localisation functor induces a functor $`G:\text{D}(R)\text{D}(R^a)`$ and, in view of corollary 2.2.10, $`MM_!`$ induces a functor $`F:\text{D}(R^a)\text{D}(R)`$. We have a natural isomorphism $`GF\text{1}_{\text{D}(R^a)}`$ and a natural transformation $`FG\text{1}_{\text{D}(R)}`$. These satisfy the triangular identities of (p.83) so $`F`$ is a left adjoint to $`G`$. If $`\mathrm{\Sigma }`$ denotes the multiplicative set of morphisms in $`\text{D}(R)`$ which induce almost isomorphisms on the cohomology modules, then the localised category $`\mathrm{\Sigma }^1\text{D}(R)`$ exists (see e.g. (Th.10.3.7)) and by the same argument we get an equivalence of categories $`\mathrm{\Sigma }^1\text{D}(R)\text{D}(R^a)`$.
Given an $`A`$-module $`M`$, we can derive the functors $`M_A`$ (resp. $`\text{alHom}_A(M,)`$, resp. $`\text{alHom}_A(,M)`$) by taking flat (resp. injective, resp. almost projective) resolutions : one remarks that bounded above exact complexes of flat (resp. almost projective) $`A`$-modules are acyclic for the functor $`M_A`$ (resp. $`\text{alHom}_A(,M)`$) (recall the standard argument: if $`F_{}`$ is a complex of flat $`A`$-modules, let $`\mathrm{\Phi }_{}`$ be a flat resolution of $`M`$; then $`\text{Tot}(\mathrm{\Phi }_{}_AF_{})M_AF_{}`$ is a quasi-isomorphism since it is so on rows, and $`\text{Tot}(\mathrm{\Phi }_{}_AF_{})`$ is acyclic since its colums are; similarly, if $`P_{}`$ is a complex of almost projective objects, one considers the double complex $`\text{alHom}_A(P_{},J^{})`$ where $`J^{}`$ is an injective resolution of $`M`$; cp. §2.7); then one uses the construction detailed in (Th.10.5.9). We denote by $`\text{Tor}_i^A(M,)`$ (resp. $`\text{alExt}_A^i(M,)`$, resp. $`\text{alExt}_A^i(,M)`$) the corresponding derived functors. If $`A=R^a`$ for some $`V`$-algebra $`R`$, we obtain easily natural isomorphisms $`\text{Tor}_i^R(M,N)^a\text{Tor}_i^A(M^a,N^a)`$ for all $`R`$-modules $`M,N`$. A similar result holds for $`\text{Ext}_R^i(M,N)`$.
###### Remark 2.3.11.
i) Clearly, an $`A`$-module $`M`$ is flat (resp. almost projective) if and only if $`\text{Tor}_i^A(M,N)=0`$ (resp. $`\text{alExt}_A^i(M,N)=0`$) for all $`A`$-modules $`N`$ and all $`i>0`$.
ii) Let $`M,N`$ be two flat (resp. almost projective) $`A`$-modules. Then $`M_AN`$ is a flat (resp. almost projective) $`A`$-module and for any $`A`$-algebra $`B`$, the $`B`$-module $`B_AM`$ is flat (resp. almost projective).
iii) Resume the notation of proposition 2.3.8. If $`M`$ is almost finitely presented, then one has also that the natural morphism $`\underset{\mathrm{\Lambda }}{\text{colim}}\text{alExt}_A^1(M,N_\lambda )\text{alExt}_A^1(M,\underset{\mathrm{\Lambda }}{\text{colim}}N_\lambda )`$ is a monomorphism. This is deduced from proposition 2.3.8(ii), using the fact that $`(N_\lambda )`$ can be injected into an inductive system $`(J_\lambda )`$ of injective almost modules (e.g. $`J_\lambda =E^{\text{Hom}_A(N_\lambda ,E)}`$, where $`E`$ is an injective cogenerator for $`A\text{-}\mathrm{𝐌𝐨𝐝}`$), and by applying alExt sequences.
###### Lemma 2.3.12.
Let $`M`$ be an almost finitely generated $`A`$-module. Consider the following properties:
i) $`M`$ is almost projective.
ii) For arbitrary $`\epsilon 𝔪`$ there exist $`n(\epsilon )`$ and $`A`$-linear morphisms
(2.3.13)
such that $`v_\epsilon u_\epsilon =\epsilon \text{1}_M`$.
iii) $`M`$ is flat.
Then (i) $``$ (ii) $``$ (iii).
###### Proof.
(ii)$``$(i): for given $`\epsilon 𝔪`$, we consider any $`A`$-module $`N`$ and we apply the functor $`\text{alExt}_A^i(,N)`$ to (2.3.13) :
which implies $`\epsilon \text{alExt}^i(M,N)=0`$ for all $`i>0`$. Since $`\epsilon `$ is arbitrary, (i) follows from remark 2.3.11(i).
(i)$``$(ii): by hypothesis, for arbitrary $`\epsilon 𝔪`$ we can find $`n=n(\epsilon )`$ and a morphism $`\varphi _\epsilon :A^nM`$ such that $`\epsilon \text{Coker}(\varphi _\epsilon )=0`$. Let $`M_\epsilon `$ be the image of $`\varphi _\epsilon `$, so that $`\varphi _\epsilon `$ factors as $`A^{n(\epsilon )}\stackrel{\psi _\epsilon }{}M_\epsilon \stackrel{j_\epsilon }{}M`$. Also $`\epsilon \text{1}_M:MM`$ factors as $`M\stackrel{\gamma _\epsilon }{}M_\epsilon \stackrel{j_\epsilon }{}M.`$ Since by hypothesis $`M`$ is almost projective, the natural morphism induced by $`\psi _\epsilon `$
is an epimorphism. Then for arbitrary $`\delta 𝔪`$ the morphism $`\delta \gamma _\epsilon `$ is in the image of $`\psi _\epsilon ^{}`$, in other words, there exists an $`A`$-linear morphism $`u_{\epsilon \delta }:MA^n`$ such that $`\psi _\epsilon u_{\epsilon \delta }=\delta \gamma _\epsilon `$. If now we take $`v_{\epsilon \delta }=\varphi _\epsilon `$, it is clear that $`v_{\epsilon \delta }u_{\epsilon \delta }=\epsilon \delta \text{1}_M`$. This proves (ii), since the $`\epsilon 𝔪`$ satisfying the assertion of (ii) form an ideal.
(ii)$``$(iii): for a given $`A`$-module $`N`$, apply the functor $`\text{Tor}_i^A(,N)`$ to the sequence (2.3.13). This yields $`\epsilon \text{Tor}_i^A(M,N)=0`$. Now the claim follows from remark 2.3.11(i). ∎
There is a converse to lemma 2.3.12 in case $`M`$ is almost finitely presented. Before stating it, we need the following lemma.
###### Lemma 2.3.14.
Let $`R`$ be any ring, $`M`$ any $`R`$-module and $`C=\text{Coker}(\varphi :R^nR^m)`$ any finitely presented (left) $`R`$-module. Let $`C^{}=\text{Coker}(\varphi ^{}:R^mR^n)`$ be the cokernel of the transpose of the map $`\varphi `$. Then there is a natural isomorphism
$$\text{Tor}_1^R(C^{},M)\text{Hom}_R(C,M)/\text{Im}(\text{Hom}_R(C,R)_RM).$$
###### Proof.
We have a spectral sequence :
$$E_{ij}^2=\text{Tor}_i^R(H_j(\text{Cone}(\varphi ^{})),M)H_{i+j}(\text{Cone}(\varphi ^{})_RM).$$
On the other hand we have also natural isomorphisms
$$\text{Cone}(\varphi ^{})_RM\text{Hom}_R(\text{Cone}(\varphi ),R)[1]_RM\text{Hom}_R(\text{Cone}(\varphi ),M)[1].$$
Hence :
$$\begin{array}{cc}\hfill E_{10}^2E_{10}^{\mathrm{}}H_1(\text{Cone}(\varphi ^{})_RM)/E_{01}^{\mathrm{}}& H^0(\text{Hom}_R(\text{Cone}(\varphi ),M))/\text{Im}(E_{01}^2)\hfill \\ \hfill & \text{Hom}_R(C,M)/\text{Im}(\text{Hom}_R(C,R)_RM)\hfill \end{array}$$
which is the claim. ∎
###### Proposition 2.3.15.
(i) Every almost finitely generated almost projective $`A`$-module is almost finitely presented.
(ii) Every almost finitely presented flat $`A`$-module is almost projective.
###### Proof.
(ii) : let $`M`$ be such an $`A`$-module. Let $`\epsilon ,\delta 𝔪`$ and pick a three term complex
such that $`\epsilon \text{Coker}(\varphi )=\delta \text{Ker}(\varphi )/\text{Im}(\psi )=0`$. Set $`P=\text{Coker}(\psi _{})`$; this is a finitely presented $`A_{}`$-module and $`\varphi _{}`$ factors through a morphism $`\overline{\varphi }_{}:PM_{}`$. Let $`\gamma 𝔪`$; from lemma 2.3.14 we see that $`\gamma \overline{\varphi }`$ is the image of some element $`_{j=1}^n\varphi _jm_j\text{Hom}_A_{}(P,A_{})_A_{}M_{}`$. If we define $`L=A_{}^n`$ and $`v:PL`$, $`w:LM_{}`$ by $`v(x)=(\varphi _1(x),\mathrm{},\varphi _n(x))`$ and $`w(y_1,\mathrm{},y_n)=_{j=1}^ny_jm_j`$, then clearly $`\gamma \overline{\varphi }=wv`$. Let $`K=\text{Ker}(\overline{\varphi }_{})`$. Then $`\delta K^a=0`$ and the map $`\delta \text{1}_{P^a}`$ factors through a morphism $`\sigma :(P/K)^aP^a`$. Similarly the map $`\epsilon \text{1}_M`$ factors through a morphism $`\lambda :M(P/K)^a`$. Let $`\alpha =v^a\sigma \lambda :ML^a`$ and $`\beta =w^a:L^aM`$. The reader can check that $`\beta \alpha =\epsilon \delta \gamma \text{1}_M`$. By lemma 2.3.12 the claim follows.
(i) : let $`P`$ be such an almost finitely generated almost projective $`A`$-module. For any finitely generated ideal $`𝔪_0𝔪`$ pick a morphism $`\varphi :A^rP`$ such that $`𝔪_0\text{Coker}(\varphi )=0`$. If $`\epsilon _1,\mathrm{},\epsilon _k`$ is a set of generators for $`𝔪_0`$, a standard argument shows that, for any $`ik`$, $`\epsilon _i\text{1}_P`$ lifts to a morphism $`\psi _i:PA^r/\text{Ker}(\varphi )`$; then, since $`P`$ is almost projective, $`\epsilon _j\psi _i`$ lifts to a morphism $`\psi _{ij}:PA^r`$. Now claim 2.3.7 applies with $`F_1=A^r`$, $`F_2=M=P`$, $`p=\varphi `$, $`q=\text{1}_P`$ and $`\psi =\psi _{ij}`$ and shows that $`\text{Ker}(\varphi )`$ has a finitely generated submodule $`M_{ij}`$ containing $`\epsilon _i\epsilon _j\text{Ker}(\varphi )`$. Then the span of all such $`M_{ij}`$ is a finitely generated submodule of $`\text{Ker}(\varphi )`$ containing $`𝔪_0^2\text{Ker}(\varphi )`$. By proposition 2.3.5(ii), the claim follows. ∎
###### Definition 2.3.16.
For an $`A`$-module $`M`$, the dual $`A`$-module of $`M`$ is the $`A`$-module $`M^{}=\text{alHom}_A(M,A)`$. $`M`$ is reflexive if the natural morphism
(2.3.17)
$$M(M^{})^{}m(ff(m))$$
is an isomorphism of $`A`$-modules.
###### Remark 2.3.18.
Notice that if $`B`$ is an $`A`$-algebra and $`M`$ any $`B`$-module, then by “restriction of scalars” $`M`$ is also an $`A`$-module and the dual $`A`$-module $`M^{}`$ has a natural structure of $`B`$-module. This is defined by the rule $`(bf)(m)=f(bm)`$ ($`bB_{}`$, $`mM_{}`$ and $`fM_{}^{}`$). With respect to this structure (2.3.17) becomes a $`B`$-linear morphism. Incidentally, notice that the two meanings of “$`M_{}^{}`$” coincide, i.e. $`(M_{})^{}(M^{})_{}`$.
###### Proposition 2.3.19.
Let $`P`$ be an almost projective $`A`$-module and denote by $`I_P`$ the image of the natural “evaluation morphism” $`P_AP^{}A`$.
i) For every morphism of algebras $`AB`$ we have $`I_{B_AP}=I_PB`$.
ii) $`I_P=I_P^2`$.
iii) $`P=0`$ if and only if $`I_P=0`$.
iv) $`P`$ is faithfully flat if and only if $`I_P=A`$.
###### Proof.
Pick an indexing set $`I`$ large enough, and an epimorphism $`\varphi :F=A^{(I)}P`$. For every $`iI`$ we have the standard morphisms $`A\stackrel{e_i}{}F\stackrel{\pi _i}{}A`$ such that $`\pi _ie_j=\delta _{ij}\text{1}_A`$ and $`_{iI}e_i\pi _i=\text{1}_F`$. For every $`x𝔪`$ choose $`\psi _x\text{Hom}_A(P,F)`$ such that $`\varphi \psi _x=x\text{1}_P`$. It is easy to check that $`I_P`$ is generated by the almost elements $`\pi _i\psi _x\varphi e_j`$ ($`i,jI`$, $`x𝔪`$). (i) follows already. For (iii), the “only if” is clear; if $`I_P=0`$, then $`\psi _x\varphi =0`$ for all $`x𝔪`$, hence $`\psi _x=0`$ and therefore $`x\text{1}_P=0`$. Next, notice that, from (i) and (iii) we derive $`P/(I_PP)=0`$, i.e. $`P=I_PP`$, so (ii) follows directly from the definition of $`I_P`$. Since $`P`$ is flat, to show (iv) we have only to verify that the functor $`MP_AM`$ is faithful. To this purpose, it suffices to check that $`P_A(A/J)0`$ for every proper ideal $`J`$ of $`A`$. This follows easily from (i) and (iii). ∎
If $`E`$, $`F`$ and $`N`$ are $`A`$-modules, there is a natural morphism :
(2.3.20)
$$E_A\text{alHom}_A(F,N)\text{alHom}_A(F,E_AN).$$
###### Lemma 2.3.21.
(i) The morphism (2.3.20) is an isomorphism in the following cases :
a) when $`E`$ is flat and $`F`$ is almost finitely presented;
b) when either $`E`$ or $`F`$ is almost finitely generated and almost projective;
c) when $`F`$ is almost projective and $`E`$ is almost finitely presented;
d) when $`E`$ is almost projective and $`F`$ is almost finitely generated.
(ii) The morphism (2.3.20) is a monomorphism in the following cases :
a) when $`E`$ is flat and $`F`$ is almost finitely generated;
b) when $`E`$ is almost projective.
(iii) The morphism (2.3.20) is an epimorphism when $`F`$ is almost projective and $`E`$ is almost finitely generated.
###### Proof.
If $`FA^{(I)}`$ for some finite set $`I`$, then $`\text{alHom}_A(F,N)N^{(I)}`$ and the claims are obvious. More generally, if $`F`$ is almost projective and almost finitely generated, for any $`\epsilon 𝔪`$ there exists a finite set $`I=I(\epsilon )`$ and morphisms
(2.3.22)
such that $`v_\epsilon u_\epsilon =\epsilon \text{1}_F`$. We apply the natural transformation
$$E_A\text{alHom}_A(,N)\text{alHom}_A(,E_AN)$$
to (2.3.22) : an easy diagram chase allows then to conclude that the kernel and cokernel of (2.3.20) are killed by $`\epsilon `$. As $`\epsilon `$ is arbitrary, it follows that (2.3.20) is an isomorphism in this case. An analogous argument works when $`E`$ is almost finitely generated almost projective, so we get (i.b). If $`F`$ is only almost projective, then we still have morphisms of the type (2.3.22), but now $`I(\epsilon )`$ is no longer necessarily finite. However, the cokernels of the induced morphisms $`\text{1}_Eu_\epsilon `$ and $`\text{alHom}_A(v_\epsilon ,E_AN)`$ are still annihilated by $`\epsilon `$. Hence, to show (iii) (resp. (i.c)) it suffices to consider the case when $`F`$ is free and $`E`$ is almost finitely generated (resp. presented). By passing to almost elements, we can further reduce to the analogous question for usual rings and modules, and by the usual juggling we can even replace $`E`$ by a finitely generated (resp. presented) $`A_{}`$-module and $`F`$ by a free $`A_{}`$-module. This case is easily dealt with, and (iii) and (i.c) follow. Case (i.d) (resp. (ii.b)) is similar : one considers almost elements and replaces $`E_{}`$ by a free $`A_{}`$-module (resp. and $`F_{}`$ by a finitely generated $`A_{}`$-module). In case (ii.a) (resp. (i.a)), for every finitely generated submodule $`𝔪_0`$ of $`𝔪`$ we can find, by proposition 2.3.5, a finitely generated (resp. presented) $`A`$-module $`F_0`$ and a morphism $`F_0F`$ whose kernel and cokernel are annihilated by $`𝔪_0`$. It follows easily that we can replace $`F`$ by $`F_0`$ and suppose that $`F`$ is finitely generated (resp. presented). Then the argument in (Ch.I §2 Prop.10) can be taken over verbatim to show (ii.a) (resp. (i.a)). ∎
###### Lemma 2.3.23.
i) Let $`P`$ be an $`A`$-module and $`B`$ an $`A`$-algebra. If either $`P`$ or $`B`$ is almost finitely generated almost projective as an $`A`$-module, then the natural morphism
(2.3.24)
$$B_A\text{alHom}_A(P,N)\text{alHom}_B(B_AP,B_AN)$$
is an isomorphism for all $`A`$-modules $`N`$.
ii) Every almost projective almost finitely generated $`A`$-module is reflexive.
###### Proof.
(i) is an easy consequence of lemma 2.3.21(i.b). To prove (ii), we we apply the natural transformation (2.3.17) to (2.3.22) : by diagram chase one sees that the kernel and cokernel of the morphism $`F(F^{})^{}`$ are killed by $`\epsilon `$. ∎
###### Lemma 2.3.25.
Let $`\{M_n;\varphi _n:M_nM_{n+1}|n\}`$ be a direct system of $`A`$-modules and suppose there exist sequences $`\{\epsilon _n|n\}`$ and $`\{\delta _n|n\}`$ of ideals of $`V`$ such that
i) $`\underset{n\mathrm{}}{\text{lim}}\epsilon _n=V`$ (convergence for the uniform structure on ideals of $`V`$) and $`_{j=0}^{\mathrm{}}\delta _j`$ is a Cauchy product;
ii) for all $`n`$ there exist integers $`N(n)`$ and morphisms of $`A`$-modules $`\psi _n:A^{N(n)}M_n`$ such that $`\epsilon _n\text{Coker}(\psi _n)=0`$;
iii) $`\delta _n\text{Coker}(\varphi _n)=0`$ for all $`n`$.
Then $`\underset{n}{\text{colim}}M_n`$ is an almost finitely generated $`A`$-module.
###### Proof.
Let $`M=\underset{n}{\text{colim}}M_n`$. For any $`n`$ let $`a_n=_{m0}(_{j=n}^{n+m}\delta _j)`$. Then $`\underset{n\mathrm{}}{\text{lim}}a_n=V`$. For $`m>n`$ set $`\varphi _{n,m}=\varphi _m\mathrm{}\varphi _{n+1}\varphi _n:M_nM_{m+1}`$ and let $`\varphi _{n,\mathrm{}}:M_nM`$ be the natural morphism. An easy induction shows that $`_{j=n}^m\delta _j\text{Coker}(\varphi _{n,m})=0`$ for all $`m>n`$. Since $`\text{Coker}(\varphi _{n,\mathrm{}})=\underset{m}{\text{colim}}\text{Coker}(\varphi _{n,m})`$ we obtain $`a_n\text{Coker}(\varphi _{n,\mathrm{}})=0`$ for all $`n`$. Therefore $`\epsilon _na_n\text{Coker}(\varphi _{n,\mathrm{}}\psi _n)=0`$ for all $`n`$. Since $`\underset{n\mathrm{}}{\text{lim}}\epsilon _na_n=V`$, the claim follows. ∎
###### Lemma 2.3.26.
Let $`\{M_n;\varphi _n:M_nM_{n+1}|n\}`$ be a direct system of $`A`$-modules and suppose there exist sequences $`\{\epsilon _n|n\}`$ and $`\{\delta _n|n\}`$ of ideals of $`V`$ such that
i) $`\underset{n\mathrm{}}{\text{lim}}\epsilon _n=V`$ and $`_{j=0}^{\mathrm{}}\delta _j`$ is a Cauchy product;
ii) $`\epsilon _n\text{alExt}_A^i(M_n,N)=\delta _n\text{alExt}_A^i(\text{Coker}(\varphi _n),N)=0`$ for all $`A`$-modules $`N`$, all $`i>0`$ and all $`n`$;
iii) $`\delta _n\text{Ker}(\varphi _n)=0`$ for all $`n`$.
Then $`\underset{n}{\text{colim}}M_n`$ is an almost projective $`A`$-module.
###### Proof.
Let $`M=\underset{n}{\text{colim}}M_n`$. By the above remark 2.3.11(i) it suffices to show that $`\text{alExt}_A^i(M,N)`$ vanishes for all $`i>0`$ and all $`A`$-modules $`N`$. The maps $`\varphi _n`$ define a map $`\varphi :_nM_n_nM_n`$ such that we have a short exact sequence $`0_nM_n\stackrel{\text{1}\varphi }{}_nM_nM0`$. Applying the long exact alExt sequence one obtains a short exact sequence (cp. (3.5.10))
$$0\underset{n}{\text{lim}}^1\text{alExt}_A^{i1}(M_n,N)\text{alExt}_A^i(M,N)\underset{n}{\text{lim}}\text{alExt}_A^i(M_n,N)0.$$
Then lemma 2.3.2(ii) implies that $`\text{alExt}_A^i(M,N)0`$ for all $`i>1`$ and moreover $`\text{alExt}_A^1(M,N)`$ is isomorphic to $`\underset{n}{\text{lim}}^1\text{alHom}_A(M_n,N)`$. Let
$$\varphi _n^{}:\text{alHom}_A(M_{n+1},N)\text{alHom}_A(M_n,N)ff\varphi _n$$
be the transpose of $`\varphi _n`$ and write $`\varphi _n`$ as a composition $`M_n\stackrel{p_n}{}\text{Im}(\varphi _n)\stackrel{q_n}{}M_{n+1}`$, so that $`\varphi _n^{}=q_n^{}p_n^{}`$, the composition of the respective transposed morphims. We have monomorphisms
$$\begin{array}{cc}& \text{Coker}(p_n^{})\text{alHom}_A(\text{Ker}(\varphi _n),N)\hfill \\ & \text{Coker}(q_n^{})\text{alExt}_A^1(\text{Coker}(\varphi _n),N)\hfill \end{array}$$
for all $`n`$. Hence $`\delta _n^2\text{Coker}(\varphi _n^{})=0`$ for all $`n`$. Since $`_{n=0}^{\mathrm{}}\delta _n^2`$ is a Cauchy product, lemma 2.3.2(iii) shows that $`\underset{n}{\text{lim}}^1\text{alHom}_A(M_n,N)0`$ and the assertion follows. ∎
###### Proposition 2.3.27.
Suppose that $`\stackrel{~}{𝔪}`$ is a flat $`V`$-module. Then for any $`V`$-algebra $`R`$ the functor $`MM_!`$ commutes with tensor products and takes flat $`R^a`$-modules to flat $`R`$-modules.
###### Proof.
Let $`M`$ be a flat $`R^a`$-module and $`NN^{}`$ an injective map of $`R`$-modules. Denote by $`K`$ the kernel of the induced map $`M_!_RNM_!_RN^{}`$; we have $`K^a0`$. We obtain an exact sequence $`0\stackrel{~}{𝔪}_VK\stackrel{~}{𝔪}_VM_!_RN\stackrel{~}{𝔪}_VM_!_RN^{}`$. But one sees easily that $`\stackrel{~}{𝔪}_VK=0`$ and $`\stackrel{~}{𝔪}_VM_!M_!`$, which shows that $`M_!`$ is a flat $`R`$-module. Similarly, let $`M,N`$ be two $`R^a`$-modules. Then the natural map $`M_{}_RN_{}(M_{R^a}N)_{}`$ is an almost isomorphism and the assertion follows from remark 2.1.3(i). ∎
### 2.4. Almost homotopical algebra
The formalism of abelian tensor categories provides a minimal framework wherein the rudiments of deformation theory can be developed.
Let $`(\text{C},,U)`$ be an abelian tensor category; we assume henceforth that $``$ is a right exact functor. Let $`A`$ be a given C-monoid. A two-sided ideal of $`A`$ is an $`A`$-sub-bimodule $`IA`$. The quotient $`A/I`$ in the underlying abelian category C has a unique C-monoid structure such that $`AA/I`$ is a morphism of monoids. $`A/I`$ is unitary if $`A`$ is. For $`I,J`$ subobjects of $`A`$ one denotes $`IJ=\text{Im}(IJAA\stackrel{\mu _A}{}A)`$. If $`I`$ is a two-sided ideal of $`A`$ such that $`I^2=0`$, then, using the right exactness of $``$ one checks that $`I`$ has a natural structure of an $`A/I`$-bimodule, unitary when $`A`$ is.
###### Definition 2.4.1.
A C-extension of a C-monoid $`B`$ by a $`B`$-bimodule $`I`$ is a short exact sequence of objects of C
(2.4.2)
such that $`C`$ is a C-monoid, $`p`$ is a morphism of C-monoids, $`I`$ is a square zero two-sided ideal in $`C`$ and the $`E/I`$-bimodule structure on $`I`$ coincides with the given bimodule structure on $`I`$. The C-extensions form a category $`\mathrm{𝐄𝐱𝐦𝐨𝐧}_\text{C}`$. The morphisms are commutative diagrams with exact rows
such that $`g`$ and $`h`$ are morphisms of C-monoids. We let $`\mathrm{𝐄𝐱𝐦𝐨𝐧}_\text{C}(B,I)`$ be the subcategory of $`\mathrm{𝐄𝐱𝐦𝐨𝐧}_\text{C}`$ consisting of all C-extensions of $`B`$ by $`I`$, where the morphisms are all short exact sequences as above such that $`f=\text{1}_I`$ and $`h=\text{1}_B`$.
We have also the variant in which all the C-monoids in (2.4.2) are required to be unitary (resp. to be algebras) and $`I`$ is a unitary $`B`$-bimodule (resp. whose left and right $`B`$-module actions coincide, i.e. are switched by composition with the “commutativity constraints” $`\eta _{B|I}`$ and $`\eta _{I|B}`$, see 2.2); we will call $`\mathrm{𝐄𝐱𝐮𝐧}_\text{C}`$ (resp. $`\mathrm{𝐄𝐱𝐚𝐥}_\text{C}`$) the corresponding category. For a morphism $`\varphi :CB`$ of C-monoids, and a C-extension $`X`$ in $`\mathrm{𝐄𝐱𝐦𝐨𝐧}_\text{C}(B,I)`$, we can pullback $`X`$ via $`\varphi `$ to obtain an exact sequence $`X\varphi `$ with a morphism $`\varphi ^{}:X\varphi X`$; one checks easily that there exists a unique structure of C-extension on $`X\varphi `$ such that $`\varphi ^{}`$ is a morphism of C-extension; then $`X\varphi `$ is an object in $`\mathrm{𝐄𝐱𝐦𝐨𝐧}_\text{C}(C,I)`$. Similarly, given a $`B`$-linear morphism $`\psi :IJ`$, we can push out $`X`$ and obtain a well defined object $`\psi X`$ in $`\mathrm{𝐄𝐱𝐦𝐨𝐧}_\text{C}(B,J)`$ with a morphism $`X\psi X`$ of $`\mathrm{𝐄𝐱𝐦𝐨𝐧}_\text{C}`$. In particular, if $`I_1`$ and $`I_2`$ are two $`B`$-bimodules, the functors $`p_i`$ ($`i=1,2`$) associated to the natural projections $`p_i:I_1I_2I_i`$ establish an equivalence of categories
(2.4.3)
$$\mathrm{𝐄𝐱𝐦𝐨𝐧}_\text{C}(B,I_1I_2)\stackrel{}{}\mathrm{𝐄𝐱𝐦𝐨𝐧}_\text{C}(B,I_1)\times \mathrm{𝐄𝐱𝐦𝐨𝐧}_\text{C}(B,I_2)$$
whose essential inverse is given by $`(E_1,E_2)(E_1E_2)\delta `$, where $`\delta :BBB`$ is the diagonal morphism. A similar statement holds for $`\mathrm{𝐄𝐱𝐚𝐥}`$ and $`\mathrm{𝐄𝐱𝐮𝐧}`$. These operations can be used to induce an abelian group structure on the set $`\text{Exmon}_\text{C}(B,I)`$ of isomorphism classes of objects of $`\mathrm{𝐄𝐱𝐦𝐨𝐧}_\text{C}(B,I)`$ as follows. For any two objects $`X,Y`$ of $`\mathrm{𝐄𝐱𝐦𝐨𝐧}_\text{C}(B,I)`$ we can form $`XY`$ which is an object of $`\mathrm{𝐄𝐱𝐦𝐨𝐧}_\text{C}(BB,II)`$. Let $`\alpha :III`$ be the addition morphism of $`I`$. Then we set $`X+Y=\alpha (XY)\delta `$. One can check that $`X+YY+X`$ for any $`X,Y`$ and that the trivial split C-extension $`BI`$ is a neutral element for $`+`$. Moreover every isomorphism class has an inverse $`X`$. The functors $`XX\varphi `$ and $`X\psi X`$ commute with the operation thus defined, and induce group homomorphisms
$$\begin{array}{c}\varphi :\text{Exmon}_\text{C}(B,I)\text{Exmon}_\text{C}(C,I)\hfill \\ \psi :\text{Exmon}_\text{C}(B,I)\text{Exmon}_\text{C}(B,J).\hfill \end{array}$$
We will need the variant $`\text{Exal}_\text{C}(B,I)`$ defined in the same way, starting from $`\mathrm{𝐄𝐱𝐚𝐥}_\text{C}(B,I)`$. For instance, if $`A`$ is an almost algebra (resp. a commutative ring), we can consider the abelian tensor category $`\text{C}=A\text{-}\mathrm{𝐌𝐨𝐝}`$. In this case the C-extensions will be called simply $`A`$-extensions, and we will write $`\mathrm{𝐄𝐱𝐚𝐥}_A`$ rather than $`\mathrm{𝐄𝐱𝐚𝐥}_\text{C}`$. In fact the commutative unitary case will soon become prominent in our work, and the more general setup is only required for technical reasons, in the proof of proposition 2.4.6 below, which is the abstract version of a well-known result on the lifting of idempotents over nilpotent ring extensions.
Let $`A`$ be a C-monoid. We form the biproduct $`A^{}=UA`$ in C. We denote by $`p_1`$, $`p_2`$ the associated projections from $`A^{}`$ onto $`U`$ and respectively $`A`$. Also, let $`i_1`$, $`i_2`$ be the natural monomorphisms from $`U`$, resp. $`A`$ to $`A^{}`$. $`A^{}`$ is equipped with a unitary monoid structure
$$\mu ^{}=i_2\mu (p_2p_2)+i_2\mathrm{}_A^1(p_1p_2)+i_2r_A^1(p_2p_1)+i_1u^1(p_1p_1)$$
where $`\mathrm{}_A`$, $`r_A`$ are the natural isomorphisms provided by (Prop. 1.3) and $`u:UUU`$ is as in loc. cit. §1. In terms of the ring $`A_{}^{}U_{}A_{}`$ this is the multiplication $`(u_1,b_1)(u_2,b_2)=(u_1u_2,b_1b_2+b_1u_2+u_1b_2)`$. Then $`i_2`$ is a morphism of monoids and one verifies that the “restriction of scalars” functor $`i_2^{}`$ defines an equivalence from the category $`A^{}\text{-}\mathrm{𝐔𝐧𝐢}.\mathrm{𝐌𝐨𝐝}`$ of unitary $`A^{}`$-modules to the category $`A\text{-}\mathrm{𝐌𝐨𝐝}`$ of all $`A`$-modules; let $`j`$ denote the inverse functor. A similar discussion applies to bimodules. Similarly, we derive equivalences of categories
for all $`A`$-bimodules $`M`$.
Next we specialise to $`A=U`$ : for a given $`U`$-module $`M`$ let $`e_M=\sigma _M\mathrm{}_M:MM`$; working out the definitions one finds that the condition that this is a module structure is equivalent to $`e_M^2=e_M`$. Let $`U\times U`$ be the product of $`U`$ by itself in the category of C-monoids. There is an isomorphism of unitary C-monoids $`\zeta :U^{}U\times U`$ given by $`\zeta =i_1p_1+i_2p_1+i_2p_2`$. Another isomorphism is $`\tau \zeta `$, where $`\tau `$ is the flip $`i_1p_2+i_2p_1`$. Hence we get equivalences of categories
The composition $`i_2^{}(\zeta ^1\tau \zeta )^{}j`$ defines a self-equivalence of $`U\text{-}\mathrm{𝐌𝐨𝐝}`$ which associates to a given $`U`$-module $`M`$ the new $`U`$-module $`M^{\text{flip}}`$ whose underlying object in C is $`M`$ and such that $`e_{M^{\text{flip}}}=\text{1}_Me_M`$. The same construction applies to $`U`$-bimodules and finally we get equivalences
(2.4.4)
$$\text{}XX^{\text{flip}}$$
for all $`U`$-bimodules $`M`$. If $`X=(0ME\stackrel{\pi }{}U0)`$ is an extension and $`X^{\text{flip}}=(0M^{\text{flip}}E^{\text{flip}}U0)`$, then one verifies that there is a natural isomorphism $`X^{\text{flip}}X`$ of complexes in C inducing $`\text{1}_M`$ on $`M`$, the identity on $`U`$ and carrying the multiplication morphism on $`E^{\text{flip}}`$ to
$$\mu _E+\mathrm{}_E^1(\pi \text{1}_E)+r_E^1(\text{1}_E\pi ):EEE.$$
In terms of the associated rings, this corresponds to replacing the given multiplication $`(x,y)xy`$ of $`E_{}`$ by the new operation $`(x,y)\pi _{}(x)y+\pi _{}(y)xxy`$.
###### Lemma 2.4.5.
If $`M`$ is a $`U`$-bimodule whose left and right actions coincide, then every extension of $`U`$ by $`M`$ splits uniquely.
###### Proof.
Using the idempotent $`e_M`$ we get a $`U`$-linear decomposition $`MM_1M_2`$ where the bimodule structure on $`M_1`$ is given by the zero morphisms and the bimodule structure on $`M_2`$ is given by $`\mathrm{}_M^1`$ and $`r_M^1`$. We have to prove that $`\mathrm{𝐄𝐱𝐦𝐨𝐧}_\text{C}(U,M)`$ is equivalent to a one-point category. By (2.4.3) we can assume that $`M=M_1`$ or $`M=M_2`$. By (2.4.4) we have $`\mathrm{𝐄𝐱𝐦𝐨𝐧}_\text{C}(U,M_2)\mathrm{𝐄𝐱𝐦𝐨𝐧}_\text{C}(U,M_2^{\text{flip}})`$ and on $`M_2^{\text{flip}}`$ the bimodule actions are the zero morphisms. So it is enough to consider $`M=M_1`$. In this case, if $`X=(0MEU0)`$ is any extension, $`\mu _E:EEE`$ factors through a morphism $`UUE`$ and composing with $`u:UUU`$ we get a right inverse of $`EU`$, which shows that $`X`$ is the split extension. Then it is easy to see that $`X`$ does not have any non-trivial automorphisms, which proves the assertion. ∎
###### Proposition 2.4.6.
i) Let $`X=(0IA\stackrel{p}{}A^{}0)`$ be a C-extension and suppose that $`e^{}A_{}^{}`$ is an idempotent element whose left action on the $`A^{}`$-bimodule $`I`$ coincides with its right action. Then there exists a unique idempotent $`eA_{}`$ such that $`p_{}(e)=e^{}`$.
ii) Especially, if $`A^{}`$ is unitary and $`I`$ is a unitary $`A^{}`$-bimodule, then every extension of $`A^{}`$ by $`I`$ is unitary.
###### Proof.
(i) : the hypothesis $`e_{}^{}{}_{}{}^{2}=e^{}`$ implies that $`e^{}:UA^{}`$ is a morphism of (non-unitary) C-monoids. We can then replace $`X`$ by $`Xe^{}`$ and thereby assume that $`A^{}=U`$, $`p:AU`$ and $`I`$ is a (non-unitary) $`U`$-bimodule and the right and left actions on $`I`$ coincide. The assertion to prove is that $`\underset{¯}{1}_U`$ lifts to a unique idempotent $`eA_{}`$. However, this follows easily from lemma 2.4.5. To show (ii), we observe that, by (i), the unit $`\underset{¯}{1}_A^{}`$ of $`A_{}^{}`$ lifts uniquely to an idempotent $`eA_{}`$. We have to show that $`e`$ is a unit for $`A_{}`$. Let us show the left unit property. Via $`e:UA`$ we can view the extension $`X`$ as an exact sequence of left $`U`$-modules. We can then split $`X`$ as the direct sum $`X_1X_2`$ where $`X_1`$ is a sequence of unitary $`U`$-modules and $`X_2`$ is a sequence of $`U`$-modules with trivial actions. But by hypothesis, on $`I`$ and on $`A`$ the $`U`$-module structure is unitary, so $`X=X_1`$ and this is the left unit property. ∎
So much for the general nonsense; we now return to almost algebras. As already announced, from here on, we assume throughout that $`\stackrel{~}{𝔪}`$ is a flat $`V`$-module. As an immediate consequence of proposition 2.4.6 we get natural equivalences of categories
(2.4.7)
whenever $`A_1`$, $`A_2`$ are $`V^a`$-algebras, $`B_i`$ is a $`A_i`$-algebra and $`M_i`$ is a (unitary) $`B_i`$-module, $`i=1,2`$.
Notice that, if $`A=R^a`$ for some $`V`$-algebra $`R`$, $`S`$ (resp. $`J`$) is a $`R`$-algebra (resp. an $`S`$-module) and $`X`$ is any object of $`\mathrm{𝐄𝐱𝐚𝐥}_R(S,J)`$, then by applying termwise the localisation functor we get an object $`X^a`$ of $`\mathrm{𝐄𝐱𝐚𝐥}_A(S^a,J^a)`$. With this notation we have the following lemma.
###### Lemma 2.4.8.
i) Let $`B`$ be any $`A`$-algebra and $`I`$ a $`B`$-module. The natural functor
(2.4.9)
$$\mathrm{𝐄𝐱𝐚𝐥}_{A_{!!}}(B_{!!},I_{})\mathrm{𝐄𝐱𝐚𝐥}_A(B,I)XX^a$$
is an equivalence of categories.
ii) The equivalence (2.4.9) induces a group isomorphism $`\text{Exal}_{A_{!!}}(B_{!!},I_{})\stackrel{}{}\text{Exal}_A(B,I)`$ functorial in all arguments.
###### Proof.
Of course (ii) is an immediate consequence of (i). To show (i), let $`X=(0IEB0)`$ be any object of $`\mathrm{𝐄𝐱𝐚𝐥}_A(B,I)`$. Using corollary 2.2.10 one sees easily that the sequence $`X_!=(0I_!E_{!!}B_{!!}0)`$ is right exact; $`X_!`$ won’t be exact in general, unless $`B`$ (and therefore $`E`$) is an exact algebra. In any case, the kernel of $`I_!E_{!!}`$ is almost zero, so we get an extension of $`B_{!!}`$ by a quotient of $`I_!`$ which maps to $`I_{}`$. In particular we get by pushout an extension $`X_!`$ by $`I_{}`$, i.e. an object of $`\mathrm{𝐄𝐱𝐚𝐥}_{A_{!!}}(B_{!!},I_{})`$ and in fact the assignment $`XX_!`$ is an essential inverse for the functor (2.4.9). ∎
###### Remark 2.4.10.
By inspecting the proof, we see that one can replace $`I_{}`$ by $`I_!=\text{Im}(I_!I_{})`$ in (i) and (ii) above. When $`B`$ is exact, also $`I_!`$ will do.
In (II.1.2) it is shown how to associate to any ring homomorphism $`RS`$ a natural simplicial complex of $`S`$-modules denoted $`𝕃_{S/R}`$ and called the cotangent complex of $`S`$ over $`R`$.
###### Definition 2.4.11.
Let $`AB`$ be a morphism of almost $`V`$-algebras. The almost cotangent complex of $`B`$ over $`A`$ is the simplicial $`B_{!!}`$-module
$$𝕃_{B/A}=B_{!!}_{(V^a\times B)_{!!}}𝕃_{(V^a\times B)_{!!}/(V^a\times A)_{!!}}.$$
Usually we will want to view $`𝕃_{B/A}`$ as an object of the derived category $`\text{D}_{}(s.B_{!!})`$ of simplicial $`B_{!!}`$-modules. Indeed, the hyperext functors computed in this category relate the cotangent complex to a number of important invariants. Recall that, for any simplicial ring $`R`$ and any two $`R`$-modules $`E,F`$ the hyperext of $`E`$ and $`F`$ is the abelian group defined as
$$𝔼\text{xt}_R^p(E,F)=\underset{np}{\text{colim}}\text{Hom}_{\text{D}_{}(R)}(\sigma ^nE,\sigma ^{n+p}F)$$
(where $`\sigma `$ is the suspension functor of (I.3.2.1.4)).
Let us fix an almost algebra $`A`$. First we want to establish the relationship with differentials.
###### Definition 2.4.12.
Let $`B`$ be any $`A`$-algebra, $`M`$ any $`B`$-module.
i) An $`A`$-derivation of $`B`$ with values in $`M`$ is an $`A`$-linear morphism $`:BM`$ such that $`(b_1b_2)=b_1(b_2)+b_2(b_1)`$ for $`b_1,b_2B_{}`$. The set of all $`M`$-valued $`A`$-derivations of $`B`$ forms a $`V`$-module $`\text{Der}_A(B,M)`$ and the almost $`V`$-module $`\text{Der}_A(B,M)^a`$ has a natural structure of $`B`$-module.
ii) We reserve the notation $`I_{B/A}`$ for the ideal $`\text{Ker}(\mu _{B/A}:B_ABB)`$. The module of relative differentials of $`\varphi `$ is defined as the (left) $`B`$-module $`\mathrm{\Omega }_{B/A}=I_{B/A}/I_{B/A}^2`$. It is endowed with a natural $`A`$-derivation $`\delta :B\mathrm{\Omega }_{B/A}`$ defined by $`b\underset{¯}{1}bb\underset{¯}{1}`$ for all $`bB_{}`$. The assignment $`(AB)\mathrm{\Omega }_{B/A}`$ defines a functor
$$\mathrm{\Omega }:V^a\text{-}\mathrm{𝐀𝐥𝐠}.\mathrm{𝐌𝐨𝐫𝐩𝐡}V^a\text{-}\mathrm{𝐀𝐥𝐠}.\mathrm{𝐌𝐨𝐝}$$
from the category of morphisms $`AB`$ of almost $`V`$-algebras to the category $`V^a\text{-}\mathrm{𝐀𝐥𝐠}.\mathrm{𝐌𝐨𝐝}`$ consisting of all pairs $`(B,M)`$ where $`B`$ is an almost $`V`$-algebra and $`M`$ is a $`B`$-module. The morphisms in $`V^a\text{-}\mathrm{𝐀𝐥𝐠}.\mathrm{𝐌𝐨𝐫𝐩𝐡}`$ are the commutative squares; the morphisms $`(B,M)(B^{},M^{})`$ in $`V^a\text{-}\mathrm{𝐀𝐥𝐠}.\mathrm{𝐌𝐨𝐝}`$ are all pairs $`(\varphi ,f)`$ where $`\varphi :BB^{}`$ is a morphism of almost $`V`$-algebras and $`f:B^{}_BMM^{}`$ is a morphism of $`B^{}`$-modules.
The module of relative differentials enjoys the familiar universal properties that one expects. In particular $`\mathrm{\Omega }_{B/A}`$ represents the functor $`\text{Der}_A(B,)`$, i.e. for any (left) $`B`$-module $`M`$ the morphism
(2.4.13)
$$\text{Hom}_B(\mathrm{\Omega }_{B/A},M)\text{Der}_A(B,M)ff\delta $$
is an isomorphism. As an exercise, the reader can supply the proof for this claim and for the following standard proposition.
###### Proposition 2.4.14.
i) Let $`B`$ and $`C`$ be two $`A`$-algebras. Then there is a natural isomorphism:
$$\mathrm{\Omega }_{C_AB/C}C_A\mathrm{\Omega }_{B/A}.$$
ii) Let $`B`$ be an $`A`$-algebra, $`C`$ a $`B`$-algebra. There is a natural exact sequence of $`C`$-modules:
$$C_B\mathrm{\Omega }_{B/A}\mathrm{\Omega }_{C/A}\mathrm{\Omega }_{C/B}0.$$
iii) Let $`I`$ be an ideal of the $`A`$-algebra $`B`$ and let $`C=B/I`$ be the quotient $`A`$-algebra. Then there is a natural exact sequence: $`I/I^2C_B\mathrm{\Omega }_{B/A}\mathrm{\Omega }_{C/A}0`$.
iv) The functor $`\mathrm{\Omega }:V^a\text{-}\mathrm{𝐀𝐥𝐠}.\mathrm{𝐌𝐨𝐫𝐩𝐡}V^a\text{-}\mathrm{𝐀𝐥𝐠}.\mathrm{𝐌𝐨𝐝}`$ commutes with all colimits. ∎
###### Lemma 2.4.15.
For any $`A`$-algebra $`B`$ there is a natural isomorphism of $`B_{!!}`$-modules
$$(\mathrm{\Omega }_{B/A})_!\mathrm{\Omega }_{B_{!!}/A_{!!}}.$$
###### Proof.
Using the adjunction (2.4.13) we are reduced to showing that the natural map
$$\varphi _M:\text{Der}_{A_{!!}}(B_{!!},M)\text{Der}_A(B,M^a)$$
is a bijection for all $`B_{!!}`$-modules $`M`$. Given $`:BM^a`$ we construct $`_!:B_!M_!^aM`$. We extend $`_!`$ to $`VB_!`$ by setting it equal to zero on $`V`$. Then it is easy to check that the resulting map descends to $`B_{!!}`$, hence giving an $`A`$-derivation $`B_{!!}M`$. This procedure yields a right inverse $`\psi _M`$ to $`\varphi _M`$. To show that $`\varphi _M`$ is injective, suppose that $`:B_{!!}M`$ is an almost zero $`A`$-derivation. Composing with the natural $`A`$-linear map $`B_!B_{!!}`$ we obtain an almost zero map $`^{}:B_!M`$. But $`𝔪B_!=B_!`$, hence $`^{}=0`$. This implies that in fact $`=0`$, and the assertion follows. ∎
###### Proposition 2.4.16.
Let $`M`$ be a $`B`$-module. There exists a natural isomorphism of $`B_{!!}`$-modules
$$𝔼\text{xt}_{B_{!!}}^0(𝕃_{B/A},M_!)\text{Der}_A(B,M).$$
###### Proof.
To ease notation, set $`\stackrel{~}{A}=V^a\times A`$ and $`\stackrel{~}{B}=V^a\times B`$. We have natural isomorphisms :
$$\begin{array}{ccc}𝔼\text{xt}_{B_{!!}}^0(𝕃_{B/A},M_!)\hfill & 𝔼\text{xt}_{\stackrel{~}{B}_{!!}}^0(𝕃_{\stackrel{~}{B}_{!!}/\stackrel{~}{A}_{!!}},M_!)\hfill & \text{by }\text{[14]}\text{ (I.3.3.4.4)}\hfill \\ \hfill & \text{Der}_{\stackrel{~}{A}_{!!}}(\stackrel{~}{B}_{!!},M_!)\hfill & \text{by }\text{[14]}\text{ (II.1.2.4.2)}\hfill \\ \hfill & \text{Der}_{\stackrel{~}{A}}(\stackrel{~}{B},M)\hfill & \text{by lemma }\text{2.4.15}\text{.}\hfill \end{array}$$
But it is easy to see that the natural map $`\text{Der}_A(B,M)\text{Der}_{\stackrel{~}{A}}(\stackrel{~}{B},M)`$ is an isomorphism. ∎
###### Theorem 2.4.17.
There is a natural isomorphism
$$\text{Exal}_A(B,M)𝔼\text{xt}_{B_{!!}}^1(𝕃_{B/A},M_!).$$
###### Proof.
With the notation of the proof of proposition 2.4.16 we have natural isomorphisms
$$\begin{array}{ccc}𝔼\text{xt}_{B_{!!}}^1(𝕃_{B/A},M_!)\hfill & 𝔼\text{xt}_{\stackrel{~}{B}_{!!}}^1(𝕃_{\stackrel{~}{B}_{!!}/\stackrel{~}{A}_{!!}},M_!)\hfill & \text{by }\text{[14]}\text{ (I.3.3.4.4)}\hfill \\ \hfill & \text{Exal}_{\stackrel{~}{A}_{!!}}(\stackrel{~}{B}_{!!},M_!)\hfill & \text{by }\text{[14]}\text{ (III.1.2.3)}\hfill \\ \hfill & \text{Exal}_{\stackrel{~}{A}}(\stackrel{~}{B},M)\hfill & \end{array}$$
where the last isomorphism follows directly from lemma 2.4.8(ii) and the subsequent remark 2.4.10. Finally, (2.4.7) shows that $`\text{Exal}_{\stackrel{~}{A}}(\stackrel{~}{B},M)\text{Exal}_A(B,M)`$, as required. ∎
Moreover we have the following transitivity theorem as in (II.2.1.2).
###### Theorem 2.4.18.
Let $`ABC`$ be a sequence of morphisms of almost $`V`$-algebras. There exists a natural distinguished triangle of $`\text{D}_{}(s.C_{!!})`$
where the morphisms $`u`$ and $`v`$ are obtained by functoriality of $`𝕃`$.
###### Proof.
It follows directly from loc. cit. ∎
###### Proposition 2.4.19.
Let $`(A_\lambda B_\lambda )_{\lambda I}`$ be a system of almost $`V`$-algebra morphisms indexed by a small filtered category $`I`$. Then there is a natural isomorphism in $`\text{D}_{}(s.\underset{\lambda I}{\text{colim}}B_{\lambda !!})`$
$$\underset{\lambda I}{\text{colim}}𝕃_{B_\lambda /A_\lambda }𝕃_{\underset{\lambda I}{\text{colim}}B_\lambda /\underset{\lambda I}{\text{colim}}A_\lambda }.$$
###### Proof.
Remark 2.2.15 gives an isomorphism : $`\underset{\lambda I}{\text{colim}}A_{\lambda !!}\stackrel{}{}(\underset{\lambda I}{\text{colim}}A_\lambda )_{!!}`$ (and likewise for $`\underset{\lambda I}{\text{colim}}B_\lambda `$). Then the claim follows from (II.1.2.3.4). ∎
Next we want to prove the almost version of the flat base change theorem (II.2.2.1). To this purpose we need some preparation.
###### Proposition 2.4.20.
Let $`B`$ and $`C`$ be two $`A`$-algebras and set $`T_i=\text{Tor}_i^{A_{!!}}(B_{!!},C_{!!})`$. If $`A`$, $`B`$, $`C`$ and $`B_AC`$ are all exact, then for every $`i>0`$ the natural morphism $`\stackrel{~}{𝔪}_VT_iT_i`$ is an isomorphism.
###### Proof.
For any almost $`V`$-algebra $`D`$ we let $`k_D`$ denote the complex of $`D_{!!}`$-modules $`[\stackrel{~}{𝔪}_VD_{!!}D_{!!}]`$ placed in degrees $`1,0`$; we have a distiguished triangle
$$\text{T}(D):\text{}$$
By the assumption, the natural map $`k_Ak_B`$ is a quasi-isomorphism and $`\stackrel{~}{𝔪}_VB_{!!}B_!`$. On the other hand, for all $`i`$ we have
$$\text{Tor}_i^{A_{!!}}(k_B,C_{!!})\text{Tor}_i^{A_{!!}}(k_A,C_{!!})H^i(k_A_{A_{!!}}C_{!!})=H^i(k_C).$$
In particular $`\text{Tor}_i^{A_{!!}}(k_B,C_{!!})=0`$ for all $`i>1`$. As $`\stackrel{~}{𝔪}`$ is flat over $`V`$, we have $`\stackrel{~}{𝔪}_VT_i\text{Tor}_i^{A_{!!}}(\stackrel{~}{𝔪}_VB_{!!},C_{!!})`$. Then by the long exact Tor sequence associated to $`\text{T}(B)\stackrel{\text{L}}{}_{A_{!!}}C_{!!}`$ we get the assertion for all $`i>1`$. Next we consider the natural map of distinguished triangles $`\text{T}(A)\stackrel{\text{L}}{}_{A_{!!}}A_{!!}\text{T}(B)\stackrel{\text{L}}{}_{A_{!!}}C_{!!}`$; writing down the associated morphism of long exact Tor sequences, we obtain a diagram with exact rows :
By the above, the leftmost vertical map is an isomorphism; moreover, the assumption gives $`\text{Ker}(i)\text{Ker}(\stackrel{~}{𝔪}V)\text{Ker}(i^{})`$. Then, since $``$ is injective, also $`^{}`$ must be injective, which implies our assertion for the remaining case $`i=1`$. ∎
###### Corollary 2.4.21.
Keep the notation of proposition 2.4.20 and suppose that $`\text{Tor}_i^A(B,C)0`$ for some $`i>0`$. Then the corresponding $`T_i`$ vanishes. ∎
###### Theorem 2.4.22.
Let $`B`$, $`A^{}`$ be two $`A`$-algebras. Suppose that the natural morphism $`B\stackrel{\text{L}}{}_AA^{}B^{}=B_AA^{}`$ is an isomorphism in $`\text{D}_{}(s.A)`$. Then the natural morphisms
$$\begin{array}{c}B_{!!}^{}_{B_{!!}}𝕃_{B/A}𝕃_{B^{}/A^{}}\hfill \\ (B_{!!}^{}_{B_{!!}}𝕃_{B/A})(B_{!!}^{}_{A_{!!}^{}}𝕃_{A^{}/A})𝕃_{B^{}/A}\hfill \end{array}$$
are quasi-isomorphisms.
###### Proof.
Let us remark that the functor $`DV^a\times D`$ : $`A\text{-}\mathrm{𝐀𝐥𝐠}(V^a\times A)\text{-}\mathrm{𝐀𝐥𝐠}`$ commutes with tensor products; hence the same holds for the functor $`D(V^a\times D)_{!!}`$ (see remark 2.2.15). Then, in view of corollary 2.4.21, the theorem is reduced immediately to (II.2.2.1). ∎
As an application we obtain the vanishing of the almost cotangent complex for a certain class of morphisms.
###### Theorem 2.4.23.
Let $`RS`$ be a morphism of almost algebras such that
$$\text{Tor}_i^R(S,S)0\text{Tor}_i^{S_RS}(S,S)\text{for all }i>0$$
(for the natural $`S_RS`$-module structure induced by $`\mu _{S/R}`$). Then $`𝕃_{S/R}0`$ in $`\text{D}_{}(S_{!!})`$.
###### Proof.
Since $`\text{Tor}_i^R(S,S)=0`$ for all $`i>0`$, theorem 2.4.22 applies (with $`A=R`$ and $`B=A^{}=S`$), giving the natural isomorphisms
(2.4.24)
$$\begin{array}{c}(S_RS)_{!!}_{S_{!!}}𝕃_{S/R}𝕃_{S_RS/S}\hfill \\ ((S_RS)_{!!}_{S_{!!}}𝕃_{S/R})((S_RS)_{!!}_{S_{!!}}𝕃_{S/R})𝕃_{S_RS/R}\hfill \end{array}$$
Since $`\text{Tor}_i^{S_RS}(S,S)=0`$, the same theorem also applies with $`A=S_RS`$, $`B=S`$, $`A^{}=S`$, and we notice that in this case $`B^{}S`$; hence we have
(2.4.25)
$$𝕃_{S/S_RS}S_{!!}_{S_{!!}}𝕃_{S/S_RS}𝕃_{S/S}0.$$
Next we apply transitivity to the sequence $`RS_RSS`$, to obtain (thanks to (2.4.25))
(2.4.26)
$$S_{!!}_{S_RS_{!!}}𝕃_{S_RS/R}𝕃_{S/R}.$$
Applying $`S_{!!}_{S_RS_{!!}}`$ to the second isomorphism (2.4.24) we obtain
(2.4.27)
$$𝕃_{S/R}𝕃_{S/R}S_{!!}_{S_RS_{!!}}𝕃_{S_RS/R}.$$
Finally, composing (2.4.26) and (2.4.27) we derive
(2.4.28)
$$𝕃_{S/R}𝕃_{S/R}\stackrel{}{}𝕃_{S/R}.$$
However, by inspection, the isomorphism (2.4.28) is the sum map. Consequently $`𝕃_{S/R}0`$, as claimed. ∎
Finally we have a fundamental spectral sequence as in (III.3.3.2).
###### Theorem 2.4.29.
Let $`\varphi :AB`$ be a morphism of almost algebras such that $`B_ABB`$ (e.g. such that $`B`$ is a quotient of $`A`$). Then there is a first quadrant homology spectral sequence of bigraded almost algebras
$$E_{pq}^2=H_{p+q}(\text{Sym}_B^q(𝕃_{B/A}^a))\text{Tor}_{p+q}^A(B,B).$$
###### Proof.
We replace $`\varphi `$ by $`\text{1}_{V^a}\times \varphi `$ and apply the functor $`BB_{!!}`$ (which commutes with tensor products by remark 2.2.15) thereby reducing the assertion to the above mentioned (III.3.3.2). ∎
## 3. Almost ring theory
### 3.1. Flat, unramified and étale morphisms
Let $`AB`$ be a morphism of almost $`V`$-algebras. Using the natural “multiplication” morphism of $`A`$-algebras $`\mu _{B/A}:B_ABB`$ we can view $`B`$ as a $`B_AB`$-algebra.
###### Definition 3.1.1.
Let $`\varphi :AB`$ be a morphism of almost $`V`$-algebras.
i) We say that $`\varphi `$ is a flat (resp. faithfully flat, resp. almost projective) morphism if $`B`$ is a flat (resp. faithfully flat, resp. almost projective) $`A`$-module.
ii) We say that $`\varphi `$ is almost finite (resp. finite) if $`B`$ is an almost finitely generated (resp. finitely generated) $`A`$-module.
iii) We say that $`\varphi `$ is weakly unramified (resp. unramified) if $`B`$ is a flat (resp. almost projective) $`B_AB`$-module (via the morphism $`\mu _{B/A}`$ defined above).
iv) $`\varphi `$ is weakly étale (resp. étale) if it is flat and weakly unramified (resp. unramified).
###### Lemma 3.1.2.
Let $`\varphi :AB`$ and $`\psi :BC`$ be morphisms of almost $`V`$-algebras.
i) Any base change of a flat (resp. almost projective, resp. faithfully flat, resp. almost finite, resp. weakly unramified, resp. unramified, resp. weakly étale, resp. étale) morphism is flat (resp. almost projective, resp. faithfully flat, resp. almost finite, resp. weakly unramified, resp. unramified, resp. weakly étale, resp. étale);
ii) if both $`\varphi `$ and $`\psi `$ are flat (resp. almost projective, resp. faithfully flat, resp. almost finite, resp. weakly unramified, resp. unramified, resp. weakly étale, resp. étale), then so is $`\psi \varphi `$;
iii) if $`\varphi `$ is flat and $`\psi \varphi `$ is faithfully flat, then $`\varphi `$ is faithfully flat;
iv) if $`\varphi `$ is weakly unramified and $`\psi \varphi `$ is flat (resp. weakly étale), then $`\psi `$ is flat (resp. weakly étale);
v) If $`\varphi `$ is unramified and $`\psi \varphi `$ is étale, then $`\psi `$ is étale;
vi) $`\varphi `$ is faithfully flat if and only if it is a monomorphism and $`B/A`$ is a flat $`A`$-module;
vii) if $`\varphi `$ is almost finite and weakly unramified, then $`\varphi `$ is unramified.
###### Proof.
For (vi) use the Tor sequences. In view of proposition 2.3.15(ii), to show (vii) it suffices to know that $`B`$ is an almost finitely presented $`B_AB`$-module; but this follows from the existence of an epimorphism of $`B_AB`$-modules $`(B_AB)_AB\text{Ker}(\mu _{B/A})`$ defined by $`xbx(\underset{¯}{1}bb\underset{¯}{1})`$. Of the remaining assertions, only (iv) and (v) are not obvious, but the proof is just the “almost version” of a well-known argument. Let us show (v); the same argument applies to (iv). We remark that $`\mu _{B/A}`$ is an étale morphism, since $`\varphi `$ is unramified. Define $`\mathrm{\Gamma }_\psi =\text{1}_C_B\mu _{B/A}`$. By (i), $`\mathrm{\Gamma }_\psi `$ is étale. Define also $`p=(\psi \varphi )_A\text{1}_B`$. By (i), $`p`$ is flat (resp. étale). The claim follows by remarking that $`\psi =\mathrm{\Gamma }_\psi p`$ and applying (ii). ∎
###### Remark 3.1.3.
i) Suppose we work in the classical limit case, that is, $`V=𝔪`$ (cp. example 2.1.1(ii)). Then we caution the reader that our notion of “étale morphism” is more general than the usual one, as defined in . The relationship between the usual notion and ours is discussed in the digression at the end of section 3.4.
ii) The naive hope that the functor $`AA_{!!}`$ might preserve flatness is crushed by the following counterexample. Let $`(V,𝔪)`$ be as in example 2.1.1(i) and let $`k`$ be the residue field of $`V`$. Consider the flat map $`V\times VV`$ defined as $`(x,y)x`$. We get a flat morphism $`V^a\times V^aV^a`$ in $`V^a\text{-}\mathrm{𝐀𝐥𝐠}`$; applying the left adjoint to localisation yields a map $`V\times _kVV`$ that is not flat. On the other hand, faithful flatness is preserved. Indeed, let $`\varphi :AB`$ be a morphism of almost algebras. Then $`\varphi `$ is a monomorphism if and only if $`\varphi _{!!}`$ is injective; moreover, $`B_{!!}/\text{Im}(A_{!!})B_!/A_!`$, which is flat over $`A_{!!}`$ if and only if $`B/A`$ is flat over $`A`$, by proposition 2.3.27.
We will find useful to study certain “almost idempotents”, as in the following proposition.
###### Proposition 3.1.4.
A morphism $`\varphi :AB`$ is unramified if and only if there exists an almost element $`e_{B/A}B_AB_{}`$ such that
i) $`e_{B/A}^2=e_{B/A}`$;
ii) $`\mu _{B/A}(e_{B/A})=\underset{¯}{1}`$;
iii) $`xe_{B/A}=0`$ for all $`xI_{B/A}`$.
###### Proof.
Suppose that $`\varphi `$ is unramified. We start by showing that for every $`\epsilon 𝔪`$ there exist almost elements $`e_\epsilon `$ of $`B_AB`$ such that
(3.1.5)
$$e_\epsilon ^2=\epsilon e_\epsilon \mu _{B/A}(e_\epsilon )=\epsilon \underset{¯}{1}I_{B/A}e_\epsilon =0.$$
Since $`B`$ is an almost projective $`B_AB`$-module, for every $`\epsilon 𝔪`$ there exists an “approximate splitting” for the epimorphism $`\mu _{B/A}:B_ABB`$, i.e. a $`B_AB`$-linear morphism $`u_\epsilon :BB_AB`$ such that $`\mu _{B/A}u_\epsilon =\epsilon \text{1}_B`$. Set $`e_\epsilon =u_\epsilon \underset{¯}{1}:AB_AB`$. We see that $`\mu _{B/A}(e_\epsilon )=\epsilon \underset{¯}{1}`$. To show that $`e_\epsilon ^2=\epsilon e_\epsilon `$ we use the $`B_AB`$-linearity of $`u_\epsilon `$ to compute
$$e_\epsilon ^2=e_\epsilon u_\epsilon (\underset{¯}{1})=u_\epsilon (\mu _{B/A}(e_\epsilon )\underset{¯}{1})=u_\epsilon (\mu _{B/A}(e_\epsilon ))=\epsilon e_\epsilon .$$
Next take any almost element $`x`$ of $`I_{B/A}`$ and compute
$$xe_\epsilon =xu_\epsilon (\underset{¯}{1})=u_\epsilon (\mu _{B/A}(x)\underset{¯}{1})=0.$$
This establishes (3.1.5). Next let us take any other $`\delta 𝔪`$ and a corresponding almost element $`e_\delta `$. Both $`\epsilon \underset{¯}{1}e_\epsilon `$ and $`\delta \underset{¯}{1}e_\delta `$ are elements of $`I_{B/A}`$, hence we have $`(\delta \underset{¯}{1}e_\delta )e_\epsilon =0=(\epsilon \underset{¯}{1}e_\epsilon )e_\delta `$ which implies
(3.1.6)
$$\delta e_\epsilon =\epsilon e_\delta \text{for all }\epsilon ,\delta 𝔪\text{.}$$
Let us define a map $`e_{B/A}:𝔪_V𝔪B_AB_{}`$ by the rule
(3.1.7)
$$\epsilon \delta \delta e_\epsilon \text{for all }\epsilon ,\delta 𝔪.$$
To show that (3.1.7) does indeed determine a well defined morphism, we need to check that $`\delta ve_\epsilon =\delta e_{v\epsilon }`$ and $`\delta e_{\epsilon +\epsilon ^{}}=\delta (e_\epsilon +e_\epsilon ^{})`$ for all $`\epsilon ,\epsilon ^{},\delta 𝔪`$ and all $`vV`$. However, both identities follow easily by a repeated application of (3.1.6). It is easy to see that $`e_{B/A}`$ defines an almost element with the required properties.
Conversely, suppose an almost element $`e_{B/A}`$ of $`B_AB`$ is given with the stated properties. We define $`u:BB_AB`$ by $`be_{B/A}(1b)`$ ($`bB_{}`$) and $`v=\mu _{B/A}`$. Then (iii) says that $`u`$ is a $`B_AB`$-linear morphism and (ii) shows that $`vu=\text{1}_B`$. Hence, by lemma 2.3.12, $`\varphi `$ is unramified. ∎
###### Corollary 3.1.8.
Under the hypotheses and notation of the proposition, the ideal $`I=I_{B/A}`$ has a natural structure of $`A`$-algebra, with unit morphism given by $`\underset{¯}{1}_{I/A}=\underset{¯}{1}_{B_AB/A}e_{B/A}`$ and whose multiplication is the restriction of $`\mu _{B_AB/A}`$ to $`I`$. Moreover the natural morphism
$$B_ABI_{B/A}Bx(x\underset{¯}{1}_{I/A}\mu _{B/A}(x))$$
is an isomorphism of $`A`$-algebras.
###### Proof.
Left to the reader as an exercise. ∎
### 3.2. Almost traces
Let $`A`$ be an almost $`V`$-algebra. For any integer $`n>0`$, the standard direct sum decomposition of $`A^n`$ determines uniquely $`A`$-linear morphisms $`A\stackrel{e_i^A}{}A^n\stackrel{\pi _j^A}{}A`$ (for $`i,j=1,\mathrm{},n`$) such that $`\pi _j^Ae_i^A=\delta _{ij}\text{1}_A`$ for all $`i,j`$ and $`_{i=1}^ne_i^A\pi _i^A=\text{1}_{A^n}`$. We can then define a natural trace homomorphism
(3.2.1)
$$\text{Tr}:\text{alHom}_A(A^n,A^n)A\varphi \underset{i=1}{\overset{n}{}}\pi _i^A\varphi e_i^A$$
which is an $`A`$-linear morphism. For any $`\varphi ,\psi \text{alHom}_A(A^n,A^n)_{}`$ we have $`\text{Tr}(\varphi \psi )=\text{Tr}(\psi \varphi )`$. It follows easily that Tr is independent of the given direct sum decomposition of $`A^n`$. More generally, suppose that $`M`$ is an almost projective almost finitely generated $`A`$-module. Then for any $`\epsilon 𝔪`$ we can find $`n=n(\epsilon )`$ and morphisms
(3.2.2)
such that $`v_\epsilon u_\epsilon =\epsilon \text{1}_M`$. Let $`E(M)=\text{alHom}_A(M,M)`$; notice that $`E(M)_{}`$ is naturally isomorphic to $`\text{Hom}_A(M,M)`$. We consider the $`A`$-linear morphism
(3.2.3)
$$t_\epsilon :E(M)A\varphi \text{Tr}(u_\epsilon \varphi v_\epsilon )\text{(}\varphi E(M)_{}\text{).}$$
Now, pick any other $`\delta 𝔪`$. We compute
$$\begin{array}{cc}\hfill \epsilon t_\delta (\varphi )=& \epsilon \text{Tr}(u_\delta \varphi v_\delta )=\text{Tr}(u_\delta v_\epsilon u_\epsilon \varphi v_\delta )\hfill \\ \hfill =& \text{Tr}(u_\epsilon \varphi v_\delta u_\delta v_\epsilon )=\delta \text{Tr}(u_\epsilon \varphi v_\epsilon )=\delta t_\epsilon (\varphi ).\hfill \end{array}$$
This allows us to define a map
$$t_{M/A}:𝔪_V𝔪_VE(M)_{}A_{}$$
by setting $`\epsilon \delta \varphi \epsilon t_\delta (\varphi )`$. We leave to the reader the verification that $`t_{M/A}`$ is well defined and does not depend on the choice of $`t_\delta `$. It induces a morphism $`E(M)A`$ that we denote again by $`t_{M/A}`$ and we call the almost trace morphism for the almost $`A`$-module $`M`$.
Let $`fM_{}^{}`$, $`mM_{}`$ and define $`\varphi _{f,m}E(M)_{}`$ by the formula $`\varphi _{f,m}(m^{})=f(m^{})m`$ for all $`m^{}M_{}`$. We have the following :
###### Lemma 3.2.4.
With the above notation : $`t_{M/A}(\varphi _{f,m})=f(m)`$.
###### Proof.
Let $`f:MA`$ and $`m:AM`$ be given. Obviously we have $`\varphi _{f,m}=mf`$ and $`f(m)=fm`$. Pick morphisms $`u_\epsilon `$ and $`v_\epsilon `$ as in (3.2.2). Using the foregoing notation, we can write :
$$\begin{array}{cc}\hfill t_\epsilon (\varphi _{f,m})=& _{i=1}^n(\pi _i^Au_\epsilon m)(fv_\epsilon e_i^A)\hfill \\ \hfill =& _{i=1}^n(fv_\epsilon e_i^A)(\pi _i^Au_\epsilon m)\hfill \\ \hfill =& fv_\epsilon u_\epsilon m=\epsilon fm\hfill \end{array}$$
from which the claim follows directly. ∎
###### Lemma 3.2.5.
Let $`M`$ and $`N`$ be almost finitely generated almost projective $`A`$-modules, and $`\varphi :MN`$, $`\psi :NM`$ two $`A`$-linear morphisms. Then :
i) $`t_{M/A}(\psi \varphi )=t_{N/A}(\varphi \psi )`$.
ii) If $`\psi \varphi =a\text{1}_M`$ and $`\varphi \psi =a\text{1}_N`$ for some $`aA_{}`$, and if, furthermore, there exist $`u\text{End}_A(M)`$, $`v\text{End}_A(N)`$ such that $`v\varphi =\varphi u`$, then $`a(t_{M/A}(u)t_{N/A}(v))=0`$.
###### Proof.
(i) is left to the reader as an exercise. For (ii) we compute using (i) : $`at_{M/A}(u)=t_{M/A}(\psi \varphi u)=t_{M/A}(\psi v\varphi )=t_{N/A}(v\varphi \psi )=at_{N/A}(v)`$. ∎
###### Proposition 3.2.6.
Let $`\underset{¯}{M}=(0M_1\stackrel{i}{}M_2\stackrel{p}{}M_30)`$ be an exact sequence of almost finitely generated almost projective $`A`$-modules, and let $`\underset{¯}{u}=(u_1,u_2,u_3):\underset{¯}{M}\underset{¯}{M}`$ be an endomorphism of $`\underset{¯}{M}`$, given by endomorphisms $`u_i:M_iM_i`$ ($`i=1,2,3`$). Then we have $`t_{M_2/A}(u_2)=t_{M_1/A}(u_1)+t_{M_3/A}(u_3)`$.
###### Proof.
Suppose first that there exists a splitting $`s:M_3M_2`$ for $`p`$, so that we can view $`u_2`$ as a matrix $`\left(\begin{array}{cc}u_1\hfill & v\hfill \\ 0\hfill & u_3\hfill \end{array}\right)`$, where $`v\text{Hom}_A(M_3,M_1)`$. By additivity of the trace, we are then reduced to show that $`t_{M_2/A}(ivp)=0`$. By lemma 3.2.5(i), this is the same as $`t_{M_3/A}(piv)`$, which obviously vanishes. In general, for any $`a𝔪`$ we consider the morphism $`\mu _a=a\text{1}_{M_3}`$ and the pull back morphism $`\underset{¯}{M}\mu _a\underset{¯}{M}`$ :
Then $`\underset{¯}{M}\mu _a`$ is a split exact sequence with the endomorphism $`\underset{¯}{u}\mu _a=(u_1,v,u_3)`$, for a certain $`v\text{End}_A(P)`$. The pair of morphisms $`(a\text{1}_{M_2},p)`$ induces a morphism $`\psi :M_2P`$, and it is easy to check that $`\varphi \psi =a\text{1}_{M_2}`$ and $`\psi \varphi =a\text{1}_P`$. We can therefore apply lemma 3.2.5 to deduce that $`a(t_{P/A}(v)t_{M/A}(u))=0`$. By the foregoing we know that $`t_{P/A}(v)=t_{M_1/A}(u_1)+t_{M_3/A}(u_3)`$, so the claim follows. ∎
Suppose now that $`B`$ is an almost finitely generated almost projective $`A`$-algebra. For any $`bB_{}`$, denote by $`\mu _b:BB`$ the $`B`$-linear morphism $`b^{}bb^{}`$. The map $`b\mu _b`$ defines a $`B`$-linear monomorphism $`\mu :BE(B)`$. The composition
$$\text{Tr}_{B/A}=t_{B/A}\mu :BA$$
will also be called the almost trace morphism of the $`A`$-algebra $`B`$.
###### Proposition 3.2.7.
Let $`A`$ and $`B`$ be as in the above discussion.
i) If $`\varphi :AB`$ is an isomorphism, then $`\text{Tr}_{B/A}=\varphi ^1`$.
ii) If $`C`$ any other $`A`$-algebra, then $`\text{Tr}_{C_AB/C}=\text{1}_C_A\text{Tr}_{B/A}`$.
iii) If $`C`$ is an almost projective almost finite $`B`$-algebra, then $`\text{Tr}_{C/A}=\text{Tr}_{B/A}\text{Tr}_{C/B}`$.
###### Proof.
(i) and (ii) are left as exercises for the reader. We verify (iii). For given $`\epsilon ,\delta 𝔪`$ pick morphisms $`B\stackrel{u_\epsilon }{}A^n\stackrel{v_\epsilon }{}B`$ and $`C\stackrel{u_\delta ^{}}{}B^m\stackrel{v_\delta ^{}}{}C`$ such that $`v_\epsilon u_\epsilon =\epsilon \text{1}_B`$ and $`v_\delta ^{}u_\delta ^{}=\delta \text{1}_C`$. If we set $`u_\epsilon ^m=u_\epsilon _A\text{1}_{A^m}`$, $`u_{\delta \epsilon }^{\prime \prime }=u_\epsilon ^mu_\delta ^{}:CA^n_AA^m`$, $`v_\epsilon ^m=v_\epsilon _A\text{1}_{A^m}`$ and $`v_{\delta \epsilon }^{\prime \prime }=v_\delta ^{}v_\epsilon ^m:A^n_AA^mC`$ then we have $`v_{\delta \epsilon }^{\prime \prime }u_{\delta \epsilon }^{\prime \prime }=\epsilon \delta \text{1}_C`$. Define
$$\begin{array}{cc}\hfill t_{\epsilon ,B/A}:& BAb\text{Tr}(u_\epsilon \mu _bv_\epsilon )\hfill \\ \hfill t_{\delta ,C/B}:& CBc\text{Tr}(u_\delta ^{}\mu _cv_\delta ^{})\hfill \\ \hfill t_{\delta \epsilon ,C/A}:& CAc\text{Tr}(u_{\delta \epsilon }^{\prime \prime }\mu _cv_{\delta \epsilon }^{\prime \prime }).\hfill \end{array}$$
Using (3.2.1) we can write
$$\begin{array}{cc}\hfill t_{\delta \epsilon ,C/A}(c)=& \underset{i,j=1}{\overset{n,m}{}}(\pi _i^A_A\pi _j^A)u_{\delta \epsilon }^{\prime \prime }\mu _cv_{\delta \epsilon }^{\prime \prime }(e_i^A_Ae_j^A)\hfill \\ \hfill =& \underset{i,j=1}{\overset{n,m}{}}(\pi _i^A\pi _j^A)u_\epsilon ^mu_\delta ^{}\mu _cv_\delta ^{}v_\epsilon ^m(e_i^Ae_j^A)\hfill \\ \hfill =& \underset{i,j=1}{\overset{n,m}{}}\pi _i^Au_\epsilon \pi _j^Bu_\delta ^{}\mu _cv_\delta ^{}e_j^Bv_\epsilon e_i^A\hfill \\ \hfill =& \underset{i=1}{\overset{n}{}}\pi _i^Au_\epsilon t_{\delta ,C/B}(c)v_\epsilon e_i^A\hfill \\ \hfill =& t_{\epsilon ,B/A}t_{\delta ,C/B}(c)\hfill \end{array}$$
which implies immediately the claim. ∎
###### Corollary 3.2.8.
Let $`AB`$ be a faithfully flat almost finitely presented and étale morphism of almost $`V`$-algebras. Then $`\text{Tr}_{B/A}:BA`$ is an epimorphism.
###### Proof.
Under the stated hypotheses, $`B`$ is an almost projective $`A`$-module (by proposition 2.3.15). Let $`C=\text{Coker}(\text{Tr}_{B/A})`$ and $`\text{Tr}_{B/B_AB}`$ the trace morphism for the morphism of almost $`V`$-algebras $`\mu _{B/A}`$. By faithful flatness, the natural morphism $`CC_AB=\text{Coker}(\text{Tr}_{B_AB/B})`$ is a monomorphism, hence it suffices to show that $`\text{Tr}_{B_AB/B}`$ is an epimorphism (here $`B_AB`$ is considered as a $`B`$-algebra via the second factor). However, from proposition 3.2.7(i) and (iii) we see that $`\text{Tr}_{B/B_AB}`$ is a right inverse for $`\text{Tr}_{B_AB/B}`$. The claim follows. ∎
It is useful to introduce the $`A`$-linear morphism
$$\text{tr}_{B/A}=\text{Tr}_{B/A}\mu _{B/A}:B_ABA.$$
We can view $`\text{tr}_{B/A}`$ as a bilinear form; it induces an $`A`$-linear morphism
$$\tau _{B/A}:BB^{}=\text{alHom}_A(B,A)$$
characterized by the equality $`\text{tr}_{B/A}(b_1b_2)=\tau _{B/A}(b_1)(b_2)`$ for all $`b_1,b_2B_{}`$. We say that $`\text{tr}_{B/A}`$ is a perfect pairing if $`\tau _{B/A}`$ is an isomorphism.
###### Theorem 3.2.9.
An almost projective and almost finite morphism $`\varphi :AB`$ of almost $`V`$-algebras is étale if and only if the trace form $`\text{tr}_{B/A}`$ is a perfect pairing.
###### Proof.
Suppose that $`\varphi `$ is étale. Let $`e_{B/A}`$ be the idempotent almost element of $`B_AB`$ provided by proposition 3.1.4. We define a morphism $`\sigma :B^{}B`$ by $`f(f_A\text{1}_B)(e_{B/A})`$. To start with, we remark that both $`\tau _{B/A}`$ and $`\sigma `$ are $`B`$-linear morphisms (for the natural $`B`$-module structure of $`B^{}`$ defined in remark 2.3.18). Indeed, let us pick any $`b,b^{},b^{\prime \prime }B_{}`$, $`fB_{}^{}`$ and compute directly
$$\begin{array}{cc}\hfill (b\tau _{B/A}(b^{}))(b^{\prime \prime })=& \tau _{B/A}(b^{})(bb^{\prime \prime })=\text{Tr}_{B/A}(bb^{}b^{\prime \prime })=(\tau _{B/A}(bb^{}))(b^{\prime \prime }).\hfill \\ \hfill b\sigma (f)=& b(f_A\text{1}_B)(e_{B/A})=(f_A\text{1}_B)((\underset{¯}{1}_{B/A}b)e_{B/A})\hfill \\ \hfill =& (f_A\text{1}_B)((b\underset{¯}{1}_{B/A})e_{B/A})=((bf)_A\text{1}_B)(e_{B/A})\hfill \\ \hfill =& \sigma (bf).\hfill \end{array}$$
Next we show that $`\sigma `$ is a left inverse for $`\tau _{B/A}`$. In fact, let $`bB_{}`$. We have
$$\begin{array}{cc}\hfill \sigma \tau _{B/A}(b)=& (\tau _{B/A}(b)_A\text{1}_B)(e_{B/A})=(\text{Tr}_{B/A}_A\text{1}_B)((b\underset{¯}{1}_{B/A})e_{B/A})\hfill \\ \hfill =& \text{Tr}_{B_AB/B}((\underset{¯}{1}_{B/A}b)e_{B/A})=b\text{Tr}_{B_AB/B}(e_{B/A}).\hfill \end{array}$$
Therefore it suffices to show that $`\text{Tr}_{B_AB/B}(e_{B/A})=\underset{¯}{1}`$. However, by hypothesis $`\varphi `$ is unramified, hence corollary 3.1.8 gives a decomposition $`B_ABBI_{B/A}`$ such that $`e_{B/A}`$ acts as the identity on the first factor and as the zero morphism on the second factor.
Now, let $`X=\text{Ker}(\sigma )`$. From the above we derive a $`B`$-linear isomorphism $`B^{}BX`$. We dualize and apply lemma 2.3.23(ii) to obtain another $`B`$-linear isomorphism
(3.2.10)
$$B(B^{})^{}(BX)^{}B^{}X^{}BXX^{}.$$
Finally, composing the isomorphism (3.2.10) with the projection on the first factor, we get a split $`B`$-linear epimorphism $`BB`$, hence a surjective $`B_{}`$-linear morphism $`B_{}B_{}`$. Such a morphism is necessarily an isomorphism, and, tracing backward, the same must hold for $`\tau _{B/A}`$.
Conversely, suppose that the trace form is a perfect pairing. By lemma 2.3.23(i) the natural morphism $`\alpha :B^{}_AB\text{alHom}_B(B_AB,B)`$ is an isomorphism and one verifies easily that $`\alpha (\tau _{B/A}_A\text{1}_B)=\tau _{B_AB/B}`$. In particular $`\tau _{B_AB/B}`$ is also an isomorphism. The multiplication gives an almost element $`\mu _{B/A}\text{alHom}_B(B_AB,B)_{}`$; let $`e=\tau _{B_AB/B}^1(\mu _{B/A})`$. We derive
(3.2.11)
$$\text{Tr}_{B_AB/B}(e)=\tau _{B_AB/B}(e)(\underset{¯}{1}_{B_AB})=\mu _{B/A}(\underset{¯}{1}_{B_AB})=\underset{¯}{1}_{B/A}.$$
Furthermore, we have already remarked that $`\tau _{B/A}`$ is a $`B`$-linear morphism, hence $`\tau _{B_AB/B}`$ is a $`B_AB`$-linear morphism. Consequently, for any almost element $`x`$ of $`B_AB`$ we have
$$\tau _{B_AB/B}(xe)=x\tau _{B_AB/B}(e)=x\mu _{B/A}=\mu _{B/A}(x)\mu _{B/A}=\mu _{B/A}(x)\tau _{B_AB/B}(e).$$
Since by hypothesis $`\tau _{B/A}`$ is an isomorphism, this implies
(3.2.12)
$$xe=\mu _{B/A}(x)e.$$
Consider the morphism $`\mu _e:B_ABB_AB`$ defined by $`xex`$; then $`\mu _e`$ is $`B`$-linear (for the $`B`$-module structure defined by the second factor). Applying (3.2.12) and lemma 3.2.4 we conclude that $`t_{B_AB/B}(\mu _e)=\mu _{B/A}(e)`$. On the other hand, (3.2.11) says that this trace is equal to $`\underset{¯}{1}_{B/A}`$, hence
(3.2.13)
$$\mu _{B/A}(e)=\underset{¯}{1}_{B/A}.$$
Let $`\beta :BB_AB`$ be defined as $`bbe.`$ From (3.2.12) we see that both $`\beta `$ and $`\mu _{B/A}`$ are $`B_AB`$-linear morphisms and from (3.2.13) we know moreover that $`\mu _{B/A}\beta =\text{1}_B`$. By lemma 2.3.12 we deduce that $`B`$ is an almost projective $`B_AB`$-module, i.e. $`\varphi `$ is unramified, as claimed. ∎
###### Definition 3.2.14.
The nilradical of an almost algebra $`A`$ is the ideal $`\text{nil}(A)=\text{nil}(A_{})^a`$ (where, for a ring $`R`$, we denote by $`\text{nil}(R)`$ the ideal of nilpotent elements in $`R`$). We say that $`A`$ is reduced if $`\text{nil}(A)0`$.
Notice that, if $`R`$ is a $`V`$-algebra, then every nilpotent ideal in $`R^a`$ is of the form $`I^a`$, where $`I`$ is a nilpotent ideal in $`R`$ (indeed, it is of the form $`I^a`$ where $`I`$ is an ideal, and $`𝔪I`$ is seen to be nilpotent). It follows easily that $`\text{nil}(A)`$ is the colimit of the nilpotent ideals in $`A`$; moreover $`\text{nil}(R)^a=\text{nil}(R^a)`$. Using this one sees that $`A/\text{nil}(A)`$ is reduced.
###### Proposition 3.2.15.
Let $`AB`$ be an étale almost finitely presented morphism of almost algebras. If $`A`$ is reduced then $`B`$ is reduced as well.
###### Proof.
Under the stated hypothesis, $`B`$ is an almost projective $`A`$-module (by virtue of proposition 2.3.15(ii)). Hence, for given $`\epsilon 𝔪`$, pick a sequence of morphisms $`B\stackrel{u_\epsilon }{}A^n\stackrel{v_\epsilon }{}B`$ such that $`v_\epsilon u_\epsilon =\epsilon \text{1}_B`$; with the notation of (3.2.3), define $`\nu _b:A^nA^n`$ by $`\nu _b=v_\epsilon \mu _bu_\epsilon `$, so that $`t_\epsilon (b)=\text{Tr}(\nu _b)`$. One verifies easily that $`\nu _b^m=\epsilon ^{m1}\nu _{b^m}`$ for all integers $`m>0`$.
Now, suppose that $`b\text{nil}(B_{})`$. It follows that $`b^m=0`$ for $`m`$ sufficiently large, hence $`\nu _b^m=0`$ for $`m`$ sufficiently large. Let $`𝔭`$ be any prime ideal of $`A_{}`$; let $`\pi :A_{}A_{}/𝔭`$ be the natural projection and $`F`$ the fraction field of $`A_{}/𝔭`$. The $`F`$-linear morphism $`\nu _{b}^{}{}_{}{}^{}_A_{}\text{1}_F`$ is nilpotent on the vector space $`F^n`$, hence $`\pi \text{Tr}(\nu _{b}^{}{}_{}{}^{})=\text{Tr}(\nu _{b}^{}{}_{}{}^{}_A_{}\text{1}_F)=0`$. This shows that $`\text{Tr}(\nu _{b}^{}{}_{}{}^{})`$ lies in the intersection of all prime ideals of $`A_{}`$, hence it is nilpotent. Since by hypothesis $`A`$ is reduced, we get $`\text{Tr}(\nu _{b}^{}{}_{}{}^{})=0`$. Finally, this implies $`\text{Tr}_{B/A}(b)=0`$. Now, for any $`b^{}B_{}`$, the almost element $`bb^{}`$ will be nilpotent as well, so the same conclusion applies to it. This shows that $`\tau _{B/A}(b)=0`$. But by hypothesis $`B`$ is étale over $`A`$, hence theorem 3.2.9 yields $`b=0`$, as required. ∎
###### Remark 3.2.16.
Let $`M`$ be an $`A`$-module. We say that an almost element $`a`$ of $`A`$ is $`M`$-regular if the multiplication morphism $`mam:MM`$ is a monomorphism. Assume (A) (cf. section 2.1) and suppose furthermore that $`𝔪`$ is generated by a multiplicative system S which is a cofiltered semigroup under the preorder structure $`(\text{S},)`$ induced by the divisibility relation in $`V`$. We say that S is archimedean if, for all $`s,t\text{S}`$ there exists $`n>0`$ such that $`s^nt`$. Suppose that S is archimedean and that $`A`$ is a reduced almost algebra. Then S consists of $`A`$-regular elements. Indeed, by hypothesis $`\text{nil}(A_{})^a=0`$; since the annihilator of S in $`A_{}`$ is $`0`$ we get $`\text{nil}(A_{})=0`$. Suppose that $`sa=0`$ for some $`s\text{S}`$ and $`aA_{}`$. Let $`t\text{S}`$ be arbitrary and pick $`n>0`$ such that $`t^ns`$. Then $`(ta)^n=0`$ hence $`ta=0`$ for all $`t\text{S}`$, hence $`a=0`$.
### 3.3. Lifting theorems
Throughout the following, the terminology “epimorphism of $`V^a`$-algebras” will refer to a morphism of $`V^a`$-algebras that induces an epimorphism on the underlying $`V^a`$-modules.
###### Lemma 3.3.1.
Let $`AB`$ be an epimorphism of almost $`V`$-algebras with kernel $`I`$. Let $`U`$ be the $`A`$-extension $`0I/I^2A/I^2B0`$. Then the assignment $`ffU`$ defines a natural isomorphism
(3.3.2)
###### Proof.
Let $`X=(0ME\stackrel{p}{}B0)`$ be any $`A`$-extension of $`B`$ by $`M`$. The composition $`g:AE\stackrel{p}{}B`$ of the structural morphism for $`E`$ followed by $`p`$ coincides with the projection $`AB`$. Therefore $`g(I)M`$ and $`g(I^2)=0`$. Hence $`g`$ factors through $`A/I^2`$; the restriction of $`g`$ to $`I/I^2`$ defines a morphism $`f\text{Hom}_B(I/I^2,M)`$ and a morphism of $`A`$-extensions $`fUX`$. In this way we obtain an inverse for (3.3.2). ∎
Now consider any morphism of $`A`$-extensions
(3.3.3)
The morphism $`u`$ induces by adjunction a morphism of $`C_0`$-modules
(3.3.4)
$$C_0_{B_0}IJ$$
whose image is the ideal $`IC`$, so that the square diagram of almost algebras defined by $`\stackrel{~}{f}`$ is cofibred (i.e. $`C_0C_BB_0`$) if and only if (3.3.4) is an epimorphism.
###### Lemma 3.3.5.
Let $`\stackrel{~}{f}:\stackrel{~}{B}\stackrel{~}{C}`$ be a morphism of $`A`$-extensions as above, such that the corresponding square diagram of almost algebras is cofibred. Then the morphism $`f:BC`$ is flat if and only if $`f_0:B_0C_0`$ is flat and (3.3.4) is an isomorphism.
###### Proof.
It follows directly from the (almost version of the) local flatness criterion (see Th. 22.3). ∎
We are now ready to put together all the work done so far and begin the study of deformations of almost algebras.
The morphism $`u:IJ`$ is an element in $`\text{Hom}_{B_0}(I,J)`$; by lemma 3.3.1 the latter group is naturally isomorphic to $`\text{Exal}_B(B_0,J)`$. By applying transitivity (theorem 2.4.18) to the sequence of morphisms $`BB_0\stackrel{f_0}{}C_0`$ we obtain an exact sequence of abelian groups
$$\text{Exal}_{B_0}(C_0,J)\text{Exal}_B(C_0,J)\text{Hom}_{B_0}(I,J)\stackrel{}{}𝔼\text{xt}_{C_{0!!}}^2(𝕃_{C_0/B_0},J_!).$$
Hence we can form the element $`\omega (\stackrel{~}{B},f_0,u)=(u)𝔼\text{xt}_{C_{0!!}}^2(𝕃_{C_0/B_0},J_!)`$. The proof of the next result goes exactly as in (III.2.1.2.3).
###### Proposition 3.3.6.
i) Let the $`A`$-extension $`\stackrel{~}{B}`$, the $`B_0`$-linear morphism $`u:IJ`$ and the morphism of $`A`$-algebras $`f_0:B_0C_0`$ be given as above. Then there exists an $`A`$-extension $`\stackrel{~}{C}`$ and a morphism $`\stackrel{~}{f}:\stackrel{~}{B}\stackrel{~}{C}`$ completing diagram (3.3.3) if and only if $`\omega (\stackrel{~}{B},f_0,u)=0`$. (i.e. $`\omega (\stackrel{~}{B},f_0,u)`$ is the obstruction to the lifting of $`\stackrel{~}{B}`$ over $`f_0`$.)
ii) Assume that the obstruction $`\omega (\stackrel{~}{B},f_0,u)`$ vanishes. Then the set of isomorphism classes of $`A`$-extensions $`\stackrel{~}{C}`$ as in (i) forms a torsor under the group $`\text{Exal}_{B_0}(C_0,J)`$ ($`𝔼\text{xt}_{C_{0!!}}^1(𝕃_{C_0/B_0},J_!)`$).
iii) The group of automorphisms of an $`A`$-extension $`\stackrel{~}{C}`$ as in (i) is naturally isomorphic to $`\text{Der}_{B_0}(C_0,J)`$ ($`𝔼\text{xt}_{C_{0!!}}^0(𝕃_{C_0/B_0},J_!)`$). ∎
The obstruction $`\omega (\stackrel{~}{B},f_0,u)`$ depends functorially on $`u`$. More exactly, if we denote by
$$\omega (\stackrel{~}{B},f_0)𝔼\text{xt}_{C_{0!!}}^2(𝕃_{C_0/B_0},(C_0_{B_0}I)_!)$$
the obstruction corresponding to the natural morphism $`IC_0_{B_0}I`$, then for any other morphism $`u:IJ`$ we have
$$\omega (\stackrel{~}{B},f_0,u)=v_!\omega (\stackrel{~}{B},f_0)$$
where $`v`$ is the morphism (3.3.4). Taking lemma 3.3.5 into account we deduce
###### Corollary 3.3.7.
Suppose that $`B_0C_0`$ is flat. Then
i) The class $`\omega (\stackrel{~}{B},f_0)`$ is the obstruction to the existence of a flat deformation of $`C_0`$ over $`B`$, i.e. of a $`B`$-extension $`\stackrel{~}{C}`$ as in (3.3.3) such that $`C`$ is flat over $`B`$ and $`C_BB_0C_0`$ is an isomorphism.
ii) If the obstruction $`\omega (\stackrel{~}{B},f_0)`$ vanishes, then the set of isomorphism classes of flat deformations of $`C_0`$ over $`B`$ forms a torsor under the group $`\text{Exal}_{B_0}(C_0,C_0_{B_0}I)`$.
iii) The group of automorphisms of a given flat deformation of $`C_0`$ over $`B`$ is naturally isomorphic to $`\text{Der}_{B_0}(C_0,C_0_{B_0}I)`$. ∎
Now, suppose we are given two $`A`$-extensions $`\stackrel{~}{C}^1,\stackrel{~}{C}^2`$ with morphisms of $`A`$-extensions
and morphisms $`v:J^1J^2`$, $`g_0:C_0^1C_0^2`$ such that
(3.3.8)
$$u^2=vu^1\text{and}f_0^2=g_0f_0^1.$$
We consider the problem of finding a morphism of $`A`$-extensions
(3.3.9)
such that $`\stackrel{~}{f}^2=\stackrel{~}{g}\stackrel{~}{f}^1`$. Let us denote by $`e(\stackrel{~}{C}^i)𝔼\text{xt}_{C_{0!!}^i}^1(𝕃_{C_0^i/B},J_!^i)`$ the classes defined by the $`B`$-extensions $`\stackrel{~}{C}^1,\stackrel{~}{C}^2`$ via the isomorphism of theorem 2.4.17 and by
$$\begin{array}{cc}\hfill v:& 𝔼\text{xt}_{C_{0!!}^1}^1(𝕃_{C_0^1/B},J_!^1)𝔼\text{xt}_{C_{0!!}^1}^1(𝕃_{C_0^1/B},J_!^2)\hfill \\ \hfill g_0:& 𝔼\text{xt}_{C_{0!!}^2}^1(𝕃_{C_0^2/B},J_!^2)𝔼\text{xt}_{C_{0!!}^2}^1(C_{0!!}^2_{C_{0!!}^1}𝕃_{C_0^1/B},J_!^2)\hfill \end{array}$$
the canonical morphisms defined by $`v`$ and $`g_0`$. Using the natural isomorphism
$$𝔼\text{xt}_{C_{0!!}^1}^1(𝕃_{C_0^1/B},J_!^2)𝔼\text{xt}_{C_{0!!}^2}^1(C_{0!!}^2_{C_{0!!}^1}𝕃_{C_0^1/B},J_!^2)$$
we can identify the target of both $`v`$ and $`g`$ with $`𝔼\text{xt}_{C_{0!!}^1}^1(𝕃_{C_0^1/B},J_!^2)`$. It is clear that the problem admits a solution if and only if the $`A`$-extensions $`v\stackrel{~}{C}^1`$ and $`\stackrel{~}{C}^2g_0`$ coincide, i.e. if and only if $`ve(\stackrel{~}{C}^1)e(\stackrel{~}{C}^2)g_0=0`$. By applying transitivity to the sequence of morphisms $`BB_0C_0^1`$ we obtain an exact sequence
$$𝔼\text{xt}_{C_{0!!}^1}^1(𝕃_{C_0^1/B_0},J_!^2)𝔼\text{xt}_{C_{0!!}^1}^1(𝕃_{C_0^1/B},J_!^2)\text{Hom}_{C_0^1}(C_0^1_{B_0}I,J^2)$$
It follows from (3.3.8) that the image of $`ve(\stackrel{~}{C}^1)e(\stackrel{~}{C}^2)g_0`$ in the group $`\text{Hom}_{C_0^1}(C_0^1_{B_0}I,J^2)`$ vanishes, therefore
(3.3.10)
$$ve(\stackrel{~}{C}^1)e(\stackrel{~}{C}^2)g_0𝔼\text{xt}_{C_{0!!}^1}^1(𝕃_{C_0^1/B_0},J_!^2).$$
In conclusion, we derive the following result as in (III.2.2.2).
###### Proposition 3.3.11.
With the above notations, the class (3.3.10) is the obstruction to the existence of a morphism of $`A`$-extensions $`\stackrel{~}{g}:\stackrel{~}{C}^1\stackrel{~}{C}^2`$ as in (3.3.9) such that $`\stackrel{~}{f}^2=\stackrel{~}{g}\stackrel{~}{f}^1`$. When the obstruction vanishes, the set of such morphisms forms a torsor under the group $`\text{Der}_{B_0}(C_0^1,J^2)`$ (the latter being identified with $`𝔼\text{xt}_{C_{0!!}^2}^0(C_{0!!}^2_{C_{0!!}^1}𝕃_{C_0^1/B_0},J_!^2)`$). ∎
For a given almost $`V`$-algebra $`A`$, we define the category $`𝐰.\stackrel{´}{𝐄}𝐭(A)`$ as the full subcategory of $`A\text{-}\mathrm{𝐀𝐥𝐠}`$ consisting of all weakly étale $`A`$-algebras. Notice that, by lemma 3.1.2(iv) all morphisms in $`𝐰.\stackrel{´}{𝐄}𝐭(A)`$ are weakly étale.
###### Theorem 3.3.12.
i) Let $`AB`$ be a weakly étale morphism of almost algebras. Let $`C`$ be any $`A`$-algebra and $`IC`$ a nilpotent ideal. Then the natural morphism
$$\text{Hom}_{A\text{-}\mathrm{𝐀𝐥𝐠}}(B,C)\text{Hom}_{A\text{-}\mathrm{𝐀𝐥𝐠}}(B,C/I)$$
is bijective.
ii) Let $`A`$ be a $`V^a`$-algebra, $`IA`$ a nilpotent ideal and $`A^{}=A/I`$. Then the natural functor
$$𝐰.\stackrel{´}{𝐄}𝐭(A)𝐰.\stackrel{´}{𝐄}𝐭(A^{})(\varphi :AB)(\text{1}_A^{}_A\varphi :A^{}A^{}_AB)$$
is an equivalence of categories.
iii) The equivalence of (ii) restricts to an equivalence $`\stackrel{´}{𝐄}𝐭(A)\stackrel{´}{𝐄}𝐭(A^{})`$.
###### Proof.
By induction we can assume $`I^2=0`$. Then (i) follows directly from proposition 3.3.11 and theorem 2.4.23. We show (ii) : by corollary 3.3.7 (and again theorem 2.4.23) a given weakly étale morphism $`\varphi ^{}:A^{}B^{}`$ can be lifted to a unique flat morphism $`\varphi :AB`$. We need to prove that $`\varphi `$ is weakly étale, i.e. that $`B`$ is $`B_AB`$-flat. However, it is clear that $`\mu _{B^{}/A^{}}:B^{}_A^{}B^{}B^{}`$ is weakly étale, hence it has a flat lifting $`\stackrel{~}{\mu }:B_ABC`$. Then the composition $`AB_ABC`$ is flat and it is a lifting of $`\varphi ^{}`$. We deduce that there is an isomorphism of $`A`$-algebras $`\alpha :BC`$ lifting $`\text{1}_B^{}`$ and moreover the morphisms $`b\stackrel{~}{\mu }(b\underset{¯}{1})`$ and $`b\stackrel{~}{\mu }(\underset{¯}{1}b)`$ coincide with $`\alpha `$. Claim (ii) follows. To show (iii), suppose that $`A^{}B^{}`$ is étale and let $`I_{B^{}/A^{}}`$ denote as usual the kernel of $`\mu _{B^{}/A^{}}`$. By corollary 3.1.8 there is a natural morphism of almost algebras $`B^{}_A^{}B^{}I_{B^{}/A^{}}`$ which is clearly étale. Hence $`I_{B^{}/A^{}}`$ lifts to a weakly étale $`B_AB`$-algebra $`C`$, and the isomorphism $`B^{}_A^{}B^{}I_{B^{}/A^{}}B^{}`$ lifts to an isomorphism $`B_ABCB`$ of $`B_AB`$-algebras. It follows that $`B`$ is an almost projective $`B_AB`$-module, i.e. $`AB`$ is étale, as claimed. ∎
We conclude with some results on deformations of almost modules. These can be established independently of the theory of the cotangent complex, along the lines of (IV.3.1.12).
We begin by recalling some notation from loc. cit. Let $`R`$ be a ring and $`JR`$ an ideal with $`J^2=0`$. Set $`R^{}=R/J`$; an extension of $`R`$-modules $`\underset{¯}{M}=(0KM\stackrel{p}{}M^{}0)`$ where $`K`$ and $`M^{}`$ are killed by $`J`$, defines a natural morphism of $`R^{}`$-modules $`u(\underset{¯}{M}):J_R^{}M^{}K`$ such that $`u(\underset{¯}{M})(xm^{})=xm`$ for $`xJ`$, $`mM`$ and $`p(m)=m^{}`$. By the local flatness criterion ( Th.22.3) $`M`$ is flat over $`R`$ if and only if $`M^{}`$ is flat over $`R^{}`$ and $`u(\underset{¯}{M})`$ is an isomorphism. One can then show the following.
###### Proposition 3.3.13.
(cp. (IV.3.1.5))
i) Given $`R^{}`$-modules $`M^{}`$ and $`K`$ and a morphism $`u^{}:J_R^{}M^{}K`$ there exists an obstruction $`\omega (R,u^{})\text{Ext}_R^{}^2(M^{},K)`$ whose vanishing is necessary and sufficient for the existence of an extension of $`R`$-modules $`\underset{¯}{M}`$ of $`M^{}`$ by $`K`$ such that $`u(\underset{¯}{M})=u^{}`$.
ii) When $`\omega (R,u^{})=0`$, the set of isomorphism classes of such extensions $`\underset{¯}{M}`$ forms a torsor under $`\text{Ext}_R^{}^1(M^{},K)`$; the group of automorphisms of such an extension is $`\text{Hom}_R^{}(M^{},K)`$. ∎
###### Lemma 3.3.14.
Let $`AB`$ be a finite morphism of almost algebras with nilpotent kernel. Let $`\varphi :MN`$ be an $`A`$-linear morphism and set $`\varphi _B=\varphi _A\text{1}_B:M_ABN_AB`$. Then there exists $`m0`$ such that
i) $`\text{Ann}_A(\text{Coker}(\varphi _B))^m\text{Ann}_A(\text{Coker}(\varphi ))`$.
ii) $`(\text{Ann}_V(\text{Ker}(\varphi _B))\text{Ann}_V(\text{Tor}_1^A(B,N))\text{Ann}_V(\text{Coker}(\varphi )))^m\text{Ann}_A(\text{Ker}(\varphi ))`$.
If $`B=A/I`$ for some nilpotent ideal $`I`$, and $`I^n=0`$, then we can take $`m=n`$ in (i) and (ii).
###### Proof.
Under the assumptions, we can find a finitely generated $`A_{}`$-module $`Q`$ such that $`𝔪B_{}QB_{}`$. By (1.1.5), there exists a finite filtration $`0=J_m\mathrm{}J_1J_0=A_{}`$ such that each $`J_i/J_{i+1}`$ is a quotient of a direct sum of copies of $`Q`$. This implies that, for every $`A`$-module $`M`$, we have
(3.3.15)
$$\text{Ann}_A(M_AB)^m\text{Ann}_A(M).$$
(i) follows easily. Notice that if $`B=A/I`$ and $`I^n=0`$, then we can take $`m=n`$ in (3.3.15). For (ii) let $`C^{}=\text{Cone}(\varphi )`$. We estimate $`H=H^1(C^{}\stackrel{\text{L}}{}_AB)`$ in two ways. By the first spectral sequence of hyperhomology we have an exact sequence $`\text{Tor}_1^A(N,B)H\text{Ker}(\varphi _B)`$. By the second spectral sequence for hyperhomology we have an exact sequence $`\text{Tor}_2^A(\text{Coker}(\varphi ),B)\text{Ker}(\varphi )_ABH`$. Hence $`\text{Ker}(\varphi )_AB`$ is annihilated by the product of the three annihilators in (ii) and the result follows by applying (3.3.15) with $`M=\text{Ker}(\varphi )`$. ∎
###### Lemma 3.3.16.
Keep the assumptions of lemma 3.3.14, let $`M`$ be an $`A`$-module and set $`M_B=B_AM`$.
i) If $`AB`$ is an epimorphism, $`M`$ is flat and $`M_B`$ is almost projective over $`B`$, then $`M`$ is almost projective over $`A`$.
ii) If $`M_B`$ is an almost finitely generated $`B`$-module then $`M`$ is an almost finitely generated $`A`$-module.
iii) If $`\text{Tor}_1^A(B,M)=0`$ and $`M_B`$ is almost finitely presented over $`B`$, then $`M`$ is almost finitely presented over $`A`$.
###### Proof.
(i) : we have to show that $`\text{Ext}_A^1(M,N)`$ is almost zero for every $`A`$-module $`N`$. Let $`I=\text{Ker}(AB)`$; by assumption $`I`$ is nilpotent, so by the usual devissage we may assume that $`IN=0`$. If $`\chi \text{Ext}_A^1(M,N)`$ is represented by an extension $`0NQM0`$ then after tensoring by $`B`$ and using the flatness of $`M`$ we get an exact sequence of $`B`$-modules $`0NB_AQM_B0`$. Thus $`\chi `$ comes from an element of $`\text{Ext}_B^1(M_B,N)`$ which is almost zero by assumption.
(ii) : let $`𝔪_0=(\epsilon _1,\mathrm{},\epsilon _m)`$ be a finitely generated subideal of $`𝔪`$. By assumption there is a map $`\varphi ^{}:B^rM_B`$ such that $`𝔪_0\text{Coker}(\varphi ^{})=0`$. For all $`jm`$ the morphism $`\epsilon _j\varphi ^{}`$ lifts to a morphism $`\varphi _j:A^rM`$. Then $`\varphi =\varphi _1\mathrm{}\varphi _m:A^{rm}M`$ satisfies $`𝔪_0^2\text{Coker}(\varphi _A\text{1}_B)=0`$. By lemma 3.3.14(i) it follows $`𝔪_0^{2n}\text{Coker}(\varphi )=0`$ for some $`n0`$. As $`𝔪_0`$ was arbitrary, the result follows.
(iii) Let $`𝔪_0`$ be as above. By (ii), $`M`$ is almost finitely generated over $`A`$, so we can choose a morphism $`\varphi :A^rM`$ such that $`𝔪_0\text{Coker}(\varphi )=0`$. Consider $`\varphi _B=\varphi _A\text{1}_B:B^rM_B`$. By lemma 2.3.6, there is a finitely generated submodule $`N`$ of $`\text{Ker}(\varphi _B)`$ containing $`𝔪_0^2\text{Ker}(\varphi _B)`$. Notice that $`\text{Ker}(\varphi )_AB`$ maps onto $`\text{Ker}(B^r\text{Im}(\varphi )_AB)`$ and $`\text{Ker}(\text{Im}(\varphi )_ABM_B)\text{Tor}_1^A(B,\text{Coker}(\varphi ))`$ is annihilated by $`𝔪_0`$. Hence $`𝔪_0\text{Ker}(\varphi _B)`$ is contained in the image of $`\text{Ker}(\varphi )`$ and therefore we can lift a finite generating set $`\{x_1^{},\mathrm{},x_n^{}\}`$ for $`𝔪_0^2N`$ to almost elements $`\{x_1,\mathrm{},x_n\}`$ of $`\text{Ker}(\varphi )`$. If we quotient $`A^r`$ by the span of these $`x_i`$, we get a finitely presented $`A`$-module $`F`$ with a morphism $`\overline{\varphi }:FM`$ such that $`\text{Ker}(\overline{\varphi }_AB)`$ is annihilated by $`𝔪_0^4`$ and $`\text{Coker}(\overline{\varphi })`$ is annihilated by $`𝔪_0`$. By lemma 3.3.14(ii) we derive $`𝔪_0^{5m}\text{Ker}(\overline{\varphi })=0`$ for some $`m0`$. Since $`𝔪_0`$ is arbitrary, this proves the result. ∎
###### Remark 3.3.17.
(i) Inspecting the proof, one sees that parts (ii) and (iii) of lemma 3.3.16 hold whenever (3.3.15) holds. For instance, if $`AB`$ is any faithfully flat morphism, then (3.3.15) holds with $`m=1`$.
ii) Consequently, if $`AB`$ is faithfully flat and $`M`$ is an $`A`$-module such that $`M_B`$ is flat (resp. almost finitely generated, resp. almost finitely presented) over $`B`$, then $`M`$ is flat (resp. almost finitely generated, resp. almost finitely presented) over $`A`$.
iii) On the other hand, we do not know whether a general faithfully flat morphism $`AB`$ descends almost projectivity. However, using (ii) and proposition 2.3.15 we see that if the $`B`$-module $`M_B`$ is almost finitely generated almost projective, then $`M`$ has the same property.
iv) However, if $`B`$ is faithfully flat and almost finitely presented as an $`A`$-module, then $`AB`$ does descend almost projectivity, as can be easily deduced from lemma 2.3.23(i) and proposition 2.3.15(ii).
###### Theorem 3.3.18.
Let $`I`$ be a nilpotent ideal of the almost algebra $`A`$ and set $`A^{}=A/I`$. Suppose that $`\stackrel{~}{𝔪}`$ is a (flat) $`V`$-module of homological dimension $`1`$. Let $`P^{}`$ be an almost projective $`A^{}`$-module.
i) There is an almost projective $`A`$-module $`P`$ with $`A^{}_APP^{}`$.
ii) If $`P^{}`$ is almost finitely presented, then $`P`$ is almost finitely presented.
###### Proof.
As usual we reduce to $`I^2=0`$. Then proposition 3.3.13(i) applies with $`R=A_{}`$, $`J=I_{}`$, $`R^{}=A_{}/I_{}`$, $`M^{}=P_!^{}`$, $`K=I_{}_R^{}P_!^{}`$ and $`u^{}=\text{1}_K`$. We obtain a class $`\omega (A_{},u^{})\text{Ext}_R^{}^2(P_!^{},I_{}_R^{}P_!^{})`$ which gives the obstruction to the existence of a flat $`A_{}`$-module $`F`$ lifting $`P_!^{}`$. Since $`P_!^{}`$ is almost projective, we know that $`𝔪\text{Ext}_R^{}^2(P_!^{},I_{}_R^{}P_!^{})=0`$, which says that $`0=\epsilon \omega (A_{},u^{})=\omega (A_{},\epsilon u^{})`$ for all $`\epsilon 𝔪`$. In other words, for every $`\epsilon 𝔪`$ we can find an extension of $`A_{}`$-modules $`\underset{¯}{P_\epsilon }`$ of $`P_!^{}`$ by $`I_{}_R^{}P_!^{}`$ such that $`u(\underset{¯}{P_\epsilon })=\epsilon \text{1}_{I_{}_R^{}P_!^{}}`$. Let $`\chi _\epsilon \text{Ext}_A_{}^1(P_!^{},I_{}_R^{}P_!^{})`$ be the class of $`\underset{¯}{P_\epsilon }`$. Notice that, for any $`\delta 𝔪`$, $`\delta \chi _\epsilon `$ is the class of an extension $`\underset{¯}{X}`$ such that $`u(\underset{¯}{X})=\delta u(\underset{¯}{P_\epsilon })=\delta \epsilon \text{1}_{I_{}_R^{}P_!^{}}`$, hence, by proposition 3.3.13(ii), $`\gamma (\delta \chi _\epsilon \chi _{\delta \epsilon })=0`$ for all $`\gamma 𝔪`$. Hence we can define a morphism
$$\chi :𝔪_V𝔪_V𝔪\text{Ext}_A_{}^1(P_!^{},I_{}_R^{}P_!^{})\epsilon \delta \gamma \delta \gamma \chi _\epsilon .$$
However, one sees easily that $`𝔪_V𝔪_V𝔪\stackrel{~}{𝔪}`$ and $`\stackrel{~}{𝔪}_VP_!^{}P_!^{}`$, hence we can view $`\chi `$ as an element of $`\text{Hom}_V(\stackrel{~}{𝔪},\text{Ext}_A_{}^1(P_!^{},I_{}_R^{}P_!^{}))`$ and moreover we have a spectral sequence
$$E_2^{pq}=\text{Ext}_V^p(\stackrel{~}{𝔪},\text{Ext}_A_{}^q(P_!^{},I_{}_R^{}P_!^{}))\text{Ext}_A_{}^{p+q}(P_!^{},I_{}_R^{}P_!^{})$$
with $`E_2^{pq}=0`$ for all $`p2`$ (this spectral sequence is constructed e.g. from the double complex $`\text{Hom}_V(F_p,\text{Hom}_A_{}(F_q^{},I_{}_R^{}P_!^{}))`$ where $`F_{}`$ (resp. $`F_{}^{}`$) is a projective resolution of $`\stackrel{~}{𝔪}`$ (resp. $`P_!^{}`$)). In particular, our $`\chi `$ is an element in $`E_2^{01}`$ which therefore survives in the abutment as a class of $`E_{\mathrm{}}^{01}`$. The latter can be lifted to an element $`\stackrel{~}{\chi }`$ via the surjection $`\text{Ext}_A_{}^1(P_!^{},I_{}_R^{}P_!^{})E_{\mathrm{}}^{01}`$. Let $`0I_{}_R^{}P_!^{}QP_!^{}0`$ be an extension representing $`\stackrel{~}{\chi }`$. Checking compatibilities, we see that $`\delta \epsilon \stackrel{~}{\chi }=\delta \chi _\epsilon `$ for every $`\epsilon ,\delta 𝔪`$. Hence $`u(\stackrel{~}{\chi }):I_{}_R^{}P_!^{}I_{}_R^{}P_!^{}`$ coincides with the identity map on the submodule $`𝔪I_{}_R^{}P_!^{}`$. Since $`𝔪P_!=P_!`$, we see that $`u(\stackrel{~}{\chi })`$ is actually the identity map. By the local flatness criterion, it then follows that $`Q`$ is flat over $`R`$, hence the $`A`$-module $`P=Q^a`$ is a flat lifting of $`P^{}`$, so it is almost projective, by lemma 3.3.16(i). Now (ii) follows from (i), lemma 3.3.16(ii) and proposition 2.3.15(i). ∎
###### Remark 3.3.19.
(i) According to proposition 2.1.10(ii), theorem 3.3.18 applies especially when $`𝔪`$ is countably generated as a $`V`$-module.
(ii) For $`P`$ and $`P^{}`$ as in theorem 3.3.18(ii) let $`\sigma _P:PP^{}`$ be the projection. It is natural to ask whether the pair $`(P,\sigma _P)`$ is uniquely determined up to isomorphism, i.e. whether, for any other pair $`(Q,\sigma _Q:QP^{})`$ for which theorem 3.3.18 holds, there exists an $`A`$-linear isomorphism $`\varphi :PQ`$ such that $`\sigma _Q\varphi =\sigma _P`$. The answer is negative in general. Consider the case $`P^{}=A^{}`$. Take $`P=Q=A`$ and let $`\sigma _P`$ be the natural projection, while $`\sigma _Q=(u^{}\text{1}_A^{})\sigma _P`$, where $`u^{}`$ is a unit in $`A_{}^{}`$. Then the uniqueness question amounts to whether every unit in $`A_{}^{}`$ lifts to a unit of $`A_{}`$. The following counterexample is related to the fact that the completion of the algebraic closure $`\overline{}_p`$ of $`_p`$ is not maximally complete. Let $`V=\overline{}_p`$, the integral closure of $`_p`$ in $`\overline{}_p`$. Then $`V`$ is a non-discrete valuation ring of rank one, and we take $`𝔪`$ as in example 2.1.1(i), $`A=(V/p^2V)^a`$ and $`A^{}=A/pA`$. Choose a compatible system of roots of $`p`$. An almost element of $`A^{}`$ is just a $`V`$-linear morphism $`\varphi :\underset{n>0}{\text{colim}}p^{1/n!}VV/pV`$. Such a $`\varphi `$ can be represented (in a non-unique way) by an infinite series of the form $`_{n=1}^{\mathrm{}}a_np^{11/n!}`$ ($`a_nV`$). The meaning of this expression is as follows. For every $`m>0`$, scalar multiplication by the element $`_{n=1}^ma_np^{11/n!}V`$ defines a morphism $`\varphi _m:p^{1/m!}VV/pV`$. For $`m^{}>m`$, let $`j_{m,m^{}}:p^{1/m!}Vp^{1/m^{}!}V`$ be the imbedding. Then we have $`\varphi _m^{}j_{m,m^{}}=\varphi _m`$, so that we can define $`\varphi =\underset{m>0}{\text{colim}}\varphi _m`$. Similarly, every almost element of $`A`$ can be represented by an expression of the form $`a_0+_{n=1}^{\mathrm{}}a_np^{21/n!}`$. Now, if $`\sigma :AA^{}`$ is the natural projection, the induced map $`\sigma _{}:A_{}A_{}^{}`$ is given by: $`a_0+_{n=1}^{\mathrm{}}a_np^{21/n!}a_0`$. In particular, its image is the subring $`V/p(V/p)_{}=A_{}^{}`$. For instance, the unit $`_{n=1}^{\mathrm{}}p^{11/n!}`$ of $`A_{}^{}`$ does not lie in the image of this map.
In the light of the above remark, the best one can achieve in general is the following result.
###### Proposition 3.3.20.
Assume (A) (see section 2.1) and keep the notation of theorem 3.3.18. Suppose moreover that $`(Q,\sigma _Q:QP^{})`$ and $`(P,\sigma _P:PP^{})`$ are two pairs as in remark 3.3.19. Then for all $`\epsilon 𝔪`$ there exist $`A`$-linear morphisms $`t_\epsilon :PQ`$ and $`s_\epsilon :QP`$ such that
PQ($`\epsilon `$)$`\begin{array}{cccc}\hfill \sigma _Qt_\epsilon =& \epsilon \sigma _P\hfill & \hfill \sigma _Ps_\epsilon =& \epsilon \sigma _Q\hfill \\ \hfill s_\epsilon t_\epsilon =& \epsilon ^2\text{1}_P\hfill & \hfill t_\epsilon s_\epsilon =& \epsilon ^2\text{1}_Q.\hfill \end{array}`$
###### Proof.
Since both $`Q`$ and $`P`$ are almost projective and $`\sigma _P,\sigma _Q`$ are epimorphisms, there exist morphisms $`\overline{t}_\epsilon :PQ`$ and $`\overline{s}_\epsilon :QP`$ such that $`\sigma _Q\overline{t}_\epsilon =\epsilon \sigma _P`$ and $`\sigma _P\overline{s}_\epsilon =\epsilon \sigma _Q`$. Then we have $`\sigma _P(\overline{s}_\epsilon \overline{t}_\epsilon \epsilon ^2\text{1}_P)=0`$ and $`\sigma _Q(\overline{t}_\epsilon \overline{s}_\epsilon \epsilon ^2\text{1}_Q)=0`$, i.e. the morphism $`u_\epsilon =\epsilon ^2\text{1}_P\overline{s}_\epsilon \overline{t}_\epsilon `$ (resp. $`v_\epsilon =\epsilon ^2\text{1}_Q\overline{t}_\epsilon \overline{s}_\epsilon `$) has image contained in the almost submodule $`IP`$ (resp. $`IQ`$). Since $`I^m=0`$ this implies $`u_\epsilon ^m=0`$ and $`v_\epsilon ^m=0`$. Hence
$$\epsilon ^{2m}\text{1}_P=(\epsilon ^2\text{1}_P)^mu_\epsilon ^m=(\underset{a=0}{\overset{m1}{}}\epsilon ^{2a}u_\epsilon ^{m1a})\overline{s}_\epsilon \overline{t}_\epsilon .$$
Define $`\overline{s}_{(2m1)\epsilon }=(_{a=0}^{m1}\epsilon ^{2a}u_\epsilon ^{m1a})\overline{s}_\epsilon `$. Notice that $`\overline{s}_{(2m1)\epsilon }=\overline{s}_\epsilon (_{a=0}^{m1}\epsilon ^{2a}v_\epsilon ^{m1a}).`$ This implies the equalities $`\overline{s}_{(2m1)\epsilon }\overline{t}_\epsilon =\epsilon ^{2m}\text{1}_P`$ and $`\overline{t}_\epsilon \overline{s}_{(2m1)\epsilon }=\epsilon ^{2m}\text{1}_Q`$. Then the pair $`(\overline{s}_{(2m1)\epsilon },\epsilon ^{2(m1)}\overline{t}_\epsilon )`$ satisfies PQ($`\epsilon ^{2m1}`$). Under (A), every element of $`𝔪`$ is a multiple of an element of the form $`\epsilon ^{2m1}`$, therefore the claim follows for arbitrary $`\epsilon 𝔪`$. ∎
### 3.4. Descent
Faithfully flat descent in the almost setting presents no particular surprises: since the functor $`AA_{!!}`$ preserves faithful flatness of morphisms (see remark 3.1.3) many well-known results for usual rings and modules extend verbatim to almost algebras. So for instance, faithfully flat morphisms are of universal effective descent for the fibred categories $`F:V^a\text{-}\mathrm{𝐀𝐥𝐠}.\mathrm{𝐌𝐨𝐝}^oV^a\text{-}\mathrm{𝐀𝐥𝐠}^o`$ and $`G:V^a\text{-}\mathrm{𝐀𝐥𝐠}.\mathrm{𝐌𝐨𝐫𝐩𝐡}^oV^a\text{-}\mathrm{𝐀𝐥𝐠}^o`$ (see definition 2.4.12: for an almost $`V`$-algebra $`B`$, the fibre $`F_B`$ (resp. $`G_B`$) is the opposite of the category of $`B`$-modules (resp. $`B`$-algebras)). Then, using remark 3.3.17, we deduce also universal effective descent for the fibred subcategories of flat (resp. almost finitely generated, resp. almost finitely presented, resp. almost projective almost finitely generated) modules. Likewise, a faithfully flat morphism is of universal effective descent for the fibred subcategories $`\stackrel{´}{𝐄}𝐭^oV^a\text{-}\mathrm{𝐀𝐥𝐠}^o`$ of étale (resp. $`𝐰.\stackrel{´}{𝐄}𝐭^oV^a\text{-}\mathrm{𝐀𝐥𝐠}^o`$ of weakly étale) algebras.
More generally, since the functor $`AA_{!!}`$ preserves pure morphisms in the sense of , and since, by a theorem of Olivier (loc. cit.), pure morphisms are of universal effective descent for modules, the same holds for pure morphisms of almost algebras.
Non-flat descent is more delicate. Our results are not as complete here as it could be wished, but nevertheless, they suffice for current applications (namely, for the cases needed in ).
Our first statement is the almost version of a theorem of Gruson and Raynaud (cp. (Part II, Th.1.2.4)).
###### Proposition 3.4.1.
A finite monomorphism of almost algebras descends flatness.
###### Proof.
Let $`\varphi :AB`$ be such a morphism. Under the assumption, we can find a finite $`A_{}`$-module $`Q`$ such that $`𝔪B_{}QB_{}`$. One sees easily that $`Q`$ is a faithful $`A_{}`$-module, so by (Part II, Th.1.2.4 and lemma 1.2.2), $`Q`$ satisfies the following condition :
(3.4.2) If $`(0NLP0)`$ is an exact sequence of $`A_{}`$-modules with $`L`$ flat, such that $`\text{Im}(N_A_{}Q)`$ is a pure submodule of $`L_A_{}Q`$, then $`P`$ is flat.
Now let $`M`$ be an $`A`$-module such that $`M_AB`$ is flat. Pick an epimorphism $`p:FM`$ with $`F`$ free over $`A`$. Then $`\underset{¯}{Y}=(0\text{Ker}(p_A\text{1}_B)F_ABM_AB0)`$ is universally exact over $`B`$, hence over $`A`$. Consider the sequence $`\underset{¯}{X}=(0\text{Im}(\text{Ker}(p)_!_A_{}Q)F_!_A_{}QM_!_A_{}Q0)`$. Clearly $`\underset{¯}{X}^a\underset{¯}{Y}`$. However, it is easy to check that a sequence $`\underset{¯}{E}`$ of $`A`$-modules is universally exact if and only if the sequence $`\underset{¯}{E}_!`$ is universally exact over $`A_{}`$. We conclude that $`\underset{¯}{X}=(\underset{¯}{X}^a)_!`$ is a universally exact sequence of $`A_{}`$-modules, hence, by condition (3.4.2), $`M_!`$ is flat over $`A_{}`$, i.e. $`M`$ is flat over $`A`$ as required. ∎
###### Corollary 3.4.3.
Let $`AB`$ be a finite morphism of almost algebras, with nilpotent kernel. If $`C`$ is a flat $`A`$-algebras such that $`C_AB`$ is weakly étale (resp. étale) over $`B`$, then $`C`$ is weakly étale (resp. étale) over $`A`$.
###### Proof.
In the weakly étale case, we have to show that the multiplication morphism $`\mu :C_ACC`$ is flat. As $`N=\text{Ker}(AB)`$ is nilpotent, the local flatness criterion reduces the question to the situation over $`A/N`$. So we may assume that $`AB`$ is a monomorphism. Then $`C_AC(C_AC)_AB`$ is a monomorphism, but $`\mu _{C_AC}\text{1}_{(C_AC)_AB}`$ is the multiplication morphism of $`C_AB`$, which is flat by assumption. Therefore, by proposition 3.4.1, $`\mu `$ is flat.
For the étale case, we have to show that $`C`$ is almost finitely presented as a $`C_AC`$-module. By hypothesis $`C_AB`$ is almost finitely presented as a $`C_AC_AB`$-module and we know already that $`C`$ is flat as a $`C_AC`$-module, so by lemma 3.3.16(iii) (applied to the finite morphism $`C_ACC_AC_AB`$) the claim follows. ∎
Next we consider the following situation. We are given a cartesian diagram of almost algebras
(3.4.4)
such that one of the morphisms $`A_iA_3`$ ($`i=1,2`$) is an epimorphism. We denote by $`\text{M}_i`$ (resp. $`\text{M}_{i,\mathrm{fl}}`$, resp. $`\text{M}_{i,\mathrm{proj}}`$) the category of all (resp. flat, resp. almost projective) $`A_i`$-modules, for $`i=0,\mathrm{},3`$. Diagram (3.4.4) induces an essentially commutative diagram for the corresponding categories $`\text{M}_i`$, where the arrows are given by the “extension of scalars” functors. There follows a natural functor
$$\pi :\text{M}_0\text{M}_1\times _{\text{M}_3}\text{M}_2$$
from $`\text{M}_0`$ to the 2-fibred products of $`\text{M}_1`$ and $`\text{M}_2`$ over $`\text{M}_3`$. Recall (cp. (Ch.VII §3)) that $`\text{M}_1\times _{\text{M}_3}\text{M}_2`$ is the category whose objects are the triples $`(M_1,M_2,\xi )`$, where $`M_i`$ is an $`A_i`$-module ($`i=1,2`$) and $`\xi :A_3_{A_1}M_1\stackrel{}{}A_3_{A_2}M_2`$ is an $`A_3`$-linear isomorphism. Given such an object $`(M_1,M_2,\xi )`$, let us denote $`M_3=A_3_{A_2}M_2`$; we have a natural morphism $`M_2M_3`$, and $`\xi `$ gives a morphism $`M_1M_3`$, so we can form the fibre product $`T(M_1,M_2,\xi )=M_1\times _{M_3}M_2`$. In this way we obtain a functor $`T:\text{M}_1\times _{\text{M}_3}\text{M}_2\text{M}_0`$, and we leave to the reader the verification that $`T`$ is right adjoint to $`\pi `$. Let us denote by $`\epsilon :\text{1}_{\text{M}_0}T\pi `$ and $`\eta :\pi T\text{1}_{\text{M}_1\times _{\text{M}_3}\text{M}_2}`$ the unit and counit of the adjunction.
###### Lemma 3.4.5.
The functor $`\pi `$ induces an equivalence of full subcategories :
having $`T`$ as essential inverse.
###### Proof.
General nonsense. ∎
###### Lemma 3.4.6.
Let $`M`$ be any $`A_0`$-module. Then $`\epsilon _M`$ is an epimorphism. If $`M`$ is flat over $`A_0`$, $`\epsilon _M`$ is an isomorphism.
###### Proof.
Indeed, $`\epsilon _M:M(A_1_{A_0}M)\times _{A_3_{A_0}M}(A_2_{A_0}M)`$ is the natural morphism. So, the assertions follow by applying $`_{A_0}M`$ to the short exact sequence of $`A_0`$-modules
(3.4.7)
$$0A_0\stackrel{f}{}A_1A_2\stackrel{g}{}A_30$$
where $`f(a)=(f_1(a),f_2(a))`$ and $`g(a,b)=g_1(a)g_2(b)`$. ∎
There is another case of interest, in which $`\epsilon _M`$ is an isomorphism. Namely, suppose that one of the morphisms $`A_iA_3`$ ($`i=1,2`$), say $`A_1A_3`$, has a section. Then also the morphism $`A_0A_2`$ gains a section $`s:A_2A_0`$ and we have the following :
###### Lemma 3.4.8.
In the above situation, suppose that the $`A_0`$-module $`M`$ arises by extension of scalars from an $`A_2`$-module $`M^{}`$, via the section $`s:A_2A_0`$. Then $`\epsilon _M`$ is an isomorphism.
###### Proof.
Indeed, in this case, (3.4.7) is split exact as a sequence of $`A_2`$-modules, and it remains such after tensoring by $`M^{}`$. ∎
###### Lemma 3.4.9.
$`\eta _{(M_1,M_2,\xi )}`$ is an isomorphism for all objects $`(M_1,M_2,\xi )`$.
###### Proof.
To fix ideas, suppose that $`A_1A_3`$ is an epimorphism. Consider any object $`(M_1,M_2,\xi )`$ of $`\text{M}_1\times _{\text{M}_3}\text{M}_2`$. Let $`M=T(M_1,M_2,\xi )`$; we deduce a natural morphism
$$\varphi :(M_{A_0}A_1)\times _{M_{A_0}A_3}(M_{A_0}A_2)M_1\times _{M_3}M_2$$
such that $`\varphi \epsilon _M=\text{1}_M`$. It follows that $`\epsilon _M`$ is injective, hence it is an isomorphism, by lemma 3.4.6. We derive a commutative diagram with exact rows :
From the snake lemma we deduce
$$\begin{array}{cc}()\hfill & \text{Ker}(\varphi _1)\text{Ker}(\varphi _2)\text{Ker}(\varphi _3)\hfill \\ ()\hfill & \text{Coker}(\varphi _1)\text{Coker}(\varphi _2)\text{Coker}(\varphi _3).\hfill \end{array}$$
Since $`M_3M_1_{A_1}A_3`$ we have $`A_3_{A_1}\text{Coker}(\varphi _1)\text{Coker}(\varphi _3)`$. But by assumption $`A_1A_3`$ is an epimorphism, so also $`\text{Coker}(\varphi _1)\text{Coker}(\varphi _3)`$ is an epimorphism. Then $`()`$ implies that $`\text{Coker}(\varphi _2)=0`$. But $`\varphi _3=\text{1}_{A_3}_{A_2}\varphi _2`$, thus $`\text{Coker}(\varphi _3)=0`$ as well. We look at the exact sequence $`0\text{Ker}(\varphi _1)M_{A_0}A_1\stackrel{\varphi _1}{}M_10`$ : applying $`A_3_{A_1}`$ we obtain an epimorphism $`A_3_{A_1}\text{Ker}(\varphi _1)\text{Ker}(\varphi _3)`$. From $`()`$ it follows that $`\text{Ker}(\varphi _2)=0`$. In conclusion, $`\varphi _2`$ is an isomorphism. Hence the same is true for $`\varphi _3=\text{1}_{A_3}_{A_2}\varphi _2`$, and again $`()`$, $`()`$ show that $`\varphi _1`$ is an isomorphism as well, which implies the claim. ∎
###### Lemma 3.4.10.
If $`(A_1\times A_2)_{A_0}M`$ is flat over $`A_1\times A_2`$, then $`M`$ is flat over $`A_0`$.
###### Proof.
Suppose that $`A_1A_3`$ is an epimorphism and let $`I`$ be its kernel. Let $`\stackrel{~}{A}=A_{1!!}\times _{A_{3!!}}A_{2!!}`$; it suffices to show that $`M_!`$ is a flat $`\stackrel{~}{A}`$-module. However, in view of proposition 2.3.27, the assumption implies that $`(A_{1!!}\times A_{2!!})_{\stackrel{~}{A}}M_!`$ is a flat $`A_{1!!}\times A_{2!!}`$-module. $`I_!`$ is the kernel of the epimorphism $`A_{1!!}A_{3!!}`$. Moreover, $`I_!`$ identifies naturally with an ideal of $`\stackrel{~}{A}`$ and $`\stackrel{~}{A}/I_!A_{2!!}`$. Then the desired conclusion follows from (lemma in loc. cit.). ∎
###### Proposition 3.4.11.
The functor $`\pi `$ restricts to equivalences :
$$\begin{array}{c}\text{M}_{0,\mathrm{fl}}\stackrel{}{}\text{M}_{1,\mathrm{fl}}\times _{\text{M}_{3,\mathrm{fl}}}\text{M}_{2,\mathrm{fl}}\\ \text{M}_{0,\mathrm{proj}}\stackrel{}{}\text{M}_{1,\mathrm{proj}}\times _{\text{M}_{3,\mathrm{proj}}}\text{M}_{2,\mathrm{proj}}.\end{array}$$
###### Proof.
The assertion for flat almost modules follows directly from lemmata 3.4.5, 3.4.6, 3.4.9 and 3.4.10. Set $`B=A_1\times A_2`$. To establish the second equivalence, it suffices to show that, if $`P`$ is an $`A_0`$-module such that $`B_{A_0}P`$ is almost projective over $`B`$, then $`P`$ is almost projective over $`A_0`$, or which is the same, that $`\text{alExt}_{A_0}^i(P,N)0`$ for all $`i>0`$ and any $`A_0`$-module $`N`$. We know already that $`P`$ is flat. Let $`M`$ be any $`A_0`$-module and $`N`$ any $`B`$-module. The standard isomorphism $`R\text{Hom}_B(B\stackrel{\text{L}}{}_{A_0}M,N)R\text{Hom}_{A_0}(M,N)`$ yields a natural isomorphism $`\text{alExt}_B^i(B_{A_0}M,N)\text{alExt}_{A_0}^i(M,N)`$, whenever $`\text{Tor}_j^{A_0}(B,M)=0`$ for every $`j>0`$. In particular, we have $`\text{alExt}_{A_0}^i(P,N)0`$ whenever $`N`$ comes from either an $`A_1`$-module, or an $`A_2`$-module. For a general $`A_0`$-module $`N`$ there is a 3-step filtration such that $`\text{Fil}_0(N)=0`$, $`\text{gr}_1(N)=\text{Fil}_1(N)=\text{Ker}(\epsilon _N)`$, $`\text{gr}_2(N)=\text{Ker}(A_1_{A_0}NA_3_{A_0}N)`$ and $`\text{gr}_3(N)=A_2_{A_0}N`$. By an easy devissage, we reduce to verify that $`\text{alExt}_{A_0}^i(P,\text{gr}_j(N))=0`$ for every $`i>0`$ and $`j=1,2,3`$. However, $`\text{gr}_2(N)`$ is an $`A_1`$-module and $`\text{gr}_3(N)`$ is an $`A_2`$-module, so the required vanishing follows for $`j=2,3`$. Moreover, applying $`_{A_0}N`$ to (3.4.7), we derive a short exact sequence :
(3.4.12)
$$0\text{Tor}_1^{A_0}(N,A_2)\frac{\text{Tor}_1^{A_0}(N,A_3)}{\text{Tor}_1^{A_0}(N,A_1)}\text{gr}_1(N)0.$$
Here again, the leftmost term of (3.4.12) is an $`A_2`$-module, and the middle term is an $`A_1`$-module, so the same devissage yields the sought vanishing also for $`j=1`$. ∎
###### Corollary 3.4.13.
In the situation of (3.4.4), denote by $`\text{A}_{i,\mathrm{fl}}`$ (resp. $`\stackrel{´}{𝐄}𝐭_i`$, resp. $`𝐰.\stackrel{´}{𝐄}𝐭_i`$) the category of flat (resp. étale, resp. weakly étale) $`A_i`$-algebras. The functor $`\pi `$ induces equivalences
$`\text{A}_{0,\mathrm{fl}}\stackrel{}{}\text{A}_{1,\mathrm{fl}}\times _{\text{A}_{3,\mathrm{fl}}}\text{A}_{2,\mathrm{fl}}\stackrel{´}{𝐄}𝐭_0\stackrel{}{}\stackrel{´}{𝐄}𝐭_1\times _{\stackrel{´}{𝐄}𝐭_3}\stackrel{´}{𝐄}𝐭_2𝐰.\stackrel{´}{𝐄}𝐭_0\stackrel{}{}𝐰.\stackrel{´}{𝐄}𝐭_1\times _{𝐰.\stackrel{´}{𝐄}𝐭_3}𝐰.\stackrel{´}{𝐄}𝐭_2.`$
Next we want to reinterpret the equivalences of proposition 3.4.11 in terms of descent data. If $`F:\text{C}V^a\text{-}\mathrm{𝐀𝐥𝐠}^o`$ is a fibred category over the opposite of the category of almost algebras, and if $`XY`$ is a given morphism of almost algebras, we shall denote by $`\mathrm{𝐃𝐞𝐬𝐜}(\text{C},Y/X)`$ the category of objects of the fibre category $`F_Y`$, endowed with a descent datum relative to the morphism $`XY`$ (cp. (Ch.II §1)). In the arguments hereafter, we consider morphisms of almost algebras and modules, and one has to reverse the direction of the arrows to pass to morphisms in the considered fibred category. Denote by $`p_i:YY_XY`$ ($`i=1,2`$), resp. $`p_{ij}:Y_XYY_XY_XY`$ ($`1i<j3`$) the usual morphisms. As an example, $`\mathrm{𝐃𝐞𝐬𝐜}(V^a\text{-}\mathrm{𝐀𝐥𝐠}.\mathrm{𝐌𝐨𝐝}^o,Y/X)`$ consists of the pairs $`(M,\beta )`$ where $`M`$ is a $`Y`$-module and $`\beta `$ is a $`Y_XY`$-linear isomorphism $`\beta :p_2^{}(M)\stackrel{}{}p_1^{}(M)`$ such that
(3.4.14)
$$p_{12}^{}(\beta )p_{23}^{}(\beta )=p_{13}^{}(\beta ).$$
Let now $`IX`$ be an ideal, and set $`\overline{X}=X/I`$, $`\overline{Y}=Y/IY`$. For any $`F:\text{C}V^a\text{-}\mathrm{𝐀𝐥𝐠}^o`$ as above, one has an essentially commutative diagram:
This induces a functor :
(3.4.15)
$$\mathrm{𝐃𝐞𝐬𝐜}(\text{C},Y/X)\mathrm{𝐃𝐞𝐬𝐜}(\text{C},\overline{Y}/\overline{X})\times _{F_{\overline{Y}}}F_Y.$$
###### Lemma 3.4.16.
With the above notation, suppose moreover that the natural morphism $`IIY`$ is an isomorphism. Then the functor (3.4.15) is an equivalence whenever C is one of the fibred categories $`V^a\text{-}\mathrm{𝐀𝐥𝐠}.\mathrm{𝐌𝐨𝐝}^o`$, $`V^a\text{-}\mathrm{𝐀𝐥𝐠}.\mathrm{𝐌𝐨𝐫𝐩𝐡}^o`$, $`\stackrel{´}{𝐄}𝐭^o`$, $`𝐰.\stackrel{´}{𝐄}𝐭^o`$.
###### Proof.
For any $`n>0`$, denote by $`Y^n`$ (resp. $`\overline{Y}^n`$) the $`n`$-fold tensor product of $`Y`$ (resp. $`\overline{Y}`$) with itself over $`X`$ (resp. $`\overline{X}`$), and by $`\rho _n:Y^n\overline{Y}^n`$ the natural morphism. First of all we claim that, for every $`n>0`$, the natural diagram of almost algebras
(3.4.17)
is cartesian (where $`\mu _n`$ and $`\overline{\mu }_n`$ are $`n`$-fold multiplication morphisms). For this, we need to verify that, for every $`n>0`$, the induced morphism $`\text{Ker}(\rho _n)\text{Ker}(\rho _1)`$ (defined by multiplication of the first two factors) is an isomorphism. It then suffices to check that the natural morphism $`\text{Ker}(\rho _n)\text{Ker}(\rho _{n1})`$ is an isomorphism for all $`n>1`$. Indeed, consider the commutative diagram
From $`IY=\varphi (Y)`$, it follows that $`p^{}`$ is an epimorphism. Hence also $`\psi `$ is an epimorphism. Since $`i`$ is a monomorphism, it follows that $`\psi `$ is also a monomorphism, hence $`\psi `$ is an isomorphism and the claim follows easily.
We consider first the case $`\text{C}=V^a\text{-}\mathrm{𝐀𝐥𝐠}.\mathrm{𝐌𝐨𝐝}^o`$; we see that (3.4.17) is a diagram of the kind considered in (3.4.4), hence, for every $`n>0`$, we have the associated functor $`\pi _n:Y^n\text{-}\mathrm{𝐌𝐨𝐝}\overline{Y}^n\text{-}\mathrm{𝐌𝐨𝐝}\times _{\overline{Y}\text{-}\mathrm{𝐌𝐨𝐝}}Y\text{-}\mathrm{𝐌𝐨𝐝}`$ and also its right adjoint $`T_n`$. Denote by $`\overline{p}_i:\overline{Y}\overline{Y}^2`$ ($`i=1,2`$) the usual morphisms, and similarly define $`\overline{p}_{ij}:\overline{Y}^2\overline{Y}^3`$. Suppose there is given a descent datum $`(\overline{M},\overline{\beta })`$ for $`\overline{M}`$, relative to $`\overline{X}\overline{Y}`$. The cocycle condition (3.4.14) implies easily that $`\overline{\mu }_2^{}(\overline{\beta })`$ is the identity on $`\overline{\mu }_2^{}(\overline{p}_i^{}(\overline{M}))=\overline{M}`$. It follows that the pair $`(\overline{\beta },\text{1}_M)`$ defines an isomorphism $`\pi _2(p_1^{}M)\stackrel{}{}\pi _2(p_2^{}M)`$ in the category $`\overline{Y}^2\text{-}\mathrm{𝐌𝐨𝐝}\times _{\overline{Y}\text{-}\mathrm{𝐌𝐨𝐝}}Y\text{-}\mathrm{𝐌𝐨𝐝}`$. Hence $`T_2(\overline{\beta },\text{1}_M):T_2\pi _2(p_1^{}M)T_2\pi _2(p_2^{}M)`$ is an isomorphism. However, we remark that either morphism $`\overline{p}_i`$ yields a section for $`\mu _2`$, hence we are in the situation contemplated in lemma 3.4.8, and we derive an isomorphism $`\beta :p_2^{}(M)\stackrel{}{}p_1^{}(M)`$. We claim that $`(M,\beta )`$ is an object of $`\mathrm{𝐃𝐞𝐬𝐜}(\text{C},Y/X)`$, i.e. that $`\beta `$ verifies the cocycle condition (3.4.14). Indeed, we can compute: $`\pi _3(p_{ij}^{}\beta )=(\rho _3^{}(p_{ij}^{}\beta ),\mu _3^{}(p_{ij}^{}\beta ))`$ and by construction we have $`\rho _3^{}(p_{ij}^{}\beta )=\overline{p}_{ij}^{}(\overline{\beta })`$ and $`\mu _3^{}(p_{ij}^{}\beta )=\mu _2^{}(\beta )=\text{1}_M`$. Therefore, the cocycle identity for $`\overline{\beta }`$ implies the equality $`\pi _3(p_{12}^{}(\beta ))\pi _3(p_{23}^{}(\beta ))=\pi _3(p_{13}^{}(\beta ))`$. If we now apply the functor $`T_3`$ to this equality, and then invoke again lemma 3.4.8, the required cocycle identity for $`\beta `$ will ensue. Clearly $`\beta `$ is the only descent datum on $`M`$ lifting $`\overline{\beta }`$. This proves that (3.4.15) is essentially surjective. The same sort of argument also shows that the functor (3.4.15) induces bijections on morphisms, so the lemma follows in this case. Next, the case $`\text{C}=V^a\text{-}\mathrm{𝐀𝐥𝐠}.\mathrm{𝐌𝐨𝐫𝐩𝐡}^o`$ can be deduced formally from the previous case, by applying repeatedly natural isomorphisms of the kind $`p_i^{}(M_YN)p_i^{}(M)_{Y_XY}p_i^{}(N)`$ ($`i=1,2`$). Finally, the “étaleness” of an object of $`\mathrm{𝐃𝐞𝐬𝐜}(V^a\text{-}\mathrm{𝐀𝐥𝐠}.\mathrm{𝐌𝐨𝐫𝐩𝐡}^o,Y/X)`$ can be checked on its projection onto $`Y\text{-}\mathrm{𝐀𝐥𝐠}^o`$, hence also the cases $`\text{C}=𝐰.\stackrel{´}{𝐄}𝐭^o`$ and $`\text{C}=\stackrel{´}{𝐄}𝐭^o`$ follow directly. ∎
Now, let $`B=A_1\times A_2`$; to an objet $`(M,\beta )`$ in $`\mathrm{𝐃𝐞𝐬𝐜}(V^a\text{-}\mathrm{𝐀𝐥𝐠}.\mathrm{𝐌𝐨𝐝}^o,B/A)`$ we assign an object $`(M_1,M_2,\xi )`$ of $`\text{M}_1\times _{\text{M}_3}\text{M}_2`$, as follows. Set $`M_i=A_i_BM`$ ($`i=1,2`$) and $`A_{ij}=A_i_{A_0}A_j`$. We can write $`B_{A_0}B=_{i,j=1}^2A_{ij}`$ and $`\beta `$ gives rise to the $`A_{ij}`$-linear isomorphisms $`\beta _{ij}:A_{ij}_{B_{A_0}B}p_2^{}(M)\stackrel{}{}A_{ij}_{B_{A_0}B}p_1^{}(M)`$. In other words, we obtain isomorphisms $`\beta _{ij}:A_i_{A_0}M_jM_i_{A_0}A_j`$. However, we have a natural isomorphism $`A_{12}A_3`$ (indeed, suppose that $`A_1A_3`$ is an epimorphism with kernel $`I`$; then $`I`$ is also an ideal of $`A_0`$ and $`A_0/IA_2`$; now the claim follows by remarking that $`IA_1=I`$). Hence we can choose $`\xi =\beta _{12}`$. In this way we obtain a functor :
(3.4.18)
$$\mathrm{𝐃𝐞𝐬𝐜}(V^a\text{-}\mathrm{𝐀𝐥𝐠}.\mathrm{𝐌𝐨𝐝}^o,B/A_0)(\text{M}_1\times _{\text{M}_3}\text{M}_2)^o.$$
###### Proposition 3.4.19.
The functor (3.4.18) is an equivalence of categories.
###### Proof.
Let us say that $`A_1A_3`$ is an epimorphism with kernel $`I`$. Then $`I`$ is also an ideal of $`B`$ and we have $`B/IA_3\times A_2`$ and $`A_0/IA_2`$. We intend to apply lemma 3.4.16 to the morphism $`A_0B`$. However, the induced morphism $`\overline{B}=B/I\overline{A}_0=A_0/I`$ in $`V^a\text{-}\mathrm{𝐀𝐥𝐠}^o`$ has a section, and hence it is of universal effective descent for every fibred category. Thus, we can replace in (3.4.15) the category $`\mathrm{𝐃𝐞𝐬𝐜}(V^a\text{-}\mathrm{𝐀𝐥𝐠}.\mathrm{𝐌𝐨𝐝}^o,\overline{B}/\overline{A}_0)`$ by $`\overline{A}_0\text{-}\mathrm{𝐌𝐨𝐝}^o`$, and thereby, identify (up to equivalence) the target of (3.4.15) with the 2-fibred product $`(\text{M}_1\times \text{M}_2)^o\times _{(\text{M}_3\times \text{M}_2)^o}\text{M}_2^o`$. The latter is equivalent to the category $`\text{M}_1^o\times _{\text{M}_3^o}\text{M}_2^o`$ and the resulting functor $`\mathrm{𝐃𝐞𝐬𝐜}(V^a\text{-}\mathrm{𝐀𝐥𝐠}.\mathrm{𝐌𝐨𝐝}^o,B/A_0)\text{M}_1^o\times _{\text{M}_3^o}\text{M}_2^o`$ is canonically isomorphic to (3.4.18), which gives the claim. ∎
Putting together propositions 3.4.11 and 3.4.19 we obtain the following :
###### Corollary 3.4.20.
In the situation of (3.4.4), the morphism $`A_0A_1\times A_2`$ is of effective descent for the fibred categories of flat almost modules and of almost projective almost modules. ∎
Next we would like to give sufficient conditions to ensure that a morphism of almost algebras is of effective descent for the fibred category $`𝐰.\stackrel{´}{𝐄}𝐭^oV^a\text{-}\mathrm{𝐀𝐥𝐠}^o`$ of weakly étale algebras (resp. for étale algebras). To this aim we are led to the following :
###### Definition 3.4.21.
A morphism $`\varphi :AB`$ of almost algebras is said to be strictly finite if $`\text{Ker}(\varphi )`$ is nilpotent and $`BR^a`$, where $`R`$ is a finite $`A_{}`$-algebra.
###### Theorem 3.4.22.
Let $`\varphi :AB`$ be a strictly finite morphism of almost algebras. Then :
i) For every $`A`$-algebra $`C`$, the induced morphism $`CC_AB`$ is again strictly finite.
ii) If $`M`$ is a flat $`A`$-module and $`B_AM`$ is almost projective over $`B`$, then $`M`$ is almost projective over $`A`$.
iii) $`AB`$ is of universal effective descent for the fibred categories of weakly étale (resp. étale) almost algebras.
###### Proof.
(i): suppose that $`B=R^a`$ for a finite $`A_{}`$-algebra $`R`$; then $`S=C_{}_A_{}R`$ is a finite $`C_{}`$-algebra and $`S^aC`$. It remains to show that $`\text{Ker}(CC_AB)`$ is nilpotent. Suppose that $`R`$ is generated by $`n`$ elements as an $`A_{}`$-module and let $`F_A_{}(R)`$ (resp. $`F_C_{}(S)`$) be the Fitting ideal of $`R`$ (resp.of $`S`$); we have $`\text{Ann}_C_{}(S)^nF_C_{}(S)\text{Ann}_C_{}(S)`$ (see (Chap.XIX Prop.2.5)); on the other hand $`F_C_{}(S)=F_A_{}(R)C_{}`$, so the claim is clear.
(iii): we shall consider the fibred category $`F:𝐰.\stackrel{´}{𝐄}𝐭^oV^a\text{-}\mathrm{𝐀𝐥𝐠}^o`$; the same argument applies also to étale almost algebras. We begin by establishing a very special case :
###### Claim 3.4.23.
Assertion (iii) holds when $`B=(A/I_1)\times (A/I_2)`$, where $`I_1`$ and $`I_2`$ are ideals in $`A`$ such that $`I_1I_2`$ is nilpotent.
Proof of the claim: First of all we remark that the situation considered in the claim is stable under arbitrary base change, therefore it suffices to show that $`\varphi `$ is of $`F`$-2-descent in this case. Then we factor $`\varphi `$ as a composition $`AA/\text{Ker}(\varphi )B`$ and we remark that $`AA/\text{Ker}(\varphi )`$ is of $`F`$-2-descent by theorem 3.3.12; since a composition of morphisms of $`F`$-2-descent is again of $`F`$-2-descent, we are reduced to show that $`A/\text{Ker}(\varphi )B`$ is of $`F`$-2-descent, i.e. we can assume that $`\text{Ker}(\varphi )0`$. However, in this case the claim follows easily from corollary 3.4.20.
###### Claim 3.4.24.
More generally, assertion (iii) holds when $`B=_{i=1}^nA/I_i`$, where $`I_1,\mathrm{},I_n`$ are ideals of $`A`$, such that $`_{i=1}^nI_i`$ is nilpotent.
Proof of the claim: We prove this by induction on $`n`$, the case $`n=2`$ being covered by claim 3.4.23. Therefore, suppose that $`n>2`$, and set $`B^{}=A/(_{i=1}^{n1}I_j)`$. By induction, the morphism $`B^{}_{i=1}^{n1}A/I_i`$ is of universal $`F`$-2-descent. However, according to (Chap.II Prop.1.1.3), the sieves of universal $`F`$-2-descent form a topology on $`V^a\text{-}\mathrm{𝐀𝐥𝐠}^o`$; for this topology, $`\{A,B\}`$ is a covering family of $`A\times B`$ and $`(AB^{}\times (A/I_n))^o`$ is a covering morphism, hence $`\{B^{},A/I_n\}`$ is a covering family of $`A`$, and then, by composition of covering families, $`\{_{i=1}^{n1}A/I_i,A/I_n\}`$ is a covering family of $`A`$, which is equivalent to the claim.
Now, let $`AB`$ be a general strictly finite morphism, so that $`B=R^a`$ for some finite $`A_{}`$-algebra $`R`$. Pick generators $`f_1,\mathrm{},f_m`$ of the $`A_{}`$-module $`R`$, and monic polynomials $`p_1(X),\mathrm{},p_m(X)`$ such that $`p_i(f_i)=0`$ for $`i=1,\mathrm{},m`$.
###### Claim 3.4.25.
There exists a finite and faithfully flat extension $`C`$ of $`A_{}`$ such that the images in $`C[X]`$ of $`p_1(X)`$,…,$`p_m(X)`$ split as products of monic linear factors.
Proof of the claim: This extension $`C`$ can be obtained as follows. It suffices to find, for each $`i=1,\mathrm{},m`$, an extension $`C_i`$ that splits $`p_i(X)`$, because then $`C=C_1_A_{}\mathrm{}_A_{}C_m`$ will split them all, so we can assume that $`m=1`$ and $`p_1(X)=p(X)`$; moreover, by induction on the degree of $`p(X)`$, it suffices to find an extension $`C^{}`$ such that $`p(X)`$ factors in $`C^{}[X]`$ as a product of the form $`p(X)=(X\alpha )q(X)`$, where $`q(X)`$ is a monic polynomial of degree $`\mathrm{deg}(p)1`$. Clearly we can take $`C^{}=A_{}[T]/(p(T))`$.
Given a $`C`$ as in claim 3.4.25, we remark that the morphism $`AC^a`$ is of universal $`F`$-2-descent. Considering again the topology of universal $`F`$-2-descent, it follows that $`AB`$ is of universal $`F`$-2-descent if and only if the same holds for the induced morphism $`C^aC^a_AB`$. Therefore, in proving assertion (iii) we can replace $`\varphi `$ by $`\text{1}_C_A\varphi `$ and assume from start that the polynomials $`p_i(X)`$ factor in $`A_{}[X]`$ as product of linear factors. Now, let $`\mathrm{deg}(p_i)=d_i`$ and $`p_i(X)=_j^{d_i}(X\alpha _{ij})`$ (for $`i=1,\mathrm{},m`$). We get a surjective homomorphism of $`A_{}`$-algebras $`D=A_{}[X_1,\mathrm{},X_m]/(p_1(X_1),\mathrm{},p_m(X_m))R`$ by the rule $`X_if_i`$ ($`i=1,\mathrm{},m`$). Moreover, any sequence $`\underset{¯}{\alpha }=(\alpha _{1,j_1},\alpha _{2,j_2},\mathrm{},\alpha _{m,j_m})`$ yields a homomorphism $`\psi _{\underset{¯}{\alpha }}:DA_{}`$, determined by the assignment $`X_i\alpha _{i,j_i}`$. A simple combinatorial argument shows that $`_{\underset{¯}{\alpha }}\text{Ker}(\psi _{\underset{¯}{\alpha }})=0`$, where $`\underset{¯}{\alpha }`$ runs over all the sequences as above. Hence the product map $`_{\underset{¯}{\alpha }}\psi _{\underset{¯}{\alpha }}:D_{\underset{¯}{\alpha }}A_{}`$ has nilpotent kernel. We notice that the $`A_{}`$-algebra $`(_{\underset{¯}{\alpha }}A_{})_DR`$ is a quotient of $`_{\underset{¯}{\alpha }}A_{}`$, hence it can be written as a product of rings of the form $`A_{}/I_{\underset{¯}{\alpha }}`$, for various ideals $`I_{\underset{¯}{\alpha }}`$. By (i), the kernel of the induced homomorphism $`R_{\underset{¯}{\alpha }}A_{}/I_{\underset{¯}{\alpha }}`$ is nilpotent, hence the same holds for the kernel of the composition $`A_{\underset{¯}{\alpha }}A/I_{\underset{¯}{\alpha }}^a`$, which is therefore of the kind considered in claim 3.4.24. Hence $`A_{\underset{¯}{\alpha }}A/I_{\underset{¯}{\alpha }}^a`$ is of universal $`F`$-2-descent. Since such morphisms form a topology, it follows that also $`AB`$ is of universal $`F`$-2-descent, which concludes the proof of (iii).
Finally, let $`M`$ be as in (ii) and pick again $`C`$ as in the proof of claim 3.4.25. By remark 3.3.17(iv), $`M`$ is almost projective over $`A`$ if and only if $`C^a_AM`$ is almost projective over $`C^a`$; hence we can replace $`\varphi `$ by $`\text{1}_{C^a}_A\varphi `$, and by arguing as in the proof of (iii), we can assume from start that $`B=_{j=1}^n(A/I_j)`$ for ideals $`I_jA`$, $`j=1,\mathrm{},n`$ such that $`I=_{j=1}^nI_j`$ is nilpotent. By an easy induction, we can furthermore reduce to the case $`n=2`$. We factor $`\varphi `$ as $`AA/IB`$; by proposition 3.4.11 it follows that $`(A/I)_AM`$ is almost projective over $`A/I`$, and then lemma 3.3.16(i) says that $`M`$ itself is almost projective. ∎
###### Remark 3.4.26.
It is natural to ask whether theorem 3.4.22 holds if we replace everywhere “strictly finite” by “finite with nilpotent kernel” (or even by “almost finite with nilpotent kernel”). We do not know the answer to this question.
We conclude with a digression to explain the relationship between our results and related facts that can be extracted from the literature. So, we now place ourselves in the “classical limit” $`V=𝔪`$ (cp. example 2.1.1(ii)). In this case, weakly étale morphisms had already been considered in some earlier work, and they were called “absolutely flat” morphisms. A ring homomorphism $`AB`$ is étale in the usual sense of if and only if it is absolutely flat and of finite presentation. Let us denote by $`𝐮.\stackrel{´}{𝐄}𝐭^o`$ the fibred category over $`V\text{-}\mathrm{𝐀𝐥𝐠}^o`$, whose fibre over a $`V`$-algebra $`A`$ is the opposite of the category of étale $`A`$-algebras in the usual sense. We claim that, if a morphism $`AB`$ of $`V`$-algebras is of universal effective descent for the fibred category $`𝐰.\stackrel{´}{𝐄}𝐭^o`$ (resp. $`\stackrel{´}{𝐄}𝐭^o`$), then it is a morphism of universal effective descent for $`𝐮.\stackrel{´}{𝐄}𝐭^o`$. Indeed, let $`C`$ be an étale $`A`$-algebra (in the sense of definition 3.1.1) and such that $`C_AB`$ is étale over $`B`$ in the usual sense. We have to show that $`C`$ is étale in the usual sense, i.e. that it is of finite presentation over $`A`$. This amounts to showing that, for every filtered inductive system $`(A_\lambda )_{\lambda \mathrm{\Lambda }}`$ of $`A`$-algebras, we have $`\underset{\lambda \mathrm{\Lambda }}{\text{colim}}\text{Hom}_{A\text{-}\mathrm{𝐀𝐥𝐠}}(C,A_\lambda )\text{Hom}_{A\text{-}\mathrm{𝐀𝐥𝐠}}(C,\underset{\lambda \mathrm{\Lambda }}{\text{colim}}A_\lambda )`$. Since, by assumption, this is known after extending scalars to $`B`$ and to $`B_AB`$, it suffices to show that, for any $`A`$-algebra $`D`$, the natural sequence
is exact. For this, note that $`\text{Hom}_{A\text{-}\mathrm{𝐀𝐥𝐠}}(C,D)=\text{Hom}_{D\text{-}\mathrm{𝐀𝐥𝐠}}(C_D,D)`$ (and similarly for the other terms) and by hypothesis $`(DD_AB)^o`$ is a morphism of 1-descent for the fibred category $`𝐰.\stackrel{´}{𝐄}𝐭^o`$ (resp. $`\stackrel{´}{𝐄}𝐭^o`$).
As a consequence of these observations and of theorem 3.4.22, we see that any finite ring homomorphism $`\varphi :AB`$ with nilpotent kernel is of universal effective descent for the fibred category of étale algebras. This fact was known as follows. By (Exp.IX, 4.7), $`\text{Spec}(\varphi )`$ is of universal effective descent for the fibred category of separated étale morphisms of finite type. One has to show that if $`X`$ is such a scheme over $`A`$, such that $`X_AB`$ is affine, then $`X`$ is affine. This follows by reduction to the noetherian case and (Chap.II, 6.7.1).
### 3.5. Behaviour of étale morphisms under Frobenius
We consider the following category B of base rings. The objects of B are the pairs $`(V,𝔪)`$, where $`V`$ is a ring and $`𝔪`$ is an ideal of $`V`$ with $`𝔪=𝔪^2`$ and $`\stackrel{~}{𝔪}`$ is flat. The morphisms $`(V,𝔪_V)(W,𝔪_W)`$ between two objects of B are the ring homomorphisms $`f:VW`$ such that $`𝔪_W=f(𝔪_V)W`$.
We have a fibred and cofibred category $`\text{B}\text{-}\mathrm{𝐌𝐨𝐝}\text{B}`$ (see (Exp.VI §5,6,10) for generalities on fibred categories). An object of $`\text{B}\text{-}\mathrm{𝐌𝐨𝐝}`$ (which we may call a “B-module”) consists of a pair $`((V,𝔪),M)`$, where $`(V,𝔪)`$ is an object of B and $`M`$ is a $`V`$-module. Given two objects $`X=((V,𝔪_V),M)`$ and $`Y=((W,𝔪_W),N)`$, the morphisms $`XY`$ are the pairs $`(f,g)`$, where $`f:(V,𝔪_V)(W,𝔪_W)`$ is a morphism in B and $`g:MN`$ is an $`f`$-linear map.
Similarly one has a fibred and cofibred category $`\text{B}\text{-}\mathrm{𝐀𝐥𝐠}\text{B}`$ of B-algebras. We will also need to consider the fibred and cofibred category $`\text{B}\text{-}\mathrm{𝐌𝐨𝐧}\text{B}`$ of non-unitary commutative B-monoids: an object of $`\text{B}\text{-}\mathrm{𝐌𝐨𝐧}`$ is a pair $`((V,𝔪),A)`$ where $`A`$ is a $`V`$-module endowed with a morphism $`A_VAA`$ subject to associativity and commutativity conditions, as discussed in section 2.2. The fibre over an object $`(V,𝔪)`$ of B, is the category of $`V`$-monoids denoted $`(V,𝔪)\text{-}\mathrm{𝐌𝐨𝐧}`$ or simply $`V\text{-}\mathrm{𝐌𝐨𝐧}`$.
The almost isomorphisms in the fibres of $`\text{B}\text{-}\mathrm{𝐌𝐨𝐝}\text{B}`$ give a multiplicative system $`\mathrm{\Sigma }`$ in $`\text{B}\text{-}\mathrm{𝐌𝐨𝐝}`$, admitting a calculus of both left and right fractions. The “locally small” conditions are satisfied (see p.381), so that one can form the localised category $`\text{B}^a\text{-}\mathrm{𝐌𝐨𝐝}=\mathrm{\Sigma }^1(\text{B}\text{-}\mathrm{𝐌𝐨𝐝})`$. The fibres of the localised category over the objects of B are the previously considered categories of almost modules. Similar considerations hold for $`\text{B}\text{-}\mathrm{𝐀𝐥𝐠}`$ and $`\text{B}\text{-}\mathrm{𝐌𝐨𝐧}`$, and we get the fibred and cofibred categories $`\text{B}^a\text{-}\mathrm{𝐌𝐨𝐝}\text{B}`$, $`\text{B}^a\text{-}\mathrm{𝐀𝐥𝐠}\text{B}`$ and $`\text{B}^a\text{-}\mathrm{𝐌𝐨𝐧}\text{B}`$. In particular, for every object $`(V,𝔪)`$ of B, we have an obvious notion of almost $`V`$-monoid and the category consisting of these is denoted $`V^a\text{-}\mathrm{𝐌𝐨𝐧}`$. The localisation functors
$`\text{B}\text{-}\mathrm{𝐌𝐨𝐝}\text{B}^a\text{-}\mathrm{𝐌𝐨𝐝}:MM^a\text{B}\text{-}\mathrm{𝐀𝐥𝐠}\text{B}^a\text{-}\mathrm{𝐀𝐥𝐠}:BB^a`$
have left and right adjoints. These adjoints can be chosen as functors of categories over B such that the adjunction units and counits are morphisms over identity arrows in B. On the fibres these induce the previously considered left and right adjoints $`MM_!`$, $`MM_{}`$, $`BB_{!!}`$, $`BB_{}`$. We will use the same notation for the corresponding functors on the larger categories. Then it is easy to check that the functor $`MM_!`$ is cartesian and cocartesian (i.e. it sends cartesian arrows to cartesian arrows and cocartesian arrows to cocartesian arrows), $`MM_{}`$ and $`BB_{}`$ are cartesian, and $`BB_{!!}`$ is cocartesian.
Let $`\text{B}/𝔽_p`$ be the full subcategory of B consisting of all objects $`(V,𝔪)`$ where $`V`$ is an $`𝔽_p`$-algebra. Define similarly $`\text{B}\text{-}\mathrm{𝐀𝐥𝐠}/𝔽_p`$, $`\text{B}\text{-}\mathrm{𝐌𝐨𝐧}/𝔽_p`$ and $`\text{B}^a\text{-}\mathrm{𝐀𝐥𝐠}/𝔽_p`$, $`\text{B}^a\text{-}\mathrm{𝐌𝐨𝐧}/𝔽_p`$, so that we have again fibred and cofibred categories $`\text{B}^a\text{-}\mathrm{𝐀𝐥𝐠}/𝔽_p\text{B}/𝔽_p`$ and $`\text{B}^a\text{-}\mathrm{𝐀𝐥𝐠}/𝔽_p\text{B}/𝔽_p`$ (resp. the same for non-unitary monoids). We remark that the categories $`\text{B}^a\text{-}\mathrm{𝐀𝐥𝐠}/𝔽_p`$ and $`\text{B}^a\text{-}\mathrm{𝐌𝐨𝐧}/𝔽_p`$ have small limits and colimits, and these are preserved by the projection to $`\text{B}/𝔽_p`$. Especially, if $`AB`$ and $`AC`$ are two morphisms in $`\text{B}^a\text{-}\mathrm{𝐀𝐥𝐠}/𝔽_p`$ or $`\text{B}^a\text{-}\mathrm{𝐌𝐨𝐧}/𝔽_p`$, we can define $`B_AC`$ as such a colimit.
If $`A`$ is a (unitary or non-unitary) B-monoid over $`𝔽_p`$, we denote by $`\varphi _A:AA`$ the Frobenius endomorphism $`xx^p`$. If $`(V,𝔪)`$ is an object of $`\text{B}/𝔽_p`$, it follows from proposition 2.1.5(ii) that $`\varphi _V:(V,𝔪)(V,𝔪)`$ is a morphism in B. If $`B`$ is an object of $`\text{B}\text{-}\mathrm{𝐀𝐥𝐠}/𝔽_p`$ (resp. $`\text{B}\text{-}\mathrm{𝐌𝐨𝐧}/𝔽_p`$) over $`V`$, then the Frobenius map induces a morphism $`\varphi _B:BB`$ in $`\text{B}\text{-}\mathrm{𝐀𝐥𝐠}/𝔽_p`$ (resp. $`\text{B}\text{-}\mathrm{𝐌𝐨𝐧}/𝔽_p`$) over $`\varphi _V`$. In this way we get a natural transformation from the identity functor of $`\text{B}\text{-}\mathrm{𝐀𝐥𝐠}/𝔽_p`$ (resp. $`\text{B}\text{-}\mathrm{𝐌𝐨𝐧}/𝔽_p`$) to itself that induces a natural transformation on the identity functor of $`\text{B}^a\text{-}\mathrm{𝐀𝐥𝐠}/𝔽_p`$ (resp. $`\text{B}^a\text{-}\mathrm{𝐌𝐨𝐧}/F_p`$).
Using the pull-back functors, any object $`B`$ of $`\text{B}\text{-}\mathrm{𝐀𝐥𝐠}`$ over $`V`$ defines new objects $`B_{(m)}`$ of $`\text{B}\text{-}\mathrm{𝐀𝐥𝐠}`$ ($`m`$) over $`V`$, where $`B_{(m)}=(\varphi _V^m)^{}(B)`$, which is just $`B`$ considered as a $`V`$-algebra via the homomorphism $`V\stackrel{\varphi ^m}{}VB`$. These operations also induce functors $`BB_{(m)}`$ on almost B-algebras.
###### Definition 3.5.1.
i) Let $`(V,𝔪)`$ be an object of $`\text{B}/𝔽_p`$; we say that a morphism $`f:AB`$ of almost $`V`$-algebras (resp. almost $`V`$-monoids) is invertible up to $`\varphi ^m`$ if there exists a morphism $`f^{}:BA`$ in $`\text{B}^a\text{-}\mathrm{𝐀𝐥𝐠}`$ (resp. $`\text{B}^a\text{-}\mathrm{𝐌𝐨𝐧}`$) over $`\varphi _V^m`$, such that $`f^{}f=\varphi _A^m`$ and $`ff^{}=\varphi _B^m`$.
ii) We say that an almost $`V`$-monoid $`I`$ (e.g. an ideal in a $`V^a`$-algebra) is Frobenius nilpotent if $`\varphi _I`$ is nilpotent.
Notice that a morphism $`f`$ of $`V^a\text{-}\mathrm{𝐀𝐥𝐠}`$ (or $`V^a\text{-}\mathrm{𝐌𝐨𝐧}`$) is invertible up to $`\varphi ^m`$ if and only if $`f_{}:A_{}B_{}`$ is so as a morphism of $`𝔽_p`$-algebras.
###### Lemma 3.5.2.
Let $`(V,𝔪)`$ be an object of $`\text{B}/𝔽_p`$ and let $`f:AB`$, $`g:BC`$ be morphisms of almost $`V`$-algebras or almost $`V`$-monoids.
i) If $`f`$ is invertible up to $`\varphi ^n`$ and $`g`$ is invertible up to $`\varphi ^m`$, then $`gf`$ is invertible up to $`\varphi ^{m+n}`$.
ii) If $`f`$ is invertible up to $`\varphi ^n`$ and $`gf`$ is invertible up to $`\varphi ^m`$, then $`g`$ is invertible up to $`\varphi ^{m+n}`$.
iii) If $`g`$ is invertible up to $`\varphi ^n`$ and $`gf`$ is invertible up to $`\varphi ^m`$, then $`f`$ is invertible up to $`\varphi ^{m+n}`$.
iv) The Frobenius morphisms induce $`\varphi _V`$-linear morphisms (i.e. morphisms in $`\text{B}^a\text{-}\mathrm{𝐌𝐨𝐝}`$ over $`\varphi _V`$) $`\varphi ^{}:\text{Ker}(f)\text{Ker}(f)`$ and $`\varphi ^{\prime \prime }:\text{Coker}(f)\text{Coker}(f)`$, and $`f`$ is invertible up to some power of $`\varphi `$ if and only if both $`\varphi ^{}`$ and $`\varphi ^{\prime \prime }`$ are nilpotent.
v) Consider a map of short exact sequences of almost $`V`$-monoids :
and suppose that two of the morphisms $`f^{},f,f^{\prime \prime }`$ are invertible up to a power of $`\varphi `$. Then also the third morphism has this property.
###### Proof.
(i): if $`f^{}`$ is an inverse of $`f`$ up to $`\varphi ^n`$ and $`g^{}`$ is an inverse of $`g`$ up to $`\varphi ^m`$, then $`f^{}g^{}`$ is an inverse of $`gf`$ up to $`\varphi ^{m+n}`$. (ii): given an inverse $`f^{}`$ of $`f`$ up to $`\varphi ^n`$ and an inverse $`h^{}`$ of $`h=gf`$ up to $`\varphi ^m`$, let $`g^{}=\varphi _B^nfh^{}`$. We compute :
$$\begin{array}{cc}gg^{}=\hfill & g\varphi _B^nfh^{}=\varphi _C^ngfh=\varphi _C^n\varphi _C^m\hfill \\ g^{}g=\hfill & \varphi _B^nfh^{}g=fh^{}g\varphi _B^n=fh^{}gff^{}\hfill \\ \hfill =& f\varphi _A^mf^{}=\varphi _B^mff^{}=\varphi _B^m\varphi _B^n.\hfill \end{array}$$
(iii) is similar and (iv) is an easy diagram chasing left to the reader. (v) follows from (iv) and the snake lemma. ∎
###### Lemma 3.5.3.
Let $`(V,𝔪)`$ be an object of $`\text{B}/𝔽_p`$.
(i) If $`f:AB`$ is a morphism of almost $`V`$-algebras, invertible up to $`\varphi ^n`$, then so is $`A^{}A^{}_AB`$ for every morphism $`AA^{}`$ of almost algebras.
(ii) If $`f:(V,𝔪_V)(W,𝔪_W)`$ is a morphism in $`\text{B}/𝔽_p`$, the functors $`f_{}:(V,𝔪_V)^a\text{-}\mathrm{𝐀𝐥𝐠}(W,𝔪_W)^a\text{-}\mathrm{𝐀𝐥𝐠}`$ and $`f^{}:(W,𝔪_W)^a\text{-}\mathrm{𝐀𝐥𝐠}(V,𝔪_V)^a\text{-}\mathrm{𝐀𝐥𝐠}`$ preserve the class of morphisms invertible up to $`\varphi ^n`$.
###### Proof.
(i): given $`f^{}:BA_{(m)}`$, construct a morphism $`A^{}_ABA_{(m)}^{}`$ using the morphism $`A^{}A_{(m)}^{}`$ coming from $`\varphi _A^{}^m`$ and $`f^{}`$. (ii): the assertion for $`f^{}`$ is clear, and the assertion for $`f_{}`$ follows from (i). ∎
###### Remark 3.5.4.
Statements like those of lemma 3.5.3 hold for the classes of flat, (weakly) unramified, (weakly) étale morphisms.
###### Theorem 3.5.5.
Let $`(V,𝔪)`$ be an object of $`\text{B}/𝔽_p`$ and $`f:AB`$ a weakly étale morphism of almost $`V`$-algebras.
(i) If $`f`$ is invertible up to $`\varphi ^n`$ ($`n0`$), then it is an isomorphism.
(ii) For every integer $`m0`$ the natural square diagram
(3.5.6)
is cocartesian.
###### Proof.
(i): we first show that $`f`$ is faithfully flat. Since $`f`$ is flat, it remains to show that if $`M`$ is an $`A`$-module such that $`M_AB=0`$, then $`M=0`$. It suffice to do this for $`M=A/I`$, for an arbitrary ideal $`I`$ of $`A`$. After base change by $`AA/I`$, we reduce to show that $`B=0`$ implies $`A=0`$. However, $`A_{}B_{}`$ is invertible up to $`\varphi ^n`$, so $`\varphi _A_{}^n=0`$ which means $`A_{}=0`$. In particular, $`f`$ is a monomorphism, hence the proof is complete in case that $`f`$ is an epimorphism. In general, consider the composition $`B\stackrel{\text{1}_Bf}{}B_AB\stackrel{\mu _{B/A}}{}B`$. From lemma 3.5.3(i) it follows that $`\text{1}_Bf`$ is invertible up to $`\varphi ^n`$; then lemma 3.5.2(ii) says that $`\mu _{B/A}`$ is invertible up to $`\varphi ^n`$. The latter is also weakly étale; by the foregoing we derive that it is an isomorphism. Consequently $`\text{1}_Bf`$ is an isomorphism, and finally, by faithful flatness, $`f`$ itself is an isomorphism.
(ii): the morphisms $`\varphi _A^m`$ and $`\varphi _B^m`$ are invertible up to $`\varphi ^m`$. By lemma 3.5.3(i) it follows that $`\text{1}_B\varphi _A^m:BB_AA_{(m)}`$ is invertible up to $`\varphi ^m`$; hence, by lemma 3.5.2(ii), the morphism $`h:B_AA_{(m)}B_{(m)}`$ induced by $`\varphi _B^m`$ and $`f_{(m)}`$ is invertible up to $`\varphi ^{2m}`$ (in fact one verifies that it is invertible up to $`\varphi ^m`$). But $`h`$ is a morphism of weakly étale $`A_{(m)}`$-algebras, so it is weakly étale, so it is an isomorphism by (i). ∎
###### Remark 3.5.7.
Theorem 3.5.5(ii) extends a statement of Faltings ( p.10) for his notion of almost étale extensions.
We recall (cp. (Chap.0, 3.5)) that a morphism $`f:XY`$ of objects in a site is called bicovering if the induced map of associated sheaves of sets is an isomorphism; if $`f`$ is squarable (“quarrable” in French), this is equivalent to the condition that both $`f`$ and the diagonal morphism $`XX\times _YX`$ are covering morphisms.
Let $`FE`$ be a fibered category and $`f:PQ`$ a squarable morphism of $`E`$. Consider the following condition:
(3.5.8) for every base change $`P\times _QQ^{}Q^{}`$ of $`f`$, the inverse image functor $`F_Q^{}F_{P\times _QQ^{}}`$ is an equivalence of categories.
Inspecting the arguments in (Chap.II,§1.1) one can show:
###### Lemma 3.5.9.
With the above notation, let $`\tau `$ be the topology of universal effective descent relative to $`FE`$. Then we have :
i) if (3.5.8) holds, then $`f`$ is a covering morphism for the topology $`\tau `$;
ii) $`f`$ is bicovering for $`\tau `$ if and only if (3.5.8) holds both for $`f`$ and for the diagonal morphism $`PP\times _QP`$.
###### Remark 3.5.10.
In (Chap.II, 1.1.3(iv)) it is stated that “la réciproque est vraie si $`i=2`$”, meaning that (3.5.8) is equivalent to the condition that $`f`$ is bicovering for $`\tau `$. (Actually the cited statement is given in terms of presheaves, but one can show that (3.5.8) is equivalent to the corresponding condition for the fibered category $`F^+\widehat{E}_U`$ considered in op.cit.) However, this fails in general : as a counterexample we can give the following. Let $`E`$ be the category of schemes of finite type over a field $`k`$; set $`P=𝔸_k^1`$, $`Q=\text{Spec}(k)`$. Finally let $`FE`$ be the discretely fibered category defined by the presheaf $`XH^0(X,)`$. Then it is easy to show that $`f`$ satisfies (3.5.8) but the diagonal map does not, so $`f`$ is not bicovering. The mistake in the proof is in (Chap.II, 1.1.3.5), where one knows that $`F^+(d)`$ is an equivalence of categories (notation of loc.cit.) but one needs it also after base changes of $`d`$.
###### Lemma 3.5.11.
(i) If $`f:AB`$ is a morphism of $`V^a`$-algebras which is invertible up to $`\varphi ^m`$, then the induced functors $`\stackrel{´}{𝐄}𝐭(A)\stackrel{´}{𝐄}𝐭(B)`$ and $`𝐰.\stackrel{´}{𝐄}𝐭(A)𝐰.\stackrel{´}{𝐄}𝐭(B)`$ are equivalences of categories.
ii) If $`AB`$ is weakly étale and $`CD`$ is a morphism of $`A`$-algebras invertible up to $`\varphi ^m`$, then the induced map $`\text{Hom}_{A\text{-}\mathrm{𝐀𝐥𝐠}}(B,C)\text{Hom}_{A\text{-}\mathrm{𝐀𝐥𝐠}}(B,D)`$ is bijective.
###### Proof.
We first consider (i) for the special case where $`f=\varphi _A^m:AA_{(m)}`$. The functor $`(\varphi _V^m)^{}:V^a\text{-}\mathrm{𝐀𝐥𝐠}V^a\text{-}\mathrm{𝐀𝐥𝐠}`$ induces a functor $`()_{(m)}:A\text{-}\mathrm{𝐀𝐥𝐠}A_{(m)}\text{-}\mathrm{𝐀𝐥𝐠}`$, and by restriction (see remark 3.5.3) we obtain a functor $`()_{(m)}:\stackrel{´}{𝐄}𝐭(A)\stackrel{´}{𝐄}𝐭(A_{(m)})`$; by theorem 3.5.5(ii), the latter is isomorphic to the functor $`(\varphi ^m)_{}:\stackrel{´}{𝐄}𝐭(A)\stackrel{´}{𝐄}𝐭(A_{(m)})`$ of the lemma. Furthermore, from remark 2.1.3(ii) and (2.2.2) we derive a natural ring isomorphism $`\omega :A_{(m)}A_{}`$, hence an essentially commutative diagram
where $`\alpha `$ and $`\beta `$ are the equivalences of remark 2.2.3. Clearly $`\alpha `$ and $`\beta `$ restrict to equivalences on the corresponding categories of étale algebras, hence the lemma follows in this case.
For the general case of (i), let $`f^{}:BA_{(m)}`$ be a morphism as in definition 3.5.1. Diagram (3.5.6) induces an essentially commutative diagram of the corresponding categories of algebras, so by the previous case, the functor $`(f^{})_{}:\stackrel{´}{𝐄}𝐭(B)\stackrel{´}{𝐄}𝐭(A_{(m)})`$ has both a left essential inverse and a right essential inverse; these essential inverses must be isomorphic, so $`f_{}`$ has an essential inverse as desired. Finally, we remark that the map in (ii) is the same as the map $`\text{Hom}_{C\text{-}\mathrm{𝐀𝐥𝐠}}(B_AC,C)\text{Hom}_{D\text{-}\mathrm{𝐀𝐥𝐠}}(B_AD,D)`$, and the latter is a bijection in view of (i). ∎
###### Remark 3.5.12.
Notice that lemma 3.5.11(ii) generalises the lifting theorem 3.3.12(i) (in case $`V`$ is an $`𝔽_p`$-algebra). Similarly, it follows from lemmata 3.5.11(i) and 3.5.2(iv) that, in case $`V`$ is an $`𝔽_p`$-algebra, one can replace “nilpotent” in theorem 3.3.12 parts (ii) and (iii) by “Frobenius nilpotent”.
In the following $`\tau `$ will denote indifferently the topology of universal effective descent defined by either of the fibered categories $`𝐰.\stackrel{´}{𝐄}𝐭^oV^a\text{-}\mathrm{𝐀𝐥𝐠}^o`$ or $`\stackrel{´}{𝐄}𝐭^oV^a\text{-}\mathrm{𝐀𝐥𝐠}^o`$.
###### Proposition 3.5.13.
If $`f:AB`$ is a morphism of almost $`V`$-algebras which is invertible up to $`\varphi ^m`$, then $`f^o`$ is bicovering for the topology $`\tau `$.
###### Proof.
In light of lemmata 3.5.9(ii) and 3.5.11(i), it suffices to show that $`\mu _{B/A}`$ is invertible up to a power of $`\varphi `$. For this, factor the identity morphism of $`B`$ as $`B\stackrel{\text{1}_Bf}{}B_AB\stackrel{\mu _{B/A}}{}B`$ and argue as in the proof of theorem 3.5.5. ∎
###### Proposition 3.5.14.
Let $`AB`$ be a morphism of almost $`V`$-algebras and $`IA`$ an ideal. Set $`\overline{A}=A/I`$ and $`\overline{B}=B/IB`$. Suppose that either
a) $`IIB`$ is an epimorphism with nilpotent kernel, or
b) $`V`$ is an $`𝔽_p`$-algebra and $`IIB`$ is invertible up to a power of $`\varphi `$.
Then we have :
i) conditions (a) and (b) are stable under any base change $`AC`$.
ii) $`(AB)^o`$ is covering (resp. bicovering) for $`\tau `$ if and only if $`(\overline{A}\overline{B})^o`$ is.
###### Proof.
Suppose first that $`IIB`$ is an isomorphism; in this case we claim that $`ICI(C_AB)`$ is an epimorphism and $`\text{Ker}(ICI(C_AB))^2=0`$ for any $`A`$-algebra $`C`$. Indeed, since by assumption $`IIB`$, $`C_AB`$ acts on $`C_AI`$, hence $`\text{Ker}(CC_AB)`$ annihilates $`C_AI`$, hence annihilates its image $`IC`$, whence the claim. If, moreover, $`V`$ is an $`𝔽_p`$-algebra, lemma 3.5.2(iv) implies that $`ICI(C_AB)`$ is invertible up to a power of $`\varphi `$.
In the general case, consider the intermediate almost $`V`$-algebra $`A_1=\overline{A}\times _{\overline{B}}B`$ equipped with the ideal $`I_1=0\times _{\overline{B}}(IB)`$. In case (a), $`I_1=IA_1`$ and $`AA_1`$ is an epimorphism with nilpotent kernel, hence it remains such after any base change $`AC`$. To prove (i) in case (a), it suffices then to consider the morphism $`A_1B`$, hence we can assume from start that $`IIB`$ is an isomorphism, which is the case already dealt with. To prove (i) in case (b), it suffices to consider the cases of $`(A,I)(A_1,I_1)`$ and $`(A_1,I_1)(B,IB)`$. The second case is treated above. In the first case, we do not necessarily have $`I_1=IA_1`$ and the assertion to be checked is that, for every $`A`$-algebra $`C`$, the morphism $`ICI_1(A_1_AC)`$ is invertible up to a power of $`\varphi `$. We apply lemma 3.5.2(v) to the commutative diagram with exact rows:
to deduce that $`AA_1`$ is invertible up to some power of $`\varphi `$, hence so is $`CA_1_AC`$, which implies the assertion.
As for (ii), we remark that the “only if” part is trivial; and we assume therefore that $`(\overline{A}\overline{B})^o`$ is $`\tau `$-covering (resp. $`\tau `$-bicovering). Consider first the assertion for “covering”. We need to show that $`(AB)^o`$ is of universal effective descent for $`F`$, where $`F`$ is either one of our two fibered categories. In light of (i), this is reduced to the assertion that $`(AB)^o`$ is of effective descent for $`F`$. We notice that $`(AA_1)^o`$ is bicovering for $`\tau `$ (in case (a) by theorem 3.3.12 and lemma 3.5.9(ii), in case (b) by proposition 3.5.13). As $`(\overline{A}A_1/I_1)^o`$ is an isomorphism, the assertion is reduced to the case where $`IIB`$ is an isomorphism. In this case, by lemma 3.4.16, there is a natural equivalence: $`\mathrm{𝐃𝐞𝐬𝐜}(F,B/A)\stackrel{}{}\mathrm{𝐃𝐞𝐬𝐜}(F,\overline{B}/\overline{A})\times _{F_{\overline{B}}}F_B`$. Then the assertion follows easily from corollary 3.4.13. Finally suppose that $`(\overline{A}\overline{B})^o`$ is bicovering. The foregoing already says that $`(AB)^o`$ is covering, so it remains to show that $`(B_ABB)^o`$ is also covering. The above argument again reduces to the case where $`IIB`$ is an isomorphism. Then, as in the proof of lemma 3.4.16, the induced morphism $`I(B_AB)IB`$ is an isomorphism as well. Thus the assertion for “bicovering” is reduced to the assertion for “covering”. ∎
## 4. Appendix
### 4.1.
In this appendix we have gathered a few miscellaneous results that were found in the course of our investigation, and which may be useful for other applications.
We need some preliminaries on simplicial objects : first of all, a simplicial almost algebra is just an object in the category $`s.(V^a\text{-}\mathrm{𝐀𝐥𝐠})`$. Then for a given simplicial almost algebra $`A`$ we have the category $`A\text{-}\mathrm{𝐌𝐨𝐝}`$ of $`A`$-modules : it consists of all simplicial almost $`V`$-modules $`M`$ such that $`M[n]`$ is an $`A[n]`$-module and such that the face and degeneracy morphisms $`d_i:M[n]M[n1]`$ and $`s_i:M[n]M[n+1]`$ $`(i=0,1,\mathrm{},n)`$ are $`A[n]`$-linear. We will need also the derived category of $`A`$-modules; it is defined as follows.
A bit more generally, let C be any abelian category. For an object $`X`$ of $`s.\text{C}`$ let $`N(X)`$ be the normalized chain complex (defined as in (I.1.3)). By the theorem of Dold-Kan ( th.8.4.1) $`XN(X)`$ induces an equivalence $`N:s.\text{C}\text{C}_{}(\text{C})`$. Now we say that a morphism $`XY`$ in $`s.\text{C}`$ is a quasi-isomorphism if the induced morphism $`N(X)N(Y)`$ is a quasi-isomorphism of chain complexes.
In the following we fix a simplicial almost algebra $`A`$.
###### Definition 4.1.1.
We say that $`A`$ is exact if the almost algebras $`A[n]`$ are exact for all $`n`$. A morphism $`\varphi :MN`$ of $`A`$-modules (or $`A`$-algebras) is a quasi-isomorphism if the morphism $`\varphi `$ of underlying simplicial almost $`V`$-modules is a quasi-isomorphism. We define the category $`\text{D}_{}(A)`$ (resp. the category $`\text{D}_{}(A\text{-}\mathrm{𝐀𝐥𝐠})`$) as the localization of the category $`A\text{-}\mathrm{𝐌𝐨𝐝}`$ (resp. $`A\text{-}\mathrm{𝐀𝐥𝐠}`$) with respect to the class of quasi-isomorphisms.
As usual, the morphisms in $`\text{D}_{}(A)`$ can be computed via a calculus of fraction on the category $`\text{Hot}_{}(A)`$ of simplicial complexes up to homotopy. Moreover, if $`A_1`$ and $`A_2`$ are two simplicial almost algebras, then the “extension of scalars” functors define equivalences of categories
$$\begin{array}{c}\text{D}_{}(A_1\times A_2)\stackrel{}{}\text{D}_{}(A_1)\times \text{D}_{}(A_2)\hfill \\ \text{D}_{}(A_1\times A_2\text{-}\mathrm{𝐀𝐥𝐠})\stackrel{}{}\text{D}_{}(A_1\text{-}\mathrm{𝐀𝐥𝐠})\times \text{D}_{}(A_2\text{-}\mathrm{𝐀𝐥𝐠}).\hfill \end{array}$$
###### Proposition 4.1.2.
(i) The functor on $`A`$-algebras given by $`B(s.V^a\times B)_{!!}`$ preserves quasi-isomorphisms and therefore induces a functor $`\text{D}_{}(A\text{-}\mathrm{𝐀𝐥𝐠})\text{D}_{}((s.V^a\times A)_{!!}\text{-}\mathrm{𝐀𝐥𝐠})`$.
(ii) The localisation functor $`RR^a`$ followed by “extension of scalars” via $`s.V^a\times AA`$ induces a functor $`\text{D}_{}((s.V^a\times A)_{!!}\text{-}\mathrm{𝐀𝐥𝐠})\text{D}_{}(A\text{-}\mathrm{𝐀𝐥𝐠})`$ and the composition of this and the above functor is naturally isomorphic to the identity functor on $`\text{D}_{}(A\text{-}\mathrm{𝐀𝐥𝐠})`$.
###### Proof.
(i) : let $`BC`$ be a quasi-isomorphism of $`A`$-algebras. Clearly the induced morphism $`s.V^a\times Bs.V^a\times C`$ is still a quasi-isomorphism of $`V`$-algebras. But by remark 2.2.13, $`s.V^a\times B`$ and $`s.V^a\times C`$ are exact simplicial almost $`V`$-algebras; moreover, it follows from corollary 2.2.10 that $`(s.V^a\times B)_!(s.V^a\times C)_!`$ is a quasi-isomorphism of $`V`$-modules. Then the claim follows easily from the exactness of the sequence (2.2.11). Now (ii) is clear. ∎
###### Remark 4.1.3.
In case $`𝔪`$ is flat, then all $`A`$-algebras are exact, and the same argument shows that the functor $`BB_{!!}`$ induces a functor $`\text{D}_{}(A\text{-}\mathrm{𝐀𝐥𝐠})\text{D}_{}(A_{!!}\text{-}\mathrm{𝐀𝐥𝐠})`$. In this case, composition with localisation is naturally isomorphic to the identity functor on $`\text{D}_{}(A\text{-}\mathrm{𝐀𝐥𝐠})`$.
###### Proposition 4.1.4.
Let $`f:RS`$ be a map of $`V`$-algebras such that $`f^a:R^aS^a`$ is an isomorphism. Then $`𝕃_{S/R}^a0`$ in $`\text{D}_{}(s.S^a)`$.
###### Proof.
We show by induction on $`q`$ that
VAN($`q;S/R`$) $`H_q(𝕃_{S/R}^a)=0.`$
For $`q=0`$ the claim follows immediately from (II.1.2.4.2). Therefore suppose that $`q>0`$ and that VAN($`j;D/C`$) is known for all almost isomorphisms of $`V`$-algebras $`CD`$ and all $`j<q`$. Let $`\overline{R}=f(R)`$. Then by transitivity ( (II.2.1.2)) we have a distinguished triangle in $`\text{D}_{}(s.S^a)`$
We deduce that VAN($`q;\overline{R}/R`$) and VAN($`q;S/\overline{R}`$) imply VAN($`q;S/R`$), thus we can assume that $`f`$ is either injective or surjective. Let $`S_{}S`$ be the simplicial $`V`$-algebra augmented over $`S`$ defined by $`S_{}=P_V(S)`$. It is a simplicial resolution of $`S`$ by free $`V`$-algebras, in particular the augmentation is a quasi-isomorphism of simplicial $`V`$-algebras. Set $`R_{}=S_{}\times _SR`$. This is a simplicial $`V`$-algebra augmented over $`R`$ via a quasi-isomorphism. Moreover, the induced morphisms $`R[n]^aS[n]^a`$ are isomorphisms. By (II.1.2.6.2) there is a quasi-isomorphism $`𝕃_{S/R}𝕃_{S_{}/R_{}}^\mathrm{\Delta }`$. On the other hand we have a spectral sequence
$$E_{ij}^1=H_j(𝕃_{S[i]/R[i]})H_{i+j}(L_{S_{}/R_{}}^\mathrm{\Delta }).$$
It follows easily that VAN$`(j;S[i]/R[i])`$ for all $`i0,jq`$ implies VAN$`(q;S/R)`$. Therefore we are reduced to the case where $`S`$ is a free $`V`$-algebra and $`f`$ is either injective or surjective. We examine separately these two cases. If $`f:RV[T]`$ is surjective, then we can find a right inverse $`s:V[T]R`$ for $`f`$. By applying transitivity to the sequence $`V[T]RV[T]`$ we get a distinguished triangle
$$(V[T]_R𝕃_{R/V[T]})^a\stackrel{u}{}𝕃_{V[T]/V[T]}^a\stackrel{v}{}𝕃_{V[T]/R}^a\sigma (V[T]_R𝕃_{R/V[T]})^a.$$
Since $`𝕃_{V[T]/V[T]}^a0`$ there follows an isomorphism : $`H_q(𝕃_{V[T]/R})^aH_{q1}(V[T]_R𝕃_{R/V[T]})^a`$. Furthermore, since $`f^a`$ is an isomorphism, $`s^a`$ is an isomorphism as well, hence by induction (and by a spectral sequence of the type (I.3.3.3.2)) $`H_{q1}(V[T]_R𝕃_{R/V[T]})^a0`$. The claim follows in this case.
Finally suppose that $`f:RV[T]`$ is injective. Write $`V[T]=\text{Sym}(F)`$, for a free $`V`$-module $`F`$ and set $`\stackrel{~}{F}=\stackrel{~}{𝔪}_VF`$; since $`f^a`$ is an isomorphism, $`\text{Im}(\text{Sym}(\stackrel{~}{F})\text{Sym}(F))R`$. We apply transitivity to the sequence $`\text{Sym}(\stackrel{~}{F})R\text{Sym}(F)`$. By arguing as above we are reduced to showing that $`𝕃_{\text{Sym}(F)/\text{Sym}(\stackrel{~}{F})}^a0.`$ We know that $`H_0(𝕃_{\text{Sym}(F)/\text{Sym}(\stackrel{~}{F})}^a)0`$ and we will show that $`H_q(𝕃_{\text{Sym}(F)/\text{Sym}(\stackrel{~}{F})}^a)0`$ for $`q>0`$. To this purpose we apply transitivity to the sequence $`V\text{Sym}(\stackrel{~}{F})\text{Sym}(F)`$. As $`F`$ and $`\stackrel{~}{F}`$ are flat $`V`$-modules, (II.1.2.4.4) yields $`H_q(𝕃_{\text{Sym}(F)/V})H_q(𝕃_{\text{Sym}(\stackrel{~}{F})/V})0`$ for $`q>0`$ and $`H_0(𝕃_{\text{Sym}(\stackrel{~}{F})/V})`$ is a flat $`\text{Sym}(\stackrel{~}{F})`$-module. In particular $`H_j(\text{Sym}(F)_{\text{Sym}(\stackrel{~}{F})}𝕃_{\text{Sym}(\stackrel{~}{F})/V})0`$ for all $`j>0`$. Consequently $`H_{j+1}(𝕃_{\text{Sym}(F)/\text{Sym}(\stackrel{~}{F})})0`$ for all $`j>0`$ and $`H_1(𝕃_{\text{Sym}(F)/\text{Sym}(\stackrel{~}{F})})\text{Ker}(\text{Sym}(F)_{\text{Sym}(\stackrel{~}{F})}\mathrm{\Omega }_{\text{Sym}(\stackrel{~}{F})/V}\mathrm{\Omega }_{\text{Sym}(F)/V})`$. The latter module is easily seen to be almost zero. ∎
###### Theorem 4.1.5.
Let $`\varphi :RS`$ be a map of simplicial $`V`$-algebras inducing an isomorphism $`R^a\stackrel{}{}S^a`$ in $`\text{D}_{}(R^a)`$. Then $`(𝕃_{S/R}^\mathrm{\Delta })^a0`$ in $`\text{D}_{}(S^a)`$.
###### Proof.
Apply the base change theorem ( II.2.2.1) to the (flat) projections of $`s.V\times R`$ onto $`R`$ and respectively $`s.V`$ to deduce that the natural map $`𝕃_{s.V\times S/s.V\times R}^\mathrm{\Delta }𝕃_{S/R}^\mathrm{\Delta }𝕃_{s.V/s.V}^\mathrm{\Delta }𝕃_{S/R}^\mathrm{\Delta }`$ is a quasi-isomorphism in $`\text{D}_{}(s.V\times S)`$. By proposition 4.1.2 the induced morphism $`(s.V\times R)_{!!}^a(s.V\times S)_{!!}^a`$ is still a quasi-isomorphism. There are spectral sequences
$$\begin{array}{c}E_{ij}^1=H_j(𝕃_{(V\times R[i])/(V\times R[i])_{!!}^a})H_{i+j}(𝕃_{(s.V\times R)/(s.V\times R)_{!!}^a}^\mathrm{\Delta })\\ F_{ij}^1=H_j(𝕃_{(V\times S[i])/(V\times S[i])_{!!}^a})H_{i+j}(𝕃_{(s.V\times S)/(s.V\times S)_{!!}^a}^\mathrm{\Delta }).\end{array}$$
On the other hand, by proposition 4.1.4 we have $`𝕃_{(V\times R[i])/(V\times R[i])_{!!}^a}^a0𝕃_{(V\times S[i])/(V\times S[i])_{!!}^a}^a`$ for all $`i`$. Then the theorem follows directly from (II.1.2.6.2(b)) and transitivity. ∎
###### Proposition 4.1.6.
Let $`AB`$ be a morphism of exact almost $`V`$-algebras. Then the natural map $`\stackrel{~}{𝔪}_V𝕃_{B_{!!}/A_{!!}}𝕃_{B_{!!}/A_{!!}}`$ is a quasi-isomorphism.
###### Proof.
By transitivity we may assume $`A=V^a`$. Let $`P_{}=P_V(B_{!!})`$ be the standard resolution of $`B_{!!}`$ (see II.1.2.1). Each $`P[n]^a`$ contains $`V`$ as a direct summand, hence it is exact, so that we have an exact sequence of simplicial $`V`$-modules $`0s.\stackrel{~}{𝔪}s.V(P_{}^a)_!(P_{}^a)_{!!}0`$. The augmentation $`(P_{}^a)_!(B_{!!}^a)_!B_!`$ is a quasi-isomorphism and we deduce that $`(P_{}^a)_{!!}B_{!!}`$ is a quasi-isomorphism; hence $`(P_{}^a)_{!!}P_{}`$ is a quasi-isomorphism as well. We have $`P[n]\text{Sym}(F_n)`$ for a free $`V`$-module $`F_n`$ and the map $`(P[n]^a)_{!!}P[n]`$ is identified with $`\text{Sym}(\stackrel{~}{𝔪}_VF_n)\text{Sym}(F_n)`$, whence $`\mathrm{\Omega }_{P[n]_{!!}^a/V}_{P[n]_{!!}^a}P[n]\mathrm{\Omega }_{P[n]/V}`$ is identified with $`\stackrel{~}{𝔪}_V\mathrm{\Omega }_{P[n]/V}\mathrm{\Omega }_{P[n]/V}`$. By (II.1.2.6.2) the map $`𝕃_{(P_{}^a)_{!!}/V}^\mathrm{\Delta }𝕃_{P_{}/V}^\mathrm{\Delta }`$ is a quasi-isomorphism. In view of (II.1.2.4.4) we derive that $`\mathrm{\Omega }_{(P_{}^a)_{!!}/V}\mathrm{\Omega }_{P_{}/V}`$ is a quasi-isomorphism, i.e. $`\stackrel{~}{𝔪}_V\mathrm{\Omega }_{P_{}/V}\mathrm{\Omega }_{P_{}/V}`$ is a quasi-isomorphism. Since $`\stackrel{~}{𝔪}`$ is flat and $`\mathrm{\Omega }_{P_{}/V}\mathrm{\Omega }_{P_{}/V}_P_{}B_{!!}=𝕃_{B_{!!}/V}`$ is a quasi-isomorphism, we get the desired conclusion. ∎
In view of proposition 4.1.4 we have $`𝕃_{(V^a\times A)_{!!}/V\times A_{!!}}^a0`$ in $`\text{D}_{}(V^a\times A)`$. By this, transitivity and localisation ( II.2.3.1.1) we derive that $`𝕃_{B/A}^a𝕃_{B_{!!}/A_{!!}}^a`$ is a quasi-isomorphism for all $`A`$-algebras $`B`$. If $`A`$ and $`B`$ are exact (e.g. if $`𝔪`$ is flat), we conclude from proposition 4.1.6 that the natural map $`𝕃_{B/A}𝕃_{B_{!!}/A_{!!}}`$ is a quasi-isomorphism.
Finally we want to discuss left derived functors of (the almost version of) some notable non-additive functors that play a role in deformation theory. Let $`R`$ be a simplicial $`V`$-algebra. Then we have an obvious functor $`G:\text{D}_{}(R)\text{D}_{}(R^a)`$ obtained by applying dimension-wise the localisation functor. Let $`\mathrm{\Sigma }`$ be the multiplicative set of morphisms of $`\text{D}_{}(R)`$ that induce almost isomorphisms on the cohomology modules. An argument as in section 2.3 shows that $`G`$ induces an equivalence of categories $`\mathrm{\Sigma }^1\text{D}_{}(R)\text{D}_{}(R^a)`$.
Now let $`R`$ be a $`V`$-algebra and $`\text{F}_p`$ one of the functors $`^p`$, $`\mathrm{\Lambda }^p`$, $`\text{Sym}^p`$, $`\mathrm{\Gamma }^p`$ defined in (I.4.2.2.6).
###### Lemma 4.1.7.
Let $`\varphi :MN`$ be an almost isomorphism of $`R`$-modules. Then $`\text{F}_p(\varphi ):\text{F}_p(M)\text{F}_p(N)`$ is an almost isomorphism.
###### Proof.
Let $`\psi :\stackrel{~}{𝔪}_VNM`$ be the map corresponding to $`(\varphi ^a)^1`$ under the bijection (2.2.2). By inspection, the compositions $`\varphi \psi :\stackrel{~}{𝔪}_VNN`$ and $`\psi (\text{1}_{\stackrel{~}{𝔪}}\varphi ):\stackrel{~}{𝔪}_VMM`$ are induced by scalar multiplication. Pick any $`s𝔪`$ and lift it to an element $`\stackrel{~}{s}\stackrel{~}{𝔪}`$; define $`\psi _s:NM`$ by $`n\psi (\stackrel{~}{s}n)`$ for all $`nN`$. Then $`\varphi \psi _s=s\text{1}_N`$ and $`\psi _s\varphi =s\text{1}_M`$. This easily implies that $`s^p`$ annihilates $`\text{Ker}\text{F}_p(\varphi )`$ and $`\text{Coker}\text{F}_p(\varphi )`$. In light of proposition 2.1.5(ii), the claim follows. ∎
Let $`B`$ be an almost $`V`$-algebra. We define a functor $`\text{F}_p^a`$ on $`B\text{-}\mathrm{𝐌𝐨𝐝}`$ by $`M(\text{F}_p(M_!))^a`$, where $`M_!`$ is viewed as a $`B_{!!}`$-module or a $`B_{}`$-module (to show that these choices define the same functor it suffices to observe that $`B_{}_{B_{!!}}NN`$ for all $`B_{}`$-modules $`N`$ such that $`N=𝔪N`$). For all $`p>0`$ we have diagrams :
(4.1.8)
where the downward arrows are localisation and the upward arrows are the functors $`MM_!`$. Lemma 4.1.7 implies that the downward arrows in the diagram commute (up to a natural isomorphism) with the horizontal ones. It will follow from the following proposition 4.1.9 that the diagram commutes also going upward.
For any $`V`$-module $`N`$ we have an exact sequence $`\mathrm{\Gamma }^2N^2N\mathrm{\Lambda }^2N0`$. As observed in the proof of proposition 2.1.5, the symmetric group $`S_2`$ acts trivially on $`^2\stackrel{~}{𝔪}`$ and $`\mathrm{\Gamma }^2\stackrel{~}{𝔪}^2\stackrel{~}{𝔪}`$, so $`\mathrm{\Lambda }^2\stackrel{~}{𝔪}=0`$. Also we have natural isomorphisms $`\mathrm{\Gamma }^p\stackrel{~}{𝔪}\stackrel{~}{𝔪}`$ for all $`p>0`$.
###### Proposition 4.1.9.
Let $`R`$ be a commutative ring and $`L`$ a flat $`R`$-module with $`\mathrm{\Lambda }^2L=0`$. Then for $`p>0`$ and for all $`R`$-modules $`N`$ we have natural isomorphisms
$$\mathrm{\Gamma }^p(L)_R\text{F}_p(N)\stackrel{}{}\text{F}_p(L_RN).$$
###### Proof.
Fix an element $`x\text{F}_p(N)`$. For each $`R`$-algebra $`R^{}`$ and each element $`lR^{}_RL`$ we get a map $`\varphi _l:R^{}_RNR^{}_RL_RN`$ by $`yly`$, hence a map $`\text{F}_p(\varphi _l):R^{}_R\text{F}_p(N)\text{F}_p(R^{}_RN)\text{F}_p(R^{}_RL_RN)R^{}_R\text{F}_p(L_RN)`$. For varying $`l`$ we obtain a map of sets $`\psi _{R^{},x}:R^{}_RLR^{}_R\text{F}_p(L_RN)`$ : $`l\text{F}_p(\varphi _l)(1x)`$. According to the terminology of , the system of maps $`\psi _{R^{},x}`$ for $`R^{}`$ ranging over all $`R`$-algebras forms a homogeneous polynomial law of degree $`p`$ from $`L`$ to $`\text{F}_p(L_RN)`$, so it factors through the universal homogeneous degree $`p`$ polynomial law $`\gamma _p:L\mathrm{\Gamma }^p(L)`$ . The resulting $`R`$-linear map $`\overline{\psi }_x:\mathrm{\Gamma }^p(L)\text{F}_p(L_RN)`$ depends $`R`$-linearly on $`x`$, hence we derive an $`R`$-linear map $`\psi :\mathrm{\Gamma }^p(L)_R\text{F}_p(N)\text{F}_p(L_RN)`$. Next notice that by hypothesis $`S_2`$ acts trivially on $`^2L`$ so $`S_p`$ acts trivially on $`^pL`$ and we get an isomorphism $`\beta :\mathrm{\Gamma }^p(L)\stackrel{}{}^pL`$. We deduce a natural map $`(^pL)_R\text{F}_p(N)\text{F}_p(L_RN)`$. Now, in order to prove the proposition for the case $`\text{F}_p=^p`$, it suffices to show that this last map is just the natural isomorphism that “reorders the factors”. Indeed, let $`x_1,\mathrm{},x_nL`$ and $`q=(q_1,\mathrm{},q_n)^n`$ such that $`|q|=_iq_i=p`$; then $`\beta `$ sends the generator $`x_1^{[q_1]}\mathrm{}x_n^{[q_n]}`$ to $`\left(\genfrac{}{}{0pt}{}{p}{q_1,\mathrm{},q_n}\right)x_1^{q_1}\mathrm{}x_n^{q_n}`$. On the other hand, pick any $`y^pN`$ and let $`R[T]=R[T_1,\mathrm{},T_r]`$ be the polynomial $`R`$-algebra in $`n`$ variables; write $`(T_1x_1+\mathrm{}+T_nx_n)^py=\psi _{R[T],y}(T_1x_1+\mathrm{}+T_nx_n)=_{r^n}T^rw_r`$ with $`w_r^p(L_RN)`$. Then $`\psi ((x_1^{[q_1]}\mathrm{}x_n^{[q_n]})y)=w_q`$ (see pp.266-267) and the claim follows easily. Next notice that $`\mathrm{\Gamma }^p(L)`$ is flat, so that tensoring with $`\mathrm{\Gamma }^p(L)`$ commutes with taking coinvariants (resp. invariants) under the action of the symmetric group; this implies the assertion for $`\text{F}_p=\text{Sym}^p`$ (resp. $`\text{F}_p=\text{TS}^p`$). To deal with $`\text{F}_p=\mathrm{\Lambda }^p`$ recall that for any $`V`$-module $`M`$ and $`p>0`$ we have the antisymmetrization operator $`a_M=_{\sigma S_p}\text{sgn}(\sigma )\sigma :^pM^pM`$ and a surjection $`\mathrm{\Lambda }^p(M)\text{Im}(a_M)`$ which is an isomorphism for $`M`$ free, hence for $`M`$ flat. The result for $`\text{F}_p=^p`$ (and again the flatness of $`\mathrm{\Gamma }^p(L)`$) then gives $`\mathrm{\Gamma }^p(L)\text{Im}(a_N)\text{Im}(a_{L_RN})`$, hence the assertion for $`\text{F}_p=\mathrm{\Lambda }^p`$ and $`N`$ flat. For general $`N`$ let $`F_1\stackrel{}{}F_0\stackrel{\epsilon }{}N0`$ be a presentation with $`F_i`$ free. Define $`j_0,j_1:F_0F_1F_0`$ by $`j_0(x,y)=x+(y)`$ and $`j_1(x,y)=x`$. By functoriality we derive an exact sequence
which reduces the assertion to the flat case. For $`\text{F}_p=\mathrm{\Gamma }^p`$ the same reduction argument works as well (cf. p.284) and for flat modules the assertion for $`\mathrm{\Gamma }^p`$ follows from the corresponding assertion for $`\text{TS}^p`$. ∎
###### Lemma 4.1.10.
Let $`A`$ be a simplicial almost algebra, $`L,E`$ and $`F`$ three $`A`$-modules, $`f:EF`$ a quasi-isomorphism. If $`L`$ is flat or $`E,F`$ are flat, then $`L_Af:L_AEL_AF`$ is a quasi-isomorphism.
###### Proof.
It is deduced directly from (I.3.3.2.1) by applying $`MM_!`$. ∎
As usual, this allows one to show that $`:\text{Hot}_{}(A)\times \text{Hot}_{}(A)\text{Hot}_{}(A)`$ admits a left derived functor $`\stackrel{\text{L}}{}:\text{D}_{}(A)\times \text{D}_{}(A)\text{D}_{}(A)`$. If $`R`$ is a simplicial $`V`$-algebra then we have essentially commutative diagrams
where again the downward (resp. upward) functors are induced by localisation (resp. by $`MM_!`$).
We mention the derived functors of the non-additive functor $`\text{F}_p`$ defined above in the simplest case of modules over a constant simplicial ring. Let $`A`$ be a (commutative) almost algebra.
###### Lemma 4.1.11.
If $`u:XY`$ is a quasi-isomorphism of flat $`s.A`$-modules then $`\text{F}_p^a(u):\text{F}_p^a(X)\text{F}_p^a(Y)`$ is a quasi-isomorphism.
###### Proof.
This is deduced from (I.4.2.2.1) applied to $`N(X_!)N(Y_!)`$ which is a quasi-isomorphism of chain complexes of flat $`A_{!!}`$-modules. We note that loc. cit. deals with a more general mixed simplicial construction of $`\text{F}_p`$ which applies to bounded above complexes, but one can check that it reduces to the simplicial definition for complexes in $`\text{C}_{}(A_{!!})`$. ∎
Using the lemma one can construct $`L\text{F}_p^a:\text{D}_{}(s.A)\text{D}_{}(s.A)`$. If $`R`$ is a $`V`$-algebra we have the derived category version of the essentially commutative squares (4.1.8), relating $`L\text{F}_p:\text{D}_{}(s.R)\text{D}_{}(s.R)`$ and $`L\text{F}_p^a:\text{D}_{}(s.R^a)\text{D}_{}(s.R^a)`$.
Acknowledgements The second author is very much indebted to Gerd Faltings for many patient explanations on the method of almost étale extensions. Next he would like to acknowledge several interesting discussions with Ioannis Emmanouil. He is also much obliged to Pierre Deligne, for a useful list of critical remarks. Finally, he owes a special thank to Roberto Ferretti, who has read the first tentative versions of this article, has corrected many slips and has made several valuable suggestions.
| Ofer Gabber | | Lorenzo Ramero |
| --- | --- | --- |
| I.H.E.S. | | Université de Bordeaux I |
| Le Bois-Marie | | Laboratoire de Mathématiques Pures |
| 35, route de Chartres | | 351, cours de la Liberation |
| F-91440 Bures-sur-Yvette | | F-33405 Talence Cedex |
| E-mail address: gabber@ihes.fr | | E-mail address: ramero@math.u-bordeaux.fr |
| | | web: http://www.math.u-bordeaux.fr/$``$ramero |
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# 1. Introduction: Few-parameter models for magnetospheric dynamics
## 1. Introduction: Few-parameter models for magnetospheric dynamics
There is growing evidence that the coupled solar wind-magnetosphere -ionosphere (SW-M-I) system, viewed as a whole, is non-equilibrium, driven, dissipative and nonlinear (Vörös, 1991). That this should be so is reasonable, given that the magnetosphere is a complex system, with multiple self-interacting phenomena, occurring on a vast range of length and time scales. A consequence of this view is that part or all of the observed magnetospheric phenomenology may be a manifestation of physics resulting from the underlying complexity of the whole system. Because of their analytical intractability, such systems in space physics are typically studied using a “large” (i.e. many-parameter) numerical simulation code. More recently, however, some systems of this type in nature have been shown to lend themselves to few-parameter descriptions (Hastings and Sugihara, 1993), which arise because, paradoxically, the complexity of the system gives rise to simplicity in some of its observable characteristics. Examples of such descriptions are shown in the top row of Table 1, adapted from figure 7.1 of Hastings and Sugihara, (1993). Starting with the simplest, “linear plus noise” description, applied to the magnetosphere by Bargatze et al., (1985), we go to low dimensional nonlinear “chaotic” models such as Baker et al. (1990)’s modified “dripping tap”. We then see fractional Brownian motion (fBm), used by Takalo et al., (1994), as a null hypothesis against which chaos could be tested and finally we have Self-Organised Criticality (SOC), the magnetospheric application of which is due to Chang, (1992). Few-parameter models of intrinsically complex systems have already demonstrated their value in space physics by their ability to describe reproducible aspects of the magnetosphere’s behaviour and to motivate nonlinear predictive filters for geomagnetic activity (see the reviews of Klimas et al., 1996 and Sharma, 1995). The extent to which such models are applicable bears directly on the extent to which magnetospheric (and other laboratory or astrophysical macroscopic plasma systems) may have predictable phenomenology. In consequence the study of few parameter models for energetically open but spatially confined plasma systems is a highly topical subject both with respect to the magnetospheric confinement system (e.g. Angelopoulos et al., 1999; Baker et al., 1999; Horton et al., 1999) and to magnetically confined laboratory plasmas (e.g. Dendy and Helander, 1997; Carreras et al., 1999; Pedrosa et al., 1999). One possible new approach of considerable current interest is the SOC paradigm introduced by Bak et al. (1987).
There is a natural hierarchy of few-parameter descriptions, ordered by the extent to which the many coupled degrees of freedom of the system manifest themselves. Broadly speaking, as we go from left to right along Table 1, we move down the hierarchy of description and the large number of degrees of freedom become increasingly explicit in the description. In consequence the importance of an underlying theory to define the model fully tends to also increase. To make effective use of the theory-model correspondence in such a table, however, theories must be falsifiable, as otherwise the parameters of the simple model may simply be “tuned” to bring it into closer and closer agreement with data. This may be a two-way process, as for example, casting a given model in falsifiable form by defining which phenomena must be tested for helps to clarify the underlying theory.
To make these abstract points more concrete, consider Table 1. The first row shows some notable examples of various simple models that have been applied both to the complex magnetospheric system, and to other such complex systems. The first column shows various properties that these models have which could be tested for in data, provided that suitable variables are measurable. If a given property can be shown not to be present in data then we can eliminate models which depend on it from consideration. In this paper we first describe the construction of Table 1 by describing the four levels of description which it encapsulates. As models based on the SOC hypothesis are of current interest for the SW-M-I coupling problem, we then specifically address the tests necessary to cast SOC models in falsifiable form.
### 2.1 Linear models with optional noise term
Column 1 of Table 1 is the “linear + noise” model, typically a linear differential equation with optionally a stochastic noise term $`\mathrm{\Delta }𝐰(t)`$, (adopting the notation of Hastings and Sugihara (1993)) to which we may also add a driving term $`𝐅(𝐭)`$:
$$\frac{d𝐱(t)}{dt}=𝐠(𝐱,t)+\mathrm{\Delta }𝐰(t)+𝐅(t)$$
(1)
where $`𝐠(𝐱,t)`$ can only be linear in the variables $`𝐱`$. Physically an input-output system is linear if the form of a system’s response closely resembles that of the forcing terms. This was the first level of approximation used in the input-output analysis of the SW-M-I system (Bargatze et al., 1985). The second and third rows of the table, labelled “Short-” and “long-term predictable” refer to the fact that in the absence of noise ($`\mathrm{\Delta }𝐰(t)=0`$) the short-and long-term behaviour of equation (1) is completely deterministic, while even if an additive stochastic noise term is present, closely-spaced initial conditions do not show exponential divergence. Such systems typically show relatively narrow-band spectral behaviour if the $`𝐠`$ term is dominant i.e. characteristic frequencies, and so we have “no” in the “global scaling” row for this model (row 4, column 2) to indicate that they would then not be scale free across the whole frequency range. They may, however, show regions of scale free behaviour in their frequency spectra, indicated in the table by the “sometimes” in the “scaling regions” row (column 2, row 5). The entry “no” for “low G-P dimension” (sixth row, column two), refers to the fact that such a system will usually appear high dimensional to the Grassberger-Procaccia (GP) algorithm (Grassberger and Procaccia, 1983), because of many degrees of freedom of the noise term.
### 2.2 Nonlinear deterministic models
Bargatze et al. (1985) confirmed the presence of nonlinearity in the $`AE`$ family of indices, ($`AE,AU`$, and $`AL`$) and hence the need for a next level of approximation. A prototype for differential equation models which exhibit nonlinear but deterministic dynamics (see the reviews of Sharma (1995) and Klimas et al. (1996)) is the “dripping faucet” of Shaw (1984), which was adapted to the magnetospheric problem by Baker at al. (1990). These models are of the form
$$\frac{d𝐱(t)}{dt}=𝐠(𝐱,t))+𝐅(t)$$
(2)
where unlike equation (1) the term $`𝐠(𝐱,t)`$ now contains nonlinear terms. In the hierarchy of Table 1, this is a nonlinear model, which can exhibit low-dimensional, chaotic dynamics (column 3). Familiar examples of such systems in nonlinear physics include the (continuous) driven nonlinear pendulum and (discrete) logistic map (see e.g. Rowlands, 1990). In the magnetosphere this description was inspired by an analogy between a dripping tap and plasmoid ejection during substorms. The analogy was developed into a simplified magnetospheric model by estimating the large-scale electrical properties of the M-I system and combining these electrical components into a driven nonlinear oscillator circuit model (Klimas et al., 1992). It has been further developed into a plasma physics model by Horton and Doxas (1996).
In the case of a dissipative, driven, autonomous low dimensional system such as the Lorenz model, the dynamics, rather than exploring all of phase space ergodically, collapses onto a low dimensional region called an attractor. This attractor has fractional dimension (i.e. it is a strange attractor, in contrast for example with the 2D ellipse described in phase space by a simple linear 1D pendulum). A time series drawn from such a system will thus also have low fractional dimension when tested with the Grassberger-Procaccia algorithm, so we write “yes” against “Low G-P dimension” in column 3, row 6 of table 1. A strange attractor has the property that closely-spaced trajectories, with initial conditions identical to within measurement error, will diverge strongly if they traverse certain regions of the attractor i.e. repulsive fixed points (see figure 1 of Palmer, (1993) for a clear illustration of this). We thus write “no” against “long-term predictability” (column 3, row 3), because in this sense, measured by a positive Lyapunov exponent (e.g. Rowlands, 1990), it is now not present. The significance of this “new” low-dimensional, deterministic, chaos is that sensitive dependence on initial conditions arises from the $`𝐠`$ term and so exists without the presence of “old fashioned” stochasticity i.e. we need no $`\mathrm{\Delta }𝐰`$ term. Such a model can generate wide-band, scale free “1/f” spectral behaviour when near a tangent bifurcation leading to intermittency (Lichtenberg and Lieberman, 1992), but this requires choice of certain values of the control parameters i.e. tuning. We thus write “sometimes” for against “global scaling” and “scaling regions”. “Tuning” in this sense has been considered a weakness in the applicability of low dimensional chaos to any complex natural system (Bak, 1997), not only the magnetosphere. A second practical question with such methods is that because the model definition usually starts from the observables, the map to which one applies nonlinear dynamics must be derived from data rather than given a priori from theory. One might however see this as a strength, and in practice this is addressed by an iterative process whereby the parameters suggested by observation and theory are being brought closer together (Klimas et al., 1996).
### 2.3 Stochastic descriptive models
Osborne and Provenzale (1989) showed that time series taken from certain random “coloured noise” processes, when tested with the G-P algorithm, would exhibit low dimensionality, and thus behave in this respect as a low dimensional chaotic system would. This led to the application by Takalo et al., (1994 and references therein) of a third type of model, fractional Brownian motion denoted by fBm in Table 1 (e.g. Malamud and Turcotte, 1999), as a hypothesis against which to test the low dimensional nonlinear models described by the previous column. The suggestion of fBm recognised the possibility that the apparent low dimensionality and fractality of the magnetospheric indices was the consequence of their being the output of an otherwise intrinsically many-degree of freedom stochastic system, identified by a particular “coloured noise” power spectrum (hence “sometimes” against “Low G-P dimension” row 6, column 4). Effectively the model is:
$$\frac{d𝐱(t)}{dt}=\mathrm{\Delta }𝐰(t)$$
(3)
A simple example is Brownian motion where the time evolution is discrete, and each step ($`\delta 𝐱=\delta t\mathrm{\Delta }𝐰`$ ) is drawn from a white Gaussian distribution (e.g. Malamud and Turcotte, 1999). We then find that neither short term nor long term prediction is possible because each step is entirely stochastic, giving us “no” against “short-” and “long-term predictable” (rows 2 and 3 of column 3). We note however that closely-spaced initial conditions diverge algebraically rather than exponentially, i.e. the impossibility of long term forecasting here arises from the external stochasticity in $`\mathrm{\Delta }𝐰(t)`$ rather than intrinsic chaos from $`𝐠`$. Global scaling (row 4, column 4) must arise, irrespective of any free parameters, because there is no time scale in such a model. More complex time evolutions, where successive steps are taken from a fractional Gaussian noise, are called fractal Brownian motions. A subset of such motions has been shown to have low G-P dimension (Osborne and Provenzale, 1989). We note that the presence of global scaling or scaling regions in the power spectra drawn from time series, or low G-P dimension, cannot distinguish between nonlinear low dimensional models (column 3) and fBm (column 4), because they are shared properties. The differences will only be unambiguous when one notes the different physical origins of the low G-P dimension between chaos and coloured noise, see Takalo et al., (1994), or when one uses another discriminator such as short term predictability.
### 2.4 Sandpile models of self organised criticality
The most recently introduced class of models (column 5) in Table 1 are those motivated by the hypothesis of Self-Organised Criticality (Chang, 1992; 1999, see also Vörös, (1991); Chen and Holland, (1993); Robinson, (1993)). SOC was first identified in (Bak et al., 1987), and can be modelled by, numerical cellular automaton “sandpile models” (Katz, 1986; Bak et al., 1987) These discrete-variable models (Consolini, 1997; Uritsky and Pudovkin, 1998) and the closely-related continuous-variable discrete-space time models (Chapman et al, 1998;2000, Takalo et al., 1999a;1999b; Watkins et al., 1999b) are currently being studied for their possible magnetospheric application. Consideration of SOC in our context was motivated in particular by the fact that SOC can account for known magnetospheric phenomenology such as low dimensionality (Chang, 1992) and scale free power spectra, while providing a framework for understanding observed properties of the magnetosphere such as bursty transport in the tail (c.f. “Bursty bulk flows” (Angelopoulos, 1996)). It may be that the long term value of SOC to plasma physics will be as a starting point for more realistic “avalanche” models of turbulent transport (see Dendy and Helander, 1997, for more on this question, as applied to laboratory plasmas). However, at this stage of its consideration with respect to understanding in magnetospheric physics it remains useful to consider Bak et al.’s original, sandpile model-based, definition of SOC in the framework of our table of observables, as it is used explicitly or implicitly by much of the current work on SOC in this and other fields.
Bak et al. (1987, henceforth BTW87) originally proposed SOC to explain the apparent ubiquity of both spatial fractals and “1/f” spectra in nature. They observed it in a class of numerical cellular automata, called “sandpile models” for which analogous continuous thresholded diffusion equations have since been shown to exist (Lu, 1995). The equations are modified from stochastic diffusion equations (Pelletier and Turcotte, 1999) which have a form such as
$$\frac{𝐠(𝐱,t)}{t}=^2𝐠(𝐱,t)+\mathrm{\Delta }w(x,t)$$
(4)
by the introduction of a thresholding process (Jensen, 1998), e. g.
$$\frac{𝐠(𝐱,t)}{t}=^2𝐠(𝐱,t)\mathrm{\Theta }(𝐠(𝐱,t)𝐠_𝐜)+\mathrm{\Delta }w(x,t)$$
(5)
where the step function $`\mathrm{\Theta }`$ initiates diffusive transport when the variable $`𝐠`$ reaches a critical gradient $`𝐠_𝐜`$. The main debate centres on how to motivate and satisfactorily introduce this ad hoc thresholding term (see Lu, 1995;Jensen, 1998) but several properties, such as the low frequency power spectrum may be dependent only on equation (4) and the nature of the boundary terms (Jensen, 1998), and so may be common to both SOC and stochastic diffusion.
The behaviour classified by BTW87 as SOC is the evolution of the medium described by the cellular automaton or differential equation models from arbitrary initial conditions to a non-equilibrium but steady state, “self-organisation”. The medium then evolves by dissipating energy on all scales via thresholded reconfiguration/energy release events called “avalanches”.
The assertion by Bak et al. (1987;1988) that the observed scale free, and hence power law, distribution of the size of these energy release events (the “avalanche distribution”) measured the arbitrary response of a self-organised fractal structure in the medium to perturbations introduced by random fuelling was the reason for their use of the term “critical”. Their analogy was with the scale free critical state associated with phase transitions in critical phenomena (Huang, 1987). The common observable features cited by Bak et al. were that both systems were globally scale free (hence we may write “yes” in column 5, row 4 of table 1), and also that a finite size scaling analysis gave a good data collapse, as it would in a bona fide critical system (Cardy,1996). The combination of a scale free response to perturbations with the release of this energy by random unloading events, was expected to give rise to a power law frequency spectrum. This was expected to be “1/f” and hence to explain the ubiquity in nature of noise with correlations on all time-scales. Unlike a “$`1/f^2`$” spectrum above a characteristic frequency which can be explained in many systems simply by random switching of levels, the appearance of “1/f” ($`f^\beta `$) spectra where the spectral index is between about 0.8 and 1.4 is a long standing problem in many branches of physics (Jensen, 1998).
In summary, in this picture, the SOC hypothesis would be that: “extended driven systems will tend to self-organise to fractal structures which dissipate energy on all scales in space and time, and hence give rise to scale-free “avalanche” energy burst distributions and “1/f” noise”. More recent sandpile algorithms which allow fuelling to continue while unloading occurs have a broken power law spectrum and so we have added “yes” to column 5, row 5 as well (see also section 3.1.2).
SOC behaviour, as diagnosed by scale-free energy release and/or “1/f” power spectra, has since been claimed for many systems in nature (see chapter 3 of Jensen (1998) for a compact review, and Rodriguez-Itube and Rinaldo (1997) for a longer exposition in the particular context of fractal river networks). At this point, we simply note that a definition of SOC in terms of what an SOC system does can only be used to identify an SOC system if no other system does exactly the same things. Identification of global scaling, shared by SOC, fractal Brownian motion and low dimensional chaotic systems when intermittent is, for example, thus not an unambiguous test. It is for this reason that we have used Table 1, as a guide to how the “footprint” of SOC may be more unambiguously defined. We have left the other rows as question marks because BTW87’s sandpile model definition of SOC was not unambiguous in these respects, and these issues are still under study. The wider definition of SOC used by Chang, (1992), predicts low dimensional behaviour i.e. “yes” in column 5, row 6, at least close to criticality, while many workers have taken the predictions for column 5, rows 1 and 2 to be “no” e.g. the remarks of Consolini, (1997): “In fact, if the magnetospheric dynamics could be the result of a low-dimensional dimensional chaotic dynamics, we could have some hope to forecast the evolution. On the contrary, the existence of a critical state removes this possibility, because the fluctuations of the system at a critical point are completely unpredictable”. This is equivalent behaviour to fBm. See also Bak (1994) on this point where sandpile models are asserted not to show sensitive dependence on initial conditions i.e. they are unpredictable on both long and short timescales but not chaotic. The similarity of SOC to fBm with regard to the phenomena in Table 1 might suggest that SOC adds nothing to fBm. However, there are differences. One is the fact that an SOC system releases energy by means of avalanches, effectively a new observable property, which we have thus indicated by adding a row in Table 1 to those used by Hastings and Sugihara (1993). We are indebted to a referee for the suggestion that avalanche models may also have different phase spectra to the (usually random) phase behaviour of noise. A second advantage is that SOC can be explained in terms of an underlying theory and encapsulated in terms of sandpile models, which begin to allow explanation in terms of the underlying plasma physics of the system. A third advantage is that the release of energy by avalanches is suggestive both of bursty transport in plasma confinement systems (e.g Carreras et al., 1999;Pedrosa et al., 1999), and, possibly, the substorm problem (Chang, 1992; Consolini, 1997).
The study of SOC in solar terrestrial-physics has proceeded initially through comparison of signatures in data, particularly the AE/AU/AL indices, with analogous signatures in “sandpile model” realisations of SOC. However, as we will now show, these signatures are not all unique to SOC, and the combinations in which they appear may be model dependent. Furthermore, we recall that the original proposal of the relevance of SOC to the SW-M-I system (Chang, 1992) was not predicated exclusively on a definition of SOC derived from sandpile models. To minimise possible confusion in this rapidly developing area, two questions are addressed. These are i) which experimental signatures will be needed to distinguish unambiguously between SOC and, for example, deterministic chaos and ii) what are the predictions of SOC models which will be robust against fluctuations in the input, or limited station or satellite coverage etc?
## 3 Towards unambiguous tests of SOC
Having decided what the predictions of SOC are which may be confidently entered in Table 1, we now go on to see what the observations currently available enable us to say. The question immediately arises as to whether, Picture A), the SOC system is seen as being the complete magnetosphere (“global SOC”), in which case $`𝐱`$ in equation (5) are the system state variables, for which the AE indices (Davis and Sugiura, 1966) have been taken as proxies; or, Picture B), SOC is more local in scope (“local SOC”), and plays a role in generating, stabilising and destabilising the magnetotail current sheet, in which case $`𝐱`$ might be a locally-measured magnetic field or the field seen in an MHD-derived sandpile simulation. Picture A) is closer to that given in Uritsky and Pudovkin (1998) and Consolini (1997; 1999), while Picture B) seems to us to be the motivation for Takalo et al., (1999a;b;c). Because any approximation will have a natural maximum scale of applicability, the idea of “local SOC” is not the contradiction it may at first appear to be.
It has been objected that if A) were true, all system-level outputs should show global scaling and that some are observed not to have this property (Borovsky and Nemzek, 1994). However Chapman et al. (1998), used the 1-dimensional avalanche model of Dendy and Helander (1998), to illustrate a system in which the internal energy release showed scaling while energy flowing out of the system (“systemwide”) did not, a feature seen in some other sandpile models (e.g. Pinho and Andrade, 1998). Until pictures A and B can either be distinguished or reconciled, care must be taken not to justify one using measurements consistent with the other and vice versa. We thus first (section 3.1) examine those system level outputs in which evidence of SOC has been claimed, and then (section 3.2) briefly consider evidence for SOC on more internal scales.
### 3.1 Remote Observations of system outputs
So far the main global dataset for testing for SOC has been the $`AE`$ indices. This is because, since Bargatze et al. (1985), a candidate dynamical variable for all the models discussed above has been the energy dissipated by the magnetosphere into the ionosphere, for which most workers have taken the Auroral Electrojet Index ($`AE`$) to be a proxy.
#### 3.1.1 Global scaling: Power law power spectra
The work of Tsurutani et al. (1990) described the power spectrum of $`AE`$ as “broken power law”, in that the high frequency behaviour is approximately $`1/f^2`$ while the lower frequency behaviour is approximately $`1/f`$, with a break at (1/5) hours<sup>-1</sup>. This “1/f” behaviour has been cited as evidence of SOC in the magnetosphere (Consolini, 1997; 1999; Uritsky and Pudovkin, 1998). Two main scenarios have been advanced, in one the “1/f” spectrum is seen as arising from interactions between correlated avalanches, which would then be interpreted as substorms (Consolini, 1997); while in the second the “1/f” behaviour (Uritsky and Pudovkin, 1998) is related in part to the input, which is allowed to modulate the threshold values in the sandpile algorithm.
The first apparent complication in this interpretation is that the $`AE`$ spectrum is “broken” i.e. a “1/f<sup>2</sup>” high frequency part has been reported, whereas criticality in BTW87’s original picture was expected to give long-period correlations and hence a “1/f” spectrum (Jensen, 1998). The resolution, due to Consolini (1997), is discussed in the section 3.1.2. The second, more fundamental, problem is that a “1/f” spectrum in the ouput of a system could only be an unambiguous indicator of SOC if this spectrum is not being passed through from the input. The fact that the input spectrum of the solar wind follows $`AE`$ closely over the low frequency “1/f” range that concerns us here (Tsurutani et al., 1990; Freeman et al., 1998) suggests that the power spectrum should not be used for this purpose. In this context, the high degree of predictability of $`AE`$ from the solar wind input is suggestive (e.g. Baker et al, 1997), as is the fact that the power spectrum of the signal from a neural network prediction of $`AE`$ has “1/f” form (Takalo et al., 1996). We return to this question in section 3.1.3 when we discuss the avalanche statistics. The possible ability of some avalanche algorithms to emulate a nonlinear filter, and show sensitivity to the distribution of the input fuelling rate (Takalo et al., 1999c); or conversely to eliminate traces of fluctuating input (Watkins et al., 1999b), increases the relevance of this question.
### 3.1.2 Scaling Regions: Spectral breaks
If the system is known to be SOC a priori or from other tests, the presence of a high frequency “1/f<sup>2</sup>” component is understandable. This is because recent work (notably Hwa and Kardar, 1992) has shown that a “running” sandpile (and hence SOC differential equation models such as the example used by (Takalo et al., 1999a)) can show this type of “broken” power spectrum. The reason is that allowing the fuelling to occur on a similar time scale to the unloading events permits a bursty “1/f<sup>2</sup>” power spectrum of individual avalanches to co-exist with the “1/f” power spectrum which is ascribed to interactions between events (see Jensen, 1998). If the bursts are identified with substorms then the break at 5 hours will be related to the maximum duration of a substorm. Furthermore, the original Bak et al. (i.e. “non-running”) sandpile model was quickly shown (Jensen et al., 1989) to produce a “1/f<sup>2</sup>” spectrum in its energy release events, illustrating that the although the pile is in critical state, shown in particular by the finite size scaling of the avalanche distribution (Bak et al., 1988), the critical state is not revealed by the energy release power spectrum.
The complications in this very appealing simple interpretation arise for two main reasons. One is that we do not know a priori that the system is SOC, so mapping an output variable of the sandpile model to the observed $`AE`$ spectrum is not a unique process. An alternative way to get a broken spectrum of the form shown by Tsurutani et al. (1990) for $`AE`$ is in a boundary driven 2D sandpile of the BTW87 type (see figure 4.6 of Jensen, 1998). In this case the variable whose spectrum is obtained is not the transport of sand over the edge of the pile but the sum of the dynamical variable $`𝐠`$ across the pile i.e
$$<g(t)>=\underset{i,j}{}g_{ij}𝑑𝐱𝐠(𝐱,t)$$
(6)
As with Tsurutani et al, 1990, the spectrum shown by Jensen (1998) is $`1/f`$ below a critical frequency and $`1/f^2`$ above, with the break set by a time scale $`T_{max}(L)`$ corresponding to the longest avalanche possible in the system. This would imply that $`T_{max}`$ is related to the system size $`L`$, and furnish a possible test if the system’s value of $`L`$ could be varied significantly. In other words, the robust property is the break itself rather than the variable whose broken spectrum is being calculated.
The second problem is that the broken power law spectrum for $`AE`$ cannot be uniquely interpreted as the output of an SOC system because other types of physical system can produce power spectra which show global scaling or scaling over a region or regions. The example of global scaling, i.e. scaling over a very wide bandwidth, discussed in section 2.3 was simple Brownian motion which has an $`f^2`$ spectrum at all values of $`f`$. A less well known example of scaling over a restricted region is the “random telegraph”, a random sequence of square pulses (i.e. states +1 or -1) switched at Poisson distributed intervals which gives $`f^2`$ for frequencies higher than the inverse correlation time but has a flat spectrum (because uncorrelated) for lower frequencies (Bendat, 1958; Jensen, 1998). It is very important to note that the “1/f<sup>2</sup>” part of the spectrum here is due entirely to the exponential autocorrelation of the pulses, and is not the same as the intrinsically scale free, and wideband, behaviour of a coloured noise source such as Brownian motion. If the lifetimes of the correlated pulses extend over two orders of magnitude in time then so will the “1/f<sup>2</sup>” spectrum, and so a test such as the second order structure function (see Takalo et al., 1994 and references therein), or a variance histogram (i.e. Fourier power spectrum) will be unable to distinguish this “trivial” apparent scaling from the “interesting” scaling resulting from coloured noise. A similar problem whereby level changes with a $`1/f^2`$ spectrum might mask an intrinsic Kolmogorov spectrum was treated for solar wind turbulence by Roberts and Goldstein (1987).
The fact that the high frequency scaling region in the spectrum of $`AE`$ might arise from a cause other than SOC is important in our application because $`AE`$ is known a priori to be a compound index which mixes driven and unloading effects (Kamide and Baumjohann, 1991). This mixed origin is reflected by its power spectrum (Freeman et al., 1998), structure function (Takalo et al., 1994), and avalanche distribution (section 3.1.3 and Freeman et al., 2000). In consequence a suitable “null hypothesis” for the power spectrum against which the SOC models so far proposed should be evaluated is that $`AE`$ consists of a solar wind driven “1/f” component - arising from the $`DP2`$ convection electrojet (Kamide and Baumjohann, 1991) - and a random unloading $`DP1`$ substorm electrojet component which looks like “1/f<sup>2</sup>” over two orders of magnitude in frequency and appears predominantly in $`AL`$. By analogy with section 3 we may call this Picture C (“no global SOC”).
A possible avenue for testing Consolini’s (1997) “interacting burst” interpretation of the “1/f” spectrum would then be to see if the correlation properties of the “1/f” part of the power spectrum of $`AE`$ differ in any way from those of Akasofu’s $`ϵ`$ parameter, which estimates the componenent of solar wind Poynting flux entering the magnetosphere. If they do, this adds support to the possibility that the bursts may be correlated with each other as a result of a process which occurs in the magnetosphere itself; rather than the “1/f” behaviour being explained by the long-period correlation already present in the solar wind’s power spectrum.
### 3.1.3 Global scaling: Avalanche distributions
Because of the above concerns, we see a better candidate for an unambiguous indicator of SOC as being the statistical distribution of energy released by individual “events”. Since the work of Bak et al. an SOC system has been expected to show a “power law” probability distribution for this quantity. Consolini, (1999) has plotted the distribution $`D(s)`$ of a burst measure $`s=_\mathrm{\Omega }(AE(t)L_{AE})𝑑t`$, formed from $`AE`$ where $`L_{AE}`$ was a running quiet time background level of $`70\pm 30`$ nT, and each integration was taken over a period $`\mathrm{\Omega }`$ where the integrand was positive (a burst). The $`AE`$ data used covered the period from 1975 and 1978-1987. The distribution obtained could be described by an exponentially cutoff power law (Consolini, personal communication, 1998) extending the result previously obtained by Consolini, (1997) for data from 1978.
The presence of such a power law suggested a magnetospheric analogue of the Gutenberg-Richter law in seismology, and has played a significant role in generating the interest in SOC in magnetospheric physics. Both its existence and its apparent robustness with activity level require confirmation and explanation whether by an SOC theory or another one. Although both a simple power law, or the above exponentially cut off power law are possible fits, Consolini (1999) has recently demonstrated that a better description for the burst size distribution of $`AE`$ is an exponentially modified power law with a small lognormal “bump” component. If the system is SOC, this requires an explanation of the “bump”, which may be found in the different behaviour of internal and systemwide dissipation in some sandpile models (Chapman et al., 1998) or in subcritical dynamics in the SOC system (Consolini, 1999).
### 3.1.4 Global scaling: Lifetime distributions
The SOC hypothesis had earlier led Takalo (1993) and Consolini (1999) to examine the distribution of lifetimes of the bursts. This is potentially a stronger indicator of SOC than the burst size, because exponentially modified power-law burst size distributions can also be generated by randomly quenched, exponentially growing instabilities in an otherwise non-critical medium (Aschwanden et al., 1998). $`AE`$ was found (Consolini,1999) to show a exponentially modified power law distribution of lifetimes, but with evidence of a “bump” at around 100 minutes. The “null hypothesis” mentioned in 3.1.2 (Picture C) led Freeman et al., (2000) to calculate the analogous burst lifetime distributions for $`AU`$, $`AL`$ and the solar wind quantities $`ϵ`$ and $`vB_s`$. They found exponentially modified power laws with very similar slopes for all quantities, but the “bump” was only found in the AL and AE, magnetospheric component. This suggests that the “bump” is of intrinsically magnetospheric origin (due to the DP1 current (Kamide and Baumjohann, 1991)) while the scale-free burst lifetime distribution may actually be of solar wind origin (Freeman et al., 2000), if the DP2 current system (Kamide and Baumjohann, 1991) is transparent to the driver.
### 3.1.5 Other tests: Predictability, low dimensionality
It will be necessary to examine the sandpile models and other realisations of SOC in more detail before we can say with certainty what they predict for the remaining rows of column 5, table 1. Bak (1994)’s arguments about long term prediction are based on the assertion that $`\delta `$, the separation of two initially infinitesimally close trajectories in a BTW87-type sandpile model grows with time quadratically, $`\delta =at^2`$, rather than exponentially, $`\delta =e^{\lambda t}`$, as would be the case in a chaotic system such as that of column 3 (where $`\lambda `$ is the Lyapunov exponent). This assertion needs to be tested in other sandpile models and in data from candidate systems. It is a potential discriminator between chaos and SOC.
The demonstration that an SOC system can show low dimensional behaviour was given by Chang (1992;1999) on the basis of a more general formulation of SOC than the sandpile-inspired one of Bak et al. It thus remains to be seen in general what sandpile models predict for dimensionality.
### 3.1.6 New tests: Intermittency and laminar time
The open questions described in the previous section and the ambiguities necessarily present in the data reviewed in sections 3.1.1 to 3.1.4 mean that at present it is not possible to completely eliminate any of the models discussed in Table 1, except the artificially simple linear model which was included for completeness. New tests are thus required. An example of such a test is the degree of intermittency present in the time series. Consolini et al. (1996) showed that the $`f^2`$ spectral regime of $`AE`$ might be described as an $`f^{1.8}`$ regime corresponding to the inertial range of a turbulent system, with an exponent modified from the Kolmogorov value by the presence of intermittency. These authors showed a good fit to the p-model of turbulence, also shown in the solar wind by Horbury and Balogh (1997). We are thus not presently able to distinguish between intermittency intrinsic to $`AE`$ and that due to the solar wind driver. Further work on this topic is likely to prove valuable (see also Vörös et al., 1998).
More recently, it has been claimed (Boffetta et al., 1999) that the probability density $`D(\tau )`$ of time intervals $`\tau `$ between bursts (the “laminar time”) can be used to distinguish an SOC system of the BTW87 type, which has exponential $`D(\tau )`$, from a shell model of turbulence, which has power law $`D(\tau )`$. It might seem that the power law $`D(\tau )`$ for AE shown by Consolini (1999) would rule out SOC in the global magnetosphere. However, as emphasised by Einaudi and Velli, (1999), the predictions for $`D(\tau )`$ are not in general known either for more realistic SOC models or for all turbulence models. The relevance of this work to the issue of “sympathetic flaring” (Boffetta et al., 1999 and references therein) in solar physics is likely to give rise to further exploration of this topic, and hence magnetospheric application.
### 3.2 Local observations of current sheet dynamics
It is fair to say that, for the above reasons, the evidence of SOC in the largest scale outputs of the magnetospheric system, measured by $`AE`$ and $`AL`$, is not yet persuasive. The main problem is that the behaviour of the $`AE`$ indices is similar to that of the solar wind in a number of respects. If an intrinsic and a solar wind component are always present, then testing for SOC in these compound indices will always be problematic, as will the interpretation of results based on them (c.f. Consolini, 1999; Freeman et al., 2000).
It may be more instructive to study regions of the magnetosphere where the effect of the solar wind input is less directly visible, and recent attention has been focused on SOC as a model of the magnetotail current sheet. So far this has been achieved by truncating the ideal MHD equations i.e. replacing the convection term in
$$\frac{𝐁(𝐱,t)}{t}=^2𝐁(𝐱,t)+(𝐯𝐁)$$
(7)
by a source term, resulting in an equation analogous to (4)
$$\frac{𝐁(𝐱,t)}{t}=^2𝐁(𝐱,t)+\mathrm{\Delta }w(x,t)$$
(8)
and then introducing one of several possible thresholding terms c.f. equation (5) (Vassiliadis et al., 1998 ; Takalo et al, 1999a;b;c) to map the problem onto a cellular automaton or a differential equation like that of (Lu, 1995). In view of the limited applicability of such non-self consistent models it is encouraging that reduced MHD simulations of turbulence are also demonstrating SOC phenomenology (Einaudi and Velli, 1999).
Consideration of SOC as a model for the magnetotail may also be motivated by the suggestion by of Zelenyi et al.,(1998) that the tail exists as a critical percolation cluster. Critical percolation, whereby an avalanche can extend exactly to the maximum scale length of a system rather than just below or just above it, was the original proposed explanation for the relevance of criticality in SOC (Bak et al., 1988). The idea has recently been further developed to explain how self-organisation occurs in sandpiles (Zapperi et al., 1997) in a picture whereby the edge fluctuations drive the system back to a critical percolation state.
## 4. What signatures of SOC are robust enough to be detectable in “real world” data ?
SOC is of particular interest to magnetospheric physics because it is robust, in the sense that the characteristic observed behaviour does not necessarily change greatly over wide ranges of parameter space. This robustness is thus distinct from the scale invariant phenomena that arise near critical points such as phase transitions and hence in restricted regions of parameter space (Huang, 1987). SOC systems are in this way also distinct from chaotic systems such as the dripping tap which frequently show radically different types of behaviour as control parameters change, a point emphasised by Bak (1997, pages 29-31). They are not however, as easily distinguishable from fBm.
However we still need to ask what aspect or aspects of the magnetosphere’s behaviour would be both sufficient to uniquely identify SOC and yet also robust enough to be seen under the wide range of activity levels exhibited by the magnetosphere and the solar wind i.e., to pick the specific robust discriminator. The slowly driven condition is of particular interest in the magnetospheric context because most SOC simulations have been conducted in the limit where the rate of inflow is “slow” relative to dissipation. Watkins et al., (1999b); and Chapman et al., (1999) have recently studied the question of how robust the magnetospherically relevant aspects of SOC are to changes in the inflow rate. They found that the power law avalanche distribution was preserved for the largest values of internal energy release, and gave arguments as to why this should be so. This result may give confidence that such a distribution, if shown on other grounds to be unique to SOC, will be able to be used as a diagnostic. Similar studies for the power spectrum are being carried out in an MHD-derived model by Takalo et al., (1999c).
## 5. Conclusions
We have attempted to identify the distinguishing observable features of different few-parameter models applied to the magnetosphere. The “linear+noise” model was abandoned because of observed nonlinearity, low dimensionality and lack of long-term predictability in the auroral index time series. Low dimensional models have been questioned because the low dimensionality is not unique to them and because their scaling properties are not robust against changes in the input parameters. An alternative, fractional Brownian motion, which gives low dimensionality and robust scaling is unsatifying because it does not lead to an underlying plasma physical description. The newest alternative, SOC, chosen for its robust scaling properties, can be seen as providing both a physical explanation for fBm and also accounting for the bursty nature of transport in the magnetosphere. SOC has yet to have its low dimensionality and predictability properties fully defined, but so far they seem to be similar to those of fBm. Thus attention must be focused on other means of distinguishing these last two, such as the observed intermittency and avalanching properties. Even so, questions about the application of fBm and SOC as models of the magnetosphere’s large scale output (picture A) rather than of its solar wind-driven aspects are raised by the similarity of input and output power spectra and burst distributions. Resolution of these issues is hampered by the narrow bandwidth of even the best available data series, which for example make it difficult to distinguish between wide-band coloured noise and random state changes as the origin of the $`f^2`$ spectrum of AE (Watkins et al., 1999a).
This seems to leave four possibilities, not all of which are mutually exclusive:
i) The global “SOC”-like properties we have referred to come from outside the magnetosphere, i.e. the magnetosphere can be quite well described by a “weakly nonlinear plus coloured noise” model; weakly nonlinear to give the necessary degree of predictability of the output from the input while giving long-term unpredictibility, but with a coloured noise input from the fBm or SOC nature of the turbulent solar wind causing the scaling properties. This scenario appears to be compatible with picture C, the data in Freeman et al. (2000), and alternative iii) below.
ii) Some SOC systems (Watkins et al, 1999b; Consolini, private communication, 1998) will destroy the information contained in their input. The scaling observed in their outputs is then independent of any present in the input, so any common scaling exponents between input and output are either coincidental or evidence of universality in certain confined plasma systems.
iii) Measuring global properties is the wrong thing to do, i.e. SOC is not an aspect of the global magnetosphere but relevant more locally to the magnetotail (compatible with picture B, and alternative i) above). This possibility is likely to be illuminated by further studies of SOC as a magnetotail model.
or iv) that another type of model is required (e.g. Chapman, 1999).
It is also important to emphasise that the extent to which SOC is observable, and distinguishable from other nonlinear physics paradigms (such as those presently used to study turbulence) is an important generic question in contemporary physics (including but going beyond, plasma physics). The diversity and quality of the existing ground-based and space-based magnetospheric databases provide a key testbed with which these intrinsically interdisciplinary questions can now be addressed; while ongoing investigations in astrophysical and laboratory confinement systems, both in plasma physics and elsewhere will continue to be applicable to the question of magnetospheric SOC.
The authors wish to acknowledge many enjoyable and valuable discussions with John Barrow, Ben Carreras, Tom Chang, Giuseppe Consolini, Patrick Diamond, Per Helander, Henrik Jensen, John King, David Newman, Christophe Rhodes, David Riley, George Rowlands, Jouni Takalo, Sunny Tam, Dave Tetreault, Vadim Uritsky, Zoltan Vörös and Dave Willis. NWW and SCC would like to acknowledge the hospitality of the MIT Center for Space Research where some of this work was carried out. SCC acknowledges a PPARC Lecturer fellowship and ROD the support of Euratom and the UK DTI.
| 1.Model: | 2.Linear | 3.Low -dimensional | 4.fBm | 5.SOC |
| --- | --- | --- | --- | --- |
| | (plus noise) | nonlinear | nonlinear | sandpile |
| 1. Property | | | | |
| 2.Short term predictable | Yes | Yes | No | ? |
| 3.Long term predictable | Yes | No | No | ? |
| 4.Global Scaling | No | Sometimes | Yes | Yes |
| 5.Scaling Regions | Sometimes | Sometimes | Yes | Yes |
| 6.Low G-P Dimension | No | Yes | Sometimes | ? |
| 7. Avalanches | No | ? | ? | Yes |
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# The Impact of Rotation on Cluster Dynamics
## 1. Background
Star clusters are self-gravitating Newtonian systems of choice where to brew complex gravitational dynamics (Meylan & Heggie 1997 for a review). Observations of old, globular, stellar clusters have led to the formulation of spherically symmetric dynamical models of equilibria. The most successful and universally studied one-integral spherical models are the King (1966) profiles. However, surveys of up to 100 Milky Way clusters have found small but significant departures from spherical symmetry (White & Shawl 1987): fits to their projected isophotes yield ellipticities $`<ϵ>1<a/b>0.07\pm 0.01`$. A study of 173 clusters in M31 found $`<ϵ>=0.09\pm 0.04`$ (Staneva, Spassova & Golev 1996). Observations of young clusters in the Large Magellanic Cloud revealed isophotal contours with ellipticities as large as $`ϵ=0.3`$ (Elson, Fall & Freeman 1987; Kontizas et al. 1990). This raises the possibility that clusters are formed as strongly flattened structures which then evolve towards rounder configurations (cf. Frenk & Fall 1982; Boily, Clarke & Murray 1999; Theis & Spurzem 1999), and brings up important theoretical issues concerning processes which may drive this evolution.
Rotation stretches any stellar association along a preferred axis: observations of the clusters $`\omega `$Centauri and M13 have shown that they are flattened by rotation (Meylan & Mayor 1986; Merritt, Meylan & Mayor 1997; Lupton, Gunn & Griffin 1987). Thus angular momentum, measured or possibly lost during evolution, offers a way to account for the morphology of clusters. Yet to date there are few evolutionary models of clusters with initial angular momentum. We have started on a project to develop three-dimensional dynamical models of rotating star clusters. In this articles results for two n-body models of isolated clusters are presented. Previous modelling of rotating clusters is reviewed first.
## 2. Gas and Fokker-Planck Models of rotating clusters
Agekian (1958) considered the effects of angular momentum diffusion on the equilibria of rotating fluid masses of uniform density. In his analysis, concentric spheroids rotating about their minor axis become rounder in time when the spheroids have initially an ellipticity $`ϵ=1a/b0.735`$, where $`a`$ and $`b`$ are the minor and major axes. Shapiro & Marchant (1976) integrated the equations of motion for this fluid in the limit of adiabatic (slow) diffusion of momentum. Angular momentum losses are driven by mass elements moving in the direction of the stream leaving the system at a rate higher than those moving in the opposite direction. The energy required for escape comes from ‘heat’, attributed to two-body encounters. Thus, angular momentum losses are accrued over a local two-body relaxation timescale, $`t_{\mathrm{col}}`$, which is inversely proportional to the mass density $`(1/\rho _{})`$. In practice this hinders applications of the results to actual clusters, which show centrally peaked density profiles (Meylan & Heggie 1997). Nevertheless, the framework set by Agekian provides a start in linking rotating bodies and observed (non-rotating) globular clusters.
With zero rotation, the central region of a cluster evolves towards a cusp in density during what is known as the gravothermal catastrophe. Does rotation stop the formation of a cusp? Hachisu (1979, 1982) discussed the time-evolution of self-gravitating cylindrical distributions of gas with angular momentum. He predicted a runaway collapse of the central region whenever angular momentum is expelled faster than a critical rate (see also Lagoute & Longaretti 1996). Hachisu dubbed this the ‘gravo-gyro catastrophe’, by analogy with the non-rotating case. These were until recently the only evolutionary models of rotating clusters. Two-dimensional orbit-averaged Fokker-Planck methods have now also been developed to address this issue. Following Goodman’s (1983) approach, Einsel & Spurzem (1999) integrated the Fokker-Planck equation in energy-momentum space $`[E,J_z]`$. Their initial configurations are truncated King models with added bulk motion. This velocity field takes the form of a Maxwellian distribution such that the mean velocity scales in proportion to radius away from the centre, then drops off at large radii. Their adopted axisymmetric distribution function (cf. Lupton, Gunn & Griffin 1987)
$$f(E,J_z)\mathrm{exp}\left(\beta \mathrm{\Omega }_oJ_z\right)\left[\mathrm{exp}\left(\beta E\right)1\right]$$
(1)
where $`\beta `$ is the inverse square central velocity dispersion and $`\mathrm{\Omega }_o`$ an angular velocity. The initial conditions are fixed by specifying the dimensionless parameters
$$\omega _o=\mathrm{\Omega }_o/\sqrt{9G\rho _c/(4\pi )}\mathrm{and}W_o,$$
ie, the scales of angular momentum and gravitational potential, respectively. The latter is the King parameter.
In the 2D Fokker-Planck models core-collapse proceeds on a much shorter timescale than in the non-rotating case, confirming Hachisu’s early intuition. However the central angular velocity does not increase at the high rates expected during the on-set of a gravo-gyro catastrophe; however near the end of core-collapse the central velocity dispersion bears the same relation to the central density as in the non-rotating self-similar collapse. This leaves open the question of what controls the final phase of evolution in these systems, ie whether or not rotation truly survives up to core-collapse. We chose to approach this problem using three-dimensional numerical integration; the setup is summarised below, followed by results and a discussion.
## 3. Basic properties
Self-consistent n-body realisations of the distribution function were obtained from the equilibrium Fokker-Planck code FOPAX developed by Christian Einsel. The models are fully specified once values are assigned to $`(W_o,\omega _o)`$. Figure 1 illustrates the properties of a set of models of 10,000 particles with $`W_o=6.0`$ and four values of $`\omega _o`$. The model clusters rotate about the z-axis and the equator lies in the x-y plane of a Cartesian coordinate system. Rotation causes the cluster in equilibrium to flatten down the z-axis and this is shown from computing the components of the inertia tensor $`I_{ij}`$ for a series of twenty concentric spherical shells of equal mass $`\mathrm{d}M`$. We define
$$\eta [r_k]1\frac{2I_{zz}}{I_{xx}+I_{yy}},\mathrm{particles}\mathrm{in}r_k\mathrm{d}r<r<r_k+\mathrm{d}r.$$
(2)
The parameter $`\eta =0`$ when the mass within a shell is distributed isotropically; $`\eta <0`$ (or, $`>0`$) when the distribution is anisotropic oblate (or, prolate). For the spherical model $`\omega _o=0`$ we found indeed near-zero values of $`\eta `$ at all radii.
Models with rotation have $`\omega _o0`$ and a range of values for $`\eta `$ increasing with it. Note that all models, save one with $`\omega _o=0.8`$, have values of $`\eta `$ compatible with sphericity at the centre. At larger radii, the models are all distinguished from one another.
Fast-rotating models need be more compact in order to sustain the accrued centrifugal force,
$$\mathrm{centrifugal}\mathrm{force}=\frac{v_\varphi ^2}{r}=r\mathrm{\Omega }^2,$$
which must always be smaller than the gravitational force, giving the condition
$$\mathrm{\Omega }^2\frac{GM}{r_s^3}.$$
(3)
Thus at constant mass $`M`$ the system radius $`r_s`$ must be smaller to allow for larger angular speed $`\mathrm{\Omega }`$. This is illustrated on figure 1, which displays $`\eta `$ and $`\mathrm{\Omega }`$ computed from the same set of particles. All models show $`\mathrm{\Omega }`$ decreasing with radius. Note that the curves are consistent with solid-body rotation in the core-region. Further out $`\mathrm{\Omega }`$ declines to near-zero, in a trend opposite that of $`\eta `$. The core remains roundish despite the large angular speed because the gravity is relatively stronger there than near the edge, and so random motion of the particles dominate over streaming motion. Overall the fraction of kinetic energy invested in streaming motion ranges from 0% to 4%, 14% and 26% in increasing order of $`\omega _o`$. For comparisons, the cluster $`\omega `$Centauri invests perhaps as much as 22 % of its kinetic energy in rotation (Merritt, Meylan & Mayor 1997).
## 4. N-body simulations
The code NBODY6++ is an Aarseth-type integration code based on a Hermite expansion of the variables in time (Aarseth 1999). It has been ported to parallel architecture (Spurzem 2000); the calculations were performed on CRAY computers linked up with MPI library. The code treats particles as point-masses and stellar evolution options were switched off. The chain-regularisation algorithm for hierarchical stellar encounters as well as the standard ‘KS’ regularisation (Mikkola & Aarseth 1998; Aarseth 1999) ensures high-precision integration during close interactions. Only simulations with N = 5,000 equal-mass particles will be discussed. There are no external tides.
Figures 2 & 3 illustrate the time-evolution of the models. The central density, total angular momentum and mean and core radii are plotted as function of time in units of the two-body relaxation time $`t_{\mathrm{col}}`$ (see Meylan & Heggie 1997; Casertano & Hut 1985). (Note: the lengths were normalised to their initial values.) The top panels show evolution for the case of $`\omega _o=0.5`$, the bottom set for $`\omega _o=0.8`$. Looking at these diagrams we find an evolution of the central density similar to the standard case with no rotation: the contraction of the central region leads to more close encounters and ejection, hence further contraction ensues, etc, until $`t5t_{\mathrm{col}}`$ when the density peaks sharply, indicating core-collapse: at the end of the simulations $`r_c0.03`$ and $`0.015`$, respectively, for $`\omega _o=0.5`$ and 0.8. At constant energy, core-contraction drives the expansion of the outer envelope and hence the mean radius expands rapidly at core-collapse. Note a subtle but noticeable difference between the two simulations, namely that the cluster with initially more rotation evolves faster; this is particularly visible in a comparison of radii at fixed time. A more convincing demonstration of the fast evolution of such clusters follows if we recall that in this unit of time, clusters without rotation reach core-collapse in around $`12t_{\mathrm{col}}`$, which is more than twice as long.
To appreciate how many stars might be lost to galactic tides, were a tidal field present, we imagine the cluster orbiting the galaxy on a circular orbit. A tidal radius may then be defined from the initial configuration, by computing the radius $`r_t`$ at which the mean density at time $`t`$ equates the initial mean density: $`r_t(t)=<r>[0]\times (M[t]/M[0])^{1/3}`$. If we label as escapers all stars found outside $`2r_t(t)`$, we obtain an estimate of the number of stars likely to leave the cluster on a timescale short compared with $`t_{\mathrm{col}}`$; the angular momentum they carry with them is deduced from summing up all the momenta of the stars left behind, and comparing with the initial value. Implementing this algorithm, we found the run with $`\omega =0.5`$ (top panels) would have lost 3.7% of the initial angular momentum over the time of evolution, but only 1.0% (51:5000) of its mass. The second model, with more rotation, would have lost 0.96% (48:5000) of its mass, but only 1.4% of its total momentum to such escapers. This shows how the cluster redistributes angular momentum efficiently within itself, such that the core evolves towards core-collapse despite added rotational support. Evolution in the core is faster, the faster the core rotates initially, since the cluster is more compact (cf. Eq. 3 and Fig. 1) which speeds up two-body effects.
Figure 3 graphs the specific angular momentum of the $`\omega _o=0.8`$ cluster as a function of radius $`r`$ for three different times. For comparison, two components are given: the z-axis component about which the cluster rotates; and the x-axis component. Dividing the cluster in ten concentric shells, we computed $`𝑳=𝒓\times 𝒗`$ and summed up the momenta in each shell: the result is the series of black squares shown on the figure. Initially the (net) z-angular momentum increases steadily from the centre, outwards; the symmetry of the figure would make the sum over $`L_x`$’s cancel out, and the black squares have been left out for this quantity. Evolution is monotonic, with the net z-momenta inside $`r=2`$pc decreasing, from which we deduce that an increasing fraction of the momentum is transferred to the volume $`>2`$pc. Notice on figure 3 that the stars form a core around $`r=0.5`$pc in the final stage of the simulation (right-most panels). It is not clear whether this is the result of an $`m=1`$ (lopsided) instability, attributable to the dynamics of the system (from a d.f. point of view), or a case of core-wandering, likely due to the small number of particles inside $`r=0.5`$pc (Sweatman 1993).
## 5. Conclusion
The faster evolution of clusters with rotation has been illustrated with two sample runs. The time to core-collapse we found from three-dimensional n-body simulations are in agreement with two-dimensional Fokker-Planck calculations (Einsel & Spurzem 1999): the collapse time of 5.4 $`t_{\mathrm{col}}`$ obtained for the $`\omega _o=0.8`$ agrees with the Fokker-Planck solution of 5.6 $`t_{\mathrm{col}}`$ for these parameters. This increases confidence in the results up to core collapse, obtained with two different algorithms.
### Acknowledgments.
Operating grants from the Neumann Institute for Computing, Jülich, and the Centre for High-Performance Computing, Stuttgart, awarded to R. Spurzem (ARI) are gratefully acknowledged. CMB was funded by research grant A/99/49003 awarded by the German DAAD in 1999. Thanks to P. Kroupa for commenting on some aspects of this paper.
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Aarseth, S.J. 1999, Cel. Mech. & Dyn. Ast., 73, 127
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## Discussion
A.Eckart: What happens to the rotation of the cluster after core-collapse?
C.M.Boily: I stopped my simulations precisely at core-collapse. What you measure in the envelope depends a lot on e.g. the galactic tide, which is absent in these simulations. In the core region I expect evolution will proceed much as in the standard case without rotation, since the cusp is isotropic and carries little net momentum. More detailed modelling is needed here.
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# Kadec-Pełczyński Decomposition for Haagerup 𝐿^𝒑-spaces
## 1. Introduction
In , Kadec and Pełczyński proved a fundamental property that if $`1p<\mathrm{}`$ then every bounded sequence $`(f_n)`$ in $`L^p[0,1]`$ has a subsequence that can be decomposed into two extreme sequences $`(g_k)`$ and $`(h_k)`$, where the $`h_k`$’s are pairwise disjoint, the $`g_k`$’s are $`L_p`$-equi-integrable that is $`\underset{m(A)0}{lim}\underset{k}{sup}\chi _Ag_k_p0`$ and $`h_kg_k`$ for every $`k1`$. This result is generally known as the Kadec-Pełczyński subsequence decomposition and has been investigated by several authors for the cases of Banach lattices and symmetric spaces (see for instance and ).
Motivated by the characterization of relatively weakly compact subsets of preduals of von Neumann algebras by Akemann , the above decomposition was studied in for non-commutative $`L^1`$-spaces associated with semi-finite von Neumann algebras equipped with distinguished, faithful, normal, semi-finite traces. A more general situation on $`E(,\tau )`$, where $`E`$ is a symmetric space of functions on $`(0,\mathrm{})`$ and $``$ is a semi-finite von Neumann algebra, was studied in . In particular, the result in was generalized for $`L^p(,\tau )`$ for all $`0<p<\mathrm{}`$.
The aim of the present paper is to provide extensions of the Kadec-Pełczyński decomposition theorem for general von Neumann algebras which are not necessarily semi-finite. There are many different methods of constructions of non-commutative $`L^p`$-spaces associated with arbitrary von Neumann algebras; for instance, those of Araki-Masuda , Haagerup , Hilsum , Izumi , Kosaki , Terp and many others. But it is known that, for a given von Neumann algebra $``$ and a fixed index $`p`$, all these $`L^p`$-spaces are isometrically isomorphic. We will consider Haagerup’s $`L^p`$-spaces since they can be viewed as spaces of operators that can be embedded as subspaces of symmetric spaces of measurable operators obtained from semi-finite von Neumann algebras via crossed product (see a brief description below). Our main result is Theorem 4.1 which roughly says that any bounded sequence in $`L^p()`$ has a subsequence that can be splitted into two sequences; one is uniformly integrable and the other consists of elements supported by decreasing projections that converges to zero. Our initial motivation is the case $`p=1`$ where $`L^1()`$ can simply be viewed as the predual of $``$. This case allows us to get informations on copies of $`\mathrm{}^1`$ in duals of $`C^{}`$-algebras. It has been known that every non-reflexive subspace of duals of $`C^{}`$-algebras contains complemented copies of $`\mathrm{}^1`$ . On the other hand, Dowling and Lennard showed in that for $`L^1[0,1]`$, these complemented copies can be chosen to be asymptotically isometric.
Using the main decomposition for the case $`p=1`$, we can conclude that every non reflexive subspace of duals of $`C^{}`$-algebras contains sequences that generate complemented copies of $`\mathrm{}^1`$ and are asymptotically isometric. As in and , these asymptotically isometric copies of $`\mathrm{}^1`$ yield self maps on convex bounded sets that fail to have any fixed points. These lead to a more general structural consequence that non-reflexive subspaces of duals of $`JB^{}`$-triples fail the fixed point property for self-maps on closed bounded convex sets.
The paper is organized as follows: in Section 2 below, we set some preliminary background on Haagerup $`L^p`$-spaces. In particular, we provide a brief discussion on its connection to the semi-finite case and define the notion of uniformly integrable sets in these $`L^p`$-spaces. Section 3 is devoted to the proof a key result which is essentially the crusial part of the paper. We present our main results in Section 4 and finally, Section 5 is where we provide all the applications on copies of $`\mathrm{}^1`$ on duals of $`C^{}`$ algebras and $`JB^{}`$-triples.
Our notation and terminology are standard as may be found in for Banach spaces, and for operator algebras.
## 2. Non-commutative $`L^p`$-spaces
In this section, we will describe different spaces involved and discuss some properties that will be crusial for the presentation. We will begin from the semi-finite case. We denote by $`𝒩`$ a semi-finite von Neumann algebra on a Hilbert space $``$, with a distinguished normal, faithful semi-finite trace $`\tau `$. The identity in $`𝒩`$ will be denoted by $`\mathrm{𝟏}`$. A closed and densely defined operator $`a`$ on $``$ is said to be affiliated with $`𝒩`$ if $`ua=au`$ for all unitary operator $`u`$ in the commutant $`𝒩^{}`$ of $`𝒩`$.
A closed and densely defined operator $`x`$, affiliated with $`𝒩`$, is called $`\tau `$-measurable if for every $`\epsilon >0`$, there exists an orthogonal projection $`p𝒩`$ such that $`p()\text{dom}(x)`$, $`\tau (\mathrm{𝟏}p)<\epsilon `$ and $`xp`$. The set of all $`\tau `$-measurable operators will be denoted by $`\stackrel{~}{𝒩}`$. The set $`\stackrel{~}{𝒩}`$ is a $``$-algebra with respect to the strong sum, the strong product and the adjoint operation. Given a self-adjoint operator $`x`$ in $`\stackrel{~}{𝒩}`$ and $`B`$ a Borel subset of $``$, we denote by $`\chi _B(x)`$ the projection $`_B1𝑑e^x`$ where $`e^x()`$ is the spectral measure of $`x`$. For fixed $`x\stackrel{~}{𝒩}`$ and $`t0`$, we recall
$$\mu _t(x)=inf\{s0:\tau (e^{|x|}(s,\mathrm{}))t\}.$$
The function $`\mu _{(.)}(x):[0,\mathrm{})[0,\mathrm{}]`$ is called the generalized singular value function (or decreasing rearrangement) of $`x`$. For a complete study of $`\mu _{(.)}`$, we refer to .
If $`E`$ is a symmetric (r.i. for short) quasi-Banach function space on $`^+`$, the symmetric space of measurable operators $`E(𝒩,\tau )`$ is defined by setting
$$E(𝒩,\tau ):=\{x\stackrel{~}{}:\mu (x)E\}$$
and
$$x_{E(𝒩,\tau )}=\mu (x)_E\text{for all}xE(𝒩,\tau ).$$
The space $`E(𝒩,\tau )`$ is a (quasi) Banach space and is often referred to as the non-commutative version of the (quasi) Banach function space $`E`$. We remark that if $`0<p<\mathrm{}`$ and $`E=L^p(^+,m)`$ then $`E(𝒩,\tau )`$ coincides with the usual non-commutative $`L^p`$-space associated to the semi-finite von Neumann algebra $`𝒩`$. We refer to , and for extensive background on the space $`E(𝒩,\tau )`$.
We now provide a short description of the Haagerup $`L^p`$-spaces. Let assume that $``$ is a general von Neumann algebra (not necessarily semi-finite). Let $`𝒩`$ be the crossed product of $``$ by the modular automorphism group $`(\sigma _t)_t`$ of a fixed semi-finite weight on $``$. The von Neumann algebra $`𝒩`$ admits the dual action $`(\theta _s)_s`$ and a normal faithful semi-finite trace $`\tau `$ satisfying, $`\tau \theta _s=e^s\tau `$, $`s`$. For $`1p<\mathrm{}`$, the Haagerup $`L^p`$-space associated with $``$ is defined by
$$L^p():=\{x\stackrel{~}{𝒩}:\theta _s(x)=e^{s/p}x,s\}.$$
It is well known that there is a linear order isomorphism $`\phi h_\phi `$ from $`_{}`$ onto $`L^1()`$. One can define a positive linear functional $`Tr`$ on $`L^1()`$ by setting
$$Tr(h_\phi )=\phi (\mathrm{𝟏}),\phi _{}.$$
For $`1p<\mathrm{}`$, the spaces $`L^p()`$ are Banach spaces with the norm defined by
$$x_p=\left(Tr(|x|^p)\right)^{\frac{1}{p}},\text{for}xL^p().$$
For complete details on the construction of $`L^p()`$, we refer to . Also, it was shown in \[12, Lemma 4.8\] that if $`xL^p()`$, $`1p<\mathrm{}`$, then
$$\mu _t(x)=t^{1/p}x_p,t>0.$$
where the singular value is relative to the canonical trace on $`𝒩`$.
We recall that if $`1p<\mathrm{}`$, then the Lorentz space $`L^{p,\mathrm{}}(^+,m)`$ is the set of (class of) all Lebesgue measurable functions on $`^+`$ with the norm
$$f_{p,\mathrm{}}=\underset{t>0}{sup}\{t^{1/p}\mu _t(f)\}.$$
It is well known that if $`1<p<\mathrm{}`$, then the space $`L^{p,\mathrm{}}(^+,m)`$ equipped with the equivalent Calderon norm given by
$$f_{(p,\mathrm{})}=\underset{t>0}{sup}\left\{t^{1/p1}_0^t\mu _s(f)𝑑s\right\},fL^{p,\mathrm{}}(^+,m),$$
is a symmetric Banach function space on $`^+`$ with the Fatou property. The following proposition is an immediate consequence of the above remarks.
###### Proposition 2.1.
If $`1<p<\mathrm{}`$, then the space $`L^p()`$ is a closed subspace of the symmetric space of measurable operators $`L^{p,\mathrm{}}(𝒩,\tau )`$. Moreover if $`1/q+1/p=1`$, then
$$x_p=qx_{(p,\mathrm{})}$$
for all $`xL^p()`$.
Let us now extend the notion of uniform integrability to the Haagerup $`L^p`$-spaces. Following , we define uniform integrability in $`L^p()`$ as in Akemann’s characterization of relatively weakly compact subsets of $`_{}`$.
###### Definition 2.2.
Let $`1p<\mathrm{}`$ and $`K`$ be a bounded subset of $`L^p()`$. We say that $`K`$ is uniformly integrable if $`\underset{n\mathrm{}}{lim}\underset{\phi K}{sup}e_n\phi e_n=0`$ for every decreasing sequence $`(e_n)_n`$ of projections in $``$ with $`e_n_n0`$.
We note that for $`p=1`$, a subset $`K`$ is uniformly integrable $`L^1()`$ if and only if it is relatively weakly compact.
Throughout, $`𝒟`$ denotes the set of all sequences of decreasing projections in $``$ that converges to zero; $`𝒟:=\left\{(e_n)_n;\text{the }e_n\text{’s are projections in }\text{ and}e_n_n0\right\}`$. Also for any subset $`K`$ of $`L^p()`$, $`|K|`$ denotes the set of all modudi of elements of $`K`$; $`|K|:=\left\{|x|;xK\right\}`$.
Fact 1. If $`xL^p()`$ and $`yL^p()`$ are such that $`xy`$ (i.e. ($`\mathrm{supp}x)(\mathrm{supp}y)`$) then $`x+y^p=x^p+y^p`$.
###### Proof.
$`x^p=Tr(|x|^p)`$ and if $`xy`$ as elements of $`𝒩`$, $`|x+y|^p=|x|^p+|y|^p`$ and therefore $`x+y^p=Tr(|x|^p+|y|^p)=x^p+y^p`$. ∎
Fact 2. If $`xL^p()`$ and $`e`$ is a projection in $``$, then $`x^pexe^p+(1e)x(1e)^p`$.
###### Proof.
Set $`u=2e1`$. It is clear that $`u`$ is unitary and $`exe+(1e)x(1e)=\frac{1}{2}\left(x+uxu^{}\right)`$. It follows that $`exe+(1e)x(1e)^px^p`$ and hence $`exe^p+(1e)x(1e)^px^p`$. ∎
We finish this section with the following two lemmas which can be proved using similar arguments as in the semi-finite case (, ) and will be used in the sequel. Details are left to the readers.
###### Lemma 2.3.
Let $`1p<\mathrm{}`$, $`(p_n)_n𝒟`$ and $`K`$ be a bounded subset of $`L^p()`$ such that for each $`n_01`$, the sets $`(\mathrm{𝟏}p_{n_0})K`$ and $`|K(\mathrm{𝟏}p_{n_0})|`$ are uniformly integrable. Then $`K`$ is uniformly integrable if and only if $`\underset{n\mathrm{}}{lim}sup_{\phi K}p_n\phi p_n=0.`$
###### Lemma 2.4.
Let $`1p<\mathrm{}`$, $`(\phi _n)_n`$ be a bounded sequence in $`L^p()`$ and $`(p_n)_n𝒟`$. Assume that $`\underset{n\mathrm{}}{lim}\underset{k}{sup}p_n\phi _kp_n=\gamma >0`$ then there exists a subsequence $`(\phi _{k_n})`$ so that $`\underset{n\mathrm{}}{lim}p_n\phi _{k_n}p_n=\gamma `$.
## 3. Preliminary Results
This section is devoted to the proof of Theorem 3.1 below which is the key result that we will use to prove our main theorem. We remark that the case of finite von-Neumann algebras can be obtained with minor changes from the proof of the commutaive case (see ).
###### Theorem 3.1.
Let $``$ be a $`\sigma `$-finite von-Neumann algebra and $`1p<\mathrm{}`$. Assume that $`K`$ is a subset of the positive part of the unit ball of $`L^p()`$ that is not uniformly integrable. Then there exists a sequence $`(\phi _n)_nK`$ and $`(f_n)_n𝒟`$ such that:
$$sup\left\{\underset{n\mathrm{}}{lim}\underset{k}{sup}e_n\phi _ke_n;(e_n)_n𝒟\right\}=\underset{n\mathrm{}}{lim}\underset{k}{sup}f_n\phi _kf_n>0.$$
###### Lemma 3.2.
Let $`𝒩`$ be a semi-finite von Neumann algebra with distinguished faithful normal semi-finite trace $`\tau `$ as above and $`E`$ be a symmetric quasi-Banach function space on $`(0,\mathrm{})`$. If $`xE(𝒩,\tau )`$ and $`u𝒩`$ then
$$xu_Ex_E^{\frac{1}{2}}u^{}|x|u_E^{\frac{1}{2}}x_E^{\frac{3}{4}}uu^{}|x|uu^{}_E^{\frac{1}{4}}.$$
###### Proof.
Let $`x=v|x|`$ be the polar decomposition of $`x`$. Then
$$\begin{array}{cc}\hfill xu=v|x|u& |x|u=|x|^{\frac{1}{2}}|x|^{\frac{1}{2}}u\hfill \\ & |x|^{\frac{1}{2}}_{E^{(2)}}|x|^{\frac{1}{2}}u_{E^{(2)}}\hfill \\ & =x_E^{\frac{1}{2}}u^{}|x|u_E^{\frac{1}{2}}.\hfill \end{array}$$
For the second inequality,
$$\begin{array}{cc}\hfill xu& x_E^{\frac{1}{2}}|x|^{\frac{1}{2}}u_{E^{(2)}}\hfill \\ & =x_E^{\frac{1}{2}}u^{}|x|^{\frac{1}{2}}_{E^{(2)}}\hfill \\ & x_E^{\frac{1}{2}}|x|^{\frac{1}{2}}uu^{}|x|^{\frac{1}{2}}_E^{\frac{1}{2}}\hfill \\ & x_E^{\frac{1}{2}}\left(|x|^{\frac{1}{2}}uu^{}_{E^{(2)}}|x|^{\frac{1}{2}}_{E^{(2)}}\right)^{\frac{1}{4}}\hfill \\ & =x_E^{\frac{1}{2}}uu^{}|x|uu^{}_E^{\frac{1}{4}}x_E^{\frac{1}{4}}\hfill \\ & =x_E^{\frac{3}{4}}uu^{}|x|uu^{}_E^{\frac{1}{4}}.\hfill \end{array}$$
Lemma 3.2 shows in particular that if $`1p<\mathrm{}`$, $`xL^p()`$ and $`u`$ then
$$xux^{\frac{3}{4}}uu^{}|x|uu^{}^{\frac{1}{4}}.$$
###### Lemma 3.3.
Let $`\gamma >0`$ and $`(\phi _k)_k`$ be a sequence in the positive part of the unit ball of $`L^p()`$. If there exists a sequence $`(a_n)_n`$ in the unit ball of $``$ with $`a_n_n0`$ and such that $`lim_n\mathrm{}sup_ka_n\phi _ka_n\gamma `$. Then for every $`\epsilon >0`$, there exists a sequence $`(s_n)_n`$ of projections with:
* $`s_ns_1`$ for every $`n1`$;
* $`s_n0`$ for the strong operator topology;
* for every $`n_0`$, $`lim_n\mathrm{}sup_k(s_{n_0}a_ns_{n_0})\phi _k(s_{n_0}a_ns_{n_0})\gamma \epsilon `$.
###### Proof.
Fix $`\delta >0`$ with $`\delta (\epsilon /8)^2`$ and define the sequence of projections as follows:
$$\{\begin{array}{cc}s_1:=\chi _{(\delta ,1)}(a_1)\hfill & and\hfill \\ s_n:=\chi _{(\delta ,1)}(s_1a_ns_1)\hfill & forn2.\hfill \end{array}$$
Clearly $`s_n`$ is a subprojection of the support of $`s_1a_ns_1`$ so $`s_ns_1`$. Also $`\delta s_ns_n(s_1a_ns_1)s_n`$, and since $`s_n`$ and $`s_1a_ns_1`$ are commuting operators, $`\delta s_ns_n(s_1a_ns_1)s_ns_1a_ns_1`$ and therefore $`s_n0`$ so $`(i)`$ and $`(ii)`$ are verified.
Claim: Let $`n_0`$ and $`nn_0`$, for every $`\phi L^p()`$ with $`\phi 1`$, $`\phi a_n(1s_{n_0})2\delta ^{1/2}`$. Similarily, $`\phi (1s_{n_0})a_n2\delta ^{1/2}`$.
To see this claim, it is enough to notice that
$$\begin{array}{cc}\hfill \phi a_n(1s_{n_0})& \phi a_ns_1(1s_{n_0})+\phi a_n(1s_1)\hfill \\ & \phi a_ns_1(1s_{n_0})_{\mathrm{}}+\phi a_n(1s_1)_{\mathrm{}}\hfill \\ & (1s_{n_0})s_1a_n^2s_1(1s_{n_0})_{\mathrm{}}^{1/2}+(1s_1)a_n^2(1s_1)_{\mathrm{}}^{1/2}\hfill \\ & (1s_{n_0})s_1a_ns_1(1s_{n_0})_{\mathrm{}}^{1/2}+(1s_1)a_n(1s_1)_{\mathrm{}}^{1/2}\hfill \end{array}$$
and since $`(a_n)_n`$ is a decreasing sequence and $`1n_0n`$, we get
$$\phi a_n(1s_{n_0})(1s_{n_0})s_1a_{n_0}s_1(1s_{n_0})_{\mathrm{}}^{1/2}+(1s_1)a_1(1s_1)_{\mathrm{}}^{1/2}2\delta ^{1/2}.$$
A similar estimate can be established for $`\phi (1s_{n_0})a_n`$ which verifies the claim.
To complete the proof, let $`\phi L^p()`$, $`\phi 1`$. For $`nn_0`$, we can write $`a_n\phi a_n`$ as:
$$a_n\phi a_n=(s_{n_0}a_ns_{n_0})\phi a_n+s_{n_0}a_n(1s_{n_0})\phi a_n+(1s_{n_0})a_n\phi a_n$$
and using the claim above, $`a_n\phi a_n(s_{n_0}a_ns_{n_0})\phi a_n+4\delta ^{1/2}`$. A similar estimate would give $`(s_{n_0}a_ns_{n_0})\phi a_n(s_{n_0}a_ns_{n_0})\phi (s_{n_0}a_ns_{n_0})+4\delta ^{1/2}`$ and combining these two estimates, we get
$$a_n\phi a_n(s_{n_0}a_ns_{n_0})\phi (s_{n_0}a_ns_{n_0})+8\delta ^{1/2}.$$
This shows that
$$\underset{n\mathrm{}}{lim}\underset{k}{sup}(s_{n_0}a_ns_{n_0})\phi _k(s_{n_0}a_ns_{n_0})\gamma 8\delta ^{1/2}\gamma \epsilon .$$
The proof is complete. ∎
The next result shows that using projections in the definition of uniform integrability is not essential. One can use elements of the positive part of the unit ball of $``$.
###### Proposition 3.4.
Let $`\gamma >0`$ and $`(\phi _k)_k`$ be a sequence in the positive part of the unit ball of $`L^p()`$. If there exists a sequence $`(a_n)_n`$ in the unit ball of $``$ with $`a_n_n0`$ and such that $`lim_n\mathrm{}sup_ka_n\phi _ka_n\gamma `$. Then for every $`\epsilon >0`$, there exists a sequence $`(p_n)_n𝒟`$ with $`p_nsupp(a_1)`$ for all $`n1`$ and such that $`lim_n\mathrm{}sup_kp_n\phi _kp_n\gamma \epsilon `$.
###### Proof.
The sequence $`(p_n)`$ will be constructed inductively. Let $`(\epsilon _j)_j`$ be a sequence in the open interval $`(0,\epsilon )`$ such that $`_{j=1}^{\mathrm{}}\epsilon _j=\epsilon `$ and $`\omega _0`$ be a faithful state in $`_{}`$.
By Lemma 3.3, one can choose a sequence of projections $`(s_n^{(1)})_n`$ with $`s_n^{(1)}s_1^{(1)}`$ for every $`n1`$, $`s_n^{(1)}0`$ (as n tends to $`\mathrm{}`$) satisfying the conclusion of Lemma 3.3 for $`(a_n)_n`$, $`\gamma `$ and $`\epsilon _1`$.
Choose $`n_11`$ such that $`\omega _0\left(s_{n_1}^{(1)}\right)1/2`$. From (iii) of Lemma 3.3,
$$\underset{n\mathrm{}}{lim}\underset{k}{sup}(s_{n_1}^{(1)}a_ns_{n_1}^1)\phi _k(s_{n_1}^{(1)}a_ns_{n_1}^{(1)})\gamma \epsilon _1.$$
Reapplying Lemma 3.3, on $`(a_n^{(2)})_n=(s_{n_1}^{(1)}a_ns_{n_1}^{(1)})_n`$, $`\gamma \epsilon _1`$ and $`\epsilon _2`$, one would get a sequence of projections $`(s_n^{(2)})_n`$ with $`s_n^{(2)}s_{n_1}^{(1)}`$ for every $`n1`$, $`s_n^{(2)}0`$ (as n tends to infinity). As above, on can choose $`n_2`$ such that $`\omega _0\left(s_{n_2}^{(2)}\right)1/2^2`$ and
$$\underset{n\mathrm{}}{lim}\underset{k}{sup}(s_{n_2}^{(2)}a_ns_{n_2}^{(2)})\phi _k(s_{n_2}^{(2)}a_ns_{n_2}^{(2)})\gamma \epsilon _1\epsilon _2.$$
The induction is clear, repeating the argument above would give a decreasing sequence of projections $`s_{n_1}^{(1)}s_{n_2}^{(2)}\mathrm{}s_{n_j}^{(j)}\mathrm{}`$ so that for every $`j1`$, $`\omega _0\left(s_{n_j}^{(j)}\right)1/2^j`$ and
$$\underset{n\mathrm{}}{lim}\underset{k}{sup}(s_{n_j}^{(j)}a_ns_{n_j}^{(j)})\phi _k(s_{n_j}^{(j)}a_ns_{n_j}^{(j)})\gamma \underset{i=1}{\overset{j}{}}\epsilon _i.$$
If for every $`j1`$, we set $`p_j=s_{n_j}^{(j)}`$ then $`(p_j)_j`$ belongs to $`𝒟`$ and
$$\underset{k}{sup}(p_ja_jp_j)\phi _k(p_ja_jp_j)\gamma \underset{i=1}{\overset{j}{}}\epsilon _i$$
which shows that
$$\underset{j\mathrm{}}{lim}\underset{k}{sup}(p_ja_jp_j)\phi _k(p_ja_jp_j)\gamma \epsilon $$
and since $`p_ja_j_{\mathrm{}}1`$, the desired conclusion follows. ∎
###### Proposition 3.5.
Let $`K`$ be as in the statement of Theorem 3.1. There exists a sequence $`(\phi _k)_k`$ in $`K`$ such that
$$sup\left\{\underset{n\mathrm{}}{lim}\underset{k}{sup}e_n\phi _ke_n;(e_n)_n𝒟\right\}=sup\left\{\underset{¯}{\mathrm{lim}}_n\mathrm{}e_n\phi _ne_n;(e_n)_n𝒟\right\}>0.$$
###### Proof.
Set $`\alpha _0:=sup\left\{lim_n\mathrm{}sup_{\phi K}e_n\phi e_n;(e_n)𝒟\right\}`$ and let $`(\epsilon _j)_j`$ be a subset of the open interval $`(0,1)`$ such that $`\mathrm{\Pi }_{j=1}^{\mathrm{}}(1\epsilon _j)>0`$.
Since $`\alpha _0>0`$, one can choose a sequence $`(y_n)_n`$ in $`K`$ and $`(e_n^{(1)})_n𝒟`$ such that
$$\underset{n\mathrm{}}{lim}\underset{k}{sup}e_n^{(1)}y_ke_n^{(1)}\alpha _0(1\epsilon _1).$$
A further subsequence $`(y_k^{(1)})_k(y_k)`$ can be chosen so that
$$\underset{n\mathrm{}}{lim}e_n^{(1)}y_n^{(1)}e_n^{(1)}\alpha _0(1\epsilon _1).$$
Set $`\alpha _1:=sup\left\{\underset{n\mathrm{}}{lim}\underset{k}{sup}e_ny_k^{(1)}e_n;(e_n)𝒟\right\}`$. It is clear that $`\alpha _1\alpha _0(1\epsilon _1)`$ and as above a sequence $`(y_k^{(2)})_{k1}(y_k^{(1)})_k`$ can be chosen so that
$$\underset{n\mathrm{}}{lim}e_n^{(2)}y_n^{(2)}e_n^{(2)}\alpha _1(1\epsilon _2).$$
Inductively, one can construct sequences $`(y_n)_n(y_n^{(1)})_n(y_n^{(2)})_n\mathrm{}(y_n^{(j)})_n\mathrm{}`$ in $`K`$ and sequences $`(e_n^{(1)})_n,(e_n^{(2)})_n,\mathrm{},(e_n^{(j)})_n,\mathrm{}`$ in $`𝒟`$ so that for every $`j1`$,
$$\underset{n\mathrm{}}{lim}e_n^{(j)}y_n^{(j)}e_n^{(j)}\alpha _{j1}(1\epsilon _j).$$
Let $`(\phi _n)_n`$ be the diagonal sequence obtained from $`(y_n^{(j)})_n,j1`$. For every $`j1`$, $`(\phi _n)_{nj}`$ is a subsequence of $`(y_n^{(j)})_{n1}`$ so
$$\underset{¯}{\mathrm{lim}}_n\mathrm{}e_n^{(j)}\phi _ne_n^{(j)}\alpha _{j1}(1\epsilon _j)$$
and
$$sup\left\{\underset{n\mathrm{}}{lim}\underset{k}{sup}e_n\phi _ke_n;(e_n)𝒟\right\}\alpha _j.$$
We note that $`\alpha _{j1}\alpha _j\alpha _{j1}(1\epsilon _j)`$ so for every $`j1`$,
$$\begin{array}{cc}\hfill sup\left\{\underset{n\mathrm{}}{lim}\underset{k}{sup}e_n\phi _ke_n;(e_n)𝒟\right\}& \alpha _j\hfill \\ & \alpha _j(1\epsilon _{j+1})\frac{1}{1\epsilon _{j+1}}\hfill \\ & \frac{1}{1\epsilon _{j+1}}\underset{¯}{\mathrm{lim}}_n\mathrm{}e_n^{(j+1)}\phi _ne_n^{(j+1)}\hfill \end{array}$$
which implies that
$$sup\{\underset{n\mathrm{}}{lim}\underset{k}{sup}e_n\phi _ke_n;(e_n)_n𝒟\}\frac{1}{1\epsilon _{j+1}}sup\{\underset{¯}{\mathrm{lim}}_n\mathrm{}e_n\phi _ne_n||;(e_n)_n𝒟\}.$$
Taking the limit as $`j`$ goes to $`\mathrm{}`$,
$$sup\left\{\underset{n\mathrm{}}{lim}\underset{k}{sup}e_n\phi _ke_n;(e_n)_n𝒟\right\}sup\left\{\underset{¯}{\mathrm{lim}}_n\mathrm{}e_n\phi _ne_n;(e_n)_n𝒟\right\}.$$
The other inequality is trivial.
To check that $`sup\left\{\underset{¯}{\mathrm{lim}}_n\mathrm{}e_n\phi _ne_n;(e_n)_n𝒟\right\}>0`$, it is plain that
$$\underset{¯}{\mathrm{lim}}_n\mathrm{}e_n^{(j)}\phi _ne_n^{(j)}\alpha _{j1}(1\epsilon _j)\alpha _1\mathrm{\Pi }_{j=2}^{\mathrm{}}(1\epsilon _j)>0.$$
The proof of the proposition is complete. ∎
Proof of Theorem 3.1
Let $`(\phi _n)`$ be the sequence in $`K`$ obtained from Proposition 3.5 ; i.e. $`(\phi _n)_n`$ is the sequence in $`K`$ satisfying:
$$sup\left\{\underset{¯}{\mathrm{lim}}_n\mathrm{}e_n\phi _ne_n;(e_n)_n𝒟\right\}=sup\left\{\underset{n\mathrm{}}{lim}\underset{k}{sup}e_n\phi _ke_n;(e_n)_n𝒟\right\}:=\alpha >0$$
Claim: $`\alpha `$ is attained.
Assume the opposite i.e. for every $`(p_n)_n𝒟`$, $`\underset{¯}{\mathrm{lim}}_n\mathrm{}p_n\phi _np_n<\alpha `$.
Inductively, we will construct sequences of integers and projections in $``$ satisfying the following conditions:
(3.1)
$$m_1m_2\mathrm{}m_j\mathrm{}\text{a sequence in };$$
(3.2)
$$n_1<n_2<\mathrm{}<n_j<\mathrm{}\text{infinite sequence in };$$
sequences $`(p_n^{(1)})_{n1},(p_n^{(2)})_{n1},\mathrm{}`$ in $`𝒟`$ such that for every $`j2`$ and every $`n2`$,
(3.3)
$$p_n^{(j)}\underset{k=1}{\overset{j1}{}}p_{n_k}^{(k)};$$
if we set $`(f_n^{(1)})_{n1}=(p_n^{(1)})_{n1}`$ and
(3.4)
$$f_n^{(j)}=\{\begin{array}{cc}f_n^{(j1)}p_{n_{j1}}^{(j)}n<n_{j1}\hfill & \\ f_n^{(j1)}+p_n^{(j)}nn_{j1}\hfill & \end{array}$$
then
(3.5)
$$\underset{n\mathrm{}}{lim}\underset{k}{sup}f_n^{(j)}\phi _kf_n^{(j)}\alpha (1\frac{1}{2^{m_j1}}),$$
(3.6)
$$\underset{k}{sup}f_{n_j}^{(j)}\phi _kf_{n_j}^{(j)}<\alpha (1\frac{1}{2^{m_j}})$$
and
(3.7)
$$\underset{¯}{\mathrm{lim}}_n\mathrm{}f_n^{(j)}\phi _nf_n^{(j)}^p\underset{¯}{\mathrm{lim}}_n\mathrm{}f_n^{(j1)}\phi _nf_n^{(j1)}^p+\frac{\alpha ^{4p}}{(2^{m_{j1}+3})^{4p}}.$$
Fix a sequence $`(p_n^{(1)})_{n1}𝒟`$ such that $`\underset{¯}{\mathrm{lim}}_n\mathrm{}p_n^{(1)}\phi _np_n^{(1)}\alpha (1\frac{1}{2^2})`$ and choose $`m_1`$ such that
$$\alpha (1\frac{1}{2^{m_11}})\underset{n\mathrm{}}{lim}\underset{k}{sup}p_n^{(1)}\phi _kp_n^{(1)}<\alpha (1\frac{1}{2^{m_1}})$$
(such $`m_1`$ exists since $`\alpha `$ is not attained).
Assume that the construction is done for $`1,2,\mathrm{},(j1)`$. By the definition of $`\alpha `$, one can choose $`(q_n)_n𝒟`$ so that $`\underset{¯}{\mathrm{lim}}_n\mathrm{}q_n\phi _nq_n>\alpha (1\frac{1}{2^{m_{j1}+1}})`$. Writing $`q_n\phi _nq_n`$ in the form
$$q_n\phi _nq_n=q_nf_{n_{j1}}^{(j1)}\phi _nf_{n_{j1}}^{(j1)}q_n+q_nf_{n_{j1}}^{(j1)}\phi _n(1f_{n_{j1}}^{(j1)})q_n+q_n(1f_{n_{j1}}^{(j1)})\phi _nq_n,$$
one can see that
$$q_n\phi _nq_nf_{n_{j1}}^{(j1)}\phi _nf_{n_{j1}}^{(j1)}+2\phi _n(1f_{n_{j1}}^{(j1)})q_n.$$
Applying Lemma 3.2 for $`x=\phi _n`$ and $`u=(1f_{n_{j1}}^{(j1)})q_n`$, we get
$$\begin{array}{cc}\hfill q_n\phi _nq_n& f_{n_{j1}}^{(j1)}\phi _nf_{n_{j1}}^{(j1)}+\hfill \\ & 2\phi _n^{\frac{3}{4}}(1f_{n_{j1}}^{(j1)})q_n(1f_{n_{j1}}^{(j1)})\phi _n(1f_{n_{j1}}^{(j1)})q_n(1f_{n_{j1}}^{(j1)})^{\frac{1}{4}}.\hfill \end{array}$$
Applying (3.6) for $`(j1)`$ gives
$$q_n\phi _nq_n\alpha (1\frac{1}{2^{m_{j1}}})+2(1f_{n_{j1}}^{(j1)})q_n(1f_{n_{j1}}^{(j1)})\phi _n(1f_{n_{j1}}^{(j1)})q_n(1f_{n_{j1}}^{(j1)})^{\frac{1}{4}}.$$
Taking the limit (as $`n`$ tends to $`\mathrm{}`$),
$$\alpha (1\frac{1}{2^{m_{j1}+1}})\alpha (1\frac{1}{2^{m_{j1}}})+2\underset{¯}{\mathrm{lim}}_n\mathrm{}(1f_{n_{j1}}^{(j1)})q_n(1f_{n_{j1}}^{(j1)})\phi _n(1f_{n_{j1}}^{(j1)})q_n(1f_{n_{j1}}^{(j1)})^{\frac{1}{4}}.$$
which implies that
$$\underset{¯}{\mathrm{lim}}_n\mathrm{}(1f_{n_{j1}}^{(j1)})q_n(1f_{n_{j1}}^{(j1)})\phi _n(1f_{n_{j1}}^{(j1)})q_n(1f_{n_{j1}}^{(j1)})\frac{\alpha ^4}{(2^{m_{j1}+2})^4}.$$
If we set $`a_n^{(j)}=(1f_{n_{j1}}^{(j1)})q_n(1f_{n_{j1}}^{(j1)})`$ then $`a_n^{(j)}_n0`$ and
$$\underset{¯}{\mathrm{lim}}_n\mathrm{}a_n^{(j)}\phi _na_n^{(j)}\frac{\alpha ^4}{(2^{m_{j1}+2})^4}.$$
Applying Proposition 3.4 for $`(\phi _n)_n`$, $`(a_n^{(j)})_n`$, $`\gamma =\frac{\alpha ^4}{(2^{m_{j1}+2})^4}`$ and $`\epsilon =\frac{\alpha ^4}{(2^{m_{j1}+2})^4}\frac{\alpha ^4}{(2^{m_{j1}+3})^4}`$ would provide a sequence $`(p_n^{(j)})𝒟`$ such that
$$\underset{¯}{\mathrm{lim}}_n\mathrm{}p_n^{(j)}\phi _np_n^{(j)}\frac{\alpha ^4}{(2^{m_{j1}+3})^4}.$$
Since $`p_n^{(j)}\mathrm{supp}(a_1^{(j)})1f_{n_{j1}}^{(j1)}`$, it is clear that $`p_n^{(j)}f_{n_{j1}}^{(j1)}`$ for every $`n1`$ so (3.3) is verified.
If we define $`(f_n^{(j)})`$ as in (3.4) then appropriate $`m_jm_{j1}`$ and $`n_j>n_{j1}`$ can be choosen so that (3.5) and (3.6) would be satisfied.
Now since $`p_n^{(j)}+f_n^{(j1)}=f_n^{(j)}`$ for $`nn_j`$,
$$\begin{array}{cc}\hfill f_n^{(j)}\phi _nf_n^{(j)}^p& f_n^{(j1)}\phi _nf_n^{(j1)}^p+p_n^{(j)}\phi _np_n^{(j)}^p\hfill \\ & f_n^{(j1)}\phi _nf_n^{(j1)}^p+\frac{\alpha ^{4p}}{(2^{m_{j1}+3})^{4p}}\hfill \end{array}$$
and (3.7) is verified. The construction is done.
To complete the proof of the theorem, we note from (3.7) that
$$\begin{array}{c}\hfill \underset{¯}{\mathrm{lim}}_n\mathrm{}f_n^{(j)}\phi _nf_n^{(j)}^p\underset{¯}{\mathrm{lim}}_n\mathrm{}f_n^{(1)}\phi _nf_n^{(1)}^p+\alpha ^{4p}\underset{k=1}{\overset{j1}{}}\frac{1}{(2^{m_k+3})^{4p}}\end{array}$$
So the series $`{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2^{m_k+3})^{4p}}}`$ is convergent. In particular $`lim_k\mathrm{}m_k=\mathrm{}`$.
We note from (3.4) that every $`j1`$ and $`nn_{j1}`$, $`f_n^{(j)}={\displaystyle \underset{k=1}{\overset{j}{}}}p_n^{(k)}=_{k=1}^jp_n^{(k)}`$; this is the case since $`(n_j)`$ is increasing so if $`nn_{j1}`$ then $`nn_{l1}`$ for all $`lj`$ and hence
$$\begin{array}{cc}\hfill f_n^{(j)}& =f_n^{(j1)}+p_n^{(j)}\hfill \\ & =f_n^{(j2)}+p_n^{(j1)}+p_n^{(j)}\hfill \\ & =\underset{k=1}{\overset{j}{}}p_n^{(k)}\hfill \end{array}$$
and all the $`(p_n^{(k)})_{1kj}`$ are mutually disjoint.
Now choose an increasing sequence $`(k_j)`$ so that $`k_j>\mathrm{max}(k_{j1},n_{j1})`$, $`\omega _0(f_{k_j}^{(j)})<\frac{1}{2^j}`$ and
(3.8)
$$\alpha (1\frac{1}{2^{m_j2}})\underset{k}{sup}f_{k_j}^{(j)}\phi _kf_{k_j}^{(j)}$$
(this last condition is possible from (3.5)).
Claim: $`\{f_{k_j}^{(j)};j1\}`$ is a commuting family of projections in $``$.
In fact for each $`jl`$, $`f_{k_j}^{(j)}=_{k=1}^jp_{k_j}^{(k)}`$ and $`f_{k_l}^{(l)}=_{k=1}^lp_{k_l}^{(k)}`$. For each $`1kl`$, $`p_{k_l}^{(k)}p_{k_j}^{(k)}`$ and $`p_{k_l}^{(k)}\underset{s=1;sk}{\overset{j}{}}p_{k_j}^{(s)}`$, hence
$$f_{k_j}^{(j)}f_{k_l}^{(l)}=f_{k_l}^{(l)}f_{k_j}^{(j)}=\underset{k=1}{\overset{l}{}}p_{k_j}^{(k)}$$
and the claim follows.
Set $`𝒮`$ to be a maximal abelian von Nemann subalgebra of $``$ generated by $`\left\{f_{k_j}^{(j)};j1\right\}`$. Since $`𝒮`$ is abelian, $`\omega _0`$ restricted to $`𝒮`$ is a faithful tracial state on $`𝒮`$.
Set $`p_n:=\underset{jn}{}f_{k_j}^{(j)}`$ (where the supremum is taken in $`𝒮`$). It is clear that
$$\omega _0(p_n)\underset{j=n}{\overset{\mathrm{}}{}}\omega _0(f_{k_j}^{(j)})\underset{j=n}{\overset{\mathrm{}}{}}\frac{1}{2^j}$$
so $`\omega _0(p_n)0`$ which shows that $`(p_n)_n𝒟`$. Moreover, since $`p_nf_{k_n}^{(n)}`$ for all $`n1`$, condition (3.8) implies that
$$\underset{k}{sup}p_n\phi _kp_n\alpha (1\frac{1}{2^{m_n2}}).$$
This would show that $`lim_n\mathrm{}sup_kp_n\phi _kp_n=\alpha `$.
This is a contradiction with the initial assumption that $`\alpha `$ is not attained. The proof is complete.
## 4. Main Result
The main results in this section are Theorem 4.1 and Theorem 4.4 which generalize the classical Kadec-Pełczyński subsequence decomposition to bounded sequences in the Haagerup $`L^p`$-spaces.
###### Theorem 4.1.
Let $``$ be a von Nemann algebra, $`1p<\mathrm{}`$ and $`(\phi _n)_n`$ be a bounded sequence in $`L^p()`$. Then there exist a subsequence $`(\phi _{n_k})_k`$ of $`(\phi _n)_n`$, two bounded sequences $`(y_k)`$, $`(z_k)`$ in $`L^p()`$ and a decreasing sequence of projections $`e_k_k0`$ in $``$ such that:
* $`\phi _{n_k}=y_k+z_k`$ for all $`k1`$;
* $`\{y_k,k1\}`$ is uniformly integrable in $`L^p()`$;
* $`z_k=e_kz_ke_k`$ for all $`k1`$.
###### Proof.
We will assume first that $``$ is $`\sigma `$-finite. Without loss of generality, we can and do assume that $`\phi _n1`$ for all $`n1`$ and $`\{\phi _n,n1\}`$ is not uniformly integrable. We will show that there exist a sequence $`(n_k)`$ in $``$ and $`(e_k)_k𝒟`$ such that the bounded set $`\{\phi _{n_k}e_k\phi _{n_k}e_k;k1\}`$ is uniformly integrable in $`L^p()`$.
By Theorem 3.1, there exists a subsequence of $`(\phi _n)_n`$ ( which we will denote again by $`(\phi _n)_n`$) and $`(e_n)_n𝒟`$ such that
$$\begin{array}{cc}\hfill sup& \{\underset{n\mathrm{}}{lim}\underset{k}{sup}f_n(|\phi _k|+|\phi _k^{}|)f_n;(f_n)𝒟\}\hfill \\ & =\underset{n\mathrm{}}{lim}\underset{k}{sup}e_n(|\phi _k|+|\phi _k^{}|)e_n=\alpha >0.\hfill \end{array}$$
Choose a subsequence $`(\phi _{n_k})_k(\phi _k)`$ so that
(4.1)
$$\underset{k\mathrm{}}{lim}e_k\phi _{n_k}e_k=\alpha .$$
Set $`u_k:=\phi _{n_k}`$ and $`v_k:=u_ke_ku_ke_k`$ for all $`k1`$.
Claim: The set $`V=\{v_k;k\}`$ is uniformly integrable in $`L^p()`$.
To see this claim, we will first prove the following intermediate lemma:
###### Lemma 4.2.
Let $`n_0`$, $`(1e_{n_0})V`$ and $`|V(1e_{n_0})|`$ are uniformly integrable subsets of $`L^p()`$.
We will show that $`|V(1e_{n_0})|`$ is uniformly integrable. Assume the opposite. There exists $`(f_n)_n𝒟`$ such that
$$\underset{n\mathrm{}}{lim}\underset{k}{sup}f_n|v_k(1e_{n_0})|f_n>0.$$
From this, there would exists $`(p_n)_n𝒟`$ with $`p_n1e_{n_0}`$ and such that
$$\underset{n\mathrm{}}{lim}\underset{k}{sup}p_n(|u_k|+|u_k^{}|p_n)>0.$$
In fact, for each $`k1`$, if we denote by $`\omega _k`$ the partial isometry in $``$ so that $`|v_k(1e_{n_0})|=\omega _kv_k(1e_{n_0})`$, then
$$\begin{array}{c}\hfill f_n|v_k(1e_{n_0})|f_n=f_n\omega _kv_k(1e_{n_0})f_nv_k(1e_{n_0})f_n\end{array}$$
Note that for $`kn_0`$, $`e_k(1e_{n_0})=0`$ so $`f_n|v_k(1e_{n_0})|f_nu_k(1e_{n_0})f_n`$ and by Lemma 3.2,
$$f_n|v_k(1e_{n_0})|f_nu_k^{\frac{3}{4}}(1e_{n_0})f_n(1e_{n_0})|u_k|(1e_{n_0})f_n(1e_{n_0})^{\frac{1}{4}}$$
which shows that
$$f_n|v_k(1e_{n_0})|f_n(1e_{n_0})f_n(1e_{n_0})|u_k|(1e_{n_0})f_n(1e_{n_0})^{\frac{1}{4}}.$$
Let $`a_n=(1e_{n_0})f_n(1e_{n_0})`$. It is clear that $`a_n_n0`$ and using Proposition 3.4, we conclude that there exists $`(p_n)_n𝒟`$ , $`p_n1e_{n_0}`$ such that $`lim_n\mathrm{}sup_kp_n|u_k|p_n>0`$. In particular:
$$\underset{n\mathrm{}}{lim}\underset{k}{sup}p_n(|u_k|+|u_k^{}|)p_n>0.$$
Now choose a subsequence $`(k_j)`$ so that there exists $`\gamma >0`$ satisfying
(4.2)
$$\underset{j\mathrm{}}{lim}p_j(|u_{k_j}|+|u_{k_j}^{}|)p_j=\gamma >0.$$
Since $`p_j1e_{n_0}`$ for all $`j`$, $`e_{k_j}p_j`$ for $`k_j>n_0`$ and therefore
$$\begin{array}{c}\hfill (e_{k_j}+p_j)\left(|u_{k_j}|+|u_{k_j}^{}|\right)(e_{k_j}+p_j)^pe_{k_j}\left(|u_{k_j}|+|u_{k_j}^{}|\right)e_{k_j}^p+p_j\left(|u_{k_j}|+|u_{k_j}^{}|\right)p_j^p\end{array}$$
and taking the limit as $`j\mathrm{}`$, (4.1) and (4.2) imply $`\alpha ^p\gamma ^p+\alpha ^p`$. This is a contradiction since $`\gamma >0`$. The proof of the lemma is complete.
To complete the proof of the theorem, assume that $`V`$ is not uniformly integrable. Using Lemma 2.2 and Lemma 4.2,
$$\underset{n\mathrm{}}{lim}\underset{k}{sup}e_nv_ke_n>0.$$
Again, choose a subsequence $`(k_n)`$ so that
(4.3)
$$\underset{n\mathrm{}}{lim}e_nv_{k_n}e_n>0.$$
Claim: $`e_nv_{k_n}e_n^24(e_ne_{k_n})\left(|u_{k_n}|+|u_{k_n}^{}|\right)(e_ne_{k_n})`$.
To see this claim, we note that since $`e_ne_{k_n}`$, $`e_nv_{k_n}e_n=e_nu_{k_n}e_ne_{k_n}u_{k_n}e_{k_n}`$ so $`e_nv_{k_n}e_n=(e_ne_{k_n})u_{k_n}e_n+e_{k_n}u_{k_n}(e_ne_{k_n})`$ and therefore
$$\begin{array}{cc}\hfill e_nv_{k_n}e_n& (e_ne_{k_n})u_{k_n}+u_{k_n}(e_ne_{k_n})\hfill \\ & u_{k_n}^{}(e_ne_{k_n})+u_{k_n}(e_ne_{k_n})\hfill \\ & u_{k_n}^{}^{\frac{1}{2}}(e_ne_{k_n})|u_{k_n}^{}|(e_ne_{k_n})^{\frac{1}{2}}+u_{k_n}^{\frac{1}{2}}(e_ne_{k_n})|u_{k_n}|(e_ne_{k_n})^{\frac{1}{2}}\hfill \end{array}$$
and since $`u_{k_n}1`$,
$$\begin{array}{cc}\hfill e_nv_{k_n}e_n& (e_ne_{k_n})|u_{k_n}^{}|(e_ne_{k_n})^{\frac{1}{2}}+(e_ne_{k_n})|u_{k_n}|(e_ne_{k_n})^{\frac{1}{2}}\hfill \\ & 2(e_ne_{k_n})\left(|u_{k_n}|+|u_{k_n}^{}|\right)(e_ne_{k_n})^{\frac{1}{2}}\hfill \end{array}$$
and the claim follows.
From the claim above and equation (4.3), there exists $`\nu >0`$ such that
(4.4)
$$\underset{¯}{\mathrm{lim}}_n\mathrm{}(e_ne_{k_n})\left(|u_{k_n}|+|u_{k_n}^{}|\right)(e_ne_{k_n})=\nu >0.$$
Observe that since $`e_{k_n}(e_ne_{k_n})`$,
$$\begin{array}{c}\hfill e_n\left(|u_{k_n}|+|u_{k_n}^{}|\right)e_n^pe_{k_n}\left(|u_{k_n}|+|u_{k_n}^{}|\right)e_{k_n}^p+(e_ne_{k_n})\left(|u_{k_n}|+|u_{k_n}^{}|\right)(e_ne_{k_n})^p.\end{array}$$
Taking the limit ( as $`n\mathrm{}`$) together with (4.1) and (4.4) would imply $`\alpha ^p\nu ^p+\alpha ^p`$. This is a contradiction since $`\nu >0`$.
By setting $`y_k:=v_k`$ and $`z_k:=e_ku_ke_k`$, the proof for the $`\sigma `$-finite case is complete.
For the general case, let $``$ be a von Neumann algebra (not necessarily $`\sigma `$-finite) and $`(\phi _n)_n`$ in $`L^p()`$ as in the theorem. Fix an orthogonal family of cyclic projections $`(e_\alpha )_{\alpha I}`$ in $``$ such that $`1=_{\alpha I}e_\alpha `$ ( see for instance, \[20, Proposition 5.5.9, p. 336\] ) .
###### Lemma 4.3.
There exists a countably decomposable projection $`e`$ such that for all $`n1`$, $`e\phi _n=\phi _ne=\phi _n`$.
For each $`n`$ and $`\epsilon >0`$, set $`E_{n,\epsilon }:=\{\alpha I;e_\alpha \phi _n>\epsilon \}`$ and $`E_n:=\{\alpha I;e_\alpha \phi _n0\}`$.
Claim: $`E_{n,\epsilon }`$ is finite (hence $`E_n`$ is countable).
To see this, assume that $`E_{n,\epsilon }`$ is infinite. Then there exists an infinite sequence $`(e_k)_k(e_\alpha )_{\alpha I}`$ such that $`e_k\phi _n>\epsilon `$ for all $`k`$. If $`J`$ is a finite subset of $``$, then
$$\begin{array}{cc}\hfill \underset{kJ}{}e_k\phi _n& =(\underset{kJ}{}e_k)\phi _n\hfill \\ & =(_{kJ}e_k)\phi _n\phi _n.\hfill \end{array}$$
So $`_{kJ}e_k\phi _n\phi _n`$ (a constant independant of $`J`$) which shows that $`_{k=1}^{\mathrm{}}e_k\phi _n`$ is a weakly unconditionally Cauchy (w.u.c.) series in $`L^p()`$ but since $`L^p()`$ does not contain any copies of $`c_0`$, $`_{k=1}^{\mathrm{}}e_k\phi _n`$ is unconditionally convergent and hence $`lim_k\mathrm{}e_k\phi _n=0`$ (see for instance p.45). This is in contradiction with the assumption $`e_k\phi _n\epsilon `$ for all $`k`$. We proved that $`E_{n,\epsilon }`$ is finite. It is clear that $`E_n=_kE_{n,\frac{1}{k}}`$ so it is at most countable. The claim is verified.
Similarly, if $`R_n=\{\alpha I,\phi _ne_\alpha 0\}`$ then $`R_n`$ is at most countable.
Let $`C=\underset{n=1}{\overset{\mathrm{}}{}}(R_nE_n)`$ ; $`C`$ is at most countable and if $`e=_{\alpha C}e_\alpha `$ then $`e`$ is the union of a countable family of disjoint cyclic projections in $``$ so $`e`$ is countably decomposable in $``$ (\[20, Proposition 5.5.19 p.340 \]). The construction of $`e`$ implies that $`e\phi _n=\phi _ne=\phi _n`$ for all $`n1`$. The lemma is proved.
To conclude the proof of the theorem, consider the von Neumann algebra $`ee`$. Since $`e`$ is countably decomposable, $`ee`$ is $`\sigma `$-finite. Let $`T:ee`$ be the map that takes $`x`$ to $`exe`$. The map $`T`$ is bounded and is weak to weak continuous so there exists a map $`S:(ee)_{}_{}`$ so that $`S^{}=T`$.
Let $`R:_{}(ee)_{}`$ be the restriction map. The operators $`T`$ and $`R`$ can be interpolated and since $`L^p()`$ (resp. $`L^p(ee)`$) is isometrically isomorphic to $`(,_{})_\theta `$ (resp. $`(ee,(ee)_{})_\theta )`$ for $`\theta =\frac{1}{p}`$, (see ), we get a bounded linear map $`T_p:L^p()L^p(ee)`$. Similarly, if one considers the inclusion map $`ee`$ and $`S:(ee)_{}_{}`$ as above, then by interpolation, we obtain a map $`S_p:L^p(ee)L^p()`$.
Apply the $`\sigma `$-finite case to the sequence $`(T_p(\phi _n))_{n1}`$ in $`L^p(ee)`$ to get a decomposition
$$T_p(\phi _{n_k})=y_k+z_kk1$$
with $`(y_k)_k`$ and $`(z_k)_k`$ satisfying the conclusion of the theorem. It is enough to consider the decomposition:
$$\phi _{n_k}=S_p(y_k)+S_p(z_k)k1.$$
The proof is complete. ∎
The theorem which follows shows that, as in the semi-finite case, the decreasing projections in Theorem 4.1 can be replaced by mutually orthogonal projections. Its proof is identical to that of the semi-finite case (, Theorem 3.7).
###### Theorem 4.4.
Let $``$ be a von Neumann algebra and $`1p<\mathrm{}`$, Let $`(\phi _n)_n`$ be a bounded sequence in $`L^p()`$ then there exists a subsequence $`(\phi _{n_k})`$ of $`\phi _n`$, bounded sequences $`(y_k)`$ and $`(\zeta _k)_k`$ in $`L^p()`$ and mutually orthogonal sequence of projections $`(e_k)_k`$ in $``$ such that:
* $`\phi _{n_k}=y_k+\zeta _k`$ for all $`k1`$;
* $`\{y_k:k1\}`$ is uniformly integrable and $`e_ky_ke_k=0`$ for all $`k1`$;
* $`(\zeta _k)_k`$ is such that $`e_k\zeta _ke_k=\zeta _k`$ for all $`k1`$.
###### Remark 4.5.
For $`1<p<\mathrm{}`$, it should be noted that since $`L^p()`$ is is a closed subspace of $`L^{p,\mathrm{}}(𝒩,\tau )`$ and $`L^{p,\mathrm{}}(^+,m)`$ has the Fatou property, one could apply the semi-finite case of the the Kadec-Pełczyński subsequence decomposition to any bounded sequence of $`L^p()`$ (viewed as bounded sequence in $`L^{p,\mathrm{}}(𝒩,\tau ))`$. However, that procedure would provide decreasing projections in $`𝒩`$ and as is noted in \[26, Remarks 3.5 (iii)\], these projections are either of finite trace or their orthogonal complements are of finite trace which guaranties that projections obtained from applying the semifinite case cannot be in $``$.
## 5. Applications
A result of Maurey (, see also ) states that every reflexive subspace of $`L^1[0,1]`$ has the fixed point property for nonexpansive mappings (FPP). Later, Dowling and Lennard showed that the converse of Maurey’s result is valid: every non-reflexive subspace $`L^1[0,1]`$ fails the FPP (). This section is for the study of generalizations to the case of duals of $`C^{}`$-algebras and requires the notion of asymptotically isometric copies of $`\mathrm{}^1`$ which was introduced by Dowling and Lennard in .
###### Definition 5.1.
A Banach space $`X`$ is said to contain asymptotically isometric copies of $`\mathrm{}^1`$ if for every null sequence $`(\epsilon _n)`$ of positive numbers, there exists a sequence $`(x_n)`$ in $`X`$ such that:
$$\underset{n=1}{\overset{\mathrm{}}{}}(1\epsilon _n)|a_n|\underset{n=1}{\overset{\mathrm{}}{}}a_nx_n\underset{n=1}{\overset{\mathrm{}}{}}|a_n|.$$
for all $`(a_n)\mathrm{}^1`$.
The following result is a generalization of .
###### Theorem 5.2.
let $`𝒜`$ be a $`C^{}`$-algebra. Every non-reflexive subspace of $`𝒜^{}`$ contains asymptotically isometric copies of $`\mathrm{}^1`$.
###### Proof.
Note that $`𝒜^{}`$ is a von Neumann algebra so subspaces of $`𝒜^{}`$ are subspaces of preduals of von Neumann algebras. The proof then follows the argument used in using Theorem 4.4. Details are left to the readers. ∎
###### Remark 5.3.
In , Bélanger proved an improved version of the Akemann’s characterization of weak compactness on preduals of von Neumann algebras. He then went on to show that non-reflexive preduals of von Neumann algebras contain complemented copies of $`\mathrm{}^1`$. This fact can also be deduced from a result of Pfitzner which states that $`C^{}`$-algebras have Pełczyński property (V) so their duals have property (V\*) (). It is plain from Theorem 4.4 that the asymptotically isometric copies of $`\mathrm{}^1`$ in Theorem 5.2 are complemented with good projection constants.
For the next extension, we recall that $`JB^{}`$-triples are all those Banach spaces whose open unit balls are bounded symmetric domains . Examples of $`JB^{}`$-triples are $`C^{}`$-algebras and Hilbert spaces. Other important examples are the so-called Cartan factors $`C^k(k=1,2,\mathrm{},6)`$ where the rectangular Cartan factor $`C^1=(H,K)`$ consists of bounded operators between Hilbert spaces, the symplectic factor $`C_n^2`$ is $`\{z(H);z=jz^{}j\}`$ where $`j:HH`$ is a conjugate linear isometric involution, the Hermitian Cartan factor $`C^3`$ is $`\{z(H);z=jz^{}j\}`$, $`C^4`$ is the spin factor, $`C^5`$ is the (finite dimentional) exceptional Cartan factor consisting of $`1\times 2`$ matrices over the complex Caley numbers $`𝐎`$ and $`C^6`$ is the set of all $`3\times 3`$ Hermitian matrices over $`𝐎`$. Dual $`JB^{}`$-triples are called $`JBW^{}`$-triples. For more informations, we refer to , and .
###### Corollary 5.4.
If $`𝒥`$ is a $`JB^{}`$-triple then every non-reflexive subspace of $`𝒥^{}`$ contains asymptotically isometric copies of $`\mathrm{}^1`$.
For the poof we will need two lemmas on stability of asymptotically isometric copies of $`\mathrm{}^1`$.
###### Lemma 5.5.
Let $`E_1`$ and $`E_2`$ be weakly sequentially complete Banach spaces so that any sequence equivalent to the unit vector basis of $`\mathrm{}^1`$ in $`E_j`$ $`(j=1,2)`$ has a normalized block that is asymptotically isometric to $`\mathrm{}^1`$ then every sequence equivalent to the unit vector basis of $`\mathrm{}^1`$ in $`E_1_1E_2`$ has a normalized block that is asymptotically isometric to $`\mathrm{}^1`$.
###### Proof.
Let $`\{U_n=(x_n,y_n)\}_{n=1}^{\mathrm{}}`$ be a sequence in $`E_1_1E_2`$ that is equivalent to $`\mathrm{}^1`$. After taking subsequences, either $`(x_n)_n`$ or $`(y_n)_n`$ is equivalent to $`\mathrm{}^1`$. Let assume that $`(x_n)_n`$ is equivalent to $`\mathrm{}^1`$. We have two cases.
Case 1: The sequence $`(y_n)_n`$ is weakly convergent. By taking normalized blocks, we can assume that $`(x_n)_n`$ is asymptotically isometric to $`\mathrm{}^1`$ and $`lim_n\mathrm{}y_n=0`$. There exists a null sequence $`(\epsilon _n)`$ of positive numbers such that:
$$\underset{n=1}{\overset{\mathrm{}}{}}(1\epsilon _n)|a_n|\underset{n=1}{\overset{\mathrm{}}{}}a_nx_n\underset{n=1}{\overset{\mathrm{}}{}}|a_n|.$$
for all $`(a_n)\mathrm{}^1`$ but since $`_{n=1}^{\mathrm{}}a_nU_n=_{n=1}^{\mathrm{}}a_nx_n+_{n=1}^{\mathrm{}}a_ny_n`$, we get that
$$\underset{n=1}{\overset{\mathrm{}}{}}(1\epsilon _n)|a_n|\underset{n=1}{\overset{\mathrm{}}{}}a_nU_n\underset{n=1}{\overset{\mathrm{}}{}}(1+y_n)|a_n|.$$
This concludes that $`(U_n)_n`$ is asymptotically isometric to $`\mathrm{}^1`$.
Case 2: The sequence $`(y_n)`$ is equivalent to $`\mathrm{}^1`$. As above, one can find a block so that both the coresponding block for $`(x_n)_n`$ and $`(y_n)_n`$ are asymptotically isometric to $`\mathrm{}^1`$. Set $`Z_n:=U_n/2=(x_n/2,y_n/2)`$. It can be easily seen that $`(Z_n)_n`$ is equivalent to an asympotically isometric copy of $`\mathrm{}^1`$ in $`E_1_1E_2`$. ∎
###### Lemma 5.6.
Let $`(\mathrm{\Omega },\mathrm{\Sigma },\lambda )`$ be a measure space and $`R`$ be a reflexive Banach space. Every sequence equivalent to the unit vector basis of $`\mathrm{}^1`$ in $`L^1(\lambda ,R)`$ has a normalized block that is asymptotically isometric to $`\mathrm{}^1`$.
###### Proof.
Let $`(f_n)_n`$ be a sequence equivalent to the $`\mathrm{}^1`$ basis. Since $`R`$ is reflexive, the sequence $`(f_n)_n`$ can not be uniformly integrable (see for instance ). Apply the classical Kadec-Pełczyński subsequence decomposition to the sequence $`(f_n())_n`$ in $`L^1(\lambda )`$ to get a pairwise disjoint sequence of measurable sets $`(A_n)_n`$ such that $`\left\{f_n()\chi _{\mathrm{\Omega }A_n},n1\right\}`$ is uniformly integrable. The space $`R`$ being reflexive implies that $`\left\{f_n\chi _{\mathrm{\Omega }A_n},n1\right\}`$ is relatively weakly compact in $`L^1(\lambda ,R)`$. We conclude the proof as in the scalar case. ∎
Proof of Corollary 5.4: Let $`𝒥`$ be a $`JB^{}`$-triple and $`X`$ be a non-reflexive subspace of $`^{}`$. Since $`𝒥^{}`$ is a $`JBW^{}`$-triple, we can assume that $`X`$ is a subspace of the predual of a $`JBW^{}`$-triple $``$. By and , $``$ admits the following form:
$$=\left(\underset{\alpha }{}C(\mathrm{\Omega }_\alpha ,C^\alpha )\right)_{\mathrm{}^{\mathrm{}}}_{\mathrm{}}J^7_{\mathrm{}}J^8,$$
where $`C(\mathrm{\Omega }_\alpha ,C^\alpha )`$ is the space of continuous functions from a hyperstonean space $`\mathrm{\Omega }_\alpha `$ to a Cartan factor $`C^\alpha `$, $`J^7=\{aM;\mathrm{\Theta }(a)=a\}`$ with $`\mathrm{\Theta }:MM`$ is a $`w^{}`$-continuous \- antiautomorphism of period $`2`$ on a von Neumann algebra $`M`$ and $`J^8`$ is a $`w^{}`$-closed right ideal of a von Neumann algebra $`N`$. The predual of $``$ is equal to the $`\mathrm{}^1`$-sum
$$_{}=\left(\underset{\alpha }{}L^1(\mathrm{\Sigma }_\alpha ,C_{}^\alpha )\right)_\mathrm{}^1_1J_{}^7_1J_{}^8.$$
By \[4, Theorem 2\], the space $`E_1=\left(_{\alpha 5,6}L^1(\mathrm{\Sigma }_\alpha ,C_{}^\alpha )\right)_\mathrm{}^1_1J_{}^7_1J_{}^8`$ is isometric to a 1-complemented subspace of the predual of a von Neumann algebra so $`E_1`$ is isometric to a subspace of the predual of such von Neumann algebra and hence satisfies the assumption of Lemma 5.3. Moreover, since $`C^5`$ and $`C^6`$ are finite dimensional, the space $`E_2=L^1(\mathrm{\Sigma }_5,C^5)_1L^1(\mathrm{\Sigma }_6,C^6)`$ satisfies (as does $`L^1`$-spaces) the assumption of Lemma 5.3. We conclude that every sequence equivalent to the unit vector basis of $`\mathrm{}^1`$ in $`_{}=E_1_1E_2`$ has a normalized block that is asymptotically isometric to $`\mathrm{}^1`$. The proof is complete. ∎
###### Corollary 5.7.
If $`𝒥`$ is a $`JB^{}`$-triple then every non-reflexive subspace of $`𝒥^{}`$ fails the fixed point property for nonexpansive self-maps on closed bounded convex sets.
Acknowledgements. This project started when the author was participating in the NSF-supported Workshop on Linear Analysis and Probability (Summer 1998) at the Department of Mathematics of the Texas A& M University. The author would like to express his gratitute to Professor W. Johnson for the invitation and warm hospitality. The author is also indebted to Professor P. Dowling for introducing him to the topic of this paper and for several useful discussions.
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# 1 Introduction
## 1 Introduction
Conformal covariance is essential for the global formulation of scale invariant theories (conformal models) on compact Riemann surfaces of any genus. It is at the heart of $`W`$-algebras which are non-linear generalizations of the two-dimensional conformal algebra, i.e. of the Virasoro algebra. Moreover, as was realized in the eighties and nineties, conformal symmetry manifests itself in several respects in two-dimensional integrable models like the KdV or Boussinesq equations.
By taking into account the underlying symmetries of a given theory, one usually gains a better understanding of this theory . From a practical point of view, these symmetries generally provide a useful tool for determining solutions or for checking results within a given theory.
Within the aforementioned theories and models, conformal symmetry manifests itself by the occurrence of conformally covariant differential operators in the time evolution equations or in the structure relations. In the present notes, we briefly review the definition and construction of these operators and of their supersymmetric extensions.
In our write-up, we have tried to maintain the informal style of the oral presentation and therefore some results are only illustrated by the simplest examples. For more details, we refer to the series of articles - and to the work cited therein. (Among the latter, we explicitly mention references which represent the basis for some parts of -.) In reference , we illustrate how conformally covariant operators enter the physical models we mentioned and we show how they constrain, or largely determine, the form of some of these theories.
## 2 Geometric framework
### 2.1 Basic definitions
The arena we will work on, is a Riemann surface $`𝚺`$, i.e. a connected, topological $`2`$-manifold which is equipped with a complex structure (or equivalently, a real, smooth, connected and oriented $`2`$-manifold which is equipped with a conformal class of metrics) . Roughly speaking, this means that any two systems of local complex coordinates, say $`z`$ and $`z^{}`$, are related by a conformal coordinate transformation,
$$z\stackrel{conf.}{}z^{}(z).$$
In the following, we will use the notation $`\frac{}{z}`$ and we will denote the complex conjugate of $`z`$ by $`\overline{z}`$. Moreover, we assume that the considered Riemann surfaces are compact so that they are characterized by their genus $`g0`$.
A conformal (or primary) field of weight $`k𝐙/2`$ on the Riemann surface $`𝚺`$ is a collection $`\{c(z,\overline{z})\}`$ of local complex-valued functions on $`𝚺`$ (one for each coordinate system $`(z,\overline{z})`$), transforming according to
$$c^{}(z^{},\overline{z}^{})=(z^{})^kc(z,\overline{z})$$
(1)
under a conformal change of coordinates. Thus, $`c`$ transforms linearly with a certain power of the Jacobian of the change of coordinates<sup>3</sup><sup>3</sup>3One can consider conformal fields which also transform with a certain power of $`\overline{}\overline{z}^{}`$, but we will not need them in the sequel.. The space of conformal fields of weight $`k`$ on $`𝚺`$ will be denoted by $`_k`$.
The Schwarzian derivative of a conformal change of coordinates $`zz^{}(z)`$ is defined by
$$S(z^{};z)=^2\mathrm{ln}z^{}\frac{1}{2}\left(\mathrm{ln}z^{}\right)^2.$$
(2)
A projective (or Schwarzian) connection on the Riemann surface $`𝚺`$ is a collection $`\{R(z,\overline{z})\}`$ of local complex-valued functions on $`𝚺`$ with the properties
(i) $`R`$ is locally holomorphic, i.e. $`_{\overline{z}}R=0`$,
(ii) $`R`$ transforms inhomogeneously with the Schwarzian derivative under a conformal change of coordinates $`zz^{}(z)`$ :
$$R^{}(z^{})=\left(z^{}\right)^2\left[R(z)S(z^{};z)\right].$$
(3)
Such connections exist globally on compact Riemann surfaces of any genus. From the physical point of view, the field $`R`$ and its complex conjugate represent the components of the energy-momentum tensor in two-dimensional conformal field theory.
### 2.2 Projective coordinates
A change of local coordinates $`ZZ^{}(Z)`$ which has the form
$$Z^{}=\frac{aZ+b}{cZ+d}\mathrm{with}a,b,c,d𝐂\mathrm{and}adbc=1,$$
(4)
is called a projective (or Möbius or fractional linear) transformation. We note that the associated Jacobian is given by
$$_ZZ^{}=(cZ+d)^2.$$
(5)
In the following, coordinates belonging to a projective atlas on the Riemann surface $`𝚺`$ will always be denoted by capital letters $`Z`$ or $`Z^{}`$.
A projective structure on $`𝚺`$ is an atlas of local coordinates for which all coordinate transformations are projective. Every Riemann surface admits such a structure. As a matter of fact, there is a one-to-one correspondence between projective structures and projective connections , see section 3.3 below.
Let $`𝚺`$ be a compact Riemann surface with a given projective structure. Then, a quasi-primary field of weight $`k𝐙/2`$ on $`𝚺`$ is a collection $`\{C_k(Z,\overline{Z})\}`$ of local complex-valued functions on $`𝚺`$ which transform linearly with the $`k`$-th power of the Jacobian (5) under a projective change of coordinates:
$$C_k^{}=(cZ+d)^{2k}C_k.$$
(6)
### 2.3 Covariant linear differential operators
Consider the local form of a linear, holomorphic differential operator of order $`n𝐍`$, which is defined on the Riemann surface $`𝚺`$:
$$L^{(n)}=a_0^{(n)}^n+a_1^{(n)}^{n1}+a_2^{(n)}^{n2}+\mathrm{}+a_n^{(n)}\mathrm{with}a_k^{(n)}=a_k^{(n)}(z).$$
If the leading coefficient $`a_0^{(n)}`$ does not vanish anywhere, we can divide it. Therefore, in the following, we will assume that $`a_0^{(n)}1`$.
###### Definition 2.1
A holomorphic, $`n`$-th order differential operator, which is locally given on the compact Riemann surface $`𝚺`$ by
$$L^{(n)}=^n+a_1^{(n)}^{n1}+a_2^{(n)}^{n2}+\mathrm{}+a_n^{(n)},$$
is called conformally covariant if it maps conformal fields (of some weight $`p𝐙/2`$) to conformal fields:
$$L^{(n)}:_p_{p+n}.$$
This requirement is equivalent to the one that $`L^{(n)}`$ transforms according to the following operatorial relation under a conformal change of coordinates $`zz^{}(z)`$:
$$L^{(n)}=(z^{})^{(p+n)}L^{(n)}(z^{})^p.$$
(7)
According to the following result, the coefficient $`a_1^{(n)}`$ of a conformally covariant operator can always be eliminated without destroying conformal covariance .
###### Theorem 2.1
Consider $`n𝐍^{}`$. On a compact Riemann surface of genus $`g>1`$, a conformally covariant operator $`L^{(n)}`$ for which the coefficient $`a_1^{(n)}`$ does not identically vanish, can only exist if it acts on conformal fields of weight
$$p=\frac{1n}{2}.$$
In this case, $`a_1^{(n)}`$ transforms linearly under a conformal change of coordinates $`zz^{}(z)`$,
$$a_1^{(n)}=(z^{})^1a_1^{(n)},$$
(8)
and thereby one can consistently impose the vanishing of this coefficient. The transformation law of $`a_2^{(n)}`$ then takes the simple form
$$a_2^{(n)}=(z^{})^2\left[a_2^{(n)}k_nS(z^{};z)\right]wherek_n=\frac{n(n^21)}{12},$$
(9)
and where $`S`$ denotes the Schwarzian derivative.
Accordingly, in the sequel, we will always consider conformally covariant operators which are normalized by $`a_0^{(n)}1,a_1^{(n)}0`$ and, for short, we will refer to these as CCO’s:
###### Definition 2.2
A CCO (conformally covariant operator) of order $`n`$ on the compact Riemann surface $`𝚺`$ is a map
$$L^{(n)}:_{\frac{1n}{2}}_{\frac{1+n}{2}}$$
(10)
with the local expression
$$L^{(n)}=^n+a_2^{(n)}^{n2}+\mathrm{}+a_n^{(n)}.$$
(11)
Here, the coefficients $`a_2^{(n)},\mathrm{},a_n^{(n)}`$ are locally holomorphic functions on $`𝚺`$ and $`a_2^{(n)}`$ is a multiple of a projective connection:
$$a_2^{(n)}=\frac{n(n^21)}{12}R.$$
The remaining coefficients $`a_3^{(n)},\mathrm{},a_n^{(n)}`$ transform in a more complicated way than $`R`$ under conformal changes of coordinates , so as to ensure the covariance (10).
## 3 CCO’s
From the conceptual point of view, CCO’s are best approached by starting from the special coordinate system where $`a_2^{(n)}=0`$ (i.e. by starting from projective coordinates $`Z`$) and then going over to generic local coordinates $`z`$ by a conformal transformation: the dependence of the operators on the projective structure then translates into a dependence on a projective connection. Therefore, we will first discuss operators on $`𝚺`$ which are covariant with respect to projective transformations.
### 3.1 Möbius covariant operators
Class 1: Operators which only depend on the projective structure
The operator $`_Z^n(\frac{}{Z})^n`$ (where $`Z`$ belongs to a projective atlas on $`𝚺`$) transforms homogeneously if it acts on a quasi-primary field of weight $`\frac{1n}{2}`$ :
###### Lemma 3.1 (Bol’s lemma)
Consider a projective atlas on $`𝚺`$ with local changes of coordinates (4). If $`C_{\frac{1n}{2}}(Z,\overline{Z})`$ is a quasi-primary field on $`𝚺`$, then $`_Z^nC_{\frac{1n}{2}}`$ also is, i.e. it transforms according to
$$\left(_Z^nC_{\frac{1n}{2}}\right)^{}=\left(cZ+d\right)^{1+n}_Z^nC_{\frac{1n}{2}}(n=\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}2},\mathrm{}).$$
(12)
Class 2: Operators which depend linearly on a quasi-primary field
For a given $`n𝐍`$ with $`n3`$, we consider linear Möbius covariant operators $`M_{W_3}^{(n)},\mathrm{},M_{W_n}^{(n)}`$ acting on quasi-primary fields of weight $`\frac{1n}{2}`$. These operators do not only depend on the projective structure, but also, in a linear way, on quasi-primary fields $`W_3,\mathrm{},W_n`$, respectively. Moreover, they are differential operators of lower order than $`_Z^n`$.
Rather than giving a general formula for all of these operators (e.g. see ), we present their explicit expression for $`n=5`$:
$`M_{W_5}^{(5)}`$ $`=`$ $`W_5,M_{W_4}^{(5)}=W_4_Z+{\displaystyle \frac{1}{2}}(_ZW_4)`$ (13)
$`M_{W_3}^{(5)}`$ $`=`$ $`W_3_Z^2+(_ZW_3)_Z+{\displaystyle \frac{2}{7}}(_Z^2W_3).`$
### 3.2 From projective to generic coordinates
Let us now go over from the projective coordinates $`Z`$ to generic holomorphic coordinates $`z`$ by a conformal transformation,
$$Z\stackrel{conf.}{}z.$$
In doing so, a quasi-primary field $`C_k`$ becomes a primary field $`c_k`$, both fields being related by
$$C_k(Z,\overline{Z})=(Z)^kc_k(z,\overline{z}).$$
(14)
Moreover, a Möbius covariant operator becomes a CCO. To discuss this passage, we consider in turn the two classes of examples introduced above.
### 3.3 Class 1: Bol operators
When passing from the projective coordinates $`Z`$ to the holomorphic coordinates $`z`$ by a conformal transformation, the $`n`$-th order derivative $`_Z^n`$ acting on a quasi-primary field $`C_{\frac{1n}{2}}`$ becomes the $`n`$-th order Bol operator denoted by $`L_n`$:
$$_Z^nC_{\frac{1n}{2}}=(Z)^{\frac{1+n}{2}}L_nc_{\frac{1n}{2}}.$$
(15)
By substituting the relation (14) with $`k=\frac{1n}{2}`$ into equation (15), we obtain the following operatorial expression for the CCO $`L_n`$:
$$L_n=(Z)^{\frac{1+n}{2}}\left(\frac{1}{Z}\right)^n(Z)^{\frac{1n}{2}}.$$
(16)
Thus, the Bol operator $`L_n`$ represents the conformally covariant version of the differential operator $`^n`$, the simplest examples being given by
$`L_0`$ $`=`$ $`1`$
$`L_1`$ $`=`$ $``$
$`L_2`$ $`=`$ $`^2+{\displaystyle \frac{1}{2}}R`$
$`L_3`$ $`=`$ $`^3+2R+(R)`$ (17)
$`L_4`$ $`=`$ $`^4+5R^2+5(R)+{\displaystyle \frac{3}{2}}\left[(^2R)+{\displaystyle \frac{3}{2}}R^2\right],`$
where
$$R_{zz}(z)S(Z;z).$$
(18)
This expression represents a projective connection, because it has the correct transformation properties thanks to the chain rule for the Schwarzian derivative. From this chain rule, it also follows that the definition (18) is not affected by a projective transformation of $`Z`$. Note that the quantity (18) is holomorphic since the change of coordinates $`zZ(z)`$ has this property. Equation (18) expresses the one-to-one correspondence between projective structures and projective connections that we already mentioned.
The basic operator $`L_2`$ (which is known as Hill operator) appears for instance in the Lax representation of the KdV equation while $`L_3`$ appears in the Poisson brackets for the Virasoro algebra or in the conformal Ward identity .
### 3.4 Class 2: Operators depending linearly on conformal fields
Upon passage $`Zz`$, the quasi-primary field $`W_k`$ becomes a primary field $`w_k`$, both fields being related by $`W_k=(Z)^kw_k`$. Moreover, the Möbius covariant operator $`M_{W_k}^{(n)}`$ becomes a CCO $`M_{w_k}^{(n)}`$ which depends linearly on $`w_k`$ and which acts on $`_{\frac{1n}{2}}`$. For instance, the $`n=5`$ operators (13) become
$`M_{w_5}^{(5)}`$ $`=`$ $`w_5,M_{w_4}^{(5)}=w_4+{\displaystyle \frac{1}{2}}(w_4)`$ (19)
$`M_{w_3}^{(5)}`$ $`=`$ $`w_3\left[^2+2R\right]+(w_3)+{\displaystyle \frac{2}{7}}\left[^23R\right]w_3.`$
### 3.5 Complete classification
Any CCO
$$L^{(n)}=^n+a_2^{(n)}^{n2}+\mathrm{}+a_n^{(n)}\mathrm{with}a_2^{(n)}=\frac{n(n^21)}{12}R$$
can be reparametrized in the following way in terms of the projective connection $`R`$ and $`n2`$ conformal fields $`w_3,\mathrm{},w_n`$:
$$L^{(n)}=L_n+M_{w_3}^{(n)}+\mathrm{}+M_{w_n}^{(n)}.$$
(20)
The relation between the coefficients $`a_3^{(n)},\mathrm{},a_n^{(n)}`$ and the conformal fields $`w_3,\mathrm{},w_n`$ is given by differential polynomials which involve $`R`$ and this relation is invertible.
The parametrization (20) of $`L^{(n)}`$ in terms of the energy-momentum tensor and some conformal fields is very helpful for the construction and formulation of $`W_n`$-algebras, see section 3.7 below.
### 3.6 Nonlinear conformally covariant operators
There exists a unique bilinear conformally covariant operator $`J(,)`$, the so-called Gordan transvectant . Here, we only note that it encompasses the CCO’s $`M_{w_k}^{(n)}`$:
$$M_{w_k}^{(n)}cJ(w_k,c).$$
(21)
The bilinear operator $`J(,)`$ as well as higher multilinear conformally covariant operators appear in the defining relations of $`W`$-algebras .
### 3.7 Matrix representation of CCO’s
The CCO’s $`L_n`$ and $`M_{w_k}^{(n)}`$ admit a matrix representation which is related to the principal embedding of the Lie algebra $`sl(2)`$ into $`sl(n)`$ . Since $`sl(2)`$ is the Lie algebra of the Möbius group, this algebraic relationship which underlies the matrix representation of CCO’s, reflects the fact that these covariant operators come from Möbius covariant ones. We will now illustrate the matrix representation for $`L^{(3)}=L_3+M_{w_3}^{(3)}=L_3+w_3`$.
Let us rewrite the scalar, conformally covariant differential equation
$$0=L^{(3)}f_3\left[^3+2R+(R)+w_3\right]f_3\mathrm{with}f_3_{\frac{1}{2}}$$
(22)
as a system of three first-order differential equations:
$$\left[\begin{array}{c}0\\ 0\\ 0\end{array}\right]=\left[\begin{array}{ccccc}& R& w_3& & \\ 1& & R& & \\ 0& 1& & & \end{array}\right]\left[\begin{array}{c}f_1\\ f_2\\ f_3\end{array}\right]\{\begin{array}{c}0=f_1+Rf_2+w_3f_3\hfill \\ f_1=f_2+Rf_3\hfill \\ f_2=f_3.\hfill \end{array}$$
(23)
Substitution of the last two equations into the first one reproduces the scalar equation (22).
Equation (23) can also be written in the form
$$\stackrel{}{0}=(𝒜)\stackrel{}{F}\mathrm{with}𝒜=\left[\begin{array}{ccccc}0& R& w_3& & \\ 1& 0& R& & \\ 0& 1& 0& & \end{array}\right],\stackrel{}{F}=\left[\begin{array}{c}f_1\\ f_2\\ f_3\end{array}\right].$$
(24)
Here, the matrix $`𝒜`$ can be viewed as the $`z`$-component of a two-dimensional gauge connection with values in the Lie algebra $`sl(3)`$. After supplementing $`𝒜`$ with a $`\overline{z}`$-component, one can derive the $`W_3`$-algebra by imposing a zero curvature condition on the connection .
## 4 $`N=1`$ supersymmetry
### 4.1 General framework
The $`N=1`$ supersymmetric generalization of the previous results has been worked out in references (see also ) by using a superspace approach. We note that $`N=1`$ superspace is locally parametrized by complex coordinates $`z`$ and $`\theta `$ which are even and odd, respectively. The transition from ordinary space to superspace can be summarized as follows:
$$\begin{array}{c}\text{Riemann surface}\hfill \\ z\hfill \\ \hfill \\ \text{conformal transformation}\hfill \\ \text{conformal field}c_k\hfill \\ \text{projective connection}R_{zz}(z)\hfill \\ L_2=^2+\frac{1}{2}R_{zz}\hfill \end{array}\}\{\begin{array}{c}\text{super Riemann surface}\hfill \\ z,\theta \hfill \\ ,D\frac{}{\theta }+\theta (D^2=)\hfill \\ \text{superconformal transf.}:Dz^{}=\theta ^{}D\theta ^{}\hfill \\ \text{superconformal field}:𝒞_k^{}=(D\theta ^{})^k𝒞_k\hfill \\ \text{superprojective connection}_{z\theta }(z,\theta )\hfill \\ _1=D^3+_{z\theta }.\hfill \end{array}$$
(25)
The odd superdifferential operator $`_1`$ acts on a superconformal field $`𝒞_1𝒞`$. By applying $`D`$ to $`_1𝒞`$ and subsequently projecting onto the lowest component of the resulting superfield, we find
$`(D_1𝒞)|`$ $`=`$ $`[D^4+(D_{z\theta })|]𝒞|_{z\theta }|(D𝒞_3)|`$
$`=`$ $`[^2+{\displaystyle \frac{1}{2}}R_{zz}]c+\rho _{z\theta }(D𝒞)|,`$
i.e. the basic Bol operator $`L_2`$ plus a fermionic contribution.
### 4.2 Matrix representation of super CCO’s
Let us rewrite the scalar, superconformally covariant differential equation
$$0=_1F_3\left[D^3+\right]F_3$$
(26)
as a system of three first-order differential equations:
$$\left[\begin{array}{c}0\\ 0\\ 0\end{array}\right]=\left[\begin{array}{ccccc}D& 0& & & \\ 1& D& 0& & \\ 0& 1& D& & \end{array}\right]\left[\begin{array}{c}F_1\\ F_2\\ F_3\end{array}\right]\{\begin{array}{c}0=DF_1+F_3\hfill \\ F_1=DF_2\hfill \\ F_2=DF_3.\hfill \end{array}$$
(27)
Analogously to the non supersymmetric theory, substitution of the last two equations into the first one reproduces the scalar equation (26).
If we now rewrite equation (27) in the form $`\stackrel{}{0}=(D𝒜)\stackrel{}{F}`$, we realize that the matrix $`𝒜`$ belongs to the Lie superalgebra $`sl(2|1)`$. However, the graded matrix $`𝒜`$ does not have the standard format which consists of arranging the even and odd matrix elements into blocks: this is an example of a nonstandard matrix format, to which we have referred as the diagonal format since there are alternatively even and odd diagonals .
This and other possible nonstandard matrix formats have been studied in a systematic way in reference . Although they are simply related to the standard format by a similarity transformation, they have many appealing features. Moreover, such formats naturally occur in various physical applications, e.g. in superconformal field theory, superintegrable models, for super $`W`$-algebras and quantum supergroups.
## 5 $`N=2`$ supersymmetry
A $`N=2`$ super Riemann surface is locally parametrized by an even complex coordinate $`z`$ and two odd complex coordinates $`\theta `$ and $`\overline{\theta }`$. There is a new feature in $`N=2`$ superspace geometry which makes this theory considerably richer and more complicated than the $`N=1`$ supersymmetric theory: the “square root” of the translation generator $``$ is not given by a single odd operator as in $`N=1`$ supersymmetry ($`D^2=`$), but it involves two odd operators,
$$D=\frac{}{\theta }+\frac{1}{2}\overline{\theta },\overline{D}=\frac{}{\overline{\theta }}+\frac{1}{2}\theta ,$$
(28)
satisfying
$$\{D,\overline{D}\}=$$
(29)
(and $`D^2=0=\overline{D}^2`$). Therefore, one has to deal with partial differential equations (involving $`D`$ and $`\overline{D}`$) rather than ordinary differential equations (only involving $`D`$). Another aspect of the algebra $`\{D,\overline{D}\}=`$ consists of the fact that it introduces a $`U(1)`$ symmetry into the theory: after projection from the super Riemann surface to the underlying ordinary Riemann surface, one thereby recovers $`U(1)`$-transformations in addition to the familiar conformal transformations. Henceforth, the Bol operators (17) acting on $`U(1)`$-neutral fields are to be generalized to conformally covariant operators acting on $`U(1)`$-charged fields. The latter as well as the original operators (17) arise from different types of $`N=2`$ super Bol operators which have been constructed and classified in reference . For a particular class of them, the so-called ‘sandwich operators’ (relating the chiral and anti-chiral subspaces of superconformal fields), one can give a matrix representation. The results following from a zero curvature condition for the operator product expansions of the $`N=2`$ super $`W_3`$-algebra coincide with those obtained by other methods .
## 6 Concluding comments
In these notes, we have tried to give a short introduction to some mathematical notions which are quite useful for the study of many physically interesting models in two dimensions. While the appearance of conformal symmetry in conformal models or in their non-linear generalizations (related to $`W`$-algebras) is quite natural, the role of conformal invariance in integrable models is less clear and still a matter of current research (see references and the contribution of M.Olshanetsky to this workshop).
Acknowledgments
I wish to thank my collaborators with whom I have worked on the topics reviewed in this talk, for our pleasant collaborations. Furthermore, I express my gratitude to Evgeny Ivanov, Sergey Krivonos and Anatoly Pashnev for their kind invitation to participate in the workshop, for nicely taking care of all practical matters and for rendering my visit to Dubna a most pleasant and rewarding one. The final version of this text was prepared while I was on sabbatical leave at the Technical University of Vienna: I wish to thank Prof.M.Schweda for his kind invitation and all the members of the Institut für Theoretische Physik for the hospitality extended to me.
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# Perpendicular transport properties of YBa2Cu3O7-δ/PrBa2Cu3O7-δ superlattices.
## Abstract
The coupling between the superconducting planes of YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> / PrBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> superlattices has been measured by c-axis transport. We show that only by changing the thickness of the superconducting YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> layers, it is possible to switch between quasi-particle and Josephson tunneling. From our data we deduce a low temperature c-axis coherence length of $`\xi _c=0.27`$ nm.
, ,
thanks: Corresponding author. E-mail: martinez@mail.uni-mainz.de
Artificial YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub>/PrBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> superlattices (Y123/Pr123) constitute ideal model systems for isolating given properties of High Temperature Superconductors. In those systems it is possible for instance to modify the c-axis tunneling properties simply by varying the periodicities of the Y123 and Pr123 layers.
Series of 200 nm thick Y123/Pr123 superlattices have been prepared by high-pressure dc-sputtering. The high quality of the samples was checked by detecting up to third order satellite peaks in x-ray $`\theta 2\theta `$ scans. Later on series of mesa structures with dimensions between $`15\times 15`$ and $`50\times 50`$ $`\mu `$m<sup>2</sup> were prepared by ion milling. In this work we present low temperature data on 2:7 (2 layers of Y123 and 7 of Pr123) and 8:8 superlattices.
We measured simultaneously the $`U`$ vs. $`I`$ characteristics and the differential conductivity $`\sigma (U)`$ by means of a standard Lock-In technique. In Fig. 1 we show the $`\sigma (U)`$ on a $`30\times 30\mu `$m<sup>2</sup> mesa done on a 2:7 superlattice at 2.0 K. No superconducting current could be detected. However the peak in $`\sigma (U)`$ corresponds to a $`c`$-axis superconducting gap. From the peak to peak voltage $`U_{pp}`$ peak, we estimate that each of the $`n=`$8 to 10 bi-layers constituting this mesa have a $`c`$-axis gap $`\mathrm{\Delta }_c=U_{pp}/4n=5.0\pm 0.5`$ meV. This value is in excellent agreement with the value of $`\mathrm{\Delta }_c`$ given in the literature which scatters between 4 and 6 meV for planar junctions .
In Fig. 1, we observe sharper features in $`\sigma (U)`$. In order to verify a quasi-periodicity, we marked each minimum by a vertical line which is associated to an integer. In Fig. 1 we plot this index as a function of the minima position. A clear zero crossing of the linear fit is obtained by choosing an index n=9 for the lowest index. From the linear fit we deduce a period of $`(11.1\pm 0.5)`$ mV which gives a periodicity $`\delta U=(1.2\pm 0.1)`$ meV for a single junction. These features are reproducible and temperature independent up to $`20`$ K.
Such a quasi-periodic structure in the density of states has been theoretically predicted by Hahn . According to this work, above the superconducting gap, additional structures should appear with a periodicity of: $`\xi _c/s=1/\pi ^2\delta U/\mathrm{\Delta }_c`$ where $`s`$ is the period of the superlattice and $`\xi _c`$ the $`c`$-axis coherence length. From our data and by taking $`s=10.5`$ nm we deduce a $`c`$-axis coherence length of $`\xi _c=0.27`$ nm. If we assume for Y123 an anisotropy of $`\gamma 5`$ we would obtain a in-plane coherence length $`\xi _{ab}=\gamma \xi _c1.4`$ nm. This value is close to the generally quoted $`\xi _{ab}=1.5`$ nm for Y123 .
In Fig. 2, we show the $`I`$ vs. $`U`$ characteristics of a $`40\times 40\mu `$m<sup>2</sup> mesa done on a 8:8 superlattice for zero field and for $`B=1`$ T. From the difference between the two curves we deduce the presence of two distinct low and high current regimes. To investigate the difference between these two regimes we measured the $`B`$ dependence of $`\sigma (U=0)`$ at zero bias current and at $`115`$ $`\mu `$A (see inserts).
At zero current $`\sigma (B)`$ shows a modulation of $`B_\mathrm{\Phi }=0.7`$ T. By considering that $`B_\mathrm{\Phi }=\mathrm{\Phi }_0/(s+t)b_{eff}`$, s=9.4 nm (Y123 thickness) and t=9.4 nm (Pr123 thickness), we deduce $`b_{eff}=0.152`$ $`\mu `$m. The low current regime can therefore be associated to structural shorts . At $`115`$ $`\mu `$A we observe a $`\sigma (B)`$ modulation of $`2.4`$ mT. This value is very similar to the $`2.5`$ mT predicted for a $`40\times 40`$ $`\mu `$m<sup>2</sup> mesa. The high current regime can therefore be associated to a c-axis Josephson Effect. The absence of a similar effect in the 2:7 superlattices is probably due to a non-fully developed superconducting order parameter in the 2 unit cells thick Y123 layers.
This work was supported by the German BMBF (Contract 13N6916) and the E.U. (ERBFMBICT972217).
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# A Model of Supernovae Feedback in Galaxy Formation
## 1 Introduction
Since the pioneering paper of White and Rees (1978), it has been clear that some type of feedback mechanism is required to explain the shape of the galaxy luminosity function in hierarchical clustering theories. The reason for this is easy to understand; if the power spectrum of mass fluctuations is approximated as a power law $`P(k)k^n`$, the Press-Schechter (1974) theory for the distribution of virialized haloes predicts a power law dependence at low masses
$$\frac{dN(m)}{dm}m^{(9n)/6}.$$
(1)
For any reasonable value of the index $`n`$ ($`n2`$ on the scales relevant to galaxy formation in cold dark matter (CDM) models), equation (1) predicts a much steeper mass spectrum than the observed faint end slope of the galaxy luminosity function, $`dN(L)/dLL^\alpha `$, with $`\alpha 1`$ (Efstathiou, Ellis and Peterson 1988, Loveday et al. 1992, Zucca et al. 1997). Furthermore, the cooling times of collisionally ionised gas clouds forming at high redshift are short compared to the Hubble time (Rees and Ostriker 1977, White and Rees 1978). Thus, in the absence of feedback, one would expect that a large fraction of the baryons would have collapsed at high redshift into low mass dark matter haloes, in contradiction with observations.
In reality, there are a number of complex physical mechanisms that can influence galaxy formation and these need to be understood if we are to construct a realistic model of galaxy formation. In the ‘standard’ cold dark matter model (i.e. nearly scale invariant adiabatic perturbations), the first generation of collapsed objects will form in haloes with low virial temperatures ($`T\stackrel{<}{}10^4`$K, characteristic circular speeds $`v_c\stackrel{<}{}20\mathrm{k}\mathrm{m}\mathrm{s}^1`$). Molecular hydrogen is the dominant coolant at such low temperatures and so an analysis of the formation of the first stellar objects requires an understanding of the molecular hydrogen abundance and how this is influenced by the ambient ultraviolet radiation field (Haiman, Rees and Loeb 1997, Haiman, Abel and Rees 1999). As the background UV flux rises, the temperature of the intergalactic medium will rise to $`10^4`$K (e.g. Gnedin and Ostriker 1997) and the UV background will reduce the effectiveness of cooling in low density highly ionized gas (Efstathiou 1992). A UV background can therefore suppress the collapse of gas in regions of low overdensity. It is this low density photoionized gas that we believe accounts for the Ly$`\alpha `$ absorption lines (Cen et al. 1994, Hernquist et al. 1996, Theuns et al. 1998, Bryan et al. 1999). Photoionization can also suppress the collapse of gas in haloes with circular speeds of up to $`v_c20`$$`30\mathrm{k}\mathrm{m}\mathrm{s}^1`$. However, numerical simulations have shown that a UV background cannot prevent the collapse of gas in haloes with higher circular speeds, though it can reduce significantly the efficiency with which low density gas is accreted onto massive galaxies (Quinn, Katz and Efstathiou 1996, Navarro and Steinmetz 1997).
To explain the galaxy luminosity function, feedback is required in galaxies with circular speeds $`v_c\stackrel{>}{}50\mathrm{kms}^1`$ with characteristic virial temperatures of $`\stackrel{>}{}10^5`$K. Energy injection from supernovae is probably the most plausible feedback mechanism for systems with such high virial temperatures. Winds from quasars might also disrupt galaxy formation (Silk and Rees 1998) or, more plausibly, limit the growth of the central black hole (Fabian 1999). Here we will be concerned exclusively with supernova driven feedback and will not consider feedback from an active nucleus. Simple parametric models of supernovae feedback were developed by White and Rees (1978) and White and Frenk (1991) and form a key ingredient of semi-analytic models of galaxy formation (e.g. Kauffmann, White and Guiderdoni 1993, Lacey et al. 1993, Cole et al. 1994, Baugh et al. 1996, 1998, Somerville and Primack 1999). In this paper, we develop a more detailed model of the feedback process itself. Previous papers on supernovae feedback include those of Larson (1974), Dekel and Silk (1986) and Babul and Rees (1992). These authors compute the energy injected by supernovae into a uniform interstellar medium (ISM) and apply a simple binding energy criterion to assess whether the ISM will be driven out of the galaxy. The feedback process in these models is explosive, operating on the characteristic timescale of $`10^6`$$`10^7`$ yrs for supernova remnants to overlap. This is much shorter than the typical infall timescale of hot gas in the halo, begging the question of how a resevoir of cold gas accumulated in the first place. The present paper differs in that we model the ISM as a two-phase medium consisting of cold clouds and a hot pressure confining medium, i.e. as a simplified version of the three-phase model of the ISM developed by McKee and Ostriker (1977, hereafter MO77). The cold component contains most of the gaseous mass of the disc and is converted into a hot phase by thermal evaporation in expanding supernovae remnants. In this type of model, the cold phase can be lost gradually in a galactic wind as it is slowly converted into a hot phase.
The main result of this paper is that low rates of star formation can expel a large fraction of the baryonic mass in dwarf galaxies over a relatively long timescale of $`1`$ Gyr. We therefore propose that effective feedback can operate in an steady, unspectacular mode; strong bursts of star formation and superwind-like phenomena (e.g. Heckman, Armus and Miley 1990) are not required, although galaxies may experience additional feedback of this sort. In fact, hydrodynamic simulations suggest that nuclear starbursts are ineffective in removing the ISM from galaxies with gas masses $`\stackrel{>}{}10^6M_{}`$ (Mac Low and Ferrara 1999, Strickland and Stevens 1999) because hot gas generated in the nuclear regions is a expelled in a bipolar outflow without coupling to the cool gas in the rest of the disc. This result provides additional motivation for investigating a “quiescent” mode of feedback. Silk (1997) describes a model which is similar, in some respects, to the model described here. However, the model described here is more detailed and allows a crude investigation of the radial properties of a disc galaxy during formation.
The layout of this paper is as follows. A simple model of star formation regulated by disc instabilities is described in Section 2. This is applied to ‘closed box’ (i.e. no infall or outflow of gas) models of disc galaxies neglecting feedback. Section 3 describes a model of the interaction of expanding supernovae shells in a two-phase ISM. This section is based on the model of MO77, but instead of focussing on equilibrium solutions that might apply to our own Galaxy, we compute the net rate of conversion of cold gas to hot gas incorporating the model for self-regulating star formation. This yields the temperature and density of the hot phase as a function of time and radius within the disc. Section 4 revisits the model of Section 2, but includes simultaneous infall and outflow of gas. This model is extended in Section 5 to include a galactic fountain, the pressure response of the cold ISM to the hot phase, and a model of chemical evolution. Section 6 describes some results from this model and discusses the effects of varying some of the input parameters. In addition, the efficiency of feedback is computed as a function of the circular speed of the surrounding dark matter halo. Our conclusions are summarized in Section 7. Although we focus on disc galaxies in this paper, a similar formalism could be applied to the formation of bulges if the assumption that gas conserves its angular momentum during collapse is relaxed.
## 2 Star Formation Regulated by Disc Instabilities
### 2.1 Rotation curve for the Disc and Halo
The dark halo is assumed to be described by the Navarro, Frenk and White (1996, hereafter NFW) profile
$$\rho _H(r)=\frac{\delta _c\rho _c}{(cx)(1+cx)^2},xr/r_v,$$
(2)
where $`\rho _c`$ is the critical density, $`r_v`$ is the virial radius at which the halo has a mean overdensity of $`200`$ with respect to the background and $`c`$ is a concentration parameter (approximately $`10`$ for CDM models). The circular speed corresponding to this profile is
$$v_H^2(r)=v_v^2\frac{1}{x}\frac{\left[\mathrm{ln}(1+cx)cx/(1+cx)\right]}{\left[\mathrm{ln}(1+c)c/(1+c)\right]},v_v^2\frac{GM_v}{r_v},$$
(3)
where $`M_v`$ is the mass of the halo within the virial radius.
We assume that the disc surface mass density distribution is described by an exponential,
$$\mu _D(r)=\mu _0\mathrm{exp}(r/r_D),M_D2\pi \mu _0r_D^2,$$
(4)
where $`M_D`$ is the total disc mass. The rotation curve of a cold exponential disc is given by (Freeman 1970)
$`v_D^2(r)=2v_c^2y^2\left[I_0(y)K_0(y)I_1(y)K_1(y)\right],`$ (5)
$`y{\displaystyle \frac{1}{2}}{\displaystyle \frac{r}{r_D}},v_c^2={\displaystyle \frac{GM_D}{r_D}}.`$
To relate the disc scale length, $`r_D`$, to the virial radius of the halo $`r_v`$, we assume that the angular momentum of the disc material acquired by tidal torques is approximately conserved during the collapse of the disc (see Fall and Efstathiou 1980). This fixes the collapse factor
$$f_{coll}=\frac{r_V}{r_D}$$
(6)
in terms of the dimensionless spin parameter $`\lambda _HJ|E|^{1/2}G^1M^{5/2}`$ of the halo component. The spin parameter is found to have a median value of $`0.05`$ from N-body simulations (Barnes and Efstathiou 1987), and for the models described here, this value is reproduced for collapse factors of around $`50`$. A more detailed calculation of the collapse factor of the disc is given in Section 4.
### 2.2 Vertical scale height of the disc
The velocity dispersion of the cold gas clouds in the vertical direction is assumed to be constant and equal to $`\sigma _g^2`$. The equations of stellar hydrodynamics then give the following solution
$$\rho (z)=\frac{\mu _g}{2H_g}\mathrm{sech}^2\left(\frac{z}{H_g}\right),$$
(7)
where $`\mu _g`$ is the surface mass density of the gas and the scale height is given by
$$H_g=\frac{\sigma _g^2}{\pi G\mu _g}.$$
(8)
Equation (8) must be modified to take into account the stellar disc. We do this approximately by assuming ‘disc pressure equilibrium’ (Talbot and Arnett 1975)
$$H_g=\frac{\sigma _g^2}{\pi G\mu _g}\frac{1}{(1+\beta /\alpha )},$$
(9)
where the quantities $`\alpha `$ and $`\beta `$ relate the vertical velocity dispersion $`\sigma _{}^2`$ and surface mass density $`\mu _{}`$ of the stars to those of the gas clouds
$`\sigma _{}=\alpha \sigma _g,`$ (10a)
$`\mu _{}=\beta \mu _g.`$ (10b)
### 2.3 Stability of a two-component rotating disc
The stability of rotating discs of gas and collisionless particles to axisymmetric modes has been analysed in classic papers by Goldreich and Lynden-Bell (1965) and by Toomre (1964). Here we use the results of Jog and Solomon (1984) who analysed the stability of a rotating disc consisting of two isothermal fluids of sound speeds $`c_1`$, $`c_2`$ and surface mass densities $`\mu _1`$ and $`\mu _2`$. These authors find that such a disc is stable to axisymmetric modes of wavenumber $`k`$ if
$$x=\frac{2\pi G\mu _1}{\kappa ^2}\frac{k}{(1+k^2c_1^2/\kappa ^2)}+\frac{2\pi G\mu _2}{\kappa ^2}\frac{k}{(1+k^2c_2^2/\kappa ^2)}<1,$$
(11)
where $`\kappa `$ is the epicyclic frequency
$$\kappa =2\omega \left(1+\frac{1}{2}\frac{r}{\omega }\frac{d\omega }{dr}\right)^{1/2}.$$
(12)
Equation (11) yields a cubic equation for the most unstable mode $`k_m`$. Solving this equation in terms of the parameters $`\alpha `$ and $`\beta `$ of equations (10), and ignoring the small differences between a gaseous and collisionless disc, we can write the stability criterion for a two-component system as
$$\sigma _g=\frac{\pi G\mu _g}{\kappa }g(\alpha ,\beta ).$$
(13)
This is identical to the Goldreich-Lynden-Bell criterion except for the factor $`g(\alpha ,\beta )`$. This factor is plotted in Figure 1 for various values of $`\alpha `$ and $`\beta `$.
### 2.4 Star formation and supernovae energy input
We assume a stellar initial mass function (IMF) of the standard Salpeter (1955) form
$`{\displaystyle \frac{dN_{}}{dm}}=Am^{(1+x)},m_l<m<m_u,x=1.35,`$ (14)
$`m_l=0.1M_{},m_u=50M_{},`$
and that each star of mass greater than $`8M_{}`$ releases $`10^{51}E_{51}`$ ergs in kinetic energy in a supernova explosion. For the IMF of equation (14), one supernova is formed for every $`125M_{}`$ of star formation. The energy injection rate is therefore related to the star formation rate by
$$\dot{E}_{sn}=2.5\times 10^{41}E_{51}\dot{M}_{}\mathrm{erg}/\mathrm{sec},$$
(15)
where $`\dot{M}_{}`$ is the star formation rate in $`M_{}`$ per year.
### 2.5 Energy dissipated by cloud collisions
We assume cold clouds of constant density $`\overline{\rho }_c=7\times 10^{23}`$ g/cm<sup>3</sup> with a distribution of cloud radii
$`{\displaystyle \frac{dN_{ca}}{da}}=N_0a^4,a_l<a<a_u,`$ (16)
$`a_l=0.5\mathrm{pc},a_u=10\mathrm{pc},`$
(MO77). Following MO77, the clouds are assumed to have an isotropic Gaussian velocity distribution with velocity dispersion independent of cloud size and that they lose energy through inelastic collisions. The rate of energy loss per unit volume is given by
$`{\displaystyle \frac{dE_{coll}}{dtdV}}=24\pi ^{3/2}\overline{\rho }_cN_{cl}^2a_l^5\sigma _g^3I_a,`$ (17)
$`I_a={\displaystyle \frac{1}{2}}{\displaystyle _1^{a_u/a_l}}{\displaystyle _1^{a_u/a_l}}{\displaystyle \frac{(x+y)^2}{(x^3+y^3)}}{\displaystyle \frac{dx}{x}}{\displaystyle \frac{dy}{y}},`$
where $`N_{cl}`$ is the local cloud density $`N_{cl}=N_0/3a_l^3`$. Integrating equation (17) over the vertical direction and using equation (9) for the scale height, the rate of energy loss per unit surface area $`\dot{E}_{coll}^\mathrm{\Omega }`$ is
$`\dot{E}_{coll}^\mathrm{\Omega }=5.0\times 10^{29}\left(1+{\displaystyle \frac{\beta }{\alpha }}\right)\mu _{g5}^3\sigma _{g5}\mathrm{erg}/\mathrm{sec}/\mathrm{pc}^2,`$ (18)
where $`\mu _{g5}`$ is the surface mass density of the gas component in units of $`5M_{}/\mathrm{pc}^2`$ and $`\sigma _{g5}`$ is the cloud velocity dispersion in units of $`5`$ km/sec. These values are close to those observed in the local solar neighbourhood. To estimate the efficiency with which supernovae accelerate the system of clouds, we normalize to the observed net star formation rate of the Milky Way. Assuming that the gas distribution has a flat surface mass density profile to $`R_{max}=14`$ kpc (Mihalas and Binney 1981), $`\beta 10`$, $`\alpha 5`$, and equating the integral of (18) to $`ϵ_c\dot{E}_{sn}`$ (equation 15), we find
$$ϵ_cE_{51}\dot{M}_{}=0.004.$$
(19)
An efficiency parameter of $`ϵ_c=0.01`$ produces a net star formation rate of $`0.4M_{}`$/yr which is reasonable for a Milky Way-like galaxy. We will therefore adopt a constant value of $`ϵ_c=0.01`$ in the models of the next subsection. The value of $`ϵ_c`$ will, of course, depend on the properties of the clouds, ISM and star formation rate. For example, in the model of MO77 the clouds are accelerated by interactions with the cold shells surrounding supernova remnants and they find efficiencies $`ϵ_c`$ of typically a few percent. We investigate the effect of varying $`ϵ_c`$ in Section 6.
### 2.6 Self-regulating models without inflow or outflow
The equations derived above allow us to evolve an initially gaseous disc and to compute the local star formation rate, cloud velocity dispersion etc. The system of stars and gas is constrained to satisfy the stability criterion of equation (13), which fixes the cloud velocity dispersion $`\sigma _g`$. There is some empirical evidence that star formation in nearby galaxies is regulated by a stability criterion of this sort (e.g. Kennicutt 1998). The energy lost in cloud collisions (equation 18) is balanced against the energy input from supernovae assuming a constant efficiency factor $`ϵ_c=0.01`$. We assume further that $`\alpha =5`$ (equation 10a) i.e. that stars are instantaneously accelerated to higher random velocities than the system of gas clouds, and that the properties of the gas clouds (mass spectrum, internal density, etc) are independent of time. These are clearly restrictive assumptions, but they allow us to generate simple models of self-regulating star formation with only one free parameter $`ϵ_c`$.
We study the evolution of two model galaxies with parameters listed in Table 1. Model MW has parameters roughly similar to those of the Milky Way and model DW has parameters similar to those of a relatively high surface brightness dwarf galaxy.
Figure 2 shows the evolution of the gas and stellar surface mass densities of the two models. The net star formation rates, gas fractions and mean gas cloud velocity dispersion are plotted in Figure 3. In model MW, the star formation rates are initially high ($`>100M_{}`$/yr) and hence the timescale for star formation is short; half the disc mass is converted into stars in $`10^7`$ years. The star formation rate declines rapidly to less than $`1M_{}`$/yr after a few Gyr. As figure 2 shows, the star formation at early times is concentrated to the inner parts of the disc which have a high surface density and hence the gas distribution develops a characteristic surface density profile with an inner ‘hole’, similar to what is seen in the HI distributions in real galaxies (see Burton 1991). The stellar disc is truncated at about the Holmberg (1958) radius ($`r/r_D5`$), in rough agreement with observations. The truncation arises because the gas disc becomes thick at large radii (equation 9) and the rate of energy lost in cloud collisions can be balanced by a very low star formation rate.
The evolution of model DW is qualitatively similar, though the star formation rate is scaled down roughly in proportion to the disc mass. Half the gas is converted to stars by $`3\times 10^7`$ yr, and the gas fraction is $`0.12`$ after $`10^{10}`$ yr, similar to the final gas fraction of $`0.13`$ in model MW.
Neither of these models is satisfactory. The star formation rate in model MW is too high at early times to be compatible with deep number counts (see e.g. Ellis 1997), which require more gentle star formation rates in typical $`L^{}`$ galaxies. Model SW converts most of its gas into stars on a short timescale and so does not solve the problem raised in the introduction of explaining the flat faint end slope of the luminosity function in CDM-like models. As we will see in later sections, infall of gas provides the solution to the former problem, since this allows the disc to build up gradually on a cooling or dynamical timescale. Outflow of hot gas heated by supernovae provides a solution to the latter problem.
## 3 Evolution of a two-phase ISM
In this Section we consider the interaction of a multiphase interstellar medium with expanding supernova remnants following the model of MO77 and discuss the conditions under which a protogalaxy can form a wind. The key ingredients of the model are as follows. Most of the cold gaseous mass is assumed to be in cold clouds with properties as given in Section 2.5. Supernovae explode and their remnants propagate evaporating some of the cold clouds and forming a low density hot phase of the ISM. The star formation rate therefore determines the evaporation rate and hence the rate of conversion of the cold phase to a hot phase. A wind from the galaxy can result if the hot phase is: (i) sufficiently pervasive (filling factor of order unity), (ii) low density (so that radiative cooling is unimportant) and (iii) the temperature of the hot phase exceeds the virial temperature of the galaxy. In this section we follow closely the theory of the ISM developed by MO77 and we use their notation where possible.
### 3.1 Evaporation of cold clouds
An expanding supernovae remnant will evaporate a mass of
$$M_{ev}540E_{51}^{6/5}\mathrm{\Sigma }^{3/5}n_h^{4/5}M_{},$$
(20)
where $`n_h`$ is the density interior to the supernovae remnant and $`\mathrm{\Sigma }`$ (in pc<sup>2</sup>) is the evaporation parameter introduced by MO77
$$\mathrm{\Sigma }=\frac{\gamma }{4\pi a_lN_{cl}\varphi _\kappa }.$$
(21)
Here the parameter $`\gamma `$ relates the blast wave velocity to the isothermal sound speed ($`v_b=\gamma c_h`$, $`\gamma 2.5`$) and the parameter $`\varphi _\kappa `$ quantifies the effectiveness of the classical thermal conductivity of the clouds ($`\kappa _{eff}=\kappa \varphi _\kappa `$) and so is less than unity if the conductivity is reduced by tangled magnetic fields, turbulence etc. Using equations (7) and (9) to estimate the mean cloud density we find
$$\mathrm{\Sigma }280\frac{\sigma _{g5}^2}{\mu _{g5}^2}\frac{1}{(1+\beta /\alpha )}\frac{1}{\varphi _\kappa }\mathrm{pc}^2=f_\mathrm{\Sigma }\mathrm{\Sigma }_{},\mathrm{\Sigma }_{}95\mathrm{pc}^2,$$
(22)
where $`\mathrm{\Sigma }_{}`$ is the evaporation parameter characteristic of the local solar neighbourhood ($`\beta /\alpha 2)`$.
Evaluating equation (20), we find
$$M_{ev}1390E_{51}^{6/5}f_\mathrm{\Sigma }^{3/5}\varphi _\kappa ^{3/5}n_{h2}^{4/5}M_{},$$
(23)
where $`n_{h2}`$ is $`n_h`$ in units of $`10^2\mathrm{cm}^2`$ (a characteristic value for the hot component). Thus, provided thermal conduction is not highly suppressed, a single supernovae remnant can evaporate a much larger mass than the $`125M_{}`$ formed in stars per supernovae for a standard Saltpeter IMF (Section 2.4). If a significant fraction of this evaporated gas can escape in a wind, then star formation will be efficiently suppressed.
### 3.2 Temperature and density of the hot phase
To compute the properties of the hot phase we assume that the disc achieves a state in which the porosity parameter $`Q`$ is equal to unity. The disc is then permeated by a network of overlapping supernovae remnants. Ignoring cooling interior to the remnants (which we will see is a reasonable approximation for an ISM with low metallicity) the age, radius and temperature of a SNR when $`Q=1`$ are given by
$`t_o=5.5\times 10^6S_{13}^{5/11}\gamma ^{6/11}E_{51}^{3/11}n_h^{3/11}\mathrm{yr},`$ (24a)
$`R_o=100S_{13}^{2/11}\gamma ^{2/11}E_{51}^{1/11}n_h^{1/11}\mathrm{pc},`$ (24b)
$`T_o=1.2\times 10^4S_{13}^{6/11}\gamma ^{6/11}E_{51}^{8/11}n_h^{8/11}\mathrm{K}.`$ (24c)
where $`S_{13}`$ is the supernova rate in units of $`10^{13}\mathrm{pc}^3\mathrm{yr}^1`$. The density of a remnant at $`t_o`$ ($`n_h^oM_{ev}/(4/3\pi R_o^3)`$), gives an approximate estimate of the density of the ambient hot phase
$$n_h^o4.3\times 10^3S_{13}^{0.36}\gamma ^{0.36}E_{51}^{0.61}f_\mathrm{\Sigma }^{0.393}\mathrm{cm}^3.$$
(25)
Inserting this estimate into equations (24) we find
$`t_o=1.2\times 10^6S_{13}^{0.36}\gamma ^{0.64}(E_{51}f_\mathrm{\Sigma })^{0.11}\mathrm{yr},`$ (26a)
$`R_o=164(S_{13}/\gamma )^{0.21}E_{51}^{0.04}f_\mathrm{\Sigma }^{0.035}\mathrm{pc},`$ (26b)
$`T_o=6.6\times 10^5(S_{13}E_{51}f_\mathrm{\Sigma }/\gamma )^{0.29}\mathrm{K},`$ (26c)
and the rate at which clouds are evaporated is
$$\dot{M}_{ev}=2.7\times 10^{10}S_{13}^{0.71}\gamma ^{0.29}E_{51}^{0.71}f_\mathrm{\Sigma }^{0.29}\mathrm{M}_{}\mathrm{pc}^3\mathrm{yr}^1.$$
(27)
Integrating equation (27) over the scale height of the disc gives the evaporated mass per unit area,
$`\dot{M}_{ev}^\mathrm{\Omega }1\times 10^7\left({\displaystyle \frac{\sigma _{g5}^2}{\mu _{g5}(1+\beta /\alpha )}}\right)\times `$
$`S_{13}^{0.71}\gamma ^{0.29}E_{51}^{0.71}f_\mathrm{\Sigma }^{0.29}\mathrm{M}_{}\mathrm{pc}^2\mathrm{yr}^1.`$ (28)
Adopting a cooling rate of $`\mathrm{\Lambda }2.5\times 10^{22}T_5^{1.4}`$ erg cm<sup>3</sup> s<sup>-1</sup> for $`10^5\stackrel{<}{}T\stackrel{<}{}10^6`$ for a gas with primordial composition, the ratio of $`t_o`$ to the cooling time $`t_{\mathrm{cool}}`$ is
$$\frac{t_o}{t_{\mathrm{cool}}}0.5T_5^{2.4}f_\mathrm{\Sigma }^{0.5},$$
(29)
Thus if the temperature of the hot phase is higher than about $`10^5`$K, the assumption that cooling can be neglected will be valid. A cooling function for a gas with primordial composition will be used throughout this paper. As the metallicity of the gas builds up, the cooling time of the hot component will shorten and more of the supernovae energy will be lost radiatively. This effect will reduce the efficiency of feedback in galaxies with high metallicity but is not included in this paper.
### 3.3 Simple self-regulating model with outflow
In this section we apply the results of the previous paragraphs to construct a simplified self-regulating model with outflow. The star formation rate is governed by the self-regulation algorithm as in Section 2.6 with the parameter $`ϵ_c=0.01`$. This provides an estimate of the local supernova rate per unit volume which we insert in equations (26) to compute the properties of the hot phase, adopting a value $`\varphi _\kappa =0.1`$ in equation (21) for the conduction efficiency parameter. The hot gas will be lost from the system if its specific enthalpy
$$\frac{1}{2}v^2+\frac{5}{2}\frac{p}{\rho }$$
exceeds to within a factor of order unity its gravitational binding energy per unit mass. If the gas has an initial isothermal sound speed of $`c_i=(kT/\mu _p)=37T_5^{1/2}\mathrm{kms}^1`$ (for a mean mass per particle of $`\mu _p=0.61m_p`$), conservation of specific enthalpy implies that the wind will reach a bulk speed of $`v_w\sqrt{5}c_i`$. Some of the thermal energy will be lost radiatively, and in fact the spherical steady wind solutions described in Appendix B suggest that a more accurate criterion for the wind to escape from a galaxy is $`v_w\sqrt{2.5}c_i>v_{esc}`$, where $`v_{esc}`$ is the escape speed from the centre of the halo (neglecting the potential of the disc). If $`\sqrt{2.5}c_i<v_{esc}`$, the hot phase is returned instantaneously to the cold phase. This type of binding energy criterion for outflow has been adopted in previous studies (e.g. Larson 1974, Dekel and Silk 1986) and is clearly oversimplified, as are the assumptions of instantaneous mass loss and return of cold gas. These points will be discussed further in Section 6, but for the moment these assumptions will be adopted to illustrate the qualitative features of the model. As gas is lost from the system, the circular speed of the disc component (equation 5) is simply rescaled by the square root of the mass of the disc that remains.
The evolution of the surface mass densities for the two disc models is illustrated in Figures 4 and 5. In model MW, the evolution is similar to that without outflow shown in Figure 2. With the simple prescription for mass loss used here, no hot gas is lost unless the temperature of the hot phase exceeds $`T_{crit}5\times 10^6`$K. This does happen at early times when the star formation rate is high, and about $`25`$% of the galaxy mass is lost within $`10^7`$ yr. Thereafter, no more mass is lost and a nearly exponential disc is built up with a gas distribution containing a central hole as in Figure 2. The star formation rate in this model declines strongly with time, exceeding $`100M_{}`$/yr in the early phases of evolution.
The behaviour of model DW is qualitatively different. Here the critical temperature for mass loss is much lower, $`T_{crit}3\times 10^5`$ K, hence half the mass of the galaxy is expelled by $`10^8`$ yr and and $`66\%`$ by $`1`$ Gyr. After $`1`$ Gyr, the temperature of the hot phase drops below $`T_{crit}`$ and the galaxy settles into a stable state with a low rate of star formation.
The wind prescription in these models, and particularly the assumption that gas below the critical temperature necessary for escape is returned instantaneously to the cold phase, is clearly oversimplified and so the mass loss fractions should not be taken too seriously. A more detailed model is developed in Section 6. A more serious deficiency of the model presented here is that the entire gas disc is assumed to have formed instantaneously at $`t=0`$. This is unrealistic and leads to high rates of star formation and gas ejection at early times. A simple infall model, similar to those adopted in semi-analytic models (White and Frenk 1991, Cole et al. 1994) is included in the next Section.
## 4 Infall Model
### 4.1 Conservation of specific angular momentum
Following Fall and Efstathiou (1980), the gas is assumed to follow the spatial distribution of the halo component with the same distribution of specific angular momentum prior to collapse. The halo is assumed to rotate cylindrically with rotation speed $`v_H^{rot}(\varpi _H)`$, where $`\varpi _H`$ is the radial coordinate in the cylindrical coordinate system. The gas is assumed to conserve its specific angular momentum during its collapse, so that the final specific angular momentum of the disc at radius $`\varpi _D`$, $`h_D=\varpi _Dv_D^{rot}`$, is equal to the specific angular momentum of the halo $`h_D=\varpi _Hv_H^{rot}`$ at the radius $`\varpi _H`$ from which the gas originated. Mass conservation relates the radii $`\varpi _H`$ and $`\varpi _D`$,
$$\frac{d\varpi _H}{d\varpi _D}=\frac{\mu _D(\varpi _D)}{\mu _H(\varpi _H)}\frac{M_H}{M_D}\frac{\varpi _D}{\varpi _H},$$
(30)
where $`M_H/M_D`$ is the ratio of the halo to disc mass interior to the maximum infall radius of the disc (see figure 6a below) and $`\mu _H`$ is the projected surface mass density of the halo
$$\mu _H(\varpi )=2_0^{\mathrm{}}\rho _H\left((\varpi ^2+z^2)^{1/2}\right)𝑑z.$$
(31)
The solution of equation (30) yields $`\varpi _D(\varpi _H)`$ and the rotation speed of the halo follows from the conservation of specific angular momentum, $`v_H^{rot}=\varpi _Dv_D^{rot}(\varpi _D)/\varpi _H`$. The results for the parameters of models MW and DW are shown in figure 6, where we have used the notation $`s=\varpi /r_D`$. When expressed in the dimensionless units of Figure 6, the solutions for models MW and DW are identical.
This prescription is guaranteed to form an exponential disc with the required parameters. The derived rotation velocity of the halo is almost independent of radius in general agreement with what is found in N-body simulations (Frenk et al. 1988, Warren et al. 1992). The upturn in the halo rotation speed at $`s_H(max)30`$ is caused by the rapid decline in the mass of the input exponential disc at large radii and is of little consequence in the discussion that follows. The values of the spin parameter quoted in Table 1 were derived from the mass and binding energy of the halo and assuming that the the halo rotation velocity is constant at $`0.095v_c`$ at large radii.
### 4.2 Mass infall rate
To determine the gas infall rate we compute the free fall time for a gas element at rest at radius $`r_i`$,
$$t_{\mathrm{ff}}=_0^{r_i}\frac{dr}{\sqrt{2}[\varphi _H(r_i)\varphi _H(r)]^{1/2}},$$
(32)
and the cooling time
$$t_{\mathrm{cool}}=\frac{3}{2}\frac{kT_v\times 1.92}{\mathrm{\Lambda }(T_v)n_e(r)},$$
(33)
where $`n_e(r)`$ is the electron density. The temperature $`T_v`$ in equation (33) is set to the virial temperature derived from the equation of hydrostatic equilibrium assuming that the temperature is slowly varying with radius
$$T_vv_H^2(r)\frac{\mu _p}{k}\frac{d\mathrm{ln}r}{d\mathrm{ln}\rho _b(r)},$$
(34)
where we assume that the baryons follow the same spatial distribution as the halo. The infall rate is given by
$$\dot{M}_{inf}=4\pi \rho _b(r_H)r_H^2\{\begin{array}{cc}dr_H(t_{\mathrm{ff}}=t)/dt\hfill & t_{\mathrm{ff}}>t_{\mathrm{cool}}\hfill \\ dr_H(t_{\mathrm{cool}}=t)/dt\hfill & t_{\mathrm{cool}}>t_{\mathrm{ff}}\hfill \end{array}.$$
(35)
Finally, conservation of specific angular momentum specifies the final radius in the disc for each gas element. Since the halo is assumed to rotate on cylinders, the gas near to the poles in an infalling shell has a lower specific angular momentum that the gas at the equator. The infalling material is therefore distributed through the disc according to
$$2\pi \varpi _D\dot{\mu }_D(\varpi _D)d\varpi _D=\dot{M}_{inf}\frac{\varpi _Hd\varpi _H}{r_H(r_H^2\varpi _H^2)^{1/2}},$$
(36)
where $`\varpi _D`$ and $`\varpi _H`$ are related by the solution of equation (30).
Equations (32) – (36) specify the infall model. The free-fall and cooling times of the two model galaxies are shown in Figure 7. In the larger galaxy, gas within $`r_H/r_D10`$ infalls on the free-fall timescale and ends up within one scale length of the final disc. The material in the outer parts of the disc infalls on the cooling timescale. In contrast, apart from a small amount of gas in the very central part of the halo with virial temperature $`<10^4`$ K, the gas in the dwarf galaxy infalls on a free-fall timescale because the cooling time is so short.
### 4.3 Simple self-regulating model with inflow and outflow
The models described in this section are exactly the same as those described in section 3.3, except that we grow the discs gradually using the infall model of sections 4.1 and 4.2. In the models described below, inflow and outflow are assumed to occur simultaneously. This is often assumed in semi-analytic models of galaxy formation (e.g. Cole et al. 1994, Somerville and Primack 1999) and may not be completely unrealistic if the infalling gas is clumpy. The dark matter haloes will contain significant sub-structure (e.g. Moore et al. 1999) which may contain pockets of cooled gas. Furthermore, if the cooling time is short compared to the dynamical time, the infalling gas will be thermally unstable (Fall and Rees 1985) and will fragment into clouds. These will fall to the centre on a free-fall timescale if they are sufficiently dense and massive that gravity dominates over the ram pressure of the wind. This requires clouds with masses
$`m_{\mathrm{cloud}}\stackrel{>}{}9.5\times 10^5M_{}\left({\displaystyle \frac{a_{\mathrm{cloud}}}{1\mathrm{k}\mathrm{p}\mathrm{c}}}\right)\left({\displaystyle \frac{r}{10\mathrm{k}\mathrm{p}\mathrm{c}}}\right)^1\times `$
$`\left({\displaystyle \frac{\dot{M}_w}{1M_{}/\mathrm{yr}}}\right)\left({\displaystyle \frac{v_w}{100\mathrm{k}\mathrm{m}/\mathrm{s}}}\right)\left({\displaystyle \frac{v_v}{100\mathrm{k}\mathrm{m}/\mathrm{s}}}\right)^2,`$ (37)
where $`a_{\mathrm{cloud}}`$ is the radius of the cloud. However, even if (37) is satisfied, the clouds may be sheared and disrupted into smaller clouds by Kelvin-Helmholtz instabilities on a timescale of a few sound crossing times as they flow through the wind (e.g. Murray et al. 1993). The wind energy will be partially thermalized in shocks with the infalling clouds and dissipated in evaporating small clouds. But for the typical mass outflow rates expected from dwarf galaxies ($`\dot{M}_w\stackrel{<}{}0.2M_{}/\mathrm{yr}`$), the rate at which energy is supplied by the wind $`\dot{E}_w=1/2\dot{M}_wv_w^2`$ is much smaller than the energy lost in radiative cooling,
$`{\displaystyle \frac{\dot{E}_{\mathrm{cool}}}{\dot{E}_w}}50\mathrm{\Lambda }_{23}\left({\displaystyle \frac{v_v}{100\mathrm{k}\mathrm{m}/\mathrm{s}}}\right)^4\left({\displaystyle \frac{r_{\mathrm{cool}}}{10\mathrm{k}\mathrm{p}\mathrm{c}}}\right)^1\times `$
$`\left({\displaystyle \frac{v_w}{100\mathrm{k}\mathrm{m}/\mathrm{s}}}\right)^2\left({\displaystyle \frac{\dot{M}_w}{1M_{}/\mathrm{yr}}}\right)^1,`$ (38)
where $`r_{\mathrm{cool}}`$ is the radius at which the cooling time is equal to the age of the system.
The qualitative picture that we propose is as follows. In galaxies with a short cooling time, clouds formed by thermal instabilities will infall ballistically if (37) is satisfied. If (37) is not satisfied, the ram pressure of the wind will drive out the infalling gas and infall will be suppressed. With infall suppressed, the star formation rate in the disc and the wind energy will decline until infall can begin again. The wind will be partially thermalised before reaching $`r_{\mathrm{cool}}`$ and completely thermalised at $`r_{\mathrm{cool}}`$, but the energy supplied by the wind will be small compared to the energy radiated by the gas at $`r\stackrel{>}{}r_{\mathrm{cool}}`$ and so cannot prevent radiative cooling. If (37) is satisfied, some of the outflowing gas may fall back down to the disc after shocking against infalling clouds. However, in the models described here the efficiency of converting infalling gas into stars is low in dwarf systems, so provided the gas does not cycle around the halo many times, neglecting return of some of the outflowing gas should not affect the qualitative features of the models. The global geometry of the system, e.g. if the wind is weakly collimated perpendicular to the disc, may also permit simultaneous inflow and outflow of gas.
The interaction of an outflowing wind with an inhomogeneous infalling gas clearly poses a complex physical problem. In reality, the process may be far from steady, with outflow occurring in bursts accompanied by infall from discrete sub-clumps containing cooled gas. In the models described below and in the rest of this paper, we will assume that the infall and outflow occur simultaneously, steadily and without any interaction between the inflowing and outflowing gas. As the discussion of the preceeding two paragraphs indicates, this is obviously an over-simplification. It should be viewed as an idealization, on a similar footing to some of the other assumptions adopted in this paper (e.g. spherical symmetry, neglect of halo substructure and merging, steady star formation rates etc) designed to give some insight into how a quiescent mode of feedback might operate.
The analogues of figures 4 and 5 for the models incorporating infall and outflow are shown in figures 8 and 9. The discs build up from the inside out, as in the models of disc formation described by Fall and Efstathiou (1980) and Gunn (1982). Most of the star formation occurs in a propagating ring containing the most recently accreted gas. The most significant differences from the models of section 3.3, are the net rates of star formation (figure 9) and the timescale of outflow. The initial high rates of star formation in the models of section 3.3 are suppressed in the models with infall, and the timescale for outflow is now much longer because it is closely linked to the gas infall timescale. Apart from these differences, the final states, gas fractions and mass-loss fractions are similar to those in the models without infall. In model MW some outflow occurs when $`t\stackrel{<}{}10^8`$yrs and the temperature of the hot gas is high enough that it can escape from the system. Thereafter, the hot component cannot escape and the disc builds up without further outflow. About 17% of the total galaxy mass is expelled in the early phases of evolution, but as we have described above, this could be an overestimate since some of this gas may be returned to the galaxy if the wind energy is thermalized before it reaches the virial radius. In contrast, model DW drives a wind until $`t1`$ Gyr and expels about 74% of its mass. About half of the gas is lost within $`3\times 10^8`$ yrs, i.e. on about the infall timescale for most of the gas in the halo (see figure 7b).
## 5 Refinements of the Model
The models described in the previous sections contain a number of simplifications, which we will attempt to refine in this section. We do not address the problem of the interaction of a wind with the infalling gas, which is well beyond the scope of this paper. Instead, we introduce some simple improvements to the infall model (§5.1), model for mass loss from the galactic disc (§5.2, §5.3) and the pressure response of the cold ISM to the hot phase (§5.4, §5.5).
### 5.1 Infall Model
In Section 4, we used a simplified model of infall that guarantees the formation of an exponential disc if angular momentum is conserved during collapse. In this section, we assume a specific functional form for the rotation velocity of the halo,
$$v_H^{rot}(s)=c_1v_c\frac{(s/c_2)}{(1+(s/c_2))(1+(s/c_3)^{c_4})},sr/r_D$$
(39)
with $`c_1=0.115`$, $`c_2=0.6`$, $`c_3=16`$, $`c_4=0.25`$. The functional form and coefficients in equation (39) have been chosen to provide a good fit to the halo rotation profile derived in section 4.1 from conservation of specific angular momentum and is plotted as the dashed line in figure (6b). As in the previous section, the gas is assumed to follow the same radial density distribution and rotation velocity as the halo component, but its final radius in the disc is computed by assuming conservation of specific angular momentum and self-consistently solving for the rotation speed of the disc component. The halo component is assumed to be rigid and the contribution of the disc component to $`v_D^{rot}`$ is computed using the Fourier-Bessel theorem (see Binney and Tremaine 1987 §2.6)
$`v^2(r)=r{\displaystyle _0^{\mathrm{}}}S(k)J_1(kr)k𝑑k,`$
$`S(k)=2\pi G{\displaystyle _0^{\mathrm{}}}J_0(kr)\mu _D(r)r𝑑r.`$ (40)
Equation (40) is time consuming to evaluate accurately and in our application $`v^2(r)`$ must be computed many times. A fast algorithm has therefore been developed as described in Appendix A. The epicyclic frequency $`\kappa `$ is required in equation (13) to compute the instantaneous star formation rate and is derived by numerically differentiating the rotation speed.
With this formulation of the infall model, the infall rate is governed by the dark matter profile and the ratio of dark to baryonic mass within the virial radius, $`M_v/M_D=(v_v/v_c)^2f_{coll}`$. For the models described here, we adopt $`M_v/M_D=10`$, consistent with the parameters listed in Table 1. The final disc surface mass density will be close to an exponential by construction, since the halo rotation velocity (39) has been chosen to match the rotation profile derived by assuming an exponential disc and conservation of specific angular momentum.
### 5.2 Galactic Fountain
In previous sections, we have assumed that gas is lost from the disc if the bulk velocity of the wind $`v_w\sqrt{2.5}c_i`$, exceeds the escape speed $`v_{esc}`$ from the halo, but is otherwise returned instantaneously to the ISM. More realistically, gas with $`v_w<v_{esc}`$ will circulate in the halo along a roughly ballistic trajectory and will cool forming a galactic fountain (Shapiro and Field 1976, Bregman 1980). In the models described in this section, hot gas with $`v_w<v_{esc}`$ is returned to the disc at the radius from which it was expelled after a time $`t_{ret}`$,
$$t_{ret}=2t_{ff}(r_{max}),v_w^2=2[\varphi _H(r_{max})\varphi _H(0)],$$
(41)
i.e. we ignore the gravity of the disc and the angular momentum of the gas in computing the ballistic trajectory of a gas element.
### 5.3 Escape Velocity of the Wind
The detailed dynamics of the hot corona itself is complicated and beyond the scope of this paper. Type II supernovae at the upper and lower edges of the gas disc will be able to inject their energy directly into the hot gas, as will Type Ia supernovae forming in the thicker stellar disc. In addition, the hot component will interact with the primordial infalling gas in a complicated way as sketched in §4.3.
In the absence of radiative cooling, the hot gas will extend high above the galactic disc in an extended corona. For an isothermal corona, the equation of hydrostatic equilibrium in the z-direction has the following approximate solution,
$`{\displaystyle \frac{\rho _H(\varpi ,z)}{\rho _H(0)}}\mathrm{sech}^{2p_g}\left({\displaystyle \frac{z}{H_g}}\right)\mathrm{sech}^{2p_s}\left({\displaystyle \frac{z}{H_s}}\right)\times `$
$`\mathrm{exp}\left({\displaystyle \frac{1}{c_i^2}}{\displaystyle _0^z}{\displaystyle \frac{zv_H^2(r)}{r^2}}𝑑z\right),r^2=\varpi ^2+z^2,`$ (42)
$`p_g={\displaystyle \frac{\mu _g\sigma _g^2\sigma _s}{c_i^2(\mu _g\sigma _s+\mu _s\sigma _g)}},p_s={\displaystyle \frac{\mu _s\sigma _s^2\sigma _g}{c_i^2(\mu _g\sigma _s+\mu _s\sigma _g)}},`$
where we have assumed that both the stars and the gas follow $`\mathrm{sech}^2`$ vertical distributions (equation 7) and $`c_i`$ is the isothermal sound speed of the hot gas. We define a characteristic scale height for the hot component, $`H_{hot}(\varpi )`$, at which the density drops by a factor $`e`$ according to equation (42). If radiative cooling were negligible, we would expect a sonic point in the flow at about $`zH_{hot}`$. It is interesting to compare some characteristic numbers for the coronal gas:
$`\dot{E}_{in}1.1\times 10^{40}T_{h6}\dot{M}_{ev}\mathrm{erg}\mathrm{s}^1`$ (43a)
$`\dot{E}_{\mathrm{SNII}}5.7\times 10^{40}ϵ_{\mathrm{SNII}}E_{51}\dot{M}_{}\mathrm{erg}\mathrm{s}^1`$ (43b)
$`\dot{E}_{\mathrm{cool}}2\times 10^{39}n_{h2}^2\mathrm{\Lambda }_{23}\left({\displaystyle \frac{H_{hot}}{1\mathrm{kpc}}}\right)\left({\displaystyle \frac{R_{hot}}{3\mathrm{kpc}}}\right)^2\mathrm{erg}\mathrm{s}^1`$ (43c)
where the rates $`\dot{M}_{ev}`$, $`\dot{M}_{}`$ are in $`M_{}/\mathrm{yr}`$. Here $`\dot{E}_{in}`$ is the thermal energy injected into the hot corona by evaporating cold gas at a rate $`\dot{M}_{ev}`$, $`\dot{E}_{\mathrm{SNII}}`$ is the energy supplied to the corona by Type II supernovae forming above and below one scale height of the cold gas layer and the parameter $`ϵ_{\mathrm{SNII}}`$ expresses the efficiency with which this energy is coupled to the hot coronal gas. $`\dot{E}_{\mathrm{cool}}`$ is the rate of energy lost by a uniform density isothermal corona of scale height $`H_{hot}`$ within a cylinder of radius $`R_{hot}`$. For a large galaxy such as the Milky Way that can sustain an evaporation rate of $`10M_{}/\mathrm{yr}`$, $`\dot{E}_{\mathrm{cool}}`$ is small compared to $`\dot{E}_{in}`$ and it is a good approximation to neglect radiative cooling in the early stages of the flow (see Appendix B). However, in a dwarf galaxy $`\dot{E}_{\mathrm{cool}}`$ is typically larger than $`\dot{E}_{in}`$. In this case, we expect that the hot component will develop a sonic point at a characteristic cooling scale height $`H_{\mathrm{cool}}`$
$$H_{\mathrm{cool}}v_wt_{\mathrm{cool}}11\left(\frac{v_w}{100\mathrm{km}/\mathrm{s}}\right)T_{h6}n_{h2}^1\mathrm{\Lambda }_{23}^1\mathrm{kpc}.$$
(44)
(see e.g. Kahn 1981 and Appendix B) and that most of the thermal energy will be converted into kinetic energy by the time the gas flows to $`H_{\mathrm{cool}}`$. We therefore ignore radiative cooling and estimate the bulk velocity of the wind at each radial shell in the disc from
$$\frac{1}{2}\dot{M}_{out}^\mathrm{\Omega }v_w^2=\dot{E}_{in}^\mathrm{\Omega }+\dot{E}_{\mathrm{SNII}}^\mathrm{\Omega },$$
(45a)
and to close the equations we assume that
$$\dot{M}_{out}^\mathrm{\Omega }=\dot{M}_{in}^\mathrm{\Omega }=\dot{M}_{ev}^\mathrm{\Omega }.$$
(45b)
(Note that the numerical coefficient in equation (43a) has been adjusted to give $`v_w=\sqrt{2.5}c_i`$ if $`\dot{E}_{\mathrm{SNII}}=0`$ so that the criterion for the wind to escape, $`v_w>v_{esc}`$, is the same as in the preceeding Section).
Energy input from Type II supernovae exploding high above the cold gas layer will make a small contribution to the thermal energy of the hot coronal gas. For values of $`ϵ_{SNII}1`$, $`\dot{E}_{SNII}^\mathrm{\Omega }`$ will be about $`20\%`$ or so of $`\dot{E}_{in}`$ and cannot be higher because $`\dot{M}_{ev}^\mathrm{\Omega }`$ and the star formation rate are nearly proportional to each other (equation 27). Type Ia supernovae will also supply energy to the corona, with a time lag of perhaps $`\stackrel{>}{}1`$ Gyr (Madau, Della Valle and Panagia 1998). However, this effect will also make a small perturbation to the energy budget of the corona and so it is ignored here. Furthermore, in the models described here much of the gas is expelled on a timescale of $`\stackrel{<}{}1`$ Gyr, thus feedback is likely to be more or less complete before energy injection from Type Ia supernovae becomes significant.
### 5.4 Pressure equilibrium and cold cloud radii
In the models of Section 3 and 4, the cold cloud radii were kept constant irrespective of the pressure of the confining hot phase. More realistically, the cold cloud radii will adjust to maintain approximate pressure equilibrium with the hot phase. Thus
$$\frac{a}{a_{}}0.53\left(\frac{T_c}{80K}\right)^{1/3}(n_{h2}T_{h6})^{1/3},$$
(46)
where $`a_{}`$ denotes the cloud radii at the solar neighbourhood with values as given in equation (2.5) and $`T_c`$ is the internal temperature of the cold clouds. In our own Galaxy, photoelectric heating of dust grains is believed to be the main heating mechanism of the cold clouds (see e.g. Wolfire et al. 1995) but other heating mechanisms may be important, for example, cosmic-ray heating (Field, Goldsmith and Habing 1969). We therefore expect that $`T_c`$ varies in a complex (and uncomputable) way as a galaxy evolves. To assess the effects of the pressure response of the cold clouds, $`T_c`$ will be assumed to remain constant at $`80`$K. The cloud radii are then determined solely by the pressure of the hot component via equation (46).
The energy lost through cloud collisions (equation 2.5) varies as $`a^2`$ (for fixed cloud masses). However, the cloud heating efficiency factor $`ϵ_c`$ will also change as the cloud radii change in response to the pressure of the hot phase. In the model of MO77, the energy acquired from momentum exchange with cooling supernovae shells varies as $`a^4`$. The net effect of these variations in the self-regulating star formation model is to introduce positive feedback, since a higher rate of star formation is required to balance energy lost through cloud collisions in regions where the ambient pressure is higher. This is modelled by assuming $`\dot{E}_{coll}(a/a_{})^2`$ and $`ϵ_c=ϵ_c(a/a_{})^4`$, where $`ϵ_c`$ is a fiducial efficiency factor.
### 5.5 Induced star formation
The maximum mass for an isothermal cloud in pressure equilibrium with the confining medium of pressure $`p_h`$ is given by the Bonner-Ebert criterion (Bonner 1956, Ebert 1955, Spitzer 1968),
$`m_{\mathrm{BE}}=1.18\left({\displaystyle \frac{kT_c}{\mu _p}}\right)^2G^{3/2}p_h^{1/2}`$
$`=433\left({\displaystyle \frac{T_c}{80K}}\right)^2(n_{h2}T_{h6})^{1/2}M_{}.`$ (47)
For our own Galaxy (MO77), $`n_h1.5\times 10^3\mathrm{cm}^3`$, $`T_h4\times 10^5`$ and $`p_h3\times 10^{12}\mathrm{dyne}\mathrm{cm}^2`$, hence $`m_{\mathrm{BE}}1700M_{}`$. This is reasonably close to the upper mass limit, $`m_u=4300M_{}`$, for the cold cloud mass spectrum adopted in this paper ($`a_u=10\mathrm{p}\mathrm{c}`$ with $`\rho _c=7\times 10^{23}\mathrm{g}/\mathrm{cm}^3`$). Gravitational stability requires $`m_um_{\mathrm{BE}}`$ and we will henceforth impose this condition in determining the upper mass limit of the cold cloud spectrum. An increase in the pressure of the hot phase will lead to a decrease in $`m_u`$ and hence to some pressure induced star formation. If the over-pressured clouds fragment into stars with an efficiency $`ϵ_{\mathrm{BE}}`$, the induced star formation rate is given by
$$\frac{dM_s^\mathrm{\Omega }}{dt}=ϵ_{\mathrm{BE}}\frac{M_g^\mathrm{\Omega }}{2p_h\mathrm{ln}(m_u/m_L)}\{\begin{array}{cc}dp_h/dt\hfill & dp_h/dt>0\hfill \\ 0\hfill & dp_h/dt0\hfill \end{array},$$
(48)
where $`m_L`$ is the lower limit to the cloud mass spectrum, $`m_L0.5M_{}`$. This provides an additional source of positive feedback, since as the pressure of the hot component rises the star formation rate of the self-regulating model is enhanced by pressure induced star formation.
### 5.6 Chemical evolution
It is straightforward to include chemical evolution in the models using the instantaneous recycling approximation. We distinguish between ‘primordial’ infalling gas accreting at a rate $`d\mu _I/dt`$ with metallicity $`Z_I`$, and processed gas from the galactic fountain of metallicity $`Z_F`$ accreted at a rate $`d\mu _F/dt`$. The equation of chemical evolution is then
$$\mu _gdZ=pd\mu _s+(Z_IZ)d\mu _I+(Z_FZ)d\mu _F,$$
(49)
(see e.g. Pagel 1997), where $`p`$ is the yield. We adopt a yield of $`p=0.02`$ and assume that the primordial gas has zero metallicity ($`Z_I=0`$). Gas ejected in a galactic fountain is assumed to have the same metallicity as the ISM at the time that the gas was ejected. Within the disc, the ISM gas is assumed to be perfectly mixed at all times. We normalize the metallicities to the solar value, for which we adopt $`Z_{}=0.02`$.
## 6 Results and Discussion
### 6.1 Variation of input parameters
In addition to the many simplifying assumptions introduced in previous sections, the model described here has $`4`$ key parameters: (i) $`\varphi _\kappa `$, determining the efficiency of heat conduction (equation 20); (ii) $`ϵ_c`$, controlling the star formation rate (equation 19); (iii) $`ϵ_{\mathrm{SNII}}`$, determining the efficiency with which energy from Type II supernovae couples directly to the gas (equation 45a) ; (iv) $`ϵ_{\mathrm{BE}}`$, setting the efficiency with which over-pressured ISM clouds collapse to make stars (equation 48). In addition, the ISM cloud radii can be allowed to vary in response to the pressure of the ISM as described in Section 5.4.
The effects of varying these parameters are summarised in Table 2. Here we have run six models of galaxies MW and DW varying the input parameters. We list the final stellar mass $`M_s`$, gaseous disc mass $`M_g`$, and ejected mass $`M_{ej}`$ after $`10`$ Gyr for model DW (there is very little evolution after this time) and after $`15`$ Gyr for model MW. The parameters $`f_{ej}`$ and $`f_{}`$ are the final ejected and stellar masses divided by the total baryonic mass ($`M_{ej}+M_{}+M_g`$). $`\tau _{ej}`$ is the time when half the final ejected mass is lost. The last three numbers list the final mean metallicities of the cold ISM, the stars and the ejected gas.
The most important result from this table is that the final parameters of the models are remarkably insensitive to variations of the input parameters. For models DW, the final stellar disc mass varies between $`4\times 10^7`$ and $`7\times 10^7M_{}`$ and the gas ejection fraction varies from $`0.59`$ to $`0.82`$. For models MW, the final stellar disc mass varies between $`2.3\times 10^{10}`$ and $`2.8\times 10^{10}M_{}`$ and the gas ejection fraction varies from $`0.12`$ to $`0.22`$. Figure 10 shows the evolution of the radial density profiles of models MW1 and DW1 and figure 11 shows the time evolution of the star formation rates, gas fractions and gas velocity dispersions. The models of Table 2 behave in similar ways, and so these two figures are representative of the behaviour of all of the models. These figures are qualitatively similar to those of the simple model of Section 4 (figure 8 and 9). The main differences are:
(a) The gas discs have a sharper outer edge. This is a consequence of the infall model; the outer edge is determined by the final time of the model which sets the maximum cooling radius within the halo (cf. figures 7).
(b) The radial profiles of models MW show oscillatory behaviour near their centres, and the star formation rates and gas fractions show oscillatory behaviour as a function of time. Both of these effects are a consequence of the galactic fountain.
In these models, the star formation rate begins to rise as the gas disc builds up from infalling gas. As the star formation rate rises, the cold ISM is converted efficiently into a hot phase and this is either driven out of the halo or becomes part of the galactic fountain. In models MW, most of the gas that escapes from the system is lost within this early ($`\stackrel{<}{}0.2`$ Gyr) period of star formation when the net star formation rate is close to its peak of $`10M_{}`$/yr. After about $`0.2`$ Gyr, the temperature of the hot phase in models MW settles to $`10^6`$ K very nearly independent of radius (cf figure 10), and so the galactic fountain cycles on a characteristic time-scale of $`4\times 10^8`$ yr. Models DW behave in much the same way as the simpler models of Section 4.3, except that infall, by construction, extends over a longer period of time. In these models, the escape criterion for the wind is satistfied over most of the lifetime of the disc and hence the model of a galactic fountain is unimportant.
We discuss briefly the effects of varying the input parameters:
$`\varphi _k`$: The evaporation rate $`\dot{M}_{ev}`$ has a weak dependence on the evaporation parameter $`\varphi _k`$ ($`\varphi _k^{0.29}`$, equation 28) and obviously decreases as $`\varphi _k`$ is reduced. However, the temperature of the hot component is proportional to $`\varphi _k^{0.29}`$ and hence rises as $`\varphi _k`$ is reduced. The net effect is that the mass of gas ejected is relatively insensitive to $`\varphi _k`$, but the mass of the final stellar disc increases as $`\varphi _k`$ is reduced.
$`ϵ_c`$: Increasing this parameter reduces the star formation rate in the self-regulating model for a fixed gas surface density and velocity dispersion (equations 15 and 18). However, a lower past star formation rate leads to a higher gas surface density which increases the star formation rate. These effects tend to cancel and so the models are insensitive to variations in $`ϵ_c`$.
$`ϵ_{\mathrm{SNII}}`$: Setting this parameter to unity increases the temperature of the hot component slightly and hence increases the efficiency of feedback. As explained in Section 5.3, energy injection by Type II supernovae at large vertical scale heights will always be small compared to the internal energy of the hot phase.
$`ϵ_{BE}`$: Values of $`ϵ_{\mathrm{BE}}0.05`$ have little effect on the evolution. Provided that $`ϵ_{\mathrm{BE}}`$ is not too large (so that it does not dominate the net star formation rate), pressure enhanced star formation is self-limiting because it increases the velocity dispersion of the cold clouds (reducing the star formation rate in the self-regulating model) and converts cold gas to hot gas. Both of these effects tend to reduce the net star formation rate.
Pressure response of cold cloud radii: Allowing the cold gas radii to respond to the pressure of the hot phase provides strong positive feedback in the very early stages of galaxy formation when the pressure of the hot phase is high. However, most of the cold ISM is ejected on a much longer timescale (cf. values of $`\tau _{ej}`$ in Table 2) when the typical pressure of the ISM is similar to that in our own Galaxy. The pressure response of the cold cloud radii therefore has little effect on the final feedback efficiency.
The models described here involve a complex set of coupled equations and a number of parameters. However, one of the most interesting aspects of this study is that the equations interact in such a way that the evolution of the models is insensitive to the parameters. Most importantly, the efficiency of feedback is insensitive to the thermal conduction parameter $`\varphi _\kappa `$. The possible severe suppression of thermal conduction by tangled magnetic fields in astrophysical environments is a long standing theoretical problem. However, our results show that even a reduction of $`\kappa `$ by a factor of $`100`$ or more will not significantly alter the efficiency of feedback.
### 6.2 Chemical evolution
In this section, we summarize some of the results relating to chemical evolution in these models. Our intention is not to present a detailed model of chemical evolution in disc systems along the lines of, for example, Lacey and Fall (1983, 1985) but to investigate some of the general features of chemical evolution with physically motivated models of inflow and outflow. The chemical evolution model is based on the instantaneous recycling approximation as described in Section 5.6. This is probably a reasonable approximation since the timescales of star formation and outflow are $`1`$Gyr, but will overestimate the gas metallicities where the gas density is low. As in the previous Section, results from models MW1 and DW1 are used to illustrate the general features of the models. The other models listed in Table 2 behave in very similar ways.
#### 6.2.1 Stellar metallicity distribution
The final mean stellar metallicities are typically $`Z_s/Z_{}0.5`$ in models MW and $`0.2`$ in models DW. Models DW have a lower stellar metallicity because a larger fraction of the ISM is expelled in a wind. The stellar metallicity distributions are shown in Figures (12) and (13). Figure (12c) is particularly interesting because this radius is close to the solar radius. This metallicity distribution is quite similar to that of G-dwarfs in the solar cylinder (see e.g. figure 8.19 of Pagel 1997), showing that the infall model solves the ‘G-dwarf’ problem of closed box models of chemical evolution. The metallicity distributions of model DW1 plotted in figure (13) also show a lack of stars with low metallicities.
#### 6.2.2 Metallicity gradients
Over most of the stellar disc, model MW1 has a fairly weak stellar metallicity gradient (figure 14a) except at the very outer edge where the stellar density and metallicity fall abrubtly. This differs from the metallicity gradients seen in large disc systems which show linear gradients (see Vila-Costas 1998 for a recent review). It is possible that this problem might be resolved by including radial gas flows in the models (Lacey and Fall 1985, Pitts and Tayler 1985). The stellar metallicity gradients in model DW1 are steeper, in qualitative with observations which indicate that the abundance gradients in Scd and Irr galaxies are steeper than those in earlier type galaxies.
The radial gas metallicity profiles are shown in figure (14b). Model DW contains a gasesous disc extending well beyond the edge of the stellar disc. This gas disc has a low metallicity in the outer parts, with $`Z/Z_{}\stackrel{<}{}10^2`$ at $`r\stackrel{>}{}2`$kpc. At these large radii, the star formation rate is always low and the gas disc can survive for much longer than a Hubble time without converting into stars. This is unlikely to happen in all galaxies for at least two reasons: (i) the energy injection from supernovae into this gas will not be uniform as assumed in this paper; (ii) the extended gas disc is susceptible to external disturbances and so could be tidally stripped or transported towards the centre of the system in a tidal interaction. Nevertheless, it is possible that dwarf galaxies at high redshift possess extended gaseous discs, some of which survive to the present day.
#### 6.2.3 Effective yields
According to the simple closed box model of chemical evolution, the metallicity of the ISM is related to the gas fraction according to
$$Z_g=p\mathrm{ln}(M_g/(M_g+M_s)).$$
(50)
It is well known that the yields derived from applying this relation to gas rich galaxies (usually dwarf systems) result in “effective yields”, $`p_{\mathrm{eff}}`$, that are much lower than the yield expected from a standard Salpeter-like IMF. For example, Vila-Costas and Edmunds (1992) find effective yields in the range $`p_{\mathrm{eff}}0.004`$$`0.02`$ and that the effective yield decreases with increasing radius.
The solid line in figure (15) shows the final gas metallicity in radial rings in model DW1 plotted against the gas fraction within each ring. The dashed line shows equation (50) with an effective yield of $`0.004`$ (i.e. one-fifth of the true yield). The strong outflows in this model suppress the effective yield well below the true yield and produce a strong radial variation of the effective yield, in qualitative agreement with observations.
#### 6.2.4 Metallicity of ejected gas
The last line in Table 2 lists the mean metallicity of the gas that escapes from the galaxy. The mean metallicity of the ejected gas is about $`0.3Z_{}`$ for model MW1 and about $`0.1Z_{}`$ for model DW1. In model DW1 this value is about $`3`$ to $`5`$ times higher than the mean metallicity of the final gas disc. The ejected gas in this model is therefore ‘metal enhanced’ relative to the gaseous disc. The mechanism for this metal enhancement is physically different to that in the models of Vader (1986, 1987) and Mac Low and Ferrara (1999). In the models of these authors, metal enhancement arises from incomplete local mixing between the supernovae ejecta and the ISM. In our models, the gas is assumed to be well mixed locally, but metal enhancement arises because the gas is lost preferentially from the central part of the galaxy, which has a higher metallicity than the gas in the outer parts of the system.
### 6.3 Connection with damped Lyman alpha systems
The column density threshold for the identification of damped Ly$`\alpha `$ systems is $`N(\mathrm{HI})\stackrel{>}{}2\times 10^{20}\mathrm{cm}^2`$ (Wolfe 1995) corresponding to a neutral gas surface mass density of $`1.6M_{}/\mathrm{pc}^2`$. Comparison with Figure 10 shows that the extended cold gasesous discs around dwarf galaxies would be detectable as damped Ly$`\alpha `$ systems. Furthermore, in CDM-like models, such extended discs around dwarf galaxies would dominate the cross-section for the identification of damped Ly$`\alpha `$ systems at high redshift because the space density of haloes with low circular speeds is high (Kauffman and Charlot 1994, Mo and Miralda-Escude 1994). If this is the case, then the metallicities of damped Ly$`\alpha `$ systems would be expected to be low at high redshift, $`Z/Z_{}\mathrm{few}\times 10^2`$, with occasional lines-of-sight intersecting the central regions of galaxies where the metallicity rises to $`Z/Z_{}\stackrel{>}{}0.1`$. At lower redshifts, the metallicities of damped systems would be expected to show a similarly large scatter, but with perhaps a trend for the mean metallicity to increase as disc systems with higher circular speeds form and the extended gaseous discs around dwarfs are disrupted by tidal encounters.
This is qualitatively in accord with what is observed (Pettini et al. 1997, Pettini et al. 1999, Pettini 1999). These authors find that the typical metallicity of a damped Ly$`\alpha `$ system at $`z2`$$`3`$ is about $`0.08Z_{}`$ with a spread of about two orders of magnitude. Comparing the metallicities of high redshift systems with those of $`10`$ damped Ly$`\alpha `$ systems with redshifts $`z=0.4`$$`1.5`$, Pettini et al. (1999) find no evidence for evolution of the column density weighted metallicity. Whether these and other properties of the damped Ly$`\alpha `$ systems can be reproduced with the feedback model described here requires more detailed ‘semi-analytic’ calculations along the lines described by Kauffmann (1996). However, the key point that we wish to emphasise is that according to the models described here, most of the cross-section at any given redshift will be dominated by largely unprocessed gas in the outer parts of galaxies that does not participate in the star formation process. The metallicity distributions and the evolution of $`\mathrm{\Omega }_{\mathrm{HI}}`$ as a function of redshift are therefore more likely to tell us about feedback processes and the outer parts of dwarf galaxies than about the history of star formation. Attempts to use the properties of damped Ly$`\alpha `$ systems to constrain the cosmic star formation history (e.g. Pei, Fall and Hauser 1999) should therefore be viewed with caution.
### 6.4 Feedback efficiency as a function of halo circular speed and semi-analytic models of galaxy formation
In this section we investigate the efficiency of feedback as a function of halo circular speed. We have adopted the parameters of models 4 in Table 1 and run a series of models varying the halo circular speed $`v_v`$. The virial radius of the halo is set to $`r_v=150(v_v/126\mathrm{kms}^1)^2\mathrm{kpc}`$, the concentration parameter $`c=10`$ and the ratio of gas to halo mass within the virial radius is set to $`1/10`$. The halo rotation speed is set by equation (39) with the fiducial disc scale length equal to $`r_v/50`$. With these parameters, the family of models has a constant value for the halo spin parameter of $`\lambda _H=0.065`$.
The retained baryonic fraction, $`1f_{ej}`$, is plotted as a function of halo circular speed in Figure 16. The dotted line shows the relation used by Cole et al. (1994, hereafter C94) in their semi-analytic models,
$$1f_{hot}=\frac{1}{1+\beta (v_v)},\beta (v_v)=\left(\frac{v_v}{v_{hot}}\right)^{\alpha _{hot}},$$
(51)
where $`f_{hot}`$ is the fraction of the cooled gas that is reheated and $`v_{hot}`$ and $`\alpha _{hot}`$ are parameters. C94 adopt a severe feedback prescription with $`\alpha _{hot}=5.5`$ and $`v_{hot}=140\mathrm{kms}^1`$ to reproduce the flat faint end slope of the B-band galaxy luminosity function in a critical density CDM model. The C94 feedback model does not agree at all well with the models described here. There is a slight ambiguity in the appropriate value of $`v_v`$ to use in equation (51) because C94 adopt an isothermal rather than an NFW halo profile; the halo circular speed at $`0.1r_v`$ may be $`20\%`$ higher than the circular speed at the virial radius, but this is far too small a difference to reconcile the C94 feedback prescription with the models of this paper.
In fact, the dashed line in Figure 16 shows that our models are reasonably well described by equation (51) with $`v_{hot}=75\mathrm{kms}^1`$ and $`\alpha _{hot}=2.5`$. Our results therefore suggest a much gentler feedback prescription than assumed in C94. Note that with the C94 parameters, a Milky Way type galaxy with $`v_v130\mathrm{kms}^1`$ would have lost about $`60\%`$ of its baryonic mass in a wind. This is well outside the range found from our models for plausible choices of the input parameters (cf. Table 2).
Recently Baugh et al. (1999) and Cole et al. (1999) describe semi-analytic models applied to $`\mathrm{\Lambda }`$-dominated CDM cosmologies that employ a gentler feedback model. The prescription for their reference model is similar to that of equation (51) with $`v_{hot}=150\mathrm{kms}^1`$ and $`\alpha _{hot}=2.0`$, but with $`v_v`$ replaced by the disc circular speed $`v_{\mathrm{disc}}`$. This model is closer to the results found here. Assuming angular momentum conservation, a halo with $`\lambda _H0.06`$ will produce a disc with a circular speed $`v_{\mathrm{disc}}1.7v_v`$ (cf. Table 1) and so their model can be approximated by equation (51) with $`v_{hot}90\mathrm{kms}^1`$ and $`\alpha _{hot}=2.0`$. With these parameters, their model gives somewhat stronger feedback than found in our models, but is well within the range of physical uncertainties. Kauffmann et al. (1993) and Kauffmann, Guiderdoni and White (1994) also adopt a much less severe feedback prescription than C94 in their semi-analytic models. For a detailed analysis of the effects of varying the feedback prescription (and other parameters) in semi-analytic models see Somerville and Primack (1999).
The change from an Einstein-de Sitter CDM cosmology in C94 to a $`\mathrm{\Lambda }`$-dominated CDM model in Cole et al. (1999) partly explains why the revised models provide a reasonable match to observations using less efficient feedback. However, the revised models predict a faint end slope for the B-band luminosity function that is consistent with the observations of Zucca et al. (1997) but not with those of other authors (e.g. Loveday et al. 1992, Maddox et al. 1998). (The earlier paper of Cole et al. 1994 attempted to reproduce the flat faint end slope of the Loveday et al. luminosity function). The observational differences in estimates of the faint end slope of the optical luminosity function are not properly understood and so it remains unclear whether a gentle feedback model, of the type proposed here and used in Cole et al. (1999), can account for galaxy formation in CDM-type models.
Table 3: Dependence of Feedback Efficiency
of Model DW on Halo Angular Momentum
| $`f_{coll}`$ | $`v_c`$ (km/s) | $`\lambda _H`$ | $`f_{ej}`$ |
| --- | --- | --- | --- |
| $`\mathrm{\hspace{0.33em}\hspace{0.33em}25}`$ | $`\mathrm{\hspace{0.33em}\hspace{0.33em}50}`$ | $`0.12`$ | $`0.64`$ |
| $`\mathrm{\hspace{0.33em}\hspace{0.33em}50}`$ | $`\mathrm{\hspace{0.33em}\hspace{0.33em}70}`$ | $`0.065`$ | $`0.59`$ |
| $`150`$ | $`120`$ | $`0.031`$ | $`0.82`$ |
With the Cole et al. (1999) parameterization the efficiency of feedback depends, by construction, on the surface density of the galaxy and hence on the angular momentum of the parent halo. In their model, higher angular momentum haloes lead to more efficient feedback because they form low surface density discs with low disc circular speeds. This is not what is found in our models. Table 3 lists the ejected gas fraction as a function of the halo spin parameter $`\lambda _H`$. Here, the halo circular speed and virial radius, $`v_v`$ and $`r_v`$, are the same as for model DW in Table 1, but the amplitude of the halo rotation speed (or equivalently the parameter $`f_{coll}`$) is adjusted to change the spin parameter of the halo. The parameters of the feedback model are the same as those for model DW1 in Table 2. The feedback efficiency depends weakly (and non-monatonically) on $`\lambda _H`$, and is greater in systems with low values of $`\lambda _H`$. This is because higher surface densities in low $`\lambda `$ galaxies result in higher star formation rates and a higher temperature hot component that can escape more easily from the halo.
The timescale for feedback in C94 and Cole et al. 1999 is closely linked to the star formation timescale which is assumed to be shorter in galaxies with high circular speeds. This is not what is found in the models of Table 2. The timescale for star formation is of order several Gyr in models MW (which have a roughly constant star formation rate at late times, see Figure 11), yet the ejection of hot gas occurs only in the initial stages of formation with a characteristic timescale of $`0.3`$ Gyr. In models DW, the situation is reversed with star formation occuring on a somewhat shorter timescale than that for outflow.
## 7 Conclusions
The main aim of this paper has been to show that supernovae driven feedback can operate in a quiescent mode and that high rates of star formation are not required to drive efficient feedback. In dwarf galaxies feedback occurs on an infalling timescale and so can extend over a period of $`1`$ Gyr. In the feedback model developed here, cold gas clouds are steadily evaporated in expanding supernovae remnants and converted into a hot component. Critically, the rate at which cold gas is evaporated can exceed the rate at which mass is converted into stars. If the temperature of the hot component is high enough, a wind will form and the hot gas can escape from the halo (provided the interaction with infalling gas can be ignored). If the temperature of the hot component is not high enough for it to escape from the halo, it will cool and fall back down to the disc in a galactic fountain. Some characteristic features of the models are as follows:
(i) In a Milky Way type system, feedback from supernovae may drive out some of the gas from the halo in the early phases of evolution ($`t\stackrel{<}{}0.3`$ Gyr) when the star formation rate is high and the temperature of the hot phase exceeds $`5\times 10^6`$ K. For plausible sets of parameters, perhaps $`20`$$`30\%`$ of the final stellar mass might escape from the galaxy. At later times, the temperature of the hot phase drops to $`T10^6`$ K and the evaporated gas cycles within the halo in a galactic fountain.
(ii) In a dwarf galaxy with a circular speed $`50\mathrm{kms}^1`$, expanding supernovae remnants can convert the cold interstellar medium efficiently into a hot component with a chacteristic temperature of a few times $`10^5`$ K. This evaporated gas can escape from the halo in a cool wind. Typically, only about $`10\%`$ of the baryonic material forms stars. Gas accreted from the halo at $`\stackrel{>}{}1`$ Gyr forms an extended gaseous disc which, according to the self-regulating star formation model used here, can survive for longer than a Hubble time without converting into stars.
(iii) The feedback model developed here is meant to provide a sketch of how feedback might operate in a multi-phase interstellar medium. The model contains a number of obvious over-simplifications. For example, we have neglected any interaction of the outflowing gas with the infalling medium, we have not addressed the origin of the cold cloud spectrum, ignored the dense molecular cloud component of the ISM and neglected any local dissipation of supernovae energy in star forming regions. These effects, and other processes, are undoubtedly important in determining the efficiency of feedback. Nevertheless, the simplified model presented here contains some interesting features. Firstly, the model shows how positive feedback (via pressure induced star formation) and negative feedback (via outflowing gas) can occur simultaneously. Secondly, the models are remarkably insensitive to uncertain physical parameters, in particular, thermal conduction would need to be suppressed relative to its ideal value by factors of much more than $`100`$ to qualitatively change the model. If thermal conduction is highly suppressed, it may be possible to construct a qualitatively similar model to the one presented here in which cold gas is converted into hot gas in shocks.
(iv) The self-regulated star formation and feedback models described here provide physically based models for the star formation timescale and feedback efficiency as a function of the parameters of the halo. The star formation timescale and feedback efficiency (or timescale) are taken as free parameters in semi-analytic models of galaxy formation (e.g. Cole et al. 1999, Kauffmann et al. 1994) and are critically important in determining some of the key predicted properties of these models, for example, the faint end slope of the galaxy luminosity function and the star formation history at high redshifts (see e.g. Somerville and Primack 1999). It is therefore important that we develop a theoretical understanding of these parameters (as attempted crudely here) and also that ways are found to constrain these parameters observationally. The results presented here show that supernovae feedback is much less effective than assumed in some earlier semi-analytic models (Cole et al. 1994, Baugh et al. 1996) but is closer to the more gentler feedback prescriptions used in more recent models (Cole et al. 1999, Somerville and Primack 1999).
The feedback model described in this paper has a number of consequences and raises some problems which are summarized below.
(i) Evidence for outflows: According to the models described here, outflows with speeds of $`200T_{h6}^{1/2}\mathrm{kms}^1`$ should be common in high redshift galaxies. There is evidence for an outflow of $`200\mathrm{kms}^1`$, a mass loss rate of $`60M_{}/\mathrm{yr}`$ and a star formation rate of $`40M_{}/\mathrm{yr}`$ in the gravitationally lensed Lyman break galaxy MS1512-cB58 (Pettini et al. 2000). The outflow velocity in this galaxy is consistent with our models, but the star formation and mass loss rates (which are highly uncertain) are high. The most likely explanations are either that MS1512-cB58 is a massive galaxy driving an outflow that will remain bound to the system, or that it is a less massive system undergoing a burst of star formation. In addition to direct detection of outflowing gas, winds may have other observational consequences. The winds from dwarf galaxies will cool rapidly (see Appendix B). Wang (1995b) has suggested that photoionized gas clouds formed in the cooling wind might contribute to the Ly$`\alpha `$ forest. Nulsen, Barcon and Fabian (1998) suggest that outflows caused by bursts of star formation in dwarf galaxies might even produce damped Ly$`\alpha `$ systems.
(ii) Damped Ly$`\alpha `$ systems: The extended gaseous discs that form around dwarf galaxies in our models have low metallicities because they have low rates of star formation. If this is correct, then this largely unprocessed gas would dominate the cross section for the formation of damped Ly$`\alpha `$ absorbers. The metallicities of most of these systems would be low, but would show a large scatter because some lines of sight will pass close to the central regions of galaxies containing gas of high metallicity. This is broadly in agreement with what is observed. Extended gaseous discs would be vulnerable in tidal interactions. Some of the gas might be stripped and some might be transported into the central regions to be converted into stars and hot gas. The evolution of $`\mathrm{\Omega }_{\mathrm{HI}}`$ determined from damped Ly$`\alpha `$ systems (e.g. Storrie-Lombardi, McMahon and Irwin 1996) might have more to do with infall, feedback and tidal disruption than with the cosmic star formation history.
(iii) Angular momentum conservation: In hydrodynamic simulations, gas is found to cool effectively in sub-units during the formation of a protogalaxy. These sub-units lose their orbital angular momentum to the halo as they spiral towards the centre and merge. Hence the gas does not conserve angular momentum during the formation of a massive galaxy. (Navarro and Benz 1991, Navarro and Steinmetz 1997, Weil Eke and Efstathiou 1998, Navarro and Steinmetz 2000). In fact, in the absence of feedback it has proved impossible to form discs with angular momenta similar to those of real disc galaxies starting from CDM initial conditions. In the models described here, it has been assumed for simplicity that the specific angular momentum of the gas is conserved during collapse. This assumption could easily be relaxed. However, the feedback model decribed here suggests that the numerical simulations miss some important physics. Firstly, it is preferentially the low angular momentum gas, infalling in the early stages of evolution, that is most likely to be ejected from a developing protogalaxy. Secondly, supernovae driven feedback may help to solve the angular momentum momentum problem by ejecting gas efficiently from sub-units. The ejected gas may then infall at later times when the halo is less sub-structured, approximately conserving its angular momentum (Weil et al.1998, Eke, Efstathiou and Wright 2000).
(iv) Implementing feedback in numerical simulations: There have been a number of attempts to implement supernovae feedback in gas dynamical numerical simulations (e.g. Katz 1992, Navarro and White 1993, Navarro and Steinmetz 2000). These involve either heating the gas around star forming regions (which is ineffective because the energy is quickly radiated away) or reversing the flow of infalling gas. An implementation of the feedback model described here is well beyond the capabilities of present numerical codes. It would require modelling several gas phases, a cold interstellar medium, a hot outflowing medium and an infalling component, including mass transfer between each phase. It might be worth attempting simpler simulations in which cold high density gas is added to the halo beyond the virial radius at a rate that is determined by the local star formation rate.
(vi) Starbursts vs quiescent feedback: It is likely that starburts are more common at high redshift because of the increased frequency of galaxy interactions. Starbursts could contribute to supernovae driven feedback in addition to the quiescent mode described here. However, at any one time, our models suggest that the cold gas component will have a mass of only $`20`$$`50\%`$ of the the mass of the stellar disk. Even if a substantial fraction of this gas is transported towards the centre of a galaxy in a tidal encounter (see e.g. Barnes and Hernquist 1996) and is subsequently ejected in a superwind, this mode of feedback will be inefficient because the mass of gas involved is a small fraction of the total gas mass ejected in the quiescent feedback mode over the lifetime of the galaxy.
(vii) Metallicity ejection: The mean metallicity of the gas ejected from a dwarf galaxy is typically about $`Z_{}/10`$ in our models, and comparable to the mean metallicity of the stars in the final galaxy. Yet typically a dwarf galaxy is predicted to expel $`5`$ to $`10`$ times its residual mass in stars. Dwarf galaxies can therefore pollute the IGM with metals to a much higher level than might be inferred from their stellar content. The high metallicity of gas in the central regions of clusters $`Z_{}/3`$ (e.g. Mushotsky and Loewenstein 1997) may require a top-heavy IMF and gas ejection from massive galaxies.
Aknowledgments: The author aknowledges the award of a PPARC Senior Fellowship.
## Appendix A Fast Computation of the Rotation Curve of a Thin Disc
We begin with the expression for the potential of a thin disc at $`z=0`$
$$\varphi (r,z=0)=2\pi G_0^{\mathrm{}}_0^{\mathrm{}}J_0(kr)J_0(kr^{})r^{}\mu (r^{})𝑑r^{}𝑑k.$$
(52)
The integral over $`k`$ is a well-known discontinuous integral (e.g. Watson 1944)
$$_0^{\mathrm{}}J_0(kr)J_0(kr^{})𝑑k=\frac{2}{\pi }\{\begin{array}{cc}(1/r)K(r^{}/r)\hfill & r^{}<r\hfill \\ (1/r^{})K(r/r^{})\hfill & r^{}>r\hfill \end{array}$$
(53)
where $`K`$ is the complete elliptic integral of the first kind. Differentiating equation (53), we find
$`v^2(r)=4Gr\{{\displaystyle _0^{rϵ}}I_<(r,r^{})r^{}\mu (r^{})dr^{}`$
$`+{\displaystyle _{r+ϵ}^{\mathrm{}}}I_>(r,r^{})r^{}\mu (r^{})dr^{}\},`$ (A54a)
$`I_<(r,r^{})={\displaystyle \frac{E(r^{}/r)}{(r^2r_{}^{}{}_{}{}^{2})}},`$ (A54b)
$`I_>(r,r^{})={\displaystyle \frac{K(r/r^{})}{rr^{}}}{\displaystyle \frac{r^{}E(r/r^{})}{r(r_{}^{}{}_{}{}^{2}r^2)}},`$ (A54c)
where $`E`$ is the complete elliptic integral of the second kind. This integral is convergent in the limit $`ϵ0`$. The functions $`I_<`$ and $`I_>`$ can be evaluated once and stored, reducing the computation of $`v^2(r)`$ to a simple integral over the surface density of the disc multiplied by the tabulated functions. We evaluate the epicyclic frequency by differentiating equation (A3) numerically.
## Appendix B Steady Spherical Winds
The equations governing a steady spherically symmetric wind are
$`{\displaystyle \frac{1}{r^2}}{\displaystyle \frac{d}{dr}}(\rho vr^2)=q(r),`$ (B55a)
$`\rho v{\displaystyle \frac{dv}{dr}}={\displaystyle \frac{dp}{dr}}\rho {\displaystyle \frac{d\mathrm{\Phi }}{dr}}q(r)v,`$ (B55b)
$`{\displaystyle \frac{1}{r^2}}{\displaystyle \frac{d}{dr}}\left[\rho vr^2\left({\displaystyle \frac{1}{2}}v^2+{\displaystyle \frac{5}{2}}{\displaystyle \frac{p}{\rho }}\right)\right]+\rho v{\displaystyle \frac{d\mathrm{\Phi }}{dr}}=𝒞,`$ (B55c)
where $`q(r)`$ is the mass density injected per unit time and $``$ and $`𝒞`$ are the heating and cooling rates per unit volume (e.g. Burke 1968, Holzer and Axford 1970). We assume that the gravitational force is given by the NFW halo potential (equation 3), $`d\mathrm{\Phi }/dr=v_H^2(r)/r`$, and rewrite these equations as two dimensionless first order equations
$`{\displaystyle \frac{dv^2}{dx}}={\displaystyle \frac{1}{2\pi x^2(c^2v^2)}}[8\pi xc^2v^2+4\pi xv^2v_H^2`$
$`+{\displaystyle \frac{4}{3}}\gamma v^2+\gamma c_i^2{\displaystyle \frac{2}{3}}\kappa ]`$ (B56a)
$`{\displaystyle \frac{dc^2}{dx}}={\displaystyle \frac{1}{6\pi x^2(c^2v^2)v^2}}[8\pi xc^2v^4+4\pi xc^2v^2v_H^2`$
$`\gamma (v^2c^2{\displaystyle \frac{3}{2}}c^4{\displaystyle \frac{5}{6}}v^4)`$ (B56b)
$`{\displaystyle \frac{3}{2}}\gamma c_i^2(c^2{\displaystyle \frac{5}{3}}v^2)+\kappa (c^2{\displaystyle \frac{5}{3}}v^2)]`$ (B56c)
where $`c`$ is the adiabatic sound speed, $`x=r/r_v`$, and all velocities are expressed in units of $`v_v`$. The quanities $`\gamma `$ and $`\kappa `$ in these equations are related to the mass injection and cooling rates according to
$`\gamma (x,v)={\displaystyle \frac{q(r)\dot{M}(r)}{\rho ^2r_vv_v^2}},\kappa (x,c)={\displaystyle \frac{\dot{M}(r)\mathrm{\Lambda }(T)n_e^2}{\rho ^2v_v^4r_v}},`$ (B56d)
$`\dot{M}(r)=4\pi {\displaystyle _0^r}q(r)r^2𝑑r.`$ (B56e)
and the injected gas is assumed to have a uniform initial isothermal sound speed of $`c_i=(kT_i/(0.61m_p))^{1/2}`$.
We illustrate the behaviour of the wind solutions by studying two regimes. Firstly, we assume that $`q=0`$ beyond an initial radius $`r_i=0.04r_v`$ defining the base of the flow (i.e. two disc scale lengths for $`f_{coll}=50`$). Equations (B2) do not have a transonic point when $`q=0`$ (Wang 1995a, see also the discussion below) and so we begin the integrations at a Mach number slightly greater than unity with $`c^2=5c_i^2/3`$. We adopt the parameters of models MW and DW given in Table 1 and integrate the equations (B1) adopting $`\dot{M}=10M_{}/\mathrm{yr}`$ for model MW and $`\dot{M}=0.2M_{}/\mathrm{yr}`$ for model DW. These mass injection rates are close to the maximum rates at times $`t\tau _{ej}`$ for the models described in Section 6. The curves in Figure 17 show solutions for initial isothermal sound speeds of $`0.75`$, $`1.0`$ and $`1.25`$ times the escape velocity from the centre of the halo ($`v_{esc}=430\mathrm{kms}^1`$ for model MW and $`107\mathrm{kms}^1`$ for model DW).
The figure shows that the criterion $`v_w\sqrt{2.5}c_i\stackrel{>}{}v_{esc}`$ is about right if the wind is to reach beyond the virial radius. For $`c_iv_{esc}`$ the wind in model MW begins at a high temperature of $`T_i1.4\times 10^7\mathrm{K}`$ and cools almost adiabatically initially, reaching a temperature of $`1.5\times 10^5\mathrm{K}`$ at the virial radius. The timescale for the flow to reach the virial radius, $`2\times 10^8\mathrm{yrs}`$, is slightly longer than the cooling time at $`r_v`$. The behaviour of models DW is quite different. For $`c_iv_{esc}`$ the initial temperature of the gas is $`T_i8\times 10^5\mathrm{K}`$ and cools to $`\stackrel{<}{}10^4\mathrm{K}`$ by $`r=0.3r_v`$. As expected from the discussion in Section 6, cooling is important in outflows from dwarf galaxies (see e.g. Kahn 1981, Wang 1995a, b).
An investigation of transonic solutions of equations (B2) require a model for $`q(r)`$. An example is illustrated in Figure 18 for model DW, using
$$q(r)=\frac{\dot{M}(\mathrm{})}{8\pi r_w^3}\mathrm{exp}(r/r_w),r_w=0.04r_v.$$
(B57)
In this solution, $`\dot{M}(\mathrm{})=0.2M_{}/\mathrm{yr}`$ and the central gas density was adjusted to obtain a critical solution for the case $`c_i=v_{esc}`$. If the gas is to escape from a dwarf galaxy the transonic point must occur before cooling sets in. For such systems, the wind parameters would adjust so that a sonic point exists at a characteristic cooling scale height as shown in Figure 18. The wind will then cool radiatively just beyond the sonic point forming a cold wind as discussed above. It is also likely that the wind will be heated to a temperature of $`T10^4`$K by photoionizing radiation from the galaxy and the general UV background. These sources of heating have not been included in the models of Figures 17 and 18.
The wind will be thermally unstable when cooling sets in, and may form clouds. However, in the absence of a confining medium, the clouds would have a filling factor of order unity so the wind is likely to maintain its integrity until it meets the surrounding IGM. The external pressure required to balance the ram pressure of the wind is
$`{\displaystyle \frac{p_{ext}}{k}}80\left({\displaystyle \frac{\dot{M}}{0.2M_{}/\mathrm{yr}}}\right)\left({\displaystyle \frac{r}{10\mathrm{kpc}}}\right)^2\times `$
$`\left({\displaystyle \frac{v_w}{100\mathrm{k}\mathrm{m}/\mathrm{s}}}\right)\mathrm{cm}^3\mathrm{K},`$ (B58)
which is about equal to the pressure of the IGM with a temperature of $`10^4`$K and an overdensity of
$`\mathrm{\Delta }4500\left({\displaystyle \frac{2}{1+z}}\right)^3T_4^1\left({\displaystyle \frac{\dot{M}}{0.2M_{}/\mathrm{yr}}}\right)\times `$
$`\left({\displaystyle \frac{r}{10\mathrm{kpc}}}\right)^2\left({\displaystyle \frac{v_w}{100\mathrm{k}\mathrm{m}/\mathrm{s}}}\right).`$ (B59)
Provided that the halo is devoid of high pressure gas, the cool wind will propagate beyond the virial radius and will be halted either by the ram pressure of infalling gas or after sweeping up a few times its own mass. As pointed out by Babul and Rees (1992), if a dwarf galaxy is embedded in a group or cluster of galaxies with a pressure exceeding $`100\mathrm{cm}^3\mathrm{K}`$, the bulk motion of the outflowing gas would be thermalized in a shock and the cooled shocked gas could fall back onto the galaxy generating a new burst of star formation. The efficiency of feedback is therefore likely to be a function of local environment.
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# References
PION INTERFEROMETRY FROM A RELATIVISTIC FLUID WITH A FIRST-ORDER PHASE TRANSITION IN CERN-SPS 158 GeV/A Pb+Pb COLLISIONS
K. MORITA, S. MUROYA, H. NAKAMURA and C. NONAKA
Department of Physics, Waseda University
169-8555 Tokyo, Japan
Tokuyama Women’s College
Tokuyama, 745-8511 Yamaguchi, Japan
Department of Physics, Hiroshima University
Higashi-Hiroshima, 739-8526 Hiroshima, Japan
## Abstract
We investigate pion source sizes through the Yano-Koonin-Podgoretskiĭ (YKP) parametrization for the Hanbury-Brown Twiss (HBT) effect in the CERN-SPS 158 GeV/A central collisions. We calculate two-particle correlation functions numerically based on a (3+1)-dimensional relativistic hydrodynamics with a first order phase transition and analyze the pair momentum dependence of the HBT radii extracted from the YKP parametrization in detail. We find that even in the case of a first order phase transition, expansion and the surface dominant freeze-out make the source in the hydrodynamical model opaque significantly. Consequently, the interpretation of the temporal radius parameter as the time duration becomes unavailable for the hydrodynamical model.
Pion interferometry is one of the promising tools to investigate the hot and dense matter created in the relativistic heavy ion collisions . It is well known as Hanbury-Brown Twiss (HBT) effect that we can obtain a size of the source from a two-particle correlation function because of the symmetry of the wave function of the emitted bosons. In the experimental analyses, the size of the source is obtained as a fitting parameter in the Gaussian fit to the correlation function and its physical meaning depends on the Gaussian fitting function. For central collisions, Yano-Koonin-Podgoretskiĭ (YKP) parametrization is used as one of the fitting functions. In the YKP parametrization, three spatial parameters are interpreted directly as the transverse (perpendicular to the collision axis), longitudinal (parallel to the collision axis) and temporal (emission duration) extents, respectively . In the case of relativistic heavy ion collisions, the reaction is highly dynamical. The collective expansion, which contains some informations of an equation of state of the matter, takes place naturally. For example, if there exists a first order phase transition between quark-gluon plasma (QGP) and hadronic gas, enlarged time duration due to the existence of the mixed phase may be observed . However, the collective expansion makes the meaning of the size ambigious. Therefore, we investigate the sizes in the YKP parametrization based on a hydrodynamical model with a first order phase transition .
Assuming a completely chaotic source for simplicity, the two-particle correlation function is given as
$$C(𝐤_\mathrm{𝟏},𝐤_\mathrm{𝟐})=1+\frac{|I(𝐤_\mathrm{𝟏},𝐤_\mathrm{𝟐})|^2}{(dN/d^3𝐤_\mathrm{𝟏})(dN/d^3𝐤_\mathrm{𝟐})},$$
(1)
where $`𝐤_\mathrm{𝟏}`$ and $`𝐤_\mathrm{𝟐}`$ are momenta of a emitted pair and $`I(𝐤_\mathrm{𝟏},𝐤_\mathrm{𝟐})`$ is the interference term.
As usual, we introduce the relative momentum $`q^\mu =k_1^\mu k_2^\mu `$ and the average momentum $`K^\mu =(k_1^\mu +k_2^\mu )/2`$ that satisfy the on-shell condition $`q_\mu K^\mu =0`$. Assuming the cylindrical symmetry, we can take $`K^\mu =(K^0,K_T,0,K_L)`$. Hence, only three components of relative momenta are independent. In the YKP parametrization, $`q_{}=\sqrt{q_x^2+q_y^2}`$, $`q_{}=q_z`$, and $`q^0`$ are taken as the independent components. Then the fitting function is given as
$`C(q^\mu ,K^\mu )`$ $`=`$ $`1+\lambda \mathrm{exp}\{R_{}^2(K^\mu )q_{}^2R_{}^2(K^\mu )[q_{}^2(q^0)^2]`$ (2)
$`[R_0^2(K^\mu )+R_{}^2(K^\mu )]\left[q_\mu u^\mu (K^\mu )\right]^2\},`$
where $`u^\mu `$ is a four-velocity which has only longitudinal component $`v(K^\mu )`$ interpreted as a expansion velocity and the fitting parameters are $`v(K^\mu )`$ and $`R_i`$ $`(i=,,0)`$. Because the three size parameters are invariant under longitudinal boosts, we may use a special frame called YKP frame, where $`v=0`$, in the discussion of the meanings of the size parameters. Introducing the source function as
$$S(x^\mu ,K^\mu )=_\mathrm{\Sigma }\frac{U_\mu (x^{})d\sigma ^\mu (x^{})}{(2\pi )^3}\frac{U_\nu (x^{})k^\nu }{\mathrm{exp}(U_\rho (x^{})k^\rho /T)1}\delta ^4(xx^{}),$$
(3)
and the weighted average
$$A(x^\mu )=\frac{d^4xA(x^\mu )S(x^\mu ,K^\mu )}{d^4xS(x^\mu ,K^\mu )},$$
(4)
the size parameters in the YKP frame are expressed as
$`R_{}^2(K^\mu )`$ $`=`$ $`(\mathrm{\Delta }y)^2,`$ (5)
$`R_0^2(K^\mu )`$ $`=`$ $`(\mathrm{\Delta }t)^2{\displaystyle \frac{2}{\beta _{}}}\stackrel{~}{x}\stackrel{~}{t}+{\displaystyle \frac{1}{\beta _{}^2}}[(\mathrm{\Delta }x)^2(\mathrm{\Delta }y)^2],`$ (6)
$`R_{}^2(K^\mu )`$ $`=`$ $`(\mathrm{\Delta }z)^2{\displaystyle \frac{2\beta _{}}{\beta _{}}}\stackrel{~}{x}\stackrel{~}{z}+{\displaystyle \frac{\beta _{}^2}{\beta _{}^2}}[(\mathrm{\Delta }x)^2(\mathrm{\Delta }y)^2],`$ (7)
where $`\stackrel{~}{x}=xx`$, $`\mathrm{\Delta }x=\sqrt{x^2x^2}`$, $`\beta _{}=K_T/K^0`$ and $`\beta _{}=K_L/K^0`$. From the expressions above, we can interprete the size parameters as the extents of the source if the last two terms in Eqs. (6) and (7) are sufficiently small. It has been already shown that the interpretation holds within a class of thermal source model , and the case of a hydrodynamical model is the point of the present paper.
Figure 1 shows $`K_T`$ dependence of the pion HBT radii. Close circles stand for the HBT radii extracted from the pion correlation function (1) through the fitting. Solid lines are space-time extensions calculated from Eqs. (5)-(7), which are expected to agree with the HBT radii when only thermal pions are considered. Dotted lines stand for the “source sizes” ($`\mathrm{\Delta }z`$ for $`R_{}`$ and $`\mathrm{\Delta }t`$ for $`R_0`$). As far as $`R_{}`$ and $`R_{}`$ are concerned, those quantities seem to be consistent and well reproduce the NA49 experimental results (open circles). As for $`R_0`$, our results are a few fm smaller than the experimental results in spite of the existence of the first order phase transition. Furthermore, the time durations $`\mathrm{\Delta }t`$ do not agree with the $`R_0`$ at large $`K_T`$, that means the failure of the interpretation of $`R_0`$ as the time duration. As explained below, this failure comes from an opaque property of the source.
Fig. 1 $`K_T`$ dependence of YKP radii. Open and closed circles show the experimental data and HBT radii, respectively. Solid lines stand for space-time extensions (5)-(7). Dotted lines stand for source sizes. ($`\mathrm{\Delta }z`$ for $`R_{}`$ and $`\mathrm{\Delta }t`$ for $`R_0`$.)
The “opaque” source emits particles dominantly from the surface and the “transparent” source emits from whole region. Figure 2 shows the source function (3) projected onto the $`xy`$ plane, $`\stackrel{~}{S}_T(x,y)=dzdtS(x^\mu ,K^\mu )`$ being normalized as $`𝑑x𝑑y\stackrel{~}{S}_T(x,y)=1`$. Note that the average momenta of the emitted pions is $`x`$ direction in the figure. The left figure shows clearly surface emission; pions are emitted from the crescent region and the source is thinner in the $`x`$ direction than in the $`y`$ direction. This is typical property of the opaque source and should be appeared in the third and fourth term in Eq. (6) and the factor $`(\mathrm{\Delta }x)^2(\mathrm{\Delta }y)^2`$ works as a measure of source opacity. On the other hand, if we neglected the transverse flow by putting $`U^r=0`$ by hand, the source function becomes as shown in the right figure of Fig. 2. In this case, the source function is proportional to the space-time volume of freeze-out hypersurface. As a consequence of restoration of the azimuthal symmetry which is shown clearly in the figure, the above measure of source opacity vanishes. Although the particle emission takes place almostly from the thin surface, the source is not opaque in this sence. When the transverse flow exists, the source function is deformed by thermal Boltzmann factor $`\mathrm{exp}(K_T\mathrm{sinh}Y_T\mathrm{cos}\varphi /T_f)`$ with $`Y_T`$ almost proportional to $`r`$. Consequently, the number of emitted particles increases in the region where $`x>0`$ (i.e., $`\mathrm{cos}\varphi >0`$) and decreases in the region where $`x<0`$ ($`\mathrm{cos}\varphi <0`$). This flow effect deforms the surface dominant volume distribution (right figure) to the crescent shape (left figure).
Fig. 2 The source functions as functions of transverse coordinates $`x`$ and $`y`$ for $`K_T=450`$ MeV and $`Y_{\pi \pi }=4.15`$ ($`Y_{\pi \pi }=\frac{1}{2}\mathrm{ln}\frac{K^0+K_L}{K^0K_L}`$). The left figure denotes the source function with the transverse flow. The right figure denotes the source function without transverse flow.
In summary, we analyze $`K_T`$ dependence of source parameters of the YKP parametrization based on the relativistic hydrodynamics for the CERN-SPS 158 GeV/A Pb+Pb collisions. We obtain the results almostly consistent with the experiment. However, the temporal source parameter, $`R_0`$, shows the behavior different from the experiment. The source opacity makes the interpretation of $`R_0`$ as the time duration hard. We found that the source opacity was caused by the transverse flow and the characteristics of the freeze-out hypersurface, the surface dominant freeze-out. The deviations of our results from the experiment would be improved by including the resonance decay and other effects.
The authors are indebted to Professor I. Ohba and Professor H. Nakazato for their helpful comments. This work was partially supported by a Grant-in-Aid for Science Research, Ministry of Education, Science and Culture, Japan (Grant No. 09740221) and Waseda University Media Network Center.
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# Four-manifolds which admit ℤ_𝑝×ℤ_𝑝 actions
## 1. Introduction
In 1970, Orlik and Raymond proved that any closed, simply connected four-manifold which admits a smooth, effective $`S^1\times S^1`$ action can be expressed as a connected sum of copies of $`S^2\times S^2`$, $`P^2`$, and $`P^2`$. Later, Fintushel and Yoshida each showed that the same conclusion holds for smooth $`S^1`$ actions. In 1995, Huck generalized this result to show that the intersection form of a closed cohomology four-manifold $`M`$ with $`H_1(M)=0`$ on which $`S^1`$ acts must split as a sum of rank $`1`$ and $`2`$ forms $`(\pm 1)`$ and $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$, provided a certain regularity condition holds near the fixed-point set of the action. Huck and Puppe subsequently generalized further by removing the restriction on $`H_1(M)`$.
Stated simply, Huck’s approach is to study the equivariant cohomology of the singular set of an $`S^1`$ action using earlier techniques of Puppe , and thereby derive a characterization of the possible intersection forms. Related methods were used independently by the author to study actions of finite nonabelian groups on four-manifolds. Our methods actually simplify somewhat when the groups are abelian, and we apply them here to prove: If $`M`$ is a closed four-manifold with $`H_1(M)=0`$ which admits a locally linear, homologically trivial action by $`_p\times _p`$ (with $`p`$ prime), then the intersection form of $`M`$ splits as a sum of copies of $`(\pm 1)`$ and $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$.
If this $`M`$ is simply connected, then by the work of Freedman , it must be homeomorphic to a connected sum of copies of $`S^2\times S^2`$, $`\pm P^2`$, and perhaps a copy of $`\pm \widehat{P^2}`$, where $`\widehat{P^2}`$ denotes the manifold homotopy equivalent to $`P^2`$, but with non-vanishing Kirby-Siebenmann invariant. We generalize an observation of Wilczyński to show (with exactly one exception) that $`KS(M)`$ must vanish. As a corollary, we obtain an analogue of Orlik and Raymond’s result for $`_p\times _p`$ actions.
Finally, we discuss the question of classifying the actions themselves. A complete classification would be very difficult, but by combining our results with those of Orlik and Raymond, we prove an equivariant version of the decomposition theorem which reduces the general question to that of classifying the possible actions on $`S^4`$.
## 2. The singular set of a $`_p\times _p`$ action
Suppose $`M`$ is a closed, connected four-manifold with $`H_1(M)=0`$ and $`b_2(M)1`$. If $`G=_p\times _p`$ acts on $`M`$, the set $`\mathrm{\Sigma }=\{xM|G_x\{0\}\}`$ is called the *singular set* of the action. We assume throughout the paper that the action is effective, locally linear, and homologically trivial. By results of Edmonds , each $`g0`$ in $`G`$ has a fixed point set consisting of isolated points and $`2`$-spheres, and each $`2`$-sphere represents a nontrivial homology class. Our first task is to understand how the fixed-point sets of the cyclic subgroups $`_pG`$ fit together to form the overall structure of $`\mathrm{\Sigma }`$.
Recall that the Borel equivariant cohomology $`H_G(X)`$ of a $`G`$-space $`X`$ is the ordinary cohomology of the balanced product $`EG\times _GX`$. This balanced product has a natural fibration over $`EG/G=BG`$, and the Leray-Serre spectral sequence of the fibration is called the *Borel spectral sequence*.
The next lemma and the proposition which follows it appeared in slightly different form in Edmonds and , but we re-state them here for convenience:
###### Lemma 2.1.
Suppose $`G`$ acts homologically trivially on a closed four-manifold $`M`$ with $`H_1(M)=0`$. If either
1. $`H^2(M)`$ contains a class $`u`$ whose square generates $`H^4(M)`$, and $`H^3(G)`$ has no $`3`$-torsion, or
2. $`b_2(M)3`$,
then the Borel spectral sequence $`E(M)`$ collapses with coefficients in $``$ or any field. It follows that $`H_G^{}(M)`$ is a free $`H^{}(G)`$ module on $`b_2(M)+2`$ generators corresponding to generators for $`H^{}(M)`$.
###### Proof.
If $`uH^2(M)`$ has nonzero square, then, since $`u^3=0`$, $`0=d_3(u^3)=3d_3(u)u^2`$. But $`E_3^{,4}`$ is a free $`H^{}(G)`$ module generated by $`u^2`$. So if $`H^3(G)`$ has no $`3`$-torsion, then $`d_3(u)`$ must be $`0`$. And then, of course, $`d_3(u^2)=0`$, as well. Thus $`E_2(M)=E_3(M)`$. Since $`d_5(u^2)=2ud_5(u)=0`$, the sequence collapses.
Now suppose $`b_2(M)3`$. Then for each generator $`uH^2(M)`$, there is a $`vH^2(M)`$ which is linearly independent of $`u`$ in $`H^2(M)`$, and such that $`uv=0`$. Since the action of $`G`$ is homologically trivial, $`E_2(M)`$ is a free $`H^{}(G)`$-module on generators corresponding to those of $`H^{}(M)`$, so in fact $`u`$ and $`v`$ must be independent in $`E_2(M)`$, as well. But $`d_3(uv)=ud_3(v)+vd_3(u)=0`$. This is only possible if $`d_3(u)=d_3(v)=0`$. $`H^4(M)`$ is generated by products of two-dimensional classes, so $`d_3^{,4}=0`$, as well. It follows that $`E_2(M)=E_3(M)=E_4(M)`$. The same argument shows that $`E_4(M)=E_5(M)=E_{\mathrm{}}(M)`$.
The conclusion about $`H_G^{}(M)`$ follows immediately from tom Dieck \[14, III.1.18\]. ∎
Whenever $`G`$ is a finite group and $`S`$ is a multiplicative, central subset of the cohomology ring $`H^{}(G)`$, we can define the “$`S`$-singular set” $`\mathrm{\Sigma }_S=\{xX|S\mathrm{ker}r_{G_x}^{}=\mathrm{}\}`$. The fundamental *Localization Theorem* (See Hsiang or tom Dieck ) then states that the localized restriction map $`S^1H_G^{}(M)S^1H_G^{}(\mathrm{\Sigma }_S)`$ is an isomorphism. Applying the Localization Theorem in specific cases requires careful choice of $`S`$, based on knowledge of the restriction maps from the cohomology of $`G`$ to that of its subgroups. But it can yield useful information about the structure of $`\mathrm{\Sigma }`$. We apply it to prove:
###### Proposition 2.2.
Let $`M`$ be a closed four-manifold such that $`H_1(M)=0`$, and suppose that either $`b_2(M)3`$, or $`b_2(M)=2`$ but the intersection form of $`M`$ is diagonalizable over $``$. If $`G`$ is a finite abelian group which acts locally linearly and homologically trivially on $`M`$, then the rank of $`G`$ is at most 2, and $`G`$ has nonempty fixed-point set. If $`G=_p\times _p`$, then $`\mathrm{Fix}(G)`$ consists of exactly $`b_2(M)+2`$ points.
###### Proof.
Suppose first that $`G=_2\times _2\times _2`$ acts on $`M`$. By Lemma 2.1, $`H_G^{}(M;_2)`$ is a free $`H^{}(G,_2)`$ module on $`b_2(M)+2`$ generators. Recall that $`H^{}(G;_2)_2[a]_2[b]_2[c]`$, where $`a,b,`$ and $`c`$ generate $`\mathrm{Hom}(G,_2)`$. Since $`H^{}(G;_2)`$ is a polynomial ring, it contains no zero-divisors, so it makes sense to localize at the set $`S`$ consisting of all of the nonzero elements. We check easily that $`S^1H^{}(G,_2)_2`$.
Now, each proper subgroup $`HG`$ is the kernel of some nonzero homomorphism $`\phi _H:G_2`$, and this $`\phi _H`$, viewed as an element of $`H^1(G;_2)`$, restricts trivially to $`H^1(H;_2)`$. So the $`S`$-singular set $`\mathrm{\Sigma }_S`$ contains only those points fixed by all of $`G`$. By the Localization Theorem, $`\mathrm{\Sigma }_S`$ is nonempty. Consideration of the isotropy representation of $`G`$ at a fixed point $`x_0`$ shows that there must be $`g,hG`$ such that $`g`$ fixes a two-dimensional subspace $`VT_{x_0}`$, while $`h|_V`$ acts by $`\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$. But $`V`$ forms part of a $`2`$-sphere $`S`$ fixed by $`g`$. If $`h`$ reverses orientation on $`V`$, it also acts by $`1`$ on $`[S]H_2(M)`$, contradicting homological triviality.
If $`G=_2\times _2`$, the same argument shows that $`\mathrm{Fix}(G)`$ contains $`b_2+2`$ points, but of course no contradiction ensues from the isotropy representation of $`G`$.
If $`p`$ is odd, a similar argument applies, except in the case where $`b_2(M)=2`$ and $`p=3`$. To ensure that $`S`$ is central in $`H^{}(G;_p)`$, we replace the one-dimensional generators of $`\mathrm{Hom}(G,_p)`$, with their two-dimensional images under the Bockstein map, which generate the polynomial part of $`H^{}(G,_p)`$.
Finally, suppose $`b_2(M)=2`$ and $`G=_3\times _3\times _3`$, with generators $`g,h`$, and $`k`$. The Lefschetz fixed-point theorem implies that $`\chi (\mathrm{Fix}(g))=4`$. Thus $`\mathrm{Fix}(g)`$ either contains at least one $`2`$-sphere, or consists of exactly four isolated points. In the first case, $`G/g`$ acts effectively on the sphere, which is impossible. In the second case, the action of $`h`$ on $`\mathrm{Fix}(g)`$ must fix at least one point $`x_0`$. But $`g,h`$ cannot act freely on the linking sphere to $`x_0`$, so some other element $`g^{}`$ fixes a $`2`$-sphere, and the argument proceeds as before.
There are indeed actions of $`_2\times _2`$ on $`S^2\times S^2`$, and pseudofree actions of $`_3\times _3`$ on $`P^2`$ and $`\widehat{P^2}`$, whose fixed point set is empty. Inspection shows that such actions are the only exceptions to the rule that $`\mathrm{Fix}(G)\mathrm{}`$, and that if $`\mathrm{Fix}(G)\mathrm{}`$, then in fact it contains $`b_2+2`$ points. Because our desired conclusion about the intersection form holds in the exceptional cases, we assume as a convenience in this section and the next that the fixed-point set is non-empty. (Later, in the geometrical analysis, the assumption will be more essential, and we will not take it for granted.)
Since the isotropy representation of $`G`$ at any fixed point splits as a sum of rank two real representations, each fixed point is included in exactly two singular $`2`$-spheres. Since $`G`$ is abelian, $`G`$ acts on $`\mathrm{Fix}(g)`$ for each $`gG`$, so each sphere has a rotation action with fixed points at its north and south poles. Thus $`\mathrm{\Sigma }`$ contains a total of $`b_2+2`$ spheres $`S_1,\mathrm{},S_{b_2+2}`$, and each path component of $`\mathrm{\Sigma }`$ is a chain of such spheres arranged in a closed loop. Since the action of $`G`$ on $`\mathrm{\Sigma }`$ is just a rotation on each sphere, $`G`$ acts trivially on $`H^{}(\mathrm{\Sigma })`$.
###### Lemma 2.3.
If $`p=2`$, each $`[S_i]`$ represents a primitive homology class in $`H_2(M,)`$.
###### Proof.
If each component of $`\mathrm{\Sigma }`$ contains at least three spheres, then each sphere intersects its neighbor geometrically once, and the claim follows. If some component contains exactly two spheres, then each intersects the other twice. One of them might, a priori, represent a multiple of two in $`H_2(M,)`$. But the theorem of Edmonds cited above implies that it must be nontrivial in $`H_2(M;_2)`$. ∎
If $`p`$ is odd, this argument does not suffice to rule out certain $`[S_i]`$ being multiples of 2. However, if $`p`$ is odd, then $`2`$-torsion will not affect the cohomology calculations of the next section. The calculations of that section will show that $`\mathrm{\Sigma }`$ is connected, and then it will follow that each $`S_i`$ does, in fact, intersect its neighbor only once.
Our next goal is to show that the inclusion $`H_2(\mathrm{\Sigma })H_2(M)`$ is (split) surjective. When we have shown this, it will follow that the intersection form of $`M`$ is represented by the geometrical intersections of the spheres in $`\mathrm{\Sigma }`$.
From the cohomology long exact sequence of the pair $`(M,\mathrm{\Sigma })`$, we extract:
$$0H^1(\mathrm{\Sigma })H^2(M,\mathrm{\Sigma })H^2(M)H^2(\mathrm{\Sigma })H^3(M,\mathrm{\Sigma })0.$$
A short diagram chase shows that $`G`$ acts trivially on the relative cohomology groups. Let $`N`$ denote the number of path components of $`\mathrm{\Sigma }`$, and $`L`$, the (integral) rank of $`\mathrm{coker}H^1(\mathrm{\Sigma })H^2(M,\mathrm{\Sigma })`$. As we have noted, each $`S_i`$ represents an “almost primitive” homology class in $`M`$. More precisely: $`H^3(M,\mathrm{\Sigma })^{L+2}T`$, where $`T=0`$ if $`p=2`$, and $`2T=0`$ if $`p`$ is odd.
We shall prove:
###### Lemma 2.4.
$`L=0`$,
From which the claim about $`H_2(\mathrm{\Sigma })H_2(M)`$ is immediate.
###### Proof.
Recall that
$$H^{}(_2\times _2;)\frac{[\alpha _2,\beta _2]P[\mu _3]}{2\alpha =2\beta =2\mu =0,\mu ^2=\alpha \beta ^2+\alpha ^2\beta },$$
while for $`p`$ odd,
$$H^{}(_p\times _p;)\frac{[\alpha _2,\beta _2]E[\mu _3]}{p\alpha =p\beta =p\mu =0}.$$
Let $`\pi `$ denote the projection $`MM/G=M^{}`$. The Borel spectral sequence of the pair $`(M,\mathrm{\Sigma })`$ has
$$E_2^{i,j}(M,\mathrm{\Sigma })=H^i(G;H^j(M,\mathrm{\Sigma }))H_G^{}(M,\mathrm{\Sigma }).$$
On the other hand, $`M\mathrm{\Sigma }`$ is a free $`G`$-space, so $`H_G^{}(M,\mathrm{\Sigma })`$ is canonically isomorphic to $`H^{}(M^{},\mathrm{\Sigma }^{})`$.
Since $`(M^{},\mathrm{\Sigma }^{})`$ is a relative manifold pair, Poincaré duality gives a commutative diagram:
$$\begin{array}{ccccc}& & H^3(M,\mathrm{\Sigma })& \stackrel{}{}& H_1(M\mathrm{\Sigma })\\ & & \pi ^{}& & \pi _{}& & \\ H_G^3(M,\mathrm{\Sigma })& \stackrel{}{}& H^3(M^{},\mathrm{\Sigma }^{})& \stackrel{}{}& H_1(M^{}\mathrm{\Sigma }^{}).\end{array}$$
But $`H_1(M\mathrm{\Sigma })`$ is generated by meridians to the spheres in $`\mathrm{\Sigma }`$, and each of these is a $`p`$-fold cover of its image in $`H_1(M^{}\mathrm{\Sigma }^{})`$. Thus $`\pi _{}`$ is multiplication by $`p`$. Since the left-hand edge homomorphism $`E_2^{0,j}E_{\mathrm{}}^{0,j}`$ of the Borel spectral sequence is induced by the fiber inclusion $`j:(M,\mathrm{\Sigma })(M_G,\mathrm{\Sigma }_G)`$, we can conclude that $`\mathrm{coker}(E_2^{0,3}E_{\mathrm{}}^{0,3})`$ has exponent $`p`$. In other words, no $``$ summand of $`E_2^{0,3}`$ supports more than one non-zero differential. In rank counting arguments we can therefore treat its integral rank as though it were a $`_p`$ rank.
Notice also that, since $`H^4(M^{},\mathrm{\Sigma }^{})`$, each nonzero class in $`E_2^{i,j}`$ with $`i+j4,i0,`$ must be mortal.
Consider the terms of $`E_2(M,\mathrm{\Sigma })`$ indicated in Table 1:
Now, elements of $`E_2^{3,1}`$ can only be killed by $`d_2^{0,2}`$, while $`E_2^{2,2}`$ can be killed either by $`d_2^{0,3}`$ or $`d_2^{2,2}`$. By the above observations, we have:
$`\mathrm{rk}E^{0,3}+\mathrm{rk}E^{4,1}`$ $`\mathrm{rk}E^{2,2}+\mathrm{rk}E^{3,1}`$
$`(2+L)+(3N3)`$ $`(2N+2L)+(N1),`$
so $`L=0`$, as claimed. ∎
We obtain a corollary which is dual, in some sense, to Edmonds’s theorem \[2, 2.5\]:
###### Corollary 2.5.
Suppose $`G=_p\times _p`$ acts as we have been assuming. Then the cohomology restriction map $`H^2(M)H^2(\mathrm{\Sigma })`$ is injective, so $`\mathrm{\Sigma }`$ represents all of $`H_2(M)`$.
To better understand the geometry of the singular set, we also prove:
###### Lemma 2.6.
$`N=1`$, so $`\mathrm{\Sigma }`$ is connected.
###### Proof.
From the homology spectral sequence of the covering $`\pi :X\mathrm{\Sigma }X^{}\mathrm{\Sigma }^{}`$, we obtain a short exact sequence $`0\times \times T\stackrel{\pi _{}}{}H_1(M^{}\mathrm{\Sigma }^{})_p\times _p0`$. As we have already seen, $`\pi _{}`$ is multiplication by $`p`$ in each factor. It follows that $`H_1(M^{}\mathrm{\Sigma }^{})H^3(M^{},\mathrm{\Sigma }^{})`$ is $`p`$-torsion-free. So in fact all classes in $`E_2^{i,j}`$ with $`i>0`$ are mortal. In particular, $`E_2^{1,2}`$ must vanish, so $`N2(N2)`$, so $`N2`$. Suppose for a contradiction that $`N=2`$. The Borel spectral sequence then takes the following form (each entry with $`i>0`$ is $`_p^k`$ for some $`k`$, so to save space, we simply indicate its rank):
There are generators $`aH^1(M,\mathrm{\Sigma })`$ and $`b,cH^2(M,\mathrm{\Sigma })`$ such that $`d_2(b)=\alpha a`$ and $`d_2(c)=\beta a`$. By the multiplicative properties of the spectral sequence, this kills the entire row $`j=1`$, *except* $`H^3(G;H^1(M,\mathrm{\Sigma })_p\mu `$.
Now, $`\mathrm{ker}d_2^{2,2}=\beta b\alpha c`$, so there is some $`eH^3(M,\mathrm{\Sigma })`$ such that $`d_2(e)=\beta b\alpha c`$. Since $`E_3^{3,1}=\mu a`$ must also perish, there is $`fH^3(M,\mathrm{\Sigma })`$, independent of $`e`$, so that $`d_3(f)=\mu a`$. But then $`d_3(\alpha f)=d_3(\beta f)=0`$, since $`\mu \alpha a`$ and $`\mu \beta b`$ were already killed by $`d_2`$. Now $`\mathrm{ker}d_2^{2,3}`$ has rank $`2`$, and $`d_3^{2,3}=0`$. But $`d_2^{0,4}`$ has rank $`1`$, so $`E_{\mathrm{}}^{2,3}`$ must have rank $`1`$. This is a contradiction, so $`N=1`$. ∎
Now that we know this, each $`S_i`$ definitely intersects each neighbor only once, so $`T=0`$ for odd $`p`$.
To summarize, we have shown:
###### Proposition 2.7.
Suppose $`M`$ is a closed, topological four-manifold with $`b_2(M)1`$ and $`H_1(M)=0`$, equipped with an effective, homologically trivial, locally linear $`_p\times _p`$ action. With the exception of fixed-point free actions which exist in the two cases,
1. $`b_2(M)=1`$, $`p=3`$, and the action is pseudofree, or
2. $`M`$ has intersection form $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ when $`p=2`$,
the singular set $`\mathrm{\Sigma }`$ consists of $`b_2(M)+2`$ spheres equipped with rotation actions, intersecting pairwise at their poles, and arranged into a single closed loop. Each sphere represents a primitive class in $`H_2(M;)`$, and together these classes generate $`H_2(M)`$.
## 3. The intersection form
Let $`\sigma _1,\mathrm{},\sigma _{b_2+2}`$ denote the fundamental classes of $`S_1,\mathrm{},S_{b_2+2}\mathrm{\Sigma }`$. As generators of $`H_2(M)`$, two can be regarded as “redundant”. If we eliminate one, we cut the loop of $`\mathrm{\Sigma }`$. By removing another, we either disconnect or shorten the remaining chain. Renumber the remaining spheres, if necessary, as $`S_1,\mathrm{},S_{b_2}`$, and call the result $`\mathrm{\Sigma }^{}`$. Let $`e_i=\sigma _i\sigma _i`$. The matrix of intersections of the spheres, and therefore the intersection form of $`M`$, as well, takes the form of one, or a sum of two, pieces of the form
$$\left(\begin{array}{ccccccc}e_1& 1& & & & & \\ 1& e_2& 1& & & & \\ & 1& e_3& 1& & & \\ & & 1& \mathrm{}& & & \\ & & & & & 1& \\ & & & & 1& e_{k1}& 1\\ & & & & & 1& e_k\end{array}\right)$$
Huck and Yoshida have already proven exactly the lemma we need about such matrices (See Huck \[9, lemma 4.2\]): Each is equivalent to a sum of rank $`1`$ and $`2`$ pieces. Thus:
###### Theorem 3.1.
Suppose $`M`$ is a closed topological four-manifold with $`H_1(M)=0`$. Let $`p`$ be prime, and suppose $`_p\times _p`$ acts effectively, locally linearly, and homologically trivially on $`M`$. Then the intersection form of $`M`$ is a sum of copies of $`(\pm 1)`$ and $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$.
## 4. Vanishing of $`KS(M)`$
Edmonds showed that when $`p`$ is a prime greater than 3, locally linear, homologically trivial $`_p`$ actions exist on every simply-connected four-manifold. The actions he constructs are *pseudofree* – i.e. with only isolated fixed points. In certain cases, this is a necessary restriction. For example, Wilczyński shows that if a homotopy $`P^2`$ admits a $`_p`$ action which fixes a two-sphere, then the two-sphere can be used to split off $`P^2`$ as a connected summand of $`M`$, and it follows that $`M`$ is homeomorphic to $`P^2`$. In other words, $`\widehat{P^2}`$ admits *only* pseudofree actions.
For the remainder of the paper, we assume $`M`$ is simply connected. We generalize Wilczyński’s construction to prove the following corollary of Theorem 3.1:
###### Theorem 4.1.
If $`M`$ is a closed, simply connected four-manifold which admits a locally linear, homologically trivial $`_p\times _p`$ action, then $`M`$ is homeomorphic to a connected sum of copies of $`\pm P^2`$ and $`S^2\times S^2`$, or if $`p=3`$ and the action is pseudofree, perhaps to a single copy of $`\pm \widehat{P^2}`$.
(In Theorem 5.2, we will establish a sharper result.)
###### Proof.
For convenience, we continue to assume the action is not pseudofree, or fixed-point-free in the case of $`S^2\times S^2`$. As above, let $`\mathrm{\Sigma }^{}`$ denote the singular set with two homologically redundant spheres removed. For simplicity of notation, we assume $`\mathrm{\Sigma }^{}`$ is connected; if it isn’t, our argument will carry through on each piece.
It follows from the work of Freedman and Quinn \[6, 9.3\] (see also \[7, 1.2\]) that each $`S^2\mathrm{\Sigma }`$ has an equivariant normal bundle. Thus $`\mathrm{\Sigma }^{}`$ has a regular neighborhood $`N(\mathrm{\Sigma }^{})`$ which is homeomorphic to the manifold obtained by plumbing together disk bundles $`E(e_i)`$, $`i=1,\mathrm{},b_2`$, over $`S^2`$ according to the graph
$$A_{b_2}=e_1\text{ }e_2\text{ }\mathrm{}\text{ }e_{b_21}\text{ }e_{b_2}.$$
The boundary of such a plumbed manifold is a lens space $`L`$, and $`|H_1(L)|`$ is given by the determinant of the intersection matrix. In our case, the matrix is unimodular, so $`L`$ is in fact a three-sphere. Thus $`M^{}=N(\mathrm{\Sigma }^{})_{S^3}D^4`$ is homeomorphic to a connected sum of copies of $`\pm P^2`$ and $`S^2\times S^2`$. But $`M^{}`$ is also a connected summand of $`M`$ which carries all of its homology. By Freedman and Quinn \[6, 10.3\], $`M^{}M`$. ∎
## 5. On classifying $`_p\times _p`$ actions
The argument of Orlik and Raymond on the classification of torus actions, specialized to the simply-connected case, can be summarized as follows: The quotient space $`M/T`$ is a surface with boundary, and since $`H_1(M)=0`$, it must be a disk. The boundary of $`D`$ consists of fixed points and arcs; the arcs can be labeled according to the corresponding isotropy subgroups of $`T`$. Each arc lifts to a singular $`S^2`$ and each interior point of the disk represents a principal orbit. They show that the quotient map in fact admits an essentially unique section; thus the singular data in the quotient space determine $`M`$ up to equivariant diffeomorphism. A calculation involving the particular isotropy groups then shows that the quotient space splits in a way which lifts to an equivariant connected sum decomposition of $`M`$.
Up to an automorphism of $`T`$, listing the one-dimensional isotropy groups is equivalent to listing the Euler classes and (signed) intersection numbers of the singular $`2`$-spheres. This information is also available for $`_p\times _p`$ actions. To what extent does it classify them? We will show:
###### Proposition 5.1.
Assume the action is not one of the exceptional fixed-point-free cases.
1. Each $`_p\times _p`$ action extends to a torus action in a regular neighborhood $`\nu (\mathrm{\Sigma })`$ of $`\mathrm{\Sigma }`$.
2. $`\nu (\mathrm{\Sigma })`$ is $`T`$-equivariantly diffeomorphic to the singular set of some smooth $`T`$-action on $`M`$, but the given $`T`$-action need not extend over $`M`$.
###### Proof.
We begin with a slight variant of the plumbing construction of the previous section: Let $`t=b_2(M)+2`$, and label the spheres consecutively around $`\mathrm{\Sigma }`$ as $`S_1,\mathrm{},S_t`$. Let $`x_i`$ denote the “north pole” of $`S_i`$. Choose orientations for each of the $`S_i`$, and let $`\sigma _i`$ denote the corresponding fundamental class. Finally, choose an orientation for $`M`$ and let $`ϵ_i=\sigma _{i1}\sigma _i`$ denote the sign of the intersection at $`x_i`$. (When we considered $`\mathrm{\Sigma }^{}`$ earlier, we implicitly chose orientations to make each $`ϵ_i=+1`$; here, because the spheres are arranged in a closed loop, this might not be possible.)
With these conventions, $`\nu (\mathrm{\Sigma })`$ is obtained by plumbing together $`D^2`$-bundles $`\xi _i`$ over $`S^2`$, each with Euler class $`e_i`$, according to the orientations given by the $`ϵ_i`$. The plumbing graph is a circle, which we parameterize as $`[\frac{1}{2},t+\frac{1}{2}]`$ with the endpoints identified. $`\nu `$ can be thought of as a torus fiber bundle over the plumbing graph with a fiber-preserving, free $`_p\times _p`$ action. It is not a priori a principal bundle, but if the $`_p\times _p`$ action on the fibers extends to a torus action, it will become one. With appropriate smoothing around the plumbing points, the torus action will extend over $`\nu (\mathrm{\Sigma })`$, establishing the first part of the proposition.
The $`T`$-bundle over $`[\frac{1}{2},t+\frac{1}{2}]`$ can be assembled by gluing copies of $`T\times [i\frac{1}{2},i+\frac{1}{2}]`$ via attaching maps $`\gamma _i`$ which incorporate the clutching functions for the $`\xi _i`$, the coordinate switches at each plumbing point, and the orientations $`ϵ_i`$. (See figure 1, which is intended to invite comparison with the diagrams in .)
The maps are determined up to isotopy by their $`\pi _1(T)`$ representations. The clutching functions take the form $`\left(\begin{array}{cc}1& 0\\ e_i& 1\end{array}\right)`$; the coordinate switches are of course $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$, and the orientation changes, $`\left(\begin{array}{cc}1& 0\\ 0& ϵ_i\end{array}\right)`$. Together, such matrices generate $`GL(2,)`$, the structure group of the bundle.
The $`_p\times _p`$ action has well-defined rotation numbers in the fiber over $`x_1`$, so it extends to a torus action in that fiber. The gluing maps in $`GL(2,)`$ define a trivialization of the bundle over $`T\times [\frac{1}{2},t+\frac{1}{2}]`$ which is equivariant with respect to the $`_p\times _p`$ action. Using them, the torus action extends along all of the fibers. The structure of the torus bundle is thus determined by the total gluing function $`\gamma :T\times \{t+\frac{1}{2}\}T\times \{\frac{1}{2}\}`$; with slight abuse of notation, we may write $`\gamma =\gamma _t\mathrm{}\gamma _2\gamma _1`$.
A compatibility condition is imposed by the existence of the $`_p\times _p`$ action – namely, that the gluing map $`\gamma GL(2,)`$ must commute with order $`p`$ rotations in each factor of $`T\times 0`$. We may analyze this requirement by lifting to the universal cover $`\pi :\stackrel{~}{T}T`$. A rotation $`r`$ lifts to a translation $`\tau `$. The requirement that $`\pi _{}(\gamma ^1\tau \gamma )=\pi _{}(\tau )=r`$ means that the line spanned by each $`\tau `$ must be
1. Normalized by $`\gamma `$, if $`p=2`$. Since the total space of the bundle is the boundary of $`\nu (\mathrm{\Sigma })`$, it is orientable, so $`\gamma `$ is one of $`\pm \left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$ or $`\pm \left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),`$
2. Centralized by $`\gamma `$, if $`p>2`$, which implies $`\gamma =\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$.
In the latter case, $`\gamma `$ clearly commutes with the entire torus action, so $`\nu `$ supports the structure of a principal bundle. Even when $`p=2`$, $`\gamma `$ must respect the base-fiber splitting of the bundle $`\xi _t`$ over $`S_t`$, so $`\pm \left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ is ruled out. We proceed to rule out $`\gamma =\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$, also. If such a bundle is realized on the singular set of a $`_2\times _2`$ action, let $`\mu `$ be a small meridional loop around $`S_t`$ in $`M\mathrm{\Sigma }`$. Then $`\mu `$ is homologous to $`\mu `$, and so $`2\mu =0`$ in $`H_1(M\mathrm{\Sigma })`$. But $`H_1(M\mathrm{\Sigma })H^3(M,M\mathrm{\Sigma })`$ is torsion-free, as we saw in section 2. It is generated by any pair of meridians to neighboring two-spheres in $`\mathrm{\Sigma }`$.
This finishes the proof that $`\nu `$ is a trivial principal $`T`$-bundle, and hence also the proof of the first claim. Our proof of the second claim is constructive, based on Orlik and Raymond’s model in the case of torus actions.
The orbit space $`A=\nu (\mathrm{\Sigma })/T`$ is an annulus. Its outer boundary component $`_1A`$ consists of $`t`$ fixed points separated by $`t`$ arcs whose stabilizers are copies of $`S^1T`$. Its inner boundary $`_2A`$ consists entirely of principal orbits. Adjoin a disk $`D`$ to $`_2A`$. Because the torus bundle is trivial over $`_2A`$, there is no obstruction to lifting this adjunction to a $`T`$-equivariant gluing of $`D^2\times T`$ to $`\nu `$. The resulting manifold, denoted $`M^{}`$, is simply connected and has the same intersection form as $`M`$, so it is homeomorphic to $`M`$.
Finally, an example of Hambleton, Lee, and Madsen () shows that $`M`$ and $`M^{}`$ need not be $`_p\times _p`$-equivariantly homeomorphic. They begin with a linear $`_p\times _p`$ action on $`P^2`$, and equivariantly connect sum a $`_p`$-orbit of counterexamples to the Smith conjecture in $`S^4`$ around one of the singular $`2`$-spheres. The resulting space is still homeomorphic to $`P^2`$, but the complement of the singular set has nonabelian fundamental group. In the linear example, $`P^2\mathrm{\Sigma }`$ has the homotopy type of a torus. ∎
Let us call a $`_p\times _p`$ action *standard* if it is the restriction of a smooth torus action. It is fair to say that the standard actions are completely understood. Proposition 5.1, together with the construction of section 4, shows that we can equivariantly split off standard summands. If the two “redundant” two-spheres are adjacent, then $`MM_{\text{standard}}\mathrm{\#}S^4`$, while if there is no such choice of adjacent spheres, a two-step splitting still yields $`MM_{\text{standard}}^{}\mathrm{\#}M_{\text{standard}}^{\prime \prime }\mathrm{\#}S^4`$. Because the standard actions extend to torus actions, Orlik and Raymond’s classification theorem applies to show that each splits further into “irreducible” pieces. This proves:
###### Theorem 5.2.
Let $`M`$ be a closed, simply-connected four-manifold with an effective, locally linear, homologically trivial $`_p\times _p`$ action. Assume the action is not one of the fixed-point-free exceptions. Then $`M`$ admits an equivariant connected sum decomposition
$$MS^4\mathrm{\#}M_1\mathrm{\#}\mathrm{}\mathrm{\#}M_k,$$
where each $`M_i`$ is one of $`S^4`$, $`S^2\times S^2`$, $`\pm P^2`$, or $`P^2\mathrm{\#}P^2`$, equipped with a standard action. The action on the first $`S^4`$ summand need not be standard.
Recall that the “fixed-point-free exceptions” are the pseudofree actions of $`_3\times _3`$ on $`P^2`$ and $`\widehat{P^2}`$, and fixed-point-free actions of $`_2\times _2`$ on $`S^2\times S^2`$. Also note that Orlik and Raymond construct examples of torus actions on $`P^2\mathrm{\#}P^2`$ which admit no *equivariant* connected sum decomposition.
As a consequence of this theorem, the general problem of classifying $`_p\times _p`$ actions on simply-connected four-manifolds reduces to the question of classifying actions on $`S^4`$. The latter is still, of course, very difficult.
## 6. Questions
Finally, we leave the reader with two questions:
1. What can constructively be said about the classification of $`_p\times _p`$ actions on $`S^4`$, in light of the possible knotting of the singular set?
2. Huck and Puppe generalized Huck’s earlier work on circle actions to the case $`H_1(M)0`$. Does Theorem 3.1 generalize similarly? It is worth noting that in the general case, the singular set need not contain spheres, as examples of free actions on $`T^4`$ easily show.
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# An X-ray Flare Detected on the M8 dwarf VB 10
## 1 Introduction
Stars at the very bottom of the main sequence have been a very popular topic of research during the last few years, including in studies of stellar activity. Fleming, Schmitt, & Giampapa (1995) demonstrated that fully-convective M dwarfs, i.e. those with masses less than 0.3 M, exhibit high levels of coronal activity, although at that time, virtually no X-ray detections had been made of stars less massive than VB 8 and LHS 3003 (both type M7.) The star VB 10 (a.k.a. GL 752 B, LHS 474; type M8) appeared to have been marginally detected once by the Einstein High Resolution Imager and was catalogued by Barbera et al. (1993). Then Linsky et al. (1995) reported the detection of a large amplitude flare on VB 10 in the far UV using the Goddard High Resolution Spectrograph. More recently, Neuhauser & Comòron (1998) have detected X-ray emission from a very young brown dwarf in Chameleon. This confirms that activity does exist for the very lowest mass stellar (or substellar) configurations.
We have used the ROSAT High Resolution Imager (HRI) to reobserve VB 10, which was not detected by the ROSAT Position Sensitive Proportional Counter (PSPC) during either the All-Sky Survey or a 7.3-ksec pointed observation (Fleming et al. 1993). By obtaining a deeper exposure, we hoped to discover the nature of the VB 10 corona, detect its quiescent X-ray luminosity, and compare its level of coronal activity to that of the more massive M dwarfs. We have succeded in detecting X-ray emission from VB 10, but only during a brief flare similar to the one observed by Linsky et al. (1995) at ultraviolet wavelengths.
In this paper, we present the latest ROSAT data on VB 10. Section 2 contains a description of the observations and data analysis. In Section 3, we discuss a hypothesis to account for the stark contrast between the flare and non-flare X-ray flux values, including the possibility of a total lack of $`10^6`$ K coronal plasma about VB 10.
## 2 Observation
The star VB 10 was observed during the period of 1997 October 27-31 for a total of 21,992 sec with the ROSAT HRI (Zombeck et al. 1990). The total image is shown in Fig. 1a. The one obvious X-ray source seen in Fig. 1 is positionally coincident with the star Wolf 1055 (type M3.5 Ve), which has a common proper motion with VB 10. At the position of VB 10, which is marked in Fig. 1a, there is only a marginal detection. However, when the data are separated into their individual Observation Intervals (OBIs), we find that all of the photons in the VB 10 detection are contained within one single OBI (Figs. 1b and 1c).
This particular OBI, which began on UT 1997 October 29 at 3:05:24 and ended at 3:25:43, had an effective total exposure time of 1,138 sec. Within an extraction radius of $`12\mathrm{}`$ around the optical position of VB 10, 11 photons were detected. By examining the rest of the image, we determined that the density of background photons for this OBI was 0.0024 cts arcs<sup>-2</sup>. The solid angle of our extraction circle about VB 10 was 452.4 arcs<sup>2</sup>, which means that we would expect one photon in the extraction circle to be from the background. Therefore, with 10 source photons, the mean count rate over the entire OBI is $`8.9(\pm 2.8)\times 10^3`$ cts s<sup>-1</sup> (this includes a deadtime and vignetting correction factor of 1.017.)
Of course, to get a better idea of the flare duration and flux, one needs to look at the temporal distribution of the source photons throughout the observation. In Fig 2, we show a histogram in arrival time for the 11 photons contained within the source extraction radius. Remember, we do not know which one is the background photon. The histogram is binned in 3-minute (180 sec) intervals over the nearly 20-minute observation. One can see that 5 photons arrived in the first 3 minutes of the observation, 5 photons arrived during the last 8 minutes, with the remaining photon arriving in between. The data are consistent with there being only one flare which is tailing off during our observation. In this case, the flare duration is at least 20 minutes and $`2.8\times 10^2`$ cts s<sup>-1</sup> (5 photons in 3 minutes) represents a lower limit to the peak count rate.
In order to get the X-ray luminosity, we adopt a conversion factor of $`2.4\times 10^{11}`$ ergs cm<sup>-2</sup> cnt<sup>-1</sup>. This comes from Table 10 of David et al. (1999), the ROSAT HRI Calibration Report (Cambridge: SAO)<sup>1</sup><sup>1</sup>1$`http://heawww.harvard.edu/rosat/rsdc_www/hricalrep.html`$, for a Raymond-Smith spectrum of 0.5 keV and negligible interstellar absorption. This yields an apparent energy flux, $`f_X=2.14(\pm 0.68)\times 10^{13}`$ ergs cm<sup>-2</sup> s<sup>-1</sup> for the mean count rate. At a distance of 5.74 pc, this gives us a luminosity, $`L_X=8.4(\pm 2.7)\times 10^{26}`$ ergs s<sup>-1</sup>. Again, this is just a mean luminosity for the observation. The peak flare luminosity would be at least $`2.65\times 10^{27}`$ ergs s<sup>-1</sup>.
We have also analyzed the remaining 20,854 sec of our HRI observation, in which no X-ray source was detected at the position of VB 10. Using the non-detection analysis software in MIDAS/EXSAS, we have calculated a $`3\sigma `$ (i.e. 99.7% confidence) upper limit of $`1.8\times 10^4`$ cts s<sup>-1</sup>. This translates into $`3\sigma `$ upper limits on the apparent X-ray flux and X-ray luminosity of $`4.21\times 10^{15}`$ ergs cm<sup>-2</sup> s<sup>-1</sup> and $`1.7\times 10^{25}`$ ergs s<sup>-1</sup>, respectively, for VB 10 outside of flare. These are upper limits on any quiescent emission, if it exists.
All of the numbers presented in this section have been tabulated in Table 1.
## 3 Discussion
This X-ray flare on VB 10 is reminiscent of the far-UV flare which was detected on VB 10 by Linsky et al. (1995). These authors observed C II,IV and Si IV emission lines only during the last five minutes of an hour-long exposure taken with the GHRS onboard HST. They concluded that the flare which they had observed indicated increased magnetic heating rates for low-mass stars near the hydrogen-burning mass limit.
In both the UV and ROSAT X-ray observations, quiescent (i.e. non-flare) emission was never detected from VB 10. For the UV flare, the emission line fluxes were an order of magnitude greater than the upper limits placed on the non-flare emission by Linsky et al. (1995). For this most recent X-ray flare, the contrast is even greater. The peak flare luminosity is at least more than two orders of magnitude greater than the non-flare value.
For ease of comparison to more massive M dwarfs, we will normalize L<sub>X</sub> by L<sub>bol</sub>, which for VB 10 is $`1.74\times 10^{30}`$ ergs s<sup>-1</sup> based on the absolute K magnitude measured by Leggett (1992) and the bolometric correction of Veeder (1974). This gives us log (L<sub>X</sub>/L$`{}_{bol}{}^{})>2.8`$ for the peak of the flare and log (L<sub>X</sub>/L$`{}_{bol}{}^{})<5.0`$ outside of flare. Using the data of Fleming et al. (1995) for the late (later than M5), presumably fully-convective M dwarfs within 7 pc of the Sun, we find that these star have values of log (L<sub>X</sub>/L<sub>bol</sub>) which are typically $`3.5`$. But during flares (e.g. AZ Cnc; Fleming et al. 1993), these stars can reach values of (L<sub>X</sub>/L$`{}_{bol}{}^{})=3.0`$ to $`2.5`$. Therefore, the magnitude of the flare on VB 10 is completely consistent with that of flares on more massive M dwarfs.
In Fig. 3 we display a plot of normalized X-ray luminosity versus absolute visual magnitude for all known M dwarfs within 7 pc that are later than M5, based on data from a volume-limited survey by Fleming et al. 1995. Inspection of Fig. 3 reveals the sharp contrast between the upper limit to the non-flare X-ray emission in VB 10 and the X-ray emission levels of other late dMe stars. In particular, the upper limit for VB 10 is 1-2 orders of magnitude less than that detected in the late dMe stars. It is comparable to the upper limit for GJ1002 in Fig. 3, a quiescent dM5.5 star that does not exhibit H$`\alpha `$ emission nor any reported flare activity. By contrast, VB 10 is a known flare star with (variable) H$`\alpha `$ line emission in its spectrum. While one cannot build a theory based on one observation, we cannot help but speculate that our result for VB 10 does indeed reflect a decline in coronal heating efficency near the H-burning mass limit. But somehow these stars are still able to flare.
We do not understand in detail the mechanisms that give rise to energetic, transient outbursts identified as “flares” in the Sun and late-type stars. However, flares do appear to be the result of instabilities that return stressed systems toward configurations that are characterized by lower potential energy (Rosner & Vaiana 1978). Since we do not have an X-ray spectrum, we cannot verify through modeling that the observed emission in VB 10 is consistent with the presence of loop-like magnetic structures. However, the energetics of the event suggest the possible occurrence of a large volume of flare plasma, implying the presence of large-scale magnetic structures in the atmosphere. In particular, we can crudely estimate the plausible range of spatial scales that characterize the flare event in the following manner.
In the absence of an actual energy spectrum (such as the pulse-height spectra that were produced by the now-defunct ROSAT PSPC), we assume some plausible flare plasma parameters. Based on X-ray observations of other flare events on M dwarfs, we adopt a flaring temperature of $`T`$ 10<sup>7</sup> K. We note that a large flare event recorded by the ROSAT PSPC on a late M dwarf star was well-described by a thermal plasma model fit characterized by a temperature of 2 – 4 $`\times `$ 10<sup>7</sup> K (Sun et al. 1999). Utilizing the XSPEC analysis package combined with the observed flux from the HRI and the adopted temperature yields a differential emission measure at the flare maximum of EM $`>1.8\times 10^{29}`$ cm<sup>-5</sup>. The corresponding volume emission measure is
$$VEM=4\pi R_{}^2EM>1.1\times 10^{50}cm^3,$$
$`(1)`$
where R = 0.102 R (Linsky et al. 1995). Given this estimate of volume emission measure along with electron densities in the range of $`n_e`$ 10<sup>10</sup> cm<sup>-3</sup> to 10<sup>11</sup> cm<sup>-3</sup>, we find characteristic linear dimensions for the flare in the range of 0.003 R to 0.30 R. Thus, the flaring plasma covers a large fraction of the stellar surface if $`n_e`$ 10<sup>10</sup> cm<sup>-3</sup>, but only a small fraction if $`n_e>`$ 10<sup>11</sup> cm<sup>-3</sup>. We note that the distance traveled by a sound wave in a 10<sup>7</sup> K plasma in the 3 minute duration of the flare maximum is about 0.7 R. These estimates, while not conclusive, are consistent with the likely occurrence of large-scale magnetic structures associated with the flare event on VB 10.
The enormous contrast between the X-ray flare luminosity and the upper limit to quiescent emission invites further consideration, especially in view of the likely existence of significant surface magnetic flux and large magnetic structures in VB 10. We note that if VB10 is characterized by quiescent X-ray emission at a level which is typical for the Sun at the maximum in its activity cycle, i.e., log($`L_x/L_{bol})`$ -6.3 (following Schmitt 1997), then we would still not have had sufficient sensitivity to detect it. The upper limit for non-flare X-ray emission in VB 10 is more comparable to the levels of emission seen in earlier, more massive M dwarfs (non-dMe flare stars) which are themselves characterized by X-ray emission levels that are in excess of solar, or log($`L_x/L_{bol})`$ -6 to -5 (following Fleming, Schmitt & Giampapa 1995).
We thus cannot exclude the possibility that VB 10 does indeed have undetected, quiescent X-ray emission at the level of the Sun or that of earlier, quiescent dwarf M stars. However, inspection of Fig. 3 suggests that the low level of non-flare X-ray emission in VB 10 is unusual with respect to other late-type, active dMe flare stars. We will therefore briefly consider a hypothesis that may account for the extremely low, or even the possible absence of, steady, quiescent heating in this very cool dwarf.
While there is no comprehensive theory for coronal heating—even in the case of the Sun—a common feature of current theories is that the origin of the nonradiative heating of the corona involves the interaction between turbulent convective motions in the photosphere and the footpoints of magnetic loops (Parker 1972, 1983a,b, 1986; van Ballegooijen 1986 and references therein). In this context, our observations of VB 10 may imply that at the photosphere there is simply not enough in the way of random motions at the loop footpoints to jostle the magnetic fields and, in turn, provide sufficient dynamical stresses to lead to detectable plasma heating to coronal ($`T`$ 10<sup>6</sup> K) temperatures. For example, a high magnetic field strength may suppress convective motions, analogous to the situation in sunspots which are also X-ray quiet. We further note that the very cool and dense photosphere of VB 10 is dominated in content by molecular hydrogen combined with neutral metals bound in molecules. Thus, there are few ions to couple the magnetic field with the photospheric gas. We estimate by extrapolation from M dwarf model photospheres (Mould 1976) to the effective temperature of VB 10 (T<sub>eff</sub> = 2600 K; Linsky et al. 1995) that the ionization fraction, $`\zeta `$, in the dense, upper photosphere of VB 10 is $`\zeta `$ 10<sup>-7</sup>. By contrast, $`\zeta `$ 10<sup>-4</sup> in the upper photosphere of the quiet Sun (Vernazza, Avrett, & Loeser 1976). In early M dwarfs we have that $`\zeta `$ 10<sup>-5</sup> (following Mould 1976). Hence, for stars in the temperature-density regime of VB 10, the upper photospheres are characterized by ionization fractions that are 2 orders of magnitude less than that of earlier M dwarfs and 3 orders of magnitude below that of the Sun. Consequently, the interaction between the field and the ambient gas occurs deeper in the photosphere where the ionization fraction is higher but where sufficient energy to produce significant coronal heating is unable to propagate outward.
The consequences for coronal heating of the occurrence of magnetic structures in very cool dwarfs such as VB 10 is summarized in the following argument due to F. Meyer (1999, private communication; see also Meyer & Meyer-Hofmeister 1999). Given that the gas pressure in the corona is typically negligible compared to the magnetic stresses, we must have that the magnetic field configuration, or “loop”, is force-free, or $`𝐅=(𝐉\times 𝐁)/\mathrm{c}`$ = 0. Thus, the electric currents must flow along the magnetic field lines. Since for steady magnetic phenomena the current density is divergence-free ($`𝐉`$ = 0), currents must also flow through the magnetic footpoints in the photosphere. In the cool and dense photosphere of stars such as VB 10, the electrical conductivity is so low that any current system rapidly decays. From this it follows that the magnetic equilibrium must be current-free everywhere. That is, $`𝐉`$ = 0 so that $`\times 𝐁`$ = 0, and $`𝐁=\mathrm{\Phi }`$. Given that the potential field is a minimum energy configuration, any further build-up and storage of magnetic field energy is excluded, implying no (or only a relatively low degree) of magnetic field-related heating of the atmosphere. At the very least, the decay rate of energy at the footpoints must be faster than any energy input derived from the interaction between motions in the upper photosphere and the magnetic loops, effectively quenching any non-radiative heating that might otherwise have occurred.
Clearly, the above picture as outlined cannot explain flare events. Instead, the transient or flare outbursts observed in these very cool stars must arise from more complex magnetic topologies where the storage of magnetic field energy occurs, but which do not have footpoints in the cool, dense underlying photosphere.
In summary, we have confirmed that the M8 dwarf VB 10 (Gl 752 B) does indeed emit X-rays. It is the lowest mass star on the main sequence which is known to do so. This emission, however, appears to be only transient in nature. The contrast between the flare X-ray flux and any possible non-flare (i.e., quiescent) X-ray emission is at least two orders of magnitude. We note that X-ray emission in VB 10 at solar levels or even at the levels seen in earlier dwarf M (non-dMe) stars would not be detected at our sensitivity limits. In comparison to dMe stars later than M5, the upper limit to the non-flare X-ray emission in VB 10 is unusually low.
We have considered a hypothesis that the low ionization fraction and, hence, low conductivity, in the photosphere of VB 10 inhibits coronal heating to the point where quiescent, $`10^6`$ K coronal plasma may not even exist. In this scenario, the transient, presumably $`10^7`$ K coronal plasma which gives rise to the observed X-ray flare would then be associated with topologically complex magnetic field structures that do not have footpoints in the cool photosphere which is, in turn, dominated by neutral atomic and molecular species. Should this scenario prove to be correct, then there may indeed be a drop in coronal activity at the bottom of the main sequence: not at the mass where stars become fully convective, as was once suggested, but at the hydrogen-burning mass limit itself.
We acknowledge insightful discussions with F. Meyer, B. Durney and A. van Ballegooijen whose ideas materially contributed to this work. TAF acknowledges support from NASA grant NAGW-3160. MSG also acknowledges support from NASA under the ROSAT Guest Observer program. The ROSAT project is supported by the German Bundesministerium für Forschung und Technologie (BMFT/DARA) and the Max Planck Gesellschaft.
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# Evidence for coexistence of the superconducting gap and the pseudo - gap in Bi-2212 from intrinsic tunneling spectroscopy
## Abstract
We present intrinsic tunneling spectroscopy measurements on small Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+x</sub> mesas. The tunnel conductance curves show both sharp peaks at the superconducting gap voltage and broad humps representing the $`c`$-axis pseudo-gap. The superconducting gap vanishes at $`T_c`$, while the pseudo-gap exists both above and below $`T_c`$. Our observation implies that the superconducting and pseudo-gaps represent different coexisting phenomena.
PACS numbers: 74.25.-q, 74.50.+r, 74.72.Hs, 74.80.Dm
The existence of a pseudo-gap (PG) in the quasiparticle density of states (DOS) in the normal state of high-$`T_c`$ superconductors (HTSC) has been revealed by different experimental techniques . For a review, see Refs. . Up-to-date, there is no consensus about the origin of the PG, the correlation between the superconducting gap (SG) and PG, or the dependencies of both gaps on material and experimental parameters. A clarification of these issues is certainly important for understanding HTSC.
From surface tunneling experiments, it was concluded that the SG is almost temperature independent . At $`T>T_c`$, it continuously evolves into the PG, which can persist up to room temperature. Furthermore, it was observed that such a superconducting gap has no correlation with $`T_c`$ and continues to increase in underdoped samples despite a reduction of $`T_c`$ . This was the basis for a suggestion that the PG-state at $`T>T_c`$ is a precursor of superconductivity . On the other hand, surface tunneling into HTSC has several drawbacks, e.g. it is sensitive to surface deterioration. The growing controversy requires further studies with alternative techniques.
Intrinsic tunneling spectroscopy has become a powerful tool in studying the quasiparticle DOS inside bulk single crystals of layered HTSC and thus avoiding the sensitivity to surface deterioration. First experiments were recently attempted to study the PG in mesas fabricated on surface of Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+x</sub> (Bi-2212) single crystals. Unfortunately, intrinsic tunneling experiments also have several problems, such as internal heating and stacking faults (defects) in the mesas. To reduce overheating, a pulse technique was applied by Suzuki et al. The result at $`TT_c`$ was essentially similar to the surface measurements , leaving the obscure relationship between SG and PG unresolved.
In this paper we present results of intrinsic tunneling spectroscopy for Bi-2212 mesas with considerably smaller areas, compared to previous studies . Smaller areas allowed us to avoid stacking faults in the mesas and to avoid mixing between the $`c`$\- axis and $`ab`$\- plane transport. As a result, clean and clear tunnel-type current-voltage (I-V) characteristics were observed, which allowed us to distinguish superconducting and pseudo - gaps in a wide range of temperatures. In contrast to surface and earlier intrinsic tunneling experiments , we have clearly traced different behaviors of the SG and PG. Thorough studies of I-V curves close to $`T_c`$ revealed that the superconducting gap does vanish, while the PG does not change at $`T=T_c`$. All this speaks in favor of different origins of the two coexisting phenomena and against the precursor-superconductivity scenario of the PG. Finally, we discuss interplay between Coulomb interaction and low dimensionality as a possible mechanism for the c-axis PG in an inherent two-dimensional (2D) system, such as the Bi-2212 single crystals.
Mesas with different dimensions from 2 to 20 $`\mu `$m were fabricated simultaneously on top of Bi-2212 single crystals. To reduce the mesa area we adopted a self-alignment technique, see Ref. for details of sample fabrication. The c-axis I-V characteristics were measured in a three-probe configuration. The contact resistance was small, about two orders of magnitude less than the total resistance of the mesa at the corresponding current. All the leads to a mesa were filtered from high-frequency electrical noise. Altogether, more than 50 mesas made on different crystals were investigated. Parameters of the mesas are listed in Table 1. The three figures in the mesa number represent the batch number, the crystal number and the number of the mesa on the crystal, respectively. Letters ”Ar” or ”Ch” indicate whether the mesa was made by Ar-ion or chemical etching. Here we present results for slightly overdoped ($`T_c=89`$ K) and optimally doped ($`T_c=9394`$ K) samples. The pristine crystals were slightly overdoped. Overdoped mesas were obtained by wet chemical etching, which does not significantly change the oxygen content. Optimally doped mesas were made by Ar-ion etching, during which mesas partly loose oxygen. Such mesas had larger $`T_c`$, $`c`$-axis resistivity, $`\rho _c`$ see Table 1, and pseudo-gap, see Fig. 4.
Our fabrication procedure provides samples with highly reproducible properties. This is illustrated in the
top inset of Fig. 1, in which current density, $`j=I/S`$, vs. voltage per junction, $`v=V/N`$, curves at $`T=4.2`$ K are shown for three mesas with different areas from different batches and crystals. Here $`S`$ is the area and $`N`$ is the number of intrinsic Josephson junctions (IJJ’s) in the mesa. We note, that all normalized I-V curves collapse into a single curve. In the following, we will denote quantities corresponding to the whole mesa by capital letters, and those related to an individual IJJ, by small letters. The subscripts ”s”, ”pg” and ”n” will correspond to the superconducting, pseudo-gap and normal state properties, respectively.
In Fig. 1, $`Iv`$ curves and in Fig. 2 the voltage dependence of the dynamic conductance $`\sigma (v)=\mathrm{d}I/\mathrm{d}v(v)`$ are shown for the optimally doped mesa 423Ar at different temperatures. Figs. 1 and 2 exhibit a typical tunnel-junction behavior. At large bias current, there is a well defined normal-state part of tunneling I-V curves with tunnel resistance $`R_n`$. $`R_n`$ decreases by merely $``$ 15% from 300 K to 4.2 K and has no feature at $`T=T_c`$, as shown in the bottom inset of Fig. 1. This is in accordance with the pure tunnel junction behavior, for which $`R_n`$ is expected to be temperature independent. The weak $`T`$-dependence of $`R_n`$ indicates an absence of mixing between $`c`$-axis and $`ab`$-plane transport in our mesas. Previously, however, a strong change of $`R_n`$ at $`T=T_c`$ has been reported for larger mesas . On the other hand, the zero bias resistance, $`R_0`$, has a strong temperature dependence, see bottom inset in Fig. 1. Below $`T_c`$, $`R_0`$ is determined by the sub-gap resistance of the first IJJ, $`R_0^1`$, At $`T<`$40K, a small critical current in the first IJJ appears, see Fig. 3 a), and $`R_0`$ drops to the contact
resistance. Such a two stage decrease of $`R_0`$ is due to a deterioration of IJJ’s at the surface of the mesa .
At low $`T`$, there is a sharp peak in $`\sigma (v)`$, which we attribute to the superconducting gap voltage, $`v_s=2\mathrm{\Delta }_s/e`$. With increasing $`T`$, the peak at $`v_s`$ reduces in amplitude and shifts to lower voltages, reflecting the decrease in $`\mathrm{\Delta }_s(T)`$. At $`T83\mathrm{K}`$ $`(<T_c93`$ K), the superconducting peak is smeared out completely and only a smooth depletion of $`\sigma (0)`$ (a dip) plus a hump in conductance at $`v=v_{pg}70`$ mV remain. The dip and the hump are correlated to each other and both flatten simultaneously with increasing $`T`$, see inset in Fig. 2. Therefore, both reflect the existence of the pseudo-gap in the tunneling DOS. The $`\sigma (0)`$ gradually increases with temperature but the I-V curves remain non-linear nearly up to room temperature. At $`T>T_c`$, the zero-bias resistance, $`R_0`$, is fairly well described by the thermal-activation formula,
$$R_0exp(T^{}/T),T^{}150\pm 20K,$$
(1)
as shown by the dashed line in bottom inset of Fig. 1.
In agreement with surface tunneling experiments, there are no sharp changes at $`T_c`$. As shown in bottom inset of Fig. 1, at $`T<T_c`$, $`R_0`$ evolves continuously into the total (all $`N`$ junctions in the resistive state) sub-gap resistance, $`R_0^N`$. This implies that the PG persists also in the superconducting state. The gradual evolution of the PG hump upon cooling through $`T_c`$ is most clearly shown in inset of Fig. 2. It is seen that the PG dip/hump feature does not change qualitatively upon cooling through $`T_c`$. Moreover, the I-V curve at $`T=`$77.7 K shows that the superconducting peak at $`v_s`$ emerges on top of the PG - features which demonstrates a coexistence of both SG and PG features. From Fig. 2 it is seen that by further decreasing the temperature, the superconducting peak
shifts to higher voltages, increases in amplitude and eventually the PG hump is washed out by the much stronger superconducting peak. For optimally doped mesas, the PG hump can be resolved at $`T>`$60 K, i.e. well below $`T_c`$. The gradual opening of $`\mathrm{\Delta }_s`$ at $`T<T_c`$, in addition to the PG, can also be seen from a steeper growth of the total sub-gap resistance, $`R_0^N`$, at $`T<T_c`$, as compared to the thermal-activation behavior of $`R_0`$ at $`TT_c`$, as shown in bottom inset in Fig.1.
At low bias and $`T<T_c`$, multiple quasiparticle (QP) branches are seen in the I-V curves, representing a one-by-one switching of the IJJ’s into the resistive state . A detailed view of multiple QP branches is shown in Fig. 3 for different $`T`$. Dots and thin lines in Fig. 3 a) represent the experimental points and a polynomial fit, correspondingly. Only the last branch, having many data-points, was actually fitted, all the other thin lines were obtained by dividing the voltages of this fit, $`V_{\mathrm{fit}}(I)`$, by the integer number $`N^{}=Nn+1`$, where $`n`$ is the number of IJJ’s in the resistive state. A good scaling of QP branches is seen, which implies that there is no significant overheating of the mesa at the operational current. If there were overheating, $`V_{\mathrm{fit}}(I)/N^{}`$ would not go
through the data points because switching of additional IJJ’s would cause a progressive increase of the internal temperature and the branches with increasing count numbers would have lower voltages due to the strong temperature dependence of $`R_0^N`$ and $`\mathrm{\Delta }_s`$.
The separation between QP branches, $`\delta v_s`$, is the additional quantity, provided by intrinsic tunneling spectroscopy, which can be used to estimate $`\mathrm{\Delta }_s`$ in a wider range of temperatures. From Fig. 3 b) it is seen that multiple QP branches are clearly distinguishable up to $`TT_c2`$ K. From Table 1 it is seen that $`\delta v_s`$ scales with $`\mathrm{\Delta }_s`$. The $`\delta v_s`$ is less than $`V_s/N`$ simply because the critical current, $`I_c`$, is less than $`V_s/R_n`$ and all IJJ’s switch to the resistive state before they reach the gap voltage, see Fig. 1. The $`\delta v_s(I=I_c)`$ continuously decreases with $`T`$ and vanishes at $`T_c`$. In principle, the temperature dependence of $`I_c`$ is also involved in $`\delta v_s(T)`$, since we measure $`\delta v_s`$ at $`II_c`$. However, the inset in Fig.3 b) reveals that $`\delta v_s(T)`$ still tends to vanish at $`TT_c`$ even if we evaluate $`\delta v_s`$ at one and the same current for all $`T`$.
In Fig. 4, the temperature dependencies of the superconducting peaks, $`v_s`$ (squares), $`\delta v_s(I=I_c)`$ (circles), and pseudo-gap humps, $`v_{pg}`$ (triangles), are shown for optimally doped (solid) and overdoped (open symbols) samples. Small solid symbols represent $`v_s`$ for the rest of the mesas listed in Table 1, and the lines are guides for the eye. In agreement with previous studies, both $`v_s`$ and $`v_{pg}`$ increase upon going from overdoped to optimally doped samples. The superconducting gap deduced from the sum-gap voltage $`V_s=2N\mathrm{\Delta }_s/e`$ is $`\mathrm{\Delta }_s(4.2K)33`$ meV for the optimally doped sample, and $`26`$ meV for the overdoped one. In contrast to surface tunneling experiments, we observe that $`\mathrm{\Delta }_s`$ decreases considerably with temperature. The robust decrease of $`\mathrm{\Delta }_s(T)`$ from 4.2K to $`T_c`$ is more than 80 % for the overdoped mesas. Moreover, we can measure $`\delta v_s(I=I_c)`$ in a wider range of $`T`$ and observe that it vanishes at $`TT_c`$.
All this brings us to the conclusion that the superconducting gap does close at $`T_c`$, in agreement with the previous observations of vanishing of the superfluid density (divergence of the magnetic penetration depth) and the Josephson plasma frequency . On the contrary, the PG is almost temperature independent and exists both above and below $`T_c`$. Therefore, the SG is not developing from the PG, and these two gaps represent different coexisting phenomena. The recently observed independence of the PG on magnetic field supports our conclusion and also casts doubts about the precursor-superconductivity origin of the PG .
One possible ”non-superconducting” PG-scenario is the formation of charge or spin density waves (CDW or SDW) . HTSC’s are composed of quasi-two dimensional electronic systems with a certain degree of Fermi-surface nesting , which can make the system unstable with respect to CDW or SDW formation . A CDW or SDW is accompanied by a PG in DOS, detectable by a surface-tunneling spectroscopy . Many similarities exist between the PG in CDW or SDW (including ARPES , optical conductivity and NMR ) and the PG in HTSC. On the other hand, an opening of the gap due to CDW or SDW is typically accompanied by a metal-insulator transition , while the $`ab`$-plane resistivity in Bi-2212 shows the opposite tendency .
We would also like to emphasize a similarity between the PG features of $`c`$-axis tunneling in HTSC and Coulomb PG for tunneling into a two-dimensional electron system (2DES). The Coulomb PG in 2DES is well studied in connection with semiconducting heterostructures . Experimental $`\sigma (v)`$ curves from the inset in Fig. 2 are strikingly similar to ”V-shaped” tunneling characteristics of 2DES . Certainly, the electron system in Bi-2212 is highly two-dimensional. Moreover, a Coulomb origin of the HTSC pseudo-gap would naturally explain the increase of PG with decreasing O-doping and carrier concentration. A large Coulomb PG in low conducting 2DES is due to unscreened long-range Coulomb interaction and/or slow charge accommodation . Large PG could also appear if tunneling occurs via intermediate low conducting BiO layers .
An attractive feature of both CDW/SDW and Coulomb PG scenaria is that the PG can persist in the superconducting state. Below $`T_c`$, SG and PG are combined into a larger overall gap. This is in agreement with a definite trend for the increase of $`v_{pg}`$ at $`T<T_c`$, see Fig. 4. This might also help in understanding of large ”superconducting” gaps seen in underdoped HTSC . Whether the CDW/SDW or Coulomb PG scenaria can explain all PG features in HTSC remains to be clarified.
In Conclusion, small mesa structures were used for intrinsic tunneling spectroscopy of Bi-2212. We were able to distinguish and simultaneously observe both superconducting and pseudo gaps in a wide range of temperatures. The superconducting gap has a strong temperature dependence and vanishes at $`T_c`$, while the pseudo-gap is almost temperature independent and exists both above and below $`T_c`$. This suggests that the pseudo-gap is not directly related to superconductivity.
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# Untitled Document
RI-1-00
CERN-TH/2000-052
hep-th/0002104
Superstring Theory on $`AdS_3\times G/H`$ and Boundary $`N=3`$ Superconformal Symmetry
Riccardo Argurio<sup>1</sup>, Amit Giveon<sup>1,2</sup> and Assaf Shomer<sup>1</sup>
<sup>1</sup> Racah Institute of Physics
The Hebrew University
Jerusalem 91904, Israel
<sup>2</sup> Theory Division, CERN
CH-1211, Geneva 23, Switzerland
argurio, shomer@cc.huji.ac.il, giveon@vms.huji.ac.il
Superstrings propagating on backgrounds of the form $`AdS_3\times G/H`$ are studied using the coset CFT approach. We focus on seven dimensional cosets which have a semiclassical limit, and which give rise to $`N=3`$ superconformal symmetry in the dual CFT. This is realized for the two cases $`AdS_3\times SU\left(3\right)/U\left(1\right)`$ and $`AdS_3\times SO\left(5\right)/SO\left(3\right)`$, for which we present an explicit construction. We also provide sufficient conditions on a CFT background to enable a similar construction, and comment on the geometrical interpretation of our results.
2/00
1. Introduction
String propagation on curved backgrounds with an $`AdS_3`$ factor has been of recent interest (see for instance \[1--6\], and for additional references). One motivation is the fact that $`AdS_3SL(2)`$ is an exact background which can be treated in string pertubation theory, and thus allows to consider the $`AdS`$/CFT correspondence beyond the supergravity limit. Some specific examples that were studied in this context include superstrings propagating on $`AdS_3\times 𝒩`$ where $`𝒩`$ was a group manifold \[3,,9\], or an orbifold of a group manifold \[10,,11\]. In this paper we study cases in which $`𝒩`$ is a coset manifold. This is an interesting generalization of the $`AdS`$/CFT correspondence which has been considered in the higher dimensional cases of type IIB string theory on $`AdS_5\times 𝒩^5`$ and of M-theory on $`AdS_4\times 𝒩^7`$ \[13,,14\], where $`𝒩^5`$ and $`𝒩^7`$ are Einstein manifolds (generically coset manifolds) preserving a fraction of supersymmetry. This type of construction allows one to consider dual supersymmetric CFTs which are not ‘orbifolds’ of the maximally supersymmetric one. The $`AdS_3\times 𝒩`$ case is somewhat different since here we have the possibility of studying $`𝒩`$ in the context of coset CFTs.
We choose to study coset CFTs which have a large radius (or large level $`k`$) semiclassical limit, corresponding to superstrings propagating on seven dimensional coset manifolds. Moreover, we focus on cases in which the dual two dimensional theory (also referred to as the spacetime CFT) has an extended superconformal symmetry.<sup>1</sup> In the following we refer to the supersymmetry of, say, the left-movers only. The supersymmetry of the other sector depends on the particular superstring theory considered. Coset models leading to $`N=2`$ can be easily realized as particular cases of the general construction of \[15\], where $`𝒩`$ decomposes as a $`U(1)`$ factor times a Kazama-Suzuki model . On the other hand, there are no seven dimensional coset manifolds leading to $`N=4`$ supersymmetry in spacetime (except of course the cases \[3,,9\] in which the cosets are actually group manifolds). Therefore, we shall be interested in the cases where the spacetime CFT has $`N=3`$ supersymmetry.
Our main result, presented in sections 3 and 4, is the construction of the spacetime $`N=3`$ superconformal algebra in the two cases:
$$AdS_3\times \frac{SU(3)}{U(1)},AdS_3\times \frac{SO(5)}{SO(3)},$$
which actually turn out to be the only coset models giving rise to $`N=3`$. An interesting by-product of the construction is related to the fact that getting $`N=3`$ depends on the choice of chiral GSO projection. Choosing the opposite projection leads to an $`N=1`$ superconformal algebra in spacetime together with an $`SU(2)`$ affine algebra acting trivially on the supercharges. This spacetime structure also appears in , where a $`Z_2`$ orbifold of the large $`N=4`$ algebra obtained by is taken.
We also go beyond the coset set-up by providing in section 6 a set of sufficient conditions on a CFT background $`𝒩`$ which allow for the construction of $`N=3`$ superconformal symmetry in spacetime from superstrings on $`AdS_3\times 𝒩`$. We recover the $`N=1`$ structure for the other GSO projection also in this general set-up. The proof elaborates on the construction of $`N=2`$ spacetime supersymmetry by , the additional ingredient being the presence in the CFT on $`𝒩`$ of an affine $`SU(2)`$. Finally, we comment in section 7 on a possible geometrical interpretation of these conditions, and relate our work to the case of M-theory compactified on $`AdS_4\times 𝒩^7`$ \[13,,14\], as well as to brane configurations.
2. Spacetime $`N=3`$ superconformal algebra
Extended superconformal algebras in two dimensions also include an affine R-symmetry algebra, which generally leads to a quantization of the central charge in unitary theories. Specifically, the $`N=3`$ superconformal algebra has an affine $`SU(2)`$ subalgebra. The central charge is related to the level $`\stackrel{~}{k}`$ of this affine $`SU(2)`$, which is an integer, by $`\stackrel{~}{c}=\frac{3}{2}\stackrel{~}{k}`$ . Therefore, a necessary condition for string theory on a background of the form $`AdS_3\times 𝒩`$ to have spacetime $`N=3`$ superconformal symmetry is the existence of an affine $`SU(2)`$ in spacetime. This is obtained when the worldsheet CFT on $`𝒩`$ has an affine $`SU(2)`$ symmetry as well . If the respective worldsheet levels of $`SL(2)`$ and of $`SU(2)`$ are $`k`$ and $`k^{}`$, the analysis of \[3,,5\] shows that in the spacetime theory we have $`\stackrel{~}{c}=6kp`$ and $`\stackrel{~}{k}=k^{}p`$, where $`p`$ is the integer number introduced in , related to the maximal number of ‘long strings’ \[5,,18\]. A further condition is thus $`k^{}=4k`$ (recall that $`k`$ is not forced to be an integer).
In the following we focus on coset manifolds $`𝒩`$ which have 7 dimensions, so that a large $`k`$ semi-classical limit is possible. Two cases which satisfy the conditions given above are $`SL(2)_k\times 𝒩`$ with:
$$𝒩_1=\frac{SU(3)_{4k}}{U(1)},$$
and
$$𝒩_2=\frac{SO(5)_{4k}}{SO(3)}.$$
Note that there are several ways of choosing the $`SO(3)`$ in $`𝒩_2`$, according to the nesting of subgroups $`SO(3)SO(3)\times SO(3)SO(4)SO(5)`$. Since we require $`𝒩_2`$ to have an unbroken $`SO(3)`$ symmetry, we are forced to mod out by one of the two $`SO(3)`$ factors of $`SO(4)`$.<sup>2</sup> Note that modding out by the diagonal $`SO(3)`$ and then by a further $`U(1)`$ would lead to a Kazama-Suzuki model . It is straightforward to show that the two models above are critical:
$$c_{sl}+c_1=\left(\frac{9}{2}+\frac{6}{k}\right)+\left(12\frac{24}{4k}\frac{3}{2}\right)=15,$$
$$c_{sl}+c_2=\left(\frac{9}{2}+\frac{6}{k}\right)+\left(15\frac{30}{4k}\frac{9}{2}+\frac{6}{4k}\right)=15.$$
We now show that these two models indeed possess $`N=3`$ superconformal symmetry in spacetime by explicit construction. Since the construction is similar in the two cases, we will focus here on the first case, and then go briefly over the second one.
3. Superstring theory on $`AdS_3\times SU(3)/U(1)`$
We first have to set some notations, starting from the $`SL(2)`$ WZW part. We mainly follow the formalism of and . For simplicity we only treat the holomorphic sector.
The $`SL(2)`$ supersymmetric WZW model is constituted of the three currents of the $`SL(2)`$ affine algebra at level $`k`$, and the three fermions implied by the $`N=1`$ worldsheet supersymmetry, satisfying the following OPEs:
$$\begin{array}{cc}\hfill J^P(z)J^Q(w)& \frac{\frac{k}{2}\eta ^{PQ}}{(zw)^2}+\frac{iϵ^{PQR}\eta _{RS}J^S(w)}{zw},\hfill \\ \hfill J^P(z)\psi ^Q(w)& \frac{iϵ^{PQR}\eta _{RS}\psi ^S(w)}{zw},\hfill \\ \hfill \psi ^P(z)\psi ^Q(w)& \frac{\frac{k}{2}\eta ^{PQ}}{zw},\hfill \end{array}$$
where $`P,Q,R,S=1,2,3`$, $`\eta ^{PQ}=(++)`$ and $`ϵ^{123}=1`$. As usual in supersymmetric WZW models, the currents can be decomposed in two pieces:
$$J^P=\widehat{J}^P\frac{i}{k}\eta ^{PQ}ϵ_{QRS}\psi ^R\psi ^S.$$
The first piece $`\widehat{J}^P`$ constitutes an affine algebra at level $`k+2`$, and has regular OPE with the fermions $`\psi ^P`$. We will thus refer to $`\widehat{J}^P`$ as the bosonic currents. The second part constitutes an affine algebra at level $`2`$, and is referred to as the fermionic part of the current.
The worldsheet stress-energy tensor and $`N=1`$ supercurrent are:
$$\begin{array}{cc}\hfill T_{sl}=& \frac{1}{k}\left(\widehat{J}^P\widehat{J}_P\psi ^P\psi _P\right),\hfill \\ \hfill G_{sl}=& \frac{2}{k}\left(\psi ^P\widehat{J}_P\frac{i}{3k}ϵ_{PQR}\psi ^P\psi ^Q\psi ^R\right).\hfill \end{array}$$
Let us now turn to the $`SU(3)/U(1)`$ coset CFT. We start from the $`SU(3)`$ affine superalgebra at level $`k^{}=4k`$ realized as follows:
$$\begin{array}{cc}\hfill K^A(z)K^B(w)& \frac{\frac{k^{}}{2}\delta ^{AB}}{(zw)^2}+\frac{if_{ABC}K^C(w)}{zw},\hfill \\ \hfill K^A(z)\chi ^B(w)& \frac{if_{ABC}\chi ^C(w)}{zw},\hfill \\ \hfill \chi ^A(z)\chi ^B(w)& \frac{\frac{k^{}}{2}\delta ^{AB}}{zw}.\hfill \end{array}$$
Here $`A,B,C,D=1\mathrm{}8`$ and the structure constants $`f_{ABC}`$ are $`f_{123}=1`$, $`f_{147}=f_{156}=f_{246}=f_{257}=f_{345}=f_{367}=\frac{1}{2}`$ and $`f_{458}=f_{678}=\frac{\sqrt{3}}{2}`$. Since the metric is $`\delta ^{AB}`$ we will not keep track of the upper or lower position of the $`SU(3)`$ indices. As before, we split the currents into their bosonic and fermionic parts:
$$K^A=\widehat{K}^A\frac{i}{k^{}}f_{ABC}\chi ^B\chi ^C.$$
The bosonic currents realize an affine algebra at level $`k^{}3`$.
We now choose to mod out the $`SU(3)`$ by the $`U(1)`$ generated by $`K^8`$. The $`SU(2)`$ subgroup generated by $`K^1,K^2,K^3`$ is orthogonal to this $`U(1)`$, and thus survives as an affine algebra in the coset CFT. The stress-energy tensor and the supercurrent of the coset CFT are built as in , using the decomposition $`T_{SU(3)}=T_{SU(3)/U(1)}+T_{U(1)}`$, and similarly for the supercurrent $`G`$. The stress-energy tensor reads:
$$\begin{array}{cc}\hfill T_{coset}=& \frac{1}{k^{}}\left(\widehat{K}^1\widehat{K}^1+\mathrm{}+\widehat{K}^7\widehat{K}^7\right)\frac{1}{k^{}}\left(\chi ^1\chi ^1+\mathrm{}+\chi ^3\chi ^3\right)\hfill \\ & \frac{1}{k^{}}\left(1\frac{3}{2k^{}}\right)\left(\chi ^4\chi ^4+\mathrm{}+\chi ^7\chi ^7\right)+\frac{2i\sqrt{3}}{k_{}^{}{}_{}{}^{2}}\widehat{K}^8\left(\chi ^4\chi ^5+\chi ^6\chi ^7\right)\hfill \\ & +\frac{6}{k_{}^{}{}_{}{}^{3}}\chi ^4\chi ^5\chi ^6\chi ^7.\hfill \end{array}$$
Our goal now is to build the spacetime supercharges. For that we would like to construct spin-fields via bosonization following . Note that since we are dealing with a coset and not with a group manifold, the fermions are generically not free. Of course since the $`SU(2)`$ is preserved as an affine symmetry, the fermions belonging to it are free. Despite the above remark, we proceed to bosonize the 10 fermions into 5 bosons. This will actually uncover the interesting structure of the above coset model. Define:
$$H_1=\frac{2}{k}\psi ^1\psi ^2,H_2=\frac{2}{k^{}}\chi ^1\chi ^2,iH_3=\frac{1}{k}\psi ^3\chi ^3,H_4=\frac{2}{k^{}}\chi ^4\chi ^5,H_5=\frac{2}{k^{}}\chi ^6\chi ^7.$$
The scalars $`H_I`$ are all canonically normalized: $`H_I(z)H_J(w)\delta _{IJ}\mathrm{log}(zw)`$. Conversely, the fermions are given by:
$$\psi ^1=\frac{\sqrt{k}}{2}\left(e^{iH_1}+e^{iH_1}\right),\psi ^2=\frac{i\sqrt{k}}{2}\left(e^{iH_1}e^{iH_1}\right),$$
and similarly for $`H_2`$, $`H_4`$ and $`H_5`$, while:
$$\psi ^3=\frac{\sqrt{k}}{2}\left(e^{iH_3}e^{iH_3}\right),\chi ^3=\frac{\sqrt{k^{}}}{2}\left(e^{iH_3}+e^{iH_3}\right),$$
recalling that $`H_3^{}=H_3`$ and $`k^{}=4k`$.
In terms of these scalars, the total stress-energy tensor is:
$$\begin{array}{cc}\hfill T=T_{sl}+T_{coset}=& \frac{1}{k}\left(\widehat{J}^1\widehat{J}^1+\widehat{J}^2\widehat{J}^2\widehat{J}^3\widehat{J}^3\right)+\frac{1}{k^{}}\left(\widehat{K}^1\widehat{K}^1+\mathrm{}+\widehat{K}^7\widehat{K}^7\right)\hfill \\ & \frac{1}{2}\left(H_1H_1+H_2H_2+H_3H_3\right)+\frac{i\sqrt{3}}{k^{}}\widehat{K}^8\left(H_4+H_5\right)\hfill \\ & \frac{1}{2}\left(1\frac{3}{2k^{}}\right)\left(H_4H_4+H_5H_5\right)+\frac{3}{2k^{}}H_4H_5.\hfill \end{array}$$
Obviously, the scalars $`H_4`$ and $`H_5`$ are not free in the coset CFT. However, it is also easy to see that there is a linear combination of them which is free. This is what will enable us to build the $`N=3`$ spacetime superalgebra.
We thus write:
$$H_\pm =\frac{1}{\sqrt{2}}(H_4\pm H_5).$$
Our final expression for $`T`$ is therefore:
$$\begin{array}{cc}\hfill T=& \frac{1}{k}\left(\widehat{J}^1\widehat{J}^1+\widehat{J}^2\widehat{J}^2\widehat{J}^3\widehat{J}^3\right)+\frac{1}{k^{}}\left(\widehat{K}^1\widehat{K}^1+\mathrm{}+\widehat{K}^7\widehat{K}^7\right)\hfill \\ & \frac{1}{2}\left(H_1H_1+H_2H_2+H_3H_3+H_{}H_{}\right)\hfill \\ & \frac{1}{2}\left(1\frac{3}{k^{}}\right)H_+H_++\frac{i\sqrt{6}}{k^{}}\widehat{K}^8H_+.\hfill \end{array}$$
We conclude that $`H_{}`$ is the fourth free scalar, namely that $`H_{}`$ is a primary field of weight 1.
We now write the worldsheet $`N=1`$ supercurrent, which will be used to enforce the BRST condition on the spin fields. The supercurrent for the coset CFT reads:
$$G_{coset}=\frac{2}{k^{}}\left(\chi ^{\overline{a}}\widehat{K}^{\overline{a}}\frac{i}{3k^{}}f_{\overline{a}\overline{b}\overline{c}}\chi ^{\overline{a}}\chi ^{\overline{b}}\chi ^{\overline{c}}\right),$$
where $`\overline{a}`$ are indices in the coset $`G/H`$. Putting together (3.1) and (3.1), substituting the structure constants of $`SU(3)`$ and taking into account the bosonization in the term trilinear in the fermions, we get the expression for $`G_{tot}=G_{sl}+G_{coset}`$:
$$\begin{array}{cc}\hfill G_{tot}=& \frac{2}{k}\left(\psi ^1\widehat{J}^1+\mathrm{}\psi ^3\widehat{J}^3\right)+\frac{2}{k^{}}\left(\chi ^1\widehat{K}^1+\mathrm{}+\chi ^7\widehat{K}^7\right)\hfill \\ & +\frac{i}{\sqrt{k}}\left\{H_1\left(e^{iH_3}e^{iH_3}\right)\frac{1}{2}\left(H_2+\frac{1}{\sqrt{2}}H_{}\right)\left(e^{iH_3}+e^{iH_3}\right)\right\}\hfill \\ & +\frac{1}{2\sqrt{k}}\left(e^{iH_2i\sqrt{2}H_{}}e^{iH_2+i\sqrt{2}H_{}}\right).\hfill \end{array}$$
Before going on to the BRST condition for the spin-fields, we write for completeness the expressions for the $`SU(2)`$ currents, which remain primary fields of weight 1 in the coset CFT and can also be considered as the upper components of the fermions $`\chi ^1,\chi ^2,\chi ^3`$. Writing $`K^\pm =K^1\pm iK^2`$ and similarly for the bosonic currents and the fermions, we get:
$$\begin{array}{cc}\hfill K^\pm =& \widehat{K}^\pm e^{iH_2}\left(e^{iH_3}+e^{iH_3}\right)\pm e^{i\sqrt{2}H_{}}\hfill \\ \hfill K^3=& \widehat{K}^3i\left(H_2+\frac{1}{\sqrt{2}}H_{}\right).\hfill \end{array}$$
Note that since these currents are primaries of weight 1, this could have been an alternative way of showing that $`H_{}`$ is a free scalar.
3.1. Physical operators and the spacetime algebra
In order to construct the spacetime superconformal algebra we need, in particular, to construct physical supercharges which we choose to write in the $`1/2`$ picture :
$$Qe^{\phi /2}u^\alpha S_\alpha (z)𝑑z.$$
Here $`S_\alpha `$ is a basis of spin-fields, $`u^\alpha `$ are constants, and $`\phi `$ is the bosonized superconformal ghost. The set of operators $`e^{\phi /2}u^\alpha S_\alpha (z)`$ should be BRST invariant and mutually local. We choose a basis of spin-fields
$$S_{[ϵ_1ϵ_2ϵ_3ϵ_{}]}=e^{\frac{i}{2}(ϵ_1H_1+ϵ_2H_2+ϵ_3H_3+ϵ_{}\sqrt{2}H_{})},$$
where $`ϵ_I=\pm 1`$. Because $`H_{}`$ is a free scalar, these 16 spin-fields are primaries of weight 5/8 and, therefore, $`e^{\phi /2}u^\alpha S_\alpha (z)`$ are primaries of weight 1, as they should be.
The super BRST condition on $`e^{\phi /2}u^\alpha S_\alpha `$ further requires that there will be no $`(zw)^{3/2}`$ singular terms in the OPE of $`u^\alpha S_\alpha `$ with the supercurrent $`G_{tot}`$ (note that the only dangerous terms in $`G_{tot}`$ are the ones trilinear in the fermions, i.e. the second and third lines in (3.1)). This leaves 8 combinations $`u^\alpha S_\alpha `$ out of the 16 spin-fields (3.1). The GSO condition, i.e. mutual locality, further leads to one of two choices of chirality: $`ϵ_1ϵ_2ϵ_3=1`$ or $`ϵ_1ϵ_2ϵ_3=1`$, under which 6 or 2 of the combinations $`u^\alpha S_\alpha `$ survive, respectively.
Explicitly, the outcome of the computation is the following. For spacetime chirality $`ϵ_1ϵ_2ϵ_3=1`$, we get 6 physical spin-fields:
$$\begin{array}{cc}\hfill S_{\frac{1}{2}}^+=& S_{[]}\hfill \\ \hfill S_{\frac{1}{2}}^{}=& S_{[+++]}\hfill \\ \hfill S_{\frac{1}{2}}^3=& \frac{1}{2}(S_{[++]}S_{[+]})\hfill \\ \hfill S_{\frac{1}{2}}^+=& S_{[++]}\hfill \\ \hfill S_{\frac{1}{2}}^{}=& S_{[+++]}\hfill \\ \hfill S_{\frac{1}{2}}^3=& \frac{1}{2}(S_{[++]}S_{[+++]}).\hfill \end{array}$$
The lower and upper labels of $`S_r^a`$ denote respectively the quantum numbers of the global $`SL(2)`$ and $`SU(2)`$ symmetries, in the $`(\mathrm{𝟐},\mathrm{𝟑})`$ representation, as can be checked by taking the OPEs with the respective currents (3.1) and (3.1). For the other spacetime chirality $`ϵ_1ϵ_2ϵ_3=1`$, we get 2 physical spin-fields:
$$\begin{array}{cc}\hfill \stackrel{~}{S}_{\frac{1}{2}}=& \frac{1}{2}(S_{[++]}+S_{[+]})\hfill \\ \hfill \stackrel{~}{S}_{\frac{1}{2}}=& \frac{1}{2}(S_{[+++]}+S_{[++]}).\hfill \end{array}$$
It can be checked that the above spin-fields $`\stackrel{~}{S}_r`$ have regular OPEs with the $`SU(2)`$ currents.
We thus see that the choice of GSO projection will lead to different amounts of supersymmetry in spacetime. Namely, in a type II background, the projection in the left and right moving sectors of the worldsheet CFT determine, respectively, the amount of supersymmetry in the left and right moving sectors of the spacetime CFT. Specifically, the different GSO projections would lead in type IIA to $`N=(3,1)`$ or $`N=(1,3)`$, and in type IIB to $`N=(3,3)`$ or $`N=(1,1)`$. In the heterotic string, the different GSO projections in the worldsheet supersymmetric sector lead to $`N=(3,0)`$ or $`N=(1,0)`$ in spacetime.<sup>3</sup> Note that these examples provide, in particular, a construction of $`N=1`$ supersymmetry in spacetime which is not a $`Z_2`$ orbifolding of the $`N=2`$ construction of . If it was such an orbifold, each of the supercharges would split into two BRST invariant pieces, leading to a total of 4 physical spin-fields, in contrast with the result (3.1).
The generators of the spacetime global $`N=3`$ superconformal algebra are the following:
$$\begin{array}{cc}\hfill L_{\pm 1}=& 𝑑zJ^\pm (z),L_0=𝑑zJ^3(z)\hfill \\ \hfill T_0^\pm =& 𝑑zK^\pm (z),T_0^3=𝑑zK^3(z)\hfill \\ \hfill Q_{\frac{1}{2}}^\pm =& 𝑑ze^{\phi /2}S_{\frac{1}{2}}^\pm (z),Q_{\frac{1}{2}}^3=𝑑ze^{\phi /2}S_{\frac{1}{2}}^3(z)\hfill \\ \hfill Q_{\frac{1}{2}}^\pm =& 𝑑ze^{\phi /2}S_{\frac{1}{2}}^\pm (z),Q_{\frac{1}{2}}^3=𝑑ze^{\phi /2}S_{\frac{1}{2}}^3(z),\hfill \end{array}$$
where we omit the normalization and the cocycle factors in the definition of the $`Q`$’s. These operators close the global part of the $`N=3`$ superconformal algebra (up to picture changing), in the NS sector:
$$\begin{array}{cc}\hfill [L_m,L_n]=& (mn)L_{m+n}\hfill \\ \hfill [T_0^a,T_0^b]=& iϵ^{abc}T_0^c\hfill \\ \hfill [L_m,T_0^a]=& 0\hfill \\ \hfill [L_m,Q_r^a]=& \left(\frac{1}{2}mr\right)Q_{m+r}^a\hfill \\ \hfill [T_0^a,Q_r^b]=& iϵ^{abc}Q_r^c\hfill \\ \hfill \{Q_r^a,Q_s^b\}=& 2\delta ^{ab}L_{r+s}+iϵ^{abc}(rs)T_{r+s}^c,\hfill \end{array}$$
where $`m,n=0,\pm 1`$, $`a,b,c=1,2,3`$ and $`r,s=\pm \frac{1}{2}`$. Of course this model reproduces the full $`N=3`$ superconformal algebra. The higher modes can be built as in \[3,,5\]. For instance, we can first construct all the $`L_n`$. Then acting with them on $`T_0^a`$ and $`Q_{\pm \frac{1}{2}}^a`$ one gets all the $`T_n^a`$ and $`Q_r^a`$ higher modes. To close the algebra an additional fermionic field is needed, all the modes of which are obtained from commutators of $`T_n^a`$ and $`Q_r^a`$. The full algebra appears for example in \[17,,11\].
For completeness, we also write the (global) algebra for the other GSO projection, that is an $`N=1`$ superconformal algebra together with an affine $`SU(2)`$ which acts trivially on the supercharges. The supersymmetry generators are given in this case by:
$$\stackrel{~}{Q}_{\pm \frac{1}{2}}=𝑑ze^{\phi /2}\stackrel{~}{S}_{\pm \frac{1}{2}}(z),$$
and the algebra is:
$$\begin{array}{cc}\hfill [L_m,L_n]=& (mn)L_{m+n}\hfill \\ \hfill [T_0^a,T_0^b]=& iϵ^{abc}T_0^c\hfill \\ \hfill [L_m,T_0^a]=& 0\hfill \\ \hfill [L_m,\stackrel{~}{Q}_r]=& \left(\frac{1}{2}mr\right)\stackrel{~}{Q}_{m+r}\hfill \\ \hfill [T_0^a,\stackrel{~}{Q}_r]=& 0\hfill \\ \hfill \{\stackrel{~}{Q}_r,\stackrel{~}{Q}_s\}=& 2\delta ^{ab}L_{r+s}.\hfill \end{array}$$
Again, using the higher modes of $`L_n`$ one can generate the higher modes of the other operators, together with the fermionic superpartners of the affine currents (see for an analogous construction).
Let us conclude this section by a brief comment on a special case, when the level of the $`SU(3)`$ is $`k^{}=3`$ (the minimal level allowed by unitarity). We can decompose the coset CFT $`[SU(3)/U(1)]`$ into the product $`[SU(2)]\times [SU(3)/(SU(2)\times U(1))]`$. The central charge of the second piece is $`c=0`$ in the $`k^{}=3`$ case, thus the whole model reduces to string propagation on $`SL(2)_{3/4}\times SU(2)_3`$. If we consider the six dimensional model $`SL(2)_k\times SU(2)_k^{}`$, criticality enforces $`\frac{1}{k}\frac{1}{k^{}}=1`$. The only combination which might allow $`N=3`$ in spacetime (i.e. which verifies $`k^{}=4k`$) is the one above, and our analysis indeed shows that it has $`N=3`$. We will comment more on this case later.
4. Superstring theory on $`AdS_3\times SO(5)/SO(3)`$
The construction of the $`N=3`$ superconformal algebra in this case follows closely the steps of the previous section. We shall therefore be more schematical, and focus on the specifics of this model.
The $`SO(5)`$ current algebra looks the same as (3.1), at the same level $`k^{}=4k`$, but now the indices are $`A,B,C,D=1\mathrm{}9,0`$ and the structure constants $`f_{ABC}`$ are $`f_{123}=f_{456}=1`$, $`f_{170}=f_{189}=f_{279}=f_{280}=f_{378}=f_{390}=\frac{1}{2}`$ and $`f_{470}=f_{489}=f_{579}=f_{580}=f_{678}=f_{690}=\frac{1}{2}`$. We work in the basis where the two orthogonal $`SO(3)`$ subgroups of $`SO(5)`$ are generated respectively by $`K^1,K^2,K^3`$ and $`K^4,K^5,K^6`$. We will mod out by the second one, leaving the first one as the R-symmetry.
As before, we straightforwardly bosonize the 10 fermions in the $`SL(2)`$ and in the coset, to get:
$$H_1=\frac{2}{k}\psi ^1\psi ^2,H_2=\frac{2}{k^{}}\chi ^1\chi ^2,iH_3=\frac{1}{k}\psi ^3\chi ^3,H_4=\frac{2}{k^{}}\chi ^7\chi ^8,H_5=\frac{2}{k^{}}\chi ^9\chi ^0.$$
We can now write the total stress-energy tensor in terms of them as (we also use (3.1)):
$$\begin{array}{cc}\hfill T=& \frac{1}{k}\left(\widehat{J}^1\widehat{J}^1+\widehat{J}^2\widehat{J}^2\widehat{J}^3\widehat{J}^3\right)+\frac{1}{k^{}}\left(\widehat{K}^1\widehat{K}^1+\mathrm{}+\widehat{K}^3\widehat{K}^3+\widehat{K}^7\widehat{K}^7+\mathrm{}+\widehat{K}^0\widehat{K}^0\right)\hfill \\ & \frac{1}{2}\left(H_1H_1+H_2H_2+H_3H_3+H_+H_+\right)\frac{1}{2}\left(1\frac{3}{k^{}}\right)H_{}H_{}\hfill \\ & +\frac{i\sqrt{2}}{k^{}}\widehat{K}^6H_{}\frac{1}{k^{}}\widehat{K}^4(e^{i\sqrt{2}H_{}}e^{i\sqrt{2}H_{}})\frac{i}{k^{}}\widehat{K}^5(e^{i\sqrt{2}H_{}}+e^{i\sqrt{2}H_{}}).\hfill \end{array}$$
Now $`H_+`$ is the free scalar.
The analogous expression for the total supercurrent is:
$$\begin{array}{cc}\hfill G_{tot}=& \frac{2}{k}\left(\psi ^1\widehat{J}^1+\mathrm{}\psi ^3\widehat{J}^3\right)+\frac{2}{k^{}}\left(\chi ^1\widehat{K}^1+\mathrm{}+\chi ^3\widehat{K}^3+\chi ^7\widehat{K}^7+\mathrm{}+\chi ^0\widehat{K}^0\right)\hfill \\ & +\frac{i}{\sqrt{k}}\left\{H_1\left(e^{iH_3}e^{iH_3}\right)\frac{1}{2}\left(H_2+\frac{1}{\sqrt{2}}H_+\right)\left(e^{iH_3}+e^{iH_3}\right)\right\}\hfill \\ & \frac{1}{2\sqrt{k}}\left(e^{iH_2i\sqrt{2}H_+}e^{iH_2+i\sqrt{2}H_+}\right).\hfill \end{array}$$
The $`SU(2)`$ currents are:
$$\begin{array}{cc}\hfill K^\pm =& \widehat{K}^\pm e^{iH_2}\left(e^{iH_3}+e^{iH_3}\right)e^{i\sqrt{2}H_+}\hfill \\ \hfill K^3=& \widehat{K}^3i\left(H_2+\frac{1}{\sqrt{2}}H_+\right).\hfill \end{array}$$
The solutions to the BRST invariance conditions on the spin-fields $`S_{[ϵ_1ϵ_2ϵ_3ϵ_+]}`$ are, for the $`ϵ_1ϵ_2ϵ_3=1`$ GSO projection:
$$\begin{array}{cc}\hfill S_{\frac{1}{2}}^+=& S_{[]}\hfill \\ \hfill S_{\frac{1}{2}}^{}=& S_{[+++]}\hfill \\ \hfill S_{\frac{1}{2}}^3=& \frac{1}{2}(S_{[++]}+S_{[+]})\hfill \\ \hfill S_{\frac{1}{2}}^+=& S_{[++]}\hfill \\ \hfill S_{\frac{1}{2}}^{}=& S_{[+++]}\hfill \\ \hfill S_{\frac{1}{2}}^3=& \frac{1}{2}(S_{[++]}+S_{[+++]}).\hfill \end{array}$$
For the other chirality, $`ϵ_1ϵ_2ϵ_3=1`$, we get:
$$\begin{array}{cc}\hfill \stackrel{~}{S}_{\frac{1}{2}}=& \frac{1}{2}(S_{[++]}S_{[+]})\hfill \\ \hfill \stackrel{~}{S}_{\frac{1}{2}}=& \frac{1}{2}(S_{[+++]}S_{[++]}).\hfill \end{array}$$
From the above spin-fields and currents, the construction of the $`N=3`$ (or $`N=1`$ according to the GSO projection) algebra proceeds in exactly the same manner as in the former case.
5. $`N=3`$ superalgebra as an enhancement of $`N=2`$
Since the $`N=3`$ superconformal algebra has the $`N=2`$ superalgebra as a subalgebra, and since general conditions for the appearance of the latter are known , it is natural and instructive to investigate the relation between the two constructions.
In it was found that a general condition for having $`N=2`$ superconformal algebra in spacetime for a background of the form $`AdS_3\times 𝒩`$, is the existence of an affine $`U(1)`$ current in $`𝒩`$, such that $`𝒩/U(1)`$ has $`N=2`$ worldsheet supersymmetry. It was noted there that the $`U(1)`$ must be chosen carefully in the cases where enhancement to $`N>2`$ is expected, in order to embed the $`N=2`$ construction into the explicit construction of the larger algebra.
We now proceed to show that in our two cases there is only one choice of the complex structure in $`𝒩/U(1)`$, where the $`U(1)`$ is the Cartan subalgebra of the $`SU(2)`$, that leads to an $`N=2`$ construction which is a subalgebra of the $`N=3`$ algebra constructed in the previous sections.
The general construction of the $`N=2`$ superconformal algebra for coset models leads to the following $`U(1)`$ R-current:
$$J_R=\frac{i}{k^{}}h_{\overline{a}\overline{b}}\chi ^{\overline{a}}\chi ^{\overline{b}}+\frac{1}{k^{}}h_{\overline{a}\overline{b}}f_{\overline{a}\overline{b}C}\left(\widehat{K}^C\frac{i}{k^{}}f_{C\overline{d}\overline{e}}\chi ^{\overline{d}}\chi ^{\overline{e}}\right).$$
The index $`C`$ can run over both $`H`$ and $`G/H`$. The complex structure $`h_{\overline{a}\overline{b}}`$ has to satisfy conditions that can be found in .
The construction of $`N=2`$ supercharges then proceeds as follows .<sup>4</sup> Note that this is not the standard construction with respect to $`N=2`$ supersymmetry on the worldsheet. We present a canonically normalized scalar $`H_0`$:
$$i\sqrt{3}H_0=J_R\frac{4}{k^{}}K^3,$$
which is used to construct the spin-fields:
$$S=e^{\frac{i}{2}(ϵ_1H_1+ϵ_3H_3+ϵ_0\sqrt{3}H_0)},$$
where $`H_1`$ and $`H_3`$ are built as before (3.1). The BRST condition will pick up 4 of the above spin-fields as physical.<sup>5</sup> Note that the BRST condition in this construction leads to a definite GSO projection (i.e. no BRST invariant spin-fields of the other chirality are explicitly constructed). This construction of $`N=2`$ will be embedded in our $`N=3`$ constructions provided that the above spin-fields (5.1) can be rewritten as special cases of (3.1).
Let us consider first the $`SU(3)/U(1)`$ case. Here we have to look for complex structures of the six dimensional manifold $`SU(3)/U(1)^2`$. We take the Cartan subalgebra of the $`SU(3)`$ to be generated by $`K^3`$ and $`K^8`$. We find three possible complex structures.
The first complex structure is given by $`h_{12}=h_{45}=h_{67}=1`$, and leads to an R-current of the form:
$$J_R^{}=iH_2+i\sqrt{2}H_++\frac{2}{k^{}}K^3+\frac{2\sqrt{3}}{k^{}}K^8.$$
Note that the above expression is such that both $`H_+`$ and $`K^8`$ will appear in the definition (5.1) of $`H_0`$, and therefore the spin-fields (5.1) cannot be matched to (3.1).
The second complex structure is given by $`h_{12}=h_{45}=h_{67}=1`$, and leads to an R-current of the form:
$$J_R^{\prime \prime }=iH_2i\sqrt{2}H_++\frac{2}{k^{}}K^3\frac{2\sqrt{3}}{k^{}}K^8.$$
The same remark as above applies here. Moreover these two currents are the sum and difference of the $`N=2`$ $`U(1)`$ R-currents that one can get by decomposing this coset CFT into $`[SU(2)/U(1)]\times [SU(3)/(SU(2)\times U(1))]`$. We will explain shortly why this direct-product decomposition cannot lead to the enhancement to $`N=3`$.
The third complex structure is given by $`h_{12}=h_{45}=h_{67}=1`$, and leads to an R-current of the form:
$$J_R=iH_2+i\sqrt{2}H_{}+\frac{4}{k^{}}K^3.$$
The boson constructed as in (5.1) now reads $`\sqrt{3}H_0=H_2+\sqrt{2}H_{}`$. The spin-fields (5.1) are thus exactly of the form (3.1) with $`ϵ_2=ϵ_{}`$. The BRST invariant ones are exactly the $`N=2`$ subalgebra generators of (3.1), namely $`S_{\pm \frac{1}{2}}^\pm `$.
Moving to the $`SO(5)/SO(3)`$ case, we have to consider the complex structure of $`SO(5)/(SO(3)\times SO(2))`$, where the $`SO(2)`$ is again the Cartan subalgebra of the remaining $`SO(3)`$, generated by $`K^3`$. This case is different, as there is only one possible complex structure, given by $`h_{12}=h_{78}=h_{90}=1`$. The associated R-current is:
$$J_R=iH_2+i\sqrt{2}H_++\frac{4}{k^{}}K^3.$$
As for (5.1), the BRST invariant spin-fields constructed according to are the $`N=2`$ supercharges $`S_{\pm \frac{1}{2}}^\pm `$ of (4.1). The presence of only one complex structure in this case (as opposed to three in the previous one) is due to the fact that the four dimensional coset $`SO(5)/SO(4)`$ has no complex structure.
6. General conditions for obtaining $`N=3`$
The above discussion leads us to present general conditions for the appearance of the $`N=3`$ superconformal algebra in the context of string theory on $`AdS_3\times 𝒩`$. Such a background leads to $`N=3`$ superconformal algebra in spacetime provided that:
(i) $`𝒩`$ has an affine $`SU(2)`$ current algebra at level $`k^{}=4k`$, where $`k`$ is the level of $`SL(2)`$.
(ii) $`𝒩/U(1)`$ has $`N=2`$ worldsheet supersymmetry, where $`U(1)`$ is the Cartan subalgebra of the above $`SU(2)`$. This condition alone allows one to construct an $`N=2`$ superconformal algebra in spacetime (for a definite GSO projection).
(iii) This spacetime $`N=2`$ algebra is enhanced to $`N=3`$ if the scalar $`H_0`$ constructed as in (5.1) can be decomposed as $`\sqrt{3}H_0=H_2+\sqrt{2}\stackrel{~}{H}_0`$, where $`H_2`$ derives from the bosonization of the two remaining charged fermions of the $`SU(2)`$, and $`\stackrel{~}{H}_0`$ is orthogonal to it.
Interestingly, these conditions imply as a by-product that for the opposite GSO projection we also get supersymmetry in spacetime, namely $`N=1`$.
Let us present the proof by constructing the $`N=3`$ superalgebra generators given the above conditions. Recall that besides the scalar (5.1), we also define the scalars $`H_1=\frac{2}{k}\psi ^1\psi ^2`$ and $`iH_3=\frac{1}{k}\psi ^3\chi ^3`$. The existence of the affine $`SU(2)`$, of which $`\chi ^3`$ is the lower component of the Cartan generator, allows us to define also $`H_2=\frac{2}{k^{}}\chi ^1\chi ^2`$. Consider now the currents $`K^3`$ and $`K^\pm `$. Since they form an $`SU(2)`$ supersymmetric WZW model (embedded inside the CFT on $`𝒩`$), they can be split into orthogonal pieces:
$$K^3=\stackrel{~}{K}^3iH_2,K^\pm =\stackrel{~}{K}^\pm \frac{2}{\sqrt{k^{}}}e^{iH_2}\chi ^3.$$
We start now by noting that condition (iii) implies the following (making use of (5.1)):
$$i\sqrt{2}\stackrel{~}{H}_0(z)K^3(w)\frac{1}{(zw)^2}.$$
This means that $`K^3`$ can be split further:
$$K^3=\widehat{K}^3\frac{i}{\sqrt{2}}\stackrel{~}{H}_0iH_2,$$
where $`\widehat{K}^3`$ has a regular OPE with $`\stackrel{~}{H}_0`$ (and of course $`H_2`$). Similarly, the currents $`K^\pm `$ also split into a ‘bosonic’ part $`\widehat{K}^\pm `$ which realizes an affine $`SU(2)_{k^{}3}`$, an $`SU(2)_1`$ part built from $`\stackrel{~}{H}_0`$ and the usual fermionic $`SU(2)_2`$ piece:
$$K^\pm =\widehat{K}^\pm e^{i\sqrt{2}\stackrel{~}{H}_0}e^{iH_2}(e^{iH_3}+e^{iH_3}).$$
We can now construct the 4 physical spin-fields as in . Note that the presence of the full $`SU(2)`$ is irrelevant in this step. Using the spin-fields of the form (5.1), we get $`S_{[ϵ_1ϵ_3ϵ_0]}=S_{[]},S_{[++]},S_{[++]},S_{[++]}`$. Splitting $`\sqrt{3}H_0=H_2+\sqrt{2}\stackrel{~}{H}_0`$, we can rewrite them as:
$$S_{[ϵ_1ϵ_2ϵ_3\stackrel{~}{ϵ}_0]}=S_{[]},S_{[+++]},S_{[++]},S_{[+++]}.$$
We now find the additional two BRST invariant supercharges, by acting on the above spin-fields with the $`SU(2)`$ ladder operators $`K^\pm `$, which are also BRST invariant (as upper components of primaries of weight $`1/2`$). The result is:
$$\frac{1}{2}(S_{[+]}+S_{[++]}),\frac{1}{2}(S_{[+++]}+S_{[++]}).$$
The set of spin-fields (6.1)-(6.1) matches exactly the ones found in the cases detailed in the previous sections (note that the apparent sign difference with (3.1) and (3.1) can be absorbed in a redefinition of fields; in section 3 we preferred to stick to the usual Gell-Mann basis of $`SU(3)`$).
Defining the spacetime operators as in (3.1), one can show that the $`N=3`$ superconformal algebra closes.
The above proof builds upon the existence of the $`N=2`$ superalgebra, enhancing it to $`N=3`$ using the $`SU(2)`$ currents. An alternative way of building the $`N=3`$ superalgebra, which also reveals the existence of the $`N=1`$ superalgebra for the other GSO projection, is to decompose the supercurrent of the CFT on $`𝒩`$ into an $`SU(2)`$ part and a $`𝒩/SU(2)`$ one. It can then be used to directly find all of the 8 physical spin-fields, 6 of one chirality and 2 of the other. The $`SU(2)`$ part of the supercurrent is:
$$G_{SU(2)}=\frac{2}{k^{}}\left(\frac{1}{2}\chi ^+\stackrel{~}{K}^{}+\frac{1}{2}\chi ^{}\stackrel{~}{K}^++\chi ^3\stackrel{~}{K}^3\frac{2i}{k^{}}\chi ^1\chi ^2\chi ^3\right).$$
Using (6.1), (6.1), (6.1), and the bosonization, the relevant part of $`G_{SU(2)}`$ for the BRST condition (i.e. the one that might lead to $`(zw)^{3/2}`$ singular terms in the OPE with the spin-fields) is:
$$G_{SU(2)}=\mathrm{}+\frac{1}{\sqrt{k^{}}}\left\{i\left(H_2+\frac{1}{\sqrt{2}}\stackrel{~}{H}_0\right)(e^{iH_3}+e^{iH_3})(e^{iH_2i\sqrt{2}\stackrel{~}{H}_0}e^{iH_2+i\sqrt{2}\stackrel{~}{H}_0})\right\}.$$
The first piece will give rise to a $`(zw)^{3/2}`$ singularity only when $`ϵ_2=\stackrel{~}{ϵ}_0`$, while the second piece will do so only when $`ϵ_2=\stackrel{~}{ϵ}_0`$. Choosing $`ϵ_2=\stackrel{~}{ϵ}_0`$, we get 4 physical spin-fields of the same chirality. For $`ϵ_2=\stackrel{~}{ϵ}_0`$ we get 4 physical spin-fields, two of each chirality.
It is worth noting that the direct product $`𝒩=SU(2)_k^{}\times 𝒩^{}`$, which manifestly fulfills conditions (i) and (ii), does not fulfill the third condition (except for one case which will be discussed shortly). The reason for this is the following. If $`H_0`$ fulfills condition (iii), it is straightforward to compute the OPE of $`iH_2`$ with $`J_R`$, the R-current of $`𝒩/U(1)`$, the result being:
$$iH_2(z)J_R(w)\frac{1\frac{4}{k^{}}}{(zw)^2}.$$
However, if $`𝒩=SU(2)_{4k}\times 𝒩^{}`$, then:
$$J_R=J_R^{SU(2)/U(1)}+J_R^{}=iH_2+\frac{2}{k^{}}K^3+J_R^{}.$$
Since $`iH_2`$ has a regular OPE with $`J_R^{}`$, its OPE with $`J_R`$ in (6.1) gives a double pole with a residue of $`\left(1\frac{2}{k^{}}\right)`$ instead of (6.1). Thus such a CFT with a direct product $`SU(2)`$ factor will not lead to $`N=3`$ in spacetime through the mechanism described above.
This seems to contradict what was noted about the limiting $`k^{}=3`$ case, which was reduced to $`SL(2)_{3/4}\times SU(2)_3`$. This is resolved by noting that we can always take $`J_R`$ to $`J_R`$ (this amounts to changing an overall sign in the complex structure, see (5.1)). Doing this, the OPE of $`iH_2`$ with $`J_R`$ gives a residue of $`\left(1\frac{2}{k^{}}\right)`$, which is equal to $`\left(1\frac{4}{k^{}}\right)`$ only if $`k^{}=3`$, that is in this particular case. For completeness, we sketch the construction in this case. Using (5.1) and (6.1), we write:
$$i\sqrt{3}H_0=J_R\frac{4}{3}K^3=iH_22K^3=iH_22\widehat{K}^3.$$
Now we use the fact that the bosonic $`SU(2)`$ WZW model at level 1 can be reformulated as the CFT of a free scalar at its self-dual radius. Denoting this scalar by $`\stackrel{~}{H}_0`$, we have $`\widehat{K}^3=\frac{i}{\sqrt{2}}\stackrel{~}{H}_0`$ and $`\widehat{K}^\pm =e^{i\sqrt{2}\stackrel{~}{H}_0}`$. It is thus clear that $`H_0`$ fulfills condition (iii), and the construction of the $`N=3`$ superalgebra then proceeds as before.
7. Comments
Let us first comment on the relation between the general construction of the previous section and the two specific coset CFTs discussed before. It is possible to check that the only 7 dimensional cosets $`𝒩`$ which have an $`SU(2)`$ symmetry (and do not factorize into a direct product $`SU(2)\times 𝒩^{}`$) are precisely $`𝒩_1`$ and $`𝒩_2`$ of eqs. (2.1) and (2.1). It is interesting that condition (i) of the previous section together with the requirement of having a semiclassical $`k\mathrm{}`$ limit automatically lead to models which fulfill the two remaining conditions. Considering 7 dimensional group manifolds, only the case $`𝒩=SU(2)_{2k}\times SU(2)_{2k}\times U(1)`$ satisfies condition (i), where the $`SU(2)_{4k}`$ is the diagonal one. It is straightforward to show that this manifold also satisfies conditions (ii) and (iii). This manifold actually possesses large $`N=4`$ superconformal symmetry , of which $`N=3`$ is a subalgebra. This fact was used in to break $`N=4`$ to $`N=3`$ through a $`Z_2`$ orbifold construction. It would be interesting to find the relation between this small set of models which share the same spacetime superconformal symmetry.
We conclude by commenting on the geometrical interpretation. It would be nice to translate the conditions we impose on the CFT on $`𝒩`$ into conditions on the geometry of the manifold. It is clear as was commented before that the $`S^3SU(2)`$ has to be non-trivially fibered over the 4 dimensional base $`𝒩/SU(2)`$. A related, but different, problem was actually discussed in the literature \[13,,14\], where conditions on Einstein 7-manifolds $`𝒩`$ are found in order to get different amounts of supersymmetry when considering 11 dimensional supergravity on $`AdS_4\times 𝒩`$. The condition for getting $`N=3`$ in $`AdS_4`$ is that $`𝒩`$ has a tri-sasakian structure (in other words, the cone over it $`C(𝒩)`$ is hyperKähler; see for instance for the notions introduced hereafter). The above geometries are considered as near horizon geometries of M2-branes at the singularity of the Ricci-flat cone $`C(𝒩)`$ over such manifolds. The tri-sasakian structure implies the presence of 3 Killing vectors forming an $`SO(3)`$ algebra which rotates the 3 Killing spinors. It turns out that the only 7-dimensional tri-sasakian manifolds (satisfying some additional regularity conditions, see also ) are exactly the cosets<sup>6</sup> The case $`S^3\times S^3\times S^1`$, which has $`N=3`$ as a subgroup of large $`N=4`$, does not appear in the classification above since its metric cannot be rescaled to become an Einstein manifold because of the flat $`S^1`$ factor. $`SU(3)/U(1)`$ and $`SO(5)/SO(3)S^7`$. In the second case, the quotient is taken as in section 4 . Note also that $`S^7`$ is trivially tri-sasakian since it has $`N=8`$ Killing spinors.
Tri-sasakian manifolds can be seen as $`S^3`$ fibrations over a base, in the above cases $`\mathrm{𝐂𝐏}^2`$ and $`S^4`$ respectively. It is interesting to note that ‘squashing’ a tri-sasakian manifold (i.e. rescaling the fiber with respect to the base) leads, for a definite value of the squashing parameter, to another Einstein manifold having one Killing spinor (instead of 3) and an unbroken $`SO(3)`$ algebra of Killing vectors which acts trivially on the spinor. This is reminiscent of our results, where the different amounts of supersymmetry, $`N=3`$ or $`N=1`$, depend however on the GSO projection.
We should nevertheless stress that in spite of these similarities, there are a few differences. For instance, superstring theory on $`AdS_3\times 𝒩`$ does not require $`𝒩`$ to be an Einstein manifold.<sup>7</sup> This might explain why in the $`SO(5)/SO(3)`$ case discussed in section 4 we do not find $`N=8`$ supersymmetry but only $`N=3`$. Recall that the metric of the coset CFT sigma model, which can be obtained by gauging the WZW model on the group $`G`$ and integrating out the gauge fields, is not the same as the metric on the homogeneous $`G/H`$ coset space. Thus presumably the direct relation between the two issues is more algebraic in nature than geometrical.
Another question regards the brane configuration which might lead to the models considered here in the near horizon limit. Since we are dealing with pure NSNS backgrounds in type II theories, we expect such a brane configuration to involve fundamental strings and NS5-branes intersecting on the string, and possibly at non-trivial angles. Indeed it is known that the $`AdS_3\times S^3\times T^4`$ and the $`AdS_3\times S^3\times S^3\times S^1`$ backgrounds are the near horizon geometries of configurations involving respectively a fundamental string within a NS5-brane , and an additional NS5-brane intersecting the other orthogonally on the string \[23,,24\]. The latter brane configuration can be generalized by introducing a non-trivial angle between the two NS5-branes, but still requiring that some supersymmetries are preserved. This problem has been studied from the supergravity solution point of view in \[25,,26\], and it can be reduced to a problem of classifying the holonomy of an 8-dimensional manifold (the manifold which is orthogonal to the string intersection). It turns out that solutions preserving a fraction of 3/16 of the supersymmetries are associated with 8-dimensional manifolds of holonomy $`Sp(4)U(2,𝐇)SO(5)SO(8)`$, which are thus hyperKähler. This configuration would lead to $`N=(3,3)`$ supersymmetry in type IIA.<sup>8</sup> The supersymmetry of the full brane configuration in type IIB is further reduced. For instance, the orthogonal intersection already preserves only 1/8 of the supersymmetries. This is still different from what we are looking for. In more general NS5-brane configurations are considered, which are related to hyperKähler manifolds with torsion. The torsion allows for the holonomy to be different for the two chiralities of the spinors (again in type IIA). In this set up, solutions which have $`N=(3,1)`$ supersymmetry are found, the associated manifold having $`Sp(4)`$ holonomy for one chirality and $`Spin(7)`$ holonomy for the other. It should be noted that the latter solutions, as presented in , do not include fundamental strings. The near horizon geometry of some of the configurations above is considered in , where it is found to be actually $`AdS_3\times S^3\times S^3\times S^1`$ like in the orthogonal case (hinting towards a large $`N=4`$ dual CFT, instead of an $`N=3`$ one).
It is not straightforward to see how our coset manifolds could arise as near horizon geometries, at least in this context. In order to investigate this problem, one might need to further characterize the 8-dimensional hyperKähler manifolds involved in the brane configurations. Note that the latter are asymptotically flat, and thus generically not in the same class as the conical ones discussed previously, which were related to the classification of tri-sasakian (coset) manifolds.
An alternative to configurations with branes at angles, is to consider the near horizon limit of NS5-branes wrapping on 4-cycles, together with fundamental strings stretched along the unwrapped direction. These are expected to be equivalent to a superstring on $`AdS_3\times 𝒩`$ where $`𝒩`$ has an affine $`U(1)`$ symmetry and $`𝒩/U(1)`$ is an $`N=2`$ SCFT related to the geometry of the 4-cycle. For particular geometries, the $`𝒩/U(1)`$ SCFT was identified with the infrared limit of $`N=2`$ Landau-Ginzburg (LG) models . In the examples considered in this work $`𝒩/U(1)`$ is $`SU(3)/U(1)^2`$ and $`SO(5)/(SO(3)\times U(1))`$. Generically, these $`N=2`$ quotients do not have a LG description (except for the lowest levels of the $`SU(3)/U(1)^2`$ case). Therefore, the relation to brane configurations along these lines requires the understanding of the duality of in models of the form $`AdS_3\times G/H`$, where $`G/(H\times U(1))`$ is a generic $`N=2`$ quotient.
Acknowledgements: We are happy to thank S. Elitzur, O. Feinerman, D. Kutasov, W. Lerche, E. Rabinovici, M. Roček and A. Tomasiello for discussions and helpful remarks. A. G. also thanks the ITP in Santa Barbara for hospitality during the initial stages of this work. This work is supported in part by the Israel Academy of Sciences and Humanities – Centers of Excellence Program, and by the BSF – American-Israel Bi-National Science Foundation.
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Witten, “The D1/D5 system and singular CFT,” JHEP 9904 (1999) 017, hep-th/9903224. relax D. Friedan, E. Martinec and S. Shenker, “Conformal Invariance, Supersymmetry And String Theory,” Nucl. Phys. B271 (1986) 93. relax T. Banks, L. J. Dixon, D. Friedan and E. Martinec, “Phenomenology And Conformal Field Theory Or Can String Theory Predict The Weak Mixing Angle?,” Nucl. Phys. B299 (1988) 613; T. Banks and L. J. Dixon, “Constraints On String Vacua With Space-Time Supersymmetry,” Nucl. Phys. B307 (1988) 93. relax T. Friedrich and I. Kath, “Seven-Dimensional Compact Riemannian Manifolds With Killing Spinors,” Commun. Math. Phys. 133 (1990) 543. relax L. Castellani, L. J. Romans and N. P. Warner, “Symmetries Of Coset Spaces And Kaluza-Klein Supergravity,” Annals Phys. 157 (1984) 394; “A Classification Of Compactifying Solutions For D = 11 Supergravity,” Nucl. Phys. B241 (1984) 429. relax P. M. Cowdall and P. K. Townsend, “Gauged supergravity vacua from intersecting branes,” Phys. Lett. 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# Non-Fermi-Liquid Scaling in Ce(Ru0.5Rh0.5)2Si2
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## Abstract
We study the temperature and field dependence of the magnetic and transport properties of the non-Fermi-liquid compound Ce(Ru<sub>0.5</sub>Rh<sub>0.5</sub>)<sub>2</sub>Si<sub>2</sub>. For fields $`H`$ 0.1 T the results suggest that the observed NFL behavior is disorder-driven. For higher fields, however, magnetic and transport properties are dominated by the coupling of the conduction electrons to critical spin-fluctuations. The temperature dependence of the susceptibility as well as the scaling properties of the magnetoresistance are in very good agreement with the predictions of recent dynamical mean-field theories of Kondo alloys close to a spin-glass quantum critical point.
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The properties of a large class of f-electrons materials show striking departures from the predictions of standard Fermi-liquid theory at low temperature . Several mechanisms leading to non-Fermi-liquid (NFL) behavior have been proposed. In systems close to a quantum phase transition such as CeCu<sub>6-x</sub>Au<sub>x</sub> or CeIn<sub>3</sub>, NFL behavior is due to the coupling of the conduction electrons to critical spin fluctuations . Anomalous properties are observed when the system is driven through the quantum critical point (QCP) by alloying or by applying pressure. In other systems, such as UCu<sub>5-x</sub>Pd<sub>x</sub> , NFL properties are thought to be a consequence of the interplay of strong structural disorder and many body effects . In this paper we report results of a study of the temperature and field dependence of the magnetic and transport properties of Ce(Ru<sub>0.5</sub>Rh<sub>0.5</sub>)<sub>2</sub>Si<sub>2</sub>. We found that this system exhibits different types of anomalies depending on the value of the applied field. At weak fields, we found signatures of disorder effects such as a diverging low-temperature susceptibility and an anomaly in the low-field magnetoresistance. Above 1kG, the $`T`$ and $`H`$-dependence of the susceptibility and the resistivity agrees with the predictions of recent mean-field theories of the spin-glass (SG) QCP . The magnetoresistance is found to exhibit universal scaling properties as predicted by the theory.
In the Ce(Ru<sub>1-x</sub>Rh<sub>x</sub>)<sub>2</sub>Si<sub>2</sub> alloy series the Ce-sublattice is preserved and the hybridization between 4f and conduction electrons varies with the concentration of the ligand 4d atoms Ru and Rh. Pure CeRu<sub>2</sub>Si<sub>2</sub> is a heavy fermion compound with a $`\gamma `$ value of about 385 mJ/mol/K<sup>2</sup> and no long range magnetic order down to 20 mK. With substitution of Rh for Ru, a spin density wave (SDW) region appears between $`x=0.03`$ and $`x=0.4`$ . While Ce(Ru<sub>1-x</sub>Rh<sub>x</sub>)<sub>2</sub>Si<sub>2</sub> at $`x=0.03`$ is a normal Fermi liquid , a NFL regime exists for $`x=0.4`$ and 0.5 . In pure CeRh<sub>2</sub>Si<sub>2</sub>, the 4f electrons are localized and the material is antiferromagnetic (AF) below $`T_N=35`$ K . With increasing Ru substitution, $`T_N`$ decreases and eventually vanishes at a critical concentration $`x_c0.5`$ . There is no direct evidence of long range AF order below $`x=0.7`$. The admixture of Ru and Rh introduces magnetic frustration effects as it leads to a competition between widely different types of magnetic short-range order . A frustrated ground state of the SG type can not be excluded slightly above $`x_c`$. Recent $`\mu `$-SR studies showed that the $`T`$-dependence of the muon relaxation rate in the $`x=0.5`$ alloy is similar to that observed in spin glasses. The muon depolarization rate decays exponentially, however, showing that the spin correlations are dynamic rather than static.
Samples of Ce(Ru<sub>0.5</sub>Rh<sub>0.5</sub>)<sub>2</sub>Si<sub>2</sub> were prepared by arc-melting of the constituents in an argon atmosphere. The ingots were remelted several times to insure homogeneity. Single crystals oriented along a and c-axis were grown by the Czochralsky method in a tri-arc furnace in argon atmosphere and parallelepiped-shaped samples of size $`0.5\times 0.5\times 4`$ mm<sup>3</sup> were obtained. The resistivity was measured by a standard ac method at 17 Hz in the range 16 mK - 4 K and in magnetic fields up to 5 T. The low-field susceptibility was measured in the dilution range $``$ 50 mK - 4 K using a standard method at 130 Hz. The static magnetization was measured in the same range of temperature with a SQUID magnetometer in a dilution setup and, above 3 K, in a commercial SQUID magnetometer.
Fig. 1 displays the resistivity $`\rho `$ measured in a field applied along the c-direction up to 5 T. The data correspond to a current flow in the a-direction. The behavior along the $`c`$-axis is similar but the resistivity is about four times smaller. The magnetoresistance is positive at low temperatures and changes sign at about 2.5 K in the a-direction and at about 1.5 K in the c-direction. In both cases the data at 5 T follow a $`T^2`$-law up to T $``$ 1 K. The range of temperatures in which FL behavior is observed decreases with field and vanishes at $`H=0`$. These results are qualitatively similar to those obtained in CeCu<sub>6-x</sub>Au<sub>x</sub> and in the stoichiometric compound CeNi<sub>2</sub>Ge<sub>2</sub> . The resistivity of our samples is much higher, however, as a consequence of the high degree of substitutional disorder present in the $`x=0.5`$ alloy.
The temperature dependence of the zero-field resistivity is $`\delta \rho T^{1.6}`$ as shown in Fig. 2(a). The resistivity exponent is close but not identical (see below) to that expected for metallic antiferromagnets and spin-glasses at the QCP in the high resistivity limit . The low-temperature magnetoresistance $`(\rho (H)\rho (0))/\rho (0)`$ along the $`a`$-direction is plotted in Fig. 2(b) for $`H<1`$ T. At high fields (not shown in the figure) it shows the classical $`H^2`$-dependence due to the bending of the electron orbits. Below a few kG, however, the low-temperature resistance increases linearly with field as $`H0`$.
The susceptibility $`\chi `$ = M/H (M is the magnetization) is represented in Fig. 3 for H<sub>ac</sub> = 1 G and H<sub>dc</sub> = 0.01 T, 0.1 T, and 1 T for $`T<`$ 10 K. At low temperatures $`\chi `$ decreases strongly with increasing field between 1 G and 1 kG but weakly above 1 kG. Above 3 K the field has no effect on $`\chi `$ up to 1 T as seen by comparison of the curves at 1 kG and 1 T in Fig. (3). At 1 G, $`\chi `$ increases sharply with decreasing $`T`$ below 2 K.
The divergence of $`\chi `$ at low-fields as well as that of $`C/T`$ suggests that, in this regime, NFL behavior may be driven by disorder . The linear (rather than quadratic) rise of the low-field magnetoresistance at low temperatures reported here may be explained by Kondo-disorder effects . Recent $`\mu `$SR experiments showed a sharp increase in the muon relaxation rate below 2 K that saturates at a $`T`$-independent value below 0.7 K. This behavior has been interpreted in terms of the Griffiths-phase scenario in which finite clusters carrying a magnetic moment fluctuate very slowly at low temperature due to quantum tunneling. Magnetic, NMR and specific heat experiments performed in a temperautre range much higher than ours have also been interpreted in terms of this mechanism .
For $`H`$ 1 kG, $`\chi `$ remains finite as $`T0`$ and depends weakly on $`H`$, suggesting that application of a moderate magnetic field quenches the mechanism that leads to the divergence observed at lower fields. Although a full understanding of this fact is still lacking, it should be noticed that in the Griffiths-phase model quantum fluctuations of the largest clusters are expected to be suppressed by a small magnetic field. Indeed, while the Zeeman energy of a cluster of size $`N`$ grows as $`\sqrt{N}`$, its tunneling energy vanishes exponentially with $`N`$. In the following we concentrate on the physics above 1 kG.
The $`T`$-dependence of the susceptibility at 1 kG is still anomalous and $`\delta \chi (T)`$ is approximately linear in $`T^{3/4}`$ (cf. Fig. 3) except at the lowest temperatures, a point that will be discussed further below . A $`T^{3/4}`$ dependence of $`\chi `$ as well as the value 3/2 for the resistivity exponent where predicted by recent dynamical mean field-theories (DMFT) of the spin-glass QCP . In view of the important role that frustration is expected to play in this compound it is tempting to try to interpret our results in terms of the fully frustrated SG model. The latter describes conduction band electrons coupled to localized $`f`$-electron spins via a local Kondo coupling, $`J_\mathrm{K}`$. There is, in addition, a residual Ising-like interaction between the localized spins. The effects of the magnetic frustration introduced by disorder are incorporated by taking random spin couplings $`J_{ij}`$ chosen from a symmetric distribution of width $`J_{ij}^2=J^2`$. Disorder in the Kondo temperature is not included in the model. From the Kondo-temperature distribution determined in Refs. and one can conclude that this effect should play a lesser role than frustration in the low-temperature range that interests us. The SG model was investigated in the framework of dynamical mean-field theory . It exhibits a zero-temperature SG transition when the typical exchange energy $`J=J_\mathrm{c}T_\mathrm{K}`$, the Kondo temperature of the underlying Kondo lattice. Monte Carlo simulations of this model showed that its critical properties are described by an effective strong-coupling theory closely related to other mean-field models. The physical properties of the system depend on the effective distance to the QCP, $`\mathrm{\Delta }(T,H)`$. It can be shown that this is
$$\mathrm{\Delta }=\mathrm{\Delta }_0+2\sqrt{\mathrm{\Delta }_0}\frac{T}{T_0}\left[\sqrt{1+\frac{T}{2\sqrt{3}T_0\mathrm{\Delta }_0}}1\right],$$
(1)
where $`\mathrm{\Delta }_0=(1J/J_\mathrm{c})+(H/H_0)^2`$ and $`H_0=J_\mathrm{c}/(g\mu _\mathrm{B})`$ ($`g`$ is the gyromagnetic ratio of the Ce ion). The scale $`T_0`$ is proportional to $`T_K`$. Numerical simulations yield $`J_\mathrm{c}1.15T_\mathrm{K}`$. The spin susceptibility is
$$J_\mathrm{c}\chi =\sqrt{1+\mathrm{\Delta }}\sqrt{\mathrm{\Delta }}.$$
(2)
The spectrum of magnetic excitation has a scaling form , $`J_c\chi ^{\prime \prime }(\omega )=\sqrt{\mathrm{\Delta }}\mathrm{\Phi }(\omega /J_c\mathrm{\Delta })`$, where the universal scaling function $`\mathrm{\Phi }(x)=x/\sqrt{2}[(1+x^2)^{1/2}+1]^{1/2}`$. The temperature-dependent contribution to the resistivity is computed from $`\delta \rho 1/\tau `$ with the inverse scattering time $`\tau ^1_0^{\mathrm{}}𝑑\omega \chi ^{\prime \prime }(\omega )/\mathrm{sinh}(\beta \omega )`$, an expression valid in the dirty limit . The resistivity has the scaling form
$$\rho (T,H)\rho (0,H)T^{3/2}\mathrm{\Psi }\left(\frac{T}{\mathrm{\Delta }T_0}\right),$$
(3)
with $`\mathrm{\Psi }(x)=x^{1/2}_0^{\mathrm{}}𝑑u\mathrm{\Phi }(ux)/\mathrm{sinh}u`$. It follows that $`\delta \chi \chi (0)\chi (T)T^{3/4}`$ and $`\delta \rho T^{3/2}`$ at the QCP. The resistivity exponent of the mean-field SG model coincides with that of the $`d=3`$, $`z=2`$ antiferromagnet . The susceptibility exponent is specific to the SG model. Away from the QCP, normal Fermi-liquid behavior as $`T0`$ is recovered with both $`\delta \chi `$ and $`\delta \rho T^2/\sqrt{\mathrm{\Delta }_0}`$ at sufficiently low $`T`$. The crossover between these limiting forms will be discussed below. The parameters of the theory can be determined from an analysis of the experimental data. At the critical concentration, $`r=1J/J_\mathrm{c}`$ vanishes. However, $`x_\mathrm{c}`$ is not known accurately and a small but finite $`r`$ can not be excluded a priori. The characteristic field $`H_0`$ may be estimated from an extrapolation to $`T=0`$ of the susceptibility per Ce atom. From the definitions above and Eq. (2) we estimate $`H_0=\mu _\mathrm{B}/\chi (0)11`$ T (we have assumed $`g=2`$). Since $`H_0`$ is very large, the measuring field can be neglected in the analysis of the magnetization at 1 kG. A fit of $`\chi `$ using Eqs. (2) and (1) with $`\mathrm{\Delta }_0=0`$ gives $`T_020`$K, which is slightly smaller than the Kondo temperature of the system estimated from the $`T`$-dependence of the resistivity.
We analyzed the $`T`$\- and $`H`$-dependence of the resistivity using Eqs. (3) and (1). The condition that all the data in Fig. 1 collapse into a single scaling curve leaves little freedom in the choice of the parameters. In particular, we found it impossible to scale all the curves with $`r=0`$. A scaling plot of the resistivity along the $`a`$-axis is shown in Fig. 4. The data points are the values of the scaled resistance $`(\rho (T,H)\rho (0,H))\times T^{3/2}`$ plotted vs the reduced variable $`t/\mathrm{\Delta }`$ ($`t=T/T_0`$) for $`T0.9`$ K and $`H5`$ T. The values of the parameters are $`r=7\times 10^3`$, $`T_0=20`$ K and $`H_0=13`$ T. The value of $`r`$ measures the distance to the true QCP, $`r=\delta J/J_\mathrm{c}`$ giving $`\delta J0.2`$ K, a very small energy compared to the other energy scales present in the problem. The characteristic field determined from this analysis is close to the theoretical estimate given above. The solid line in Fig. 4 is the theoretical scaling function $`\mathrm{\Psi }(x)`$. There are no adjustable parameters other than an amplitude that fixes the vertical scale. The agreement between theory and experiment is very good except for the data for $`H`$=0 (the empty squares in Fig. 4) which lie slightly above the scaling curve. The slope of the curve, that measures the effective resistivity exponent, is correctly reproduced by the theory. We can now understand that the deviation of the resistivity exponent (cf. Fig. 2a) with respect to its value at the QCP (1.5) is due the small but finite value of $`r`$. The effective exponent only reaches 1.5 for $`t/\mathrm{\Delta }\mathrm{}`$, i.e. for $`r`$=0. We ascribe the excess amplitude for $`H=0`$, represented by the vertical shift, to additional scattering processes dues to disorder. We can compare $`A`$, the amplitude of the $`T^2`$ term in the resistivity in the FL region, with the theoretical prediction, $`A1/\sqrt{\mathrm{\Delta }_0}`$. The inset in Fig. 1 shows the field dependence of $`A`$ as determined from the initial slope $`d\rho /dT^2`$ of the resistivity and the theoretical prediction. There is good agreement.
The susceptibility can also be calculated and compared with the data. The solid line in Fig. (3) is the theoretical result for $`H=1`$ kG. The data (and the theoretical curve) deviate from a pure $`T^{3/4}`$ law as $`T0`$. This is due to the finite value of $`r`$ which results in normal FL behavior below a crossover temperature $`T_{\mathrm{FL}}T_0\mathrm{\Delta }_0`$. $`T_{\mathrm{FL}}`$ increases with field and can be estimated as $``$ 0.25 K for $`H`$=1 T. The crossover to $`T^2`$ behavior in $`\chi (H=1\mathrm{T})`$ can be seen in Fig. (3). The low-field results can be described by adding to the dynamical mean-field result an additional diverging contribution. The presence of a paramagnetic phase giving rise to a $`T^1`$ divergence of $`\chi `$ was suggested in Ref. . However, to suppress such a contribution at 3 K by a field as small as 1 kG one would need that the impurities carry huge moments ($`>`$ 3 $`\mu _B`$). Furthermore, this hypothesis would not explain the divergence of $`\gamma `$ . The inset in Fig. 3 shows a fit of the ac data for $`H`$=1 G to the expression $`\chi (T)=\chi _{\mathrm{MF}}(T)+aT^{0.8}`$. This can be interpreted in terms of the Griffiths-phase model that predicts a power-law divergence $`\delta \chi T^{1+\lambda }`$ with an exponent $`\lambda `$ that vanishes at the QCP. The value $`\lambda =0.2`$ that comes out of our analysis is consistent with this picture for a system close to the QCP. We also found that the $`\gamma `$ data can be accurately described by the analogous expression $`\gamma (T)=\gamma _{\mathrm{MF}}(T)+a^{}T^{0.8}`$ where $`\gamma _{\mathrm{MF}}(T)1b\sqrt{T}`$ is the DMFT prediction for $`C/T`$. The equality of the exponents describing the divergent parts of $`\chi `$ and $`C/T`$ is one of the predictions of the Grifith’s phase model .
In summary, we have shown that the NFL properties of Ce(Ru<sub>0.5</sub>Rh<sub>0.5</sub>)<sub>2</sub>Si<sub>2</sub> are determined by disorder and proximity to a QCP. The effects of the two mechanisms can be disentangled by applying a small magnetic field. At low fields, disorder effects dominate. Above 1kG, however, the temperature and field dependence of the susceptibility are well described by the dynamical mean-field theory of the spin glass QCP. The $`T`$\- and $`H`$-dependent resistivity is found to obey a universal scaling law.
This work was supported in part by the Grant-in-Aid for Scientific Research (B) and the Monbusho International Scientific Research Program. One of us (DRG) thanks the Newton Institute for its kind hospitality and A. H. Castro Neto for useful correspondence.
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# HARTREE-FOCK-BOGOLIUBOV FOR DEFORMED NEUTRON-RICH NUCLEI
## 1 Introduction
In neutron-rich nuclei, the effects of continuum states on the pairing correlation are expected to play an important role. As for nuclei near the neutron drip line, it is obvious that the continuum states are strongly involved in the pairing correlation because the Fermi level $`\lambda _\mathrm{n}`$ and the pairing gap $`\mathrm{\Delta }_\mathrm{n}`$ have similar sizes $`\pm 1`$ MeV. The effect of such a continuum-state pairing may be so strong that the neutron drip line can be pushed outward by several nuclei.
For nuclei with $`\lambda _\mathrm{n}\stackrel{>}{}5`$ MeV, too, the continuum states should be taken into account explicitly because the pairing-active space for HFB calculations should be larger than one major shell, i.e., $`ϵ_{\mathrm{cut}}`$ $`>`$ $`\lambda `$ $`+`$ $`\mathrm{}\omega `$, in order to take into account a correct size of shell effects.
The pairing correlation has been treated usually in the BCS approximation, which relies on an assumption that the pair-scattering matrix elements $`\psi _i\psi _{\overline{ı}}|V_\mathrm{p}|\psi _j\psi _{\overline{ȷ}}`$ do not depend on the form of the wavefunctions $`\psi _i`$ and $`\psi _j`$, e.g. a constant. This assumption results in a situation that the nucleus is surrounded by unphysical dilute neutron gas when $`ϵ_{\mathrm{cut}}>0`$. Therefore, to include the positive energy states in the pairing correlation, one needs to switch from the HF(Hartree-Fock)+BCS to the Hartree-Fock-Bogoliubov (HFB) scheme in coordinate space.
So far, several methods have been presented to solve the HFB equation. For spherical nuclei, there are three methods:
1) Radial differential equations (Dobaczewski et al., 1984),
2) Finite-element method for finite-range pairing forces (Ring et al., 1997),
3) Natural-orbital representation (Reinhard et al, 1997).
For deformed nuclei, two methods have been presented. Although spherical cases can be solved easily with present computers, deformed cases in coordinate space are still a challenge.
4) Diagonalization in the harmonic-oscillator basis (Gogny et al., 1980),
5) Two-basis method (Heenen et al., 1994).
In this paper, we describe the formulations in the subsequent two sections and then study two subjects: First, in section 4, we examine the two-basis method concerning the precision of the low-density tail at large radius as a function of the cutoff energy, i.e., the number of HF orbitals included in the diagonalization. We show that more precise solutions require higher cutoff energies, which result in HFB matrices of larger dimensions. Second, in section 5, we test the natural-orbital method for deformed cases in a three-dimensional (3D) Cartesian mesh representation.
## 2 HFB in coordinate space
In this section, we formulate the HF and the HFB in the coordinate-space representation in order to elucidate a difficulty of the HFB and to suggest its possible solution in terms of the natural-orbital representation. For the sake of simplicity, in this section, we consider only one kind of nucleons and designate the number of nucleons by $`N`$. The spin of a nucleon is represented by $`s`$.
### 2.1 HF
In the HF, one should minimize $`\mathrm{\Psi }|H|\mathrm{\Psi }`$ for single Slater-determinant states,
$`|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}a_i^{}|0,`$ (1)
$`a_i^{}`$ $`=`$ $`{\displaystyle \underset{s}{}}{\displaystyle 𝑑\stackrel{}{r}\psi _i(\stackrel{}{r},s)a_{\stackrel{}{r}s}^{}\text{ : single-particle state}},`$ (2)
by varying $`\{\psi _i\}_{i=1,\mathrm{},N}`$ under orthonormality conditions $`\psi _i|\psi _j`$ = $`\delta _{ij}`$.
### 2.2 HFB
In the HFB, the state takes the following form,
$`|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{\mathrm{\#}\mathrm{basis}}{}}}b_i|0,`$ (3)
$`b_i`$ $`=`$ $`{\displaystyle \underset{s}{}}{\displaystyle 𝑑\stackrel{}{r}\left\{\varphi _i^{}(\stackrel{}{r},s)a_{\stackrel{}{r}s}+\phi _i(\stackrel{}{r},s)a_{\stackrel{}{r}s}^{}\right\}\text{ : quasi-particle state}},`$ (4)
where “#basis” is the number of basis states of the employed representation: For a 3D-mesh representation, it is the number of the mesh points (times four when spin-orbit potentials are included) and typically $`10^4`$-$`10^5`$. One should vary $`\{\varphi _i,\phi _i\}_{i=1,\mathrm{},\mathrm{\#}\mathrm{basis}}`$ under appropriate orthonormality conditions.
The essential difference between the HF and the HFB is that one has to consider only $`N`$ $`10^2`$ wavefunctions in the former while one has to treat explicitly as many single-particle wavefunctions as the number of the basis in the latter.
### 2.3 HFB with the two-basis method
In this method, the HFB equation is solved by diagonalizing the HFB matrix in a single-particle basis $`\{\psi _i\}_{i=1,\mathrm{},K}`$ consisting of the eigenstates of the mean-field hamiltonian $`h`$ (excluding the pairing potential): $`h\psi _i`$ = $`ϵ_i\psi _i`$. The number of the basis, $`K`$, is determined by a cutoff energy $`ϵ_{\mathrm{cut}}`$ as $`ϵ_1`$ $`\mathrm{}`$ $`ϵ_K`$ $``$ $`ϵ_{\mathrm{cut}}`$ $``$ $`ϵ_{K+1}`$ $`\mathrm{}`$. We will show in section 4 that $`KN`$ to obtain high-precision density tails.
### 2.4 HFB in natural-orbital basis
Owing to the Bloch-Messiah theorem, the state (3) can be expressed as,
$`|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{\mathrm{\#}\mathrm{basis}}{}}}\left(u_i+v_ia_i^{}a_{\overline{ı}}^{}\right)|0,`$ (5)
$`a_i^{}`$ $`=`$ $`{\displaystyle \underset{s}{}}{\displaystyle 𝑑\stackrel{}{r}\psi _i(\stackrel{}{r},s)a_{\stackrel{}{r}s}^{}\text{ : natural orbital (canonical basis)}}.`$ (6)
One should vary $`\{\psi _i,u_i\}_{i=1,\mathrm{},\mathrm{\#}\mathrm{basis}}`$ under constraints for orthonormality,
$$\psi _i|\psi _j=\delta _{ij}(1ij\text{\#basis}),$$
(7)
and for the expectation value of the number of nucleons, $`2_{i=1}^Kv_i^2=N`$.
Reinhard et al. regard the advantage of the representation (5) over (3) to be that one has to consider only a single set of wavefunctions $`\{\psi _i\}`$ unlike a double set $`\{\varphi _i,\phi _i\}`$. However, we expect more benefit from the natural-orbital representation. Namely, $`i`$ may be truncated as $`iK`$ = $`𝒪(N)`$ $``$ $`\mathrm{\#}\mathrm{basis}`$ to a very good approximation. It is because $`\psi _i`$ should be a localized function while the orthogonality does not allow many wavefunctions to exist in the vicinity of the nucleus. This situation is illustrated in Fig. 1. For 3D-mesh representations, # basis is proportional to the volume of the cavity (box) while $`K`$ is proportional to the volume of the nucleus. The latter is $`10^1`$-$`10^2`$ times as small as the former.
## 3 HFB with density-dependent zero-range forces
### 3.1 Interaction
We employ density-dependent zero-range interactions in the rest of this paper for the sake of simplicity. There will not be essential differences in the formulation if we use the full-form Skyrme force. The force is expressed in the parameterization of the Skyrme force as,
$$V(\stackrel{}{r}_1,s_1;\stackrel{}{r}_2,s_2)=\left(t_0+\frac{1}{6}t_3\rho \left(\stackrel{}{r}_1\right)^\alpha \right)\delta \left(\stackrel{}{r}_1\stackrel{}{r}_2\right).$$
(8)
We adopt $`t_0`$ = $`1099`$ MeV fm<sup>3</sup>, $`t_3`$ = $`17624`$ MeV fm<sup>3+3α</sup>, and $`\alpha `$ = 0.98 (it is not 1 only for the sake of a test of the code) when the force is used to construct the mean-field (HF) potential. We use a different strength to make the pairing potential. We express the force in the parameterization of Ref.:
$$V_\mathrm{p}(\stackrel{}{r}_1,s_1;\stackrel{}{r}_2,s_2)=v_\mathrm{p}\frac{1P_\sigma }{2}\left(1\frac{\rho (\stackrel{}{r}_1)}{\rho _\mathrm{c}}\right)\delta \left(\stackrel{}{r}_1\stackrel{}{r}_2\right).$$
(9)
We use $`\rho _\mathrm{c}`$ = 0.32 fm<sup>-3</sup> (to roughly vanish the volume-changing effect).
### 3.2 Hamiltonian density
For the sake of simplicity, we treat $`N`$=$`Z`$ nuclei without Coulomb interaction in the rest of this paper. Then, the state of the nucleus is expressed as,
$$|\mathrm{\Psi }=\underset{i=1}{\overset{K}{}}\left(u_i+v_ia_i^{}a_{\overline{ı}}^{}\right)_{\mathrm{proton}}\left(u_i+v_ia_i^{}a_{\overline{ı}}^{}\right)_{\mathrm{neutron}}|0.$$
(10)
With the interactions (8) and (9), the total energy for state (10) is written as,
$`E`$ $`=`$ $`\mathrm{\Psi }|H|\mathrm{\Psi }={\displaystyle \left(\stackrel{}{r}\right)𝑑\stackrel{}{r}},`$ (11)
$``$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2m}}\tau +{\displaystyle \frac{3}{8}}t_0\rho ^2+{\displaystyle \frac{1}{16}}t_3\rho ^{2+\alpha }+{\displaystyle \frac{1}{8}}v_\mathrm{p}\left(1{\displaystyle \frac{\rho }{\rho _\mathrm{c}}}\right)\stackrel{~}{\rho }^2,`$ (12)
where $`\tau `$ is the kinetic energy density,
$$\rho (\stackrel{}{r})=4\underset{i=1}{\overset{K}{}}v_i^2|\psi _i(\stackrel{}{r})|^2,\stackrel{~}{\rho }(\stackrel{}{r})=4\underset{i=1}{\overset{K}{}}u_iv_i|\psi _i(\stackrel{}{r})|^2.$$
(13)
The stationary condition $`\delta =0`$ leads to a mean-filed and a pairing potentials:
$`h`$ $`={\displaystyle \frac{\delta E}{\delta \rho }}`$ $`={\displaystyle \frac{\mathrm{}^2}{2m}}\stackrel{}{}^2+{\displaystyle \frac{3}{4}}t_0\rho +{\displaystyle \frac{2+\alpha }{16}}t_3\rho ^{1+\alpha }{\displaystyle \frac{v_\mathrm{p}}{8\rho _\mathrm{c}}}\stackrel{~}{\rho }^2,`$ (14)
$`\stackrel{~}{h}`$ $`={\displaystyle \frac{\delta E}{\delta \stackrel{~}{\rho }}}`$ $`={\displaystyle \frac{1}{4}}v_\mathrm{p}\left(1{\displaystyle \frac{\rho }{\rho _\mathrm{c}}}\right)\stackrel{~}{\rho }.`$ (15)
## 4 Test of the accuracy of the two-basis method for HFB
We have developed from scratch an HFB program for spherical cases based on the two-basis method explained in sect. 2.3. It is used to examine the precision of the low-density tail at large radius as a function of the cutoff energy $`ϵ_{\mathrm{cut}}`$.
In Fig. 2 we show the density profiles for $`(Z,N)=(38,68)`$ calculated with the Braghin-Vautherin force and $`v_\mathrm{p}`$ = $`400`$ MeV. One can see that (i) The density is localized with accuracy $`10^6`$ fm<sup>-3</sup>. (ii) The cutoff makes ripples in the tail. The wave-length of the ripples agrees with $`2\pi \mathrm{}(2mϵ_{\mathrm{cut}})^{1/2}`$, i.e., half of the de Broglie wavelength for $`ϵ_{\mathrm{cut}}`$.
To obtain a more precise density tail, one has to increase $`ϵ_{\mathrm{cut}}`$, which results in an increase of $`K`$, the number of single-particle wavefunctions to be considered explicitly. $`K`$ increases only slowly as $`ϵ_{\mathrm{cut}}^{1/2}`$ for spherical case for each angular momentum, while it grows rapidly as $`ϵ_{\mathrm{cut}}^{3/2}`$ for deformed case. The rapid grow of the latter case makes the two-basis method practically inapplicable to deformed nuclei because the bottle-necks in computation of mean-field methods on 3D meshes are the parts which spend computation time proportional to $`K^2`$ like the orthogonalization. Therefore, one need a different method to solve the HFB for deformed nuclei with enough high cutoff energies.
## 5 Implementation of the natural-orbital HFB method on 3D mesh
The natural-orbital method explained in sect. 2.4 was originally introduced for spherical nuclei. We have implemented the method to treat deformed nuclei in a 3D-mesh representation.
First, let us present a summary of the formulation, which has some differences from Ref. Instead of minimizing $`E`$ = $`\mathrm{\Psi }|H|\mathrm{\Psi }`$ with $`|\mathrm{\Psi }`$ given by Eq. (10) under constraints of Eq. (7), one may introduced a Routhian $`R`$,
$$R=Eϵ_{\mathrm{Fermi}}4\underset{i=1}{\overset{K}{}}v_i^2\underset{i=1}{\overset{K}{}}\underset{j=1}{\overset{K}{}}\lambda _{ij}\left\{\psi _i|\psi _j\delta _{ij}\right\},\lambda _{ij}=\lambda _{ji}^{},$$
(16)
and minimize it without constraints. In the above definition, in order to make $`R`$ real for the sake of convenience, $`K^2`$ Lagrange multipliers $`\lambda _{ij}`$ obeying hermiticity are introduced instead of $`\frac{1}{2}K(K+1)`$ independent multipliers. Note that $`\delta _{ij}`$ is subtracted from $`\psi _i|\psi _j`$ in order to treat $`\lambda _{ij}`$ not as constants like $`ϵ_{\mathrm{Fermi}}`$ but as functionals of the wavefunctions.
Stationary conditions of $`R`$ result in the following equations.
$`{\displaystyle \frac{R}{v_i}}=0`$ $``$ $`v_i^2={\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{h_{ii}\lambda }{\sqrt{(h_{ii}\lambda )^2+\stackrel{~}{h}_{ii}^2}}},\text{assuming}\stackrel{~}{h}_{ii}0,`$ (17)
$`{\displaystyle \frac{\delta R}{\delta \psi _i^{}}}=0`$ $``$ $`_i\psi _i{\displaystyle \underset{j=1}{\overset{K}{}}}\lambda _{ij}\psi _j{\displaystyle \underset{j=1}{\overset{K}{}}}{\displaystyle \underset{k=1}{\overset{K}{}}}{\displaystyle \frac{\delta \lambda _{jk}}{\delta \psi _i^{}}}\left\{\psi _j|\psi _k\delta _{jk}\right\}=0,`$ (19)
$`_i=v_i^2h+u_iv_i\stackrel{~}{h}.`$
For HF, the orthogonalization conditions are easily realized because $`\psi _i`$ are eigenstates of the same hermite operator $`h`$ and thus are orthogonal at the solution: The orthogonalization procedure is needed only because the orthogonality is unstable. On the other hand, for the natural-orbital HFB method, the orthogonalization is essential because the single-particle hamiltonians $`_i`$ differs from state to state. Therefore, the determination of the explicit functional form of $`\lambda _{ij}`$ is the most important part of the method. Reinhard et al. have proposed
$$\lambda _{ij}=\frac{1}{2}\psi _j|\left(_i+_j\right)|\psi _i.$$
(20)
We can justify their choice as follows: From the requirement that Eq. (19) must hold at the solution (where $`\psi _i|\psi _j=\delta _{ij}`$), one deduces,
$$\lambda _{ij}=\psi _j|_i|\psi _i\text{at the solution}.$$
(21)
Eqs. (20) and (21) are equivalent at the solution because $`\lambda _{ij}`$ is defined to be hermite. Since this hermiticity must hold at any points, one should not adopt Eq. (21) but Eq. (20). Because what is needed is Eq. (21) and the hermiticity, one may use more complex forms like,
$$\lambda _{ij}=\frac{1}{2}\psi _j|\left(_i+_j\right)|\psi _i\left\{1+f_2\left(\psi _j|_i|\psi _i\psi _j|_j|\psi _i\right)^2\right\},$$
(22)
where $`f_2`$ is a parameter to maximize the convergence speed.
To obtain the HFB solutions in the natural-orbital formalism, one can utilize the gradient method, which includes the imaginary-time evolution method:
$$\psi _i\psi _i\mathrm{\Delta }\tau \frac{\delta R}{\delta \psi _i^{}}$$
(23)
We have developed a 3D-mesh natural-orbital HFB program from scratch according to the above formulation. First of all, we test the feasibility of the method for <sup>40</sup>Ca. The wavefunctions are expressed with $`39\times 39\times 39`$ mesh points with mesh spacing of 0.8 fm. Note that the requirement of precision is higher for HFB than for HF because one has to treat larger momentum components than the Fermi momentum in HFB. We employed the 17-point finite-difference approximation to the Laplacian. The vanishing boundary conditions are imposed on the boundary (0th and 40th mesh points) and the wavefunctions are anti-symmetrically reflected in the boundary to apply the finite-difference formula. We considered $`K=20`$ natural orbitals, which can contain 80 (=$`2\times A`$) nucleons.
We show an example of the convergence history in Fig. 3. In this calculation, we set $`f_2=0`$ in Eq. (22) and $`\mathrm{\Delta }\tau `$ = $`10^{24}`$ sec in Eq. (23) . We neglect $`\delta \lambda /\delta \psi ^{}`$ in Eq. (19). Instead, at every 50 steps, $`\{\psi _i\}`$ are Gram-Schmidt orthogonalized in the ascending order of $`h_{ii}`$ and then the HFB hamiltonian is diagonalized in the basis to renew $`\{\psi _i,v_i\}`$.
In the left-hand portion, the error of Eq. (19), i.e., $`\mathrm{max}_{i=1,\mathrm{},K}|_i\psi _i_j\lambda _{ij}\psi _j|`$ are plotted versus the evolution step. The corresponding quantity for HF, $`\mathrm{max}_{i=1,\mathrm{},A/4}|h\psi _i\psi _i|h|\psi _i\psi _i|`$ , is also plotted. The figure demonstrates that one can indeed obtain HFB solutions with the natural-orbital HFB method in the 3D-mesh representation. We obtained similar convergence curves for the error of the orthogonality and for the inconsistency between the potential and the densities. The right-hand portion shows the error of the total energy (estimated as the difference of the total energy from the convergent value), which also shows an exponential convergence pattern.
The speed of the convergence is, however, about ten times as slow as the HF case. We are now trying to improve the convergence speed, the stability of the evolution, and the robustness of the method.
## 6 Influence of pairing correlation on deformation
As examples of the applications of our natural-orbital HFB program, we show in Fig. 4 the change of shape due to pairing correlation. The root-mean-square values of $`x`$, $`y`$, and $`z`$ are plotted as functions of the strength of the pairing interaction $`v_\mathrm{p}`$ for <sup>32</sup>S and <sup>60</sup>Zn. Both nuclei are triaxial when the pairing correlation is off (pairing gap $`\mathrm{\Delta }=0`$). However, as soon as the pairing correlation sets in, the axial symmetry is restored.
The pairing correlation does not always restore symmetric shapes but sometimes break a symmetry present without pairing: For <sup>60</sup>Zn, the expectation value of $`\frac{1}{2}(1\tau _z)Y_{30}`$ is zero before the pairing sets in while it is about half the Weiskopf unit when the pairing correlation is present.
## Acknowledgments
The author thanks to Prof. J. Dobaczewski, Prof. P.-H. Heenen, Prof. D. Brink and Prof. N. Onishi for discussions during and after the symposium.
This paper has been published in the proceedings of the XVII RCNP international symposium on Innovative Computational Methods in Nuclear Many-Body Problems –Towards a new generation of physics in finite quantum systems– (INNOCOM97), Osaka, Japan, November 10-15, 1997, edited by H. Horiuchi, M. Kamimura, H. Toki, F. Fujiwara, M. Matsuo, and Y. Sakuragi, (1998) World Scientific (Singapore), pp. 343-351. The present address of the author is: Department of Applied Physics, Fukui University, Bunkyo 3-9-1, Fukui, 910-8507, Japan E-mail: tajima@quantum.apphy.fukui-u.ac.jp
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# Robustness of Wave Functions of Interacting Many Bosons in a Leaky Box
## Abstract
We study the robustness, against the leakage of bosons, of wave functions of interacting many bosons confined in a finite box, by deriving and analyzing a general equation of motion for the reduced density operator. We identify a robust wave function that remains a pure state, whereas other wave functions, such as the Bogoliubov’s ground state and the ground state with a fixed number of bosons, evolve into mixed states. Although these states all have the off-diagonal long-range order, and almost the same energy densities, we argue that only the robust state is realized as a macroscopic quantum state.
When a quantum system is subject to perturbations from its environment, most wave functions decohere, and only an exceptional wave function(s) remains pure. For quantum systems with a single degree of freedom ($`f=1`$), this robust wave function is a coherent state . For example, when a coherent state $`|\alpha `$ of single-mode photons passes through an absorptive medium, the final state is also a coherent state $`|\alpha ^{}`$, which was attenuated ($`|\alpha ^{}|<|\alpha |`$) by the absorption . It was also argued for a $`f=1`$ system that coherent states produce the least entropy in the environment, thus being stable . Since these conclusions are based on analyses on $`f=1`$ systems, a natural question is: Are they applicable to macroscopic systems, i.e., to $`f1`$ interacting systems? Moreover, we must identify which coherent states are robust, because there are many choices of the coordinate (among many degrees of freedom) by which a coherent state is defined. Furthermore, for massive bosons the superselection rule (SSR) forbids superpositions of states with different numbers of bosons. Hence, we must clarify how coherent states can be compatible with the SSR. The purpose of this Letter is to answer these questions for condensates of interacting many bosons, which (or, equivalents of which) are observed in many physical systems such as liquid He , quantum Hall systems , excitons , and trapped atoms . We also discuss the symmetry breaking in view of the robustness.
We consider many bosons which interact with each other repulsively. We assume that the bosons are confined in a large, but finite box of volume $`V`$, which is placed in a huge room of volume $`𝒱V`$, which we call the environment. Suppose that the potential of the walls of the box is not high enough, so that the box and the environment exchange bosons via tunneling processes at a small rate (flux) $`J`$. Let $`t_{\mathrm{eq}}`$ denote the time scale after which the total system, the box plus the environment, reaches the equilibrium state. We are not interested in the time region $`tt_{\mathrm{eq}}`$ because the equilibrium state is just a uniform state that is determined solely by the initial state of the environment (because $`𝒱V`$). We therefore examine the transient region for which $`tt_{eq},`$ in order to discuss the robustness of an initial state $`|\varphi (0)`$, which is prepared at $`t=0`$, of the box. Depending on the choice of the initial state $`|\varphi _\mathrm{E}(0)_\mathrm{E}`$ of the environment, the box state may be affected either drastically or moderately. For example, a moderate situation is that $`|\varphi _\mathrm{E}(0)_\mathrm{E}`$ has the same density $`n`$ of bosons as $`|\varphi (0)`$. In such a case, $`n`$ of the box will be kept constant for all $`t`$. To discuss the robustness, however, we consider the severest situation where the environment is initially in the vacuum state $`|0_\mathrm{E}`$ of bosons, so that bosons escape from the box continuously. If a box state is robust in this severest case, it would also be robust in other cases. Hence, the total wave function at $`t=0`$ is $`|\mathrm{\Phi }(0)_{\mathrm{total}}=|\varphi (0)|0_\mathrm{E}.`$ We decompose (the $`𝐫`$ dependence of) the boson field that is defined on $`V+𝒱`$ as $`\widehat{\psi }_{\mathrm{total}}(𝐫)=\widehat{\psi }(𝐫)+\widehat{\psi }_\mathrm{E}(𝐫).`$ Here, $`\widehat{\psi }(𝐫)`$ localizes in the box, whereas the low-energy component of $`\widehat{\psi }_\mathrm{E}(𝐫)`$ localizes in the environment . Accordingly, the Hamiltonian of the total system is decomposed as $`\widehat{H}_{\mathrm{total}}=\widehat{H}+\widehat{H}_\mathrm{E}+\widehat{H}_{\mathrm{SE}}.`$ Here, $`\widehat{H}`$ ($`\widehat{H}_\mathrm{E}`$) is a function of $`\widehat{\psi }`$ ($`\widehat{\psi }_\mathrm{E}`$) only, describing interacting bosons in the box (environment). On the other hand, $`\widehat{H}_{\mathrm{SE}}`$ includes both $`\widehat{\psi }`$ and $`\widehat{\psi }_\mathrm{E}`$, describing the $`\widehat{\psi }`$-$`\widehat{\psi }_\mathrm{E}`$ interaction. If the leakage flux $`J`$ is small, the probability of finding two or more bosons simultaneously in a wall of the box is negligible, and thus the dominant term of $`\widehat{H}_{\mathrm{SE}}`$ takes the following form:
$$\widehat{H}_{\mathrm{SE}}=\lambda d^3𝐫\widehat{\psi }_\mathrm{E}^{}(𝐫)w(𝐫)\widehat{\psi }(𝐫)+\mathrm{H}.\mathrm{c}.$$
(1)
Here, $`w(𝐫)`$ represents the shapes of the walls ($`w1`$ in the walls, $`w=0`$ in other regions), whose potential height is characterized by a parameter $`\lambda `$. Details of $`w`$ are irrelevant because they are all absorbed in the value of $`j`$, Eq. (11). In the time region of interest ($`tt_{eq}`$), the $`\widehat{\psi }_\mathrm{E}`$-$`\widehat{\psi }_\mathrm{E}`$ interaction should be unimportant because $`n`$ of the environment remains zero. On the other hand, we must treat the $`\widehat{\psi }`$-$`\widehat{\psi }`$ interaction appropriately. For this purpose, we use the decomposition formula for $`\widehat{\psi }`$ :
$$\widehat{\psi }=\widehat{\mathrm{\Xi }}+\widehat{\psi }^{},$$
(2)
where $`\widehat{\mathrm{\Xi }}`$ is an operator satisfying
$$\widehat{\mathrm{\Xi }}|N,\mathrm{G}=\sqrt{N}\xi |N1,\mathrm{G},$$
(3)
where $`|N,\mathrm{G}`$ denotes the ground state that has exactly $`N`$ bosons, which we call the number state of interacting bosons (NSIB), and
$$\xi N1,\mathrm{G}|\widehat{\psi }|N,\mathrm{G}/\sqrt{N}$$
(4)
is a hallmark of the condensation: $`\sqrt{N}\xi =𝒪(1)`$ for condensed states, whereas $`\sqrt{N}\xi =𝒪(1/\sqrt{V})`$ for normal states . We here consider the condensed states. Since $`\widehat{\psi }`$ alters $`N`$ exactly by 1, Eq. (4) means that $`N\mathrm{\Delta }N,\mathrm{G}|\widehat{\psi }|N,\mathrm{G}=\sqrt{N}\xi \delta _{\mathrm{\Delta }N,1}`$ for all $`\mathrm{\Delta }N`$ such that $`|\mathrm{\Delta }N|N`$. It then follows from Eqs. (2) and (3) that
$$N\mathrm{\Delta }N,\mathrm{G}|\widehat{\psi }^{}|N,\mathrm{G}=0\text{(for }|\mathrm{\Delta }N|N\text{)}.$$
(5)
Namely, $`\widehat{\psi }^{}`$ transforms $`|N,\mathrm{G}`$ into excited states.
For weakly-interacting bosons, the explicit forms of the NSIB were given in Refs. , and that of $`\widehat{\mathrm{\Xi }}`$ was given in . Because of the boson-boson interaction, they are rather complicated functions of bare operators $`\widehat{a}_𝐤`$: $`|N,\mathrm{G}=(1/\sqrt{N!})e^{i\widehat{G}}(\widehat{a}_0^{})^N|0`$ and $`\widehat{\mathrm{\Xi }}=e^{i\phi }\sqrt{n_0/nV}e^{i\widehat{G}}\widehat{a}_0e^{i\widehat{G}}`$. Here, $`n_0=\widehat{N}d^3𝐫\widehat{\psi }^{}\widehat{\psi }^{}/V`$, $`\widehat{G}(i/2nV)\widehat{a}_0^{}\widehat{a}_0^{}_{𝐪\mathrm{𝟎}}y_q\widehat{a}_𝐪\widehat{a}_𝐪+\mathrm{H}.\mathrm{c}.`$, $`\phi `$ is an arbitrary phase, and $`y_q`$ is given in Ref. . Using these expressions, we can show that
$$[\widehat{\mathrm{\Xi }},\widehat{\mathrm{\Xi }}^{}],[\widehat{\mathrm{\Xi }},\widehat{\psi }^{}],[\widehat{\mathrm{\Xi }},\widehat{\psi }^{}]=𝒪(1/V).$$
(6)
Lifshitz and Pitaevskii (LP) claimed that Eq. (6) is applicable even when the interaction is stronger. Their discussion is somewhat controversial because LP started from, instead of Eq. (3), the assumption that $`\widehat{\mathrm{\Xi }}`$ could be defined by $`\widehat{\mathrm{\Xi }}|N,\nu =\mathrm{\Xi }|N1,\nu ,`$ where $`|N,\nu `$ denotes any eigenstate that has exactly $`N`$ bosons. However, we note that for weakly-interacting bosons we have not used this assumption in the derivation of Eq. (6). We thus expect that Eq. (6) also holds for bosons with stronger interaction, even if LP’s assumption was too strong. If this is the case, the following results are applicable not only to weakly interacting bosons but also to bosons with stronger interaction, because the results will be derived only from Eqs. (1)-(6).
Since we are studying the robustness against weak perturbations, we assume that $`\lambda `$ is small, so that $`J`$ is very small. In this case, we have to consider transitions only among $`|N,\mathrm{G}`$’s with different $`N`$’s (i.e., we can neglect transitions to excited states). Hence, the reduced density operator $`\widehat{\rho }`$ can be generally written as
$$\widehat{\rho }(t)=\underset{N,M}{}\rho _{NM}(t)|N,\mathrm{G}M,\mathrm{G}|.$$
(7)
It seems almost obvious that quantum coherence between $`|N,\mathrm{G}`$ and $`|M,\mathrm{G}`$ with large $`|NM|`$ would be destroyed by the interaction with the environment. We therefore study the most interesting case where $`\rho _{NM}`$ is localized in the $`N`$-$`M`$ plane in such a way that $`\sqrt{\delta N^2}N`$. If this relation is satisfied at $`t=0`$, it is also satisfied for all $`tt_{\mathrm{eq}}`$. We also assume that $`𝒱`$ is large enough, so that the boson density in the environment ($`=[N(0)N(t)]/𝒱`$) is negligibly small for all $`tt_{\mathrm{eq}}`$. Under these conditions, we can calculate the time evolution of $`\rho _{NM}(t)`$, using Eqs. (1)-(6), as
$`\rho _{NM}(t+\mathrm{\Delta }t)=e^{i(NM)\mu (n(t))\mathrm{\Delta }t/\mathrm{}}`$ (8)
$`\times [\rho _{NM}(t)(1(N+M)j(n(t))\mathrm{\Delta }t/2)`$ (9)
$`+\rho _{N+1,M+1}(t)\sqrt{(N+1)(M+1)}j(n(t))\mathrm{\Delta }t]+𝒪(\lambda ^4),`$ (10)
for a finite time interval $`\mathrm{\Delta }t`$ that satisfies $`\mathrm{}/E_c\mathrm{\Delta }t<1/Nj(n)`$, where $`E_c`$ is the energy scale over which the matrix elements of $`\widehat{H}_{\mathrm{SE}}`$ are non-negligible. Here, $`nN/V`$, and $`\mu (n)`$ ($`>0`$ for a condensate of interacting bosons ) denotes the chemical potential of bosons in the box. Furthermore,
$$j(n)=K\frac{2\pi }{\mathrm{}}\frac{n_0}{n}|\lambda |^2\frac{v^2}{V}D(\mu (n)),$$
(11)
where $`D(\mu )`$ is the density of states per unit volume of the environment at energy $`\mu `$, $`v`$ is the total volume of the walls of the box, and $`K`$ is a constant of order unity. Both $`\mu `$ and $`j`$ depend on $`N`$ through $`n`$, but this dependence is very weak because a change of $`N`$ by 1 only causes the change of $`n`$ by $`1/V`$. Note that our basic equation (9) has only two parameters, $`\mu `$ and $`j`$. Namely, all model-dependent parameters (details of $`\widehat{H}`$, $`\widehat{H}_\mathrm{E}`$ and $`\widehat{H}_{\mathrm{SE}}`$) are absorbed in these two parameters. Therefore, the following results are general and model independent.
Using Eq. (9), we first calculate the time evolutions of the expectation value $`N`$ and the fluctuation $`\delta N^2`$ of the number $`N`$ of bosons in the box. We find
$$\frac{\mathrm{d}}{\mathrm{d}t}N=j(n)N.$$
(12)
Hence, $`N`$ decreases gradually because of the leakage flux $`J=j(n)N`$. For $`\delta N^2`$, on the other hand, we find
$$\frac{\mathrm{d}}{\mathrm{d}t}F=j(n)[1F],$$
(13)
where $`F\delta N^2/N`$ is the “Fano factor” . It is seen that a robust state must have $`F=1`$, whereas any states with $`F1`$ are fragile in the sense that their $`F`$ evolves with time, approaching unity. For example, the ground-state wave function in the Bogoliubov approximation, $`|\mathrm{Bog},\mathrm{G}`$, has $`F>1`$ . Hence, it is fragile. The ground-state wave function with a fixed number of bosons, $`|N,\mathrm{G}`$, is also fragile because $`F=0`$.
Since the evaluation of $`F`$ is easy, $`F`$ is a convenient tool for the investigation of the robustness. However, since $`F`$ is only related to the diagonal elements of $`\widehat{\rho }`$, it does not distinguish between pure and mixed states. Therefore, we now solve the basic equation (9) for various initial states to investigate the robustness of the wave functions in more detail. When the initial state is a pure state of the NSIB, i.e., $`\widehat{\rho }(0)=|N,\mathrm{G}N,\mathrm{G}|`$, then $`\rho _{NM}`$ after a short interval $`\mathrm{\Delta }t`$ is evaluated as $`\rho _{NN}(\mathrm{\Delta }t)=(1Nj(n(t))\mathrm{\Delta }t)`$, $`\rho _{N1,N1}(\mathrm{\Delta }t)=Nj(n(t))\mathrm{\Delta }t`$, and other elements are zero. Therefore, $`\widehat{\rho }`$ becomes a classical mixture of $`|N,\mathrm{G}`$ and $`|N1,\mathrm{G}`$ at $`t=\mathrm{\Delta }t`$, in consistency with the above result that states with $`F=0`$ are fragile. By evaluating the evolution at later times, we find that $`\widehat{\rho }`$ evolves toward a Poissonian mixture of $`|N,\mathrm{G}`$’s , consistent with $`F1`$. In a similar manner, we can show that the pure state of Bogoliubov’s ground state $`\widehat{\rho }(0)=|\mathrm{Bog},\mathrm{G}\mathrm{Bog},\mathrm{G}|`$, which has $`F>1`$, also evolves into a mixed state. We can also show that the number-phase squeezed state of interacting bosons (NPIB), which was found in Ref. as a number-phase minimum uncertainty state with $`0<F<1`$, also evolves into a mixed state. These examples show that $`F`$ is indeed a simple measure of the robustness: A pure state with $`F1`$ is unlikely to remain pure. Note, however, that a pure state with $`F=1`$ is not necessarily robust. For example, we can show that the coherent state of free bosons (CSFB) evolves into a mixed state, although it has $`F=1`$. Hence, $`F=1`$ is only a necessary condition for the robustness.
Among many states with $`F=1`$, we have successfully found a very special state that is robust in the sense that it remains pure when it is weakly perturbed by the environment. The state is given by
$$\widehat{\rho }(t)=|\alpha (t),\mathrm{G}\alpha (t),\mathrm{G}|.$$
(14)
Here, $`\alpha (t)`$ is a time-dependent complex number given by
$$\alpha (t)=e^{i\phi (t)}\sqrt{N(t)},$$
(15)
where $`N(t)`$ is the solution of Eq. (12), and
$$\phi (t)=\phi (0)\frac{i}{\mathrm{}}_0^t\mu (n(t))dt.$$
(16)
Here, the initial phase $`\phi (0)`$ is arbitrary, and $`n(t)N(t)/V`$. Furthermore,
$$|\alpha ,\mathrm{G}e^{|\alpha |^2/2}\underset{M=0}{\overset{\mathrm{}}{}}\frac{\alpha ^M}{\sqrt{M!}}|M,\mathrm{G},$$
(17)
which we call the coherent state of interacting bosons (CSIB). It has the same form as the CSFB except that $`|M,\mathrm{G}`$ is the NSIB. Because of this difference, simple relations for the CSFB do not hold for the CSIB. For example, $`\alpha ,\mathrm{G}|\widehat{\psi }|\alpha ,\mathrm{G}\alpha /\sqrt{V}`$, and, moreover, $`|\alpha ,\mathrm{G}`$ is not an eigenstate of $`\widehat{\psi }`$. Nevertheless, $`\alpha ,\mathrm{G}|\widehat{N}|\alpha ,\mathrm{G}=\alpha ,\mathrm{G}|\delta \widehat{N}^2|\alpha ,\mathrm{G}=|\alpha |^2`$, hence $`F=1`$ exactly, as in the case of CSFB. Since the NSIB has a complicated wave function, so does the CSIB. \[For weakly-interacting bosons, its explicit form was given in Ref. .\] Although complicated, the wave function of the CSIB is robust against weak perturbations from the environment: It keeps the same form, whose parameter $`\alpha (t)`$ evolves slowly (except for the phase rotation), and remains a pure state, in contrast to other wave functions which soon evolve into mixed states. In fact, Eqs. (14)-(17) yield
$`\rho _{NM}(t)`$ $`=`$ $`e^{N(t)}{\displaystyle \frac{e^{i(NM)\phi }}{\sqrt{N!M!}}}N(t)^{\frac{N+M}{2}},`$ (18)
$`\rho _{NM}(t+\mathrm{\Delta }t)`$ $`=`$ $`e^{N(t)(1j\mathrm{\Delta }t)}{\displaystyle \frac{e^{i(NM)(\phi \mu \mathrm{\Delta }t/\mathrm{})}}{\sqrt{N!M!}}}`$ (20)
$`\times N(t)^{\frac{N+M}{2}}(1j\mathrm{\Delta }t)^{\frac{N+M}{2}},`$
which indeed satisfy Eq. (9).
We now discuss the compatibility with the SSR, which might raise the objection that the CSIB would not be realized because superpositions between states with different values of $`N`$ are forbidden for massive bosons. To show that this intuitive objection is wrong, it is sufficient to give one counterexample. Suppose that there is another box, which also contains condensed bosons, in the same room. The total system consists of two boxes and the environment. According to the SSR, the wave function of the total system $`|\mathrm{\Phi }_{\mathrm{total}}`$ should be a superposition of states that have the same number of bosons, $`N_{\mathrm{total}}=N+N^{}+N_\mathrm{E}=`$ fixed, where $`N^{}`$ denotes the number of bosons in the second box. Consider the following state, which satisfies this constriction;
$`|\mathrm{\Phi }_{\mathrm{total}}`$ $`=`$ $`{\displaystyle \underset{N,N^{},\mathrm{}}{}}e^{|\alpha |^2/2|\alpha ^{}|^2/2}\alpha ^N\alpha _{}^{}{}_{}{}^{N^{}}C_{\mathrm{}}/\sqrt{N!N^{}!}`$ (21)
$`\times `$ $`|N,\mathrm{G}|N^{},\mathrm{G}^{}|N_{\mathrm{total}}NN^{},\mathrm{}_\mathrm{E}.`$ (22)
Here, $`\alpha =|\alpha |e^{i\phi }`$, $`\alpha ^{}=|\alpha ^{}|e^{i\phi ^{}}`$, and $`C_{\mathrm{}}`$ is a complex number, where $`\mathrm{}`$ is a quantum number labeling states of the environment $`|M,\mathrm{}_\mathrm{E}`$ which has $`M`$ bosons. Regarding the phases $`\phi `$ and $`\phi ^{}`$, only the relative value $`\phi \phi ^{}\theta `$ has a physical meaning. We thus take $`\phi ^{}=0`$ henceforth. Equation (22) yields the reduced density operator of the first box as $`\widehat{\rho }=_Ne^{|\alpha |^2}(|\alpha |^{2N}/N!)|N,\mathrm{G}N,\mathrm{G}|`$. It is easy to show that this is identical to
$$\widehat{\rho }=_\pi ^\pi \frac{d\theta }{2\pi }||\alpha |e^{i\theta },\mathrm{G}|\alpha |e^{i\theta },\mathrm{G}|.$$
(23)
Although this $`\widehat{\rho }`$ represents a mixed state of CSIB’s, we note that it does not contain the maximum information on the state in the box, whereas the best density operator should have the maximum information. The lacking information is that the phase relative to the condensate in the second box is $`\phi `$. Hence, the maximum information is Eq. (23) with the restriction $`\theta =\phi `$. This combined information is concisely expressed as $`\widehat{\rho }=||\alpha |e^{i\phi },\mathrm{G}|\alpha |e^{i\phi },\mathrm{G}|,`$ which agrees with Eq. (14). Namely, Eq. (14) is better than Eq. (23) because it contains more information. This example demonstrates that Eq. (14) can be compatible with the SSR in realistic cases where the box exchanges bosons with the environment. Only in the limiting case where the box is completely closed, should the SSR be crucial, and the NSIB would be realized if the temperature $`T0`$ .
We have established that the CSIB is a robust pure state of interacting many bosons. We finally discuss its implications. The robustness of the present work should not be confused with the “stiffness of macroscopic wave functions” , which only referes to the stability of an order parameter in a mean field approximation. For example, $`|\mathrm{Bog},\mathrm{G}`$ has the stiffness , whereas it is fragile as we have shown. The robustness is generalization of the robustness of coherent states of $`f=1`$ systems . It is thus natural to expect for $`f1`$ systems that some coherent state would be robust. However, it was not known which coherent state is robust: there are many choices of the coordinate by which a coherent state is defined. Since Eq. (3) yields
$$(\widehat{\mathrm{\Xi }}/\xi )|\alpha ,\mathrm{G}=\alpha |\alpha ,\mathrm{G},$$
(24)
the present work has revealed that the robust coherent state is the one defined by $`\widehat{\mathrm{\Xi }}/\xi `$. In this sense, $`\widehat{\mathrm{\Xi }}+\widehat{\mathrm{\Xi }}^{}`$ is the “natural coordinate” of interacting many bosons. The condensation of bosons are often characterized by the off-diagonal long-range order (ODLRO) that is defined by $`\widehat{\psi }^{}(𝐫)\widehat{\psi }(𝐫^{})=`$ finite for $`|𝐫𝐫^{}|V^{1/3}`$ . Using Eqs. (2)-(5) we can show that the CSIB, NSIB, NPIB, and the Bogoliubov’s ground state all have the ODLRO. Hence, the present work has revealed that the ODLRO does not necessarily imply the robustness. Furthermore, all of these states have almost the same energy densities, i.e., the differences of $`\widehat{H}/V`$ are only $`𝒪(1/V)`$ for the same value of $`N`$ . For example, if we let $`E_{N,G}`$ be the eigenenergy of the NSIB, $`\widehat{H}|N,G=E_{N,G}|N,G`$, we can then easily show from Eq. (17), neglecting terms of $`𝒪(1/V)`$, that $`\alpha ,\mathrm{G}|\widehat{H}|\alpha ,\mathrm{G}/V=E_{|\alpha |^2,G}/V=E_{N,G}/V`$. Therefore, the robustness of the CSIB is not due to an energy difference, but to natures of wave functions. Since interactions with the environment are finite in most physical systems, we argue that only the robust state, CSIB, should be realized as a macroscopic pure state. Since the (relative) phase of the CSIB is almost definite , the global gauge symmetry is then broken. Although $`V`$ is finite, we are thus led to the symmetry breaking by considering the robustness. This suggests that quantum phase transitions may have more profound origins than singularities that are developed as $`V\mathrm{}`$. A conventional trick to get symmetry breaking states for boson condensates is to introduce a symmetry breaking field $`\eta `$, which couples to $`\widehat{\psi }`$ as $`\widehat{H}_\eta =d^3𝐫(\eta ^{}\widehat{\psi }+\eta \widehat{\psi }^{})`$. However, $`\eta `$ is usually considered as an unphysical field , and it was sometimes argued that symmetry breaking states were meaningless because they look against the SSR . In contrast, the present work gives a physical reasoning for the symmetry breaking, assuming only physical interactions, and shows the compatibility with the SSR.
This work has been supported by the CREST program of Science and Technology Corporation of Japan.
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# The CP-PACS Project and Lattice QCD Results
## 1 Introduction
Lattice QCD is a fundamental theory of quarks and gluons which are constituents of hadrons such as protons and pions. Numerical studies of lattice QCD have developed significantly during the past decade in parallel with the development of computers. Of particular importance in this regard has been the construction of dedicated QCD computers (see for reviews Ref.?) and the move of commercial vendors toward parallel computers in recent years. In Japan the first dedicated QCD computer was developed in the QCDPAX project . The QCDPAX computer with a peak speed of 14GFLOPS is actually the 5th computer in the PAX project , which pioneered the development of parallel computers for scientific and engineering applications in Japan.
The CP-PACS project was conceived as a successor of the QCDPAX project in the early summer of 1991. The project name CP-PACS is an acronym for Computational Physics by a Parallel Array Computer System. The aim of the project was to develop a massively parallel computer for carrying out research in computational physics with primary emphasis on lattice QCD.
In this article after a brief description of lattice QCD and the background of the project in Sec.2, we present an overview of the CP-PACS project in Sec.3, and describe characteristics of the CP-PACS computer in Sec.4. The performance of the computer for lattice QCD applications as well as for the LINPACK benchmark are also given. Main results in lattice QCD are given in Sec.5. Sec.6 is devoted to conclusions.
## 2 Lattice QCD and Background of the Project
Lattice QCD is a fundamental theory of quarks and gluons defined in terms of the path-integral formalism of quantum theory on a 4-dimensional hyper-cubic lattice. The lattice spacing plays a role of an ultra-violet cutoff. The infinite volume limit and the continuum limit should be taken in order to get physical quantities.
As we have to treat quarks and gluons relativistically, we have a problem in 4-dimension in stead of 3-dimension as in solid-state problems. However, except this difference of dimensionality, it is a statistical system. Quarks are defined on sites, while gluons on bonds of a 4-dimensional hyper-cubic lattice. Numerical methods we employ are a Monte Carlo method, a molecular dynamics and a hybrid method of combination of these methods. However, due to this dimensionality, we need a lot of CPU time and a large memory size.
Because of this requirement of high performance computers for numerical simulations in lattice QCD, dedicated machines have been constructed in USA, Europe and Japan. There are two additional reasons why dedicated parallel computers were widely developed for lattice QCD: First there is an incentive to perform first-principle calculations without introducing any approximations based on the fundamental law. Second there is a spiritual atmosphere in high-energy physics community to construct a special purposed equipment like an accelerator. A massively parallel computer is an accelerator for numerical simulations.
Fig. 1 shows the recent development of the computers in terms of the theoretical peak speed versus the year when the computer was shipped or constructed. Small open symbols are for vector-type supercomputers, while large and small filled symbols are for dedicated parallel and commercial parallel computers, respectively. Open circles with dot are for QCDPAX and filled large circles are for CP-PACS. We clearly observe that the rate of the progress for parallel computers is roughly double that of vector computer and that a crossover in the peak speed took place from vector to parallel computers around 1991. For this development dedicated machines for lattice QCD made important roles.
## 3 CP-PACS Project
The CP-PACS Project aimed at developing a massively parallel computer designed to achieve high performance for numerical research of the major problems of computational physics, and it further aimed at significant progress in the solution of these problems through the application of the parallel computer upon completion of its development.
The Project formally started in April of 1992, and continued for five years, until March of 1997. The Project received about 2.2 billion Yen spread over the five year period. The Center for Computational Physics was founded in April 1992 at University of Tsukuba to carry out the Project, as well as to promote research in computational physics and parallel computer science. The Center is an inter-university facility open to researchers in academic institutions in Japan.
The Project members consist of 15 computer scientists and 18 physicists, as listed in Table I. As Table I clearly shows, the CP-PACS Project is a multi-disciplinary effort toward the advancement of computational physics encompassing not only several branches of physics but also computer science to develop parallel computers best suited for such applications. The Projected was headed by Y. Iwasaki. The development of the CP-PACS computer was led by K. Nakazawa.
A unique feature of the Project is its emphasis on cross-disciplinary research involving both physicists and computer scientists. This is a tradition carried over from the QCDPAX Project , which is the predecessor and stepping stone for the CP-PACS Project. A close collaboration of researchers from the two disciplines has been both important and fruitful in reaching a design for the CP-PACS computer which best balances the computational needs of physics applications with the latest of computer technologies.
Development of a massively parallel computer requires advanced semiconductor technology. We selected Hitachi Ltd. as the industrial parter through a formal bidding process in the early summer of 1992, and we worked in a close collaboration for the hardware and software development of the CP-PACS computer. The first stage of the CP-PACS computer consisting of 1024 processing units with a peak speed of 307 GFLOPS was completed in March 1996. An upgrade to a 2048 system with a peak speed of 614GFLOPS was completed at the end of September 1996
## 4 CP-PACS Computer
### 4.1 Architecture
The CP-PACS computer is an MIMD (Multiple Instruction-streams Multiple Data-streams) parallel computer with a theoretical peak speed of 614GFLOPS and a distributed memory of 128 Gbytes. The system consists of 2048 processing units (PU’s) for parallel floating point processing and 128 I/O units (IOU’s) for distributed input/output processing. These units are connected in an 8$`\times `$17$`\times `$16 three-dimensional array by a three-dimensional crossbar network. The specification of the CP-PACS computer is summarized in Table II. A well-balanced performance of CPU, network and I/O devices supports the high capability of CP-PACS for massively parallel processing.
The basic strategy we adopted for the design is the usage of a fast RISC micro-processor for high arithmetic performance at each node and a linking of nodes with a flexible network so as to be able to handle a wide variety of problems in computational physics. The unique features of the CP-PACS computer reflecting these goals are represented by the special node processor architecture called pseudo vector processor based on slide-windowed registers (PVP-SW) and the choice of a three-dimensional crossbar network.
### 4.2 Node processor
Each PU of the CP-PACS has a custom-made superscalar RISC processor with an architecture based on PA-RISC 1.1. In large scale computations in scientific and engineering applications on a RISC processor, the performance degradation occurring when the data size exceeds the cache memory capacity is a serious problem. For the processor of CP-PACS, an enhancement of the architecture called the PVP-SW was developed to resolve this problem, while still maintaining upward compatibility with the PA-RISC architecture.
### 4.3 Network
The 2048 processors are arranged in a three-dimensional $`8\times 16\times 16`$ array. The Hyper Crossbar network is made of crossbar switches in the $`x,y`$ and $`z`$ directions, connected together by an Exchanger at each of the three-dimensional crossing points of the crossbar array. Each exchanger is connected to a PU or IOU. Thus any pattern of data transfer can be performed with the use of at most three crossbar switches. Since the network has a huge switching capacity due to the large number of crossbar switches, the sustained data transfer throughput in general applications is very high.
### 4.4 Performance
The most CPU consuming part of lattice QCD calculations is the inversion of a linear equation. We developed a hand-optimized assembler code for the core part of the solver. The performance of the calculation part is 186 MFLOPS per node, which is 62% of the peak speed. The percentage of the communication in the total is 23 %, which makes the sustained speed for the solver 148 MFLOPS. This is about a half of the theoretical peak speed.
We also measured the performance of the LINPACK benchmark. The sustained speed for the case of 2048 PU’s is 368.2 GFLOPS, which is 59.9% of the theoretical peak speed. This performance was ranked as number one of TOP 500 Supercomputers announced in November 1996.
## 5 Physics Results
### 5.1 Hadron Spectrum in Quenched QCD
Deriving the hadron spectrum from lattice QCD is a milestone to verify that QCD is the fundamental theory of quarks and gluons. Therefore, much effort has been paid to calculate the hadron spectrum since 1981 when the first attempt of the hadron spectrum calculation was made.
A simulation of QCD without approximation requires an enormous computer time. Therefore, as the first step, the quenched approximation, in which pair creations and annihilations of quarks in the vacuum are ignored, has been employed in major simulations of QCD. However, even in the quenched approximation, it is not easy to obtain precise values of the hadron spectrum. We have to first control and then estimate various systematic errors characteristic of lattice QCD, i.e., the errors due to the infinite volume limit and the continuum limit. Moreover, it is technically difficult to simulate directly at the realistic values of light u and d quark masses, as the CPU time is proportional to the inverse of the quark mass. Therefore we have to extrapolate results obtained at relatively heavy quark masses to the light quark mass. This introduces another source of systematic errors.
In early works, it was difficult to employ large enough lattices with small enough lattice spacings, mainly due to limitation of computer power. In particular, all simulations before 1988 employed lattices much smaller than 2 fm which is the size of typical hadrons. Therefore old calculations suffer from large systematic errors. Simulations at light enough quark mass were also difficult due to algorithm adopted and the speed of computers at that time.
The best calculation prior to the CP-PACS was performed by the GF11 collaboration in 1992-1993 using their dedicated computer GF11. Performing systematic extrapolations in terms of quark masses and lattice spacing, supplemented by corrections from the finite lattice size, they determined the quenched hadron spectrum in the continuum limit. They concluded that the hadron masses in the quenched QCD are consistent with experiment within their errors, which is typically about 10%.
As the first physics project on the CP-PACS, we aimed to obtain final results for the hadron spectrum in the quenched QCD with errors of a few % level and thereby clarify the long standing issue of the magnitude of quenching errors. Simulation parameters were chosen by taking this goal into consideration.
From these simulations together with detailed systematic analyses, we succeed to determine the quenched hadron spectrum with errors about 1-2 % for mesons and 2-3 % for baryons. We were also able to much reduce various systematic errors and estimate them. This is crucial to obtain reliable numerical results. Thus we are able to establish the hadron spectrum in the quenched QCD.
In Fig.2, our results for the quenched spectrum together with experiment are shown. The experimental values of the $`\pi ,\rho `$ and $`K`$ or $`\varphi `$ masses are employed to fix the physical scale and the light quark masses.
Our results unambiguously establish a discrepancy between the quenched hadron masses and the experimental values, with up to 7$`\sigma `$ for several particles. On the other hand, the magnitude of the discrepancy is at most 10%, which is consistent with phenomenological estimates of the quenching error.
### 5.2 Hadron Spectrum in Full QCD
Since the quenched hadron mass spectrum exhibits deviation from experiment, the next step is to perform calculation of QCD without the quenched approximation (the full QCD calculation). As a step toward this goal, we have started QCD simulations taking into account of effects of pair creation and annihilation of light u, d quarks. We treat the heavier s quark in the quenched approximation.
Simulations in full QCD need computer power at least 100 times larger than that in the quenched QCD. Therefore, it is impossible to simply repeat the simulation in full QCD like that described above in the quenched QCD. In order to overcome this problem, we adopt an improved action, which is a lattice action modified in such a way that systematic errors due to finite lattice spacing is reduced.
We first made a pilot study to investigate the effects of improving using various improved actions and found that the combination of the renormalization-group improved action for gluons and the clover action significantly reduces errors due to the finite lattice discretization over the standard action. We adopt this combination of improved actions in our production runs.
A systematic study of the mass spectrum in full QCD is in progress. We have already found several interesting effects of dynamical quarks in the hadron spectrum. In Fig.3 we compare meson masses in full QCD with those in the quenched QCD. It clearly shows that in the continuum limit (the point where the lattice spacing a is zero) the discrepancies of $`K^{}`$ and $`\varphi `$ meson masses from experiment observed in the quenched QCD are significantly reduced in full QCD.
### 5.3 Quark Masses
The masses of quarks are the very fundamental parameters in nature like the mass of the electron. However, because quarks are confined in hadrons, one cannot determine their masses directly from experiment. Usually, their values have been theoretically inferred from experimental hadron masses using phenomenological models of QCD. Lattice QCD is the only known way to determine the masses of quarks from first principles.
We made systematic calculations of quark masses both in the quenched QCD and in full QCD. In Figs. 4 and 5 we show the lattice spacing dependence of the average u, d quark mass and the s quark mass, respectively.
On the lattice there are alternative definitions of the quark mass. Although the values of the quark mass differ depending on the definition at finite lattice spacing, they extrapolate to a common value in the continuum limit. The verification of the unique value in the continuum was first made in the quenched QCD in Ref.?. This verification is important because the quark mass should be the fundamental parameter in QCD.
The s quark mass is determined using experimental values of either $`K`$ meson mass or $`\varphi `$ meson mass. The s quark mass in the quenched approximation depends on the choice of input. This reflects a systematic error of quenching.
The discrepancy is found to be much reduced in our full QCD calculations. The values of the s quark mass from $`K`$ meson mass or $`\varphi `$ meson mass are consistent within one standard deviation; 90(10) MeV. This value is significantly smaller than that in the quenched QCD; 120-140 MeV. The value 90(10) MeV for the s quark mass has a significant implication for the analysis of the CP violation.
For the clarification of the CP violation in nature we need a theory like the Kobayashi-Maskawa theory, an experiment result like that from a B factory, and also numerical results from lattice QCD. This is a typical example of cases where results from three fields of theoretical physics, experimental physics and computational physics are necessary to solve a problem.
## 6 Conclusions
It was successful to develop a massively parallel computer CP-PACS with a peak speed of 614 GFLOPS due to a close collaboration among physicists, computer scientists and a vendor. The performance of the computer for physics application is as high as 50 % of the peak speed in the case of the core part of lattice QCD programs.
We are able to obtain intersting and important results in lattice QCD using the CP-PACS computer: 1) The hadron spectrum in the quenched QCD has been established. Our results unambiguously clarify a discrepancy between the quenched hadron masses and the experimental values, with up to 7$`\sigma `$ for several particles. On the other hand, the magnitude of the discrepancy is at most 10%, which is consistent with phenomenological estimates of the quenching error. 2) The discrepancies of meson masses from experiment observed in the quenched QCD are significantly reduced in full QCD. 3)We have systematically calculate the masses of light quarks in the quenched QCD and in full QCD. In particular, the mass of the s quark in full QCD is 90(10) MeV, which is much smaller than that previously estimated phenomenologically.
## ACKNOWLEDGEMENTS
I would like to thank the members of the CP-PACS Project and the members of the CP-PACS Collaboration for lattice QCD simulations, in particular, K. Nakazawa and A. Ukawa for valuable discussions. This work is supported in part by the Grand-in-Aid of the Ministry of Education, Science and Culture (Nos. 08NP0101 and 09304029).
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# The Analysis of Cosmic Ray Data
## 1 Introduction
A statement of statistical belief not uncommon in cosmic ray work is: ”you need five sigmas to convince me”. This has some justification, in that the history of cosmic rays contains many instances when a source or effect is claimed but not subsequently substantiated. Frequently this has been due to incorrect application of some statistical technique, often a failure to account fully for the ’degrees of freedom’.
Most of the present body of statistical knowledge has been developed for specific problems, few of which occur in cosmic rays, although one of the most useful of texts was produced for experimental particle physicists. The analyser of cosmic ray data has particular problems: cosmic ray data requires great effort in collection and they are unlikely, once analysed, to be repeated. The numbers are frequently small and there are usually data missing and frequently there is significant contamination by noise. Ideally, a statistical measure should be developed specifically for each application. This is the only way in which all of the parameters of the experiment can be allowed for in the analysis. It is more usual for a general statistical tool to be applied, for example $`\chi ^2`$, which may not be optimal for the purpose, and for some experimental variables to be ignored.
The focus of this review will be on methods of determining the presence of a signal rather than estimating some parameters of the data. The aim is to gather together the recent developments in methods of analysis of the temporal and spatial features of cosmic ray data, especially where the methods used are not ’traditional’.
Several new methods have been published recently which depend on Bayesian ideas, and these ideas have been introduced before the description of the methods.
## 2 On/Off Counts
### 2.1 Introduction
The subject of detecting the presence of a source in counting rate data, using off-source control data has appeared many times. Despite these numerous airings, erroneous statistical significances are occasionally still being published. In principle the question is easy to pose: if $`N_{ON}`$ counts are detected when an instrument is pointed at a source and there is also a background counting rate, and $`N_{OFF}`$ counts are detected when it is collecting background counts only under otherwise identical conditions, what is the likelihood that there is a genuine source?
A common treatment is to give for the significance of the excess counts:
$$N_{SIGMA}=\frac{N_{ON}\alpha N_{OFF}}{\sqrt{N_{ON}+\alpha ^2N_{OFF}}}$$
(1)
where $`\alpha `$ is the ratio of time on-source to time off-source, $`t_{ON}=\alpha t_{OFF}`$. This is based on the supposition that the best estimate of the observed ’signal’ is $`N_{ON}\alpha N_{OFF}`$, the variances of the ON and OFF counts for a Poisson distribution are $`N_{ON}`$ and $`N_{OFF}`$ and the variance of the difference between $`N_{ON}`$ and $`\alpha N_{OFF}`$ is the weighted sum of the variances. The statistic used is $`Student^{}s`$ $`t`$ which in the limit of large numbers is Gaussian. Since the distributions of $`N_{ON}`$ and $`N_{OFF}`$ are Poissonian, this expression should be used only if the numbers of events is sufficiently large for a Gaussian approximation to Poissonian to be valid.
It is an example of only one type of statistic which could be used in $`ON/OFF`$ situations - a *goodness-of-fit* statistic to determine whether the observed data could have arisen from an *a priori* distribution. Other statistics could have been used, for example $`\chi ^2`$, which in this instance would have one degree of freedom. Asymptotically they should have the same result, that is they both should reject or accept the null hypothesis equally. In these tests the *null hypothesis* is that the observations were both samples from the same population and that any difference arose merely by chance. There is no explicit alternative hypothesis, but an implicit one: that if the difference between the counts was unlikely to be due to chance, it arose because of a genuine source, strength unspecified.
### 2.2 Likelihood Analysis
An optimal test exists for the intermediate case where there are two completely specified hypotheses: $`H_0`$: the *null hypothesis* as described above, and $`H_1`$: a hypothesis involving another model, usually including a specific ’signal’. In this rare (in cosmic rays) case, the Neyman/Pearson theorem shows that the likelihood ratio is optimal for any distribution function for the errors.
In the more usual case, $`H_1`$ is not fully specified, but has one or more free parameters. The null hypothesis $`H_0`$ is that $`N_{OFF}`$ and $`N_{ON}`$ are both samples of the same population for which the source strength $`S=0`$. The alternative hypothesis $`H_1`$ is that $`N_{ON}`$ contains an *unknown* source component, $`S>0`$. In this case there is no optimal test, except that for errors of the exponential family, such as a Gaussian, the likelihood ratio is expected to be near-optimal.
The problem was discussed at length twenty five years ago by O’Mongain and Hearn but was not solved satisfactorily, at least in this field, until the maximum likelihood treatment of Gibson et al. and Dowthwaite et al. and later by a similar treatment by Li and Ma. In these treatments the observed $`ON`$ and $`OFF`$ counts are due to (i) an unknown background $`B`$ plus an unknown source $`S`$ and (ii) the same unknown background $`B`$ alone. The likelihood ratio is maximised with respect to the possible source counts:
$`\lambda `$ $`=`$ $`\left({\displaystyle \frac{P\left(N_{ON},N_{OFF}S=0\right)}{P\left(N_{ON},N_{OFF}S=N_{ON}\alpha N_{OFF}\right)}}\right)`$ (2)
$`=`$ $`\left[{\displaystyle \frac{\alpha }{1+\alpha }}\left({\displaystyle \frac{N_{ON}+N_{OFF}}{N_{ON}}}\right)\right]^{N_{ON}}\left[{\displaystyle \frac{1}{1+\alpha }}\left({\displaystyle \frac{N_{ON}+N_{OFF}}{N_{OFF}}}\right)\right]^{N_{OFF}}`$
A standard result is that the probability of obtaining a given $`\lambda `$ is obtained from
$$2\mathrm{ln}(\lambda )\chi ^2(1)$$
(3)
### 2.3 Comparison of methods
Both equations 1 and 3 are valid asymptotically: only for large values of $`N_{OFF}`$ and $`N_{ON}`$. Equation 1 assumes that the error distributions of $`N_{ON}`$ and $`N_{OFF}`$ are gaussian and equation 3 assumes that $`2\mathrm{log}(\lambda )`$ is distributed as $`\chi _2(1)`$. To check the region of validity, random data sets have been generated for each of a number of values of $`N_{OFF}`$. For each data set $`\alpha `$ has been set to $`1.0`$ and a value of $`S_3=3\sqrt{2B+S_3}`$ has been calculated, which is the $`3\sigma `$ value of $`S=S_3`$ assuming the validity of equation 1. At each value of $`N_{OFF}`$, $`10^6`$ data sets were generated using a poissonian random number algorithm. The fraction of samples $`N`$ of $`N_{OFF}`$ where $`N>N_{OFF}+S_3`$ is used as an estimate of the true probability of obtaining $`N_{ON}=N_{OFF}+S_3`$ by chance. The results are shown in figure 1 where both equations 1 and 3 are shown to overestimate the probability near the $`3\sigma `$ level almost equally likely, the former slightly less so. It is evident that, near the $`3\sigma `$ level, there is little to choose and both equations are adequate for values of $`N_{OFF}`$ and $`N_{ON}`$ of a few hundred or more. Since good algorithms are available for Poissonian random number generation it is likely to be better to determine the probability of $`N_{ON}`$ and $`N_{OFF}`$ for values less than $`100`$ using Monte Carlo methods tailored for the exact values of $`N_{ON}`$ and $`N_{OFF}`$.
## 3 Time Series
### 3.1 Introduction
Time series analysis has been the subject of very many books and articles and has been applied in very many fields. The term covers a wide range of concepts, including Change Point Analysis, Fourier Analysis and Trend Analysis. In cosmic ray studies, there are several areas of application, such as sidereal/solar effects on low energy cosmic rays on the ground, periodicity in data from point sources, either from satellite X- and $`\gamma `$-ray data, or from ground-based Čerenkov detectors, and sporadic emission of a wide range of cosmic ray energies. In these cases, the raw data is usually in the form of time-tagged events.
### 3.2 Bursts of Events
This section will be concerned with the problem of deciding whether the counting rate of a detector has deviated from the expected rate due to a real outburst of events. The problem is usually most difficult in data comprising time-tagged events. An initial analysis could start with binning the data and looking for a deviation from the expected Poissonian distribution of the counts. One problem with this approach is that in the model of a single Poisson process generating the counts, each bin is independent, the experimenter often has the freedom to place the bins, both in position and width, arbitrarily. This alters the ’degrees of freedom’ and experience suggests that more bursts have been ’detected’ in the past than could have been justified from the data.
The problem mentioned above is a specific one but in general most statistical problems associated with sporadic emission relate to the lack of a specific model for the form, duration and amplitude of emission, and the feeling is often that, given a free hand with the parameters, any pure noise series could made to disclose a ’burst’. A recent paper by Scargle suggests that existing methods for searching for rapid variability in $`X`$ray and $`\gamma `$-ray astronomy do not fully extract all of the information contained in photon counts. The reasons given included ’binning fallacies’, in that the data were widely binned and the size of the bins must be large enough to give ’good statistics’. Further, global methods such as autocorrelations and power spectra used on large data sets dilute the effects of sporadic bursts. The Bayesian response to these problems is discussed later.
The problem at first sight does not seem insolvable using classical statistical theory. The statistical treatment of *point processes*: data occurring as points on the real line, or as discrete times, is covered by several texts, for example Cox and Isham. The general treatment covers a variety of statistical processes, including Poisson (which is of most application here), doubly stochastic Poisson (where the average Poisson rate is itself a variable) and renewal processes where the distribution function for intervals between points is not exponential. In analysing data in the form of time-tagged photons without appreciable dead time, classical statistics would look for a powerful goodness-of-fit test of the pure Poisson process, if possible avoiding the loss of information and the arbitrary choices associated with binning.
Given such a series of times, the problem posed here is: is there evidence for ’bunching’ or ’bursts’? Alternatively, are the data consistent with a uniform distribution in time which, for events not affected by counter dead-time, would be governed by a pure Poisson process? Some recent papers such as McLaughlin et al. use just this assumption to classify sources into ’steady’ or ’variable’. Others use *ad hoc* methods to estimate the probability of bursts.
#### 3.2.1 The Scan Statistic
The test statistic postulated above exists: the Scan Statistic has been extensively studied by Parzen, Barton and David, Huntington and Naus, Neff and Naus, Naus, Glaz, Wallenstein, Naus and Glaz, Chen and Glaz and Månsson . It is a statistic for detecting clustering in time or one dimension in space. It is usually described as the maximum (or minimum) number of events which can be found in a window of fixed duration scanning smoothly through a much longer interval containing discrete events following some random process, for example Poissonian. An example of the scan statistic is shown in figure 1 for window lengths of 1% and 10% of the duration of the data. The random test data has a constant mean rate except for the third quarter which has double the rate.
The scan statistic $`S`$ has a probability $`P(S)`$ which depends on the rate of events, the duration and the width of the scanning window. Some exact solutions for the probability $`P(S)`$ have been provided. One of them, by Huntington and Naus, provides the probability of a related statistic:
$$P\left(a_na\right)=1\underset{Q}{}Rdet1/h_{ij}!det1/l_{ij}!$$
(4)
where $`a_n`$ is the smallest interval $`a`$ containing $`n`$ events in the range $`[0,1]`$, where this range contains $`N`$ events in all. The summation extends over the set $`Q`$ of all partitions of $`N`$ into $`2L+1`$ integers satisfying $`n_i+n_{i+1}<n,i=1,\mathrm{},2L`$ and $`R=N!b^M(ab)^{NM}`$ with $`M=_{k=0}^Ln_{2k+1}`$, and
$`h_{ij}`$ $`=`$ $`{\displaystyle \underset{k=2j1}{\overset{2i1}{}}}n_k(ij)nL+1ij1`$
$`=`$ $`{\displaystyle \underset{k=2i}{\overset{2j2}{}}}n_k+(ji)n\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}i<jL+1`$
$`l_{ij}`$ $`=`$ $`{\displaystyle \underset{k=2j}{\overset{2i}{}}}n_k(ij)nLij1`$
$`=`$ $`{\displaystyle \underset{k=2i+1}{\overset{2j1}{}}}n_k+(ji)n\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}i<jL`$
Equation 4, although exact, is computationally expensive for large $`N`$ and small $`a`$, that is a large data set with a small scanning window, but several approximations have been provided which are designed to be valid for certain combinations of parameters.
#### 3.2.2 Newell-Ikeda Approximation for the Scan Statistic
The Newell-Ikeda asymptotic formula is suitable for small probabilities. It gives the probability of finding a section of length $`t`$ in a data set of length $`T`$, given a Poisson process of average rate $`\lambda `$:
$$P(n;\lambda T,t/T)1\mathrm{exp}\left(\lambda ^nt^{n1}T/(n1)!\right)$$
(5)
As shown in table 1, it significantly overestimates larger probabilities.
Better approximations, although not as easy to calculate, are available, for example Conover, Bement and Iman and Naus.
#### 3.2.3 Naus Approximation for the Scan Statistic
The more exact treatment of Naus will be given without derivation. For an average rate of events $`\lambda `$, data of total duration $`T`$ and scanning window of duration $`t`$, define $`L=T/t`$. Then the probability that the number of events in a scanning window never exceeds $`n`$ is $`Q^{}(n;\lambda L,1/L)`$ and is accurately approximated by:
$$Q^{}(n;\lambda L,1/L)Q^{}(n;2\lambda ,\frac{1}{2})\left[Q^{}(n;3\lambda ,\frac{1}{3})/Q^{}(n;2\lambda ,\frac{1}{2})\right]^{L2}$$
(6)
Note that this approximation is valid for a wide range of types of distribution for the time between events. Exact formulae for $`Q^{}(n;2\lambda ,\frac{1}{2})`$ and $`Q^{}(n;3\lambda ,\frac{1}{3})`$ are given for a Poisson process:
$`Q^{}(n;2\lambda ,{\displaystyle \frac{1}{2}})`$ $`=`$ $`F_{n1}^2\left(n1\right)p_np_{n2}\left(n1\lambda \right)p_nF_{n3}`$
$`Q^{}(n;3\lambda ,{\displaystyle \frac{1}{3}})`$ $`=`$ $`F_{n3}^2A_1+A_2+A_3A_4`$
where
$`A_1`$ $`=`$ $`2p_nF_{n1}\left(\left(n1\right)F_{n2}\lambda F_{n3}\right)`$
$`A_2`$ $`=`$ $`0.5p_n^2\left(\left(n1\right)\left(n2\right)F_{n3}2\left(n2\right)\lambda F_{n4}+\lambda ^2F_{n5}\right)`$
$`A_3`$ $`=`$ $`{\displaystyle \underset{r=1}{\overset{n1}{}}}p_{2nr}F_{r1}^2`$
$`A_4`$ $`=`$ $`{\displaystyle \underset{r=2}{\overset{n1}{}}}p_{2nr}p_r\left(\left(r1\right)F_{r2}\lambda F_{r3}\right)`$
and $`p_i`$,$`F_n`$ are the Poisson probability and distribution functions: $`p_i=e^{\lambda t}(\lambda t)^i/i!`$ and $`F_n=_{i=0}^np_i`$.
Tight bounds for $`Q^{}(n)`$ have been given by Glaz and Naus and a recursive method proposed for calculating $`Q^{}(n;2t)`$ and $`Q^{}(n;3t)`$ for situations where the random quantity $`X_i`$ may take on values other than $`0,1`$, that is situations where an ’event’ cannot be given as either present or absent but only with a non-zero probability.
Other approximations for the tail of the scan statistics and the moments of its distribution have been given by Glaz and Chen and Glaz. Sample tables of the scan statistic have been given for $`n500`$ by Glaz.
This treatment of the scan statistic is for an interval of length $`t`$, specified in advance. When searching for a ’burst’ of events, an *a priori* length cannot always be specified. An extension to the treatment above has been described by Nargawalla in which the length need not be pre-assigned.
#### 3.2.4 Alm Approximation for the Scan Statistic
A new approximation has been given recently by Alm which is accurate and easy to calculate for large values of $`T/t`$ and $`\lambda t`$. This treatment examines the distribution of upcrossings, that is occurrences where the number of events in the scanning window increases by 1 as the window is moved. By separating these events into *primary* and *secondary* upcrossings, the dependence of the second type from the first (almost) independent events allows significant simplifications. If each window of length $`t`$ were independent, the expected number of events would be $`\lambda t`$ with a Poisson probability function $`F_{\lambda t}(n)`$ and distribution function $`p_{\lambda t}(n)`$. The approximation based on the ideas above gives the simple modification:
$$P\left(Nn\right)=1F_{\lambda t}(n)\mathrm{exp}\left[\left(1\frac{\lambda t}{n+1}\right)\lambda \left(Tt\right)p_{\lambda t}(n)\right]$$
(7)
Equation 7 has been tested for $`\lambda t=40`$ and $`T/t=3600`$ using 10000 Monte Carlo simulations. The results are shown in table 2 for $`13n23`$. It can be seen that equation 7 is a good approximation within the sampling errors.
#### 3.2.5 Other Approximations for the Scan Statistic
Other methods have been published, for example the Burst Expectation Search by Giles and CUSUM by VanStekelenborg and Petrakis. The first follows earlier work which used binned times of events and calculates Poisson probabilities of bin counts from a running average of a sample of bins. The BES inverts this process and, for each possible bin count from zero to several hundreds, calculates the mean rate below which the possible count could be a significant burst using a fixed sample of bins around the trial bin. The aim of keeping a fixed sample was to avoid problems arising from a step function edge entering a moving average.
#### 3.2.6 Bursts: Summary
In summary, of possible methods suggested for searching for bursts using classical statistics, the Scan Statistic is recommended, both for time-tagged data and for time-binned data. For small data sets or large window sizes, equation 4 provides an exact probability of the largest number in any window arising due to chance. Many approximate formulae are available, depending on whether the probability of the scan statistic is expected to be large or small. In most practical cases in cosmic rays the statistic is used to search for a possible outburst and so the probability of a given value of the scan statistic arising due to chance will be small in order to be useful. It becomes a matter of computational convenience which of the formulae above is used but equation 7 delivers a good approximation over a wide range of probability values and is easy to calculate. It also has the advantage, in terms of understanding the principles, of starting from the naive initial Poissonian formula with non-overlapping (independent) windows. Its use is therefore recommended here.
### 3.3 Periodicity
Most of the methods for time-series analysis, including trend analysis and auto-regressive moving average (ARMA), have been developed for fields other than cosmic rays, for example . Fourier methods suitable for data at equally spaced times are well developed but are usually not suitable, although these have been extended to discuss unequal intervals and missing data. A bibliography of astronomical time series analysis has been given by Koen.
#### 3.3.1 The Rayleigh test and Dependants
The spur for the introduction of the Rayleigh test into $`\gamma `$-ray astronomy was the unsatisfactory nature of the statistics being used before. Early tests on $`\gamma `$-ray data used epoch-folding to produce a histogram in phase, and $`\chi ^2`$ as a statistic for goodness of fit to a uniform distribution. This suffers from several disadvantages:
1. the freedom to select the number of bins,
2. the freedom to define the starting phase,
3. the failure to use the information contained in the order of the bins.
This last problem can be overcome to some degree by using the Run Test, which is independent and therefore whose probability may be combined with that from $`\chi ^2`$. The result of the freedoms listed above is that different authors could return quite different chance probabilities, given the same data, despite using the same test statistic. The analysis of $`\gamma `$-ray data from the Crab pulsar by Gibson et al. contained the first known use of the Rayleigh statistic in cosmic ray work. It is still a goodness-of-fit statistic, which has no explicit hypothesis as an alternative to the null hypothesis.
The time of each event is treated as a unit vector in the plane, with an angle equal to the pulsar phase. If $`N`$ unit vectors of random orientations (random phases) are added, the distribution of the resultant $`R`$ may be obtained from the distribution of the orthogonal components of the vectors, $`\mathrm{sin}\varphi _i`$ , $`\mathrm{cos}\varphi _i`$ where $`\varphi _i`$ is the phase of the $`i^{th}`$ vector. The means of these components are :
$$C=\frac{1}{N}\underset{i=1}{\overset{N}{}}\mathrm{cos}\varphi _iS=\frac{1}{N}\underset{i=1}{\overset{N}{}}\mathrm{sin}\varphi _i$$
From the Central Limit Theorem (CLT) means of samples of $`C`$ are distributed, for large $`N`$, as a Gaussian with $`var_C=\sigma _C^2/N`$. For vectors uniformly covering the circle:
$$\sigma _C^2=_0^{2\pi }\mathrm{cos}^2(\varphi )𝑑\varphi /2\pi =0.5$$
therefore $`var_C=var_S=1/2N`$. The quantities $`C`$ and $`S`$ are asymptotically uncorrelated and have zero means. The statistic $`2NR^2=2NC^2+2NS^2`$ is therefore the sum of the squares of two zero-mean, unit-variance uncorrelated variables and is distributed as $`\chi ^2`$ with 2 degrees of freedom . The probability distribution function (pdf) of $`R`$ is :
$$f(R)dR=2NRe^{NR^2}dR$$
(8)
and its cumulative probability distribution is :
$$F(R)=e^{NR^2}$$
(9)
The quantity $`NR^2`$ is known as the Rayleigh power.
If a data set spans a time interval $`T`$ the number of independent frequency trials in the frequency range $`f_1`$ to $`f_2`$ is $`\nu =T\left(f_1f_2\right)`$ if $`T>>f_1`$,$`f_2`$, with the independent frequencies separated by $`1/T`$. In practice allowance must be made for leakage: the possible effect of a signal at frequency $`f_0`$ on trial frequencies $`f`$ with $`ff_0>1/T`$, and oversampling: the possibility of obtaining a larger value of $`NR^2`$ by varying the frequency between adjacent independent frequencies. This has been done using by de Jager et al. by Monte Carlo techniques and analytically by Orford. Both methods agree that the number of trials is $`nT\left(f_1f_2\right)`$ where $`n`$ is a slowly varying function of $`F(R)`$ in the range $`2`$ to $`4`$, with a value approximately $`3`$ for $`F(R)10^3`$.
#### 3.3.2 The $`Z_n^2`$ Test
The $`Z_n^2`$ test is the extension of the Rayleigh test to include harmonics. IF n separate harmonics are included with independent coefficients, the statistic is
$$Z_n^2=2N\underset{i=1}{\overset{n}{}}R^2(i\omega )$$
(10)
where $`2NR^2(i\omega )`$ is the Rayleigh power for the $`i^{th}`$ harmonic. $`Z_n^2`$ is distributed as $`\chi ^2(2n)`$. Variations on this technique depend on the method used to select the number and weighting of harmonics. A similar principle is used in radio astronomy where a pulse of width $`W`$ is searched for using $`P/2W`$ harmonics which improves the signal to noise by a factor of up to $`(P/2W)^{0.5}`$. A search for $`\gamma `$-ray emission from radio pulsars proposed the use of $`Z_2^2`$ as a relatively powerful but general test for periodicity. The power of the Rayleigh test for light curves from sinusoids to $`\delta `$-functions was explored by Protheroe . A variant of $`Z_n^2`$ is the H-test in which the value of $`n`$ is obtained objectively from the data and $`Z_n^2`$ is suitably rescaled. This last test is most suitable for multi-mode light curves.
#### 3.3.3 Limitations of the Rayleigh and associated statistics
The foregoing results for the Rayleigh ($`Z_1^2`$) and $`Z_{n>1}^2`$ tests are for the asymptotic case, that is: uncorrelated $`C`$ and $`S`$ with zero means. In most practical applications, these conditions are not strictly met. Ground-based gamma-ray observations of long-period pulsars are limited by:
1. being only a few hours in duration and
2. variations in zenith angle, producing changing counting rates.
The requirement for large sample size is usually met - typical counting rates are $``$ 1 per second over several hours. The result is an enhancement of $`\chi ^2`$ in pure noise data for longer test periods - red noise. The first limitation listed above may be overcome by truncating the dataset so that only integral multiples of the trial period are tested - see Carramiñana et al. and Raubenheimer and Ögelman. As a result, the two trigonometric terms have zero expectations, given a constant mean counting rate. This truncation is easy to accomplish, but results in a variable data selection depending on the test period and therefore all periods are not accorded the same treatment. Since the periodogram is the convolution of the power spectral density with the Fourier transform of the data window, any spectral estimate based on a truncated data set is biased. Further, any correlation introduced by the second limitation above will not be removed this way. An attempt to remove the results of the counting rate variation has been made by Raubenheimer et al. by fitting a parameter $`a`$ in an *ad hoc* modification of the Rayleigh probability distribution:
$$F(R)=e^{2aNR^2}$$
(11)
to random data sets containing no signal, but with the same parameters as the test data set. For data taken on Vela X-1 (period $``$ 5 minutes) they found that equation 11 with $`a=0.4`$ (as opposed to 0.5 from simple theory) gave a probability distribution which was a good fit to the distribution in $`\chi ^2`$ for noise at periods near to 5 minutes in simulated data sets.
#### 3.3.4 Modified Rayleigh Statistic
If the expectations of $`C`$ and $`S`$, their variances and their covariance are not assumed to be zero, $`\frac{1}{2N}`$ and zero respectively, but are calculated for a specific dataset, then the asymptotic probability equation 9 may be valid, given a sufficiently large number of events.
The expression for $`\chi ^2`$ in the case of samples of two correlated variables $`C`$ and $`S`$ is :
$`\chi ^2`$ $`=`$ $`\left[\begin{array}{c}\overline{C}E(C)\\ \overline{S}E(S)\end{array}\right]^T\left[\begin{array}{cc}\sigma _C^2& cov_{S,C}\\ cov_{S,C}& \sigma _S^2\end{array}\right]^1\left[\begin{array}{c}\overline{C}E(C)\\ \overline{S}E(S)\end{array}\right]`$ (18)
For any data set, the substitution of the actual values of $`E(C)`$, $`E(S)`$, $`\sigma _C`$, $`\sigma _S`$ and $`cov_{S,C}`$ will result in a value of $`\chi ^2`$ corrected for the correlation of the variables $`C`$ and $`S`$ and with a probability distribution, for large sample size, given by $`\mathrm{exp}(\chi ^2/2)`$.
In the case of a box-car data set with a constant average counting rate, a starting time $`t_1`$ and ending time $`t_2`$ with $`T=t_2t_1`$ and a trial period $`P=2\pi /\omega `$ :
$`E(C)`$ $`=`$ $`{\displaystyle \frac{\left[\mathrm{sin}\omega t\right]_{t_1}^{t_2}}{\omega T}}`$
$`E(S)`$ $`=`$ $`{\displaystyle \frac{\left[\mathrm{cos}\omega t\right]_{t_1}^{t_2}}{\omega T}}`$
$`Nvar_C`$ $`=`$ $`{\displaystyle \frac{1}{2}}+{\displaystyle \frac{\left[\mathrm{sin}\omega t\mathrm{cos}\omega t\right]_{t_1}^{t_2}}{2\omega T}}[E(C)]^2`$
$`Nvar_S`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\left[\mathrm{sin}\omega t\mathrm{cos}\omega t\right]_{t_1}^{t_2}}{2\omega T}}[E(S)]^2`$
$`Ncov_{S,C}`$ $`=`$ $`{\displaystyle \frac{\left[\mathrm{sin}^2\omega t\right]_{t_1}^{t_2}}{2\omega T}}E(C)E(S)`$
These depend solely on $`\omega `$, $`t_1`$ and $`t_2`$ and their substitution in equation 18 gives a $`\chi ^2`$ value corrected for the finite length of the data set. If it is known that there is no secular change in counting rate the substitution of the above equations into equation 18 would give the correct formal probability of chance occurrence, even if the duration of the data set is less than the trial period, as long as the number of events was high enough for the CLT to be valid. It is more usually the case in ground-based gamma-ray observations that the box-car function is only an approximation. Monte Carlo simulations of data sets have been carried out to test the validity of data set truncation and the above formulations for the case of secular variations of counting rate superimposed on noise. In order to test the validity of the probability distribution equation 9 down to probabilities $`10^6`$, data sets were generated using a multiplicative congruential algorithm with shuffling, chosen to avoid serial correlations. The repeat period is longer than $`2\times 10^{18}`$. The time of each event $`t_i`$ was generated from the previous event: $`t_i=t_{i1}\mathrm{\Delta }(t)\mathrm{ln}(rnd)`$ where $`\mathrm{\Delta }(t)`$ is the mean separation of events as a function of time. A group of $`10^6`$ data sets of duration 8000 s were simulated with a counting rate profile $`R=R_0(10.3t/4000)`$ and $`R_0=1s^1`$. Each data set was tested for periodicity at a trial period of 295s by finding $`C`$ and $`S`$ with reference to the time of the first event. These values were substituted into equation 18 for various assumptions about the form of $`E(C)`$, $`E(S)`$, $`var_C`$, $`var_S`$ and $`cov_{S,C}`$. The probability of chance occurrence was calculated from $`e^{\chi ^2/2}`$.
The resulting cumulative frequency distributions for $`e^{\chi ^2/2}`$ have been calculated for the cases of (a) a truncation of the data set to integral multiple of the trial period, (b) box-car function and (c) a linear fit to the counting rate profile. The ratios of the observed to expected frequencies of occurrence of $`\chi ^2`$ chance probabilities is shown in figure 3, as functions of $`\mathrm{log}(\chi ^2probability)`$. Note that the duration of the data set is corrected for equally well by (b) and (c).
The boxcar and truncated statistics both make corrections for the finite length of a data set, but give a residue which may be identified as being caused by the change of counting rate during the trial period. Longer trial periods or greater rates of change in counting rate would amplify their biases. The truncation method has a distribution which may be fitted, for this simulated data set, by a form such as equation 11 with $`a=0.45`$. The linear fit model is seen to be a good representation of the noise spectrum down to chance probabilities of $`10^5`$.
#### 3.3.5 Other Tests
Leahy et al. pointed out that the unmodified Rayleigh test was powerful for detecting wide peaks in a light curve, in fact it is identical with a likelihood ratio test of a sinewave plus uniform against a uniform phase distribution . In addition, for a light curve of a von Mises form (the circular generalisation of the Gaussian), the Rayleigh statistic exhausts the data’s information on periodicity if the concentration parameter $`\kappa `$ is allowed to vary freely .
Narrow periodic pulse detection, with significant power in the higher harmonics, is bound to be quite difficult because the number of degrees of freedom increases with the trial frequency range. Protheroe proposed a test statistic
$$T_n=\frac{2}{n(n1)}\underset{i=1}{\overset{n1}{}}\underset{j=i+1}{\overset{n}{}}(\mathrm{\Delta }_{ij}+1/n)^1$$
(19)
which looked for close clustering of points on the circle. In this statistic $`\mathrm{\Delta }_{ij}`$ is the distance between the angles $`x_i`$ and $`x_j`$ of two events on the circle:
$$\mathrm{\Delta }_{ij}=0.5\left[(x_ix_j)0.5\right]$$
The null distribution was found using Monte Carlo methods for $`n200`$ and critical values given. The context of the test was the search for ultra-high energy $`\gamma `$-rays from Cygnus X-3 and was therefore not designed for large $`n`$. In this limitation it is similar to the exact expression scan statistic described above. Others have suggested variants which are designed to be powerful for certain classes of pulsed emission .
The Scan Statistic may also be used for searching for non-uniformity in phase. For narrow windows its probability distribution is well approximated by the scan statistic on the line. No systematic work on the use of the Scan Statistic in periodic analysis has been traced.
#### 3.3.6 Searching for a Periodicity
It is usually only the case that a unique periodic ephemeris is available for high energy photons from an isolated radio pulsar. In other cases, a search must be made in period, and the test used must allow for the freedoms associated with the trial period range. A rule-of-thumb arising from the number of ’degrees of freedom’ implicit in a periodicity search using the Rayleigh test is that the search should be at intervals in period of $`(IFI)/3`$ ($`IFI`$ = Independent Fourier Interval). For a period of $`P`$ in a data set of duration $`T`$ this corresponds to a trial period step of $`P^2/3T`$. This step size has the advantage that the number of degrees of freedom to be used to interpret the peak periodic amplitude found is approximately the number of periods tried. If harmonics of the test period are to be included, the spacing would be correspondingly reduced and hence the number of trial periods increased. For the $`Z_n^2`$ test, the reduction in period step, and consequent increase in both computation and the degrees of freedom to be accounted for, is by a factor $`n`$.
When searching for pulsed emission from some sources, in particular binary sources, there is frequently poor knowledge of the both the pulsar period and period derivative. In this case the light curve will be narrow only if the correct period $`P`$ and period derivative $`\dot{P}`$ is offered to the test. A nearby, but not correct, trial period and the ignorance of a period derivative will smear the light curve. If a true light curve were a $`\delta `$ function at period $`P=P_o`$ and the trial period was $`P=P_o+\mathrm{\Delta }`$, the light curve would be a rectangular distribution in phase of length $`T\mathrm{\Delta }/P_o^2`$. This effectively limits the number of harmonics which may be realistically added to $`P_o^2/T\mathrm{\Delta }1`$.
Some searches for periodicity are combined with a search for a DC excess. This is common in Čerenkov telescope searches where ON-source data is compared with OFF-source control data to detect any DC component. The combined analysis of this situation was proposed by Lewis in which a statistic $`\alpha `$ is defined as the sum of the Rayleigh statistic and the square of equation 1, distributed as $`\chi ^2`$(3). The assumption in this case is that all of the excess is pulsed; if there is an unpulsed component the test statistic will be biased. Again, the presence of a possible unpulsed component could be built into a Bayesian analysis.
#### 3.3.7 Conclusion on Periodicity
The question of the best classical test for the presence of periodicity is a complex one. The selection of the most sensitive test requires a knowledge of
* the pulsar ephemeris
* the light curve shape
* the background noise distribution
If all of these are known in advance, a most powerful test, based on the $`Z_n^2`$ extension of the Rayleigh test is likely to be close to optimum. Frequently some or all of these will be unknown or poorly known. In this case, some allowance must be made for the lack of knowledge and the test selected should not contain any assumption which causes a significant bias. It has sometimes been claimed that the Rayleigh test is ’biased’ towards broad light curves and that a test which is more sensitive to narrow light curves should be used when such a light curve is suspected. This raises the problem, discussed in the previous section, of the smearing of a light curve if the pulsar’s ephemeris is uncertain.
Protheroe has suggested that if one has no information about the nature of the phase distribution one should be conservative and adopt the Rayleigh test. A rational for this is that if one is searching for an unknown period and an unknown light curve, which is quite common in $`\gamma `$-ray work, and there is no significant power in the fundamental, then a test involving the addition of an unknown number of higher harmonics is unlikely to be successful. This point will be revisited later in discussing a Bayesian method of searching for periodicity.
A simple suggestion made before and reiterated here is that if $`Z_n^2`$ does not show evidence for periodicity, that is: there is no significant power in the fundamental or the first harmonic, either in addition or separately, then it is unlikely that the data will contain a strong periodic signal.
## 4 Spatial Analysis
### 4.1 Introduction
Spatial analysis of arrival direction data is of great interest for X- and $`\gamma `$-rays from satellites and for cosmic rays of the highest energies, which may not be greatly deflected in the galactic magnetic field.
Simple methods rely on a grid placed on the events and counts in the grid cells taken as independent Poisson-distributed events. If the cells are fixed absolutely, there is no problem in ascribing a suitable Poissonian probability to the largest number detected in any cell. If there is freedom to incrementally move the cell containing the largest count, a larger number is generally found. In this case the new cells created are correlated and the assumption of independence is incorrect: simple application of Poissonian probabilities is inappropriate. The problem of having the freedom to move the boundaries of the cells was pointed out for cosmic ray ’sources’ by Hillas who suggested a conservative number of ’sigmas’. Large scale anisotropy in gamma ray bursts were sought using dipole and quadrapole analysis. A ’pair matching’ statistic was used by Bennett and Rhie to check for gamma ray burst repeaters rather than ’nearest neighbour’ methods used by others and criticised by Nowak. Many methods have been used which are based upon a known point-spread function (PSF). Amongst these are Maximum Entropy methods such as those used for satellite X-ray imaging, maximum likelihood and Hough Transforms.
In the next section it is suggested that the scan statistic is a powerful and general statistic for which good approximations exist for the chance probabilities. It has recently been extended to two dimensions by Loader, Chen and Glaz and Alm. Kulldorf has extended this further to higher dimensional searches.
### 4.2 2-Dimensional Scan Statistic
This is the two-dimensional development of the Scan Statistic introduced above. It will be introduced using a notion of ’elemental’ cells from which a two-dimensional scanning window is constructed. In effect, the scanning window may be moved by discrete steps of the size of the elemental cell. Assume that a two-dimensional square region $`R=[0,L]\times [0,L]`$ of side $`L`$ is inspected for occurrences of ’sources’. The region is partitioned into $`n\times n`$ elemental cells so that the size of a cell $`h=L/n`$. The contents of each of the $`n^2`$ cells are independent.
For $`1in`$ and $`1jn`$, define a random variable $`Y_{i,j}`$ as the number of events in the elemental cell $`[(i1)h,ih]\times [(j1)h,jh]`$. A square box of $`m\times m`$ small cells is scanned over the whole of region $`R`$. There will be $`\nu `$ such boxes, partially dependent if $`m2`$, with $`\nu =0,(nm+1)^2`$.
Define
$$S(i_1,i_2)=\underset{j=i_2}{\overset{i_2+m1}{}}\underset{i=i_1}{\overset{i_1+m1}{}}Y_{i,j}$$
to be the number of events in the square box of $`m^2`$ adjacent cells starting at $`i=i_1`$, $`j=i_2`$. If, during the scanning of the $`m`$box, $`S(i_1,i_2)`$ exceeds a particular value $`k`$, a ’source’ has been detected. For $`1i_1,i_2nm+1`$ define an ’event’ $`A_{i_1,i_2}`$ as an occurrence of $`S(i_1,i{}_{2}{}^{})k`$ and as a member of the set $`A`$ of all such occurrences.
The two-dimensional scan statistic is defined as:
$$S_m=\mathrm{max}\{S(i_1,i_2);1i_1nm+1,1i_2nm+1\}$$
and the probability that $`S_m`$ has at least a value $`k`$ is:
$`P(S_mk)`$ $`=`$ $`P\left({\displaystyle \underset{i_1=1}{\overset{nm+1}{}}}{\displaystyle \underset{i_2=1}{\overset{nm+1}{}}}A_{i_1,i_2}\right)`$
$`=`$ $`P\left({\displaystyle \underset{i_1=1}{\overset{nm+1}{}}}{\displaystyle \underset{i_2=1}{\overset{nm+1}{}}}A_{i_1,i_2}^c\right)`$
where $`A_{i_1,i_2}^c`$ is the occurrence of $`S(i_1,i_2)<k`$.
#### 4.2.1 Glaz Approximation to 2-D Scan Statistic
For a fixed value of $`1i_1nm+1`$ the one-dimensional approximation holds:
$$P\left(\underset{i_2=1}{\overset{nm+1}{}}A_{i_1i_2}^c\right)q_{2m}\left(\frac{q_{2m}}{q_{2m1}}\right)^{n2m}$$
and since $`nm+1`$ square regions of $`m\times m`$ are scanned, a reasonable approximation is:
$$P(S_mk)1q_{2m1}\left(\frac{q_{2m}}{q_{2m1}}\right)^{(n2m+1)(nm+1)}$$
(20)
For a Poissonian distribution of events the following expression was found to be a good approximation:
$$P(S_mk)1\mathrm{exp}(\lambda ^{})$$
(21)
where the approximate mean for the asymptotic Poisson distribution is
$$\lambda ^{}=1q_{2m2}+(n2m+2)(nm+1)(q_{2m2}q_{2m1})$$
and
$$q_{m+l1}=P(A_{1,1}^cA_{1,1}^c\mathrm{}A_{1,l}^c)$$
Tables are given of this and other approximations, for the Poisson model of $`m20`$ and $`n500`$.
#### 4.2.2 Alm Approximation to 2-D Scan Statistic
A recent paper has given an approximation based on a modification of the method of counting *upcrossings* used in equation 7, which is easy to calculate and moreover is given for a more generally useful rectangular scanning window $`[0,a]\times [0,b]`$ in a rectangular region $`[0,S]\times [0,T]`$. The scan statistic $`L`$ is the maximum content of a scanning window with a two-dimensional Poissonian process $`X`$ with event density $`\lambda `$:
$$L=L(\lambda ,a,b,S,T)=\underset{atT,bsS}{\mathrm{max}}X([ta,t]\times [sb,s]$$
The probability of observing at least $`n`$ events in a scanning window is:
$$P(Ln)1F_{N(a)}e^{\gamma _{n+1}}$$
(22)
where
$$\gamma _{n+1}\left(1\frac{\lambda ab}{n+1}\right)\left(Ta\right)b\lambda \left(\mu _n\mu _{n+1}\right)e^{\mu _{n+1}}$$
$$\mu _n\left(1\frac{\lambda ab}{n}\right)\lambda a\left(Sb\right)p_{\lambda ab}(n1)$$
and
$$F_{N(a)}F_{\lambda ab}\mathrm{exp}\left[1\left(1\frac{\lambda ab}{n+1}\right)\lambda a\left(Sb\right)p_{\lambda ab}(n)\right]$$
$`p_\mu `$ and $`F_\mu `$ are the Poisson probability and cumulative probability distributions.
The predictions of equation 22 have been compared with the results of $`10^5`$ Monte Carlo simulations in table 3 for $`S=T=20`$, $`a=b=2`$ and $`N=800`$. The agreement is good, allowing for the errors inherent in the Monte Carlo results. For interest, the final column shows the Poissonian probability obtained if the cells were treated as independent. In this particular case, the result of assuming independence of the cells would be a fairly consistent overestimate of the significance of the ’source’ by about ’$`3\sigma `$’. The precise amount of underestimate of the chance probability will depend on the number of elemental cells in the scanning window. Finally, the treatment of has been extended to other shapes of scanning window, such as circular.
#### 4.2.3 Summary
In summary, the 2-D Scan Statistic is a preferred general statistic for those cases where events are located randomly on a plane, within fixed bounds, and where there is no *a priori* expectation such as a known source with known instrumental spread function. In most practical situations a good approximation is obtained by using equation 22.
## 5 Bayesian Methods
### 5.1 Introduction
For many workers in cosmic rays, Bayesian methods are relatively novel and the following section attempts to summarise the main ideas and methods. A much fuller development of the ideas discussed below is given by Loredo.
#### 5.1.1 Statistics
The term ’statistics’ arises from the concept of a *statistic*. A statistic is a number derived from observed data and which obeys certain rules, some of which depend on a hypothesis about the system under observation, some of which are extraneous. From this number, one can say how likely it is that the data was drawn from a population obeying rules specified by the particular hypothesis, assuming that all extraneous quantities are allowed for. From that, by an inversion of logic, it is inferred how likely is the hypothesis. In many cases, a particular statistic is used because the experimental results appear to be presented, or may be rearranged to be presented, in a form which allows an easy calculation of that statistic. An example is the epoch-folding of time-tagged photon times above, followed by $`\chi ^2`$ calculated from the binned phases. As was pointed out in that example, some arbitrary choices had to be made which rendered the results unsatisfactory. Also, the aspects of the experimental data which were used to calculate a statistic may not be all that are available, or the most discriminating aspects. This should, with careful design, be evident from a consideration of the statistic’s ’power’ but not necessarily. It is the claim of Bayesians that such problems are inherent in ’classical’ statistics and derive from a misunderstanding of the meaning of Probability.
#### 5.1.2 The Meaning of Probability
There were at the beginnings of the subject, and still are, two schools of thought. The first school maintains that the term *probability* is a statement of the frequency of occurrence of data, such as that taken in a very large number of repeats of an experiment, under the assumption that random factors are at work causing the possibility that the results could be different every time. Take, as an example, coin-tossing: the probability of heads is obviously 0.5 in a single toss. Everyone would agree that, assuming no trickery, an unbiased coin would land equally likely as ’heads’ or ’tails’. But there are forces at work which affect the way a coin would land - all amenable to analysis. In fact a coin-tossing machine could be made which obtained ’heads’ or ’tails’ every time. We regard coin-tossing as a random activity only because we expect humans to apply unconscious variability to the force and direction of the flip which is much greater than that needed in the initial conditions to obtain one more extra turns before landing. This illustrates an important point: unless the hypothesis is clear and specifies all pertinent factors there could be an apparent randomness. That is not to say that true randomness does not exist, only that it is often used as an alibi for lack of knowledge or precision in stating the experimental conditions.
An alternative definition of *probability* is ’a measure of belief in a certain hypothesis’. This is, and was, a much easier idea to grasp but one which was felt from an early date not to be capable of exact or scientific analysis. One consequence of this idea is that for a unique set of data, perhaps taken on a naturally-occurring phenomenon, the idea of a very large set of repeated experiments to plot out the ’frequency distribution necessary to use a ’statistic’ was unrealistic. It is this definition which underlies Bayesian thought, and indeed is the definition which more closely accords with the questions for which measurements are made. Interestingly, this meaning of probability explains the frequency version as a special case using de Finetti’s representation theorem for exchangeable sequences of events.
The main difference between the two philosophical approaches is how the data are related to the hypotheses.
1. The ’Frequentist’ approach: We obtain $`P(DH)`$, the conditional probability of obtaining the observed data, given a particular hypothesis. The hypothesis is frequently a model $`M`$ which has a parameter space $`\theta `$ and $`P(DM,\theta )`$ is the ’sampling distribution’ for the data, given the model. A frequently met hypothesis is the ’null’ hypothesis $`H_o`$ in which the parameters are set to zero. For any hypothesis, a statistic is formed which is ’locally most powerful’ or even better ’uniformly most powerful’ and the probability of observing the data is assigned from a knowledge of the statistic’s distribution function. This is now used to give a range of values in which the value of a statistic may fall by chance, with given probability (i.e. frequency). As an interesting aside, it is most usually the case for continuous measures, and frequently for discrete measures, for a range of values of the statistic, including that observed (but also including many values *not* observed), to be used to derive the probability. This is frequently performed by integrating $`P(DM,\theta )`$ over the sample (data) space. That is to say, the probability of a hypothesis is determined by the data taken, plus a whole range of values of data which were *not* observed. This curious situation is not often questioned by ordinary users of ’statistics’.
2. The ’Bayesian’ approach: We obtain $`P(HD)`$, the probability of a hypothesis, given the data - apparently a more difficult matter. However Bayes Theorem gives:
$`P(HD)P(D)`$ $`=`$ $`P(DH)P(H)\text{leading to}`$
$`P(HD)`$ $`=`$ $`P(DH)P(H)/P(D)`$
$`P(DH)`$ is the likelihood function, $`P(D)`$ is the global likelihood, usually treated as an ignorable normalising constant, $`P(H)`$ is the ’prior’ probability of the hypothesis. In addition to the extra terms not used in frequentist analysis, a crucial difference between the approaches is that Bayesian methods would integrate the likelihood function over the parameters space, rather than the sampling (data) space. Bayesian methods have been criticised for the inclusion of an apparently subjective quantity $`P(H)`$ but a trivial example demonstrates that frequentist analyses are not free from this. Frequentists would determine if the hypothesis $`H_o`$ of a histogram having all the cells identical were true by taking $`\chi ^2`$ as a statistic. They would use no prior information or knowledge. But we know that histogram cells cannot contain negative numbers, and so some relevant background information is ignored when using $`\chi ^2`$.
#### 5.1.3 An Example
An example of the different approaches is an experiment in which a coin is tossed $`N`$ times. It lands heads $`H`$ times. The question is: is it biased? In the Frequentist approach a hypothesis (the ’null’ hypothesis) is formed that the coin is unbiased and that the result is a function of randomness only. A sufficiently low probability which is obtained for a suitable statistic would be evidence that the ’null hypothesis’ should be abandoned. The binomial distribution describes the result of such discrete, bounded experiments. The probability of $`H`$ heads in $`N`$ tosses is $`0.5^H\times 0.5^{NH}\times C_H^N`$. One then calculates the probability of obtaining $`0,1,2,\mathrm{}H1`$ heads and add them to the probability of $`H`$ heads and say: ’if the null hypothesis is true, getting $`H`$ *or fewer* heads in $`N`$ tosses can occur due to chance, with a given probability. This could be interpreted as some evidence against the ’null hypothesis’, hence evidence that the coin is biased.
#### 5.1.4 Stopping Rules
Setting aside for the moment the fact that we did not see $`0,1,2,\mathrm{}H1`$ heads, the last conclusion supposes that the coin was tossed, irrespective of the result, $`N`$ times and that the number of heads $`H`$ was the random variable. But suppose the coin was actually tossed by a person until $`H`$ heads were obtained, and that happened to occur after $`N`$ tosses. In this case the number of tosses $`N`$ is the random variable and the number of heads $`H`$ is fixed. The probability is then derived by combining the individual probabilities of obtaining $`H`$ heads from $`hH`$ tosses, and may be significantly different from the first probability calculated.
Loredo gives a more apposite example: a theorist predicts that $`f=10\%`$ of the stars in a cluster should be of type A. An observer reports 5 stars of type A out of 96 observed. The theorist calculates as follows: $`N=96`$ and $`f=0.1`$ gives $`9.6`$ predicted type A stars. The probability of the value of $`\chi ^2=2.45`$ of this information is $`P=0.073`$, which is acceptable at the $`95\%`$ level. The observer however decided in advance to stop when he had found 5 stars of type A. The expected value of $`N`$ is then $`5/0.1=50`$ with variance $`5(1f)/f^2=450`$. The probability of the value of $`\chi ^2=4.7`$ of this information is $`P=0.032`$, which is *not* acceptable at the $`95\%`$ level. This ambiguity arises because of the stopping rule used by the experimenter that is - what data sets *might* have been observed.
The stopping rule can therefore be important in classical statistical analysis, and ignorance of the actual rule used may lead to an erroneous or at least ambiguous conclusion. Knowledge of the exact stopping rule is less important in Bayesian analysis, but is valuable in particular when it contains useful information about the unknown quantities. In other cases, the stopping rule could be important if the existence of some data is unknown to the analyser, perhaps because its analysis did not show it to be significant and it was suppressed by the experimenter. The message from this example is that for Frequentist analysis to be possible, an experiment must be precisely defined and if the execution is different in any way from the plan, the data could be worthless.
#### 5.1.5 Conclusion
In summary, frequentist methods establish $`P(D\theta M)`$ as the sampling distribution of the data, given a model $`M`$ with parameters $`\theta `$ and perform integration over the data space. Bayesian methods start with the same function $`P(D\theta M)`$ but treat it is a likelihood with integrations performed over the parameter space of the model. In particular, parameters which are necessary for the specification of the model but are not of interest (for example the phase when looking for a periodic signal) are integrated out, or marginalised.
In Bayesian theory, the notion of a ’random variable’ is absent so ambiguity does not arise for many types of stopping rule and there is no need for a ’reference set’ of hypothetical data. This state of affairs results from the need in Bayesian methods to be specific about all the hypotheses, or to integrate away any unspecifiable variable. Taking again the example of a histogram, and the question of whether its cells are consistent with uniformity using $`\chi ^2`$İf the null hypothesis is the only hypothesis available, the use of $`\chi ^2`$ is as a ’goodness-of-fit’ test for the supposition of flatness. The number observed in the $`i^{th}`$ bin of $`n`$ bins is $`x_i`$. The number expected in each bin, under the null hypothesis, is $`avg=_{i=1}^nx_i/n`$ and $`\chi ^2=_{i=1}^n(x_iavg)^2/avg`$, assuming $`avg`$ is large enough (usually 10 or so) for asymptotic normality. The probability of $`\chi ^2`$ for $`n1`$ degrees of freedom is interpreted as supporting or otherwise the null hypothesis. This statistic suffers from a major problem in that it ignores information - the order of the bins may be significant, and so it implicitly assumes a class of alternative models in which the order is unimportant. This can be partially rectified by applying an independent test which is only determined by the order of the bins - the Run Test. This is only applying a patch, since the Run test is most powerful against monotonicity and not other patterns. Frequentists acknowledge this problem in general by using the idea of the power of a statistic, that is its ability to identify correctly a true model from a particular alternative. Both approaches have subjective factors: Bayesian in assigning prior probabilities to hypotheses, Frequentist in the notion of randomness and its applicability in a mathematical sense to cover for a lack of knowledge of the exact experimental conditions. A consequence of this is that different experts in both fields may come to different conclusions given the same data. Another way of putting this is that the result of analysing data will be a conclusion within a range, depending on (a) the Bayesian priors, or (b) the estimate of the degrees of freedom and unknown factors.
### 5.2 Bayesian On/Off Analysis
The Bayesian ideas in the above section have recently been applied to the ON/OFF problem treated earlier. As in all Bayesian analyses, some judgement must be made of the priors to be used, but in the cases discussed here the results do not depend critically on how these priors are chosen.
An initial Bayesian analysis of the problem of detecting a source in an ON/OFF counting experiment has been given by Loredo. Using the same notation as in the ON/OFF section above, the probability of the background rate $`b`$ (the posterior density) from the OFF-source data is:
$$p\left(bN_{OFF}\right)=p\left(N_{OFF}b\right)\frac{p\left(b\right)}{p\left(N_{OFF}\right)}$$
The Poisson likelihood for $`N_{OFF}`$ is:
$$p(N_{OFF}b)=\frac{(bT)^{N_{OFF}}e^{bT}}{N_{OFF}!}$$
The parameter $`b`$ is unknown and so the ’prior’ probability would appear to be a matter of guesswork. If the range of $`b`$ were pre-specified in some non-arbitrary way, at least the scale of $`b`$ would be known, and a flat prior would be reasonable. If even the scale of $`b`$ is unknown, the ’least informative’ prior for $`b`$ is $`p(b)=1/b`$, which is uniform in $`\mathrm{log}b`$, and then
$$p(N_{OFF})=\frac{T^{N_{OFF}}}{N_{OFF}!}_0^{\mathrm{}}b^{N_{OFF}1}e^{bT}𝑑b$$
This leads to
$$p\left(bN_{OFF}\right)=\frac{T_{OFF}\left(bT_{OFF}\right)^{N_{OFF}1}e^{bT_{OFF}}}{\left(N_{OFF}1\right)!}$$
Note that the expectation of the background $`\widehat{b}=N_{OFF}/T`$ and that the assumption of $`p(b)=1/b`$ does not strongly affect the result, Loredo pointing out that a prior uniform in $`b`$ only marginally alters the expectation $`\widehat{b}=(N_{OFF}+1)/T`$. The joint probability of the background rate $`b`$ and a source rate $`s`$, given $`N_{ON}`$ and $`N_{OFF}`$, is:
$$p\left(sbN_{ON}\right)=p\left(sb\right)p\left(b\right)\frac{p\left(N_{ON}sb\right)}{p\left(N_{ON}\right)}$$
The probability of the source rate $`s`$ is obtained by marginalising $`b`$, that is $`p\left(sN_{ON}\right)=p\left(sbN_{ON}\right)𝑑b`$:
$$p\left(sN_{ON}\right)=\underset{i=1}{\overset{N_{ON}}{}}C_i\frac{T_{ON}\left(sT_{ON}\right)^{i1}e^{sT_{ON}}}{\left(i1\right)!}$$
(23)
where
$$C_i=\frac{\left(1+\frac{1}{\alpha }\right)^i\frac{\left(N_{ON}+N_{OFF}i1\right)!}{\left(N_{ON}i\right)!}}{_{j=1}^{N_{ON}}\left(1+\frac{1}{\alpha }\right)^j\frac{\left(N_{ON}+N_{OFF}j1\right)!}{\left(N_{ON}j\right)!}}$$
This result is formally correct for all positive values of $`N_{ON}`$ and $`N_{OFF}`$ and is particularly useful for small values when the asymptotic treatments fail.
Its main value is to illustrate the completely different approach and result of the application of Bayesian ideas. However, there are some computational problems for values of $`N_{ON}`$ and $`N_{OFF}`$ which exceed $`100`$. For values of $`N_{ON}`$ and $`N_{OFF}`$ which are less than $`100`$ evaluation of equation 2 and equation 23 shows small differences in the derived probabilities.
### 5.3 Bayesian Change Point Analysis - Bursts
A recent paper by Scargle has used Bayesian methods to analyse structure in photon counting data. It is worth noting that the ON/OFF problem dealt with above is a special case of change point analysis, where there is only one change point and its location is known in advance. The principles are the same as those outlined above, with the added simplicity of having simpler alternatives to the uniform model. The uniform counting rate model $`M_1`$ assumes a constant intensity over a particular time interval $`T`$. An alternative model $`M_2`$ has the interval $`T`$ broken into two regions $`T_1`$ and $`T_2`$, $`T=T_1+T_2`$, each with a different counting rate. In general, a model $`M_k`$ may be constructed with $`k`$ regions. Bayes Theorem give the probability of a model
$$p\left(M_kD,I\right)=\frac{p\left(DM_k,I\right)p\left(M_kI\right)}{p\left(DI\right)}$$
Dropping the explicit appearance of the background information $`I`$, the odds ratio $`O_{kj}`$ between two competing models $`M_k`$ and $`M_j`$ is then
$$\frac{p(M_kD)}{p(M_jD)}=\frac{p(DM_k)p(M_k)}{p(DM_j)p(M_j)}$$
The parameter $`\theta `$ or vector of parameters $`\stackrel{}{\theta }`$ of the model $`M_k`$ enter when $`p(DM_k)`$ is calculated
$$p(DM_k)=p(D\stackrel{}{\theta },M_k)p(\stackrel{}{\theta }M_k)𝑑\stackrel{}{\theta }$$
The odds ratio $`O_{kj}`$ is then
$`O_{kj}`$ $`=`$ $`{\displaystyle \frac{p(M_kD)}{p(M_jD)}}`$ (24)
$`=`$ $`{\displaystyle \frac{p(M_k)}{p(Mj)}}{\displaystyle \frac{p(D\stackrel{}{\theta },M_k)p(\stackrel{}{\theta }M_k)𝑑\stackrel{}{\theta }}{p(D\stackrel{}{\theta },M_j)p(\stackrel{}{\theta }M_j)𝑑\stackrel{}{\theta }}}`$
$`=`$ $`{\displaystyle \frac{p(M_k)}{p(Mj)}}{\displaystyle \frac{(M_k,D)}{(M_j,D)}}`$
where $`(M_k,D)`$ is the global likelihood of model $`M_k`$.
For the constant-rate model $`M_1`$, $`N`$ events arrive in a time $`T`$ which is treated as being divided into $`M`$ intervals of duration $`\delta t`$, the justification being that photon counting apparatus always has a resolving time. Note that the number of events in any particular interval $`\delta t`$ can be $`0`$ or $`1`$ only. The author shows that the global likelihood for this constant-rate model of Time Tagged Events (TTE) is
$$(M_1TTE)=\frac{\mathrm{\Gamma }(N+1)\mathrm{\Gamma }(MN+1)}{\mathrm{\Gamma }(M+2)}$$
If the data is time-binned into $`M`$ equal bins, but such that any number of events may occur in any bin, given an overall rate $`\lambda =N/T`$ and mean number per bin of $`\mu =\lambda T/M`$, the global likelihood is:
$$(M_1Binned)=\frac{\mathrm{\Gamma }(N+1)}{(M+1)^{N+1}}$$
Note that the bins are fixed and may not be scanned to maximise $``$.
The alternative model $`M_k`$ has a likelihood which is the product of the likelihoods of the individual constant-rate regions of $`T`$. For a two-rate model with the time of the change of rate being $`t_{cp}`$
$`(M_2D)`$ $`=`$ $`{\displaystyle 𝑑t_{cp}𝑑\mathrm{\Lambda }_1𝑑\mathrm{\Lambda }_2p_{cp}(t_{cp})\times p[D_1M_1(\mathrm{\Lambda }_1,T_1)]p(\mathrm{\Lambda }_1)}`$
$`\times p[D_2M_2(\mathrm{\Lambda }_2,T_2)]p(\mathrm{\Lambda }_2)`$
where $`\mathrm{\Lambda }=\lambda \delta t`$, $`P(\mathrm{\Lambda })`$ is the prior for the rate $`\mathrm{\Lambda }`$ and $`P_{cp}`$ is the prior for the change-point time $`t_{cp}`$. For time-tagged data with resolution $`\delta t`$ the integrals are sums and the change-point location is $`m_{cp}\delta t`$. Since the change-point can be tested only at the arrival time of a photon, the photon number of the change-point $`n_{cp}`$ is used as an index. The number of events in the first section, up to the change-point, is $`N_1=n_{cp}`$, $`N_2=NN_1`$ and $`M_1=m_{n_{cp}}`$ The global likelihood is then
$`(M_2D)`$ $`=`$ $`{\displaystyle \underset{n_{cp}}{}}{\displaystyle \frac{\mathrm{\Gamma }(n_{cp}+1)\mathrm{\Gamma }(m_{n_{cp}}n_{cp}+1)}{\mathrm{\Gamma }(n_{cp}+2)}}`$
$`\times {\displaystyle \frac{\mathrm{\Gamma }(Nn_{cp}+1)\mathrm{\Gamma }(m_{Nn_{cp}}(Nn_{cp})+1)}{\mathrm{\Gamma }(Nn_{cp}+2)}}\mathrm{\Delta }t_{n_{cp}}`$
The paper gives a coding in a popular mathematical package to implement the above ideas.
### 5.4 Bayesian Periodicity Analysis
#### 5.4.1 Introduction
Frequentist statistical theory allows more than one test to be applied to any situation. Any statistic, or function of the data, may be defined and the ’best’ is selected depending on its ’power’ or likelihood of selecting the ’correct’ hypothesis. One of the problems of the frequentist approach to looking for evidence of periodicity is that, in the absence of a specific light curve, the alternative hypothesis (to one of uniformity in the phase distribution) is unknown and the power of a statistical test is difficult to specify except for a narrow class of alternative light curves. The Rayleigh statistic, $`Z_1^2`$, is powerful only for the fundamental period and is formally the most powerful test for alternatives to uniformity from the Von Mises distribution - the circular equivalent of the Gaussian on the line. The $`Z_n^2`$ test allows the addition of $`n1`$ harmonics but needs a protocol to decide when to stop adding harmonics and therefore degrees of freedom (the $`H`$ test mentioned above suggests such a protocol). Finally, Protheroe’s test is powerful for very narrow light curves. Each could be tried in succession to look for evidence of periodicity, but a method which is indifferent to the shape of the light curve, without any penalty, would be of great advantage.
#### 5.4.2 Gregory & Loredo Method
Such a method based on Bayesian analysis, is claimed by Gregory and Loredo. The essence of the method is to compare a uniform model for the distribution in phase at a trial frequency with a periodic model. The great difference between this and other methods is how the periodic model is proposed and how the necessary uncertainties and their associated ’degrees of freedom’ of classical theory are accounted for. In particular, since an arbitrary postulated light curve may be of any shape, the method automatically applies Ockham’s razor, in that models with fewer variables are automatically favoured unless the evidence from the data more than compensates. More complicated light curves (not necessarily with small number of harmonics, a $`\delta `$-function is uncomplicated in this context) are penalized for their greater complexity.
Bayes Theorem is used to compare the probabilities of two parameterised models of the phase distribution. In the notations of the authors, the probability that a model $`M`$ describes the data, given the data $`D`$ and any background information $`I`$ is
$$p\left(MD,I\right)=p\left(MI\right)\frac{p\left(DM,I\right)}{p\left(DI\right)}$$
(25)
The first term on the right, $`p\left(MI\right)`$, is the prior probability of the model $`M`$, which may seem to be subjective but may be estimated in some cases from the permissible range of the parameters. The numerator in the second term, $`p\left(DM,I\right)`$, is the sampling probability of the data $`D`$, or the likelihood of the model $`M`$. The denominator, $`p\left(DI\right)`$, is the global likelihood of the entire class of models. If the model contains a parameter $`\theta `$, or in the case two or more parameters a vector $`\stackrel{}{\theta }`$, the likelihood of the model can be calculated:
$$p\left(DM\right)=_\stackrel{}{\theta }p\left(D\stackrel{}{\theta },M\right)p\left(\stackrel{}{\theta }M\right)$$
For time-tagged photon data with $`N`$ events detected over a time $`T`$, the probability of $`D`$ for a particular rate model $`r(t)`$ can be calculated. For the time $`T`$ divided into very small intervals of length $`\mathrm{\Delta }t`$, the probability of $`n`$ events in $`\mathrm{\Delta }t`$ is:
$$p_n=\frac{\left[r(t)\mathrm{\Delta }t\right]^ne^{r(t)\mathrm{\Delta }t}}{n!}$$
If $`\mathrm{\Delta }t`$ is small enough for $`p_i=0,i2`$ then the sequence of $`T/\mathrm{\Delta }t`$ time samples will contain $`N`$ containing one event and $`Q=T/\mathrm{\Delta }tN`$ containing no event. The likelihood is then:
$$p\left(Dr,I\right)=\underset{i=1}{\overset{N}{}}p_1(t_i)\underset{k=1}{\overset{Q}{}}p_0(t_k)$$
Using $`p_0(t)=e^{r(t)\mathrm{\Delta }t}`$ and $`p_1(t)=r(t)\mathrm{\Delta }te^{r(t)\mathrm{\Delta }t}`$ the likelihood function is
$$p\left(Dr,I\right)=\mathrm{\Delta }t^N\left[\underset{i=1}{\overset{N}{}}r(t_i)\right]\mathrm{exp}\left[\underset{k=1}{\overset{N+Q}{}}r(t_k)\mathrm{\Delta }t\right]$$
In the case of a periodic model, the non-uniformity in phase is characterised by the varying contents of the phase bins. Although the number of phase bins needed to detect any light curve and the origin of phase are unknowns, these will be marginalised or integrated out. If there are $`m`$ phase bins the average rate $`A=\frac{1}{m}_{j=1}^mr_j`$ and the fraction of the total rate per period in phase bin $`j`$ is $`f_j=\frac{r_j}{mA}`$. The likelihood function is shown to reduce to
$$P\left(D\omega ,\varphi ,A,𝐟,M_m\right)=\mathrm{\Delta }t^N(mA)^Ne^{AT}\left(\underset{j=1}{\overset{m}{}}f_j^{n_j}\right)$$
where $`\omega `$ is the postulated angular frequency, $`\varphi `$ the starting phase, $`𝐟`$ the set of $`m`$ values of $`f_j`$ and $`n_j`$ being the number of events occurring in bin $`j`$.
The joint prior density for the parameters $`\omega ,\varphi ,A,𝐟`$ is
$$p(\omega ,\varphi ,A,bffM_m)=p\left(\omega M_m\right)p\left(\varphi M_m\right)p\left(AM_m\right)p\left(𝐟M_m\right)$$
The prior densities are:
1. $`p(\varphi M_m)=1/2\pi `$, this assumes that any starting phase is equally likely,
2. $`p(AM_m)=1/A_{max}`$, this assumes that $`A`$ does not change during the observation and any value of $`A`$ from $`A=0`$ to $`A=A_{max}`$ is possible,
3. $`p(\omega M_m)=\omega \mathrm{ln}(\omega _{hi}/\omega _{lo})`$, where $`[\omega _{hi},\omega _{lo}]`$ is a prior range for $`\omega `$,
4. $`p(𝐟)=(m1)!\delta \left(1_{j=1}^mf_j\right)`$.
The assignment of the priors of the models themselves is all that is needed before comparing the likelihoods of the models. The two models are equally likely *a priori* and so the prior likelihood of the non-periodic model ($`M_1`$), $`p(M_1)=1/2`$ and that for the periodic model ($`M_m`$, $`m=2,m_{max}`$), $`p(M_mI)=1/2\nu `$ where $`\nu =m_{max}1`$.
The final result for the odds $`O`$ against a uniform model of phase and in favour of a periodic model with phase and period unknown (a common case) is:
$$O=\frac{1}{2\pi \nu \mathrm{ln}\left(\omega _{hi}/\omega _{lo}\right)}\frac{N!(m1)!}{(N+m1)!}_{\omega _{lo}}^{\omega _{hi}}\frac{d\omega }{\omega }_0^{2\pi }𝑑\varphi \frac{m^N}{W_m(\omega ,\varphi )}$$
(26)
where $`W_m(\omega ,\varphi )`$ is the number of ways that the set of $`n_j`$ observed counts can be made by distributing $`N`$ counts in $`m`$ bins:
$$W_m(\omega ,\varphi )=\frac{N!}{_{j=1}^mn_j!}$$
and $`n_j`$, the number of events placed in the $`j^{th}`$ phase bin depends on $`\omega `$, $`\varphi `$ and $`m`$.
If the period is known, this reduces to:
$$O(\omega )=\frac{1}{2\pi \nu }\frac{N!(m1)!}{(N+m1)!}_0^{2\pi }𝑑\varphi \frac{m^N}{W_m(\omega ,\varphi )}$$
(27)
In order to illustrate the difference between the information available from this treatment and from the Rayleigh test, a data set has been generated containing time-tagged random events with a constant mean rate, plus a periodic component. The results are shown in figure 5. A particular point to note is that although the Rayleigh power is always positive, even for pure noise, in the case of LOG(Bayesian odds), peaks do not become ’interesting’ until they become positive. This is because the ’degrees of freedom’ have been accounted for automatically and cause the offset seen in figure 5 so that peaks falling below $`\mathrm{log}(Odds)=0`$ are just those expected from noise. It can be seen much more clearly in the Bayesian Odds diagram that there is only one significant peak, at the period simulated.
The method has been used to detect a weak pulsar signal from SNR 0540-693 in ROSAT data, which could not be detected using a standard FFT technique. Moreover, the precision of determining the frequency was much higher for the Bayesian method than for $`\chi ^2`$ using epoch-folding. The frequency precision of the latter is determined mainly by the duration of the data and is not strongly influenced by the number of photons. Gregory and Loredo show that the Bayesian method obtains greater precision in parameter estimation with more photons. The method has also been used to detect 1600 day modulation in the long-term radio emission of an X-ray binary, with very non-uniformly sampled data and a Gaussian noise of unknown magnitude . There has been a recent independent use of a Bayesian method to calculate the upper limit to a pulsed flux at a known period, independent of pulse width and pulse phase.
## 6 Conclusions
The hope of this review is that the more commonly met data analysis problems may be approached by the cosmic ray worker with a more consistent and up to date approach. There have been a number of advances in recent years in the tools, and more importantly in the methods, available to cosmic ray experimenters to ensure that the maximum use is made of hard-won data. The traditional statistical methods have resulted in a measure of agreement on the ’correct’ way to look for sources from ON/OFF data, change points (bursts) in 1- and 2-dimensions and in periodicity. The application of these methods requires care to ensure that the ’degrees of freedom’ are kept under control and properly accounted for: many of the criticisms of claimed sources have been based on the latter.
New Bayesian methods of testing hypotheses have recently been proposed. A central theme of these methods is that classical methods often cloak ignorance in a way which distorts the results. There are claimed to be significant benefits to the use of Bayesian methods which derive from the requirement to be absolutely specific about the hypotheses and the methodology of marginalising nuisance parameters. In contrast to classical statistical methods, where various *statistics* may be generated from the same data, each with different assumptions, *degrees of freedom* and power, Bayesian methods provide a framework for describing completely the data and allow the direct comparison of specified hypotheses. A practical result of the philosophical differences between the approaches is that, rather than relying on a relatively easy-to-use, pre-packaged test statistic, with the accompanying dangers of hidden degrees of freedom, a Bayesian method requires the data interpreter to model the hypotheses precisely. The obvious disadvantages of this are claimed to be more than compensated by the directness of the link between the hypotheses and the data. Bayesian methods may require some time to become accepted in the field, in that the methodologies and ideas have not traditionally been part of the training of physicists; indeed may not have been as useful if physicists’ training in classical methods had been better.
## 7 Acknowledgements
The author would like to acknowledge useful discussions with P.S.Craig and M.Goldstein.
## 8 References
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# 1 Introduction
## 1 Introduction
The lepton flavor violating (LFV) neutrinoless conversion of a bound 1s-muon to electron in the field of a nucleus -
$$(A,Z)+\mu _b^{}e^{}+(A,Z)^{},$$
(1)
is known as one of the best probes to search for the hypothetical muon and electron flavor non-conservation . So far, the experiments - seeking for $`\mu e`$ conversion events have only succeeded to put upper bounds on the branching ratio
$`R_{\mu e^{}}=\mathrm{\Gamma }(\mu ^{}e^{})/\mathrm{\Gamma }(\mu ^{}\nu _\mu ),`$ (2)
i.e. the ratio of the muon-electron conversion rate relative to the total rate of the ordinary muon capture. The best upper limits have been extracted at PSI by the SINDRUM II experiments in the values (at 90% confidence level)
$$R_{\mu e^{}}<7.0\times 10^{13}\text{ for }^{48}\text{Ti target }\text{[5]},$$
(3)
$$R_{\mu e^{}}<4.6\times 10^{11}\text{ for }^{208}\text{Pb target }\text{[6]}.$$
(4)
These limits are improvements over the previous limits set at TRIUMF a decade ago using the same targets.
At present, two $`\mu e`$ conversion experiments are launched, the ongoing experiment at PSI using <sup>48</sup>Ti target , and the planned MECO experiment at Brookhaven using <sup>27</sup>Al target. The expected sensitivity on $`R_{\mu e^{}}`$ in the PSI experiment is $`10^{14}`$ while in the Brookhaven experiment it will be roughly
$$R_{\mu e^{}}<2\times 10^{17}\text{for}^{27}\text{Al target }\text{[8]},$$
(5)
which implies an improvement over the existing limits by about four orders of magnitude.
The MECO experiment is going to be conducted in a new $`\mu `$ beam line at the AGS, where the muons are produced using a pulsed proton beam . The proton energy will be chosen in the range of 8-20 GeV to optimize the $`\mu ^{}`$ flux per unit time. Furthermore, the number of electrons with energy $`E_e=104`$ MeV, equal to the energy for the coherent peak in <sup>27</sup>Al, is very much suppressed. This is in contrast with the case of $`\mu ^{}e^{}\gamma `$, where the electron flux from $`\mu ^{}e^{}\nu \overline{\nu }`$ decay is peaked at the energy of the electrons from $`\mu ^{}e^{}\gamma `$.
For such experiments the knowledge of nuclear transition matrix elements for all accessible $`\mu ^{}e^{}`$ channels of the targets employed are of significant importance -. In this work we use the transition matrix elements calculated for the aforementioned isotopes in the coherent mode to constrain the lepton flavor violating parameters of various Lagrangians predicting this exotic process (e.g. scalar and vector couplings, neutrino mixing angles and masses -, supersymmetric R-parity violating couplings - etc.). To this aim the recent experimental data and the expected sensitivity of the MECO experiment.
As is well known, only the coherent rate can be measured because it is free from background events from bound muon decay and radiative muon capture followed by a fully asymmetric $`e^+e^{}`$ pair creation . On the other hand, previous studies of $`\mu e`$ conversion rates - have shown that for all mechanisms the coherent mode dominates the process (1) which means that this is the most important channel. The incoherent reaction leading to excited nuclear states is suppressed due to Pauli blocking effects and it is much harder to calculate but its knowledge is also useful in order to determine the experimentally interesting quantity of the ratio of the coherent to the total $`\mu ^{}e^{}`$ rate (see Refs. -).
## 2 $`\mu ^{}e^{}`$ conversion within common extensions of Standard Model
The family quantum numbers $`L_e`$, $`L_\mu `$, $`L_\tau `$ are conserved within the standard model (SM) in all orders of perturbation theory. However this is an accidental consequence of the SM field content and gauge invariance. Physics beyond the SM can easily spoil this property. Processes like $`\mu e`$ conversion, which is forbidden in the standard model by muon and electron quantum number conservation, play an important role in the study of flavor changing neutral currents and possible physics beyond the SM.
On the particle physics side , there are many mechanisms of the $`\mu e`$ conversion constructed in the literature (see - and references therein). All these mechanisms fall into two different categories: photonic and non-photonic. Mechanisms from different classes significantly differ from the point of view of the nucleon and nuclear structure calculations. This stems from the fact that they proceed at different distances and, therefore, involve different details of the structure.
The long-distance photonic mechanisms are mediated by the photon exchange between the quark and the $`\mu e`$-lepton currents. These mechanisms resort to the lepton-flavor non-diagonal electromagnetic vertex which is presumably induced by the non-standard model physics at the loop-level. The $`\mu e`$ lepton-flavor violating loop can appear as the $`\nu W`$ loop \[Fig.1(a)\] with the massive neutrinos $`\nu _i`$ and the loop with the supersymmetric particles such as the neutralino(chargino)-slepton(sneutrino) \[Fig.1(d)\]. In the R-parity violating SUSY models there are also lepton-slepton and quark-squark loops generated by the superpotential couplings $`\lambda LLE^c`$ and $`\lambda ^{}LQD^c`$ respectively.
The short-distance non-photonic mechanisms contain heavy particles in intermediate states and can be realized at tree level \[Fig.2\], at 1-loop-level \[Fig.1(a,b,d,e)\] or at the level of box diagrams \[Fig.1(c)\]. The tree-level diagrams can be constructed in the R-parity violating SUSY models with the virtual Z-boson, squarks $`\stackrel{~}{u},\stackrel{~}{d}`$ and sneutrinos $`\stackrel{~}{\nu }_i`$ \[Fig.2\]. The 1-loop diagrams of the non-photonic mechanisms include the diagrams similar to those for the photonic mechanisms but with the Z-boson instead of the photon \[Fig.1(a,d)\] as well as additional Z-boson couplings to the neutrinos and neutralinos \[Fig.1(b,e)\]. The box diagrams are constructed of the W-bosons and massive neutrinos \[Fig.1(c)\] as well as similar boxes with neutralinos and sleptons or charginos and sneutrinos.
Our purpose is to calculate the contribution of the above-described mechanisms to the $`\mu ^{}e^{}`$ conversion branching ratio (2). At the first step we construct the effective $`\mu ^{}e^{}`$ conversion Hamiltonian for these mechanisms in terms of nucleon degrees of freedom of a nucleus involved in the process. This will allow us to accomplish the calculation of $`R_{\mu e^{}}`$ by applying the conventional nuclear structure methods using non-relativistic impulse approximation.
For the photonic diagrams of Fig. 1(a,d) the hadronic vertex represents the usual nucleon electromagnetic current
$$J_\lambda ^{(1)}=\overline{N}_p\gamma _\lambda N_p=\overline{N}\gamma _\lambda \frac{1}{2}(1+\tau _3)N(photonic)$$
(6)
$`N`$ is the nucleon isospin doublet $`N^T=(N_p,N_n)`$ with $`N_p`$, $`N_n`$ being the proton and neutron spinors.
In the case of the non-photonic mechanisms of Fig. 1 the hadronic currents are either the SM neutral currents coupled to Z-boson \[Fig. 1(a,b,d,e)\] or the effective currents derived after an appropriate Fierz transformation from the effective operators of the box diagrams \[Fig. 1(c)\] involving heavy particles. These currents can be written in terms of the nucleon doublet field as
$$J_\lambda ^{(2)}=\overline{N}\gamma _\lambda \frac{1}{2}\left[(3+f_V\beta \tau _3)(f_V\beta ^{^{\prime \prime }}+f_A\beta ^{^{}}\tau _3)\gamma _5\right]N(nonphotonic)$$
(7)
where $`f_V`$, $`f_A`$ represent the vector, axial vector static nucleon form factors ($`f_A/f_V=1.24`$) and the parameters $`\beta `$, $`\beta ^{}`$, $`\beta ^{\prime \prime }`$, for the models adopted take the values given in Ref. . The parameter $`\beta `$ is defined as $`\beta =\beta _1/\beta _0`$, with $`\beta _0`$ $`(\beta _1)`$ the isoscalar (isovector) couplings at the quark level. We should note that, in general, the parameters $`\beta `$, $`\beta ^{^{}}`$ and $`\beta ^{^{\prime \prime }}`$ are functions of $`sin^2\theta _W`$. For example, in the case of Z-exchange we have $`\beta ^{^{}}=3/2sin^2\theta _W=6.90`$.
The corresponding leptonic currents for the photonic and non-photonic mechanism of Figs. 1 are given in Ref. .
In the minimal SUSY model with a most general form of R-parity violation (for review see, for instance, ) all the possible tree-level diagrams for the $`\mu ^{}e^{}`$ conversion process are shown in Fig. 2 ,. To write down the low-energy $`\mu e`$ conversion Lagrangian at quark level one starts from the well known R-parity violating superpotential and integrates out heavy intermediate fields. For the diagrams of Fig. 2 the corresponding 4-fermion low-energy effective Lagrangian at the quark level takes the form (first order of perturbation theory)
$`_{eff}^q=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}j_\mu \left[\eta _L^{ui}J_{uL(i)}^\mu +\eta _R^{ui}J_{uR(i)}^\mu +\eta _L^{di}J_{dL(i)}^\mu +\eta _R^{di}J_{dR(i)}^\mu \right]`$ (8)
$`+{\displaystyle \frac{G_F}{\sqrt{2}}}\left[\overline{\eta }_R^{di}J_{dR(i)}j_L+\overline{\eta }_L^{di}J_{dL(i)}j_R\right].`$
($`i`$ runs over generations so that $`q_i=u_i,d_i`$ with $`u_i=u,c,t`$ and $`d_i=d,s,b`$) where the coefficients $`\eta _k^q`$ contain the $`R_p/`$ SUSY parameters . The color singlet currents $`J_{q_{L/R}(i)}^\mu `$ and $`J_{d_{L/R}(i)}`$ at the quark level are written as
$$J_{q_{L/R}(i)}^\mu =\overline{q}_i\gamma ^\mu P_{_{L/R}}q_i,J_{d_{L/R}(i)}=\overline{d}_iP_{_{L/R}}d_i.$$
where $`P_{_{L,R}}=(1\gamma _5)/2`$. The leptonic currents are written as
$$j^\mu =\overline{e}\gamma ^\mu P__L\mu ,j_{_{L/R}}=\overline{e}P_{_{L/R}}\mu $$
At the next step we rewrite Eq. (8), specified at the quark level, in terms of the nucleon degrees of freedom. This is usually achieved by utilizing the on-mass-shell matching condition and gives
$$_{eff}^N=\frac{G_F}{\sqrt{2}}\left[\overline{e}\gamma _\mu (1\gamma _5)\mu J^\mu +\overline{e}\mu J_++\overline{e}\gamma _5\mu J_{}\right].$$
(9)
where now the hadronic (nucleon) currents are
$`J^\mu `$ $`=`$ $`\overline{N}\gamma ^\mu \left[(\alpha _V^{(0)}+\alpha _V^{(3)}\tau _3)+(\alpha _A^{(0)}+\alpha _A^{(3)}\tau _3)\gamma _5\right]N,`$ (10)
$`J^\pm `$ $`=`$ $`\overline{N}\left[(\alpha _{\pm S}^{(0)}+\alpha _{\pm S}^{(3)}\tau _3)+(\alpha _{\pm P}^{(0)}+\alpha _{\pm P}^{(3)}\tau _3)\gamma _5\right]N,`$
The coefficients $`\alpha _K^{(0,3)}`$, with $`K=S,V,A,P`$, in Eq. (10) include the nucleon form factors (functions of the momentum transfer $`𝐪^2`$) for scalar, vector, axial vector and pseudoscalar contributions, respectively (for their definition see Ref. ). Since, however, the maximum momentum transfer $`𝐪^2`$ in $`\mu e`$ conversion is much smaller than the typical scale of nucleon structure ($`|𝐪|m_\mu /c`$, with $`m_\mu =105.6MeV`$, the muon mass), we can safely neglect the $`𝐪^2`$-dependence of the nucleon form factors.
## 3 Nucleon and nuclear structure dependence of the $`\mu ^{}e^{}`$ conversion branching ratio
One of the most interesting quantities in $`\mu ^{}e^{}`$ conversion, both from theoretical and experimental points of view, is the branching ratio $`R_{\mu e^{}}`$ defined in Eq. (2). The expression which gives $`R_{\mu e^{}}`$ in the case of the dominant coherent channel has been written as
$$R_{\mu e^{}}=\rho \gamma ,$$
(11)
where $`\rho `$ is nearly independent of nuclear physics and contains the flavor-violating parameters mentioned before. Thus, e.g. for photon-exchange $`\rho `$ is given by
$$\rho =(4\pi \alpha )^2\frac{|f_{M1}+f_{E0}|^2+|f_{E1}+f_{M0}|^2}{(G_Fm_\mu ^2)^2}$$
(12)
This expression contains the electromagnetic form factors $`f_{E0}`$, $`f_{E1}`$, $`f_{M0}`$, $`f_{M1}`$ parametrized in a specific elementary model .
The function $`\gamma (A,Z)`$ of Eq. (11) includes about all the nuclear information. By assuming that the total rate of the ordinary muon capture rate is described by the Goulard-Primakoff function $`f_{GP}`$, $`\gamma (A,Z)`$ is defined as
$$\gamma (A,Z)\gamma =\frac{E_ep_e}{m_\mu ^2}\frac{M^2}{G^2Zf_{GP}(A,Z)},$$
(13)
where $`G^26`$. Thus, the nuclear aspects of the $`\mu ^{}e^{}`$ branching ratio $`R_{\mu e^{}}`$ are mainly included in the matrix elements $`M^2`$ , which in the proton-neutron representation, are written as
$$M^2=(M_p+M_n)^2$$
(14)
In general, the transition matrix elements $`_{p,n}`$ in Eq. (14) depend on the final nuclear state populated during the $`\mu e`$ conversion. For ground state to ground state transitions ($`\mathrm{𝑔𝑠}\mathrm{𝑔𝑠}`$) in spherically symmetric nuclei the following integral representation is valid
$$_{p,n}=4\pi j_0(p_er)\mathrm{\Phi }_\mu (r)\rho _{p,n}(r)r^2𝑑r$$
(15)
where $`j_0(x)`$ the zero order spherical Bessel function and $`\rho _{p,n}`$ the proton (p), neutron (n) nuclear density normalized to the atomic number $`Z`$ and neutron number $`N`$ of the participating nucleus, respectively. The space dependent part of the muon wave function $`\mathrm{\Phi }_\mu (𝐫)`$, a spherically symmetric function, can be obtained by solving numerically the Schrödinger and Dirac equations with the Coulomb potential.
For the coherent rate in light nuclei the factorization approximation (see Ref. ) is very good and $`M_{\mathrm{coh}}^2`$ can be expressed in terms of the nuclear form factors $`F_Z(q^2)`$ (for protons) and $`F_N(q^2)`$ (for neutrons) which are easily estimated (for this channel only the ground state wave function of the studied nucleus is required). These form factors are defined as
$$F_Z=\frac{1}{Z}\underset{j}{}\widehat{j}(jj_0(qr)j)\left(V_j^p\right)^2,$$
(16)
$$F_N=\frac{1}{N}\underset{j}{}\widehat{j}(jj_0(qr)j)\left(V_j^n\right)^2,$$
(17)
and contain the single-particle orbit occupancies $`\left(V_j\right)^2`$ for the evaluation of which one must use a nuclear model in the proton-neutron representation, e.g. QRPA (for <sup>48</sup>Ti and <sup>208</sup>Pb see Ref. ), shell-model (for <sup>27</sup>Al see Ref. ) etc.
By using the Lagrangian (9) one can deduce a similar expression to that of Eq. (11) for the $`\mu e`$ conversion branching ratio as
$$R_{\mu e^{}}=\frac{G_F^2}{2\pi }\frac{p_eE_e(_p+_n)^2}{G^2Zf_{GP}(A,Z)}𝒬$$
(18)
where
$`𝒬=\mathrm{\hspace{0.17em}2}|\alpha _V^{(0)}+\alpha _V^{(3)}\varphi |^2+|\alpha _{+S}^{(0)}+\alpha _{+S}^{(3)}\varphi |^2+|\alpha _S^{(0)}+\alpha _S^{(3)}\varphi |^2`$
$`+2\mathrm{Re}\{(\alpha _V^{(0)}+\alpha _V^{(3)}\varphi )[\alpha _{+S}^{(0)}+\alpha _S^{(0)}+(\alpha _{+S}^{(3)}+\alpha _S^{(3)})\varphi ]\}`$ (19)
with
$$\varphi =(_p_n)/(_p+_n)$$
(20)
The quantity $`𝒬`$ in Eq. (19) depends weakly on the nuclear parameters determining the factor $`\varphi `$. In fact, the terms depending on $`\varphi `$ are small and in practice the nuclear dependence of $`𝒬`$ can be neglected. Thus, the corresponding upper bounds on $`𝒬`$ determine the sensitivity to the $`R_p/`$SUSY signals arriving at the detector in various $`\mu e`$ conversion experiments (see results below).
A usual approximate expression for the ratio $`\varphi `$ of Eq. (20), containing the nuclear parameters $`A`$, the atomic weight, and $`Z`$ the total charge of the nucleus, is the following
$$\varphi \stackrel{~}{\varphi }=(A2Z)/A,$$
(21)
The latter equation can be obtained assuming that $`M_pZF_Z\mathrm{\Phi }`$ and $`M_nNF_N\mathrm{\Phi }`$ where $`\mathrm{\Phi }`$ represents the mean value of the muon wave function $`\mathrm{\Phi }_\mu (r)`$. For light and medium nuclei $`F_ZF_N`$ and Eq. (21) is a good approximation. For isoscalar nuclei, i.e. $`A=2Z`$, $`\varphi \stackrel{~}{\varphi }=0`$ (see Sect. 4).
## 4 Results and Discussion.
The pure nuclear physics calculations needed for the $`\mu e`$ conversion studies involve mainly the integrals of Eq. (15). For the currently interesting nuclei $`Al`$, $`Ti`$ and $`Pb`$ the results for $`M_p`$ and $`M_n`$ are shown in Table 1. They have been calculated using proton densities $`\rho _p`$ from the electron scattering data and neutron densities $`\rho _n`$ from the analysis of pionic atom data . We have employed an analytical form for the muon wave function obtained by solving the Schrödinger equation using the Coulomb potential produced by the charge densities discussed before. In this way the nucleon finite size was taken into consideration. Moreover, we included vacuum polarization corrections as in Ref. .
In solving the Schrödinger equation we have used modern neural networks techniques which give the wave function $`\mathrm{\Phi }_\mu (r)`$ in the analytic form of a sum over sigmoid functions. Thus, in Eq. (15) only a simple numerical integration is finally required. To estimate the influence of the non-relativistic approximation on the muon wave function $`\mathrm{\Phi }_\mu (𝐫)`$, we have also determined it by solving the Dirac equation. The results do not significantly differ from those of the Schrödinger picture.
In Table 1 we also show the muon binding energy $`ϵ_b`$ (obtained as output of the Dirac and Schrödinger solution) and the experimental values for the total rate of the ordinary muon capture $`\mathrm{\Gamma }_{\mu c}`$ taken from Ref. .
Using values for $`M_p`$ and $`M_n`$ for a set of nuclei throughout the periodic table one can estimate the nuclear structure dependence of the quantities $`\rho `$ and $`𝒬`$. In Table 2 we show the results obtained for the ratio $`\varphi `$ of Eq. (20) and its approximate expression $`\stackrel{~}{\varphi }`$ of Eq. (21). We see that $`\varphi `$ is very small ($`0.17\varphi 0.04`$) and, because the terms of the quantity $`𝒬`$ depending on $`\varphi `$ are also small, we conclude that the nuclear dependence of $`𝒬`$ can be ignored as we have discussed in Sect. 2.
Since there is at maximum $`15\%`$ difference between $`\varphi `$ and $`\stackrel{~}{\varphi }`$ and because the isovector terms of the quantity $`𝒬`$ are small, especially for light nuclear systems, one can also use $`\stackrel{~}{\varphi }`$ in the expression (A,Z) dependence of $`𝒬`$.
The results of Table 1 can be exploited for setting constraints on the parameters of a specific gauge model predicting the $`\mu e`$ process. In Table 3 we quote the upper bounds for the quantities $`\rho `$ and $`𝒬`$ derived by using the recent experimental data on the branching ratio $`R_{\mu e^{}}`$ given in Eqs. (3), (4) and the expected experimental sensitivity of the Brookhaven experiment, $`R_{\mu e^{}}<2\times 10^{17}`$ \[see Eq. (5)\]. The limits of $`\rho `$ and $`𝒬`$ quoted in Table 3 are improvements by about four orders of magnitude over the previous ones.
We should stress that the limits on the quantities $`\rho `$ of Eq. (12) and $`𝒬`$ of Eqs. (19), are the only constraints imposed by the $`\mu e`$ conversion on its underlying elementary particle physics. One can extract upper limits on the individual lepton flavor violation parameters (couplings of scalar, vector currents, neutrino masses etc. ) under certain assumptions like the commonly assumed dominance of only one component of the $`\mu e`$ conversion Lagrangian which is equivalent to constrain one parameter or product of the parameters at a time. For example, in the case of $`R_p/`$SUSY mechanisms, in order to learn about the size and the regularities of possible violation of R-parity, one needs some information on the parameters $`\lambda ,\lambda ^{},\mu _i`$ and the sneutrino vacuum expectation values $`\stackrel{~}{\nu }_i`$ . Using the upper limits for $`𝒬`$ given in Table 3 we can derive under the above assumptions the constraints on $`\alpha _K^{(\tau )}`$ of Eq. (19) and the products of various $`R_p/`$parameters. Thus, the bounds obtained for the scalar current couplings $`\alpha _{\pm S}^{(0)}`$ in the $`R`$-parity violating Lagrangian for the <sup>27</sup>Al target are $`|\alpha _{\pm S}^{(0)}|<7\times 10^{10}`$. The limit for $`\alpha _{\pm S}^{(0)}`$ obtained with the data of Ti target is $`|\alpha _{\pm S}^{(0)}|<1.1\times 10^7`$, i.e. more than two orders of magnitude weaker than the limit of <sup>27</sup>Al.
With these limits it is straightforward to derive constraints on the parameters of the initial LFV Lagrangian. In Table 4 we list the most stringent constraints on the products of the trilinear $`R_p/`$couplings which we obtained from the experimental limit on $`\mu e`$ conversion in $`{}_{}{}^{48}Ti`$ \[see Eq. (3)\] and from the expected experimental sensitivity of MECO detector using <sup>27</sup>Al as a target material \[see Eq. (5)\]. The corresponding constraints for $`{}_{}{}^{208}Pb`$ are significantly weaker and they are not presented in Table 4. In this table $`B`$ denotes a scaling factor defined as
$$B=(R_{\mu e}^{exp}/7.010^{13})^{1/2},$$
which can be used for reconstructing the limits for the other experimental upper limits on the branching ratio $`R_{\mu e}^{exp}`$. As seen from Table 4, the $`\mu e`$-conversion limits on the products $`\lambda \lambda ^{}`$ are significantly more stringent than those previously known in the literature and given in the 2nd column. In Refs. it was shown that except few cases, the constraints on $`\lambda ^{}\lambda ^{}`$, $`\lambda \lambda ^{}`$ and $`\lambda \lambda `$ obtained from $`\mu e`$ conversion data are better than those derived from any other process.
As we have mentioned at the beginning, significantly better improvement on these limits is expected from the ongoing experiments at PSI and even better from the MECO experiment at Brookhaven . This would make the $`\mu e`$ conversion constraints better than those from the other processes in all the cases.
Before closing we should note that the last four limits for $`\lambda ^{}\lambda `$ in Table 4 originate from the contribution of the strange nucleon sea. These limits are comparable to the other $`\mu e`$ constraints derived from the valence quarks contributions.
## 5 Summary and Conclusions.
The transition matrix elements of the flavor violating $`\mu ^{}e^{}`$ conversion are of notable importance in computing accurately the corresponding rates for each accessible channel of this exotic process. Such calculations provide useful nuclear-physics inputs for the expected new data from the PSI and MECO experiments to put more severe bounds on the muon-number-changing parameters (isoscalar couplings, etc.) determining the effective currents in various models that predict the exotic $`\mu ^{}e^{}`$ process.
In the case of the R-parity violating interactions discussed here we have studied all the possible tree-level contributions to the $`\mu e`$ conversion in nuclei taking into account the nucleon and nuclear structure effects. We found new important contribution to $`\mu ^{}e^{}`$ originating from the strange quark sea in the nucleon which is comparable with the usual contribution of the valence $`u,d`$ quarks.
From the existing data on $`R_{\mu e^{}}`$ in $`{}_{}{}^{48}Ti`$ and $`{}_{}{}^{208}Pb`$ and the expected sensitivity of the designed MECO experiment we obtained stringent upper limits on the quantities $`\rho `$ and $`𝒬`$ introduced in Eqs. (12) and (19). They can be considered as theoretical sensitivities of a particular experiment to the $`\mu ^{}e^{}`$ conversion signals for various targets employed. Thus, these quantities are helpful for comparing different $`\mu e`$ conversion experiments. We also extracted the upper limits on the products of the trilinear $`R_p/`$parameters of the type $`\lambda \lambda ^{}`$ which are significantly more severe than those existing in the literature. As to the other products, the following observation was formulated in Refs. . If the ongoing experiments at PSI and Brookhaven will have reached the quoted sensitivities in the branching ratio $`R_{\mu e^{}}`$ then the $`\mu e`$ constraints on all the products of the $`R_p/`$parameters $`\lambda \lambda `$, $`\lambda ^{}\lambda `$, $`\lambda ^{}\lambda ^{}`$ measurable in $`\mu e`$ conversion will become more stringent than those from any other processes. This is especially important in view of the fact that no comparable improvements of the other experiments testing these couplings are expected in the near future.
T.S.K would like to acknowledge partial support of this work by Ioannina University grants and the NANP-99 Conference Organizers for hospitality at Dubna.
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# Precision Measurement of the Lifetime of the 3𝑑²𝐷_{5/2} state in 40Ca+
## 1 Introduction
In this paper we present a measurement of the natural lifetime $`\tau `$ of the $`3d^2D_{5/2}`$ metastable level in singly-ionised calcium, using quantum jumps of a single cold calcium ion in a linear Paul trap. The $`3D_{5/2}`$ level is of interest as the source of an optical frequency standard at 729 nm with $`1/\tau 1`$ Hz natural linewidth , as a means of testing atomic structure calculations , and as a diagnostic in astrophysics . In addition, our experiments offer insights into the diagnostics on the performance of the Paul trap, and highlight the need to take into account the spectral properties of semiconductor diode laser emission when such devices are used in atomic physics experiments.
Our result is $`\tau =1.168\pm 0.007\mathrm{s}`$, and is shown together with other recent measurements and theoretical predictions in figure 1. Our result is higher (a longer lifetime) than all the previous measurements, and differs by several standard deviations from most of them. It is possible that this discrepancy is at least partly due to a previously unrecognised source of systematic error, namely the presence of light of wavelength in the vicinity of 854 nm in the beam produced by a semiconductor diode laser emitting predominantly at 866 nm. This is discussed in section 4.1.
Precise knowledge of atomic structure for atoms or ions with a single electron outside closed shells is currently in demand for the analysis of atomic physics tests of electro-weak theory, especially measurements of parity violation in Cs . There has been a long-standing discrepancy at the 2% level between measured and theoretically predicted rates for electric dipole transitions . Up until now measurements of the metastable lifetimes (electric quadrupole transition rates) in Ca<sup>+</sup> have not been sufficiently precise to provide an independent test of ab initio calculations at this level of precision. Our measurement precision is $`0.6`$ %. Furthermore the electric quadrupole transition rate is harder to calculate accurately (recent calculated values have a 15% spread), so our result is of particular interest to atomic structure theory.
The paper is organized as follows. First we discuss, in section 2, the central features of our experimental method to measure $`\tau `$. Details of the apparatus are provided in section 3. This is a completely new apparatus which has not been described elsewhere, so we give a reasonably full description.
Section 4 presents our study of systematic effects in the experiment. These include collisions with the background gas, off-resonant excitation of the 854 nm $`3D_{5/2}`$$`4P_{3/2}`$ transition and of the 850 nm $`3D_{3/2}`$$`4P_{1/2}`$ transition, heating of the trapped ion, and noise in the fluorescence signal. Section 5 contains a clean demonstration of photon anti-bunching, using the random telegraph method, and section 6 presents the final accurate measurements of $`\tau `$.
## 2 Experimental method
The experimental method was identical, in principle, to that adopted by Block et al. ; preliminary experiments were carried out somewhat differently (see section 5). A single ion of <sup>40</sup>Ca<sup>+</sup> is trapped and laser-cooled to around 1 mK. The transitions of interest are shown in figure 2. Laser beams at 397 nm and 866 nm continuously illuminate the ion, and the fluorescence at 397 nm is detected by a photomultiplier. The photon count signal is accumulated for periods of duration $`t_b=10.01`$ ms (of which 2.002 ms is dead time), and logged. In our studies of systematic effects, and for our demonstration of photon antibunching, $`t_b`$ was set at $`22.022`$ ms. A laser at 850 nm drives the $`3D_{3/2}`$$`4P_{3/2}`$ transition. The most probable decay route from $`4P_{3/2}`$ is to the $`4S_{1/2}`$ ground state; alternatively, the ion can return to $`3D_{3/2}`$. However, about 1 decay in 18 occurs to $`3D_{5/2}`$, the metastable “shelving” level of interest. At this point the fluorescence abruptly disappears and the observed photon count signal falls to a background level. A shutter on the 850 nm laser beam remains open for 100 ms before it is closed, which gives ample time for shelving. Between 5 and 10 ms after the shutter has closed, we begin one “observation”, i.e., we start to record the photomultiplier count signal (see figure 3a). We keep observing the photon count, in the 10 ms bins, until it abruptly increases to a level above a threshold. This is set between the background level and the level observed when the ion fluoresces continuously. The signature for the end of a dark period is taken to be ten consecutive bins above threshold. We record the number of 10 ms bins in the observed ‘dark’ period. The 100 ms period of fluoresence also serves to allow the 397 nm and 866 nm lasers to cool the ion. After this we re-open the shutter on the 850 nm laser. This process is repeated for long periods of time (1 to 8 hours), the laser intensities being also monitored and the frequencies servo-controlled. Subsequent analysis of the large collection of dark times consists primarily of gathering them into a histogram, and fitting the expected exponential distribution, in order to derive the decay rate from the shelved state (see figure 4).
Note that we do not measure the length of the dark period from when the ion is first shelved. This is not necessary, since the probability of decay is independent of how long the ion has been in the metastable state. This gives us time to block the 850 nm light, in order to prevent both off-resonant excitation of the ion by this light, and the possibility of missed quantum jumps should this light rapidly (in less than 5 ms) re-shelve a decayed ion.
The data from a given run were analysed as follows. The raw data consist of a series of counts indicating the average fluorescence level in each bin of duration $`t_b`$. A threshold is set, typically at $`(2S_{\mathrm{dark}}+S_{\mathrm{bright}})/3`$, where $`S_{\mathrm{dark}}`$ is the mean count observed during dark periods, and $`S_{\mathrm{bright}}`$ the mean count observed during fluorescing periods. This setting is chosen because bright periods have more noise than dark periods. The number of consecutive bins below threshold is a single dark-time measurement $`x_i`$. The $`x_i`$ are expected to be distributed according to the exponential decay law, with Poissonian statistics describing the departures from the mean. It is appropriate to use a Poissonian fitting method, rather than least squares, because of the small numbers involved in part of the distribution (at large $`t`$). If $`n(t_i)`$ is the number of $`x_i`$ equal to $`t_i/t_b`$ then the cost function is
$`\mathrm{ln}\left[{\displaystyle \underset{i=0}{\overset{m}{}}}P\left(n(t_i)\right)\right]`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{m}{}}}[Ae^{\gamma t_i}+\mathrm{ln}(n(t_i)!)`$ (1)
$`n(t_i)\mathrm{ln}A+n(t_i)\gamma t_i]`$
where $`m`$ is the number of bins, and $`A`$ and $`\gamma `$ are two fitted parameters (obtained by minimising the cost function); they are the amplitude and decay rate in the assumed exponential decay $`A\mathrm{exp}(\gamma t)`$ of population of $`D_{5/2}`$. The residuals shown in figure 4 indicate that our data are well fitted by an exponential function. Only dark times of duration less than 5 s are included in the fit, since our data collection procedure misses some dark times longer than this. We found that the statistical error in the fitted parameters was consistent with the expected $`\sqrt{N}/N`$ value, where $`N`$ is the number of $`x_i`$ in the whole data set.
## 3 Apparatus
We use a linear radio-frequency (r.f.) Paul trap, combined with an all-diode laser system, to isolate and cool a single ion of <sup>40</sup>Ca<sup>+</sup>. The electrodes are made from stainless steel rods and mounted on two supports made from machinable ceramic (Macor); an end view is shown in figure 5. The radial r.f. electrodes, of diameter 1.2 mm, are centred at the corners of a square of side $`2.6`$ mm. The d.c. endcap electrodes, of diameter 1.0 mm are centred on the $`z`$-axis and positioned 7.2 mm apart. In addition to the electrodes which comprise the trap, there are a further four 1.6 mm diameter electrodes positioned in a similar configuration to that of the r.f. electrodes, but centred at the corners of a 8.4 mm square. These electrodes allow potentials to be added to compensate for stray electric fields in the trapping region.
The complete experimental apparatus is shown schematically in figure 6, to which we refer in the remainder of this section. The trap electrodes lie at the centre of a hexagonal stainless steel vacuum chamber. A high-voltage a.c. source RF supplies a drive voltage of frequency $`6.2`$ MHz and peak-to-peak amplitude 135 V for the radial electrodes, while high-voltage d.c. supplies DC provide the voltages for the endcaps (95 V in the present experiments) and compensation electrodes (typically around 60 V are applied to the upper two compensation electrodes, while the lower two are grounded).
The central chamber is pumped by a 25 l/s ion pump IP and a 30 l/s getter pump GP. An ion gauge IG on the opposite side of the main chamber monitors the pressure; this is below $`2\times 10^{11}`$ torr, the limit of the gauge’s sensitivity. To produce calcium ions in the trapping region, we use a calcium oven c and electron gun e. The oven is a thin-walled stainless steel tube filled with calcium granules, closed by crimping at each end and with a small hole at the centre pointing towards the trap region. The oven is heated by passing a 6 A current along its length, which produces a beam of calcium atoms. The electron gun consists of a tungsten filament, also resistively heated, enclosed within a grounded stainless steel “grid” with respect to which it is negatively biased by 50 V. To load the trap, the oven and electron gun are heated for a few minutes, the latter being left on for 10 s after the oven has been turned off. We capture a small cloud of approximately 10 calcium ions using this procedure. This cloud is reduced to a single ion by applying to one of the endcaps a low-voltage “tickle” oscillation, close to the axial resonance frequency of the trap, expelling ions until only one remains in the trap.
Violet light at 397 nm is generated by frequency-doubling 794 nm light from a master-slave diode laser system. A grating-stabilized master diode 794M, with a linewidth below 1 MHz, is locked to a stabilized low-finesse reference cavity RC and used to inject a slave diode 794S. Light from the slave is frequency-doubled by a 10 mm long Brewster-cut lithium triborate crystal LBO in an external enhancement cavity; Hänsch-Couillaud polarization analysis of light reflected by the cavity provides a feedback signal for a piezo-mounted mirror PZT used to lock the cavity length to the fundamental light. The measured characteristics of the doubling cavity are: finesse 130, enhancement 43, mode-matching efficiency 94%, input power 90 mW at 794 nm, output power 0.50 mW at 397 nm. Correcting for losses at the exit face of the crystal and the output coupler gives an internal crystal efficiency $`\gamma =54(5)\mu `$W/W<sup>2</sup>, some 20% greater than previously reported values and about 75% of the theoretical optimum efficiency calculated using a Boyd-Kleinman analysis and an effective non-linear coefficient for LBO of $`d_{\text{eff}}=0.855`$ pm/V . The 397 nm beam passes through a $`\lambda /2`$ waveplate and polarizing beam splitter (PBS) cube to provide intensity control, and a lens focuses it to a spot size of $`200\times 30\mu `$m (measured with a CCD camera) at the centre of the trap. The beam power used is typically 0.2 mW.
A grating-stabilized diode laser 866 provides repumping light at 866 nm whose intensity is adjusted by a $`\lambda /2`$ waveplate and PBS cube. A heated iodine bromide vapour cell IBr provides an absolute frequency reference—unfortunately not suitable for reliable locking—and a triple-pass acousto-optic modulator AOM shifts the laser light some 650 MHz into resonance with the $`3D_{3/2}4P_{1/2}`$ calcium transition. Spontaneous light near 854 nm emitted by the laser diode (see section 4.1.2) is rejected by a diffraction grating DG and aperture A. The transmitted 866 nm light is superimposed on the 397 nm beam using a PBS cube and focused onto the trap. The spot size used in the final lifetime measurements was $`250\times 130\mu `$m; the maximum beam power was 2 mW.
The light at 850 nm to shelve the ion is provided by a third grating-stabilized laser diode 850 situated on a separate optical table, and is directed into the ion trap via a polarization-preserving monomode optical fibre. The intensity is controlled by a $`\lambda /2`$ waveplate before the optical isolator. A mechanical shutter S before the fibre allows complete extinction of this light. The spot size at the trap is $`450\times 450\mu `$m and the maximum power $`0.5`$ mW.
For clarity in figure 6 the detection optics are shown in the plane of the diagram: they are actually vertically above the ion trap. A wide-aperture compound lens gathers fluorescence emitted by the trapped ion and images it onto an aperture to reject scattered light; further lenses re-image the light, via a violet filter, onto a photomultiplier PMT connected to a gated photon counter. The net collection efficiency of the detection system, including the 16% quantum efficiency of the PMT at 397 nm, is approximately $`0.12`$%. The peak photon count rate above background for a single cooled ion is typically 32 kHz.
A personal computer PC is used for data acquisition and control of the experiment; in particular it provides timing, logs PMT count data, controls the shutter S, and eliminates long-term drift of the 866 nm laser by locking it to the fluorescence signal from the trapped ion at the end of each 20 s acquisition period.
## 4 Searches for systematic effects
We now consider effects which alter the measured shelving periods from those appropriate to an unperturbed ion subject only to spontaneous decay. These are of two distinct types: those which alter the shelving periods themselves, because the ion is perturbed in some way, and those which cause systematic error in the process of measurement. We consider the two in turn. The significant difference between our final result and previous work calls for a detailed discussion here.
### 4.1 Perturbations to the ion
When the ion is fluorescing, it is cycling between the levels $`4^2S_{1/2},\mathrm{\hspace{0.33em}4}^2P_{1/2}`$ and $`3^2D_{3/2}`$. We denote this system of levels by $`\mathrm{\Lambda }`$. During the shelving period, the ion is subject to radiation from two of the lasers (397 nm and 866 nm), to thermal radiation, and to the fields associated with the trap. There may also be collisions with the background gas. Any of these perturbations can transfer the ion to the $`\mathrm{\Lambda }`$ system. We consider them in turn.
#### 4.1.1 Electric field
The ion experiences a static electric field because of imperfect compensation. The most significant effect on the internal state of the ion is to mix the $`3D_{5/2}`$ and $`4P_{3/2}`$ levels, so that the measured lifetime is shortened by the presence of the induced $`4P_{3/2}`$$`4S_{1/2}`$ strong electric dipole transition. In fact, using the known matrix elements, one finds that the effect is negligible; the induced transition probability is $`9.0\times 10^{14}E^2`$ s<sup>-1</sup>, so that a field of 300 kV m<sup>-1</sup> (three orders of magnitude larger than the typical compensating fields used) would be necessary to produce a 1% reduction in the lifetime. We note, however, that there is another effect associated with an imperfectly compensated field, that of heating during the shelving period; we consider this in section 4.2.1.
#### 4.1.2 Laser radiation
From the standard theory of atom-light interaction, one finds that transitions from a lower level (1) to a higher level (2) in the multi-level ion can be stimulated by radiation of intensity $`I`$ and angular frequency $`\omega _L`$ at a rate $`R_{12}`$ (averaged over all Zeeman components) given by
$$R_{12}=\frac{2J_2+1}{2J_1+1}\frac{\pi ^2c^3}{\mathrm{}\omega _{12}^3}A_{21}\frac{I}{c}g(\omega _L\omega _{12})$$
(2)
where $`J_1`$ and $`J_2`$ are the total angular momenta of the levels, $`\omega _{12}`$ is the atomic resonance angular frequency, and $`g(\omega _L\omega _{12})`$ is the normalised lineshape function. In our case we may assume Lorentzian lineshapes, so that
$$g(\omega _L\omega _{12})=\frac{\mathrm{\Gamma }/(2\pi )}{(\omega _L\omega _{12})^2+\mathrm{\Gamma }^2/4}$$
(3)
with the linewidth $`\mathrm{\Gamma }=1/\tau _2`$ determined by the sum of all the decay processes from the upper level.
Since Block et al. reported a significant dependence of the shelving time on the intensity of the repumping laser, we first consider excitation from $`3D_{5/2}`$ to $`4P_{3/2}`$ by light at 866 nm. The lifetime $`\tau _2`$ of the $`4P_{3/2}`$ level has been measured to be $`6.924\pm 0.019`$ ns . We use the value $`A_{21}=(7.7\pm 0.3)\times 10^6`$ s<sup>-1</sup> from ab initio and semiempirical atomic structure calculations. The $`4`$% uncertainty is our own estimate, based on the variation among the published calculations, and on the fact that the calculations produce other electric dipole matrix elements in agreement with experiment to better than this level of precision. In any case a 10% error in the value of $`A_{21}`$ would have a negligible influence on our final result.
To excite the $`3D_{5/2}`$$`4P_{3/2}`$ transition on resonance would require radiation of wavelength 854 nm. From equation (2), we find that with 866 nm light the rate is $`9.9\times 10^5`$ s$`{}_{}{}^{1}/`$mW mm<sup>-2</sup>. The probability that the ion will subsequently decay back to $`3D_{5/2}`$ is the branching ratio $`b=A_{21}\tau _2=0.053`$, so the rate at which it will be transferred to the $`\mathrm{\Lambda }`$ system is $`(1b)R_{12}=9.4\times 10^5`$ s$`{}_{}{}^{1}/`$mW mm<sup>-2</sup>. In our experiments, it is convenient to choose the 866 nm intensity such that the $`3D_{3/2}4P_{1/2}`$ transition is saturated. This requires $`I_{866}0.08`$ mW mm<sup>-2</sup> if the laser frequency is set on resonance (taking this $`A_{21}`$ coefficient $`8.4\times 10^6`$ s<sup>-1</sup> ), but to avoid the need to control this frequency precisely a much higher intensity was used, typically $`1.5`$ mW in a spot size of $`250\times 130\mu `$m, or 30 mW mm<sup>-2</sup>. This results in a contribution to the depopulation rate of the $`3D_{5/2}`$ level of around 0.3% of that due to spontaneous emission to the ground level.
This is far smaller than the experimental result of Block et al. However, in our own preliminary work we also found that there was a significant dependence of the shelving time on the intensity of the repumper laser, of the same order of magnitude as that reported by Block et al., and some 200 times larger than the theoretical value given above. This suggested that the 866 nm laser was emitting some radiation much closer in wavelength to 854 nm which was primarily responsible for the shortening of the apparent lifetime of the $`3D_{5/2}`$ level. This laser is a semiconductor diode device, operated with an extended cavity by use of a Littrow-mounted diffraction grating. Without the grating, the laser would operate close to 854 nm. Since any laser produces spontaneous emission over its gain profile, as well as stimulated emission at the lasing wavelength, there was in our preliminary work radiation incident on the ion in the vicinity of 854 nm. One might expect the intensity of this radiation to be greatly reduced because of the long (3m) beam path. In fact, this is not the case, because the emitting region in the laser is small (dimensions of order microns) so the spontaneous component is well collimated; like the coherent component, it is transported to the trap with little loss.
We therefore investigated the spectrum of the light from the laser by means of a diffraction grating, and found it to contain a broad background from 840 nm to 870 nm. When the total laser power near the ion trap was 2 mW, the power in the broad component was 8 $`\mu `$W, and the spot sizes were similar. This implies a mean spectral density around 854 nm sufficient to cause de-shelving rates considerably higher than those we observed; we ascribe the lower observed rates to the structure of the background. This was not resolved in our investigation, but is expected on the basis of other experimental work to include a long series of spikes separated by the longitudinal mode spacing of the diode (50 GHz). These occur because the gain is so high in these devices that there is some amplification right across the gain profile; evidence that this occurs in our laser is provided by the high degree of polarization (90%) of the background. It is thus reasonable to conclude that the $`3D_{5/2}`$$`4P_{3/2}`$ atomic resonance falls between peaks in the 50 GHz spaced comb.
To study the effects of this background radiation we reduced its intensity at the position of the ion. This was done in two stages: first, we reduced it by a factor of 25 using an interference filter centred on 866 nm. The observed de-shelving rate fell from $`1.85\pm 0.06`$ s<sup>-1</sup> to $`0.87\pm 0.02`$ s<sup>-1</sup>, the two measurements being taken at the same 866 nm intensity. This provided convincing evidence that the de-shelving was indeed caused by the postulated mechanism. We therefore replaced the filter with a diffraction grating and iris aperture, arranged to pass light at 866 nm. The power transmitted by the system in the vicinity of 854 nm was then reduced compared with the unfiltered laser by three and a half orders of magnitude (for a given intensity at 866 nm), and the remaining light was scattered, thus increasing the illuminated area in the vicinity of the ion by a factor measured to be 40, making a net intensity reduction $`8\times 10^6`$. Under these conditions transitions stimulated by this radiation become much less probable than those due to the 866 nm light itself. The observed de-shelving rate was $`0.858\pm 0.007`$ s<sup>-1</sup>. De-shelving rates $`\gamma `$ at various settings of the intensity $`I`$ of the 866 nm laser beam as measured in these various experiments are shown in figure 7a. A straight line of the form $`\gamma =\gamma _0+\alpha I`$ fitted to the points for which the grating and iris system were in place gives $`\gamma _0=0.857\pm 0.016`$ s<sup>-1</sup>, $`\alpha =(1.5\pm 6)\times 10^3`$ s<sup>-1</sup>/mW mm<sup>-2</sup>. The value of $`\gamma _0`$ is consistent with our final more accurate measurements, given in section 6 below, while the slope is consistent with zero (and with the very small theoretical value for de-shelving by 866 nm radiation given above). The intercept with unfiltered light is greater than that given by our final data because our method of varying the laser intensity did not alter the unpolarized component of the background radiation.
Spontaneous emission from the repumper laser appears also to be likely to account for the observations of Block et al. It is possible that it was present in previous work on the $`3D_{5/2}`$ lifetime, and unaccounted for, explaining the lower values obtained by all earlier workers. This does not immediately apply to Block et al., however, because their measurement involved an extrapolation to zero laser intensity using neutral density filters.
The 866 nm laser might also generate radiation near 854 nm if it went multimode owing to a degradation of the alignment of its own external cavity. If intermittent and at a low level, this effect could occur undetected. However, multimode operation was found in practice to have an all-or-nothing character: if it occurred at all then it was obvious from a large increase in the noise of the fluorescence signals, and in such a case the data were discarded.
The grating and iris system were in place for our final data sets. Figure 7b shows our final rate measurements plotted against the intensity of the repumper laser. The observations are consistent with the theoretical value for the slight dependence on 866 nm intensity. Our final result for the lifetime was obtained with a slope fixed at the theoretical value (see section 6).
The radiation at 397 nm is obtained by frequency-doubling and so does not contain any significant background light. There is no transition from $`3D_{5/2}`$ near enough to this wavelength to cause significant de-shelving. A check for this or some other (unidentified) effect was nevertheless carried out, where we changed the power of this light by a factor 2, and we observed, to ten percent precision, no effect on the deshelving rate.
#### 4.1.3 Thermal radiation
The rate of de-shelving is $`(1b)B_{12}\rho (\omega _{21})`$, where $`\rho (\omega )`$ is the energy density per unit frequency interval in the thermal radiation, and $`B_{12}=g_2\pi ^2c^3A_{21}/(g_1\mathrm{}\omega _{21}^3)`$. In a thermal cavity we obtain the rate
$$(1b)\frac{2J_2+1}{2J_1+1}A_{21}e^{\mathrm{}\omega _{21}/k_\mathrm{B}T}$$
(4)
for $`\mathrm{}\omega _{21}k_\mathrm{B}T`$. For the 854 nm transition in a room-temperature cavity, the rate is of order $`10^{18}`$ s<sup>-1</sup> so is negligible. Non-negligible rates can be obtained, however, when we consider the radiation produced by room lights or the filament of an ion gauge. The spectral energy density is reduced from the value in equation (4) by a geometrical factor approximately equal to $`S/(4\pi r^2)`$ where $`S`$ is the area of the hot filament and $`r`$ is its distance from the ion, if the ion is within line of sight of the filament. For example, taking $`T=1700`$ K, $`S=20`$ mm<sup>2</sup>, $`r=30`$ cm, we obtain a rate $`5\times 10^3`$ s<sup>-1</sup>. In our vacuum system, although the ion gauge is at this distance from the trapped ion, it is not in line of sight, so we expect the rate of this process to be well below this value.
#### 4.1.4 Collisional effects
The ion in the shelved level can undergo a collision with an atom of the background gas. Either or both of two processes may then occur: the ion may gain a significant amount of kinetic energy, and it may be transferred to another state. This is a source of error since in both cases the apparent shelving time — the interval during which fluorescence is not observed — will be affected. To investigate the nature and frequency of collisional effects we monitored the fluorescence from the ion for 8 hours with the 397 nm and 866 nm radiation present but with no laser operating at 850 nm to take the ion to the $`4P_{3/2}`$ level. The diffraction grating and aperture were in place. During this period we observed abrupt disappearance of the fluorescence (within the resolution of the 22 ms bins) on 17 occasions. It reappeared after times of the order of a second, with 6 “dark periods” as short as a few tens of milliseconds. The reappearance was generally abrupt, but in 5 of the longer periods it was more gradual, occurring over several bins.
One non-collisional effect which can lead to loss of fluorescence in this test is shelving in the $`3D_{5/2}`$ level. Shelving caused by the 866 nm radiation exciting the $`3D_{3/2}`$ to $`4P_{3/2}`$ transition is negligible, because the spontaneous decay to $`3D_{5/2}`$ has such a low branching ratio (the excitation rate can be calculated from the branching ratio $`b`$ and the $`3D_{3/2}4P_{3/2}`$ coefficient $`A_{21}=0.91\times 10^6`$ s<sup>-1</sup> ). In contrast, it is not possible to rule out excitation by the spontaneously emitted light from the 866 nm laser at 850 nm because the rate depends on the spectral distribution of the spontaneous emission at 850 nm; if a peak happened to be close to the frequency of the 850 nm transition the process could be responsible for a significant number of the observed events. However, we would then expect an abrupt reappearance of the fluorescence on a time scale of the order of a second, and the data then suggest that the upper limit for events of this type is around eight. These could alternatively be caused by fine-structure changing collisions; the ion in the $`3D_{3/2}`$ level can be transferred to $`3D_{5/2}`$ by a long range collision which may not transfer significant kinetic energy. At our working pressure of order $`10^{11}`$ mbar we would expect about 8 such events on the basis of an approximate value of the collisional mixing rate which has been determined for conditions similar to ours . Our observations thus provide a rough upper limit on the background gas pressure in our system directly at the location of the ion, to confirm our ion gauge reading. We note that the rates for $`D_{5/2}D_{3/2}`$ and $`D_{3/2}D_{5/2}`$ are approximately equal at room temperature, therefore the present experiment does give an indication of the de-shelving rate due to fine-structure changing collisions in our $`D_{5/2}`$ lifetime measurements.
The other events are more obviously characteristic of collisions. In particular, the more gradual reappearance of the fluorescence after a long dark period is likely to be associated with the ion being cooled again after acquiring a significant amount of kinetic energy. On the basis of our measured pressure and reasonable estimates of cross-sections, we do not expect more than two or three collisions of this type. Some of the very short dark periods are likely to be much smaller perturbations by relatively distant collisions. If such a perturbation were to occur during the lifetime measurement itself, fluorescence would not be observable for a period whatever state the ion was in after the collision. Fortunately, for the purposes of estimating the uncertainty in the lifetime measurement, detailed interpretation of the events is not necessary. An upper limit to the error introduced can be found by assuming that all 17 events are due to collisions, giving a rate of $`6\times 10^4`$ s<sup>-1</sup>, and that such a collision occurring in the experimental runs themselves while the ion was in the shelved level would have delayed (or hastened) the reappearance of fluorescence by an average of order 1 second. The net contribution to the measured rate has thus an upper limit of $`\pm 6\times 10^4`$ s<sup>-1</sup> which is an order of magnitude below the statistical error.
### 4.2 Systematic effects in the measurement process
Both the fluorescence and the background signal on which it is superposed are subject to fluctuations. As well as the random error this introduces into the measurement, because the instant at which fluorescence resumes is subject to statistical uncertainty, there are systematic effects. For example, a large fluctuation in the background can suggest that the ion has decayed while it is still shelved; our data analysis procedure had to be developed and tested to minimise errors due to such effects. Further, the ion is not cooled during the shelved period, and if there is significant heating the fluorescence may be reduced for a time after the decay.
#### 4.2.1 Heating during the shelved periods
In our preliminary experiments, we found that when an electric field was applied in the vertical direction the fluorescence reappeared only gradually after shelving, over periods of 10–50 ms. We ascribe this to heating. When the d.c. electric field in the trap is not zero, the ion experiences the r.f. driving field, which is much more noisy than the d.c. field, and so the ion motion heats during the shelved periods when it is not laser cooled. We calculate that if the ion heats up to room temperature, it would take the lasers approximately 50 ms to cool the motion down again and thus for the fluorescence to reappear. Evidence supporting this interpretation is given by the fact that the non-abrupt reappearance of fluorescence was correlated to the duration of the dark period, being more likely for longer dark periods. We therefore took care to ensure this phenomenon was not present in the runs used for our final data set. This was done by nulling the vertical field carefully before each run.
In our linear trap geometry the r.f. field is 2-dimensional. To null the vertical field we adjusted the voltage on one of the d.c. field compensation electrodes, so as to minimise the linewidth of the ion fluorescence as a function of 397 nm laser frequency (at lowered 397 nm laser power). The 397 nm laser beam used for shelving measurements enters along the direction $`(\sqrt{3},0,1)`$ (where the $`xz`$ plane is horizontal); this is sensitive to horizontal micromotion and hence the vertical field. We were able to null this field to $`\pm 3`$ V/m, and most runs were carried out with nulling to $`\pm 10`$ V/m. We measured the remaining heating rate, when this field was as well compensated as we could make it, by blocking the cooling laser beams (397 nm and 866 nm) for long periods, and looking for a non-abrupt return of the fluorescence when the beams were unblocked. No delay was observed (the limit of sensitivity being the bin size of 22 ms) unless the cooling lasers were blocked for more than 10 minutes, some 600 times longer than the shelving times occurring during the measurements themselves. The consequent systematic error in the lifetime can thus be safely neglected; heating is likely to be relatively slow when the ion is first shelved, but even assuming that the delay of fluorescence is linear with shelving time the error is only of order 30 $`\mu `$s. We note that there were no soft edges in the runs used for our final data set.
We found that a horizontal field gave a much smaller heating effect, and indeed some of our final data points were taken in the presence of a horizontal field of order 200 V/m. However, for our final experimental run, a further beam was introduced along $`(0,2,1)`$ to allow horizontal field compensation; this beam was blocked once the compensation was optimised. Although not important for the present lifetime measurements, accurate compensation is necessary for our future work in which the ion is to be cooled below the Doppler limit.
#### 4.2.2 Tests of the reliability of the data analysis
We consider three possible sources of systematic error in the data analysis in turn. First, it is straightforward to show analytically that approximating the continuous exponential distribution of dark times as a histogram causes no systematic errors in the fitted value of $`\tau `$. Second, numerically simulated data was used to show that varying the threshold level used in the data analysis, did not systematically change the fitted value of $`\tau `$ by more than $`0.1`$ ms. Our real data has a further property: intensity fluctuations. The $`\tau `$ value deduced from a real data set was found to vary as a function of threshold by amounts of order $`\pm 0.5`$ ms. We estimate this may introduce a systematic error of order $`0.5`$ ms. Third, repeated fits to numerically simulated data sets were used to show that any systematic error arising purely from the exponential fitting procedure was less than $`0.25`$ ms. They also permitted us to verify the uncertainty (one standard deviation) in the fitted value of $`\tau `$, derived from the cost function surface.
These effects give a total uncertainty of $`\pm 0.6`$ ms.
## 5 Demonstration of photon antibunching
The quantum jump method provides a convenient means to demonstrate photon antibunching , that is, generation of a light field whose second order coherence $`g^{(2)}(t)`$ falls below 1 as $`t0`$. The second-order coherence is the normalised autocorrelation function of the intensity. A classical definition of intensity leads to $`g^{(2)}(t)<g^{(2)}(0)`$ and $`g^{(2)}(0)>1`$, whereas a definition in terms of quantum electric field operators allows all values . Therefore observation of a $`g^{(2)}(0)`$ value below 1 is a direct signature of the quantum nature of the radiation field. Quantum jump observations yield $`g^{(2)}`$ easily, since whenever a de-shelving jump is observed, we may deduce, with close to 100% reliability, that one photon at 729 nm has been emitted by the ion . Therefore we deduce the presence of a 729 nm radiation field which has $`g^{(2)}(t)=n(t)n(0)/\overline{n}^2`$, where $`n(t)`$ and $`n(0)`$ are the numbers of jumps observed in two equal time intervals separated by a time $`t`$. This quantity is plotted in figure 8 for a typical data set in our experiment. For this study, we did not use the method described in section 2, in which the 850 nm light was blocked when the ion was shelved. Instead, this light was left permanently on, and the fluorescence monitored continuously. We therefore obtained the well-known ‘random telegraph’ signal, figure 3b. Since we detect very close to all the 729 nm photons, we can be confident the second order coherence does not vary at time-scales too short for us to detect. Therefore the antibunching signature $`g^{(2)}(0)<1`$ is clearly demonstrated.
We can understand the complete form of $`g^{(2)}(t)`$ by solving the rate equations for the populations of the ion’s levels. All the relevant processes are fast except for two, namely spontaneous decay from $`D_{5/2}`$ to $`\mathrm{\Lambda }`$ and excitation from $`\mathrm{\Lambda }`$ to $`D_{5/2}`$, so the problem reduces to a two-level system. If we take $`t=0`$ to be the centre of a short time interval (one bin) during which a jump occurred, so that $`n(0)=1`$, then the probability for the atom to emit a 729 nm photon at time $`t`$ is proportional to the population of $`D_{5/2}`$ at $`t`$, given that it was zero at $`t=0`$. We thus obtain
$$g^{(2)}(t)=1e^{(R+\gamma )t}$$
(5)
where $`R`$ is the excitation rate, $`\gamma `$ the decay rate. This curve is plotted on figure 8 with no free parameters ($`R`$ is taken as the inverse of the mean duration of bright periods in the same data set). The agreement between data and theory is evidence that we have a good understanding of our experimental situation.
## 6 Result and discussion
Our final data set consisted of four 8-hour runs, and two 2-hour runs, all using the experimental method of section 2 (850 nm laser blocked during dark periods). The de-shelving rates observed in these runs, as obtained by the analysis described above, are shown in figure 7b. The size of each error bar is equal to the statistical uncertainty emerging from the analysis. The straight line through the data is a single-parameter weighted least-squares fit. The line is given the theoretically expected slope, and the best fit intercept is found to be $`\gamma _0=0.856\pm 0.004`$ s<sup>-1</sup>.
The de-shelving rate $`\gamma `$ obtained in the experiment may be written
$$\gamma =\frac{1}{\tau }+\underset{i}{}\gamma _i$$
(6)
where $`\tau `$ is the natural lifetime we wish to measure, and the $`\gamma _i`$ are due to other processes which contribute to the measured rate, chiefly laser excitation of transitions in the ion and collisional processes. The effect of off-resonant excitation by 866nm radiation is already accounted for in our fitted intercept $`\gamma _0`$. However, uncertainty ($`\pm 30\%`$) in the light intensity at the ion makes this accounting imprecise, leading to a further systematic error in $`\gamma _0`$ at the level of 0.001 s<sup>-1</sup>. This and other contributions to the systematic error, from collisional processes and remaining background light from the 866nm laser, are shown in Table 1. The statistical error from the data fitting dominates. The systematic effects are independent of one another, so we add their errors in quadrature. Adding this total systematic error linearly to the statistical error, the value we obtain for the natural lifetime of the $`D_{5/2}`$ level is
$$\tau =1.168\pm 0.007\text{ s}$$
(7)
We have already noted the significant difference between this result and earlier measurements. Here we compare it with theory. Our result differs from all the reported ab initio calculations by amounts large compared with our $`0.6`$ % experimental uncertainty. The relative difference $`(\tau _{\mathrm{theor}.}\tau _{\mathrm{meas}.})/\tau _{\mathrm{meas}.}`$ is $`10`$% for a recent calculation based on the Brueckner approximation , $`+5.6`$% for relativistic many-body perturbation theory , and $`2.6`$% for multi-configurational Hartree Fock (MCHF) calculations . Of these calculations, those which produced the smallest discrepancy with our measured $`D`$ state lifetime produced the largest discrepancy with the measured $`P`$ state lifetimes obtained by Jin and Church . A semi-empirical calculation based on MCHF and core polarization , and a calculation using related methods , give a value close to our measured result (relative difference $`0.7`$%). The most natural interpretation of these observations is that the ab initio calculations of $`\tau `$ carried out so far have a precision of at best a few percent, and success at calculating electric dipole matrix elements does not guarantee the same degree of success with other parameters such as electric quadrupole matrix elements.
The ratio between the lifetimes of $`D_{3/2}`$ and $`D_{5/2}`$ is much less sensitive to imperfections in the calculations, and is given primarily by the frequency factor $`\omega ^5`$, which leads to a ratio $`1.022`$. The factor emerging in recent calculations of the two lifetimes is $`1.0335`$ , $`1.0283`$ , $`1.0175`$ . We take the standard deviation of these results to indicate the theoretical uncertainty on their mean, giving $`\tau _{3/2}/\tau _{5/2}=1.026\pm 0.007`$. Combining this with our measurement of $`\tau _{5/2}`$ gives $`\tau _{3/2}=1.20\pm 0.01`$ s for the lifetime of the $`3d^2D_{3/2}`$ level in <sup>40</sup>Ca<sup>+</sup>.
## 7 Conclusion
We have described a new linear ion trap apparatus and its use to measure the lifetime of the $`D_{5/2}`$ level in <sup>40</sup>Ca<sup>+</sup> by quantum jump measurements on a single trapped ion. Our result is more precise than previous measurements, and significantly larger. We believe this discrepancy is mostly explained by previously unrecognised systematic errors, which tend to make the lifetime appear shorter than it in fact is. Our measurement provides a new precise test of ab initio atomic structure calculations; it is in only moderate agreement with currently reported calculations.
We have discussed the spectral distribution of light emitted by a semiconductor diode laser operated in an extended cavity, emphasizing that in experiments where unwanted excitation of allowed atomic transitions is to be avoided, it is necessary to take into account the weak broad component in the spectrum, which extends many nanometres away from the main lasing wavelength. This is significant in many experiments in atomic physics, especially coherent atom optics, ‘dark’ optical lattices, frequency standards and quantum information processing.
Finally, we have discussed various diagnostics on the performance of the ion trap, made possible by careful analysis of long periods of operation of the trap. These include an upper limit for the pressure in the system, and the heating rate in the trap. We demonstrated photon anti-bunching, finding it in agreement with theoretical expectations.
Our uncertainty is dominated by statistics, so it would be possible to obtain significantly higher precision by accumulating more data, though for single-ion work this would require integration for several weeks. To go further, one would need better knowledge of the laser intensities and collisional processes.
## 8 Acknowledgements
We would like to acknowledge helpful correspondence with G. Werth, and a preprint of . L. Favre, E. Hodby, M. McDonnell and J. P. Stacey contributed to development of the apparatus, D. T. Smith designed the high voltage r.f. supply, and we are grateful for general technical assistance from G. Quelch. We thank R. Blatt and the Innsbruck group (especially F. Schmidt-Kaler, C. Roos, M. Schulz) for many helpful discussions and guidance, and similarly R. C. Thompson and D. N. Segal. This work was supported by EPSRC (GR/L95373), the Royal Society, Riverland Starlab and Oxford University (B(RE) 9772).
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# 1 Introduction
## 1 Introduction
In AdS/CFT (see for a review), bulk geometries which are only asymptotically $`AdS_5`$ near the timelike boundary are dual to relevant deformations of the CFT or to non-conformal vacua. Far from the timelike boundary, various singularities might arise. There must be some restrictions on the type of singularity which is allowed: for instance, we take it as obvious that negative mass Schwarzschild should be excluded on the grounds of vacuum stability . The goal of this paper is to address the question of what singularities are allowed in geometries of the form
| $`ds^2`$ | $`=e^{2A(r)}(dt^2+d\stackrel{}{x}^2)+dr^2`$ |
| --- | --- |
| $`\stackrel{}{\phi }`$ | $`=\stackrel{}{\phi }(r),`$ |
(1)
where by assumption $`A(r)r/L`$ and $`\stackrel{}{\phi }(r)(const)`$ for large $`r`$. $`L`$ is the radius of curvature of the asymptotically $`AdS_5`$ region. Geometries of the form (1) arise in the gauged supergravity description of certain vacuum states on the Coulomb branch of $`𝒩=4`$ super-Yang-Mills theory , and are also suspected to be the supergravity duals of most renormalization group flows arising from relevant deformations of the lagrangian of the CFT .
The classical equations of motion, plus either positivity of the scalar kinetic terms or the more general “null energy” condition, together imply $`A^{\prime \prime }(r)0`$ . Since we assume $`A(r)r/L`$ with $`L>0`$ for large $`r`$, $`A(r)`$ must in fact be monotonically increasing with $`r`$. Thus a curvature singularity can arise for the metric (1) if $`A(r)\mathrm{}`$ at some finite $`r_0`$.<sup>1</sup><sup>1</sup>1Singularities where $`A(r)\mathrm{}`$ at finite $`r_0`$ are the only type that will be encountered in this paper, but for the sake of completeness let us review three other possibilities. First, it could be that $`A(r)`$ is defined for all $`r`$, but $`A^{}(r)\mathrm{}`$ as $`r\mathrm{}`$. Supersymmetric examples with this behavior arise in certain $`U(1)`$ gauged supergravities. Second, an $`S^5`$ shell of D3-branes produces a near-horizon geometry which is $`AdS_5\times S^5`$ outside the shell and flat $`𝐑^{10}`$ inside the shell . The shell itself produces a curvature singularity at finite $`A(r)`$. Similar singularities arise in Horava-Witten theory compactified on a Calabi-Yau three-fold with M5-branes in the bulk (see for example ): in the five-dimensional description, there is a jump in $`A^{}(r)`$ across each M5-brane. This sort of singularity also arises in type I with D8-branes in the bulk. Finally, there can be a codimension one orientifold plane at finite $`A(r)`$ on which the five-dimensional spacetime ends. Again the examples are type I and Horava-Witten theory.
The UV-IR relation suggests that curvature singularities that arise as $`A(r)\mathrm{}`$ represent non-trivial infrared physics in the dual field theory. (Recall that $`A(r)=\frac{1}{2}\mathrm{log}|g_{tt}|`$ is roughly the log of the energy scale). This interpretation is natural from the point of view of the holographic renormalization group . Suppose the field theory vacuum is one which can support finite temperature: by this we mean that the vacuum can be approached smoothly in a $`T0`$ limit, with no zero temperature phase transition. Then a putative dual description in supergravity should also be able to support finite temperature, in the form of a black hole horizon which “cloaks” the curvature singularity. As the temperature is lowered to zero, the horizon should retreat until the original singular geometry is recovered. Thus we arrive at a weak form of Cosmic Censorship: the only curvature singularities allowed in geometries of the form (1) are those which can be obtained as limits of regular black holes.
Even this statement is too strong: there are interesting vacua in field theory which cannot support finite temperature. Examples include states on the Coulomb branch of $`𝒩=4`$ gauge theory. The dual geometries are typically singular in five dimensions. Any finite temperature, no matter how small, would draw a Coulomb branch state to the origin of moduli space, where conformal invariance is recovered in the zero temperature limit. Correspondingly, finite temperature in the bulk should pull the dual singular geometries back to AdS-Schwarzschild.
As detectives in search of a useful rule for distinguishing good singularities from bad, let us then ask the following two questions. (1) Do the singular five-dimensional geometries corresponding to Coulomb branch states share any common features with geometries which are limits of regular black holes? (2) Can we find other clear examples which could help us establish the identifying marks of healthy singularities as opposed to pathological ones? The main goal of this paper is to answer both questions affirmatively, and in fact to argue that
| Large curvatures in geometries of the form (1) are allowed only if |
| --- |
| the scalar potential is bounded above in the solution. |
(2)
This conjecture is peculiar in that it is the opposite of the naive intuition from classical relativity. The dominant energy condition, for example, constrains the scalar potential to be non-negative—not just in solutions, but in the lagrangian. (This of course rules out an $`AdS_5`$ solution, since the scalar potential is the cosmological constant). Nevertheless, it will be proven that (2) is indeed a necessary condition for a bulk geometry to support finite temperature in the form of a black hole horizon. It will also be shown that (2) correctly predicts some non-trivial aspects of the vacuum structure of mass-deformed $`𝒩=4`$ gauge theory.
It is important to bear in mind that (2) is a constraint on solutions, not on the five-dimensional lagrangian. In point of fact, the scalar potential for $`d=5`$ $`𝒩=8`$ gauged supergravity (from which all our examples will be drawn) is unbounded both above and below. There are solutions where scalars diverge near the singularity in such a way that the scalar potential goes to $`+\mathrm{}`$. These are bad, according to (2). There are others where the scalar potential goes to $`\mathrm{}`$. These are good. Thanks to AdS/CFT, we can not only demonstrate that the good singularities have sensible field theory duals, but in most cases even trace the sickness of the bad singularities back to some definite pathology in the dual field theory.
A complementary line of attack, which will not be pursued in this work but is potentially very fruitful, is to resolve singularities in metrics of the form (1) by lifting them to ten dimensions and finding a configuration of branes that produces the ten-dimensional geometry (possibly in some near-horizon, large $`N`$ limit). This is in general a difficult task, but when it can be carried out, it provides a satisfying account of how string theory resolves a singularity. A goal of the present paper is to try to anticipate which five-dimensional singularities will admit a brane resolution. The methods should be equally applicable to Calabi-Yau compactifications of Horava-Witten theory where one end-of-the-universe brane has retreated to infinite redshift. This is a case where our understanding of the microscopic eleven-dimensional theory is so incomplete that supergravity methods, combined with consistency conditions such as anomaly inflow, may be the only tools available.
One may legitimately inquire why any method based on five-dimensional supergravity should have predictive power regarding the nature of singular solutions. This is where AdS/CFT helps: confusing singularities on the AdS side are often transparent in the field theory dual. In the end, however, the hope is that (2) is a robust result that will survive beyond the approximations in which we have studied it, and even beyond AdS/CFT. The nature of a local singularity should be independent of the geometry far from it, so having an asymptotically $`AdS_5`$ region may be just a crutch for calculation. The physics of the naked singularity should be as independent of the asymptotically $`AdS_5`$ region as infrared effective theories are of their detailed microscopic origins.
The organization of the paper is as follows. Section 2 is a more technical introduction in which conventions are set and a method for generating solutions is reviewed. In section 3 we discuss the motivations for (2), and we explain some related conditions which are useful when looking at examples. We also discuss the infrared asymptotics of generic solutions and make some speculations regarding phase transitions. In section 4 we review three examples of asymptotically $`AdS_5`$ geometries which arise as deformations of $`𝒩=4`$ super-Yang-Mills theory by a relevant operator. In section 5 we consider linearized fluctuations around a general asymptotically $`AdS_5`$ background and propose a generalization of the Breitenlohner-Freedman bound. Concluding remarks and some further conjectures are presented in section 6. In the appendix we remark on brane-world scenarios in which there is a curvature singularity parallel to our world, at some finite proper distance from it.
## 2 Background
All the geometries we will consider are classical solutions to the equations of motion following from an action
$$S=d^5x\sqrt{g}\left[\frac{1}{4}R\frac{1}{2}(\stackrel{}{\phi })^2V(\stackrel{}{\phi })\right].$$
(3)
Usually it will be assumed that the potential $`V(\stackrel{}{\phi })`$ has a local maximum at $`\stackrel{}{\phi }=0`$. Our main example of such an action is the gravity-plus-scalars sector of $`d=5`$ $`𝒩=8`$ gauged supergravity . Here the $`42`$ scalars $`\stackrel{}{\phi }`$ parametrize the coset $`E_{6(6)}/USp(8)`$, and $`V(\stackrel{}{\phi })`$ is a $`SO(6)\times SL(2,𝐑)`$-invariant function on this coset. The kinetic term would be more accurately represented as $`G_{IJ}(\stackrel{}{\phi })\phi ^I\phi ^J`$, where $`G_{IJ}(\stackrel{}{\phi })`$ is the natural sigma model metric for the coset. In the general discussion that follows, we will suppress factors of $`G_{IJ}`$ and its inverse. Restoring them does not make any difference to the essential points of the story as long as $`G_{IJ}`$ is a smooth function of $`\stackrel{}{\phi }`$, which is true for $`d=5`$ $`𝒩=8`$ gauged supergravity. In the examples of section 4 we will work on submanifolds of $`E_{6(6)}`$ where the metric is flat, so one may find a basis of scalars where $`G_{IJ}=\delta _{IJ}`$. We will call such scalars canonically normalized.
The $`40`$ scalars that participate in $`V(\stackrel{}{\phi })`$ are tachyonic at the maximally supersymmetric point representing unperturbed $`𝒩=4`$ super-Yang-Mills theory. In AdS/CFT they correspond to dimension $`2`$ operators in the $`\mathrm{𝟐𝟎}^{}`$ of $`SO(6)`$ and dimension $`3`$ operators in the $`\mathrm{𝟏𝟎}+\overline{\mathrm{𝟏𝟎}}`$. The remaining scalars are the exactly massless dilaton and axion, corresponding to the exactly marginal complex operator $`tr(F^2+iF\stackrel{~}{F}+\mathrm{})`$. These massless scalars will not participate in the geometries we consider. Our remarks can be applied more generally to less supersymmetric AdS/CFT duals, the key point being that the scalars in question are tachyonic at $`\stackrel{}{\phi }=0`$.
It is often useful to take a potential of the form
$$V(\stackrel{}{\phi })=\frac{1}{8}\left(\frac{W}{\stackrel{}{\phi }}\right)^2\frac{1}{3}W(\stackrel{}{\phi })^2,$$
(4)
where $`W`$ is some appropriate function of the scalars. In $`U(1)`$ gaugings of $`d=5`$ $`𝒩=2`$ supergravity, $`W`$ can be identified with the superpotential. In it was shown that the potential of $`d=5`$ $`𝒩=8`$ gauged supergravity also assumes the form (4) on restricted subspaces of the scalar manifold, with $`W`$ identified as an eigenvalue of an $`SO(6)\times SL(2,𝐑)`$-invariant $`USp(8)`$ matrix $`W_{ab}`$, defined on $`E_{6(6)}/USp(8)`$, which serves as a generalized superpotential in the sense that the gravitino variation has the form $`\delta \psi _{\mu a}=𝒟_\mu ϵ_a\frac{1}{6}gW_{ab}\gamma _\mu ϵ^b`$ (see for further details).
In , the form (4) was found to be useful, independent of supersymmetry, for generating solutions of the form (1) and investigating their properties. Indeed, given $`V(\stackrel{}{\phi })`$, the equation (4) can be regarded as a PDE for $`W`$, to be solved with arbitrary but specified initial conditions. It was observed in that the $`W`$ so obtained is related to Hamilton’s principle function. Of the possible forms for $`W`$ satisfying (4), only one (or at most discretely many) can be related to the supersymmetry transformations. However, if (4) can be solved systematically, then it is possible to generate any solution of the form (1) (supersymmetric or not) to the equations of motion by solving the first order equations
$$\stackrel{}{\phi }^{}=\frac{1}{2}\frac{W}{\stackrel{}{\phi }}A^{}=\frac{1}{3}W(\stackrel{}{\phi }).$$
(5)
Except in section 5, primes will indicate derivatives with respect to $`r`$. It can happen that $`W`$ is not a single-valued function of $`\stackrel{}{\phi }`$, but rather a multi-sheeted graph in the allowed region of $`W`$$`\stackrel{}{\phi }`$ space where $`V(\stackrel{}{\phi })+\frac{1}{3}W(\stackrel{}{\phi })^20`$. In such a case, (5) must be supplemented with selection rules for going from one branch to another when the boundary of the allowed region is reached (see for a more detailed discussion). This subtlety will not be a concern in this paper.
The prima facia evidence that the solution space of (4) together with (5) coincides with the solution space of the equations of motion following from (3) is that the number of integration constants is the same. In view of the realization that (5) can be recovered in Hamilton-Jacobi theory, there can be no question that (4) together with (5) are nothing more nor less than a fancy rewriting of the equations of motion. Nevertheless, it will be instructive to follow through the counting of integration constants. Suppose there are $`n`$ scalars. The equations of motion consist of $`n`$ second order differential equations, one for each scalar, plus one first order constraint from the $`G^r_r`$ Einstein equation. That means $`2n+1`$ integration constants. In the first order formalism, $`n`$ of these integrations constants are tucked into $`W`$ as constants of the motion; $`n`$ more come from the first order equations for the scalars; and the last comes from the equation for $`A^{}`$ in (5), which is equivalent to the $`G^r_r`$ Einstein equation if the first order equations for the scalars hold.
There are various reasons to think that the space of solutions of the form (1) to the equations of motion following from (3) is too big—that not all solutions are physical. First, from a field theory perspective, it is possible that the vacuum is ill-defined. This would happen, for instance, if the effective potential in the field theory were unbounded below. It was suggested in that exactly this problem was responsible for the unphysical nature of the $`SO(5)`$ critical point. To flow to the $`SO(5)`$ point, it is necessary to deform the $`𝒩=4`$ lagrangian by a term $`m^2tr(X_1^2X_2^2X_3^2X_4^2X_5^2+5X_6^2)`$. This means the effective potential $`V_{\mathrm{eff}}`$ in the field theory is unbounded below. This pathology manifests itself in five-dimensional gauged supergravity as a violation of the Breitenlohner-Freedman bound at the $`SO(5)`$ critical point. Another pathology which one could imagine in field theory is attempting to assign a negative vacuum expectation value (VEV) to a positive definite operator. In section 4 we will encounter examples where this problem arises.
Generic solutions of the form (1) to the equations of motion following from (3) have curvature singularities at finite $`r`$. A very few, such as $`AdS_5`$ itself and the RG flows of , have curvatures everywhere bounded. We believe we are on safe ground when we assert that not all naked singularities are physical, even if we demand $`3+1`$-dimensional Poincaré invariance and/or some unbroken supersymmetry. Indeed, some of the Coulomb branch flows explored in have an unphysical property when they were lifted to ten dimensions: the brane distributions involve branes both of positive tension and positive charge, and of negative tension and negative charge. A similar problem was beautifully resolved in , but for Coulomb branch states it is obscure what enhancement of symmetry could play a relevant role. In section 4 we will encounter other pathologies in geometries which violate (2). In short, it seems hard to make a convincing case that all the nakedly singular, Poinaré-invariant geometries that we can generate as solutions to the equations of motion are physical.
We would however take the position that there are physically allowable naked singularities in AdS/CFT, representing interesting infrared physics in the dual field theory—only supergravity has limited calculation power in their vicinity. The repulson/enhançon of is one example; the dual to $`𝒩=1`$ super-Yang-Mills theory is another, which we will have more to say about later. In general, it is a non-trivial exercise to lift five-dimensional solutions to ten dimensions and find a brane configuration which realizes the geometry in some appropriate limit. And it is painful to proceed example-by-example, looking for known pathologies, like negative VEV’s for positive definite operators, or imaginary ten-dimensional metrics. Thus it is useful to have a criterion such as (2), which is easy to apply and seems fairly reliable.
## 3 Motivating the criterion
In this section we will provide two motivations for (2). The first is that it is a necessary condition for the existence of near-extremal generalizations of a solution to (5). We will prove this assertion in section 3.1. By near-extremal generalizations we mean a family of black hole solutions, with horizons hiding the naked singularity, which converge to the Poincaré invariant solution in an appropriate topology. The Hawking temperature of the horizon is identified with finite temperature in the dual field theory.
Finite temperature in field theory serves as an infrared cutoff in the sense that it masks physics at scales lower than the temperature. A naked singularity that can be hidden behind a black hole horizon is a signal of non-trivial but sensible infrared physics in the dual field theory. A naked singularity that cannot be so hidden may indicate a pathology—the absence of a well-defined dual field theory, or a field theory in an unphysical vacuum state. If we approach a “good” naked singularity by taking a limit of regular black holes, the dual picture is that the infrared cutoff (finite temperature) is being removed. The calculational power of AdS/CFT may be limited in such a circumstance because of the large curvatures near the singularity.
The second motivation for (2) is an examination of states on the Coulomb branch of $`𝒩=4`$ super-Yang-Mills theory. Finite temperature draws Coulomb branch states to the origin of moduli space, where all operators have vanishing vacuum expectation values. Thus if one can demonstrate (2) for Coulomb branch states, it is a check of the conjecture which is orthogonal to finite temperature considerations. In fact there is a class of Coulomb branch vacua which can be described in $`d=5`$ $`𝒩=8`$ gauged supergravity, and most of them have naked singularities. These naked singularities can be regarded as artifacts of the Kaluza-Klein reduction from ten to five dimensions provided the ten-dimensional geometry avoids the aforementioned pathology of D3-branes with negative tension and negative charge. As we will show in section 3.2, the pathological cases for Coulomb branch flows are precisely those where (2) is violated.
### 3.1 Finite Temperature
The finite temperature generalization of the ansatz (1) is
| $`ds^2`$ | $`=e^{2A(r)}(h(r)dt^2+d\stackrel{}{x}^2)+{\displaystyle \frac{dr^2}{h(r)}}`$ |
| --- | --- |
| $`\stackrel{}{\phi }`$ | $`=\stackrel{}{\phi }(r).`$ |
(6)
Plugging (6) into the equations of motion following from (3), we obtain the following:
| $`\text{ }\text{ }\text{ }\text{ }\stackrel{}{\phi }`$ | $`=h(\stackrel{}{\phi }^{\prime \prime }+4A^{}\stackrel{}{\phi }^{})+h^{}\stackrel{}{\phi }^{}={\displaystyle \frac{V}{\stackrel{}{\phi }}}`$ |
| --- | --- |
| $`G^t{}_{t}{}^{}G^x_x`$ | $`={\displaystyle \frac{1}{2}}(h^{\prime \prime }+4A^{}h^{})=0`$ |
| $`G^t{}_{t}{}^{}G^r_r`$ | $`=3hA^{\prime \prime }=2h\stackrel{}{\phi }^2`$ |
(7)
and the “zero energy” constraint
$$G^r{}_{r}{}^{}=\frac{3}{2}A^{}(h^{}+4A^{}h)=h\stackrel{}{\phi }^22V(\stackrel{}{\phi }).$$
(8)
Differentiating (8) with respect to $`r`$ yields an equation which follows algebraically from (7), as $`r`$-reparametrization invariance demands.
We can solve for $`A(r)`$ and $`h(r)`$ in terms of $`\stackrel{}{\phi }^{}`$ using only quadratures:
| $`A(r)`$ | $`={\displaystyle \frac{r}{L}}{\displaystyle _r^{\mathrm{}}}𝑑r_1{\displaystyle _{r_1}^{\mathrm{}}}𝑑r_2{\displaystyle \frac{2}{3}}\stackrel{}{\phi }^{}(r_2)^2`$ |
| --- | --- |
| $`h(r)`$ | $`=1{\displaystyle _r^{\mathrm{}}}𝑑r_1{\displaystyle \frac{4B}{L}}e^{4A(r_1)}.`$ |
(9)
In (9) we have specified some boundary conditions at $`r\mathrm{}`$. Choosing $`lim_r\mathrm{}h(r)=1`$ near the boundary removes the freedom to rescale $`r`$ multiplicatively. The zero energy constraint fixes $`lim_r\mathrm{}\left(h(r)A^{}(r)^2+\frac{1}{3}V(\stackrel{}{\phi }(r))\right)=0`$, so if one defines $`L`$ through the equation $`lim_r\mathrm{}V(\stackrel{}{\phi }(r))=3/L^2`$, the normalization of the linear term in $`A(r)`$ in (9) is fixed. We are free to add a constant to $`A(r)`$, or to shift $`r`$ itself by a constant. One of these freedoms can be absorbed by making a uniform dilation on $`t`$ and $`\stackrel{}{x}`$; the other amounts to a conformal transformation which scales all dimensionful quantities uniformly according to their dimensions. The only free integration constant in (9) with physical meaning is $`B`$, which is essentially a non-extremality parameter. If $`B`$ is negative, $`h(r)`$ would never be zero, and we would get solutions which are variants of negative mass AdS-Schwarzschild. These we regard as manifestly unphysical.
Suppose now that $`B>0`$. The function $`h(r)`$ can have at most one zero, because it is monotonic. If a zero exists, it indicates the presence of a black hole horizon. Its location, $`r=r_H`$, must satisfy $`r_H\frac{L}{4}\mathrm{log}B`$, with equality if and only if no scalars are excited (that is, for AdS-Schwarzschild only). The Hawking temperature and Bekenstein-Hawking entropy (per unit volume in the $`\stackrel{}{x}`$-directions) are
| $`T`$ | $`={\displaystyle \frac{e^{A(r_H)}h^{}(r_H)}{4\pi }}={\displaystyle \frac{B}{\pi L}}e^{3A(r_H)}`$ |
| --- | --- |
| $`{\displaystyle \frac{S}{V}}`$ | $`=\pi e^{3A(r_H)}.`$ |
(10)
In a Euclidean version of (6), with Euclidean time $`\tau =it`$, the good coordinates at $`r=r_H`$ are
$$y=2\sqrt{\frac{rr_H}{h^{}(r_H)}}\theta =e^{A(r_H)}\frac{h^{}(r_H)}{2}\tau .$$
(11)
In terms of these variables, the $`\tau `$$`r`$ part of the metric near $`r=r_H`$ is just $`dy^2+y^2d\theta ^2`$. Taking the $`h0`$ limit of the scalar equation of motion, one arrives at the following horizon boundary conditions $`r=r_H`$:
$$h^{}\stackrel{}{\phi }^{}=\frac{V}{\stackrel{}{\phi }}.$$
(12)
The other equations of motion do not lead to any further boundary conditions in the sense of specifying a well-defined boundary value problem. The condition (12) is actually more informative than the obvious requirement that $`\stackrel{}{\phi }`$ has to be smooth in the good variables (11). In those variables, (12) becomes $`\frac{2}{y}\frac{d\stackrel{}{\phi }}{dy}=\frac{V}{\stackrel{}{\phi }}`$.
Existence of solutions of the form (6) with regular black hole horizons does not contradict no-hair theorems. The reason is that the scalars $`\stackrel{}{\phi }`$ are by assumption tachyons in the asymptotically $`AdS_5`$ region. For a particular scalar (call it $`\phi `$) the linearized wave equation near the boundary of $`AdS_5`$ has two solutions (call them $`\phi _{(1)}`$ and $`\phi _{(2)}`$). Both of them decay as one approaches the boundary. The black hole horizon boundary conditions, (12), lead to $`n`$ constraints if there are $`n`$ scalars. But there are $`2n`$ integration constants for the scalars, leaving us with the expectation of an $`n`$ parameter family of solutions.
In this section we will obtain necessary conditions for the existence of black hole solutions of the form (6). In particular, given a Poincaré invariant solution which violates (2), we will show that there does not exist a family of black hole solutions which can be made to approximate the Poincaré invariant solution arbitrarily closely. This amounts to the statement that the Poincaé invariant solution does not have near-extremal generalizations. In fact, more stringent conditions than (2) are probably necessary in order for a singular geometry to have near-extremal generalizations. Conjectures regarding what form those conditions might take will be presented in sections 3.3 and 6.
Suppose one starts with a Poincaré invariant, singular solution to (5), which can be proven not to have near-extremal generalizations in five-dimensional classical supergravity. We wish to conclude that the dual field theory cannot support finite temperature. This is a tricky claim to make, because the curvature singularity interferes with our ability to compute reliably in classical supergravity. Two points in favor of it are as follows. 1) Assuming there is a “microscopic” description of finite temperature at a singularity which applies beyond the applicability of supergravity, there should be a smooth match from the microscopic description to the black hole horizon description at some crossover point . 2) Numerical studies of two examples for which the field theory dual is reasonably well-understood indeed exhibit a smooth limit in which black hole solutions converge to a Poincaré invariant solution. We will briefly report on the numerical work in section 6, where there is also a further discussion of the potential pitfalls of identifying solutions that can support finite temperature with those which admit near-extremal generalizations. For now we will content ourselves with an analysis of what the existence of a static black hole horizon implies.
Let us start with a regular black hole solution of the form (6). For the moment let us not assume that the solution is asymptotically $`AdS_5`$. Evaluating the zero-energy constraint (8) at the horizon ($`h=0`$), we see that $`A^{}h^{}=\frac{4}{3}V(\stackrel{}{\phi })`$. This implies $`V(\stackrel{}{\phi }(r_H))<0`$. The reason is that both $`A^{}`$ and $`h^{}`$ are everywhere positive. In the case of $`A^{}`$, this follows from the property $`A^{\prime \prime }0`$ (obvious from the third line of (7)), and the assumption that $`A^{}(r)`$ is positive somewhere outside the horizon. In the case of $`h^{}`$, it follows once one has ruled out the solutions akin to negative mass Schwarzschild: $`e^{4A}h^{}`$ is a constant, and to avoid negative mass it has to be a positive constant. This derivation scarcely relies on the matter lagrangian: $`A^{\prime \prime }0`$ is in fact implied by the null energy condition, just as in the extremal solution . The conclusion, then, from the zero energy constraint, is that the bulk cosmological constant cannot be positive at the black hole horizon. The only way to saturate this bound is to take $`A^{}`$ and $`h^{\prime \prime }`$ identically $`0`$. This gives a rather trivial solution to the vacuum Einstein equations with zero cosmological constant.
It is in fact possible to improve on the bound $`V(\stackrel{}{\phi }(r_H))<0`$ if one is willing to use properties of the matter lagrangian and the assumption that the spacetime is asymptotically $`AdS_5`$. Note that
| $`{\displaystyle \frac{d}{dr}}\left(V(\stackrel{}{\phi }(r)){\displaystyle \frac{1}{2}}h(r)\stackrel{}{\phi }^{}(r)^2\right)`$ | $`={\displaystyle \frac{V}{\stackrel{}{\phi }}}\stackrel{}{\phi }^{}{\displaystyle \frac{1}{2}}h^{}\stackrel{}{\phi }^2h\stackrel{}{\phi }^{}\stackrel{}{\phi }^{\prime \prime }`$ |
| --- | --- |
| | $`={\displaystyle \frac{1}{2}}(h^{}+8A^{}h)\stackrel{}{\phi }^20\text{if }h0\text{,}`$ |
(13)
where to obtain the second equality we have used the scalar equation of motion. At the black hole horizon, $`h=0`$ and $`\stackrel{}{\phi }^{}`$ is finite. At the boundary of $`AdS_5`$, $`h1`$ and $`\stackrel{}{\phi }^{}0`$. So $`h\stackrel{}{\phi }^2=0`$ both at the black hole boundary and at the horizon of $`AdS_5`$. Integrating (13), we find
$$V(\stackrel{}{\phi }(r_H))V_{UV},$$
(14)
where $`V_{UV}`$ is the value of $`V(\stackrel{}{\phi })`$ at the ultraviolet fixed point—usually at $`\stackrel{}{\phi }=0`$. The inequality is saturated only for AdS-Schwarzschild. A crucial ingredient in obtaining these inequalities is that the scalar kinetic term is positive definite. This is implied by the null energy condition. The same positivity properties were the basis of the proofs of the c-theorem in .
In view of (9), we can specify a solution to (7) completely through the pair $`(\stackrel{}{\phi }(r),B)`$. Suppose we have a sequence of regular black hole solutions, specified as pairs $`(\stackrel{}{\phi }_n(r),B_n)`$, with horizon radii $`r_n`$, which converges uniformly to a Poincaré invariant solution $`\stackrel{}{\phi }_0(r)`$ to (5). More precisely, the convergence criteria are $`B_n0`$ and $`sup_{rr_n}|\stackrel{}{\phi }_n(r)\stackrel{}{\phi }_0(r)|0`$ as $`n\mathrm{}`$. Since $`B_n0`$, the horizon radii $`r_n`$ will retreat as $`n\mathrm{}`$ to locations with more and more negative $`A(r)`$—either to the naked singularity of the Poincaré invariant solution, or to $`r=\mathrm{}`$ if there is no naked singularity.<sup>2</sup><sup>2</sup>2A naked singularity has $`A(r)\mathrm{}`$ at finite $`r`$. If $`A(r)\mathrm{}`$ only as $`r\mathrm{}`$, curvatures could still grow arbitrarily large. Either way, the geometry “ends” when $`A\mathrm{}`$, modulo issues of analytic continuation, which can be largely disregarded if we think of Wick rotating with respect to Poincaré time rather than some modification of global $`AdS_5`$ time. With minimal assumptions of smoothness on $`\stackrel{}{\phi }_0(r)`$ and $`V(\stackrel{}{\phi })`$, we can deduce from (14) that
$$V(\stackrel{}{\phi }_0(r))V_{UV}\text{in the infrared.}$$
(15)
This is now a constraint on the extremal solution.
By “in the infrared,” we mean that the inequality should hold for all values of $`r`$ such that $`A(r)`$ is sufficiently negative, or at least asymptotically as $`A\mathrm{}`$. As before, $`V_{UV}`$ is the limit of $`V(\stackrel{}{\phi }(r))`$ as we approach the boundary of $`AdS_5`$. The condition (15) implies (2), given that $`V`$ is a continuous function of $`\stackrel{}{\phi }`$. It is almost true that (2) implies (15): the simplest way in which one could imagine (2) holding and (15) failing is for $`\stackrel{}{\phi }`$ to approach some constant value in the infrared, corresponding to an infrared fixed point in the dual field theory, such that $`V_{IR}>V_{UV}`$—but precisely this possibility is ruled out by the c-theorem of . We will encounter several examples which violate both (2) and (15) by having $`V(\stackrel{}{\phi }(r))+\mathrm{}`$ in the infrared. All of these examples involve naked singularities, and these particular naked singularities seem very unphysical.
Before turning to Coulomb branch flows, a remark on different topologies on the space of field configurations is perhaps in order. Because $`A(r)`$ and $`h(r)`$ can be expressed in terms of $`\stackrel{}{\phi }(r)`$ and $`B`$ through quadratures, it seems natural to define a norm in terms of $`\stackrel{}{\phi }`$ and $`B`$ alone. Given two configurations, $`(\stackrel{}{\phi },B)`$ and $`(\stackrel{~}{\stackrel{}{\phi }},\stackrel{~}{B})`$, a natural Sobolev norm on the difference $`\stackrel{}{\phi }\stackrel{~}{\stackrel{}{\phi }}`$ is
$$\stackrel{}{\phi }\stackrel{~}{\stackrel{}{\phi }}_{H_s^p}=\underset{k=0}{\overset{s}{}}\stackrel{}{\phi }^{(k)}\stackrel{~}{\stackrel{}{\phi }}^{(k)}_{L^p}=\underset{k=0}{\overset{s}{}}\left(_{r_H}^{\mathrm{}}𝑑re^{4\stackrel{~}{A}(r)}\left|\stackrel{}{\phi }^{(k)}\stackrel{~}{\stackrel{}{\phi }}^{(k)}\right|^p\right)^{1/p}$$
(16)
for integer $`s`$. Here $`\stackrel{}{\phi }^{(k)}`$ indicates the $`k`$-th derivative with respect to $`r`$, and we have assumed that the horizon radius $`r_H`$ for the configuration $`(\stackrel{}{\phi },B)`$ (defined as the location where $`h=0`$) is larger than the horizon radius $`\stackrel{~}{r}_H`$ for $`(\stackrel{~}{\stackrel{}{\phi }},\stackrel{~}{B})`$. We have chosen the measure $`dre^{4\stackrel{~}{A}(r)}`$ because it descends naturally from the five-dimensional volume element $`d^5x\sqrt{g}`$. Other measures could be considered if they seem more convenient. For $`C^{\mathrm{}}`$ functions, the $`H_s^{\mathrm{}}`$ norm is particularly simple:
$$\stackrel{}{\phi }\stackrel{~}{\stackrel{}{\phi }}_{H_s^p}=\underset{k=0}{\overset{s}{}}\underset{rr_H}{sup}\left|\stackrel{}{\phi }^{(k)}(r)\stackrel{~}{\stackrel{}{\phi }}^{(k)}(r)\right|.$$
(17)
Thus the topology of uniform convergence, which we used to derive (15), follows from the $`H_0^{\mathrm{}}`$ norm, more commonly known as the $`L^{\mathrm{}}`$ norm. $`H_s^{\mathrm{}}`$ norms are very permissive near the boundary of $`AdS_5`$: any relevant deformation of the lagrangian has finite norm. The $`H_s^2`$ norms rule out such deformations for $`𝒩=4`$ gauge theory, but also they rule out VEV’s for dimension $`2`$ operators. It is easy to construct norms which rule out the $`re^{2r/L}`$ behavior of a dimension $`2`$ deformation but permit the $`e^{2r/L}`$ behavior of a dimension $`2`$ VEV. An example is
$$\stackrel{}{\phi }\stackrel{~}{\stackrel{}{\phi }}_f=_2^{\mathrm{}}𝑑pf(p)\stackrel{}{\phi }\stackrel{~}{\stackrel{}{\phi }}_{H_s^p}$$
(18)
where $`f(p)`$ is a function which is $`1`$ near $`p=2`$ and whose integral over $`(2,\mathrm{})`$ is finite. It would be interesting to see if the topology induced by such a norm could be used in place of the topology of uniform convergence for deriving (15) as a necessary condition.
### 3.2 The Coulomb branch of $`𝒩=4`$ super-Yang-Mills
Although it is nice that one can start with an arbitrary black hole solution and prove definite statements like (14) that support the view that the scalar potential must remain bounded above near a curvature singularity, one may still retain some doubt about the validity of (2). In particular, might it not be that singularities which violate (2) simply correspond to vacua in the dual field theory which cannot support finite temperature? In this section, we will argue against this possibility by considering explicit examples of such vacua where the dual geometries are well understood, singular, and in line with (2).
Next to the maximally supersymmetric $`AdS_5`$ geometry, the best-understood solutions to $`d=5`$ $`𝒩=8`$ gauged supergravity with $`3+1`$-dimensional Poincaré invariance are those which involve only the metric and the scalars in the coset $`SL(6,𝐑)/SO(6)`$ . The CFT duals (when they are well-defined) are states on the Coulomb branch of $`𝒩=4`$ super-Yang-Mills which can be parametrized by the VEV’s of operators $`trX_{(I}X_{J)}`$. These operators, transforming in the $`\mathrm{𝟐𝟎}^{}`$ of $`SO(6)`$, are precisely the ones dual to the $`20`$ scalars parametrizing $`SL(6,𝐑)/SO(6)`$.
There are VEV’s for higher dimension operators as well as for those in the $`\mathrm{𝟐𝟎}^{}`$, and one can calculate them systematically from the ten-dimensional geometry. The magic of consistent truncation is that the profiles of only finitely many scalars capture infinitely many VEV’s. Roughly speaking, these higher dimension VEV’s are sourced by arbitrary powers of the scalars in the $`\mathrm{𝟐𝟎}^{}`$. The corresponding statement in consistent truncation is that the ten-dimensional fields are determined non-linearly in terms of the five-dimensional fields.
The ten-dimensional geometry is the near-horizon limit of the background of a continuous distribution of D3-branes:
| $`ds^2`$ | $`={\displaystyle \frac{1}{\sqrt{H}}}\left(dt^2+d\stackrel{}{x}^2\right)+\sqrt{H}{\displaystyle \underset{I=1}{\overset{6}{}}}dy_I^2`$ |
| --- | --- |
| $`H`$ | $`={\displaystyle d^6\xi \sigma (\stackrel{}{\xi })\frac{L^4}{|\stackrel{}{y}\stackrel{}{\xi }|^4}},`$ |
(19)
where the distribution $`\sigma `$ integrates to $`1`$. All the solutions whose Kaluza-Klein reductions involve only the five-dimensional metric and the scalars in $`SL(6,𝐑)/SO(6)`$ arise from distributions which are limits of either of the following the following two distributions :
| $`\sigma _6(\stackrel{}{\xi })`$ | $`={\displaystyle \frac{1}{\pi ^3\mathrm{}_1\mathrm{}\mathrm{}_6}}\delta ^{}\left(1{\displaystyle \underset{I=1}{\overset{6}{}}}{\displaystyle \frac{\xi _I^2}{\mathrm{}_I^2}}\right)`$ |
| --- | --- |
| $`\sigma _5(\stackrel{}{\xi })`$ | $`={\displaystyle \frac{1}{2\pi ^3\mathrm{}_1\mathrm{}\mathrm{}_5}}[(1{\displaystyle \underset{I=1}{\overset{5}{}}}{\displaystyle \frac{\xi _I^2}{\mathrm{}_I^2}})^{3/2}\mathrm{\Theta }(1{\displaystyle \underset{I=1}{\overset{5}{}}}{\displaystyle \frac{\xi _I^2}{\mathrm{}_I^2}})`$ |
| | $`+2(1{\displaystyle \underset{I=1}{\overset{5}{}}}{\displaystyle \frac{\xi _I^2}{\mathrm{}_I^2}})^{1/2}\delta (1{\displaystyle \underset{I=1}{\overset{5}{}}}{\displaystyle \frac{\xi _I^2}{\mathrm{}_I^2}})]\delta (\xi _6).`$ |
(20)
The distribution $`\sigma _5(\stackrel{}{\xi })`$ is the limit of $`\sigma _6(\stackrel{}{\xi })`$ as $`\mathrm{}_60`$. There is a sort of shell theorem to the effect that shifting all the $`\mathrm{}_I^2`$ by a constant makes no difference to the geometry away from the branes (see for evidence of this theorem, and for a crisp statement). Both $`\sigma _5`$ and $`\sigma _6`$ fail to be everywhere positive. This is the signal that there are “ghost” D3-branes of negative tension and negative charge, as well as normal D3-branes with positive tension and positive charge. In the absence of an understanding of the geometries involving “ghosts” along the lines of , we are inclined to regard them as unphysical. If we send $`\mathrm{}_50`$ in $`\sigma _5`$, the resulting distribution is a delta function shell in the shape of an ellipsoid:
$$\sigma _4(\stackrel{}{\xi })=\frac{1}{\pi ^2\mathrm{}_1\mathrm{}\mathrm{}_4}\delta \left(1\underset{I=1}{\overset{4}{}}\frac{\xi _I^2}{\mathrm{}_I^2}\right)\delta (\xi _5)\delta (\xi _6).$$
(21)
This is a positive definite distribution, and so are its limits with one or more of the remaining $`\mathrm{}_I`$ taken to zero.
Given a distribution $`\sigma (\stackrel{}{\xi })`$, vacuum expectation values of gauge singlet operators can be computed by the rule
$$trX_{(I_1}\mathrm{}X_{I_{\mathrm{}})}=d^6\xi \sigma (\stackrel{}{\xi })\xi _{I_1}\mathrm{}\xi _I_{\mathrm{}}.$$
(22)
Conversely, a distribution of the form (20) (or any limit thereof) can be reconstructed from the VEV’s of the $`SO(6)`$ irrep operators $`trX_{(I_1}\mathrm{}X_{I_{\mathrm{}})}`$, up to an ambiguity which amounts to shifting all the $`\mathrm{}_I^2`$ by a constant. Thus the requirement that the distribution can be made positive definite (by appropriate choice of this constant shift) must translate into a series of inequalities among the VEV’s for $`trX_{(I_1}\mathrm{}X_{I_{\mathrm{}})}`$ which follow from the hermiticity of the individual $`X_I`$. The point of all this is that the “ghost” D3-branes are in this case really unphysical: no local resolution of the geometry which left the far field form unchanged would change the problematic VEV’s.
Let us first take $`n`$ of the $`\mathrm{}_I`$ equal in (20), and set the other $`6n`$ to zero. These cases were studied in detail in . The $`SO(6)`$ global symmetry is broken to $`SO(n)\times SO(6n)`$. The five-dimensional geometry involves only one scalar. More precisely, the trajectory on the scalar manifold is a geodesic with respect to the sigma-model metric, and we can parametrize it with a canonically normalized scalar, $`\mu `$. The operator dual to $`\mu `$ is $`𝒪_2=tr\left[(6n)_{In}X_I^2n_{I>n}X_I^2\right]`$. Near the boundary of $`AdS_5`$, we have $`\mu e^{2r/L}`$, corresponding to a positive VEV for $`𝒪_2`$. Far from the boundary, $`\mu `$ increases without bound. For large $`\mu `$, one finds $`W(\mu )e^{\zeta \mu }`$. The constant $`\zeta `$ has values $`2/\sqrt{15}`$, $`\sqrt{2/3}`$, $`\sqrt{4/3}`$, $`\sqrt{8/3}`$, and $`10/\sqrt{15}`$ for $`n=1`$, $`2`$, $`3`$, $`4`$, and $`5`$ respectively (see (15) of ). The geometric origin of these numbers will become clear shortly.
For $`n4`$ one finds $`V(\mu )\mathrm{}`$ as $`\mu \mathrm{}`$, whereas for $`n=5`$ one finds $`V(\mu )+\mathrm{}`$. Thus (2) rules out the pathological $`n=5`$ case but accepts the physical $`n=4`$ cases. This success, which seems to have nothing to do with finite temperature, is our second motivation for proposing (2) as a good criterion for distinguishing physical from unphysical geometries.
It is fairly straightforward to work out the cases with unequal $`\mathrm{}_I`$. Five scalars are involved in general, corresponding to the five independent eigenvalues of the matrix $`M=SS^T`$, where $`SSL(6,𝐑)/SO(6)`$. In an appropriate basis, $`M=diag\{e^{2\beta _1},e^{2\beta _2},\mathrm{},e^{2\beta _6}\}`$, and $`_k\beta _k=0`$. Following , we introduce scalars $`\phi _1`$ through $`\phi _5`$ via the relation
$$\left(\begin{array}{c}\beta _1\\ \beta _2\\ \beta _3\\ \beta _4\\ \beta _5\\ \beta _6\end{array}\right)=\left(\begin{array}{ccccc}1/\sqrt{2}& 1/\sqrt{2}& 1/\sqrt{2}& 0& 1/\sqrt{6}\\ 1/\sqrt{2}& 1/\sqrt{2}& 1/\sqrt{2}& 0& 1/\sqrt{6}\\ 1/\sqrt{2}& 1/\sqrt{2}& 1/\sqrt{2}& 0& 1/\sqrt{6}\\ 1/\sqrt{2}& 1/\sqrt{2}& 1/\sqrt{2}& 0& 1/\sqrt{6}\\ 0& 0& 0& 1& \sqrt{2/3}\\ 0& 0& 0& 1& \sqrt{2/3}\end{array}\right)\left(\begin{array}{c}\phi _1\\ \phi _2\\ \phi _3\\ \phi _4\\ \phi _5\end{array}\right).$$
(23)
The sigma model metric on the five-dimensional space parametrized by the $`\phi _I`$ is $`\delta _{IJ}`$. The potential, $`V=\frac{1}{8L^2}\left[(trM)^22trM^2\right]`$, and the superpotential, $`W=\frac{1}{4}trM`$, can be written explicitly as sums of exponentials:
| $`V(\stackrel{}{\phi })`$ | $`={\displaystyle \underset{\alpha }{}}v_\alpha e^{\stackrel{}{\eta }_\alpha \stackrel{}{\phi }}`$ |
| --- | --- |
| $`W(\stackrel{}{\phi })`$ | $`={\displaystyle \underset{\alpha }{}}w_\alpha e^{\stackrel{}{\zeta }_\alpha \stackrel{}{\phi }}.`$ |
(24)
All the constants $`w_\alpha `$ are negative, but the same is not true of the $`v_\alpha `$. The six vectors $`\stackrel{}{\zeta }_\alpha `$ are of equal length, and they sum to zero. The level sets of $`W(\stackrel{}{\phi })`$ for large $`|\stackrel{}{\phi }|`$ are asymptotic to concentric “hexahedra” in $`𝐑^5`$. A hexahedron is the six-sided polyhedron which is the five-dimensional generalization of the tetrahedron. It has $`\left(\genfrac{}{}{0pt}{}{6}{n}\right)`$ $`(n1)`$-dimensional “faces,” where a zero-dimensional “face” is a vertex, a one-dimensional “face” is an edge, and so on up to $`n=6`$. The four-dimensional faces of the hexahedra are normal to the $`\stackrel{}{\zeta }_\alpha `$. The geodesic gradient flow trajectories which preserve $`SO(n)\times SO(6n)`$ symmetry are straight lines in $`𝐑^5`$ which pass through the center of an $`n1`$-dimensional face of each hexahedron; alternatively we may describe them as the rays emanating from the origin of $`𝐑^5`$ parallel to a sum of $`6n`$ of the vectors $`\stackrel{}{\zeta }_\alpha `$. Other gradient trajectories of $`W(\stackrel{}{\phi })`$, which lift to ten-dimensional geometries with unequal $`\mathrm{}_I`$, must run asymptotically parallel to one of the geodesics we have described. Those which run asymptotically parallel to some $`\stackrel{}{\zeta }_\alpha `$ correspond to cases with five non-zero $`\mathrm{}_I`$, and these are the ones we wish to rule out on account of the pathology of “ghost” D3-branes. This is exactly what the condition (2) does. The region $`V(\stackrel{}{\phi })<3/L^2`$ includes all the faces of the hexahedra of dimension less than four, but it excludes more and more of the four-dimensional faces as $`|\stackrel{}{\phi }|`$ increases, so that asymptotically none of the trajectories perpendicular to these faces are allowed.
All this may be easier to visualize by extension from figure 1, which shows a two-dimensional analog with three-fold symmetry to the five-dimensional example (24) with six-fold symmetry.
### 3.3 Asymptotics and phase transitions
In all of the examples in section 4, we will find $`V(\stackrel{}{\phi })`$ and $`W(\stackrel{}{\phi })`$ with the same form as (24), but in some cases some of the $`w_\alpha `$ will be positive. A gradient flow trajectory of $`W`$ which does not end at a saddle point will have an asymptotic direction, which we may specify by a unit vector $`\stackrel{}{n}`$. Let us parametrize the trajectory by a canonically normalized scalar $`\mu `$. From (5) one can extract a necessary condition on $`\stackrel{}{n}`$, namely
$$\underset{\alpha }{}w_\alpha \stackrel{}{\zeta }_\alpha e^{\stackrel{}{\zeta }_\alpha \mu \stackrel{}{n}}\left|\right|\stackrel{}{n}\text{asymptotically as }\mu \mathrm{}\text{,}$$
(25)
where $`\stackrel{}{u}||\stackrel{}{v}`$ means that $`\stackrel{}{u}`$ is parallel to $`\stackrel{}{v}`$: one is a positive multiple of the other. There is a straightforward geometrical construction to identify the solutions to (25). For each vector $`\stackrel{}{\zeta }_\alpha `$, define $`H_\alpha =\{\stackrel{}{\phi }:\stackrel{}{\phi }\stackrel{}{\zeta }_\alpha <1\}`$. The intersection $`H=_\alpha H_\alpha `$ is convex. It is a compact polyhedron in all examples we will encounter. Mathematically, it is similar in construction to a Brillouin zone, only the vectors $`\stackrel{}{\zeta }_\alpha `$ do not form a lattice. Now define a function $`\sigma :H\{1,0,+1\}`$ as follows. Given a point $`P`$ on the boundary $`H`$ of $`H`$, consider all those $`H_\alpha `$ for which $`P`$ is also on the boundary of $`H_\alpha `$. If all the associated $`w_\alpha `$ are negative, let $`\sigma (P)=1`$; if all are positive, let $`\sigma (P)=+1`$; and if some are positive and some are negative, let $`\sigma (P)=0`$. Define an apothem of $`H`$ as a vector with its tail at the origin and its head on $`H`$, such that the vector is orthogonal to the lowest-dimensional face of $`H`$ which it touches. We are again thinking of faces in the generalized sense: they can be vertices, edges, two-dimensional faces, etc. The apothems of $`H`$ include the vectors $`\stackrel{}{\zeta }_\alpha /|\stackrel{}{\zeta }_\alpha |^2`$, and also certain linear combinations of these vectors. If $`\stackrel{}{\kappa }`$ is an apothem of $`H`$ such that $`\sigma =1`$ at the head of the apothem, then there is a solution $`\stackrel{}{n}`$ to (25) which is parallel to $`\stackrel{}{\kappa }`$. And any solution $`\stackrel{}{n}`$ to (25) must be parallel to an apothem of $`H`$ for which $`\sigma =1`$ or $`0`$. Along a gradient flow trajectory which is asymptotically parallel to an apothem $`\stackrel{}{\kappa }`$, we will have $`We^{\zeta \mu }`$ for large $`\mu `$, where $`\zeta =1/|\stackrel{}{\kappa }|`$. Since
$$V\frac{1}{8}\left(\frac{W}{\mu }\right)^2\frac{1}{3}W^2,$$
(26)
we learn that $`V\mathrm{}`$ if $`\zeta <\sqrt{8/3}`$, and $`V+\mathrm{}`$ if $`\zeta >\sqrt{8/3}`$. Thus the condition (2) rules out directions in which $`|\stackrel{}{\kappa }|<\sqrt{3/8}`$. In the cases where $`\sigma =0`$ at the head of $`\stackrel{}{\kappa }`$ or where $`|\stackrel{}{\kappa }|=\sqrt{3/8}`$, the behaviors of $`W`$ and $`V`$ have to be checked explicitly.
If we imagine $`W(\stackrel{}{\phi })`$ as a mountain, and ourselves as skiers, the vectors $`\stackrel{}{\kappa }`$ indicate the possible asymptotic directions that we could take if we started at $`\stackrel{}{\phi }=0`$ and followed the fall line all the way down. Vectors with $`|\stackrel{}{\kappa }|<\sqrt{3/8}`$ indicate directions which are “too steep” for safe skiing and violate (2) as a result. When all the $`w_\alpha `$ are negative, the level curves of $`W(\stackrel{}{\phi })`$ for large $`|\stackrel{}{\phi }|`$ approximate concentric dilated copies of $`H`$. When some of the $`w_\alpha `$ are positive and some are negative, the level sets are more complicated, but for large negative or large positive values of $`W`$ they include the faces of dilations of $`H`$ with $`\sigma =1`$ or $`\sigma =+1`$, respectively.
Just as we constructed the convex polyhedron $`H`$ and the sign map $`\sigma `$ starting with $`W`$ in the form of a sum of exponentials, we can construct a convex polyhedron $`K`$ and a sign map $`\rho `$ starting with $`V`$. Let $`L`$ be the union of rays which start at the origin and pass through a point on $`K`$ where $`\rho =0`$. The codimension of $`L`$ is at least one. The zero contour of $`V`$ must be asymptotic to $`L`$ for large $`|\stackrel{}{\phi }|`$. It is not guaranteed that every ray in $`L`$ is asymptotically close to the zero contour of $`V`$, but this is guaranteed for a ray in $`L`$ associated with a point $`P`$ on $`K`$ such $`\rho `$ takes both positive and negative values close to $`P`$. The set $`L`$ divides the space of scalars into wedges which do or do not violate (2) according as $`\rho =+1`$ or $`\rho =1`$ for the face(s) of $`K`$ lying in that wedge. Gradient flow trajectories which run on or asymptotic to $`L`$ must be checked explicitly.
It is transparent from the foregoing discussion that trajectories which pass (2) for potentials of the form (24) will also satisfy
$$\frac{W}{\stackrel{}{\phi }}\left|\right|\frac{V}{\stackrel{}{\phi }}\text{asymptotically as }A(r)\mathrm{}.$$
(27)
This condition can also be motivated (although not rigorously) by finite temperature considerations. In the far region of a near-extremal geometry, where $`h1`$, the equations (7) are approximately solved by a solution to the first order equations, (5). In the near-horizon region, (12) should apply to a good approximation. It is natural to think that the far and near regions can be matched onto each other provided $`\stackrel{}{\phi }^{}`$ is pointing roughly in the same direction in both regions. Modulo some (fairly strong) technical assumptions, it is possible to show that (27) is a necessary condition for the existence of a sequence of non-extremal solutions which converge to an extremal solution in the $`H_1^{\mathrm{}}`$ norm. However, the $`H_1^{\mathrm{}}`$ norm is a very strong topology in the infrared, and it may be that near-extremal solutions do not in general converge to an extremal limit in this norm.
In exploring examples, it is often convenient to look first at a more restricted category of trajectories than (2) allows. Let us define
$$𝒫=\{\stackrel{}{\phi }:\pm \frac{W}{\stackrel{}{\phi }}|\left|\frac{V}{\stackrel{}{\phi }}\right\}.$$
(28)
Given $`W`$ and $`V`$ in closed form, it is particularly straightforward to find $`𝒫`$ because it is the solution set to the equations $`ϵ^{i_1\mathrm{}i_n}(W/\phi ^{i_1})(V/\phi ^{i_2})=0`$. Generically, $`𝒫`$ consists of several intersecting curves. Some of them are asymptotic to gradient flow trajectories of $`W`$ in the limit where $`A\mathrm{}`$. They can sometimes be excellent approximations to particular gradient flow trajectories even at finite $`A`$. In the Coulomb branch example of section 3.2, $`𝒫`$ consists of lines through the origin passing perpendicularly through every face (of every dimension) of the hexahedron. It is thus precisely the union of all the gradient flow trajectories preserving some $`SO(n)\times SO(6n)`$ symmetry. Every trajectory runs asymptotically parallel to some line in $`𝒫`$ in this case. $`𝒫`$ will not have quite such nice properties in the examples of section 4, but it remains true (as will be seen case by case) that every trajectory which passes the criterion (2) is also asymptotically parallel to some curve in $`𝒫`$. It would be nice to know if there is any deep reason for this.
It is straightforward to categorize the possible behaviors in the infrared, provided the scalars’ evolution has an asymptotic direction. As explained above, this assumption is the generic expectation given a potential $`V(\stackrel{}{\phi })`$ of the form (4). Let us therefore consider in more detail the choice $`W(\mu )=e^{\zeta \mu }`$. The variable $`\mu `$ is understood as the canonically normalized scalar along the asymptotic direction of the flow. Our only aim is to obtain the correct scaling behavior, so we drop overall factors, including dimensionful ones like the asymptotic radius of curvature of $`AdS_5`$. From (5) one easily obtains
| $`ds^2`$ | $`=(rr_0)^{\frac{4}{3\zeta ^2}}(dt^2+d\stackrel{}{x}^2)+dr^2`$ |
| --- | --- |
| $`e^{\zeta \mu (r)}`$ | $`={\displaystyle \frac{\zeta ^2}{2}}(rr_0),`$ |
(29)
where $`r_0`$ is an integration constant which we readily identify as the location of the singularity. Curvature invariants involving $`n`$ derivatives diverge like $`(rr_0)^{2n}`$ as $`rr_0`$. The causal structure of the geometry (29) depends on the value of $`\zeta `$, and is more transparent if one introduces a new variable $`z`$ such that $`dr=e^Adz`$. Let us assume that the geometry (29) merges into asymptotically anti-de Sitter spacetime at large $`r`$. Then it is straightforward to show that if $`\zeta >\sqrt{2/3}`$ the Penrose diagram is a strip, as in figure 2(a), while if $`\zeta \sqrt{2/3}`$ the Penrose diagram is a wedge, as in figure 2(b). Clearly the singularity gets worse as $`\zeta `$ gets larger: a naked timelike singularity, as in figure 2(a), is more problematic than a null singularity of the type in figure 2(b), which is similar to the singular horizon of the D2-brane.
When $`\zeta >\sqrt{8/3}`$, $`V+\mathrm{}`$ as $`rr_0`$, and the conjecture (2) is that the singularity is unphysical. Certainly it fails to admit near-extremal generalizations. When $`\zeta <\sqrt{8/3}`$, $`V\mathrm{}`$, and (2) is satisfied. In this case, $`|V(\mu )|`$ grows slower than $`e^{\sqrt{32/3}\mu }`$. Of the solutions $`W(\mu )`$ to
$$\frac{1}{8}\left(\frac{W}{\mu }\right)^2\frac{1}{3}W^2=V=\left(\frac{\zeta ^2}{8}\frac{1}{3}\right)e^{2\zeta \mu },$$
(30)
only $`W(\mu )=\pm e^{\zeta \mu }`$ grows more slowly than $`e^{\sqrt{8/3}\mu }`$ as $`\mu \mathrm{}`$. All the other solutions are asymptotically of the form $`W(\mu )(const)e^{\sqrt{8/3}\mu }`$. It is possible to integrate (30) analytically, but the result is an complicated implicit equation relating $`W`$ and $`\mu `$. However the asymptotics is obvious provided $`|V(\mu )|`$ grows slower than $`e^{\sqrt{32/3}\mu }`$. This implies that the generic Poincaré-invariant solution to the equations of motion following from the action (3) is of the form (29) with $`\zeta =\sqrt{8/3}`$—provided $`|V(\stackrel{}{\phi })|`$ grows slower than $`e^{\sqrt{32/3}|\stackrel{}{\phi }|}`$. This result has been obtained previously by a direct analysis of the equations of motion .
We have already demonstrated in two different ways that (2) is a necessary condition for the existence of near-extremal solutions, once by using the zero energy constraint plus the null energy condition, and once by using the scalar equation of motion. It will be instructive to present one more proof—inferior to the previous two because it is less rigorous and less general, but suggestive of interesting physics at finite temperature. Start with a Poincaré invariant solution generated using a $`W`$ which is asymptotic to $`e^{\zeta \mu }`$ with $`\zeta >\sqrt{8/3}`$. To obtain near-extremal generalizations, one would solve (7) and (8) with a small value of $`B`$ in (9). Let us proceed on the assumption that for small $`B`$, the solution for $`A(r)`$ and $`\stackrel{}{\phi }(r)`$ changes only slightly. Setting $`L=1`$ and neglecting factors of order unity, the expression for $`h(r)`$ is
| $`h(r)`$ | $`=1B{\displaystyle _r}𝑑r_1e^{4A(r_1)}`$ |
| --- | --- |
| | $`1B{\displaystyle _r}𝑑r_1(r_1r_0)^{\frac{8}{3\zeta ^2}}.`$ |
(31)
The integral in the second line remains finite as $`rr_0`$, so it is indeed true that small $`B`$ uniformly suppresses the effects of $`h`$ in the equations of motion. Thus the approximation used in the second line of (31) is reliable for sufficiently small $`B`$. The conclusion is that $`h(r)`$ has no zero for small $`B`$: there is no horizon! This indeed indicates that there are no near-extremal generalizations. For sufficiently large $`B`$, and assuming the criterion (14) can be met, it is plausible that solutions with horizons will exist, and that they have a temperature which diverges as $`B\mathrm{}`$. Such solutions might be dubbed “far-extremal,” since they are not close to the original Poincaré invariant solution. Assuming the spacetime is asymptotically $`AdS_5`$, solutions sufficiently far from extremality (in the sense of $`B`$ being large) would be nearly AdS-Schwarzschild. Let $`B_0`$ be the infimum of the set of $`B`$’s for which there is a horizon. It is again plausible that $`e^{2A(r)}`$ and $`h(r)`$ vanish simultaneously when $`B=B_0`$, and that for $`B`$ larger than $`B_0`$ but sufficiently close to it, $`h(r)=0`$ at a location $`r=r_H`$ where $`e^{2A(r)}`$ can be made as small as we wish. Since $`B_0`$ is fixed and finite, the formula $`T=\frac{B}{\pi L}e^{3A(r_H)}`$ implies that the temperature diverges as $`BB_0`$. Thus there must be some $`B_c>B_0`$ where $`T`$ attains its minimum. (Let us assume for simplicity that $`T`$ as a function of $`B`$ has only one extremum). The mass above extremality is related monotonically to $`B`$. Generically one expects that this relation is approximately linear near $`B=B_c`$, and that $`TT_c+(BB_c)^2`$. Then as we approach the critical point from above (both in energy and temperature), $`EE_c(TT_c)^{1/2}`$, which implies a divergence in the specific heat, $`C(TT_c)^{1/2}`$. This is the signal of a second order phase transition at finite temperature. On the field theory side, if one makes the generic assumption that the inverse square of the correlation length goes to zero like $`TT_c`$ (which is to say, $`\nu =1/2`$), then the specific heat exponent $`\alpha =1/2`$ emerges from the hyperscaling relation, $`\alpha =2d\nu `$ with $`d=3`$. Why the strong interactions should not change the story drastically is unclear.
If $`\zeta <\sqrt{8/3}`$, the integral in the second line of (31) diverges as $`rr_0`$, and one might hope to argue on this basis that near-extremal generalizations of the Poincaré invariant solution do exist. But when $`h`$ deviates significantly from $`1`$, it is hard to justify the approximation of keeping the same $`A(r)`$ and $`\stackrel{}{\phi }(r)`$ as in the Poincaré invariant solution. Thus we cannot make any rigorous claims. Let us nevertheless attempt to obtain an equation of state, assuming the asymptotics
| $`A(r)`$ | $`\alpha _1\mathrm{log}(rr_0)+A_0`$ |
| --- | --- |
| $`h(r)`$ | $`1B{\displaystyle _r}𝑑r_1e^{4A(r_1)},`$ |
(32)
with $`\alpha _1<1/4`$. Again we set $`L=1`$ and neglect factors of order unity, which is OK since the goal is only to obtain a scaling relation. Equation (10) may be used to obtain
| $`B`$ | $`=e^{4A_0}(4\alpha _11)(r_Hr_0)^{4\alpha _11}`$ |
| --- | --- |
| $`T`$ | $`=e^{A_0}(4\alpha _11)(r_Hr_0)^{\alpha _11}`$ |
| $`{\displaystyle \frac{S}{V}}`$ | $`=e^{3A_0}(r_Hr_0)^{3\alpha _1}.`$ |
(33)
The factors of $`e^{A_0}`$ indicate the how each quantity scales under a conformal transformation. $`B`$ indeed scales as energy density. Eliminating $`r_Hr_0`$ leads to
$$\frac{S}{V}T^\eta \text{with}\eta =\frac{3\alpha _1}{\alpha _11}=\frac{6}{23\zeta ^2}$$
(34)
where in the last line we have assumed that $`\alpha _1=\frac{2}{3\zeta ^2}`$ as in the Poincaré invariant flow, (29). The calculation is probably trustworthy when $`\zeta `$ is sufficiently small. Note that as $`\zeta 0`$ one recovers the scaling $`ST^3`$ which applies to a conformal field theory. It is interesting to note that in the range where $`(\text{34})`$ predicts positive specific heats, it also predicts that $`F/T^4`$ is an increasing function of temperature. This is consistent with the Appelquist-Cohen-Schmaltz conjecture .<sup>3</sup><sup>3</sup>3Despite efforts by the current author and by R. Myers, there is no general proof of the monotonicity of $`F/T^4`$ in solutions of the form (1). I thank R. Myers for an extensive correspondence on this topic. For $`\zeta >\sqrt{2/3}`$, the evidence of (34) is that $`T\mathrm{}`$ as $`r_Hr_0`$. It could be that instead $`T`$ remains finite in this limit. Either of these two outcomes would signal a finite temperature phase transition. Or it is possible that $`T0`$ as $`r_Hr_0`$, and $`S`$ as a function of $`T`$ falls off faster than any power, indicating a mass gap.
It would be very interesting to find sufficient conditions for the existence of black hole solutions arbitrarily close to a given Poincaré invariant flow in an appropriately weak topology. It would also be interesting to find more stringent necessary conditions than (2). For instance, one might conjecture that a necessary condition for a gradient flow to have finite temperature generalizations which converge to the original geometry in the topology induced by the norm (18) is that the gradient flow should either terminate at a saddle point or converge to a curve in $`𝒫`$.
Section 4 is concerned with additional examples of Poincaré invariant flows where the field theory physics is more or less clear. It is useful to categorize the flows according to the value of $`\zeta `$ which controls their far infrared behavior. This is done in figure 3.
The responses to finite temperature indicated in figure 3 are based partly on calculation, partly on AdS/CFT examples where the field theory physics is clear, and partly on conjecture. The arguments (31) consitute fairly good evidence that for flows dominated by some direction in which $`\zeta >\sqrt{8/3}`$, there is a finite temperature phase transition. On the other hand, the AdS-Schwarzschild solution demonstrates that there are near-extremal geometries with arbitrarily small temperature when $`\zeta =0`$. The analysis leading to (34) suggests that for $`\zeta <\sqrt{2/3}`$ there are still near-extremal geometries with arbitrarily small temperature. It was found in that minimally coupled scalars exhibit no gap for geometries of the form (1) when $`\zeta <\sqrt{2/3}`$. This fits with the power law behavior $`ST^\eta `$ found in (34). Note that $`\zeta =\sqrt{2/3}`$ is also the value where the Penrose diagram for the Poincaré invariant flow changes from figure 2(b) to figure 2(a): for $`\zeta <\sqrt{2/3}`$ the singularity is null, while for $`\zeta >\sqrt{2/3}`$ it is timelike. What happens for $`\sqrt{2/3}<\zeta <\sqrt{8/3}`$ is harder to understand. Minimal scalars exhibit a mass gap, so one must expect that $`S`$ falls off faster than any power of $`T`$ for sufficiently small $`T`$. It could be however that curvatures at the horizon become Planckian before this behavior can be observed. If so, the conclusion is simply that a string theory resolution of the singularity is necessary before the low-temperature physics can be fully understood.
The range of $`\zeta `$’s for the Coulomb branch is puzzling. How can this simple system exhibit such a range of behaviors? One possible explanation was proposed in . It is that the supergravity geometries specify only a continuous distribution of D3-branes in ten dimensions, and at finite $`N`$ the physics is most likely an ensemble average of all discrete distributions which sufficiently closely approximate the continuous one. Roughly speaking, any D3-brane can wiggle as far as its nearest neighbor. Performing the ensemble average first in the path integral could lead to extra terms in an effective lagrangian, for instance terms of the form $`(\mathrm{\Sigma }^{IJ}trX_{(I}X_{J)})^2`$ where $`\mathrm{\Sigma }^{IJ}`$ is some tensor in the $`\mathrm{𝟐𝟎}^{}`$. These extra effective interactions get larger as the dimensionality of the brane distribution increases, and could perhaps explain the peculiar properties that supergravity seems to ascribe to these configurations. A naive estimate of the strength of the double trace terms relative to the original lagrangian is $`N^{12/p}`$ where $`p`$ is the dimensionality of the distribution. Another possibility (also speculative) is that the low-energy theory of open strings on a $`p`$-dimensional continuous distribution of D3-branes is effectively $`p+4`$ dimensional, simply because the open strings are allowed to end on any D3-brane. Then the presence or absence of a mass gap can be understood in terms of the higher dimensional theory.<sup>4</sup><sup>4</sup>4This interpretation was suggested to me by J. Polchinski and J. Maldacena.
## 4 Examples
The examples of sections 4.1 and 4.3 have appeared previously in the literature ; the example in section 4.2 is based on work with K. Pilch and N. Warner . On the supergravity side, the superpotential $`W`$ is associated with supersymmetry transformations: the first order equations (5) are precisely the conditions for some fraction of $`d=5`$ $`𝒩=8`$ supersymmetry to remain unbroken. As explained after (4), $`W`$ is an eigenvalue of a $`SO(6)\times SL(2,𝐑)`$-invariant matrix $`W_{ab}`$. Which eigenvalue is the correct one to describe a given relevant deformation can usually be deduced from how much supersymmetry is unbroken, plus known asymptotics near the boundary of $`AdS_5`$.
On the field theory side, all the examples come from relevant deformations of $`𝒩=4`$ supersymmetric Yang-Mills theory which preserve at least $`𝒩=1`$ supersymmetry. The on-shell spectrum of $`𝒩=4`$ super-Yang-Mills is a vector, four chiral fermions, and six real bosons, all in the adjoint of the gauge group:
| $`A_\mu `$ |
| --- |
| $`\lambda _1\lambda _2\lambda _3\lambda _4`$ |
| $`X_1X_2X_3X_4X_5X_6.`$ |
(35)
These can be grouped into $`𝒩=1`$ multiplets in various equivalent ways. The one we will have in mind is $`(A_\mu ,\lambda _4)`$ for the $`𝒩=1`$ vector multiplet and $`(\lambda _1,X_1,X_2)`$, $`(\lambda _2,X_3,X_4)`$, and $`(\lambda _3,X_5,X_6)`$ for the adjoint chiral multiplets—also denoted by chiral superfields $`\mathrm{\Phi }_1`$, $`\mathrm{\Phi }_2`$, $`\mathrm{\Phi }_3`$, whose complex scalar components are $`\varphi _1`$, $`\varphi _2`$, $`\varphi _3`$. The superpotential of the undeformed $`𝒩=4`$ theory is
$$W=tr\mathrm{\Phi }_1[\mathrm{\Phi }_2,\mathrm{\Phi }_3].$$
(36)
Supergravity scalars will typically be denoted $`\phi _\mathrm{\Delta }`$, where $`\mathrm{\Delta }`$ indicates the dimension of the corresponding gauge singlet operator $`𝒪_\mathrm{\Delta }`$. Units will be chosen such that $`L`$, the radius of curvature in the asymptotically $`AdS_5`$ region, is $`1`$.
### 4.1 One massive adjoint chiral
Our first example is the category of flows that arise from giving a mass to a single adjoint chiral superfield in the $`𝒩=4`$ langrangian . The mass deformation is specified by profiles for two scalars, $`\phi _2`$ and $`\phi _3`$, dual to a boson mass term $`𝒪_2`$ of dimension $`2`$ and a fermion mass term $`𝒪_3`$ of dimension $`3`$. Explicitly,
| $`𝒪_2`$ | $`=tr\left(X_1^2X_2^2X_3^2X_4^2+2X_5^2+2X_6^2\right)`$ |
| --- | --- |
| $`𝒪_3`$ | $`=tr\left(\lambda _4\lambda _4+\varphi _1[\varphi _2,\varphi _3]\right)+h.c.`$ |
(37)
The operator $`𝒪_2`$ in (37) looks like the wrong mass deformation: adding a positive multiple of it to the lagrangian will result in negative mass squared for four of the six real scalars. What we really aim to do is add $`d^2\theta \frac{1}{2}mtr\mathrm{\Phi }_3^2+h.c.`$ to the lagrangian. The operator $`𝒪_3`$ in (37) is precisely the dimension $`3`$ term in this deformation. But the dimension $`2`$ term is just $`tr\left(X_5^2+X_6^2\right)`$. What saves the day is that $`\phi _3^2`$ can couple to the $`SO(6)`$ invariant combination $`tr_{I=1}^6X_I^2`$. In a supersymmetric background involving the $`e^r`$ part of $`\phi _3`$ and the $`re^{2r}`$ part of $`\phi _2`$, the $`\phi _3^2`$ coupling to the $`SO(6)`$ singlet mass term must precisely cancel out the negative terms in the $`𝒪_2`$ of (37), leaving only the desired $`tr\left(X_5^2+X_6^2\right)`$. Similar quadratic couplings will be important in sections 4.2 and 4.3.
The scalars $`\phi _2`$ and $`\phi _3`$ are canonically normalized, and the potential and superpotential read
| $`V(\phi _2,\phi _3)`$ | $`=\left[\rho ^2+{\displaystyle \frac{\mathrm{cosh}2\phi _33}{4\rho ^4}}+{\displaystyle \frac{\rho ^8}{8}}\left(\mathrm{cosh}2\phi _31\right)\right](\mathrm{cosh}2\phi _3+1)`$ |
| --- | --- |
| $`W(\phi _2,\phi _3)`$ | $`={\displaystyle \frac{1}{2}}\mathrm{cosh}2\phi _3\left(\rho ^4{\displaystyle \frac{2}{\rho ^2}}\right)\left({\displaystyle \frac{3}{2}}\rho ^4+{\displaystyle \frac{1}{\rho ^2}}\right)`$ |
(38)
where we have defined $`\rho =e^{\phi _2/\sqrt{6}}`$. It is consistent to set all other scalars to $`0`$: locally $`V`$ is stationary perpendicular to the $`\phi _2`$-$`\phi _3`$ plane. There is a one-parameter family of flows emanating from the origin, which is the maximally symmetric point representing unperturbed $`𝒩=4`$ gauge theory. The trajectories have only been found numerically. Their asymptotics for large $`r`$ (the ultraviolet) was described in :
$$\phi _3a_0e^r\phi _2\frac{4}{\sqrt{6}}a_0^2re^{2r}+a_1e^{2r}.$$
(39)
Unsurprisingly, the leading $`re^{2r}`$ behavior of $`\phi _2`$ is fixed in terms of the leading $`e^r`$ behavior of $`\phi _3`$: this reflects the fact that the boson mass term and the fermion mass term are related by supersymmetry. One can parametrize the gradient flow trajectories with the quantity
$$\stackrel{~}{a}_1=\frac{a_1}{a_0^2}+\frac{4}{\sqrt{6}}\mathrm{log}a_0,$$
(40)
which is invariant under additive shifts of $`r`$. It is a dimensionless measure of $`𝒪_2`$ in units set by the scale of the mass deformation. $`W`$ and $`V`$ are symmetric under $`\phi _3\phi _3`$, and from now on we will simplify our discussion by considering only flows with $`\phi _30`$.
Curiously, it was found in that $`\stackrel{~}{a}_1=\stackrel{~}{a}_1^{(c)}=1.4694\mathrm{}`$ corresponds to the conformal vacuum of the theory, where $`𝒪_2=0`$ and the supergravity flow terminates at the critical point found in . By comparison with flows describing states on the Coulomb branch of $`𝒩=4`$ gauge theory, one can establish that $`\stackrel{~}{a}_1>\stackrel{~}{a}_1^{(c)}`$ corresponds to negative $`𝒪_2`$. Positive $`𝒪_2`$ should be impossible because we have made $`X_5`$ and $`X_6`$ massive. To put it another way, the infrared fixed point theory has a Coulomb branch parametrized by color singlet combinations of $`X_1`$ through $`X_4`$. The states with negative $`𝒪_2`$ are $`SU(2)_{\mathrm{global}}\times U(1)_R`$-symmetric states on that Coulomb branch. States with positive $`𝒪_2`$ are unphysical. In figure 4(b), the $`𝒪_2<0`$ trajectories are the ones to the right of the critical trajectory which ends up at the one-quarter supersymmetric fixed point. The ones to the left are supposed to be ruled out.
Felicitously, the condition (2) accepts trajectories with $`\stackrel{~}{a}_1\stackrel{~}{a}_1^{(c)}`$ but rules out those with $`\stackrel{~}{a}_1<\stackrel{~}{a}_1^{(c)}`$. The horizontal line along the $`\phi _2`$ axis in the solution set $`𝒫`$ (defined in (28)) represents Coulomb branch flows of the $`𝒩=4`$ gauge theory. When these flows are lifted back to ten dimensions, the $`\phi _2>0`$ direction corresponds to an $`S^3`$ shell of D3-branes in the hyperplane spanned by $`x_1`$ through $`x_4`$, while the $`\phi _2<0`$ direction corresponds to a uniform disk of D3-branes in the plane spanned by $`x_5`$ and $`x_6`$. Because the contours of $`W`$ form a trough around the $`\phi _2`$ axis for large positive $`\phi _2`$, the $`\stackrel{~}{a}_1>\stackrel{~}{a}_1^{(c)}`$ trajectories are drawn to this axis. The very low energy Coulomb branch physics of these flows does not depend on whether the adjoints $`X_5`$ and $`X_6`$ are given a mass. In contrast to this case, the contours of $`W`$ form a ridge around the negative $`\phi _2`$ axis, and the $`\stackrel{~}{a}_1<\stackrel{~}{a}_1^{(c)}`$ trajectories are repelled from this axis. These trajectories are asymptotically parallel to one another, and they are ruled out by (2) because $`V(\stackrel{}{\phi }(r))+\mathrm{}`$ as $`A(r)\mathrm{}`$. The flow exactly along the negative $`\phi _2`$ axis represents sensible $`𝒩=4`$ Coulomb branch physics in the $`X_5`$ and $`X_6`$ directions; but this part of the moduli space is lifted completely when $`X_5`$ and $`X_6`$ are given a mass, and the trajectories are appropriately defocused as a result.
It turns out, as claimed after (28), that all the trajectories permitted by (2) run asymptotically parallel to some curve in $`𝒫`$. In fact, all trajectories converge to some curve in $`𝒫`$ in the far infrared. The segment of $`𝒫`$ stretching from the origin to the one-quarter supersymmetric saddle point is an excellent approximation to a gradient flow trajectory. The part of the curve in $`𝒫`$ which extends northeast from the one-quarter supersymmetric point does not describe a deformation of $`𝒩=4`$ super-Yang-Mills.
The flows along the positive and negative $`\phi _2`$ axes were first considered in ten dimensions as zero-temperature, zero-spin limits of near-extremal D3-brane solutions with some angular momenta turned on. Specifically, the positive $`\phi _2`$ trajectory comes from solutions with equal $`J_{12}`$ and $`J_{34}`$ turned on, while the negative $`\phi _2`$ trajectory comes from solutions with $`J_{56}`$. In five dimensions these angular momenta are interpreted as charges under the $`U(1)^3`$ Cartan subalgebra of the $`SO(6)`$ gauge group, and the spinning D3-branes are represented as charged black holes . The time components of the gauge fields in five-dimensional supergravity act as “voltages,” or “chemical potentials,” for conserved $`R`$-currents in the gauge theory. The positive and negative $`\phi _2`$ flows are examples of states which can be approached by tuning these chemical potentials to zero in fixed ratio with an appropriate power of the temperature. There is a thermodynamic instability that sets in at some point in this limiting process , but that is somewhat outside our current scope of interest.
### 4.2 Two massive adjoint chirals
Our next example is $`𝒩=4`$ super-Yang-Mills broken to $`𝒩=2`$ by a mass term for an adjoint hypermultiplet. The analysis is based on joint work with K. Pilch and N. Warner . In $`𝒩=1`$ language, we are giving equal but opposite masses to two adjoint chiral multiplets, which we take to be $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$—exactly the chiral adjoints which we left massless in the previous example. The field theory was studied in . The supergravity involves the same scalar $`\phi _2`$ in the $`\mathrm{𝟐𝟎}^{}`$ of $`SO(6)`$ as we had in the previous example, but a different scalar $`\stackrel{~}{\phi }_3`$ in the $`\mathrm{𝟏𝟎}+\overline{\mathrm{𝟏𝟎}}`$. The dual operators are
| $`𝒪_2`$ | $`=tr\left(X_1^2X_2^2X_3^2X_4^2+2X_5^2+2X_6^2\right)`$ |
| --- | --- |
| $`\stackrel{~}{𝒪}_3`$ | $`=tr\left(\lambda _1\lambda _1\lambda _2\lambda _2+\text{(scalar trilinear)}\right)+h.c.`$ |
(41)
As in the previous example, the square of $`\stackrel{~}{\phi }_3`$ can couple to the $`SO(6)`$ singlet mass operator. The combination of this operator and $`𝒪_2`$ that supersymmetry demands is proportional to $`tr\left(X_1^2+X_2^2+X_3^2+X_4^2\right)`$. Both $`\phi _2`$ and $`\stackrel{~}{\phi }_3`$ are canonically normalized, and the potential and superpotential read
| $`V(\phi _2,\stackrel{~}{\phi }_3)`$ | $`={\displaystyle \frac{1}{4}}e^{\frac{4}{\sqrt{6}}\phi _2}\left(48e^{\sqrt{6}\phi _2}\mathrm{cosh}\sqrt{2}\stackrel{~}{\phi }_3+e^{2\sqrt{6}\phi _2}\mathrm{sinh}^2\sqrt{2}\stackrel{~}{\phi }_3\right)`$ |
| --- | --- |
| $`W(\phi _2,\stackrel{~}{\phi }_3)`$ | $`=2e^{\frac{2}{\sqrt{6}}\phi _2}e^{\frac{4}{\sqrt{6}}\phi _2}\mathrm{cosh}\sqrt{2}\stackrel{~}{\phi }_3.`$ |
(42)
The asymptotics at large $`r`$ is
$$\stackrel{~}{\phi }_3b_0e^r\phi _2\frac{2}{\sqrt{6}}b_0^2re^r+b_1e^{2r},$$
(43)
and the invariant which parametrizes the one-parameter family of flows emanating from the origin is
$$\stackrel{~}{b}_1=\frac{b_1}{b_0^2}\frac{2}{\sqrt{6}}\mathrm{log}b_0.$$
(44)
In figure 5 we exhibit contours of $`V`$ and $`W`$, the solution set $`𝒫`$, and some characteristic gradient flow trajectories (obtained numerically). Because of the $`\stackrel{~}{\phi }_3\stackrel{~}{\phi }_3`$ symmetry, we can restrict our attention to flows with $`\stackrel{~}{\phi }_30`$. When $`\stackrel{~}{b}_1=\stackrel{~}{b}_1^{(c)}0.125\mathrm{}`$,<sup>5</sup><sup>5</sup>5Random jitter near the maximally supersymmetric point makes this number (and also $`\stackrel{~}{a}_1^{(c)}`$) surprisingly hard to pin down numerically. the flow is asymptotic to the curve in $`𝒫`$ which follows the obvious northwest ridgeline of $`W`$ (see figure 5). The condition (2) allows this trajectory and all others to the left of it: $`\stackrel{~}{b}_1\stackrel{~}{b}_1^{(c)}`$. The trajectories with $`\stackrel{~}{b}_1<\stackrel{~}{b}_1^{(c)}`$ run asymptotically parallel to the horizontal axis: this is an explicit example where (27) is satisfied without the trajectories actually approaching any curve in $`𝒫`$. (States on the Coulomb branch which fail to preserve an $`SO(n)\times SO(6n)`$ global symmetry are another example). The $`\stackrel{~}{b}_1=\stackrel{~}{b}_1^{(c)}`$ trajectory is well approximated by a curve in $`𝒫`$. The trajectories with $`\stackrel{~}{b}_1>\stackrel{~}{b}_1^{(c)}`$ proceed down the northeast face of $`W`$. As is evident from comparing figure 5(a) with figure 5(b), this is a case where (27) is violated. No gradient flow trajectory of $`W`$ ends up at the saddle points near the northeast and southeast corners of figure 3(a): these points correspond to $`AdS_5`$ vacua which break all supersymmetry (see Table 1, entry (iv) of ).
There is a natural field theory interpretation of these flows, inspired by , though it must be regarded as speculative in the absence of more detailed tests. The far infrared theory is large-$`N`$ Seiberg-Witten theory, and the Seiberg-Witten curve is
$$y^2=\underset{k=1}{\overset{N}{}}(x\xi _k)^2\mathrm{\Lambda }^{2N},$$
(45)
where $`\mathrm{\Lambda }`$ is the dynamically generated scale and $`\xi _k`$ are parameters determining the position on the moduli space. In it was suggested that the enhançon singularity corresponds to Seiberg-Witten theory at the “origin” of moduli space $`\xi _k=0`$, and that the branch points $`x_k=\mathrm{\Lambda }e^{i\pi k/N}`$ correspond to the positions of branes in a ring around the enhançon. The $`\stackrel{~}{b}_1=\stackrel{~}{b}_1^{(c)}`$ flow seems most likely to be dual to this same point in moduli space, where now the complex variable $`x`$ parametrizes the directions of the adjoints $`X_5`$ and $`X_6`$ which remain massless. This interpretation is supported by the fact that the supergravity geometry preserves a $`U(1)`$ symmetry which rotates $`X_5`$ and $`X_6`$. Because the $`\stackrel{~}{b}_1>\stackrel{~}{b}_1^{(c)}`$ flows “interpolate” between the ridgeline flow and the negative $`\phi _2`$ axis, and because the negative $`\phi _2`$ axis is known to correspond to a state on the Coulomb branch of the undeformed $`𝒩=4`$ theory where the D3-branes are arranged in a uniform disk in the $`X_5`$$`X_6`$ direction, the most obvious candidate for a field theory dual of the $`\stackrel{~}{b}_1>\stackrel{~}{b}_1^{(c)}`$ flows is a point on moduli space where the $`\xi _i`$ are uniformly distributed on a disk. When the radius of that disk is much larger than $`\mathrm{\Lambda }`$, the low-energy physics approaches the Coulomb branch physics of $`𝒩=4`$ gauge theory.
On this interpretation, the trajectories with $`\stackrel{~}{b}_1<\stackrel{~}{b}_1^{(c)}`$ are unphysical. Without the mass deformation, the disk of D3-branes in the $`X_5`$$`X_6`$ directions can shrink to a point and then re-expand as an $`S^3`$ shell in the $`X_1`$ through $`X_4`$ directions. This is exactly the picture found in . But when we make $`X_1`$ through $`X_4`$ massive, this part of the moduli space disappears, and we are in the position of attempting to give a negative VEV to $`tr(X_5^2+X_6^2)`$. All this fits in perfectly with (2) which indeed rules out the $`\stackrel{~}{b}_1<\stackrel{~}{b}_1^{(c)}`$ trajectories. To lend further support to the picture we have presented, it would nice if one could show that finite temperature in the field theory draws the vacuum state to the point where all $`\xi _k=0`$. Correspondingly in the supergravity, one would hope that a horizon with small but finite temperature draws the scalars onto the $`\stackrel{~}{b}_1=\stackrel{~}{b}_1^{(c)}`$ trajectory.
### 4.3 Three massive adjoint chirals
Our last example is the category of flows believed to be dual to $`𝒩=4`$ super-Yang-Mills deformed by a uniform mass for all three chiral adjoint superfields : $`WW+\frac{m}{2}tr_i\mathrm{\Phi }_i^2`$. Naively, the infrared theory is $`𝒩=1`$ super-QCD. In fact, the model has many vacua, corresponding to non-commutative VEV’s for the massive scalars $`\varphi _i`$ which form a representation of $`SU(2)`$ .
The deformation of the lagrangian consists of a dimension $`3`$ operator, which preserves $`SO(3)`$ and gives a uniform mass to three adjoint fermions; and a dimension $`2`$ operator, which is just the $`SO(6)`$ symmetric mass term for all the scalars. The supergravity scalar dual to the dimension $`3`$ mass term will be denoted $`\phi _m`$. We will also be interested in the dimension three operator $`𝒪_3`$ which includes the $`𝒩=1`$ gaugino bilinear. The dual scalar will be denoted $`\phi _3`$. This is the same scalar which was used in section 4.1 to introduce a mass for one species of fermion. Both $`\phi _m`$ and $`\phi _3`$ are canonically normalized. The potential and the superpotential are
| $`V(\phi _m,\phi _3)`$ | $`={\displaystyle \frac{3}{8}}\left(\mathrm{cosh}^2{\displaystyle \frac{2\phi _m}{\sqrt{3}}}+4\mathrm{cosh}{\displaystyle \frac{2\phi _m}{\sqrt{3}}}\mathrm{cosh}2\phi _3\mathrm{cosh}^22\phi _3+4\right)`$ |
| --- | --- |
| $`W(\phi _m,\phi _3)`$ | $`={\displaystyle \frac{3}{2}}\left(\mathrm{cosh}{\displaystyle \frac{2\phi _m}{\sqrt{3}}}+\mathrm{cosh}2\phi _3\right).`$ |
(46)
A striking feature of this case as compared to the previous two is that there is only one scalar that takes care of the mass deformation. In fact this is entirely appropriate: there is no supergravity scalar which couples linearly to the $`SO(6)`$-invariant boson mass term, $`tr_IX_I^2`$. The scalar $`\phi _m`$ couples linearly to the $`SO(3)`$-invariant fermion mass term, and it can couple quadratically to the $`tr_IX_I^2`$. Demanding supersymmetry guarantees that the boson and fermion mass terms will come in the right ratio.
There is a one-parameter family of gradient flow trajectories for the superpotential in (46), all of which can be displayed in analytic form :
| $`A`$ | $`={\displaystyle \frac{1}{2}}\mathrm{log}\left[2\mathrm{sinh}(rC_1)\right]+{\displaystyle \frac{1}{6}}\mathrm{log}\left[2\mathrm{sinh}(3rC_2)\right]`$ |
| --- | --- |
| $`\phi _m(r)`$ | $`={\displaystyle \frac{\sqrt{3}}{2}}\mathrm{log}{\displaystyle \frac{1+e^{r+C_1}}{1e^{r+C_1}}}`$ |
| $`\phi _3(r)`$ | $`={\displaystyle \frac{1}{2}}\mathrm{log}{\displaystyle \frac{1+e^{3r+C_2}}{1e^{3r+C_2}}}.`$ |
(47)
The behavior $`\phi _m(r)e^r`$ for large $`r`$ indicates that there is a relevant perturbation of the lagrangian. The parameter $`C_1`$ amounts to setting the scale of this mass deformation. On the other hand, $`\phi _3(r)e^{3r}`$ signals a VEV for the operator $`𝒪_3`$. The invariant quantity which parametrizes the trajectories is $`lim_r\mathrm{}\phi _3(r)/\phi _m(r)^3`$. It is related to the VEV of $`𝒪_3`$:
$$\underset{r\mathrm{}}{lim}\frac{\phi _3(r)}{\phi _m(r)^3}=\frac{1}{\sqrt{27}}e^{C_23C_1}=\frac{c_1}{N^2}\frac{𝒪_3}{m^3},$$
(48)
where $`\mathrm{\Lambda }`$ is the dynamical mass scale and $`c_1`$ is a constant of order unity which can be determined once the normalization of $`𝒪_3`$ and $`\mathrm{\Lambda }`$ are unambiguously specified. To understand the factor of $`1/N^2`$ in (48), recall that there is an overall factor of $`1/G_5`$ in front of the classical supergravity action. When this action is used to compute a correlator of boundary operators, three powers of the radius of $`AdS_5`$ combine with $`1/G_5`$ to give an overall factor of $`N^2`$. An example is the stress tensor two-point function: suppressing Lorentz indices, the correlation function that one computes in the classical supergravity approximation is $`T(x)T(0)c/x^8`$ with $`cN^2`$. More generally, the operators to which canonically normalized supergravity scalars couple with strength one (that is, without explicit factors of $`N`$ in the coupling) have $`2`$-point functions which scale as $`N^2`$. An order $`1`$ profile for a supergravity scalar is dual to either an order $`1`$ coefficient for the dual operator added to the lagrangian, or an order $`N^2`$ VEV for that operator. In particular, when $`\phi _3`$ and $`\phi _m`$ of order $`1`$, $`𝒪_3`$ is of order $`N^2`$. Thus for (48) to be consistent with $`c_1`$ of order $`1`$ it is necessary to include the explicit $`1/N^2`$.
On the supergravity side, we can try to use (2) to rule out some trajectories. Contour plots of $`V`$ and $`W`$, the solution set $`𝒫`$, and some typical gradient flow trajectories are shown in figure 6.
We can take advantage of the $`\phi _m\phi _m`$ and $`\phi _3\phi _3`$ symmetries to restrict our attention to flows in the northeast corner of the plots. The condition (2) allows flows with $`C_23C_1`$: the others violate the steepness condition $`\zeta \sqrt{8/3}`$. The curves in $`𝒫`$ give good approximations to gradient flow trajectories in the infrared, but the approximation breaks down more badly than in the previous examples as one proceeds toward the ultraviolet. The saddle points on the vertical axis of (a) are the unstable $`SU(3)`$ critical points.
The field theory is more complicated than in the previous examples, and it is correspondingly more difficult to identify the vacua corresponding to different trajectories in figure 6(b). After the mass deformation, the superpotential in the field theory is
$$W=tr\left(\frac{m}{2}\underset{i=1}{\overset{3}{}}\mathrm{\Phi }_i^2+\mathrm{\Phi }_1[\mathrm{\Phi }_2,\mathrm{\Phi }_3]\right),$$
(49)
and the F-flatness conditions read
$$\frac{W}{\varphi _k}=m\varphi _k+[\varphi _{k+1},\varphi _{k+2}]=0,$$
(50)
where we identify the index $`i`$ on $`\varphi _i`$ modulo $`3`$. Thus the $`\varphi _i`$ should form some $`N`$-dimensional representation of $`SU(2)`$ (not necessarily irreducible). Imposing the D-flatness conditions and identifying gauge-equivalent configurations can be achieved simultaneously by modding out by the action of the complexified gauge group. Roughly speaking, one can think of the representations of $`SU(2)`$ as non-commutative spheres, whose radii are the quadratic Casimirs of the representations. The radius is largest for the $`N`$-dimensional irreducible representation.<sup>6</sup><sup>6</sup>6Similar non-commutative configurations were studied in the test-brane approximation in . In fact, one can see a hint of the physics of (50) by putting D3-branes in a constant background field $`G_{IJK}`$, where $`G_{(3)}`$ is the complex three-form field strength of type IIB supergravity, and $`IJK`$ are indices in the directions perpendicular to the D3-branes. $`G_{(3)}`$ has to be involved in the ten-dimensional version of the solutions specified in (47) (indeed in the ten-dimensional version of all the solutions presented in this section with the exception of the Coulomb branch trajectories of the undeformed $`𝒩=4`$ theory). If $`_6G_{(3)}=iG_{(3)}`$, where $`_6`$ denotes the Hodge dual in the $`IJK`$ directions, then in the approach of there remain commutative flat directions for the D3-branes to spread out in. Possibly an improved treatment with non-constant $`G_{(3)}`$ could approximate the physics of (50) more closely. In this case,
$$tr\varphi _1[\varphi _2,\varphi _3]\frac{N^3}{12}m^3.$$
(51)
Going to the opposite extreme, one could satisfy (50) trivially by setting all the $`\varphi _i`$ to $`0`$. In this case, the gauge group is completely unbroken, and there is a gaugino condensate, which according to the standard analysis is
$$tr\lambda _4\lambda _4c_2e^{2\pi ik/N}N\mathrm{\Lambda }^3,$$
(52)
where $`c_2`$ is another constant of order unity, $`\mathrm{\Lambda }`$ is the scale of the low-energy $`𝒩=1`$ theory, and $`k`$ can run from $`1`$ to $`N`$ (see for example ). Again up to factors of order unity, $`\mathrm{\Lambda }=m\mathrm{exp}\left(\frac{8\pi ^2}{3g^2N}\right)`$ where $`g`$ is the gauge coupling at the $`𝒩=4`$ ultraviolet fixed point. The overall factor of $`N`$ in (52) ensures that domain walls between adjacent vacua have tension scaling as $`N`$. To see this recall that the exact superpotential has one more power of $`N`$ than the gaugino condensate; but domain wall tensions are given by the change of $`W`$ across the domain wall, and $`W`$ changes by a phase $`e^{2\pi i/N}`$ from one vacuum to the next .
Naively, the operator $`𝒪_3`$ is $`tr\left(\lambda _4\lambda _4+\varphi _1[\varphi _2,\varphi _3]\right)+h.c.`$, as in (37). Both the fermion bilinear and the scalar trilinear are manifestly $`SU(3)`$ invariant. The VEV of the first term in the trivial vacuum is order $`N`$, while the VEV of the second is order $`N^3`$ in the irreducible vacuum. Classical supergravity effects correspond to VEV’s of order $`N^2`$, so the first behavior appears unobservably small, while the second seems unobservably large! We suggest however that
$$𝒪_3=tr\left(\lambda _4\lambda _4+\underset{i=1}{\overset{3}{}}\varphi _i\frac{W}{\varphi _i}\right)+h.c.,$$
(53)
where $`W`$ is the deformed superpotential, (49), and we are not attempting to be precise about factors of order unity on either the fermion or scalar terms. Classically, the scalar terms have vanishing VEV in any vacuum. However, the Konishi relation gives
$$\underset{i=1}{\overset{3}{}}tr\varphi _i\frac{W}{\varphi _i}=Ntr\lambda _4\lambda _4.$$
(54)
Thus if $`tr\lambda _4\lambda _4`$ is order $`N`$, as the standard field theory analysis indicates, then $`𝒪_3`$ is order $`N^2`$.
Because the standard field theory analysis of the trivial confining vacua indicates a VEV for $`𝒪_3`$ which is of the right order to be observed as a classical effect in supergravity, the natural guess is that these vacua (or one of them, say the $`k=0`$ vacuum) should be identified as the $`C_2=3C_1`$ trajectory in figure 6(b). This trajectory is distinguished in that it is the last one permitted by the condition (2). It has a smaller value of $`\zeta `$ than all the other trajectories, and it is asymptotic to a curve in $`𝒫`$, which makes it the most plausible candidate for a flow with near-extremal generalizations which do not involve deforming the lagrangian.
The other trajectories permitted by (2) could have one of several interpretations, and without further investigation we cannot judge which interpretation is correct. First, they might be in some way unstable, and so correspond to no well-defined field theory vacuum. Second, they might be massless vacua corresponding to embeddings of $`SU(2)`$ in $`SU(N)`$ which leave factors of $`U(1)`$ unbroken. Third, they might be massive vacua where $`SU(N)`$ is broken down to some $`SU(N/p)`$ subgroup by a representation of $`SU(2)`$ which can be written as $`N/p`$ blocks of dimension $`p\times p`$. A recent analysis indicates yet one more possibility. The claim of is that the standard analysis (52) of the gaugino condensate is dramatically altered at strong coupling, and that in the large $`N`$, large $`g_{YM}^2N`$ limit the regular $`N`$-gon of vacua for the trivial representation of $`SU(2)`$ degenerates to points on a line segment. These points are not evenly distributed; rather,
$$tr\lambda _4\lambda _4=Nm^3/j^2\text{where }j=1,2,3,\mathrm{}.$$
(55)
We have rescaled $`m`$ to conform with our conventions, set $`g_{YM}=1`$, and dropped various factors of order unity. Since $`tr\lambda _4\lambda _4`$ still scales as $`N`$, the previous arguments that $`𝒪_3`$ should be observable as a classical effect in supergravity still hold, and we might still expect the $`C_2=3C_1`$ trajectory to correspond to the $`j=1`$ vacuum (which is the large $`\lambda `$ limit of the $`k=0`$ vacuum in the standard analysis (52)). Then some of the $`C_2<3C_1`$ trajectories might be $`j>1`$ vacua.
The trajectories with $`C_2>3C_1`$ violate (2) and run asymptotically parallel to the $`\phi _3`$-only trajectory. This trajectory seems to be a particularly clear case of an unphysical geometry: the lagrangian is undeformed, a gaugino condensate has no reason to form, the flat directions are $`[X_I,X_J]=0`$, and yet we are trying to give a nonzero VEV to $`𝒪_3=tr(\lambda _4\lambda _4+\varphi _1[\varphi _2,\varphi _3])`$. By extension it is plausible to rule out all the $`C_2>3C_1`$ trajectories, since they have the same behavior in the infrared. Incidentally, the “universal solution” of Horava-Witten theory compactified on a Calabi-Yau manifold develops the same singularity as the $`C_2>3C_1`$ flows if the negative tension “hidden sector” end-of-the-world brane is taken to infinite redshift. The low-energy limit of the hidden sector gauge theory is pure $`𝒩=1`$ super-Yang-Mills theory. The unpleasant singularity could be due to neglect of the gaugino condensate in the derivation of this solution.
To sum up, what tells in favor of the condition (2) in this example is that the field theory analysis indicates VEV’s for $`𝒪_3`$ which are of order $`N^2`$ and bounded in magnitude. This qualitative feature is reproduced by (2) because it implies $`C_23C_1`$. It would be interesting to study the symmetries of the problem more carefully. In particular, is there a $`U(1)`$ symmetry in the supergravity, broken by the choice of $`\phi _3`$, which arises from the large $`N`$ limit of the discrete $`𝐙_{2N}`$ symmetry? According to this symmetry only pertains in the small $`g_{YM}^2N`$ limit.
## 5 Fluctuations
Obvious cases where the criterion (2) fails to guarantee a physical interpretation on the field theory side are flows to critical points of $`V`$ which violate the Breitenlohner-Freedman bound . Two well-known examples are flows to the $`SO(5)`$ symmetric critical point of $`d=5`$ $`𝒩=8`$ gauged supergravity, and flows to the $`SU(3)`$ critical point . We would like to have a generalization of the Breitenlohner-Freedman bound which applies to cases which are only asymptotically $`AdS_5`$ in the ultraviolet. The obvious one is to demand that the spectral decomposition of the Lorentzian two-point function of an arbitrary color singlet operator should involve only timelike momenta. This statement can be phrased in the jargon of “AdS/QCD” as the absence of tachyonic glueballs. More generally we can think of it as stability of the field theory vacuum: a gauge singlet operator $`𝒪(0)`$, acting at the origin of position space on the vacuum of the dual gauge theory, should produce only states with timelike momenta. On Wick rotating to Euclidean signature, the requirement amounts to having two-point functions which are non-oscillatory at large distances, but instead fall off like powers or exponentials. This makes it clear that the criterion is indeed a generalization of the Breitenlohner-Freedman bound, since violations of that bound translate precisely into gauge singlet operators with complex dimension.
The restriction to timelike momenta is straightforward to implement in AdS/CFT because the spectrum of momenta-squared is identical to the spectrum of an appropriately defined one-dimensional Schrodinger hamiltonian. This identification is a direct consequence of the prescription of , and it is the basis of all AdS-glueball calculations following . Since it has been discussed, implicitly or explicitly, in many other places in the literature (for instance ) we can afford to be brief. Our aim is simply to make the discussion of more systematic.
Consider a Poincaré invariant flow solution of the form (1). It is convenient to introduce a new radial variable, $`z`$, such that $`dr=e^Adz`$. The solution assumes the form
| $`ds^2`$ | $`=e^{2A(z)}\left(dt^2+d\stackrel{}{x}^2+dz^2\right)`$ |
| --- | --- |
| $`\stackrel{}{\phi }`$ | $`=\stackrel{}{\phi }(z).`$ |
(56)
Linearizing the equations of motion around this solution is tricky because linear fluctuations of the graviton couple to linear fluctuations of the scalars involved in the flow. However, scalar fluctuations $`\delta \stackrel{}{\phi }(x^\mu ,z)`$ satisfying $`\delta \stackrel{}{\phi }(x^\mu ,z)\stackrel{}{\phi }^{}(z)`$ nearly decouple from the graviton. That decoupling would be exact if $`\stackrel{}{\phi }^{}(z)`$ didn’t depend on $`z`$, but we may assume that $`\stackrel{}{\phi }^{}`$ changes direction only slowly in the deep infrared as $`A(z)\mathrm{}`$. Since we are mainly interested in discerning infrared properties, little is lost in the approximation that the orthogonal $`\delta \stackrel{}{\phi }(x^\mu ,z)`$ excitations decouple from the graviton. Linearizing the wave equation $`\text{ }\text{ }\text{ }\text{ }\text{ }\stackrel{}{\phi }=V/\stackrel{}{\phi }`$ leads to
$$\text{ }\text{ }\text{ }\text{ }\text{ }\delta \stackrel{}{\phi }=\frac{^2V}{\stackrel{}{\phi }\stackrel{}{\phi }}|_{\stackrel{}{\phi }(z)}\delta \stackrel{}{\phi }.$$
(57)
We can solve via separation of variables: setting $`\delta \stackrel{}{\phi }=e^{\frac{3}{2}A(z)}e^{ik_\mu x^\mu }\stackrel{}{\psi }(z)`$, one obtains directly
$$\left[_z^2+U(z)\right]\stackrel{}{\psi }=k^2\stackrel{}{\psi },$$
(58)
where
$$U(z)=\frac{3}{2}A^{\prime \prime }(z)+\frac{9}{4}A^{}(z)^2+e^{2A(z)}\frac{^2V}{\stackrel{}{\phi }\stackrel{}{\phi }}|_{\stackrel{}{\phi }(z)}.$$
(59)
$`U(z)`$ is a matrix acting in the space of scalars orthogonal to the flow. The Schrodinger problem (58) determines the spectrum of $`k^2`$, and in our mostly plus conventions the requirement of timelike momenta is $`k^20`$. There are two caveats to be kept in mind. First, we used the approximation of decoupling from the graviton to derive (58). Second, if the flow geometry has a curvature singularity, supergravity gives at best an approximation to the spectrum. The curvature singularity indicates an infrared problem, so in physical cases it should be possible to smooth it out—for instance by going to finite temperature and Wick rotating so that the manifold is entirely non-singular. If the smoothing is done in the far infrared, it should only affect the behavior of $`U(z)`$ near the radius where $`A(z)\mathrm{}`$. One can hope that both caveats change the spectrum only in a controllably small way: for instance, they might make the infimum of the spectrum slightly negative rather than strictly zero.
Let us now check that we recover the Breitenlohner-Freedman bound for AdS geometries. Here $`e^{A(z)}=L/z`$, no scalars are excited, and curvatures are bounded, so we can impose the sharp inequality
$$\left[_z^2+U(z)\right]=\left[_z^2+\frac{15/4+m^2L^2}{z^2}\right]0.$$
(60)
This leads to $`m^2L^24`$ by a standard Bessel function analysis.
For geometries with naked singularities, the positivity of the Schrodinger operator in (58) is essentially a condition proposed in almost twenty years ago, and also considered in .<sup>7</sup><sup>7</sup>7I thank G. Horowitz for a discussion of and of the basic criterion. With trivial adaptations appropriate to the current context, the contents of can be summarized as follows. The Schrodinger operator on the left hand side of (58) is the radial part of a linearized wave operator in a static background with a naked singularity. Let us use $`𝒜`$ to denote the Schrodinger operator. $`𝒜`$ is symmetric with respect to an appropriate inner product $`(,)`$ for the wave-functions. If it is also positive, then there is a natural extension (the Fredholm extension) of $`𝒜`$ to a self-adjoint operator on the Hilbert space of functions which is the closure under the norm $`(\stackrel{}{\psi },\stackrel{}{\psi })+(\stackrel{}{\psi },𝒜\stackrel{}{\psi })`$ of smooth functions whose support excludes both the naked singularity and the boundary of $`AdS_5`$. The existence of a self-adjoint extension of $`𝒜`$ is a natural criterion because it means that it is possible to define a unitary evolution equation for linearized fluctuations around the static background. The crucial step is positivity of $`𝒜`$, and in AdS/CFT this is equivalent to the stability of the field theory vacuum. The bulk spacetime analysis looks like it could proceed for $`𝒜`$ bounded below (not necessarily by zero), provided one is willing to introduce a sufficiently large $`(\stackrel{}{\psi },\stackrel{}{\psi })`$ component in defining the norm. But, modulo the caveats mentioned above, the natural AdS/CFT requirement is positivity.
Suppose there is a naked singularity in the geometry at $`z=z_0`$. Let $`u(z)`$ denote the smallest (i.e. most negative) eigenvalue of $`U(z)`$. Then through almost the same Bessel function analysis that one uses to translate (60) into the Breitenlohner-Freedman bound, one finds that
$$\underset{zz_0}{lim\; inf}(zz_0)^2u(z)\frac{1}{4}$$
(61)
is a necessary condition for $`𝒜`$ to be bounded below. (This is almost true; to be rigorous one must exclude the possibility that $`u(z)`$ oscillates rapidly).
The conditions we have discussed should generalize to supergravity fields with spin, and to other dimensions. It was only to avoid technical difficulties that we restricted our attention to scalar fluctuations perpendicular to the flow. The general problem of diagonalizing the linearized fluctuations of all supergravity fields is computationally challenging, but the final requirement is that that the two-point functions extracted from the diagonalized equations should involve only timelike momenta in their spectral decomposition.
## 6 Discussion
Applications of AdS/CFT to non-conformal theories appears at present to be a subject of particulars, with only a few general truths. The c-theorem is one such truth ; the possibility of replacing second order equations of motion with first order gradient flow equations is another . We have attempted to find a third, namely a criterion for what sorts of singular behavior are allowed far from the boundary of AdS. Our conjectured criterion, (2), fares reasonably well when confronted by examples in the literature. With a few exceptions, it correctly distinguishes pathological from non-pathological. The exceptions are flows to critical points which violate the Breitenlohner-Freedman bound. We have suggested that the Breitenlohner-Freedman bound is a special case of a more general requirement, namely that the spectral decomposition of two-point functions in AdS/CFT should involve only timelike momenta. The methods proposed for ruling out unphysical singularities—namely, (2) and well-behaved two-point functions—are applicable independent of supersymmetry. The examples considered in section 4 all preserve some fraction of supersymmetry, but this is only in order to have a clear understanding of the dual field theory.
The flows to critical points violating the Breitenlohner-Freedman bound are pathological cases which (2) fails to rule out. In the current literature, there are no clean exceptions to (2) in the other direction—that is, clearly physical solutions which violate (2). It would be very interesting if one could find such an exception, and to ask how and whether it supports finite temperature. In view of the examples in this paper, one’s first question when addressing a putative exception must be, “Is the theory in a physical vacuum state?”
Ideally, if the AdS/CFT map were perfectly understood, we would be able to say precisely which singularities are resolved by non-trivial infrared physics and which are not. A complete understanding of this question might lead to a complicated set of constraints on the five-dimensional bulk. The condition (2) is a simple semi-empirical rule that captures some non-trivial aspects of these constraints. A particularly clear example is a VEV for the $`SU(3)`$ singlet operator $`𝒪_3=tr\left(\lambda _4\lambda _4+\varphi _1[\varphi _2,\varphi _3]\right)`$. On the field theory side, such a VEV is forbidden unless the $`𝒩=4`$ lagrangian is deformed. The second term of $`𝒪_3`$ can’t have a VEV because the $`𝒩=4`$ effective potential is flat only along directions where the $`\varphi _i`$ commute. The Konishi anomaly relates the VEV for the first and second terms of $`𝒪_3`$, so that the vanishing of one implies the vanishing of the other. In supergravity, one is certainly free to perturb $`AdS_5`$ by a scalar profile indicating a VEV for $`𝒪_3`$, but at the expense of violating (2). The second order bulk equations have too many solutions if arbitrary singularities are allowed; (2) is a useful tool for weeding out the good from the bad.
What precisely is bad about the field theory dual to singularities that violate (2) has to be addressed case by case. We have already indicated the problem for $`𝒪_30`$ without a mass deformation of the $`𝒩=4`$ lagrangian. In sections 4.1 and 4.2, the problem with trajectories violating (2) was that positive definite operators have negative VEV’s. In section 3.2, (2) ruled out multi-center distributions of D3-branes which involved “ghosts” with negative tension and negative charge. If one formally uses these “ghostly” distributions to compute VEV’s of the form $`trX_{(I_1}\mathrm{}X_{I_{\mathrm{}})}`$, the results again violate inequalities among these VEV’s that one can prove using the hermiticity of the fields $`X_I`$.
It would be interesting to describe more fully the subset $`𝒮E_{6(6)}/USp(8)`$ on which the scalar potential is less than its value at the origin of the coset. $`𝒮`$ is far from being a uniform blob or half-space: it has tendrils which come arbitrarily close to subspaces of positive co-dimension as one proceeds further and further from the origin of the coset. We have proven that a static black hole horizon can only form at a location where the scalars lie in $`𝒮`$. The conjecture (2) essentially says that curvature singularities are allowed only if the scalars remain in $`𝒮`$ as one approaches the singularity. (Technically, we have not ruled out situations where $`|\stackrel{}{\phi }|\mathrm{}`$ and $`V(\stackrel{}{\phi })`$ approaches a constant value which is larger than $`V`$ at the origin of the coset. No example in the current literature has this feature, and it may be excludable on general grounds.)
A black hole horizon is the only purely geometrical expression of finite temperature. Black holes are also in the generic expression of finite temperature in a theory with gravity, in the sense that thermal excitations with enough energy should form a black hole. It sounds reasonable to claim that any finite temperature applied to a bulk geometry with a singularity where $`g_{tt}0`$ will result in a black hole horizon cloaking the singularity, simply because the proper temperature diverges near the singularity.<sup>8</sup><sup>8</sup>8This line of argument was suggested to me by L. Susskind. But it might happen that the horizon forms in a region very close to the singularity where curvatures are Planckian. Then low-energy supergravity will have limited value; indeed the very notion of horizon might have to be modified.
In an AdS/CFT context, it is straightforward to estimate when the five-dimensional supergravity treatment breaks down. In units where the radius of curvature of $`AdS_5`$ is unity, the five-dimensional Planck length scales as $`N^{2/3}`$. In the case where the dual CFT is $`𝒩=4`$ super-Yang-Mills theory, $`N`$ is the rank of the gauge group; more generally, $`N^2`$ is the central charge of the dual CFT. Curvature invariants near a singularity of the form (29) become Planckian when $`rr_0\begin{array}{c}<\hfill \\ \hfill \end{array}N^{2/3}`$. What is the entropy of a black hole horizon which forms just before curvatures become Planckian? The horizon area per unit volume at $`r=r_H`$ is $`e^{3A(r_H)}`$. The entropy formula is actually more complicated than $`S=A/4G_5`$ when higher derivative terms become important , but to obtain a rough-and-ready estimate we will stick with the simple area law. (For near-extremal D3-branes with no evolving scalars, the evidence of AdS/CFT is that all the $`\alpha ^{}`$ corrections only change the area law by $`4/3`$ in the $`\alpha ^{}/L^2\mathrm{}`$ limit . However a “phase transition” at finite $`\alpha ^{}/L^2`$ cannot be ruled out .) In units where $`L=1`$, we have $`1/G_5=N^2`$, so the entropy per unit volume is roughly $`N^2e^{3A(r_H)}`$. Assuming that near-extremal solutions have nearly the same $`A(r_H)`$ as their Poincaré invariant limit, one finds that the entropy per unit volume for a horizon at nearly Planckian curvatures scales as $`N^{2\frac{4}{3\zeta ^2}}`$. This is to be compared with the expected scaling on the field theory side. If the field theory confines, then at low energy there are $`O(1)`$ degrees of freedom. A black hole horizon should have an $`O(1)`$ entropy if it purports to describe the theory at a temperature much smaller than the confinement scale. This is possible if $`\zeta <\sqrt{2/3}`$, but precisely for this range there is no gap for excitations of a minimal scalar! We conclude that a five-dimensional supergravity treatment of near-extremal generalizations of singularities with $`\zeta <\sqrt{2/3}`$ should suffice to describe gapless field theory duals down to some minimum temperature which scales as an inverse power of $`N`$, but that the supergravity approximation is insufficient to describe the confining phase of a gauge theory. The original proposal of avoids this problem by arriving at a confining gauge theory as the low-energy limit of a higher-dimensional gauge theory at finite temperature. However one is then left with a description where above the confinement scale the theory is higher dimensional.
We would like to think that the considerations of this paper will generalize to a string theoretic setting; that the conjecture (2) will carry over naturally to some more refined restriction on the nature of highly curved solutions; and that string theory considerations will provide a detailed “resolution” of the supergravity singularities which are allowed. In particular, D3-branes in an $`H`$-field should resolve the singularities found in section 4.3 (entries (F) and (G) in figure 3). With the supergravity solutions and the brane resolutions both in hand, one might make an interesting study of the transition between the deconfined phase (described by a black hole horizon in the supergravity solution) and the confined phase (described in terms of the brane resolution). It is conceivable that the transition will involve some dramatic alteration of the supergravity solution which we are simply not equipped to compute without a full understanding of the brane resolution. But we prefer the hopeful view that the low-energy theory, useful in so many ways as a guide to the qualitative features of the full string theory, will match more or less smoothly onto the microscopic description.
Matching a given brane resolution back onto supergravity seems sure to impose definite boundary conditions on the bulk fields near the singularity, but the process of finding these resolutions and extracting their supergravity limits could be tedious. Black hole solutions of the form (6) provide a quick and dirty route to establishing the boundary conditions at the singularity without going to the trouble of finding the brane resolution—provided one is interested in physics that is robust under finite temperature, and provided one is willing to grant that the low-energy supergravity description of finite temperature should match onto a microscopic description smoothly enough to distinguish good singularities from bad.
Here is the strategy. Suppose we have $`n`$ tachyonic scalars, corresponding to $`n`$ relevant operators that one might add to the microscopic lagrangian. Fix the more singular behavior of all these tachyons near the boundary of $`AdS_5`$: this amounts to fixing the relevant deformations, but not the VEV’s. Impose finite temperature in the form of a black hole horizon. The horizon boundary conditions (12) amount to $`n`$ more constraints on the solution. Assuming the form (9) for $`A`$ and $`h`$ (with $`B`$ left free, to be determined in terms of the temperature) amounts to using residual coordinate invariance to fix one boundary condition on $`A`$ and one on $`h`$ as $`r\mathrm{}`$. In total, if $`T`$ is regarded as fixed, there are $`2n+3`$ constraints. There are $`2n+4`$ integration constants of integration in the bulk equations of motion, (7), but the zero-energy constraint (8) fixes one. So the boundary value problem is exactly determined: generically there are solutions with a given temperature $`T`$, but only discretely many.<sup>9</sup><sup>9</sup>9At low temperatures it may become advantageous to use the non-extremality parameter $`B`$ as a control variable instead of temperature, since there are probably cases where $`B`$ is a multiple valued function of $`T`$ but not vice versa. $`B`$ scales as energy density above extremality.
As $`T`$ is lowered, there may be phase transitions, both first and second order, corresponding to discontinuities or bifurcations in the space of all finite temperature solutions. Sorting out such finite-temperature phase transitions would be a fascinating end in itself, which we have only delved into superficially with the prediction that $`\alpha =1/2`$ at a generic bifurcation. But let us suppose we have gotten past all the phase transitions, if necessary by switching our control parameter from $`T`$ to $`B`$ (see the footnote). As $`B0`$, the horizon retreats into the singularity, but the scalars’ evolution is still perfectly determined. In short, black holes of a given temperature impose just enough boundary conditions to be unique up to discrete choices, and singular Poincaré solutions which are limits of black holes inherit the same property.
In the case of a single scalar $`\mu `$, we can offer a definite conjecture regarding the nature of the condition at the singularity. Assume that $`\mu `$ becomes large and positive in solutions with small $`T`$ (or, more properly, small $`B`$), and that $`Ve^{\eta \mu }`$ as $`\mu \mathrm{}`$ for some $`\eta <\sqrt{32/3}`$. Then the Poincaré limit of black hole solutions will have the form (29) with $`\zeta =\eta /2`$. This is a very restrictive condition: the generic Poincaré invariant solution is of the form (29) with $`\zeta =\sqrt{8/3}`$. Requiring the exceptional $`\zeta =\eta /2`$ behavior amounts to imposing one boundary condition at the curvature singularity. We conjecture that Poincaré invariant solutions of this kind and only of this kind will have near-extremal generalizations. This is a much stronger statement than (2), and the evidence is slimmer. It could be phrased as a weak form of Cosmic Censorship: static nakedly singular solutions are allowed, but a theory with any thermalizable excitations will seek out the singularity with the smallest possible $`\zeta `$. The smallest possible $`\zeta `$ will always be less than $`\sqrt{8/3}`$ for potentials of the form (4).
Unfortunately, there is no solution involving only one scalar where we have a really clean understanding of the physics on the dual field theory side. However, we have run numerics to check the “Weak Cosmic Censorship” conjecture of the previous paragraph for two interesting cases. The first case is a deformation
$$+\frac{1}{2}m^2tr\left(X_1^2X_2^2X_3^2X_4^2+2X_5^2+2X_6^2\right),$$
(62)
with $`m^2>0`$. The corresponding $`V`$ goes as $`e^{\sqrt{2/3}\mu }`$ as $`\mu \mathrm{}`$ (see ). Black hole solutions appear to exist for any $`T`$, and as $`T0`$ one does see a scaling region where the form (29) indeed appears with $`\zeta =\sqrt{1/6}`$ to $`0.05\%`$ accuracy. Not unexpectedly, the VEV for $`𝒪_2=tr\left(X_1^2X_2^2X_3^2X_4^2+2X_5^2+2X_6^2\right)`$ diverges as $`T0`$. The second case is a deformation by a uniform mass for all three chiral adjoints—the case studied in section 4.3—but with only the scalar dual to the mass deformation excited, and not the scalar dual to the gaugino condensate. The scalar evolution is along the horizontal axis of figure 6(b). One might expect the Poincaré invariant solution to be only metastable, with a perturbative instability toward forming a gaugino condensate. In fact the normal modes of linear fluctuations of the supergravity scalar dual to the condensate do not include tachyons: the Schrodinger equation (58) in that case has the form of supersymmetric quantum mechanics with $`𝒬=_z+3\mathrm{cot}2z`$, and the supersymmetric ground state corresponds to shifting to a neighboring trajectory in the plane of figure 6(b). Provisionally let us allow this case as a geometry which satisfies all the constraints we were able to put on solutions which should have field theory duals, even though the correct vacuum state on the field theory side remains obscure. At the least the solution is a candidate for verifying the conjectures of the previous paragraph, because $`V(\mu )e^{\sqrt{16/3}\mu }`$ and the Poincaré invariant, supersymmetric solution is of the form (29) with $`\zeta =\sqrt{4/3}`$. The numerics is somewhat stiffer in this case, but small $`T`$ solutions reveal scaling regions where $`\zeta =\sqrt{4/3}`$ can be verified to $`3\%`$ accuracy.
In light of these numerical results, we feel entitled to a speculation regarding the marginal case $`\zeta =\sqrt{8/3}`$. This is in some sense the generic case, because if $`Ve^{\eta \mu }`$ with $`\eta <\sqrt{32/3}`$, the generic Poincaré invariant solution is of the form (29) with $`\zeta =\sqrt{8/3}`$, whereas if $`Ve^{\eta \mu }`$ with $`\eta >\sqrt{32/3}`$, the condition (2) rules out all solutions. In the supersymmetric examples in this paper, we encountered $`\zeta =\sqrt{8/3}`$ twice: once in a Coulomb branch state of $`𝒩=4`$ gauge theory where the D3-branes were arranged in a perfect $`S^3`$ shell in the transverse dimensions; and again in a Coulomb branch state of a $`𝒩=1`$ mass deformation which had essentially the same $`S^3`$ shell interpretation. In both cases, the field theory dual makes it clear that the vacuum won’t support finite temperature. If one does turn on finite temperature at the same time as introducing a term in the lagrangian of the form (62), then the trajectory may be nearly the same in the space of scalars, but the dependence on proper distance will be very different near the singularity: $`\zeta =\sqrt{1/6}`$ in the infrared, rather than $`\sqrt{8/3}`$. It is natural to speculate that the generic solutions with $`\zeta =\sqrt{8/3}`$ are allowed precisely when they correspond to exploring flat directions in the dual field theory, and that they can never support finite temperature. This unfortunately tells against some supergravity constructions which have been claimed to exhibit properties of confinement , but the case for confinement in those examples was less than airtight.
It is particularly easy to show that the geometries of , involving only the metric and the dilaton, cannot support a black hole horizon. The scalar equation of motion plus the horizon boundary conditions imply that the dilaton is everywhere constant. The only static black hole geometry involving only the metric and the dilaton is AdS-Schwarzschild. This simple argument is another clue that singularities with $`\zeta =\sqrt{8/3}`$ cannot support finite $`T`$: they have the same scaling as the solutions of .
A feature of AdS/CFT which was essential to the proper interpretations of the examples in this paper is that a profile for a given scalar field can indicate either a deformation of the lagrangian or a VEV for some gauge singlet operator. This dual role of scalar profiles complicates any attempt to identify the value of a scalar at a given radius in $`AdS_5`$ with a coupling in an effective lagrangian for the dynamics at the corresponding scale. We must ask: 1) How do we disentangle the “VEV” part of the scalar profile from the “deformation” part? 2) Is such a disentanglement really necessary, or can the effective lagrangian be defined so as to pick out exactly the combination of VEV and deformation that AdS/CFT prescribes for its effective coupling? The issues might become clearer if one focused on solutions which are limits of black holes with prescribed behaviors for the scalar profiles corresponding to deformations of the microscopic lagrangian. The nearly unique specification of vacuum state in these solutions may simplify the interpretation of the holographic renormalization group proposed in . For instance, the identification of beta functions with gradients of scalars seems odd in Coulomb branch states of $`𝒩=4`$ gauge theory, but possibly more natural for mass deformations in some preferred vacuum. A true acid test for the subject would be to reproduce some known quantitative features of the field theory RG, such as the NSVZ exact beta function.
## Acknowledgements
This work would not have been possible without extensive communications with N. Warner, K. Pilch, and D. Freedman. I have also profited greatly from a number of exchanges with R. Myers regarding near-extremal geometries, and with H. Verlinde regarding the renormalization group in AdS/CFT. I also thank T. Banks, G. Horowitz, S. Kachru, I. Klebanov, B. Kol, V. Periwal, B. Pioline, J. Polchinski, S. Shenker, E. Silverstein, M. Strassler, L. Susskind, and E. Witten for useful discussions, and R. Myers for comments on an early draft. This research was supported in part by DOE grant DE-FG02-91ER40671 and by the Harvard Society of Fellows.
## Appendix
The purpose of this appendix is to remark on brane-world models in which there is a curvature singularity parallel to the brane on which visible sector matter exists, at some finite proper distance from it. Such constructions were first considered in the context of Horava-Witten theory in . They have also appeared as a variant of the Randall-Sundrum construction (see for instance ). Following in the spirit of , it was recently proposed by two Stanford groups that such geometries could help solve the cosmological constant problem. In this appendix we will consider only the minimal case where the “Planck brane” is at one end of the spacetime, where $`g_{tt}`$ is finite, and the singularity is at the other, where $`g_{tt}0`$. In the simplest constructions, where $`g_{tt}`$ is monotonic, gravity is in some sense localized on the Planck brane (hence the name). At least in most Horava-Witten theory constructions, the visible sector fields live on the Planck brane. The singularity then represents (or is resolved by) hidden sector fields which couple only gravitationally to the visible sector.
A crucial issue, both in the general scheme of and in the specific proposal of , is the boundary conditions at the singularity. At the level of classical five-dimensional gravity, the construction of is a boundary value problem, where one boundary is the Planck brane and the other boundary is a naked singularity. Similar analyses have appeared in many places, for example . In , spacetime ends on some orbifold fixed plane before $`g_{tt}0`$, and boundary conditions are imposed both there and on the Planck brane. By contrast, in , no boundary conditions are imposed at the singularity where $`g_{tt}0`$. In (and other similar treatments) the boundary value problem is sufficiently determined to fix the four-dimensional cosmological constant, and there is no reason for it to be small. In , there are fewer boundary conditions. The boundary value problem is underdetermined if the four-dimensional cosmological constant is left unspecified; it becomes determined (up to an multiplicative shift on the warp factor) if one specifies the four-dimensional cosmological constant. In particular, there do exist solutions with $`3+1`$-dimensional Poincaré invariance.
The issue which falls within the purview of this paper is whether the free boundary conditions used in are reasonable. Because the solutions there purport to be at least a toy model for real world physics, they must be able to support finite temperature. We will interpret this as synonymous with being able to form a black horizon with finite Hawking temperature close to the singularity, without drastically changing the rest of the geometry. (Some potential pitfalls of this interpretation have been discussed in section 6). It is straightforward to show that a horizon is impossible in classical gravity if the bulk scalar is free, unless in fact the bulk scalar does not vary at all in the solution. Essentially the relevant observation has already been made in section 6, but we will go into slightly more detail here to make the case clear. In order to have finite temperature, we seek a generalization of the Poincaré invariant solution of the form
| $`ds^2`$ | $`=e^{2A(r)}(h(r)dt^2+d\stackrel{}{x}^2)+dr^2/h(r)`$ |
| --- | --- |
| $`\phi `$ | $`=\phi (r)`$ |
(63)
where $`h(r)`$ is a function which has a simple zero at $`r=r_H`$, the black hole horizon. The temperature is related to $`h^{}(r_H)`$ (see (10)). A sketch of the geometry is presented in figure 7.
The scalar equation of motion is (by assumption) $`\text{ }\text{ }\text{ }\text{ }\text{ }\phi =0`$: explicitly,
$$e^{4A}_re^{4A}h_r\phi =0.$$
(64)
The geometry and the scalars should be perfectly regular at the horizon: an infalling observer must not notice anything special as he crosses through it. Thus in particular $`\phi ^{}`$ and $`e^{4A}`$ are finite at the horizon (as usual, primes denote derivatives with respect to $`r`$). So $`e^{4A}h\phi ^{}=0`$ at $`r=r_H`$. In view of (64), $`e^{4A}h\phi ^{}=0`$ everywhere. This is only possibility if $`\phi `$ is in fact constant. Thus the original singular geometry of , supported by a free scalar in the bulk, cannot be recovered as a limit of finite temperature black hole solutions. In , an intuitive picture of a “mechanism” by which the four-dimensional cosmological constant could be tuned to zero was presented in terms of a conserved current associated with the shift symmetry $`\varphi \varphi +\delta \varphi `$ which. The current in some sense carries off visible sector vacuum energy into the bulk. This is an appealing picture, but in a geometry of the form (63) the current has the form $`J^r=e^{4A}h_r\phi `$, which vanishes at the horizon and therefore throughout the geometry.<sup>10</sup><sup>10</sup>10I thank M. Peskin and R. Sundrum for comments on this point.
We now wish to contemplate some variations on the basic Stanford proposal, still in the general framework of . First, it is straightforward to work with any number of scalars with any bulk potential $`V(\stackrel{}{\phi })`$ which can be written in the form
$$V(\stackrel{}{\phi })=\frac{1}{8}\left(\frac{W}{\stackrel{}{\phi }}\right)^2\frac{1}{3}W(\stackrel{}{\phi })^2.$$
(65)
Then a bulk geometry with $`3+1`$-dimensional Poincaré invariance can be generated from the first order equations (5), and it is possible that the naked singularities that generically arise can be obtained as limits of finite temperature black hole solutions provided $`V`$ is bounded above. Let us assume a Planck brane action of the form
$$S_{\mathrm{brane}}=d^4\xi \sqrt{g_{(\mathrm{induced})}}\lambda _{\mathrm{Pl}}(\stackrel{}{\phi }).$$
(66)
The equations of motion that determine the possible embeddings of the Planck brane in the bulk geometry are
| $`\theta _{ij}\theta g_{ij}^{(\mathrm{induced})}`$ | $`=\lambda _{\mathrm{Pl}}(\stackrel{}{\phi })g_{ij}^{(\mathrm{induced})}`$ |
| --- | --- |
| $`_n\stackrel{}{\phi }`$ | $`={\displaystyle \frac{1}{2}}{\displaystyle \frac{\lambda _{\mathrm{Pl}}}{\stackrel{}{\phi }}}.`$ |
(67)
Here $`\theta _{ij}`$ is the extrinsic curvature of the Planck brane, $`\theta `$ is its trace, $`g_{ij}^{(\mathrm{induced})}`$ is the induced metric, and $`_n`$ is the normal derivative at the Planck brane. Equation (67) admits a solution where the induced metric is flat precisely if $`\pm \lambda _{\mathrm{Pl}}(\stackrel{}{\phi })`$ is tangent to the $`W(\stackrel{}{\phi })`$ used to generate the bulk geometry. Essentially this condition was obtained in .
The generalization of is to demand that $`W(\stackrel{}{\phi })=\pm \lambda _{\mathrm{Pl}}(\stackrel{}{\phi })`$ identically solves (65). When $`V(\phi )=0`$ for a single scalar $`\phi `$, this reduces to the special exponential $`\lambda _{\mathrm{Pl}}(\phi )`$ considered explicitly in . In Horava-Witten theory compactified on a Calabi-Yau manifold, with some embedding of the spin connection in one $`E_8`$, the tensions of the two ends of the universe do coincide, as functions of the Kahler moduli, with an appropriately defined bulk superpotential. This seeming coincidence is a result of supersymmetry. The corresponding no-force condition indeed guarantees the existence of solutions with $`3+1`$-dimensional Poincaré invariance. But it seems implausible that the tension of an end-of-the-universe brane after supersymmetry breaking would solve the equation (65) identically.
The generalization of is to demand only that $`W(\stackrel{}{\phi })`$ and $`\pm \lambda _{\mathrm{Pl}}(\stackrel{}{\phi })`$ are tangent at some point. This requires exactly one fine-tuning if $`W(\stackrel{}{\phi })`$ is regarded as a fixed function of $`\stackrel{}{\phi }`$. If there are $`n`$ scalars, then there is an $`n`$-parameter family of solutions to (65). If those $`n`$ parameters are left free, then there is an $`(n1)`$-parameter family of solutions with $`3+1`$-dimensional Poincaré invariance both in the bulk and on the brane. Thus one “postpones” the cosmological constant problem from a single fine-tuning of the parameters of the theory to a single fine-tuning to the actual state of the universe (specified in this case by a choice of integration constants for (67)). It is not clear that there is any dynamical adjustment mechanism to push the universe toward a flat-space solution. (The current mentioned after (64) is only conserved in the case of a constant potential for some scalar).
On one hand, we should note that the conjecture (2), the fluctuation analysis in section 5, and the examples in sections 3.2 and 4 all point toward conditions at curvature singularities which take the form of inequalities rather than equalities. If this is the true state of affairs, and not just an indication of flat directions in the dual field theory or of an insufficiently precise understanding of the microscopic physics at the singularities, then the Stanford construction, as well as variations of it in the spirit of , are probably tenable. The only remaining proviso is that there is not in general a clear mechanism for preferring nearly flat solutions above all others.
On the other hand, if we believe that the bulk geometry should be the limit of a black hole solution (which seems particularly reasonable in a cosmological context, where the current state of the universe was arrived at through a long cooling process from very hot initial conditions), then the natural expectation is that boundary conditions at the naked singularity are inherited from its near-extremal generalizations. For instance, in the case of a free bulk scalar, the boundary condition in a static black hole geometry can be phrased as zero shift current across the horizon. Imposing this condition on naked singularities caused by a divergence in the bulk scalar simply rules out the existence of static singular solutions. At best these singularities could appear as transient states, probably with a lifetime comparable to the five-dimensional Planck time.
The construction of is an excellent illustration of points we have argued. The singularity is of exactly the same type as the ones in . In our language, the singularity is characterized by $`\zeta =\sqrt{8/3}`$. The only difference between and is that in there is a potential for the bulk scalar which goes to $`\mathrm{}`$ at the singularity. The physics is transparent if we are willing to take the view that the bulk represents a cutoff quantum field theory, coupled to gravity by the existence of the Planck brane . The bulk is precisely the five-dimensional representation of a state on the Coulomb branch of $`𝒩=4`$ gauge theory, corresponding to a ten-dimensional geometry where the D3-branes are arranged in an $`S^3`$ shell. This is a state which seems obviously incapable of supporting finite temperature. Because of the bulk potential, one might hope that the arguments following (64) have no force and that there are near-extremal generalizations. This is in fact one of the cases which we studied numerically (see discussion following (62)), with the conclusion that there are near-extremal geometries, but their extremal limit is a singularity with $`\zeta =\sqrt{1/6}`$. The non-generic value of $`\zeta `$ amounts to having precisely one boundary condition at the singularity, which is bad news because the problem of fitting a flat Planck brane to the bulk geometry is once again fine-tuned by one real parameter!
According to what was referred to the “Weak Cosmic Censorship” conjecture in section 6, the situation of the previous paragraph is generic. If weak cosmic censorship is right, then singularities of the form (29) with $`\zeta \sqrt{8/3}`$ cannot be obtained as limits of regular black holes. This seems like only a slightly stronger statement than what we proved in sections 3.1 and 3.3, but in fact $`\zeta =\sqrt{8/3}`$ is both a hard case to settle and a common one in examples. Our evidence for ruling it out is 1) numerics on two examples, 2) the calculations following (64), and 3) the field theory intuition of the previous paragraph. This is not an airtight case. However, the conjecture that singularities with $`\zeta \sqrt{8/3}`$ do not admit near-extremal generalizations is mathematically well-posed and capable of proof or disproof.
From the point of view of a boundary value problem, the boundary conditions at the Planck brane, (67), have essentially the same character as the boundary conditions at the true boundary of $`AdS_5`$ that correspond to fixing the microscopic lagrangian. Replacing the naked singularity either with a negative tension brane or with a black hole horizon imposes enough boundary conditions to completely determine the boundary value problem. The cosmological constant is not free, but fixed, and it still seems accidental that its value in four-dimensional Planck units is as small as $`10^{120}`$. Going on the AdS/CFT intuition that the naked singularity represents interesting infrared dynamics of a quantum field theory, it seems obvious that supersymmetry breaking should induce a positive cosmological constant. It would be fascinating if this intuition were somehow wrong, and if the naked singularity could be somehow associated with a breaking of supersymmetry that does not lead to a four-dimensional cosmological constant. To realize this hope in a convincing and coherent model would require some new ideas.
Constructions such as those in , and near-extremal generalizations of them are nevertheless interesting. If the Planck brane’s properties can somehow be fine-tuned so that a flat solution can be obtained, then black hole generalizations of the Poincaré invariant solution allow one to obtain interesting spatially flat FRW metrics. The mass density above extremality of the black hole will make a contribution to $`(\dot{a}/a)^2`$ in the Friedmann equation. This was demonstrated in detail in for the case of an $`AdS_5`$ bulk, and has been considered in more generality in . The black hole geometry represents hidden sector degrees of freedom at finite temperature.
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# On existence of mini-boson stars
## 1 Introduction
Boson stars are compact gravitationally bound soliton-like equilibrium configurations of bosonic fields. The simplest kind of boson star, which is made up of a self-gravitating free complex massive scalar field, was conceived over thirty years ago by Kaup and Ruffini and Bonazzola who found numerically the ground state solution to the spherically symmetric Einstein-Klein-Gordon (EKG) equations. A decade later the systematic numerical analysis of these equations was performed by Friedberg, Lee, and Pang who rediscovered and extended the results of , in particular they found a countable sequence of excited states.
The aim of this paper is to give a rigorous proof of existence of solutions found in . In the physics literature these solutions are usually referred to as mini-boson stars (’mini’ because they are tiny objects with mass $``$ $`\frac{1}{Gm}`$, where $`m`$ is the boson mass). What are the motivations for studying such objects? Let us mention three possible reasons varying from physical to purely mathematical. First, most theories of elementary particles predict the existence of massive bosons which interact weakly with baryonic matter. To the extent one believes in these models, one should accept their consequences, like boson stars. From this standpoint, the recent surge of interest in boson stars is largely due to the suggestion that the dark matter could be bosonic since then some fraction of the missing mass of the universe would float around in the form of boson stars. Second, even if massive scalar fields do not exist in nature, they provide one of the simplest fundamental matter sources for the Einstein equations and, as such, are ideal theoretical “laboratories” for studying the dynamics of gravitational collapse. Mathematically, these studies amount to the analysis of the Cauchy problem for the EKG equations. Boson stars play an important role in this context as candidates for intermediate or final attractors of dynamical evolution. Finally, and admittedly most interestingly for us, mini-boson stars are non-perturbative solutions of the EKG equations in the sense that they have no regular flat-spacetime limit (one manifestation of this property is the fact mentioned above that the total mass of a mini-boson star is inversely proportional to the gravitational constant $`G`$). In this respect mini-boson stars are similar to the Bartnik-McKinnon solutions of the Einstein-Yang-Mills equations . However, in contrast to the Bartnik-McKinnon solutions, the mini-boson stars are not static: although the metric and the stress-energy tensor of the scalar field are time-independent, the scalar field iself has the form of a standing wave $`\varphi (r,t)=e^{i\omega t}\stackrel{~}{\varphi }(r)`$. This fact has an important consequence at the ode level, namely the lapse function does not decouple from the Klein-Gordon equation and the hamiltonian constraint which means that we have to deal with a 4-dimensional (nonautonomous) dynamical system<sup>1</sup><sup>1</sup>1For comparison, the static spherically symmetric Einstein-Yang-Mills equations reduce (within the purely magnetic ansatz) to a 3-dimensional (nonautonomous) dynamical system.. Below we analyze this system using a shooting method which is similar in spirit (but quite different in implementation) to the proof of existence of the Bartnik-McKinnon solutions .
The paper is organized as follows. In Section 2 we derive the field equations together with the boundary conditions and discuss some basic properties of solutions. We also formulate the main theorem and sketch the heuristic idea of its proof. In Section 3 we prove the local existence of solutions near the origin. In Sections 4 and 5 we discuss the limiting behavior of solutions for small and large values of the shooting parameter, respectively. In Section 6 we derive the asymptotics of globally regular solutions. Section 7 contains some technical results concerning the behavior of singular solutions. Finally, in Section 8, using the results of Sections 4-7, we complete the proof of the main theorem by a shooting argument.
## 2 Preliminaries
The action for the EKG system is given by
$$S=d^4x\sqrt{g}\left(\frac{R}{16\pi G}\frac{1}{2}_a\varphi ^{}^a\varphi \frac{1}{2}m^2\varphi ^{}\varphi \right),$$
(2.1)
where $`R`$ is the scalar curvature of the spacetime metric $`g_{ab}`$, $`\varphi `$ is the complex scalar field, and $`m`$ is a real constant called the boson mass. The associated field equations are the Einstein equations
$$R_{ab}\frac{1}{2}g_{ab}R=8\pi GT_{ab}$$
(2.2)
with the stress-energy tensor of the scalar field
$$T_{ab}=\frac{1}{2}(_a\varphi ^{}_b\varphi +_a\varphi _b\varphi ^{})\frac{1}{2}g_{ab}(g^{cd}_c\varphi ^{}_d\varphi +m^2\varphi ^{}\varphi ),$$
(2.3)
and the Klein-Gordon equation
$$(\mathrm{}m^2)\varphi =0,$$
(2.4)
where $`\mathrm{}`$ is the d’Alembertian operator associated with the metric $`g_{ab}`$. Now, we assume that the fields are spherically symmetric. We write the metric using areal radial coordinate and polar slicing
$$ds^2=e^{2\delta }Adt^2+A^1dr^2+r^2d\mathrm{\Omega }^2,$$
(2.5)
where $`d\mathrm{\Omega }^2`$ is the standard metric on the unit two-sphere, and $`A`$ and $`\delta `$ are functions of $`(t,r)`$. In this parametrization the (relevent components of) Einstein equations have the particularly simple form
$`_rA`$ $`=`$ $`{\displaystyle \frac{1A}{r}}8\pi GrT_{00},`$ (2.6)
$`_r\delta `$ $`=`$ $`4\pi GrA^1(T_{00}+T_{11}),`$ (2.7)
$`_tA`$ $`=`$ $`8\pi Gre^\delta AT_{01},`$ (2.8)
where the components of stress-energy tensor $`T_{ab}`$ are expressed in the orthonormal frame determined by the metric (2.5) ($`e_0=e^\delta A^{1/2}_t`$, $`e_1=A^{1/2}_r`$). From (2.3) we obtain
$`T_{00}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(A|_r\varphi |^2+A^1e^{2\delta }|_t\varphi |^2+m^2|\varphi |^2),`$ (2.9)
$`T_{11}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(A|_r\varphi |^2+A^1e^{2\delta }|_t\varphi |^2m^2|\varphi |^2),`$ (2.10)
$`T_{01}`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^\delta (_r\varphi ^{}_t\varphi +_r\varphi _t\varphi ^{}).`$ (2.11)
The remaining components of Einstein’s equations are equivalent to the Klein-Gordon equation.
For the scalar field $`\varphi `$ we assume the standing wave ansatz $`\varphi (r,t)=\mathrm{exp}(i\omega t)\stackrel{~}{\varphi }(r)`$, where $`\omega `$ is a real constant. Then, due to the $`U(1)`$ symmetry of the action, the stress-energy tensor and the metric are time-independent. Morever, $`T_{01}=0`$ so Eq. (2.8) is trivially satisfied. In terms of the dimensionless variables
$$x=mr,f(x)=\sqrt{4\pi G}\stackrel{~}{\varphi }(r)$$
(2.12)
and the auxiliary variable
$$C(x)=\frac{\omega }{m}A^1e^\delta ,$$
(2.13)
Eqs. (2.4),(2.6) and (2.7) reduce to the following system of ordinary differential equations (hereafter prime denotes $`\frac{d}{dx}`$)
$`\left({\displaystyle \frac{x^2f^{}}{C}}\right)^{}`$ $`=`$ $`{\displaystyle \frac{x^2}{AC}}(1AC^2)f,`$ (2.14a)
$`A^{}`$ $`=`$ $`{\displaystyle \frac{1A}{x}}x(Af_{}^{}{}_{}{}^{2}+AC^2f^2+f^2),`$ (2.14b)
$`C^{}`$ $`=`$ $`{\displaystyle \frac{C}{xA}}(A1+x^2f^2).`$ (2.14c)
Instead of $`A`$, it is sometimes convenient to use the “mass” function $`M(x)`$ defined by $`A(x)=12M(x)/x`$. From (2.14b) we have
$$M(x)=\frac{1}{2}_0^xs^2(Af_{}^{}{}_{}{}^{2}+AC^2f^2+f^2)𝑑s.$$
(2.15)
A spacetime is said to be asymptotically flat if $`\delta (\mathrm{})`$ is finite and
$$\underset{x\mathrm{}}{lim}M(x)=M_{\mathrm{}}<\mathrm{}.$$
(2.16)
The limiting value $`M_{\mathrm{}}`$ is interpreted as the total mass of a solution (in our case it is measured in units $`\frac{1}{Gm}`$). In Section 6 we will show that the finiteness of mass implies that $`C`$ has a finite limit and $`f`$ decays exponentially as $`x\mathrm{}`$.
Besides the singularity at infinity, the field equations (2.14) have the fixed singular point at $`x=0`$ and a moving singularity at $`\overline{x}`$, where $`A(\overline{x})=0`$. Regularity of solutions at $`x=0`$ requires the following behavior
$$f(x)=a+O(x^2),A(x)=1+O(x^2),C(x)=\alpha +O(x^2),$$
(2.17)
where $`a=f(0)`$ and $`\alpha =C(0)`$ are arbitrary parameters (assumed positive without loss of generality). In Sect. 3 we will show that these parameters determine uniquely a smooth local solution to Eqs.(2.14).
###### Definition 2.1.
The solution of Eqs.(2.14) starting at $`x=0`$ with the behavior (2.17) is called the $`𝐚\mathrm{𝐨𝐫𝐛𝐢𝐭}`$.
In the following whenever we write ’a solution’ we always mean the $`a`$-orbit. Also when we write that some property holds for all $`x`$ we always mean for all $`x0`$. We will frequently refer to the behavior of $`a`$-orbits in the $`(f,f^{})`$-plane; when we write, say, that the $`a`$-orbit enters the first quadrant (Q1 for brevity), we mean that the projection of the $`a`$-orbit in the $`(f,f^{})`$-plane does so.
###### Definition 2.2.
The $`a`$-orbits which exist for all $`x`$ and are asymptotically flat are called globally regular.
Now, we are ready to formulate our main result:
###### Theorem 2.1.
For each $`\alpha >1`$, there is a decreasing sequence of parameters $`a_n`$ ($`n=0,1,2,\mathrm{}`$) such that the corresponding $`a_n`$-orbits are globally regular. The index $`n`$ labels the number of nodes of the function $`f(x)`$.
This theorem makes rigorous the numerical results obtained in . Notice that although the $`a`$-orbits are determined by two parameters, only the parameter $`a`$ has to be fine-tuned so the shooting is essentially one-dimensional.
In order to prepare the ground for the proof of Theorem 2.1 we discuss now some elementary global properties of $`a`$-orbits.
###### Lemma 2.2.
$`A(x)<1`$ for all $`x>0`$ unless $`f(x)0`$.
###### Proof.
From (2.14b), $`A^{}(x_0)<0`$ if $`A(x_0)=1`$ so $`A`$ cannot cross $`1`$ from below. Since $`A(x)<1`$ for small $`x`$, the lemma follows. ∎
###### Lemma 2.3.
An $`a`$-orbit exists as long as $`A(x)>0`$.
###### Proof.
If $`A(x)>0`$ for $`x<\overline{x}<\mathrm{}`$, then $`lim_{x\overline{x}}M(x)`$ exists (because $`M^{}>0`$ and $`M(x)<x<\overline{x}`$) so $`lim_{x\overline{x}}A`$ exists as well. We will show that the orbit can be continued beyond $`\overline{x}`$ provided that $`lim_{x\overline{x}}A(x)>0`$. Since $`0<A<1`$, the only obstruction to extending the solution is the possibility that $`C`$,$`f`$, or $`f^{}`$ might be unbounded. To see that $`f`$ is bounded we note that $`(xA)^{}<1x^2Af_{}^{}{}_{}{}^{2}`$. Choose $`ϵ>0`$ such that $`A(x)>A(\overline{x})/2`$ for $`\overline{x}ϵ<x<\overline{x}`$; then $`(xA)^{}<1(\overline{x}ϵ)^2A(\overline{x})f_{}^{}{}_{}{}^{2}/2`$ and integrating from $`\overline{x}ϵ<x`$ gives that $`_{\overline{x}ϵ}^xf^{}(x)^2<\mathrm{}`$ and hence, by the Cauchy-Schwarz inequality, $`_{\overline{x}ϵ}^x|f^{}(x)|<\mathrm{}`$. Thus, $`f`$ is bounded. This implies by Eq.(2.14c) that $`(\mathrm{ln}C)^{}`$ is also bounded so both $`C`$ and $`1/C`$ are bounded. Now, (2.14a) says that $`x^2f^{}/C`$ is bounded so $`f^{}`$ is bounded. ∎
###### Remark.
It follows from Lemma 2.3 that the only possible obstruction to extendability of $`a`$-orbits to arbitrarily large $`x`$ is $`lim_{x\overline{x}}A=0`$ for some $`\overline{x}`$. If that happens we will say that the solution crashes at $`\overline{x}`$.
Let us define the function $`g=1AC^2`$. The following two properties of this function will play an important role in our discussion.
###### Lemma 2.4.
We have
If $`g(x)0`$, then $`g^{}(x)>0`$;
If $`g(x_0)0`$, then $`g(x)>0`$ for all $`x>x_0`$.
###### Proof.
A simple calculation yields
$$g^{}=C^2\left(\frac{1A}{x}+xAf_{}^{}{}_{}{}^{2}xgf^2\right).$$
(2.18)
The part (a) follows immediately from (2.18). To prove the part (b) note that $`g^{}(x_1)>0`$ if $`g(x_1)=0`$, so $`g`$ cannot cross zero from above. ∎
The restriction $`\alpha >1`$ in Theorem 2.1 can be easily seen as follows. Suppose that there is a globally regular solution with $`\alpha 1`$. Since $`g(0)=1\alpha ^2`$, it follows from Lemma 2.4 that $`g(x)`$ is positive for all $`x`$. Multiplying Eq.(2.14a) by $`f`$ and integrating by parts we get that $`ff^{}>0`$ for all $`x`$, hence $`f^2`$ is monotone increasing which is obviously impossible for globally regular solutions (in fact such solutions crash at finite $`x`$ as follows easily from Eq.(2.14b)). Thus we have
###### Lemma 2.5.
There are no (nontrivial) globally regular solutions for $`\alpha 1`$.
Note that Lemma 2.5 implies in particular that there are no static ($`\alpha =0`$) globally regular solutions. In view of Lemma 2.5 from now on we always assume that $`\alpha >1`$.
###### Definition 2.3.
The rotation function $`\theta (x,a)`$ of an $`a`$-orbit is defined by $`\theta (0,a)=0`$, $`\mathrm{tan}\theta (x,a)=f^{}(x)/f(x)`$ and $`\theta (x,a)`$ is continuous in $`x`$. We will drop the second argument of $`\theta `$ if there is no danger of confusion.
Now we list the basic properties of the rotation function of $`a`$-orbits which we will need below.
###### Lemma 2.6.
For any nonnegative integer $`n`$ we have:
If $`\theta (x_1)>(n+1/2)\pi `$ for some $`x_1`$, then $`\theta (x)>(n+1/2)\pi `$ for all $`x>x_1`$.
If $`\theta (x_1)<n\pi `$ for some $`x_1`$ and $`g(x_1)0`$, then $`\theta (x)<n\pi `$ for all $`x>x_1`$.
There are at most two values of $`x`$ with $`\theta (x)=n\pi `$.
###### Proof.
(a) We note that $`\theta ^{}(x)=(f_{}^{}{}_{}{}^{2}ff^{\prime \prime })/(f^2+f_{}^{}{}_{}{}^{2})`$, so $`\theta ^{}(x)=1`$ if $`\theta (x)=(n+1/2)\pi `$.
(b) If $`x>x_1`$ and $`g(x_1)0`$, then $`g(x)>0`$ by Lemma 2.4. Next, we note that $`\theta ^{}(x)=g(x)/A(x)<0`$ if $`\theta (x)=n\pi `$ and $`g(x)>0`$. If $`g(x)=0`$ then $`\theta ^{}(x)=0`$ but $`\theta ^{\prime \prime }(x)=g^{}(x)/A(x)<0`$ since $`g^{}(x)>0`$ when $`g(x)=0`$ by Eq.(2.18).
(c) The function $`\theta (x)n\pi `$ changes sign at each zero for which $`g(x)0`$. From Lemma 2.4, $`g`$ changes sign at most once. Thus, for $`n>0`$, $`\theta (0)n\pi <0`$ and at $`x_1`$, the first zero of $`\theta (x)n\pi `$, if $`g(x_1)0`$ then by part b) $`\theta (x)n\pi <0`$ for all $`x>x_1`$. If $`g(x_1)<0`$ then $`\theta (x)n\pi `$ changes sign at $`x_1`$, and hence, at $`x_2`$, the next zero of $`\theta (x)n\pi `$, $`g(x_2)0`$ and hence $`\theta (x)n\pi <0`$ for all $`x>x_2`$. For $`n=0`$, $`\theta (0)n\pi =0`$, $`\theta (x)>0`$ near $`x=0`$ and if $`\theta (x_1)=0`$ then $`g(x_1)0`$, hence, $`\theta (x)<0`$ for all $`x>x_1`$. ∎
Before going into details of Sections 3–8, let us outline the main idea of the proof of Theorem 2.1. According to this theorem there exists a countable family of globally regular solutions distinguished by nodal class. We first show (section 3) that there is a continuous one-parameter family of local solutions depending on $`a=f(0)`$; we all these solutions $`a`$-orbits. In Section 6 we show that an $`a`$-orbit that has bounded rotation and that is defined for all $`x`$ is a globally regular solution, that is, it has the correct asymptotic behavior as $`x\mathrm{}`$. The existence of $`a`$-orbits with bounded rotation that are defined for all $`x`$ is proven in each nodal class by an inductive application of a shooting argument. The zeroeth solution we construct has $`\theta (x,a_0)<\pi /2`$ for all $`x`$; the first solution has $`\theta (x,a_1)<3\pi /2`$ (and greater than $`\pi /2`$ for large $`x`$), etc. This is shown in Fig. 1. The crucial step of our argument is the control of behavior of $`a`$-orbits for large and small values of the parameter $`a`$. In Section 4 we show that for sufficiently small $`a`$ the $`a`$-orbit has arbitrarily large rotation; more precisely, there is a number $`b_n`$ such that $`\theta (x,a)>n\pi `$ for some $`x`$ if $`a<b_n`$. In contrast, we show in Section 5 that for $`a>>1`$ the $`a`$-orbit exits Q4 directly to Q1 (see Fig. 1).
Now, to prove the existence of a globally regular solution in the zeroeth nodal class we let $`a_0=inf\{a|\theta (x,a)<\pi /2`$ for all $`x`$ for which the $`a`$-orbit is defined}. Note that $`a_0b_1>0`$. We then prove that the $`a_0`$-orbit is the globally regular solutions in the zeroeth nodal class. It is clear that the $`a_0`$-orbit has rotation $`\theta (x,a_0)\pi /2`$ for otherwise all nearby orbits would have rotation $`>\pi /2`$ which contradicts the definition of $`a_0`$. It is also easy to see that the $`a_0`$-orbit cannot exit Q4 to Q1 because again, nearby orbits would also do so which contradicts the definition of $`a_0`$. Hence, the $`a_0`$-orbit must stay in Q4; it either crashes or is defined for all $`x`$ and is a globally regular solutions in the zeroeth nodal class. Thus, it remains to show that the $`a_0`$-orbit does not crash. The (technical) crash lemma of Section 7 shows that if an orbit crashes in Q4 then nearby orbits either crash in Q4 or exit Q4 to Q1. Thus the $`a_0`$-orbit cannot crash because nearby orbits would all be in $`\{a|\theta (x,a)<\pi /2`$ for all $`x`$ for which the $`a`$-orbit is defined} and $`a_0`$ would not be the infimum of that set.
To show the existence of globally regular solutions in higher nodal classes we proceed as above. We let $`a_n=inf\{a|\theta (x,a)<(n+1/2)\pi `$ for all $`x`$ for which the $`a`$-orbit is defined}. We then show that $`\theta (x,a_n)<(n+1/2)\pi `$. We again use the crash lemma as we did in the $`n=0`$ case to show the $`a_n`$-orbit is defined for all $`x`$. The only difference is that we must show that $`\theta (x,a_n)>n\pi `$. That fact follows easily from lemmas 2.6b and 6.3.
## 3 Local existence
###### Proposition 3.1.
There exists a two-parameter family of local solutions of Eqs.(2.14) near $`x=0`$ satisfying the initial conditions (2.17).
###### Proof.
The proof is standard so we just outline it. We introduce new variables $`w=f^{}`$, $`z=\mathrm{ln}(C)`$, $`B=(1A)/x`$, and rewrite Eqs.(2.14) as the first order system
$`f^{}`$ $`=`$ $`w,`$ (3.1a)
$`(x^2w)^{}`$ $`=`$ $`{\displaystyle \frac{x^2}{A}}\left((xf^2B)w+f(1AC^2)\right),`$ (3.1b)
$`(x^2B)^{}`$ $`=`$ $`x^2(Aw^2+f^2AC^2+f^2),`$ (3.1c)
$`z^{}`$ $`=`$ $`{\displaystyle \frac{xf^2B}{A}}.`$ (3.1d)
We will use the sup norm throughout this discussion: $`h`$ means the $`sup\{|h(x)|:0xr\}`$.
Consider the space $`X`$ of quadruples of functions $`(f,y,B,z)`$ where $`fa1,w1,BM`$, and $`z\mathrm{ln}(\alpha )1`$ and each of the four functions is in $`C^0([0,r])`$, the space of continuous functions defined on the interval $`0xr`$ with the sup norm. $`X`$ is a complete metric space if we take as metric the maximum of the four components. We define a map $`T:XX`$ by $`T(f,w,B,z)=(T_1,T_2,T_3,T_4)`$ where
$`T_1`$ $`=`$ $`a+{\displaystyle _0^x}w𝑑s,`$ (3.2a)
$`T_2`$ $`=`$ $`{\displaystyle \frac{1}{x^2}}{\displaystyle _0^x}{\displaystyle \frac{s^2}{A}}\left((sf^2B)w+f(1AC^2)\right)𝑑s,`$ (3.2b)
$`T_3`$ $`=`$ $`{\displaystyle \frac{1}{x^2}}{\displaystyle _0^x}s^2(Aw^2+f^2AC^2+f^2))ds,`$ (3.2c)
$`T_4`$ $`=`$ $`\mathrm{ln}\alpha +{\displaystyle _0^x}{\displaystyle \frac{1}{A}}(sf^2B)𝑑s.`$ (3.2d)
One verifies easily that $`T`$ does in fact take $`X`$ to $`X`$ and that $`T`$ is a contracting map if $`r`$ is sufficiently small, and that a fixed point of $`T`$ is a solution to our equations. The proof that the solution depends continuously on $`a`$ is also routine. ∎
## 4 Behavior of solutions for small $`a`$
In this section we show that the rotation $`\theta (x,a)`$ of the $`a`$-orbit is arbitrarily large if $`a`$ is sufficiently small and $`x`$ is sufficiently large.
###### Proposition 4.1.
For any $`n>0`$, there exists a $`b_n`$ such that for $`a<b_n`$ there is an $`x`$ with $`\theta (x,a)>n\pi `$.
###### Proof.
Let $`\stackrel{~}{f}=f/a`$. Then, Eqs.(2.14) become
$`\left({\displaystyle \frac{x^2\stackrel{~}{f}^{}}{C}}\right)^{}`$ $`=`$ $`{\displaystyle \frac{x^2}{AC}}(1AC^2)\stackrel{~}{f},`$ (4.1a)
$`A^{}`$ $`=`$ $`{\displaystyle \frac{1A}{x}}a^2x(A\stackrel{~}{f^{}}^2+AC^2\stackrel{~}{f}^2+\stackrel{~}{f}^2),`$ (4.1b)
$`C^{}`$ $`=`$ $`{\displaystyle \frac{C}{xA}}(A1+a^2x^2\stackrel{~}{f}^2)`$ (4.1c)
with the behavior at the origin
$$\stackrel{~}{f}(0)=1,A(0)=1,C(0)=\alpha .$$
(4.2)
For $`a=0`$ (decoupling of gravity) Eqs.(4.1bc) with conditions (4.2) have constant flat-spactime solutions $`A1`$, $`C\alpha `$. Inserting these solutions into Eq.(4.1a) gives the Bessel equation
$$(x^2\stackrel{~}{f}^{})^{}+x^2(\alpha ^21)\stackrel{~}{f}=0,$$
(4.3)
whose unique solution satisfying (4.2) is
$$\stackrel{~}{f}(x)=\frac{\mathrm{sin}\sqrt{\alpha ^21}x}{\sqrt{\alpha ^21}x}.$$
(4.4)
This solution has infinite rotation as $`x\mathrm{}`$. If $`x>\frac{n\pi }{\sqrt{\alpha ^21}}`$ then $`\theta (x,0)>n\pi `$ so for $`a`$ close to $`0`$, say $`a<b_n`$, we have $`\theta (x,a)>n\pi `$ because solutions of Eqs.(4.1) are continuous in $`a`$ and $`x`$. This concludes the proof of Proposition 4.1. ∎
## 5 Behavior of solutions for large $`a`$
###### Proposition 5.1.
The $`a`$-orbits with sufficiently large $`a`$ exit Q4 directly to Q1.
We define new variables
$$y=ax,\stackrel{~}{v}(y)=a(f(x)a),\stackrel{~}{A}(y)=A(x),\stackrel{~}{C}(y)=C(x).$$
(5.1)
Then, Eqs.(2.14) become (where now the prime denotes the derivative with respect to $`y`$)
$`\left({\displaystyle \frac{y^2v^{}}{\stackrel{~}{C}}}\right)^{}`$ $`=`$ $`{\displaystyle \frac{y^2}{\stackrel{~}{A}\stackrel{~}{C}}}(1\stackrel{~}{A}\stackrel{~}{C}^2)(1+{\displaystyle \frac{\stackrel{~}{v}}{a^2}}),`$ (5.2a)
$`\stackrel{~}{A}^{}`$ $`=`$ $`{\displaystyle \frac{1\stackrel{~}{A}}{y}}y\left({\displaystyle \frac{1}{a^2}}A\stackrel{~}{v^{}}^2+(1+\stackrel{~}{A}\stackrel{~}{C}^2)\left(1+{\displaystyle \frac{\stackrel{~}{v}}{a^2}}\right)^2\right),`$ (5.2b)
$`\stackrel{~}{C}^{}`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{C}}{y\stackrel{~}{A}}}\left(\stackrel{~}{A}1+y^2(1+{\displaystyle \frac{\stackrel{~}{v}}{a^2}})^2\right).`$ (5.2c)
The initial conditions at $`y=0`$ are
$$\stackrel{~}{A}(0)=1,\stackrel{~}{C}(0)=\alpha >1,\stackrel{~}{v}(0)=0,\stackrel{~}{v}^{}(0)=0.$$
(5.3)
As $`a\mathrm{}`$, the solutions of Eqs.(5.2) tend uniformly on compact intervals to the solutions of the following limiting system
$`\left({\displaystyle \frac{y^2v^{}}{C}}\right)^{}`$ $`=`$ $`{\displaystyle \frac{y^2}{AC}}(1AC^2),`$ (5.4a)
$`A^{}`$ $`=`$ $`{\displaystyle \frac{1A}{y}}y(1+AC^2),`$ (5.4b)
$`C^{}`$ $`=`$ $`{\displaystyle \frac{C}{yA}}(A1+y^2).`$ (5.4c)
satisfying the same initial conditions as in (5.3). The rest of this section is devoted to the analysis of Eqs.(5.4). Our goal is to show that $`v^{}(y)`$ becomes positive at a point $`y_1<\overline{y}`$. This would imply that $`v(y)`$ is bounded below, i.e., there is $`d>0`$ such that $`v(y)>d`$ for $`y<y_1`$, and therefore $`\stackrel{~}{v}(y)>d1`$ if $`a`$ is sufficiently large. Then $`f(x)>a(d+1)/a`$ and $`f^{}(x)>0`$ for $`x=y_1/a`$, hence, if $`a>\sqrt{d+1}`$, the $`a`$-orbit exits Q4 to Q1 directly without entering Q3.
Note that the function $`v`$ decouples from Eqs.(5.4bc) for the metric coefficients – this fact considerably simplifies the analysis.
###### Lemma 5.2.
The solution of Eqs.(5.4) crashes at some $`\overline{y}`$, that is, $`A(\overline{y})=lim_y\mathrm{}A(y)=0`$. Moreover, $`1<\overline{y}<\sqrt{3}`$.
###### Proof.
Note that $`(yA)^{}<1y^2`$ so integrating gives $`A<1y^2/3`$. Therefore, $`A(\overline{y})=0`$ for $`\overline{y}<\sqrt{3}`$. To show that $`\overline{y}>1`$ assume $`\overline{y}1`$. Then $`(\frac{y}{C})^{}=\frac{1y^2}{AC}0`$ so if $`0<\tau <y<\overline{y}`$ we have $`y/C(y)>\tau /C(\tau )`$ or $`C(y)<C(\tau )/\tau `$ so $`C`$ is bounded. Since $`(AC)^{}=yAC^3`$, $`(\mathrm{ln}(AC))^{}=yC^2`$ is bounded below. Thus, by integrating one concludes that $`lim_{y\overline{y}}AC(y)>0`$. But $`(AC)^2=A(AC^2);AC^2\alpha ^2`$ and $`lim_{y\overline{y}}A(y)=0`$ so $`lim_{y\overline{y}}(AC(y))^2=0`$. This contradicts $`lim_{y\overline{y}}AC(y)>0`$ so we must have $`\overline{y}>1`$. ∎
###### Proof of Proposition 5.1.
In order to prove that $`v^{}(y)`$ becomes positive at some point $`y_1<\overline{y}`$, we will show that $`v^{}(\overline{y})>0`$. By Eq.(5.4a) we have $`v^{}(y)=\frac{C}{y^2}_0^y\frac{z^2(1AC^2)}{AC}𝑑z`$, so we must show that $`_0^{\overline{y}}\frac{y^2(1AC^2)}{AC}𝑑y>0`$.
The proof of this fact is divided into two cases: (i) $`\overline{y}^23/2`$, and (ii) $`\overline{y}^2<3/2`$. Before considering these cases we list some useful properties of the function $`g=1AC^2`$.
###### Lemma 5.3.
We have:
$`g^{}=(1Ay^2g)C^2/y`$;
if $`g(y_0)0`$, then $`g(y)>0`$ for all $`y>y_0`$;
$`g^{}>0`$ if $`g1/3`$;
###### Proof.
Part (a) is a calculation. For (b) note that $`g^{}(g=0)>0`$ so $`g`$ cannot cross zero from above. For (c) we have $`(yA)^{}<1y^2`$ so integrating gives $`1A>y^2/3`$ and hence, $`g^{}>y(1/3g)C^2`$. ∎
We return now to the proof that $`_0^{\overline{y}}\frac{y^2g}{AC}𝑑y>0`$. We first consider the case (i) $`\overline{y}^2>3/2`$.
A calculation shows that $`y^2g=(2y^3/3+yAy)^{}`$, hence $`_0^{\overline{y}}y^2g𝑑y=2\overline{y}^3/3\overline{y}>0`$ if $`\overline{y}^2>3/2`$. Since $`g(0)=1\alpha ^2<0`$, this implies that $`g(\sigma )=0`$ for some $`\sigma <\overline{y}`$ and therefore $`g(y)>0`$ for $`y>\sigma `$. Note that $`AC`$ is monotone decreasing because $`(AC)^{}=yAC^3<0`$. Thus
$$\frac{y^2g(y)}{A(y)C(y)}\frac{y^2g(y)}{A(\sigma )C(\sigma )}\text{for}0y\overline{y}$$
(5.5)
and therefore
$$_0^{\overline{y}}\frac{y^2g}{AC}𝑑y\frac{1}{A(\sigma )C(\sigma )}_0^{\overline{y}}y^2g𝑑y>0.$$
(5.6)
Now we consider the case (ii) $`\overline{y}^23/2`$.
###### Lemma 5.4.
Define the function $`p=1+y^2gy^2`$. If $`y^23/2`$, then $`p(y)>0`$.
###### Proof.
Note that $`p(0)=1`$. Let $`y_1`$ be the first zero of $`p`$, that is, $`p(y_1)=0`$ and $`p^{}(y_1)0`$. If $`g(y_1)>1/3`$ then $`p=y^2g+1y^2>y^2/3+1y^2=12y^2/3`$. Thus $`p`$ can have a zero for $`y_1^23/2`$ only if $`g(y_1)1/3`$. Then, from Lemma (5.3), $`g^{}(y)>0`$ for all $`yy_1`$. Define a function $`k(y)=23Ay^2`$. A calculation gives $`y^3g^{}=(y(k+p))^{}`$ so by integrating we get $`k(y_1)>0`$. On the other hand we have $`p^{}=yC^2(kp)`$, so $`k(y_1)0`$; contradiction. ∎
To show that $`_0^{\overline{y}}\frac{y^2g}{AC}𝑑y>0`$, we rewrite it as
$$_0^{\overline{y}}\frac{y^2g}{AC}𝑑y=_0^{\overline{y}}\frac{p1+y^2}{AC}𝑑y=_0^{\overline{y}}\frac{p}{AC}𝑑y_0^{\overline{y}}\frac{1y^2}{AC}𝑑y.$$
(5.7)
The first term on the right hand side of (5.7) is positive because $`p`$ is positive. To compute the second term, note that
$$(\frac{y}{C})^{}=\frac{1y^2}{AC},$$
(5.8)
hence $`L=lim_{y\overline{y}}(y/C)`$ exists and is finite since $`\overline{y}>1`$ by Lemma 5.2. If $`L>0`$ then $`lim_{y\overline{y}}C(y)=\overline{y}/L<\mathrm{}`$ so $`C`$ is bounded. Since $`lim_{y\overline{y}}A(y)=0`$ we conclude that $`lim_{y\overline{y}}AC(y)=0`$. But $`(\mathrm{ln}AC)^{}=yC^2`$ is bounded so $`\mathrm{ln}AC`$ is bounded below and hence $`limAC0`$. This contradiction shows that $`L=0`$. Thus, the second term on the right hand side of (5.7) is zero. This concludes the proof of Proposition 5.1.
## 6 Asymptotics of globally regular solutions
In this section we derive the leading asymptotic behavior of globally regular solutions. We use $`lim`$ to denote $`lim_x\mathrm{}`$.
###### Proposition 6.1.
An $`a`$-orbit which exists for all $`x`$ and has bounded rotation is asymptotically flat. The leading asymptotic behavior for $`x\mathrm{}`$ is
$$A(x)1\frac{2M_{\mathrm{}}}{x},C(x)C_{\mathrm{}}e^{\frac{2M_{\mathrm{}}}{x}},f(x)f_{\mathrm{}}e^{bx},$$
(6.1)
where $`0<M_{\mathrm{}}<\mathrm{}`$, $`0<C_{\mathrm{}}<1`$, and $`b=\sqrt{1C_{\mathrm{}}^2}`$.
To prove this proposition we need several partial results.
###### Lemma 6.2.
An $`a`$-orbit which exists for all $`x`$ and has bounded rotation is ultimately in the second (Q2) or fourth (Q4) quadrant.
###### Proof.
If $`\theta (x)`$ is bounded above then there is an integer $`n0`$ such that $`\theta (x)<(n+1/2)\pi `$ for all $`x`$ but $`\theta (x_1)>(n1/2)\pi `$ for some $`x_1`$ and hence, by Lemma 2.6a for all $`x>x_1`$, $`(n1/2)\pi <\theta (x)<(n+1/2)\pi `$. We next show that there is an $`x_2`$ such that for all $`x>x_2`$, $`n\pi <\theta (x)<(n+1/2)\pi `$ (that is, the orbit is ultimately in Q2 or Q4). Note that, by Lemma 2.6c the orbit must satisfy either $`n\pi <\theta (x)<(n+1/2)\pi `$ or $`(n1/2)\pi <\theta (x)<n\pi `$, that is the orbit must lie in Q3 or Q2 if $`n`$ is odd and in Q1 or Q4 if $`n`$ is even. We must rule out the possibility that the orbit is in Q1 or Q3. Assume that the orbit lies in Q1 or Q3 for all $`x>x_1`$. Then $`f(x)f^{}(x)>0`$ for all $`x>x_1`$, so $`f^2(x)f^2(x_1)`$ for all $`x>x_1`$. From Eq.(2.15b) we have $`(xA)^{}=1x^2Af_{}^{}{}_{}{}^{2}x^2f^2AC^2x^2f^2`$ so $`(xA)^{}<1x^2f^2<1x^2f^2(x_1)`$ and hence $`A`$ goes to zero in finite $`x`$. This contradiction concludes the proof. ∎
###### Lemma 6.3.
Under the assumptions of Proposition 6.1 the function $`g=1AC^2`$ is eventually positive.
###### Proof.
Suppose that $`g(x)0`$ for all $`x`$. We claim that this implies $`limA=1`$. To see this, suppose that $`lim\; infA=14ϵ`$ for some $`ϵ>0`$. Let $`\beta =limg0`$ which exists because $`g^{}>0`$. Note that $`g(x)<\beta `$ for all $`x`$. Choose an $`x_1`$ such that $`g(x_1)>\beta ϵ`$. If $`A(x_2)<13ϵ`$ for some $`x_2>x_1`$, then by (2.18) $`g^{}(x)>C^2(1A)/x>(1A)/x=x(1A)/x^2>x_2(1A(x_2))/x^2>3ϵx_2/x^2`$ for $`x>x_2`$, where the last but one inequality follows from the fact that $`x(1A(x))`$ is monotone increasing. Integrating this inequality from $`x_2`$ to $`2x_2`$ say, we get $`g(2x_2)>g(x_2)+3ϵ/2>\beta ϵ+3ϵ/2>\beta `$; contradiction. Thus, $`lim\; infA=1`$ and hence $`limA=1`$. Since $`limg=lim(1AC^2)`$ exists, $`limC`$ also exists and is finite. Next, from Lemma 6.2 we know that the $`a`$-orbit is ultimately in Q2 or in Q4. For concreteness we consider the case of Q4 (the proof of the Q2-case is identical), that is $`f(x)>0`$ and $`f^{}(x)<0`$ for sufficiently large $`x`$. Then, from (2.14a), $`lim(x^2f^{}/C)`$ exists so $`lim(x^2f^{})=\tau <0`$ exists as well (where $`\tau `$ might be infinite; the point is that $`\tau 0`$). Now, by L’Hôpital’s rule, $`limxf=lim(x^2f^{})=\tau `$. But (2.14c) says $`(\mathrm{ln}C)^{}>\tau ^2/4x`$ which implies $`limC=\mathrm{}`$, a contradiction. ∎
###### Proof of Proposition 6.2.
From the previous lemma we know that there exists an $`x_1`$ such that $`g(x)>0`$ for $`x>x_1`$. Let $`u=ACf/g`$ for $`x>x_1`$. A calculation shows that $`u^{}=AC(fC^2(1A)/xf^{}g+xff_{}^{}{}_{}{}^{2})/g^2`$ so $`u^{}<0`$ if $`g>0`$. Multiplying Eq.(2.14a) by $`u`$ we obtain
$$(x^2Aff^{}/g)^{}=x^2f^2+x^2f^{}u^{}/C.$$
(6.2)
The right hand side is positive for $`x>x_1`$ so $`x^2Aff^{}/g`$ is negative and increasing, hence it has a finite non-positive limit. This implies that $`x^2f^2`$ is integrable. Similarly, multiplying Eq.(2.14a) by $`f`$ we obtain
$$x^2ff^{}/C=(x^2f^2g+Ax^2f_{}^{}{}_{}{}^{2})/(AC).$$
(6.3)
The right hand side is positive for $`x>x_1`$ so $`x^2ff^{}/C`$ is negative and increasing, hence it has a finite non-positive limit. This implies that $`Ax^2f_{}^{}{}_{}{}^{2}`$ is integrable (recall that $`AC`$ is monotone decreasing). The integrability of $`x^2f^2`$ and $`Ax^2f_{}^{}{}_{}{}^{2}`$ implies via Eq.(2.15) that $`limM=M_{\mathrm{}}<\mathrm{}`$ exists. This concludes the proof that $`A(x)12M_{\mathrm{}}/x`$.
Having $`limA=1`$ we can strengthen Lemma 6.3 by showing that $`limg=g_{\mathrm{}}>0`$ exists. To see this choose an $`x_1`$ such that $`g(x_1)>0`$. Then $`AC^2(x_1)<1`$, hence $`AC(x_1)<1`$. Since $`AC`$ is monotone decreasing, we have $`AC(x)<AC(x_1)`$ for $`x>x_1`$ and thus $`limAC<1`$. Hence, $`limAC^2=(limAC)^2/limA<1`$. Since $`g=1AC^2`$, $`limg`$ exists and $`limg>0`$.
Now we have all we need to derive the asymptotics of $`f`$. Let $`r=f^{}/f`$. Then $`r^{}=f^{\prime \prime }/fr^2=r(1+Ax^2f^2)/(xA)+g/A=g_{\mathrm{}}r^2+ϵ(x)`$, where $`limϵ=0`$. Let $`\sigma (x_2)=\mathrm{max}(|ϵ(x)|)`$ for $`x>x_2`$ and assume that $`x_2`$ is sufficiently large so that $`g_{\mathrm{}}>\sigma (x_2)`$. If $`r(x_2)>\sqrt{g_{\mathrm{}}\sigma (x_2)}`$, then clearly $`r`$ becomes eventually positive which contradicts that the orbit is eventually in Q2 or Q4. If $`r(x_2)<\sqrt{g_{\mathrm{}}+\sigma (x_2)}`$, then $`limr=\mathrm{}`$ – this is impossible because then by L’Hôpital’s rule $`limr=limf^{\prime \prime }/f^{}=limg/r=0`$. Therefore $`r(x_2)`$ must be sandwiched in the interval $`\sqrt{g_{\mathrm{}}+\sigma (x_2)}<r(x_2)<\sqrt{g_{\mathrm{}}\sigma (x_2)}`$. Since $`x_2`$ is arbitrarily large and $`lim\sigma =0`$, we conclude that $`limr=\sqrt{g_{\mathrm{}}}`$. The asymptotics of $`f`$ given in (6.1) follows immediately from this.
Finally, inserting the derived leading asymptotic behavior of $`A`$ and $`f`$ into Eq.(2.14c), we obtain $`C^{}/C2M_{\mathrm{}}/x`$, from which the asymptotics of $`C`$ follows trivially. ∎
## 7 Solutions that crash
###### Proposition 7.1.
If the $`b`$-orbit crashes at some $`\overline{x}`$ then $`g(x)>0`$ for $`x`$ near $`\overline{x}`$.
###### Proof.
Suppose that $`g(x)<0`$ for all $`x<\overline{x}`$, so $`AC^2(x)>1`$ for all $`x<\overline{x}`$. We have from (2.18) that $`g^{}>AC^2xf_{}^{}{}_{}{}^{2}>xf_{}^{}{}_{}{}^{2}`$. Integrating this inequality from some $`x_1>0`$ to some $`x_2<\overline{x}`$, we obtain
$$x_1_{x_1}^{x_2}f_{}^{}{}_{}{}^{2}𝑑x<_{x_1}^{x_2}xf_{}^{}{}_{}{}^{2}𝑑x<g(x_2)g(x_1)<\alpha ^21,$$
(7.1)
which implies (by the Cauchy-Schwartz inequality) that $`f`$ is bounded.
Next, $`A(\overline{x})=0`$, $`AC^2>1`$, implies that $`lim_{x\overline{x}^{}}C=\mathrm{}`$; moreover, by (2.14c) $`(\mathrm{ln}C)^{}<xf^2/A`$, hence $`xf^2/A`$ is not integrable near $`\overline{x}`$. Since $`f`$ is bounded, this shows that $`1/A`$ is not integrable near $`\overline{x}`$. But from (2.18), $`g^{}>C^2(1A)/x=AC^2(1A)/(xA)>1/(2xA)`$, so $`g^{}`$ is not integrable near $`\overline{x}`$, which contradicts the fact that $`g`$ is a bounded function. ∎
The importance of Proposition 7.1 derives from Lemma 2.6b which says that if $`g>0`$ then rotation stops. The main result of this section is the crash theorem which states that if an orbit has bounded rotation and crashes, then nearby orbits also have similarly bounded rotation. The precise statement is given in Proposition 7.2. Since we consider more than one orbit in this section, we use the notation $`A(x,a)`$ to denote the value of $`A`$ at $`x`$ for the $`a`$-orbit, etc.
###### Proposition 7.2 (Crash Theorem).
If the $`b`$-orbit crashes at $`x=\overline{x}`$ and
if $`(k1/2)\pi <\theta (x,b)<k\pi `$, $`k1`$, for $`x`$ near $`\overline{x}`$, then nearby orbits have rotation $`<k\pi `$ for $`x\overline{x}`$;
if $`k\pi <\theta (x,b)<(k+1/2)\pi `$, then nearby orbits have rotation $`<(k+1/2)\pi `$.
###### Proof.
Part (a): Suppose the $`b`$-orbit crashes in Q3 or Q1. By Proposition 7.1, $`g(x_1,b)>0`$ for some $`x_1<\overline{x}`$ with $`(k1/2)\pi <\theta (x_1,b)<k\pi `$; hence, for $`a`$ sufficiently near $`b`$ we have $`g(x_1,a)>0`$ with $`(k1/2)\pi <\theta (x,a)<k\pi `$. By Lemma 2.6b, $`\theta (x,a)<k\pi `$ for all $`x>x_1`$. Part (b): This case is much more difficult and will require several auxiliary results. It follows from part (a) that nearby orbits have rotation $`<(k+1)\pi `$; we must prove a much more difficult result, namely that nearby orbits have rotation $`<(k+1/2)\pi `$. ∎
###### Remark.
It is clear from numerical observations that no $`a`$-orbit crashes in Q2 or Q4; however, that appears to be quite difficult to prove. Moreover, one can easily construct orbit segments that start, for example, at $`x=1`$ with $`f=5,f^{}=0,A=0.2,C=3`$, say, that crash in Q4. Such orbit segments have $`lim_{x\overline{x}^{}}f^{}(x)=\mathrm{}`$. Nevertheless, the next lemma shows that $`Af_{}^{}{}_{}{}^{2}`$ remains bounded at crash.
###### Lemma 7.3.
If an $`a`$-orbit is defined for $`x<x_2,ff^{}(x)<0`$ for $`x_1<x<x_2`$, $`f^2(x_1)<B`$, and $`f^{}(x_1)=0`$ then $`Af_{}^{}{}_{}{}^{2}(x)\mathrm{max}(B,\alpha ^2/3)`$. In particular, if an orbit crashes in Q2 or Q4, $`lim_{x\overline{x}^{}}A(x)f^{}(x)=0`$.
###### Proof.
We set $`q=Af_{}^{}{}_{}{}^{2}`$ and then compute that
$$xq^{}=(3+x^2f_{}^{}{}_{}{}^{2}+x^2C^2f^2)qf_{}^{}{}_{}{}^{2}+2xff^{}+x^2f^2f_{}^{}{}_{}{}^{2}2AC^2xff^{}.$$
(7.2)
Note that $`q0`$ and all terms on the right side of (7.3) are negative except for the last two. If $`q>B`$, we combine the term $`qx^2f_{}^{}{}_{}{}^{2}`$ with $`x^2f^2f_{}^{}{}_{}{}^{2}`$; clearly, $`x^2f^2f_{}^{}{}_{}{}^{2}qx^2f_{}^{}{}_{}{}^{2}=(f^2q)x^2f_{}^{}{}_{}{}^{2}0`$. Next, we combine the term $`qx^2f_{}^{}{}_{}{}^{2}C^2`$ with $`2xff^{}AC^2`$ to get $`AC^2(y^22y)`$ where $`y=xff^{}`$; the maximum value of this expression occurs when $`y=1`$ and that value is $`AC^2\alpha ^2`$ by Lemma 2.4. Hence, if $`q\alpha ^2/3`$, then $`q(x^2f_{}^{}{}_{}{}^{2}C^2)3q2xff_{}^{}{}_{}{}^{2}AC^20`$. Thus, $`q\mathrm{max}(B,\alpha ^2/3)`$ implies that $`q^{}<0`$; consequently, $`Af_{}^{}{}_{}{}^{2}(x)\mathrm{max}(B,\alpha ^2/3)`$. Since $`AAf_{}^{}{}_{}{}^{2}=(Af^{})^2`$, and $`Af_{}^{}{}_{}{}^{2}`$ is bounded and $`limx\overline{x}^{}A(x)=0`$, $`lim_{x\overline{x}^{}}(A(x)f^{}(x))^2=0`$, hence $`lim_{x\overline{x}^{}}A(x)f^{}(x)=0`$. ∎
We can now discuss the strategy of the proof of part (b) of Proposition 7.2. We want to show that if an orbit is sufficiently close to an orbit that crashes in Q4 then it must either crash or exit Q4 to Q1 (the case in which the orbit crashes in Q2 is completely symmetric). To that end, let $`v(x)=A(x)f^{}(x)`$. We will prove that $`v(x,a)`$ goes to $`0`$ if $`a`$ is sufficiently close to $`b`$ and $`f(x,a)>0`$. This means either $`f^{}=0`$ and hence the orbit is exiting Q4 to Q1, or $`A=0`$, that is, the orbit is crashing in Q4. Note that $`v^{}(x)=(2Af^{}xf+xAC^2f+x^2Af_{}^{}{}_{}{}^{2}+x^2f^2f^{}AC^2)/x=v(2+x^2f_{}^{}{}_{}{}^{2}+x^2f^2C^2)/x+fg>fg`$. We know that $`v(x,b)`$ goes to $`0`$ at crash so nearby orbits will also have $`v`$ small for $`x`$ near $`\overline{x}`$. We will show that $`f`$ and $`g`$ are both uniformly bounded away from $`0`$ in an interval about $`\overline{x}`$. That is, the size of the interval and the bounds work for all $`a`$ near $`b`$. That is enough to force $`v`$ positive. The most technical part of the proof involves showing that nearby orbits stay in Q4 long enough to have $`v`$ go positive. Since $`f^{}`$ goes to $`\mathrm{}`$ at crash, nearby orbits have $`f^{}`$ large also. Now, (2.14a) can be written as $`xAf^{\prime \prime }+(1+Ax^2f^2)f^{}xgf=0`$; moreover, to get to Q3 orbits must pass through $`xf(x)<1`$ which means that the coefficient of $`f^{}`$, $`(1+Ax^2f^2)`$, is positive. That is enough to bound $`f^{}`$.
The details of the proof, especially Lemma 7.5, are tedious. We will restrict ourselves to an interval $`0.99\overline{x}<x<1.01\overline{x}`$ and replace $`x`$ by $`\overline{x}`$ (whenever justified) in making estimates.
We show next that if the $`b`$-orbit crashes at $`x=\overline{x}`$ with rotation $`k\pi <\theta (x,b)<(k+1/2)\pi `$ then $`|\overline{x}f(\overline{x})|1`$.
###### Lemma 7.4.
If the $`b`$-orbit crashes at $`x=\overline{x}`$ with $`\theta (x,b)<(k+1/2)\pi `$ for all $`x<\overline{x}`$ and $`\theta (x,b)>k\pi `$ for $`x`$ near $`\overline{x}`$, then $`|\overline{x}f(\overline{x})|1`$, in particular $`f(\overline{x})0`$.
###### Proof.
The assumption on $`\theta (x,b)`$ tells us that the orbit lies in Q2 or Q4 for $`x`$ near $`\overline{x}`$. For simplicity of exposition we only discuss the case of Q4, i.e., $`f(x)0,f^{}(x)0`$. In particular, $`f`$ is a monotone function and hence has a limit at $`\overline{x}`$. Thus, $`h(x)=xf(x)`$ is continuous; in particular, if we suppose that $`\overline{x}f(\overline{x})<1`$, then $`h(x)<1`$ for $`x`$ near $`\overline{x}`$. Since $`A(\overline{x})=0`$, we get from (2.14c) that $`xAC^{}=C(A1+x^2f^2)<0`$ for $`x`$ near $`\overline{x}`$. We conclude that $`C`$ is bounded above, hence $`lim_{x\overline{x}^{}}AC^2=0`$ and $`lim_{x\overline{x}^{}}g=1`$. Since $`g>0`$, the right hand side of Eq.(2.14a) is positive and hence $`x^2f^{}/C`$ is bounded and since $`C`$ is bounded we conclude that $`f^{}`$ is bounded; thus $`lim_{x\overline{x}^{}}Af_{}^{}{}_{}{}^{2}=0`$. Then, from (2.14b), $`xA^{}=1Ax^2f^2x^2(Af_{}^{}{}_{}{}^{2}+AC^2f^2)`$, we see that $`A^{}>0`$ near $`\overline{x}`$ so there is no crash. This is a contradiction so we conclude that $`\overline{x}f(\overline{x})1`$ and hence $`f(\overline{x})>0`$. ∎
###### Lemma 7.5.
There is a $`\gamma >0`$ such that $`h(x,a)=xf(x,a)>1/4`$ for all $`a`$ sufficiently near $`b`$ and $`\overline{x}<x<\overline{x}+\gamma `$.
###### Proof.
If the $`b`$-orbit crashes at $`x=\overline{x}`$ with rotation $`\theta (x,b)>k\pi `$, then there is a $`y`$ such that $`\theta (y,b)=k\pi `$. Let $`B=(f(y,b)+1)^2`$. By Proposition 7.3, if $`a`$ is sufficiently close to $`b`$, $`Af_{}^{}{}_{}{}^{2}`$ (along the $`a`$-orbit) is bounded in Q4 by $`D=\mathrm{max}(\alpha ^2/3,B)`$; $`D`$ is a uniform bound on $`Af_{}^{}{}_{}{}^{2}`$ in Q4 for all $`a`$ sufficiently near $`b`$. Next, choose $`x_1`$ such that $`0.99\overline{x}<x_1<\overline{x}`$ and such that $`A(x_1,b)<0.01,g(x_1,b)=2\tau >0`$, and $`h(x_1,b)>0.9`$; this is possible by Lemma 7.4 and Proposition 7.1. Then, for $`a`$ sufficiently near $`b`$ we have $`A(x_1,a)<0.02,g(x_1,a)>\tau >0,f(x_1,a)<f(x_1,b)+0.01/\overline{x}`$ and $`h(x_1)>3/4`$. We shall find a $`\gamma (0,0.01\overline{x})`$ that works for all $`a`$, that is, it satisfies $`h(x,a)>1/4`$ for all $`a`$ sufficiently near $`b`$ and $`\overline{x}<x<\overline{x}+\gamma `$. So let $`a`$ satisfy: i) $`Af_{}^{}{}_{}{}^{2}`$ (along the $`a`$-orbit) is bounded by $`D`$, ii) $`A(x_1,a)<0.02`$, iii) $`h(x_1,a)>3/4`$, and iv) $`g(x_1,a)>\tau >0`$. If $`h(x,a)>1/4`$ for all $`x<1.01\overline{x}`$ and all $`a`$ near $`b`$ we are done – let $`\gamma =0.01\overline{x}`$. Otherwise, we define $`x_2=x_2(a)`$, etc. by $`h(x_2)=3/4,h(x_3)=1/2,h(x_4)=1/4`$, where $`x_2,x_3`$, and $`x_4`$ are the largest values of $`x<1.01\overline{x}`$ with that property. For $`x>x_2`$ we have from (2.14a) $`xAf^{\prime \prime }=xgf(1+Ah^2)f^{}(1+Ah^2)f^{}f^{}/4`$ since $`h3/4`$ so $`f^{\prime \prime }f_{}^{}{}_{}{}^{3}/(4xAf_{}^{}{}_{}{}^{2})f_{}^{}{}_{}{}^{3}/(41.01\overline{x}D)`$ or $`f^{\prime \prime }/f_{}^{}{}_{}{}^{2}f^{}/(4.04\overline{x}D)`$. We now integrate the above from $`x_2`$ to $`x>x_3`$ to get
$$\frac{1}{f^{}(x)}\frac{1}{f^{}(x)}+\frac{1}{f^{}(x_2)}_{x_2}^x\frac{f^{\prime \prime }}{f_{}^{}{}_{}{}^{2}}𝑑x_{x_2}^x\frac{f^{}}{4.04\overline{x}D}𝑑x=\frac{f(x)f(x_2)}{4.04\overline{x}D}.$$
(7.3)
Now, $`f(x)f(x_3)`$, so $`f^{}(x)\frac{5\overline{x}D}{f(x_2)f(x_3)}\frac{5\overline{x}^2D}{h(x_2)h(x_3)}=20\overline{x}^2D`$. Using the uniform bound on $`f^{}`$ in the interval $`x_3xx_4`$, we have $`x_4x_3=(f(x_4)f(x_3))/f^{}(\xi )`$ for some $`\xi [x_3,x_4]`$. But $`(f(x_4)f(x_3))/f^{}(\xi )(h(x_4)h(x_3))/(\overline{x}f^{}(\xi ))1/80\overline{x}^3D`$ and hence we may take $`\gamma =1/80\overline{x}^3D`$. ∎
###### Lemma 7.6.
In the interval $`x_1<x<\overline{x}+\gamma `$, $`g(x,a)>\mathrm{min}(\tau ,0.9/h^2(\overline{x},b))`$.
###### Proof.
From (2.18) we have $`xg^{}=C^2(1A+x^2Af_{}^{}{}_{}{}^{2}x^2gf^2)C^2(1Ax^2gf^2)`$. Moreover, since $`A(x_1,a)<0.02`$ and $`xA^{}<1,A(x,a)=A(x_1,a)+A^{}(z)(xx_1)<0.02+1/z(0.02\overline{x})<0.04`$, so if $`g<0.96/h^2(x)`$ then $`g^{}>0`$. Since $`f(x_1,a)<f(x_1,b)+0.01/\overline{x},h(x,a)1.01\overline{x}f(x_1,a)<1.01(\overline{x}f(x_1,b)+0.01)<1.02\overline{x}f(x_1,b)`$, we have $`g^{}>0`$ if $`g(x_1,a)<0.9/h^2(\overline{x},b)`$. Thus, if $`\tau <g(x_1,a)<0.9/h(\overline{x},b),g^{}>0`$, and $`g(x,a)>\tau `$ in the interval $`x_1<x<\overline{x}+\gamma `$; if $`g(x_1,a)>0.9/h^2(\overline{x},b)`$, then $`g(x,a)>0.9/h^2(\overline{x},b)`$ for all $`x`$ in the interval $`x_1<x<\overline{x}+\gamma `$ because $`g`$ cannot cross that value from above. ∎
Note that the above lower bound on $`g`$ is uniform – it applies to all $`a`$ satisfying the conditions i) $`Af_{}^{}{}_{}{}^{2}`$ (along the $`a`$-orbit) is bounded by $`D`$, ii) $`A(x_1,a)<0.02`$, iii) $`h(x_1,a)>3/4`$, and iv) $`g(x_1,a)>\tau >0`$.
###### Lemma 7.7.
For all $`a`$ sufficiently near $`b`$, $`v(x,a)`$ goes to $`0`$ for some $`x<\overline{x}+\gamma `$.
###### Proof.
To show that $`v(x,a)`$ goes to $`0`$, we note that $`h(x,a)1/4`$ for all $`a`$ near $`b`$ and $`\overline{x}<x<\overline{x}+\gamma `$ by Lemma 7.5. Hence, $`f(x,a)=h(x,a)/x>1/4\overline{x}`$. By Lemma 7.6, $`g(x,a)>\mathrm{min}(\tau ,1/h(\overline{x})`$, hence $`v^{}1/4\overline{x}\mathrm{min}(\tau ,1/h(\overline{x})=\eta >0`$ for $`\overline{x}<x<\overline{x}+\gamma `$. Thus, $`v(\overline{x}+\gamma )v(\overline{x})=_{\overline{x}}^{\overline{x}+\gamma }v^{}𝑑x_{\overline{x}}^{\overline{x}+\gamma }\eta 𝑑x\eta \gamma `$. Let $`x_1`$ be chosen so that $`v(x_1,a)>\eta \gamma /2`$. Then, if $`a`$ is sufficiently close to $`b`$, $`|v(x_1,a)v(x_1,b)|>\eta \gamma /2`$ so $`v(x_1,a)>\eta \gamma `$. For such $`a`$ we then have $`v(\overline{x}+\gamma ,a)>v(\overline{x},a)+\eta \gamma `$ and $`v(\overline{x},a)>v(x_1,a)>\eta \gamma `$ because $`v^{}>fg>0`$; thus, $`v(\overline{x}+\gamma ,a)>0`$. ∎
We now complete the proof of Proposition 7.2.
###### Proof of Proposition 7.2 b).
Suppose that the $`b`$-orbit crashes at $`x=\overline{x}`$ with $`\theta (x,b)<(k+1/2)\pi `$ for all $`x<\overline{x}`$ and $`\theta (x,b)>k\pi `$ for $`x`$ near $`\overline{x}`$. For $`a`$ near $`b`$ there is an $`x<\overline{x}+\gamma `$ with $`v(x,a)=0`$ by Lemma 7.7. Since $`x<\overline{x}+\gamma `$, $`h(x)>1/4`$, i.e., $`f(x,a)>0`$, so the $`a`$-orbit crashes, $`A(x,a)=0`$, or exits Q4 to Q1 (or Q2 to Q3), $`f^{}(x,a)=0`$, never to return. In either case, the $`a`$-orbit has rotation $`\theta (x,a)<(k+1/2)\pi `$. ∎
## 8 Proof of the main theorem
###### Proof of Theorem 2.1.
Let $`X_n=\{a>0|\theta (x,a)<(n+1/2)\pi `$ for all $`x`$ for which the a-orbit is defined}. Note that $`X_{n1}X_n`$ and $`X_0\mathrm{}`$ by Proposition 5.1 and hence, $`X_n\mathrm{}`$. Also note that $`b_{n+1}>0`$ is a lower bound for $`X_n`$ by Proposition 4.1; hence, $`X_n`$ has a greatest lower bound $`a_n=inf(X_n)b_{n+1}>0`$. We will show that the $`a_n`$-orbit is a globally regular solution and $`n\pi <\theta (x,a_n)<(n+1/2)\pi `$ for large $`x`$.
We first show that $`a_nX_n`$, i.e., $`a_n`$ is the smallest element in $`X_n`$. If $`\theta (x,a_n)>(n+1/2)\pi `$ for some $`x`$ then $`\theta (x,a)>(n+1/2)\pi `$ for all $`a`$ near $`a_n`$ so $`aX_n`$ for these $`a`$’s and this contradicts the fact that $`a_n`$ is the greatest lower bound of $`X_n`$. Thus, $`a_nX_n`$. In particular, the $`a_n`$-orbit has bounded rotation.
Next we show that the $`a_n`$-orbit does not crash. Recall from Proposition 7.1 that if the $`a_n`$-orbit crashes at $`x=\overline{x}`$ then $`g(x,a_n)>0`$ for $`x`$ near $`\overline{x}`$. If the $`a_n`$-orbit crashes in Q1 or Q3, that is, if $`\theta (x,a_n)<n\pi `$ for $`x`$ near $`\overline{x}`$ then $`\theta (x,a)<n\pi `$ and $`g(x,a)>0`$ for all $`a`$ near $`a_n`$ which implies by Lemma 2.6b that the $`a`$-orbit must have $`\theta (x,a)<n\pi `$ for all $`x`$. Thus, $`aX_n`$ for all $`a`$ near $`a_n`$ and this contradicts the fact that $`a_n`$ is the greatest lower bound of $`X_n`$.
Similarly, if the $`a_n`$-orbit crashes in Q2 or Q4, that is, at some $`\overline{x}`$ with $`(n+1/2)\pi >\theta (\overline{x},a_n)>n\pi `$, then by the crash lemma $`(n+1/2)\pi >\theta (x,a)`$ for all $`x`$ in the domain of definition of the $`a`$-orbit for all $`a`$ near $`a_n`$ and this contradicts the fact that $`a_n`$ is the greatest lower bound of $`X_n`$.
Thus, the $`a_n`$-orbit is defined for all $`x`$ and hence is a globally regular solution by Propositions 6.1. Also, by Proposition 6.2, the $`a_n`$-orbit is in Q2 or Q4 for large $`x`$. It remains to prove that $`\theta (x,a_n)>n\pi `$ for large $`x`$. Suppose that $`\theta (x,a_n)<n\pi `$ for large $`x`$. By Lemma 6.3 we have that $`g(x,a_n)>0`$ for large $`x`$ and hence, $`g(x,a)>0`$ for all $`a`$ near $`a_n`$. Then, by Lemma 2.6b the $`a`$-orbit must have $`\theta (x,a)<n\pi `$ for all $`x`$ and thus $`aX_n`$, and this contradicts the fact that $`a_n`$ is the greatest lower bound for $`X_n`$. This completes the proof of Theorem 2.1. ∎
## Acknowledgment
We would like to thank the Mathematisches Forschungsinstitut in Oberwolfach for supporting this project under of the Research in Pairs program. P.B. was supported in part by the KBN grant 2 P03B 010 16.
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# Compton dragged gamma–ray bursts: the spectrum
## 1 Introduction
We have recently proposed (Lazzati et al. 2000, hereafter Paper I) that the gamma–ray burst (GRB) phenomenon originates from the interaction of a relativistic fireball with a dense photon environment, leading to Compton drag. On one hand this is an inevitable effect if the progenitors of GRBs are massive stars which are about to explode or have just exploded as supernovae; on the other hand this mechanism greatly alleviates the efficiency problem faced by the standard internal shock scenario (Lazzati, Ghisellini & Celotti 1999; Panaitescu, Spada & Meszaros 1999; Kumar 1999). In Paper I we have discussed the basic Compton drag scenario, showing how this process can convert bulk motion energy directly into radiation with a remarkable high efficiency and, on the basis of simple estimates, how the resulting spectrum should peak, in a $`\nu F(\nu )`$ representation, around $``$1 MeV, as observed.
Here we quantitatively and self–consistently estimate the predicted spectrum, assuming that the fireball propagates in a funnel inside a massive star, and show that, independently of the details of the model, it satisfactorily resembles what observed. Since the funnel walls emit a blackbody spectrum and the scattered photons are boosted by the square of the Lorentz factor ($`\mathrm{\Gamma }`$) of the fireball, the local spectrum has a blackbody shape, at a temperature enhanced by $`\mathrm{\Gamma }^2`$. However, the observed spectrum, convolution of all the locally emitted spectra, is not a blackbody, due to four main effects: i) the funnel walls would not be at a uniform temperature, but there should be a gradient between the internal and external parts; ii) if the Compton drag process is efficient, the fireball decelerates; iii) the very high energy emission produced in the internal regions can interact with the ambient photons, producing electron–positron pairs; iv) the fireball may become optically thin to scattering outside the funnel, where the ambient photons are characterized by the same temperature, but their energy density is progressively diluted with distance.
## 2 Basic assumptions
We postulate that the fireball propagates with a bulk Lorentz factor $`\mathrm{\Gamma }`$ inside a funnel cavity, whose walls emit blackbody radiation at a temperature $`T`$, of conical shape with semi–aperture angle $`\psi `$. The calculation starts at the distance $`z_0`$, assumed to be the end of the acceleration phase and, for consistency, we verify that the power emitted at $`z<z_0`$ is negligible.
We assume that the fireball is and remains cold in the comoving frame. At $`z_0`$, in fact, the internal energy has been already used to accelerate the fireball, and thus protons are sub–relativistic. On the other hand leptons might be still hot at $`z_0`$ and/or being re-heated when the bulk scattering process starts to be efficient. However, in a few (Compton) cooling timescales they would reach the (sub-relativistic) Compton temperature. It is thus reasonable to treat also the leptonic component as cold in the estimate of both the dynamics and resulting spectrum.
The initial Lorentz factor (at $`z_0`$) is indicated as $`\mathrm{\Gamma }_0`$, and the fireball energy is therefore $`E_f=\mathrm{\Gamma }_0M_fc^2`$, where $`M_f`$ is its rest mass. For simplicity, the dependence of the temperature on $`z`$, between $`z_0`$ and the radius of the star $`z_{}`$, has been parameterized by a power law:
$$T(z)=T_0\left(\frac{z}{z_0}\right)^b=T_{}\left(\frac{z}{z_{}}\right)^b$$
(1)
where $`T_{}`$ is the temperature at the top of the funnel.
Inside it, we approximate the local radiation energy density of the ambient photons as $`U(z)=aT^4(z)`$. Beyond $`z_{}`$, and in the region where the fireball remains optically thick (i.e. for $`z<z_T`$, see below), $`U(z)`$ is characterized by the same temperature, but decreases. As the relevant quantity is the amount of radiation which is indeed scattered by the fireball, we parameterize the dependence on $`z`$ of the product $`U(z)\times `$ (the scattering rate) as $`(z/z_{})^g`$.
We consider $`g`$ a free parameter. Inside the funnel $`g=0`$, while outside it a value $`g>2`$ can account for a decrease in the scattering rate due to the changing of the typical scattering angle (photons come preferentially at smaller angles as $`z`$ increases). As the scattering rate is $`(1\beta \mathrm{cos}\theta )`$, where $`\theta `$ is the angle between the photon and the electron directions, far from the star surface $`(1\beta \mathrm{cos}\theta )(z/z_{})^2`$, corresponding to $`g4`$. Furthermore some of the radiation produced by the massive star could be reflected and re–isotropized by scattering material, of unknown radial density profile, likely surrounding the massive star progenitor. In particular if this forms a wind with a $`z^2`$ profile, the energy density of the re–isotropized radiation scales as $`z^3`$, and dominates the seed photon distribution at large distances. In this case $`U(z)\times `$(the scattering rate) can have a complex profile, being flat in the vicinity of the surface of the star, then decreasing as $`z^2`$ and as $`z^4`$ for increasing $`z`$, to become flatter when the component associated with the re–isotropized photons dominates. It is also possible that, as a result of intermittent stellar activity, the stellar wind is not continuous. In this case a single shell may dominate the scattering, producing a homogeneous and isotropic scattered radiation field, dominating the total radiation energy density beyond some critical distance.
The distance $`z_T`$ at which the fireball becomes optically thin to scattering is
$$z_T=\left(\frac{\sigma _TE_f}{\pi \psi ^2m_pc^2\mathrm{\Gamma }_0}\right)^{1/2}3.7\times 10^{14}\psi _1^1E_{f,51}^{1/2}\mathrm{\Gamma }_{0,2}^{1/2}\mathrm{cm}$$
(2)
where the conventional representation $`Q=Q_x10^x`$ and c.g.s. units are adopted. It is then likely that the fireball becomes transparent at $`z>z_{}`$ (since the radius of red supergiants is $`z_{}\stackrel{<}{}10^{13}`$ cm).
As long as the fireball is opaque to scattering, the interaction with photons boosts their energy by a factor $`2\mathrm{\Gamma }^2`$. Therefore the (local) total energy emitted by the fireball through the Compton drag process (over a distance $`dz`$) is
$$dE(z)=\mathrm{\hspace{0.17em}2}\pi \psi ^2z^2aT_0^4\left(\frac{z}{z_0}\right)^{4b}\mathrm{\Gamma }^2dzz<z_{}$$
(3)
$$dE(z)=\mathrm{\hspace{0.17em}2}\pi \psi ^2z^2aT_{}^4\left(\frac{z}{z_{}}\right)^g\mathrm{\Gamma }^2dzz>z_{}.$$
(4)
The factor 2 in front of the RHS of these equations takes into account that the preferred scattering angle is $`90^{}`$, corresponding to an average energy boost of $`2\mathrm{\Gamma }^2`$.
Let us now consider the spectral shape. For this it is convenient to use dimensionless photon energies and temperatures, defined as $`xh\nu /(m_ec^2)`$ and $`\mathrm{\Theta }kT/(m_ec^2)`$, respectively.
The resulting Compton spectrum has a blackbody shape, of effective temperature $`T_c=2\mathrm{\Gamma }^2T`$ (or $`\mathrm{\Theta }_c=2\mathrm{\Gamma }^2\mathrm{\Theta }`$), i.e. the local spectral distribution produced within $`dz`$ is given by:
$`dE(z,x)`$ $`=`$ $`\pi ^2\psi ^2{\displaystyle \frac{z^2}{\mathrm{\Gamma }^6}}m_ec^2\left({\displaystyle \frac{m_ec}{h}}\right)^3{\displaystyle \frac{x^3}{e^{x/\mathrm{\Theta }_c}1}}dz;`$ (5)
$`z<z_{}`$
$`dE(z,x)`$ $`=`$ $`\pi ^2\psi ^2{\displaystyle \frac{z^2(z/z_{})^g}{\mathrm{\Gamma }^6}}m_ec^2\left({\displaystyle \frac{m_ec}{h}}\right)^3{\displaystyle \frac{x^3}{e^{x/\mathrm{\Theta }_{c,}}1}}dz;`$ (6)
$`z>z_{},`$
where $`\mathrm{\Theta }_{c,}=2\mathrm{\Gamma }^2\mathrm{\Theta }_{}`$. Equations (5) and (6) are correctly normalized, i.e. the integrated energies correspond to those expressed in (3) and (4).
## 3 The fireball dynamics
As long as the fireball remains optically thick for scattering and this occurs in the Thomson regime, the dynamics (deceleration) of the fireball due to the radiative drag, obeys:
$$M_fc^2\frac{d\mathrm{\Gamma }}{dz}=\mathrm{\hspace{0.17em}2}\pi \psi ^2z^2aT^4\mathrm{\Gamma }^2.$$
(7)
Assuming the temperature profile of equation (1) we obtain:
$`\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }_0}{1+2\pi \psi ^2aT_0^4\mathrm{\Gamma }_0^2z_0^3[(z/z_0)^{34b}1]/[E_f(34b)]}};`$ (8)
$`z_0<z<z_{},`$
and thus the deceleration radius, $`z_d`$, defined as the distance at which $`\mathrm{\Gamma }`$ is halved, corresponds to:
$`z_d`$ $`=`$ $`z_o\left[1+{\displaystyle \frac{E_f(34b)}{2\pi \psi ^2aT_0^4\mathrm{\Gamma }_0^2z_0^3}}\right]^{1/(34b)}`$ (9)
$`=`$ $`z_o\left[1+{\displaystyle \frac{E_f(34b)}{2\pi \psi ^2aT_{}^4\mathrm{\Gamma }_0^2(z_{}/z_0)^{4b}z_0^3}}\right]^{1/(34b)}`$
Beyond $`z_d`$, the Lorentz factor decreases with distance as a power law, whose slope is determined by the temperature profile.
Outside the star radius ($`z>z_{}`$) the Lorentz factor follows:
$`\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }_{}}{1+2\pi \psi ^2aT_{}^4\mathrm{\Gamma }_0\mathrm{\Gamma }_{}z_{}^3[(z/z_{})^{3g}1]/[E_f(3g)]}};`$ (10)
$`z_{}<z<z_T.`$
Note that Klein-Nishina effects are important for incoming photon energies such that $`x\mathrm{\Gamma }>1`$, i.e. when $`\mathrm{\Theta }>1/(3\mathrm{\Gamma })`$. For simplicity, we neglect interactions in this regime when calculating $`\mathrm{\Gamma }(z)`$, but we assume no scattering events when $`\mathrm{\Theta }>1/(3\mathrm{\Gamma })`$ in calculating the spectrum. This simplification is justified as long as most of the fireball energy is lost in the Thomson scattering regime (see Fig. 2, which shows that $`\mathrm{\Gamma }`$ starts to decrease at distances where the temperature is small enough to ensure scatterings entirely in the Thomson regime).
When the fireball becomes optically thin, the amount of scattered photons is correspondingly reduced, and the process becomes less efficient. As shown by equation (2), this is likely to happen at some distance from the star surface, where the photon density is also reduced, thus further decreasing the efficiency of the process. In the numerical calculations we have however included the optically thin scattering regime, and one can see its contribution in Fig. 1 (dotted line).
## 4 Pair production
A further effect which may strongly affect both the observed spectrum and the dynamics of the fireball is the production of electron–positron pairs through photon-photon interactions. Let us thus consider in turn the role of scattered and funnel radiation as seed photons for this process.
### 4.1 Interaction among photons in the beam
The threshold energy for interaction between photons of energies $`x`$ and $`x_T`$ is $`x_T>2/[x(1\mathrm{cos}\theta )]4\mathrm{\Gamma }^2/x`$, where all quantities are calculated in the observer frame. The latter expression takes into account that the high energy photons produced by the Compton drag are highly collimated, within a typical angle $`\mathrm{sin}\theta 1/\mathrm{\Gamma }`$. As the bulk of the scattered photons have energies $`x2\mathrm{\Gamma }^2(3\mathrm{\Theta })`$, pair production would occur if $`\mathrm{\Gamma }\mathrm{\Theta }>1/3`$.
However this also implies that the scattering process is in the Klein Nishina regime, and we can therefore conclude that photon–photon collisions among photons in the beam can only affect the high energy tail of the spectrum produced at each radius, while the emission at the peak is unaltered. We therefore neglect this effect.
### 4.2 Interaction between beam photons and funnel radiation
The interaction between the $`\gamma `$–rays produced by the Compton drag process and photons emitted by the funnel walls would occur at large angles, resulting in an average energy threshold $`x_T>1/x`$. Since $`x\mathrm{\Gamma }_0`$, this absorption mechanism would be significant as long as the funnel walls produce a sufficient number of photons with energies $`x_T>1/\mathrm{\Gamma }_0`$.
Let us then estimate the photon–photon optical depth $`\tau _{\gamma \gamma }`$, by integrating the product of the photon–photon cross section $`\sigma _{\gamma \gamma }(x,x_T)`$ and the photon density above threshold $`n_\gamma (x)`$ over the $`\gamma `$–ray path, i.e. from the site of creation, $`z_1`$, to infinity, and over the photon energies:
$$\tau _{\gamma \gamma }(z_1,x)=_{x_T}^{\mathrm{}}𝑑x^{}_{z_1}^{\mathrm{}}\sigma _{\gamma \gamma }(x^{},x)n_\gamma (z,x^{})𝑑z.$$
(11)
Since $`\sigma _{\gamma \gamma }(x^{},x)`$ is peaked at the threshold energy, equation (11) can be simplified (Svensson 1984, 1987) as
$`\tau _{\gamma \gamma }(z_1,x)={\displaystyle \frac{\sigma _T}{5m_ec^2}}{\displaystyle _{z_1}^{\mathrm{}}}x_TU(z,x_T)𝑑z,`$ (12)
where $`U(z,x_T)=m_ec^2n_\gamma (z,x_T)`$ is the photon energy density at threshold, at the location $`z`$, i.e.
$$U(z_1,x_T)=\frac{8\pi h}{c^3}\left(\frac{m_ec^2}{h}\right)^4\frac{x_T^3}{\mathrm{exp}[x_T/\mathrm{\Theta }(z_1)]1}.$$
(13)
The radiation flux produced at the location $`z_1`$ is then decreased by the factor $`\mathrm{exp}[\tau _{\gamma \gamma }(z_1,x)]`$ while crossing the funnel.
The absorbed radiation will be reprocessed by the pairs, and re–distributed in energy. Each electron and positron will have an energy $`\gamma x/2`$ at birth, and will cool due to the Compton drag process. The positrons will then annihilate in collisions with the electrons in the fireball, producing a Doppler blueshifted annihilation line at $`x\mathrm{\Gamma }`$. We have neglected these reprocessing mechanisms, since, as can be seen in Fig. 1, the amount of energy absorbed in $`\gamma `$$`\gamma `$ collisions is small, amounting to a few per cent at most.
## 5 The spectrum
The observed total spectrum can be computed by integrating equations (3) and (4) over $`z`$, taking into account photon–photon absorption. The contribution produced within the star is given by:
$`E(x)`$ $`=`$ $`\pi ^2\psi ^2m_ec^2\left({\displaystyle \frac{m_ec}{h}}\right)^3{\displaystyle _{z_0}^z_{}}{\displaystyle \frac{z^2}{\mathrm{\Gamma }^6}}{\displaystyle \frac{x^3e^{\tau _{\gamma \gamma }(z,x)}}{e^{x/\mathrm{\Theta }_c}1}}𝑑z;`$ (14)
$`z_0<z<z_{},`$
while beyond $`z_{}`$ the number of target photons able to interact with high energy $`\gamma `$–rays to produce pairs is negligible, and thus, ignoring photon–photon absorption, we obtain:
$`E(x)`$ $`=`$ $`\pi ^2\psi ^2m_ec^2\left({\displaystyle \frac{m_ec}{h}}\right)^3{\displaystyle _z_{}^{z_T}}{\displaystyle \frac{z^2(z/z_{})^gx^3}{\mathrm{\Gamma }^6(e^{x/\mathrm{\Theta }_{c,}}1)}}𝑑z;`$ (15)
$`z_{}<z<z_T.`$
In Fig. 1 we show three examples of the predicted spectrum corresponding to different values of the initial bulk Lorentz factors. To illustrate the main features of the model and the importance of photon–photon absorption, this is calculated both with and without the photon–photon absorption term. Together with the total spectrum, the separate contributions for $`z<z_{}`$ and for $`z_T<z<z_{}`$ are reported. In Fig. 2 we show the corresponding $`\mathrm{\Gamma }`$ profiles. The effect of the star surface temperature (and of the entire funnel, since the parameter $`b`$ is assumed to be the same for all cases) can be clearly seen in Fig. 3. Note the $`\nu ^{1/2}`$ power law shape in the X–ray band for the high temperature case. The extension of this power law branch depends on the value of $`g`$. In the case shown ($`g=2`$) the radiation energy density outside the funnel remains sufficiently large to cause the deceleration of the fireball, and this is responsible for the power law tail between 10 and 100 keV. For larger $`g`$ the extension of this power law would decrease. This effect can also be seen for the high $`\mathrm{\Gamma }_0`$ case in Fig. 1.
In order to determine the general features of the predicted spectrum and thus assess its robustness against the parameters of the model, we also derived analytical (although approximated) expressions for the spectral energy distribution.
### 5.1 Analytical approximations
First, let us approximate the blackbody spectral form with its Rayleigh–Jeans part, and let us neglect photon–photon absorption. In this case, for $`x<6\mathrm{\Theta }\mathrm{\Gamma }^2`$ we have:
$`dE(z,x)`$ $``$ $`{\displaystyle \frac{T}{\mathrm{\Gamma }^4}}z^2dz;\mathrm{for}z_o<z<z_{}`$ (16)
$``$ $`{\displaystyle \frac{T}{\mathrm{\Gamma }^4}}z^{2g}dz;\mathrm{for}z_{}<z<z_T.`$
Three regimes occur at different distances:
$`𝐳_0<z<z_d`$: — in this case $`\mathrm{\Gamma }=`$ const, and integration over $`z`$ yields:
$$E(x)x^{(33b)/b};\mathrm{for}z>z_d$$
(17)
which, for $`b=0.5`$, gives $`E(x)x^3`$.
$`𝐳_d<z<z_{}`$: — here $`\mathrm{\Gamma }`$ decreases as $`(z/z_0)^{(34b)}`$ and thus:
$$E(x)x^{3(1b)/(67b)}\mathrm{for}z_d<z<z_{};$$
(18)
which, for $`b=0.5`$, results in $`E(x)x^{3/5}`$.
$`𝐳_{}<z<z_T`$: — at these distances the ambient radiation energy density decreases as $`(z/z_0)^g`$. If $`\mathrm{\Gamma }`$ remains constant (= $`\mathrm{\Gamma }_{}`$), the spectrum $`E(x)x^2`$, while, for $`\mathrm{\Gamma }`$ decreasing as $`\mathrm{\Gamma }(z/z_0)^{(3g)}`$
$$E(x)x^{1/2}\mathrm{for}z_{}<z<z_T,$$
(19)
which is independent of $`g`$.
In conclusion, in the case of efficient Compton drag, and independently of the particular choice of parameters, the predicted spectrum is always characterized (in order of decreasing energy) by: a steep high energy tail; a first break flagging the deceleration of the fireball; a second break corresponding to radiation produced at the top of the funnel – above which the temperature of the ambient photons remains constant; a third break, below which the spectrum $`x^{1/2}`$, corresponding to the deceleration of the fireball due to the isothermal photon bath; and finally a fourth break, below which the spectrum $`F(x)x^2`$. One obtains such a hard spectrum, instead of the familiar slope $`F(x)x`$ corresponding to scatterings of isotropically distributed electrons and seed photons, because only the photons scattered along the forward direction are observed <sup>1</sup><sup>1</sup>1This can be seen by integrating Eq. 7.23 of Rybicki & Lightman (1979), in the angle range \[$`0<\theta _1<1/\mathrm{\Gamma }`$\].
## 6 Discussion
If the fireball propagates in a dense photon environment the Compton drag effect must necessarily be taken into account, and it may even be the dominant emission mechanism, able to decelerate the fireball without the need of internal shocks and without invoking the build–up of large magnetic fields.
In this letter we have shown that the predicted spectrum, rather than being simply a black body spectrum boosted in energy, has a complex shape, with power law segments corresponding to the decrease in temperature of the funnel, deceleration of the fireball, and dilution of the radiation energy density as the fireball propagates outside the funnel while remaining optically thick.
The general features of the predicted spectrum qualitatively agree with observations, since they can explain the steep power law high energy tail, the peak of the emission, and a hard tail in the X–ray band. The latter feature is particularly interesting, since other models made different predictions. In the standard internal shock synchrotron model, in fact, the spectrum cannot be harder than $`\nu ^{1/3}`$ in the thin part, and it is very unlikely that self–absorption can take place in the X–ray band (Granot, Piran & Sari 2000). This would in fact imply a huge density of relativistic particles, making the inverse Compton effect largely dominate the total radiation output. This radiation would be emitted at higher and yet unobserved frequencies, and would then worsen the already severe efficiency problem.
In the quasi–thermal Comptonization model, on the other hand, the typical predicted spectral shape in the X–ray band is $`\nu ^0`$, down to the typical frequencies of the seed soft photons, i.e. the IR–optical band (Ghisellini & Celotti 1999; Meszaros & Rees 2000).
The existing observations of a significant fraction of burst spectra harder than $`\nu ^{1/3}`$ (Preece et al., 1999a,b; Crider et al., 1997) are therefore already a challenge to existing models, and may suggest a Compton drag origin of this portion of the spectrum. However the situation is not already a clear–cut because, to receive enough photons to study the spectral shape, integration times are much longer than the dynamical time–scales of the system, with the spectrum rapidly evolving in time. More sensitive instruments, such as the Burst Alert Telescope (BAT, a coded mask detector more sensitive than BATSE) onboard the foreseen Swift satellite will probably overcome this limitation.
We must also stress that the Compton drag scenario is not alternative to the more conventional internal shock one. Indeed, the front of the fireball will decelerate first, plausibly causing subsequent un-decelerated parts to shock even if the central engine is working in a continuous way. This would produce additional radiation, either by the synchrotron and inverse Compton processes or by quasi–thermal Comptonization, depending on the details of the particle acceleration mechanism (see Ghisellini & Celotti 1999). We then expect spectral evolution: since the latter radiation mechanisms produce a steeper low energy tail, a hard–to–soft transition (i.e. from $`\nu ^2`$ to $`\nu ^{1/3}`$ or $`\nu ^0`$) would occur.
In this paper, we have considered the illustrative case of a single fireball moving out through an extended stellar envelope, along a funnel which is empty of matter but pervaded by thermal radiation from the funnel walls. The fireball itself (for typical parameters) remains optically thick until it expands beyond the stellar surface. A burst with complex time-structure could be modeled by a series of fireballs or expanding shells. However, in this more general case, the later shells would suffer less drag, since not enough time may have elapsed to replenish the entire funnel cavity with seed photons. Indeed one expects the spikes to be more powerful the longer is the time interval between them, as more seed photons could pervade the cavity. This, besides causing internal shocks with the first shell which has been efficiently decelerated by Compton drag, will also result in a distribution of $`\mathrm{\Gamma }`$–factors: they will become greater on axis, where few seed photons can efficiently Compton drag the shells, and smaller towards the border of the funnel, where seed photons can be replenished by the funnel walls.
We plan to investigate these possibilities and their consequences on the associated predicted afterglows in future work.
## Acknowledgments
AC acknowledges the Italian MURST for financial support.
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# Dissipative dynamics of Bose condensates in optical cavities
## I Introduction
The experimental realization of Bose-Einstein condensation (BEC) in dilute atomic gases a couple of years ago as a consequence of improved cooling and trapping techniques has dramatically boosted the study of ultracold atoms. Today, BEC is a widespread tool and a huge range of new phenomena has been investigated experimentally and theoretically, see e.g. Refs. for recent overviews. In this context the interaction of BECs with laser light and optical lattices has been studied intensively and effects such as the reduction of the speed of light by many orders of magnitude and the occurrence of superradiance have been found. Recently the optical creation of vortices has been demonstrated and many more intriguing effects have been theoretically predicted.
In parallel, the field of cavity quantum electrodynamics, which studies the interaction of matter with one or a few single light modes, has reached such a level of sophistication that the interaction of light with the internal and external degrees of freedom of a single neutral particle can be observed and controlled in a very precise way . For optical fields trapping and cooling of a single atom in a cavity mode has been demonstrated . It is thus a logical next step to combine these two successful techniques and study the interaction of a BEC with a high finesse optical cavity, that is, the strong coupling of a far detuned optical mode to the dynamics of a condensate described by its macroscopically occupied wave function . Several groups have already been working along these lines and, for example, predicted the amplification of matter waves and the occurrence of dressed condensates . In the extreme limit one could envisage a large atomic cloud trapped and manipulated with a single photon.
In this work we extent our recently proposed scheme for cooling one or a few atoms in high-quality optical cavities to the case of a BEC. This requires a quantum mechanical treatment of the external degrees of freedom and the inclusion of atom-atom interactions. The system under investigation is a Bose-Einstein condensate interacting with the mode (the two modes) of a driven standing wave cavity (ring cavity). The properties of such a system as a measuring device for condensates have been discussed previously .
We will investigate the ground state and collective excitations of the coupled condensate-light field system in the optical potential of the cavity. Because of the strong coupling of the condensate wave function to the cavity modes, an oscillation of the condensate also leads to an oscillation of the intracavity light fields. Accordingly, the oscillation frequencies of the collective excitations are shifted with respect to an external optical potential of fixed depth, i.e., the optical potential formed by a free space standing wave. Furthermore, for appropriately chosen parameters the dissipative dynamics of the cavity due to cavity decay gives rise to damping of the condensate excitations without incoherent spontaneous emission. We analyze this damping mechanism and study its parameter dependence by numerical solutions of the coupled equations of motion as well as by analytic solutions of a simplified model based on the Lamb-Dicke expansion.
This paper is organized as follows. In Sec. II we discuss the case of a condensate interacting with the single mode of a standing wave cavity. After presenting the set of coupled non-linear equations of motion for the condensate and the light field, we discuss the numerical and analytical solutions for the ground state and the collective excitations. Section III investigates the more complicated situation of a condensate coupled to the two independent modes of an optical ring cavity. In Sec. IV we discuss the influence of binary collisions between the atoms within the condensate on the excitation frequencies and damping rates. Finally, we summarize our results in Sec. V.
## II BEC in standing wave cavity
Let us first consider the case of a Bose-Einstein condensate interacting with a single standing wave mode. The cavity mode is assumed for all times to be in a coherent state $`|\alpha (t)`$ and the condensate is described by a single wave function $`|\psi (t)`$ for all $`N`$ particles, which is a good approximation at zero temperature. This means that we factorize the quantum state of the system and thus neglect any entanglement between the condensate and the cavity field which might build up in the course of the time evolution. This simplification is only justified in the limit of a large photon number $`|\alpha |^2`$. To avoid spontaneous emission, we assume a very large detuning of the light field from the atomic resonance. More precisely, we assume that the cavity decay is the dominant incoherent process in the system and we thus require that $`\kappa N\mathrm{\Gamma }s`$, where $`\kappa `$ is the cavity decay rate, $`N`$ the number of atoms in the condensate, $`\mathrm{\Gamma }`$ the spontaneous decay rate of the atoms, and $`s`$ the atomic saturation parameter. After adiabatic elimination of the atomic excited states we obtain the following equations of motion,
$`{\displaystyle \frac{d}{dt}}\alpha (t)=\left[i\mathrm{\Delta }_ciNU(\widehat{x})\kappa \right]\alpha (t)+\eta ,`$ (2)
$`i{\displaystyle \frac{d}{dt}}\psi (x,t)=`$ (3)
$`\left\{{\displaystyle \frac{\widehat{p}^2}{2m}}+|\alpha (t)|^2U(x)+Ng_{coll}|\psi (x,t)|^2\right\}\psi (x,t).`$ (4)
Here $`\mathrm{\Delta }_c`$ is the detuning of the pump field from the cavity mode, $`U(x)=U_0\mathrm{cos}^2(kx)`$ is the optical potential formed by a single cavity photon, and $`\eta `$ describes the action of the driving laser. The expectation value $`U(\widehat{x})`$ has to be taken with respect to the momentary wave function $`|\psi (t)`$. Equation (4) is the well known Gross-Pitaevskii equation (GPE) for a condensate in an external field which in our case depends on the momentary field intensity $`|\alpha (t)|^2`$. The last term in the GPE models the interaction of atoms within the condensate where $`g_{coll}`$ is related to the s-wave scattering length $`a`$ by $`g_{coll}=4\pi \mathrm{}^2a/m`$.
Equations (II) are two coupled nonlinear equations describing the dynamics of the compound system formed by the condensate and the cavity field . The most interesting effects occur for parameters where the coupling between these equations significantly changes the system behaviour. We thus impose the condition $`NU_0\kappa `$, which guarantees that the presence of the condensate shifts the cavity frequency efficiently into or out of resonance with the driving field. At the same time the optical potential depth $`|\alpha |^2U_0`$ should be large enough to provide at least a few bound states for the atoms. Limitations and the interesting parameter regimes for this model have been discussed in Ref. .
### A Ground state
In order to obtain the ground state of the compound condensate-cavity system, we have to find the stationary solution of the system of coupled nonlinear equations (II). This can be done by elimination of $`|\alpha (t)|^2`$ in (4) using (2), and a subsequent numerical solution of the resulting non-linear equation for the ground state wave function with the method of steepest descent. This consists of a numerical propagation of the GPE in imaginary time $`\tau =it`$ until the wave function converges to a stationary state.
In this work we will concentrate on the case of $`U_0>0`$ where the potential minima coincide with the antinodes of the field (low-field seeking atoms). The ground state wave function will thus be localized at the field antinodes, thereby minimizing the coupling of the condensate to the light field. For a cavity resonant with the driving field, this means that the photon number is maximum for the stationary ground state. Any excitation of the condensate will then lead to a smaller cavity field.
As expected we find that the ground state wave function becomes better localized for stronger driving fields $`\eta `$ and larger optical potentials $`U_0`$. On the other hand, a strong atom-atom repulsion (large positive values of $`g_{coll}`$) increases the width of the BEC wave function and thus counteracts the confining effect of the potential. This in due course leads to an increased coupling of the BEC to the cavity field and hence a smaller field intensity. A more detailed analysis of the ground state wave function $`\psi _0(x)`$, its energy $`\mu `$, and the stationary field intensity $`|\alpha _0|^2`$ has been given in Ref. .
### B Collective excitations
Let us now turn to weak excitations of the condensate from the ground state. First, we will calculate the spectrum of collective excitations of the condensate. In contrast to fixed external fields, the trapping potential in the cavity depends on the BEC wave function. Hence, excitations include small deviations of the wave function and the cavity field $`\alpha `$ from their respective stationary state. We may thus write $`\psi (x,t)=\mathrm{exp}(i\mu t)[\psi _0(x)+\delta \psi (x,t)]`$ and $`\alpha (t)=\alpha _0+\delta \alpha (t)`$. For convenience we have already included the ground state time evolution into the ansatz for the wave function here. Inserting this into Eqs. (II) and linearizing in $`\delta \psi `$ and $`\delta \alpha `$ we obtain
$`i{\displaystyle \frac{d}{dt}}\delta \alpha `$ $`=`$ $`[\mathrm{\Delta }_c+N\psi _0|U(\widehat{x})|\psi _0i\kappa ]\delta \alpha `$ (7)
$`+N\alpha _0\delta \psi |U(\widehat{x})|\psi _0+N\alpha _0\psi _0|U(\widehat{x})|\delta \psi ,`$
$`i{\displaystyle \frac{d}{dt}}\delta \psi `$ $`=`$ $`\left\{{\displaystyle \frac{\widehat{p}^2}{2m}}+|\alpha _0|^2U+2Ng_{coll}|\psi _0|^2\mu \right\}\delta \psi `$ (9)
$`+Ng_{coll}\psi _0^2\delta \psi ^{}+\alpha _0U\psi _0\delta \alpha ^{}+\alpha _0^{}U\psi _0\delta \alpha .`$
For large $`\kappa `$ (more precisely, for $`1/\kappa `$ much smaller than the time scale of the condensate motion), the cavity field follows adiabatically the changes of the wave function and thus $`\delta \alpha `$ can be adiabatically eliminated. In this case one recovers the limit of Ref. .
In general the linearized time evolution couples the deviations $`\delta \psi `$ and $`\delta \alpha `$ also to their complex conjugates. In order to obtain excitation eigenstates, i.e., periodic solutions, we thus have to use the simultaneous ansatz
$`\delta \psi `$ $`=`$ $`e^{\gamma t}\left[e^{i\nu t}\delta \psi _+(x)+e^{i\nu t}\delta \psi _{}(x)^{}\right],`$ (10)
$`\delta \alpha `$ $`=`$ $`e^{\gamma t}\left[e^{i\nu t}\delta \alpha _++e^{i\nu t}\delta \alpha _{}^{}\right].`$ (11)
The collective excitations are thus defined as the solutions of the eigenvalue problem
$$\omega \left(\begin{array}{c}\delta \alpha _+\\ \delta \alpha _{}\\ \delta \psi _+(x)\\ \delta \psi _{}(x)\end{array}\right)=𝐌\left(\begin{array}{c}\delta \alpha _+\\ \delta \alpha _{}\\ \delta \psi _+(x)\\ \delta \psi _{}(x)\end{array}\right),$$
(12)
where $`𝐌`$ is easily obtained from Eqs. (II B) as a non-Hermitian matrix. The complex eigenvalues have the form $`\omega _n=\nu _ni\gamma _n`$, where $`\nu _n`$ is the oscillation frequency of the $`n`$th collective excitation and $`\gamma _n`$ the corresponding damping rate. Note that, depending on the parameters, negative damping rates are possible, leading to an exponential growth of the collective excitations. In this case the assumption of small deviations from the ground state imposed above only holds for very short times. Hence by changing some cavity parameters we can switch between stable and unstable cases and generate controlled excitations of the condensate and study their decay. In the following we will, however, concentrate on the case of positive $`\gamma _n`$ and therefore damped excitations.
Physically this damping arises from a kind of Sisyphus mechanism. For cavity damping rates $`\kappa `$ of the order of the oscillation frequencies $`\nu _n`$, the cavity field follows with a certain delay the changes of the condensate wave function. By properly choosing the system parameters, it can be achieved that on average the wave function has to climb up the potential hills at higher cavity field intensities and runs down at lower intensities. The condensate thus loses potential energy which is carried away by the cavity output field without an intrinsic decoherence of the condensate.
Furthermore it should be emphasized that the appearance of a damping rate in the linearized equations (II B) is a purely quantum feature related to the width of the atomic wave function. In the semiclassical limit of a point-like particle, the self-consistent ground state yields a particle exactly located at the antinodes of the cavity and hence all expectation values in Eq. (7) vanish. Thus the cavity field decouples from the atomic degrees of freedom and no damping of the atomic motion occurs to lowest order in the elongation $`x`$. This is in contrast to the case of a ring cavity as will be shown in the following section.
In Fig. 1 we show the oscillation frequencies and damping rates of the lowest collective excitations obtained numerically by calculating the eigenvalues of Eq. (12) on a spatial grid. The eigenvalues are plotted as a function of the cavity decay rate $`\kappa `$. Note that in order to keep the optical potential constant, we also have to scale the driving field $`\eta ^2`$ and the optical potential per photon $`U_0`$ proportional to $`\kappa `$.
We see that there exists one single eigenvalue $`\omega _f=\nu _fi\gamma _f`$ which scales approximately proportional to $`\kappa `$ in contrast to all of the other eigenvalues. This specific excitation of the system corresponds to an eigenmode where mainly the cavity field oscillates and the condensate wave function is only weakly perturbed. In fact, equation (7) shows that in the case where the atoms are well localized at the antinodes of the field (semiclassical limit), the cavity mode decouples from the matter wave function and the eigenvalue is given by $`\omega _f=\mathrm{\Delta }_ci\kappa `$.
Second, we notice that out of the other modes $`\nu _n`$, $`n=1,2,\mathrm{}`$, the ones with odd indices are independent of $`\kappa `$ and their damping rates vanish. This effect is due to the spatial symmetry of the problem considered here. For all parameters we find that the ground state wave function $`\psi _0`$ is symmetric in the position $`x`$. Thus for all antisymmetric excitations the expectation values in Eq. (7) vanish and the light field decouples. Therefore these odd ($`n=1,3,\mathrm{}`$) excitations are the same as for a trap of constant light intensity and hence there is no Sisyphus damping mechanism at work. Consequently only the lowest symmetric collective excitations are significantly altered by the interaction with the damped cavity mode. We will discuss the parameter dependence of the excitation $`n=2`$ by using an approximate analytic solution in the next section.
Let us finally emphasize that the oscillation frequency and the damping rate of the symmetric collective excitations can be monitored nondestructively via the cavity output intensity.
### C Harmonic oscillator approximation
In order to gain more insight into the parameter dependence of this damping mechanism, we will now analytically solve an approximate model of our system. To this end we expand the optical potential $`U(x)=U_0\mathrm{cos}^2(kx)`$ with $`U_0>0`$ up to second order around the antinodes of the field, i.e., we set $`U(x)=U_0(kx)^2`$ and assume $`g_{coll}=0`$, i.e., no atom-atom interaction. For simplicity we will also assume $`\mathrm{\Delta }_c=0`$.
The ground state of the Schrödinger equation (4) is thus the well known harmonic oscillator ground state which depends on the cavity field $`|\alpha |^2`$ in a parametric way. After inserting this wave function in the expectation value in Eq. (2) we obtain an equation for the self-consistent cavity field with the solution
$$|\alpha _0|^2=\frac{\eta ^2}{\kappa ^2}\frac{N^2U_0\omega _R}{4\kappa ^2}.$$
(13)
The corresponding harmonic oscillator frequency is then
$$\omega _0=2\omega _R\sqrt{|\alpha _0|^2U_0/\omega _R}$$
(14)
and the ground state energy is $`\mu =\omega _0/2`$.
For the collective excitations we now have to solve Eqs. (II B) with the harmonic potential. The last expectation value of Eq. (7) thus reads
$`\psi _0|U(\widehat{x})|\delta \psi `$ $`=`$ $`U_0\psi _0|(k\widehat{x})^2|\delta \psi `$ (15)
$`=`$ $`U_0{\displaystyle \frac{\omega _R}{\omega _0}}\psi _0|(aa^{})^2|\delta \psi ,`$ (16)
where we have used the standard relation between the position operator $`\widehat{x}`$ and the ladder operators $`a`$ and $`a^{}`$ of the harmonic oscillator. From this we see that the cavity field only couples to wave function deviations $`\delta \psi `$ containing the ground state $`\psi _0`$ and/or the second excited state $`\psi _2`$ of the harmonic oscillator. Most of the harmonic oscillator excited states are thus unperturbed and we find the collective excitations of the form $`\delta \psi _+(x)=\psi _n`$, $`\delta \alpha _+=\delta \alpha _{}=\delta \psi _{}(x)=0`$ with the positive eigenvalues $`\omega =n\omega _0`$ for $`n=1`$ and $`n3`$. Analogously there exist excitations with negative eigenvalues $`\omega =n\omega _0`$ of the form $`\delta \psi _{}(x)=\psi _n`$, $`\delta \alpha _+=\delta \alpha _{}=\delta \psi _+(x)=0`$. Hence, in addition to the antisymmetric states already found to decouple previously, also the higher lying symmetric states decouple in the harmonic approximation. Therefore these symmetric excitations are only damped due to the anharmonicity of the potential and due to atomic collisions in the full model.
The remaining (and most interesting) collective excitations are finally found by restricting the wave functions $`\delta \psi _\pm `$ in Eq. (12) to the 2-dimensional Hilbert space spanned by $`\psi _0`$ and $`\psi _2`$. The resulting 6x6 matrix has two zero eigenvalues and the other four eigenvalues have to be found by solving the 4th order polynomial equation
$`(i\kappa +NU_0{\displaystyle \frac{\omega _R}{\omega _0}}\omega )(i\kappa NU_0{\displaystyle \frac{\omega _R}{\omega _0}}\omega )(\omega ^24\omega _0^2)`$ (17)
$`4(NU_0\omega _R)^2=0.`$ (18)
This gives us the (complex) eigenvalues $`\omega _f`$ and $`\omega _2`$ and their counterparts of negative frequency. Although an analytic solution of (18) is possible in principle, the resulting expressions are rather long and do not provide much insight. Instead, we calculate the eigenvalue $`\omega _2`$ in the limit of large $`\kappa `$ as in Fig.1, i.e., by keeping $`U_0/\kappa `$ and $`\eta ^2/\kappa `$ constant. The zeroth order in this expansion in $`\omega /\kappa `$ yields the leading order of the frequency
$$\nu _2=2\omega _0\sqrt{1\frac{N^2U_0\omega _R}{4\eta ^2}}$$
(19)
and the first order gives the leading order of the decay rate
$$\gamma _2=\frac{4N^2U_0^2\omega _R^2}{\kappa ^3}\left(1\frac{N^2U_0\omega _R}{4\eta ^2}\right)^2.$$
(20)
Equation (19) gives a quantitative explanation for the frequency shift of $`\nu _2`$ according to the coupling of the BEC and the cavity mode as compared to the value $`2\omega _0`$ for the case of a harmonic oscillator potential of fixed photon number. We also see that the small variation of $`\nu _2`$ in Fig.1(a) for small values of $`\kappa `$ are in fact of the order $`1/\kappa ^2`$. Equation (20) leads to the asymptotic behaviour like $`1/\kappa `$ for the decay rate $`\gamma _2`$ in Fig.1(b). In the limit of a strong driving field the frequency $`\nu _2`$ of the second collective excitation approaches the harmonic oscillator value. Simultaneously the damping rate $`\gamma _2`$ tends towards a constant non-vanishing value which is proportional to the square of the atom number $`N`$. A higher condensate density thus significantly increases the damping of the collective excitation.
## III BEC in a ring cavity
In this section we will now discuss the case of a BEC in a ring cavity. In this case the condensate is coupled to the two independent travelling wave modes $`\alpha _\pm `$. For simplicity we will assume in the following that both modes of the cavity are driven with the same pumping rate $`\eta `$ . Therefore the equations of motion read
$`{\displaystyle \frac{d}{dt}}\alpha _\pm (t)`$ $`=`$ $`\left[i\mathrm{\Delta }_ciNU_0\kappa \right]\alpha _\pm (t)`$ (23)
$`iNU_0e^{2ik\widehat{x}}\alpha _{}+\eta ,`$
$`i{\displaystyle \frac{d}{dt}}\psi (x,t)`$ $`=`$ $`\{{\displaystyle \frac{\widehat{p}^2}{2m}}+Ng_{coll}|\psi (x,t)|^2`$ (26)
$`+U_0|\alpha _+(t)e^{ikx}+\alpha _{}(t)e^{ikx}|^2\}`$
$`\times \psi (x,t).`$
Hence in general the condensate will scatter light between the left and right running waves and induce a strong coupling. This gives additional degrees of freedom to the system compared to the standing wave case. For example, in addition to intensity shifts the condensate can also induce a relative phase shift between the two modes, which changes the position of the potential wells. Analogously, the minima can also be controlled externally by the relative phase of the two driving fields which allows to selectively excite antisymmetric excitations.
Considering the important role which the spatial symmetry plays for the standing wave cavity, we will now change to the description of the cavity modes by $`\alpha _s=\alpha _++\alpha _{}`$ and $`\alpha _a=\alpha _+\alpha _{}`$ which have symmetric and antisymmetric mode functions, respectively. In this new basis Eqs. (III) read
$`{\displaystyle \frac{d}{dt}}\alpha _s(t)`$ $`=`$ $`\left[i\mathrm{\Delta }_ciNU_0iNU_0\mathrm{cos}(2k\widehat{x})\kappa \right]\alpha _s(t)`$ (29)
$`+NU_0\mathrm{sin}(2k\widehat{x})\alpha _a(t)+2\eta ,`$
$`{\displaystyle \frac{d}{dt}}\alpha _a(t)`$ $`=`$ $`\left[i\mathrm{\Delta }_ciNU_0+iNU_0\mathrm{cos}(2k\widehat{x})\kappa \right]\alpha _a(t)`$ (31)
$`NU_0\mathrm{sin}(2k\widehat{x})\alpha _s(t),`$
$`i{\displaystyle \frac{d}{dt}}\psi (x,t)`$ $`=`$ $`\{{\displaystyle \frac{\widehat{p}^2}{2m}}+Ng_{coll}|\psi (x,t)|^2`$ (34)
$`+U_0|\alpha _s(t)\mathrm{cos}(kx)+i\alpha _a(t)\mathrm{sin}(kx)|^2\}`$
$`\times \psi (x,t).`$
Note that because of the assumption of a single pumping rate $`\eta `$ for $`\alpha _+`$ and $`\alpha _{}`$, in the new basis only the symmetric mode $`\alpha _s`$ is pumped. The antisymmetric mode $`\alpha _a`$ only contains photons which have been scattered by the condensate out of $`\alpha _s`$.
### A Ground state
For the calculation of the ground state of the compound system formed by the BEC and the cavity modes we will again assume the case $`U_0>0`$. We then find that the ground state wave function is localized at the antinodes of the driven mode $`\alpha _s`$ and is symmetric in $`x`$. Thus the expectation values of $`\mathrm{sin}(2kx)`$ in Eqs. (III) vanish and Eq. (31) decouples. The stationary state of the antisymmetric mode is therefore given by $`\alpha _{a,0}=0`$. Equations (29) and (34) then reduce to the equations (II) for the standing wave cavity if one identifies the parameters
$`2U_0^\text{r}`$ $`=`$ $`U_0^\text{s},`$ (35)
$`|\alpha _s^\text{r}|^2/2`$ $`=`$ $`|\alpha ^\text{s}|^2,`$ (36)
$`\sqrt{2}\eta ^\text{r}`$ $`=`$ $`\eta ^\text{s},`$ (37)
for the ring cavity and the standing wave cavity, respectively. The ground state of the system can thus be obtained by using our previous results for the standing wave cavity and all of the discussions there equally apply to the ground state in the ring cavity.
### B Collective excitations
The collective excitations are calculated with the same method as in the preceding section by linearization of the equations of motion (III) in small deviations of $`|\psi `$, $`\alpha _s`$ and $`\alpha _a`$ from their stationary states $`|\psi _0`$, $`\alpha _{s,0}`$ and $`0`$. Choosing the ground state wave function to be real and taking its symmetry into account we obtain
$`i{\displaystyle \frac{d}{dt}}\delta \alpha _s`$ $`=`$ $`[\mathrm{\Delta }_c+2NU_0\psi _0|\mathrm{cos}^2(k\widehat{x})|\psi _0i\kappa ]\delta \alpha _s`$ (40)
$`+2NU_0\alpha _{s,0}[\delta \psi |\mathrm{cos}^2(k\widehat{x})|\psi _0+c.c.],`$
$`i{\displaystyle \frac{d}{dt}}\delta \alpha _a`$ $`=`$ $`[\mathrm{\Delta }_c+2NU_0\psi _0|\mathrm{sin}^2(k\widehat{x})|\psi _0i\kappa ]\delta \alpha _a`$ (42)
$`NU_0\alpha _{s,0}[\delta \psi |\mathrm{sin}(2k\widehat{x})|\psi _0+c.c.],`$
$`i{\displaystyle \frac{d}{dt}}\delta \psi `$ $`=`$ $`\{{\displaystyle \frac{\widehat{p}^2}{2m}}+|\alpha _{s,0}|^2U_0\mathrm{cos}^2(k\widehat{x})+2Ng_{coll}|\psi _0|^2`$ (46)
$`\mu \}\delta \psi +Ng_{coll}\psi _0^2\delta \psi ^{}`$
$`+U_0\mathrm{cos}^2(k\widehat{x})\psi _0(\alpha _{s,0}\delta \alpha _s^{}+c.c.)`$
$`{\displaystyle \frac{i}{2}}U_0\mathrm{sin}(2k\widehat{x})\psi _0(\alpha _{s,0}\delta \alpha _a^{}c.c.).`$
From these equations we see that the behaviour of the excitation eigenstates strongly depends on their spatial symmetry.
For symmetric excitations $`\delta \psi (x)`$ the last expectation value in (42) vanishes and the antisymmetric cavity mode decouples from the wave function. Hence in this case we find $`\delta \alpha _a=0`$. The equations of motion for $`\delta \psi `$ and $`\delta \alpha _s`$ then reduce to their standing wave counterpart discussed in the previous section if one rescales the parameters as in (37). The symmetric collective excitations are thus the same as those in a standing wave cavity.
Analogously, for antisymmetric excitations $`\delta \psi (x)`$ the symmetric cavity mode decouples and therefore $`\delta \alpha _s=0`$. We then find a new set of coupled equations for $`\delta \psi `$ and $`\delta \alpha _a`$. Thus, in contrast to the case of a standing wave cavity, also the antisymmetric excitations are damped in a ring cavity. However, the damping mechanism is of completely different physical origin. Instead of the Sisyphus mechanism discussed above, here the coherent scattering of photons from the $`\alpha _s`$ cavity mode into the $`\alpha _a`$ mode is responsible for the damping. This leads to less severe requirements for the cavity parameters as we will see in the following subsection.
Figure 2 shows the spectrum of collective excitations of a Bose condensate in a ring cavity which is obtained from the numerical solution of Eqs. (III B). First, we note that in contrast to the case of the standing wave cavity we now find two modes with eigenvalues which scale proportional to the cavity decay rate $`\kappa `$. In the semiclassical limit (atoms well localized), these correspond to pure oscillations of the symmetric and antisymmetric field mode, respectively, and are thus labelled $`\omega _s=\nu _si\gamma _s`$ and $`\omega _a=\nu _ai\gamma _a`$. The semiclassical limits of these eigenfrequencies are obtained from Eqs. (III B) as $`\omega _s=\mathrm{\Delta }_ci\kappa `$ and $`\omega _a=\mathrm{\Delta }_c+2NU_0i\kappa `$. Although the damping rates of these modes are equal, we see that the different spatially depending coupling to the atoms leads to a large difference in the oscillation frequencies.
For the parameters chosen in Fig. 2 the other oscillation frequencies $`\nu _n`$, $`n1`$, are mainly independent of $`\kappa `$. However, whereas all frequencies with $`n2`$ are equally spaced and hence very well described by harmonic oscillator states, the lowest frequency $`\nu _1`$ is significantly shifted downwards. For the damping rates we find that only the two lowest excitations exhibit relevant damping. However, the dependence of these damping rates on the system parameters is qualitatively different according to the different damping mechanisms. We will return to the discussion of these features in the following subsection where we derive analytic approximations for the eigenvalues.
### C Harmonic oscillator approximation
Let us now calculate analytic estimates for the lowest oscillation frequencies and damping rates along the lines of Sec. II C. We will thus again assume $`\mathrm{\Delta }_c=g_{coll}=0`$.
As we have already seen, the calculation of the ground state and of the symmetric collective excitations can be reduced to the problem of the standing wave if the appropriate identification of the system parameters (37) is made. We can therefore use our previous results to obtain the self-consistent cavity field
$$|\alpha _{s,0}|^2=\frac{4\eta ^2}{\kappa ^2}\frac{N^2U_0\omega _R}{\kappa ^2}$$
(47)
and the corresponding harmonic oscillator frequency
$$\omega _0=2\omega _R\sqrt{|\alpha _{s,0}|^2U_0/\omega _R}.$$
(48)
For the lowest symmetric excitation, the expansion for large values of $`\kappa `$ yields
$$\nu _2=2\omega _0\sqrt{1\frac{N^2U_0\omega _R}{4\eta ^2}}$$
(49)
and
$$\gamma _2=\frac{16N^2U_0^2\omega _R^2}{\kappa ^3}\left(1\frac{N^2U_0\omega _R}{4\eta ^2}\right)^2.$$
(50)
Analogously, we can calculate the lowest antisymmetric excitation by expanding the expectation values in Eqs. (42) and (46) to lowest order in $`k\widehat{x}`$. To this order only $`\delta \psi `$ and $`\delta \psi ^{}`$ proportional to the first harmonic oscillator wave function $`\psi _1`$ couple to the cavity field $`\delta \alpha _a`$ and $`\delta \alpha _a^{}`$ and we thus have to find the eigenvalues of a 4x4 matrix, that is, we must solve the characteristic polynomial
$`[i\kappa +2NU_0(1{\displaystyle \frac{\omega _R}{\omega _0}})\omega ][i\kappa 2NU_0(1{\displaystyle \frac{\omega _R}{\omega _0}})\omega ]`$ (51)
$`\times (\omega ^2\omega _0^2)4(NU_0\omega _0)^2(1{\displaystyle \frac{\omega _R}{\omega _0}})=0.`$ (52)
In the limit of $`\kappa \mathrm{}`$ (with constant $`U_0/\kappa `$ and $`\eta ^2/\kappa `$) this yields the oscillation frequency
$$\nu _1=\omega _0\sqrt{1\frac{4N^2U_0^2(1\frac{\omega _R}{\omega _0})}{\kappa ^2+4N^2U_0^2(1\frac{\omega _R}{\omega _0})^2}}.$$
(53)
The first order correction in $`1/\kappa `$ gives the dominant term of the corresponding damping rate
$$\gamma _1=\omega _0^2\kappa \frac{4N^2U_0^2(1\frac{\omega _R}{\omega _0})}{[\kappa ^2+4N^2U_0^2(1\frac{\omega _R}{\omega _0})^2]^2}.$$
(54)
We can now compare the behaviour of the two lowest eigenvalues as a function of the system parameters. As an example, let us consider the case of a relatively strong pump, $`\eta ^2N^2U_0\omega _R`$. In this limit, the second excitation frequency $`\nu _2`$ is only weakly shifted from the harmonic oscillator frequency $`2\omega _0`$. On the other hand, the frequency shift of the lowest excitation $`\nu _1`$ mainly depends on the ratio $`NU_0/\kappa `$. Since this ratio has to be larger than one in order to yield a significant frequency shift of the cavity by the atoms, equation (53) implies that $`\nu _1`$ is strongly shifted towards zero. Simultaneously we find for the damping rates that $`\gamma _2`$ becomes independent of $`\eta `$ in this limit, in contrast to $`\gamma _1`$ which is proportional to $`\omega _0`$ and thus proportional to $`\eta ^2`$. Therefore the damping rate of the first antisymmetric excitation can be increased arbitrarily by increasing the intensity of the pump field. The damping rate of the first symmetric excitation is much harder to manipulate because it is mainly governed by the optical potential per photon and thus by the quality of the cavity. On the other hand, we note that $`\gamma _2`$ scales proportional to $`N^2`$ whereas $`\gamma _1`$ is inversely proportional to $`N^2`$. The number of atoms thus provides another handle to change the relative size of the damping rates $`\gamma _1`$ and $`\gamma _2`$.
Another point is worth a comment here. We emphasized in the previous section that the damping mechanism for the collective excitations in a standing wave cavity is crucially related to the width of the matter wave function and vanishes in the semiclassical limit where the atoms are treated as point particles. In contrast to this we find that in the travelling wave cavity the damping mechanism still exists in the semiclassical limit. In fact, our results for the oscillation frequency (53) and the damping rate (54) agree with the semiclassical results if one takes formally the limit $`\omega _R/\omega _00`$.
In Fig. 3 we show the excitation frequencies $`\nu _{1,2}`$ and the damping rates $`\gamma _{1,2}`$ as a function of the pump strength $`\eta ^2`$ for both the numerical solution and the analytic approximations. We see that for the chosen parameters the approximations fit quite well apart from the values of $`\gamma _2`$. This comes from the fact that we obtained the complex eigenvalues $`\omega _n`$ from an expansion of Eqs. (18) and (52) for small values of $`|\omega _n|/\kappa 1`$. As we see from Fig. 3(a) this is well fulfilled for $`\omega _1`$ for the chosen parameters, but $`|\omega _2|/\kappa `$ is of the order of one. However, in the limit of a strong pump the lowest order term for the frequency $`\nu _2`$ already gives the correct value, namely twice the harmonic oscillator frequency. Thus, only the imaginary part (the damping rate $`\gamma _2`$) of the analytic approximation deviates from the exact solution in Fig. 3. In parameter regions where $`|\omega _2|/\kappa 1`$ we find a much better agreement of the two solutions.
## IV Interacting Bose gas
In the discussion so far we have omitted the effects of atomic interactions, as described by the collision rate $`g_{coll}`$ in the GPE, on the energies and damping rates of the collective excitations. Neglecting this has allowed us to obtain analytical expressions and therefore to discuss the parameter dependence of our results explicitly. However, atomic collisions are known to play a crucial role in experimental realizations of Bose-Einstein condensates. We will now discuss the changes of the collective excitations according to collisions in a numerical example of a condensate in a ring cavity.
We show in Fig. 4 the excitation frequencies $`\nu _n`$ and the corresponding damping rates $`\gamma _n`$ as a function of the collision rate $`g_{coll}`$ with all other parameters fixed. The main effect on the stationary ground state wave function of a repulsive interaction between the condensed atoms is to increase the width of the wave function. Since this larger width also changes the coupling to the cavity field, we find that the steady state photon number decreases with increasing collision rate. Consequently, the optical potential becomes more shallow and the excitation frequencies decrease. However, as we can see from Fig. 4 this argument thus not hold for the lowest (antisymmetric) excitation. Here the atomic collisions counteract the strong frequency shift which we found in the previous section and $`\nu _1`$ slightly increases with $`g_{coll}`$. Above a certain threshold value for $`g_{coll}`$ the atom-atom repulsion gets stronger than the confining effect of the optical potential. In this case the ground state wave function is no longer localized and the spectrum of excitations changes into that of unbound particles where each excitation frequency is doubly degenerate.
Figure 4(b) shows that collisional effects have an even more important influence on the damping rates of the collective excitations. We see that the effect differs for the damping rates $`\gamma _1`$ and $`\gamma _2`$. While $`\gamma _1`$ weakly decreases with increasing $`g_{coll}`$, $`\gamma _2`$ increases significantly over a broad range of values of $`g_{coll}`$. This is related to the fact that the damping of the symmetric excitation $`\gamma _2`$ depends crucially on the width of the ground state wave function whereas the damping mechanism of the antisymmetric excitation does not, as already emphasized before. Since the major effect of atom collisions is to broaden the wave function, the resulting changes of the damping rates occur predominantly for the symmetric excitations. Note also that the atomic collisions and the stronger anharmonicity of the potential according to the lower field intensity enhance the damping of higher collective excitations, as can be seen from the damping rates $`\gamma _3`$ and $`\gamma _4`$ in Fig. 4(b).
## V Conclusions
In summary we have studied in detail the interaction of a Bose-Einstein condensate with one or two single modes in a high-finesse optical cavity. We have solved the coupled set of non-linear equations of motion for the joint dynamics of the condensate and the light field numerically and compared it analytically with a simplified model based on the Lamb-Dicke expansion. We find that, even without atom-atom interaction, the oscillation frequencies are shifted with respect to their values in a fixed external potential.
For a finite cavity response time the collective excitations are damped or amplified depending on the cavity detuning, which can be easily controlled externally. We identify two distinct mechanisms depending on the spatial symmetry of the excitations. The damping mechanism in standing wave cavities and for spatially symmetric excitations in ring cavities is due to a Sisyphus type effect, which leads to larger cavity fields at times when the condensate runs up potential hills than at times when the condensate runs down. On average this effect extracts kinetic energy from the condensate which is carried away by the cavity field. On the other hand, the damping mechanism for the spatially antisymmetric excitations is only present in a ring cavity due to the scattering of cavity photons between the two counterpropagating waves. This creates an intensity imbalance, which is counteracted by the cavity damping and hence leads to momentum dissipation.
The two damping mechanisms exhibit very distinct parameter dependences. Our analytical approximations show that in the limit of strong cavity pumping the damping rate of the spatially symmetric excitations becomes independent of the pump but scales proportional to the square of the atom number $`N^2`$ while the damping rate of the antisymmetric excitations is proportional to the pump field intensity and inversely proportional to $`N^2`$, which implies less stringent requirements to cavity technology.
The difference between the damping/amplification rates of excitations with different spatial symmetry could be used to manipulate a Bose-Einstein condensate in a controlled fashion. In addition, in a ring cavity setup we can also excite oscillations by external phase and amplitude shifts of the pump light. All the effects could of course be enhanced by tailored feedback of the measured transmitted intensity onto the pump. This might give rise to useful applications of such a system in the context of quantum information and quantum computation in analogy to other recently proposed systems making use of particles in optical lattices .
## ACKNOWLEDGMENTS
We thank J. I. Cirac and P. Zoller for stimulating discussions. This work was supported by the Austrian Science Foundation FWF (Project P13435).
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# The Progenitor Masses of Wolf-Rayet Stars and Luminous Blue Variables Determined from Cluster Turn-offs. I. Results from 19 OB Associations in the Magellanic Clouds
## 1 Introduction
Conti (1976) first proposed that Wolf-Rayet (WR) stars might be a normal, late stage in the evolution of massive stars. In the modern version of the “Conti scenario” (Maeder & Conti 1994), strong stellar winds gradually strip off the H-rich outer layers of the most massive stars during the course of their main-sequence lifetimes. At first the H-burning CNO products He and N are revealed, and the star is called a WN-type WR star; this stage occurs either near the end of core-H burning or after core-He burning has begun, depending upon the luminosity of the star and the initial metallicity. Further mass-loss during the He-burning phases exposes the triple-$`\alpha `$ products C and O, and results in a WC-type WR star. Since the fraction of mass that a star loses during its main-sequence evolution depends upon luminosity (mass), we would expect that at somewhat lower masses evolution proceeds only as far as the WN stage. At still lower masses a star never loses sufficient mass to become a Wolf-Rayet at all, but spends its He-burning life as a red supergiant (RSG). Mass-loss rates also scale with metallicity as the stellar winds are driven by radiation pressure acting through highly ionized metal lines. Thus the mass-limits for becoming WN or WC stars should vary from galaxy to galaxy, and with location within a galaxy that has metallicity variations.
Studies of mixed-age populations in the galaxies of the Local Group have confirmed some of the predictions of the Conti scenario. For instance, the number ratio of WC and WN stars is a strong function of metallicity (Massey & Johnson 1998 and references therein), with proportionally more WC stars seen at higher metallicities, suggesting that the mass-limit for becoming WC stars is somewhat lower in these galaxies. Similarly the relative number of WRs and RSGs is correlated with metallicity, and there is a paucity of high luminosity RSGs at high metallicities (Massey 1998a), suggesting that these high luminosity stars have become WRs rather than RSGs.
However, fundamental questions remain concerning the evolution of massive stars:
(1) What is the role of the luminous blue variables (LBVs)? These stars are highly luminous objects that undergo photometric “outbursts” associated with increased mass-loss (Humphreys & Davidison 1994). Are LBVs a short but important stage in the lives of all high mass stars that occur at or near the end of core-H burning? Recent efforts have linked some of the LBVs to binaries, as Kenyon & Gallagher (1985) first suggested. The archetype LBV, $`\eta `$ Car, may be a binary with a highly eccentric orbit (Damineli, Conti, & Lopes 1997), but whether its outbursts have anything to do with the binary nature remains controversial (Davidson 1997), as does the orbit itself (Davidson et al. 2000). Similarly, the WR star HD 5980 in the SMC underwent an “LBV-like” outburst (Barba et al 1995); this star is also believed to be a binary with an eccentric orbit, although the nature (and multiplicity?) of the companion(s) remains unclear (Koenigsberger et al. 1998; Moffat 1999).
The Ofpe/WN9 type WRs, and the high-luminosity B\[e\] stars have recently been implicated in the LBV phenomenon. The former have spectral properties intermediate between “Of” and “WN” (Bohannan & Walborn 1989). One of the prototypes of this class, R 127, underwent an LBV outburst in 1982 (Walborn 1982; Stahl et al. 1983; see discussion in Bohannan 1997). Similarly some B\[e\] stars have been described as having LBV-like outbursts. Var C, a well-known LBV in M 33, has a spectrum indistinguishable from B\[e\] stars: compare Fig. 8a of Massey et al. (1996) with Fig. 8 of Zickgraf et al. (1986). Do all B\[e\] stars undergo an LBV phase or not? Conti (1997) has provided an insightful review.
(2) What is the evolutionary connection between WN and WC stars? We expect only the highest mass stars become WCs, while stars of a wider range in mass become WNs. The changing proportion of WCs and WNs within the galaxies of the Local Group have been attributed to the expected dependence of these mass ranges on metallicity. However, the relative time spent in the WN and WC stages may also change with metallicity, complicating the interpretation of such global measures drawn from mixed-age populations.
(3) Is there any evolutionary significance to the excitation subtypes? Both WN and WC stars are subdivided into numerical classes, or more coarsely into “early” (WNE, WCE) or “late” (WNL, WCL) based upon whether higher or lower excitation ions dominate. Recent modeling by Crowther (2000) suggests that the distinction between WNL and WNE is not actually due to temperature differences but primarily metal abundance. Armandroff & Massey (1991) and Massey & Johnson (1998) have argued that this true for the WC excitation classes based upon the metallicity of the regions where these stars are found.
If we knew the progenitor masses of LBVs and the various kinds of WRs we would have our answers to the above. However, here recourse to stellar evolution models fails us. Stellar evolutionary models show that a star’s path in the HRD during core-He burning is strongly dependent upon the amount of mass-loss that has preceded this stage. Thus the nature of the LBV phenomenon becomes very important in understanding where WRs come from, as the amount of mass ejected by LBVs is large, but given the episodic nature of LBVs, hard to include in the evolutionary models. In addition, the locations of WRs and LBVs in the H-R diagram are highly uncertain. LBVs have pronounced UV-excesses and “pseudo-photospheres” (Humphreys & Davidson 1994). For WR stars, neither the effective temperatures nor bolometric corrections are established, as none of the standard assumptions of stellar atmospheres hold in the non-LTE, rapidly expanding, “clumpy” stellar winds where both the stellar continua and emission-lines arise (e.g., Conti 1988). While the WR subtypes represent some sort of excitation sequence in the stellar winds, the relationship, if any, to the effective temperature of the star remains unclear.
There has been recent success in modeling WR atmospheres, with convincing matches to the observed line profiles and stellar continua from the UV to the near-IR. These models have the potential for determining the bolometric luminosities and effective temperatures. The “standard WR model” (Hillier 1987, 1990) assumes a spherical geometry and homogeneity, and then iteratively solves the equations for statistical equilibrium and radiative equilibrium for an adopted velocity law, mass-loss rate, and chemical composition. (See also Hillier & Miller 1998, 1999.) Comparison with observations then permits tweaking of the parameters. Although the solutions may not be unique, good agreement is often achieved with observations, and in a series of papers, Crowther and collaborators have offered the “fundamental” parameters (effective temperatures, luminosities, chemical abundances, mass-loss rates, etc.) of WN stars obtained with this model (Crowther, Hillier & Smith 1995a, 1995b; Crowther, Smith, & Hillier 1995c; Crowther et al. 1995d; Crowther, Smith, & Willis 1995e; Crowther & Smith 1997; Bohannan & Crowther 1999).
Here we utilize a complementary, observational approach to the problem, one that can not only answer the question of the progenitor masses of LBVs and WRs, but also provide data on the BCs that can help constrain and evaluate the WR atmosphere models.
### 1.1 The Use of Cluster Turn-offs
A time-honored method of understanding the nature of evolved stars is to determine the turn-off luminosities in clusters containing such objects (Johnson & Sandage 1955; Schwarzschild 1958). This was first applied by Sandage (1953) to determine the masses of RR Lyrae stars in the globular clusters M 3 and M 92, with a result that was at variance with that given by theory (Sandage 1956). Similarly, the turn-off masses of intermediate-age open clusters were used by Anthony-Twarog (1982) to determine the progenitor masses of white dwarfs. However, it is one thing to apply this to clusters with ages of $`10^{10}`$ yr, as was done for the RR Lyrae stars, or to clusters whose ages are $`2\times 10^7`$$`7\times 10^8`$ yr, as was done for white dwarfs. Can we safely extend this to clusters whose ages are only of order 3–5$`\times 10^6`$ yr in order to determine the progenitor masses of WRs and LBVs?
When stars form in a cluster or association, stars of intermediate mass appear to form over a significant time span—perhaps over several million years (Hillenbrand et al. 1993; Massey & Hunter 1998). However, modern spectroscopic and photometric studies have shown that the massive stars tend to form in a highly coeval fashion. For instance, in their study of the stellar content of NGC 6611, Hillenbrand et al. (1993) found a maximum age spread of 1 Myr for the massive stars, and noted that the data were consistent with no discernible age spread. for all one could tell “the highest-mass stars could have all been born on a particular Tuesday.” Similarly, the high mass stars in the R136 cluster have clearly formed over $`\mathrm{\Delta }\tau <1`$ Myr, given the large number of O3 V stars and the short duration that stars would have in this phase (Massey & Hunter 1998).
Such short time scales for star formation are consistent with recent studies by Elmegreen (1997, 2000a, 2000b), who argues that star formation takes place not over tens of crossing times but over one or two. For regions with large spatial extent (such as 100 pc diameter OB associations) star formation in the general region may occur over a prolonged time ($``$10 Myr). However, large OB associations can contain subgroups that have formed independently (Blaauw 1964), and are small enough so that a high degree of coevality ($`<12`$ Myr) is expected. The stars from such a subgroup need not be spatially coincident. Rather, a star with a random motion of 10 km s<sup>-1</sup> will have traveled 30 pc in just 3 Myr. Thus in an OB association we may find intermediate-mass stars which have formed from a number of subgroups over time, but massive stars which may have formed from a single subgroup and hence are coeval—even though these massive stars may now be spread out throughout the OB association. Or, it may be that massive stars of different ages are present, in which case the “turn-off mass” will not be relevant to the evolved object. We take an optimistic approach in our search for turn-off masses, but will insist that coevality be established empirically for the massive stars in the region in question.
For massive stars, the mass-luminosity relationship is much flatter than for solar-type stars ($`LM^{2.4}`$ for 30 $`_{}`$ and $`LM^{1.5}`$ for 120 $`_{}`$). As a result, the lifetimes of massive stars do not change as drastically with mass as one might expect. A 120 $`_{}`$ will have a main-sequence lifetime of 2.6 Myr, a 60 $`_{}`$ still will have a main-sequence lifetime of 3.5 Myr, and a 25 $`_{}`$ star will have a main-sequence lifetime of 6.4 Myr. (These numbers are based on the $`z=0.02`$ models of Schaller et al. 1992.)
Thus it should be possible to use clusters and OB associations to pin down the “minimum mass” of various unevolved massive stars. If the highest mass star still on the main-sequence is 60$`_{}`$, and its associated stellar aggregate contains a WC-type WR star, then we might reasonably conclude that the progenitor mass of the WC star was at least 60 $`_{}`$. Of course, if coevality does not hold, then this answer may be wrong—the WC star might have come from a 25 $`_{}`$ that formed earlier. But were that the case, it would have to have formed much earlier—at least 3 Myr earlier, according to the lifetimes given above, and such an age spread should be readily apparent.
We can in principle also find the BCs from the cluster turn-offs. It is straightforward to determine the absolute visual magnitude of the WR, making some modest correction for the emission lines. Since massive stars evolve at nearly constant bolometric luminosity, we expect that the bolometric luminosity of the WR will be at least as great as the bolometric luminosity of the highest mass main-sequence object. With modern stellar models we can improve on this by making first-order correction for modest luminosity evolution.
We are, of course, not the first to have trod on this ground. Schild & Maeder (1984) attempted to provide links between the different WR subtypes using this sort of analysis of Galactic clusters, concluding that stars with masses as low as $`18_{}`$ became WN stars, while WC stars came from stars of $`35_{}`$ and higher, and proposing various evolutionary relationships between the various subtypes. Humphreys, Nichols, & Massey (1985) also used data drawn from the literature on (mostly the same) Galactic clusters, and found a considerably higher minimum mass for becoming a WR star (30 $`_{}`$), with no difference between the masses required to become a WN or a WC. They were also the first to apply this method to determining the minimum bolometric corrections for WR stars, concluding that WNE stars have BCs $`<5.5`$ mag, WNL stars have BCs $`<3.5`$ mag, and WCs have BCs $`<5.0`$ mag. (These BCs are considerably more negative than had been commonly assumed.) Smith, Meynet, & Mermilliod (1994) re-addressed the issue of BCs by analyzing the same data from the literature on what was also mostly the same clusters, finding BCs for WNs that were typically $`4`$ mag (WNL) to $`6`$ mag (WNE), and $`4.5`$ for WCs, essentially unchanged from the Humphreys et al. findings.
There were problems, however, with these earlier studies. The most significant one was the reliance upon (the same) literature data for the spectral types of the main-sequence stars in these clusters and associations. Over the past decade we have examined the stellar content of numerous clusters and OB associations in the Milky Way, and invariably discovered stars of high mass that had been previously missed either due to reddening or simple oversight (Massey, Johnson, & DeGioia-Eastwood 1995a). A related problem is that some of the literature spectral types were “outdated” for the O-type stars, particularly for stars of type O7 and earlier, which would lead to an incorrect assignment of bolometric corrections and hence luminosities and masses. In addition, our understanding of massive star evolution has improved to the point where we can do a considerably better job in assigning masses, and in particular understand the errors associated with this procedure (see, for example, Massey 1998b). Another problem was that the spectral information was sufficiently sparse that no test of coevality could be applied to the cluster. In addition, poor photometry—often photographic—led to poor reddening corrections. And, finally, a significant limitation in these earlier studies was that all were restricted to the Milky Way. It would be most interesting to understand the origin of evolved massive stars as a function of metallicity; for this, extension to the Magellanic Clouds is a logical step.
We have attempted to rectify these problems by carrying out a modern analysis of OB associations containing WR and other evolved massive stars in galaxies of the Local Group, obtaining new spectroscopic and photometric data where warranted, and combining this with studies drawn form the recent literature. In this first paper we will determine the progenitor masses of WR and LBVs in 19 associations of the Magellanic Clouds. These two galaxies have abundances which are low compared to the solar neighborhood. In the next paper we will compare these to new results obtained for OB associations in our own Galaxy. In a third paper we will combine HST photometric and spectroscopic data with large-aperture ground-based studies to extend this work to the more distant members of the Local Group as an addition check on metallicity effects.
Throughout this paper we will assume the true distance modulus of the SMC is 18.9, and that of LMC is 18.5 (Westerlund 1997; van den Bergh 2000).
## 2 Sample Selection and Observing Strategy
In selecting this sample, we first compared the locations of known WRs and LBVs to that of the cataloged OB associations in the SMC and LMC. The probability of a chance supposition of a rare evolved object against one of these associations is, of course, low.
There are nine known WR stars in the SMC (Azzopardi & Breysacher 1979; Morgan, Vassiliadis, & Dopita 1991). Four of these are within three of the OB associations identified by Hodge (1985). We list these in Table 1. The WR star HD 5980 underwent an “LBV-like outburst” in 1994 (Barba et al. 1995). This star is located in NGC 346, which is included in our study. Three other SMC stars described as LBV-like in some way are R 40, which is not a member of any association: R 4, a B\[e\] star with “brightness variations typical for LBVs” (Zickgraf et al. 1996), located in Hodge 12, but not included here, and AV 154 (aka S 18), another B\[e\] star tied to LBVs (Morris et al. 1996), located just outside of Hodge 35, also not included here. One other high luminosity B\[e\] star, R 50 (aka S 65=Sk 193), is listed by Zickgraf et al. (1986), but is well outside any OB association.
For the LMC, Breysacher (1981) cataloged 100 Wolf-Rayet stars; an occasional additional one has been found spectroscopically (e.g., Conti & Garmany 1983; Testor, Schild, & Lortet 1993), plus components of R 136 and other crowded clusters have been successfully resolved, which brought the total of known WR stars in the LMC to 134 (Breysacher, Azzopardi, & Testor 1999). As part of the present study, we discovered a new WR star, Sk$`69^{}`$ 194, located in LH 81. We compared the positions of WRs against the Lucke-Hodge OB associations (Lucke & Hodge 1970; Lucke 1972), using only those associations with “A1” classifications. Not all were included in the current study; we list in Table 1 the 16 associations that are, along with their WR stars.
Next we considered the LMC LBVs. Six are listed by Bohannan (1997): S Dor, R 71, R 127, HD 269582, R 110, and R 143. To this list we propose that R 85 be considered a seventh, based upon our discovery here of spectral variability (Section 3.1.1.1) and a recent characterization of its photometric variability (van Genderen, Sterken, & de Groot 1998; see also Stahl et al. 1984). Of these seven, S Dor and R 85 are in LH 41, which is included here, and R 143 is in LH 100, which is not. We argue later that one of the LH 85 stars may also be an LBV based upon its spectral similarity to other LBVs, but further monitoring is needed to establish variability; we include it in Table 1 as a previously unknown, high luminosity B\[e\] star. Three other “LBV candidates” are listed by Parker (1997) : R 99, S 61 (BE 153=Sk$`67^{}266`$), and S 119 (HD 269687=Sk$`69^{}175`$). Of these, only one is located near an OB association (R 99 near LH 49), and it is not included here. Finally, we also considered the location of the high luminosity B\[e\] stars (Table 1 of Lamers et al. 1998; see also Zickgraf et al. 1986, Zickgraf 1993, and in particular Fig. 10 in Gummersbach, Zickgraf, & Wolf 1995). Only S 134, is found in one of our regions (LH 104), although several are found in other OB associations; i.e., S 22 in LH 38 and R 82 in LH 35.
We have referred to all of these stellar aggregates as “OB associations”, although the distinction between an OB association, and a bona-fide “cluster” young enough to contain O-type stars, is hard to quantify. The classical distinction, that clusters are gravitationally bound, is hard to establish, as it requires a census down to the low-mass components, plus detailed radial velocity studies. Semantics aside, our primary concern is to what degree these regions are coeval. Certainly most of the OB associations studied as part of our efforts to determine the IMFs are (Massey et al. 1995b). For the new ones studied here, we will establish the degree of coevality directly from the data.
Our observing strategy had similarities to our work that determined the initial mass functions in the LMC (e.g., Massey et al. 1989a, 1995b). It is possible to infer masses of main-sequence O- and B-type stars using their position in the physical H-R diagram ($`\mathrm{log}T_{\mathrm{eff}}`$ vs. $`M_{bol}`$) and comparing these with modern evolutionary models. There may be systematic problems with the masses thus inferred, although there is good agreement with the overlap of masses determined directly from spectroscopic binaries up to 25$`_{}`$ (Burkholder, Massey, & Morrell 1997), above which mass there is a scarcity of suitable data on binaries. Massey (1998b) discusses the errors in the inferred mass with temperature; since the BC is a steep function of the effective temperature, accurate knowledge of the latter is needed for this procedure to work. Sufficient accuracy cannot be achieved from photometry alone, but knowledge of the spectral type of the star yields adequate information in most cases. The sort of error bars associated with this can be found in Fig. 1(c) and 1(d) of Massey et al. (1995b). We will revisit this issue in Section 4.3.
For this project we considered relying simply on the photographic photometry or aperture photoelectric photometry that was available; e.g., Lucke (1972) or Azzopardi & Vigneau (1982), for the Large and Small Clouds respectively. After all, for the stars with spectroscopy (and hence accurate BC determinations) an error of 0.1 mag in the $`BV`$ color will lead to a 0.3 mag error in $`M_V`$, given $`A_V=3.1\times E(BV)`$. An error of 0.3 mag in $`M_V`$ translates to an error of 15% in the derived mass (see details in Massey 1998b). (For comparison, if we were relying upon the colors alone and were dealing with a 0.1 mag uncertainty in $`BV`$ we would have a 2 mag uncertainty in the BC, and thus a 0.4 dex uncertainty in the log of the mass (i.e., a factor of 2.5 uncertainty in the mass of the star).
For determining the IMF, it is necessary to pursue spectroscopy down the main-sequence until spectral-type of early B or later, after which good photometry provides as accurate information. Yet, in the case of determining the turn-off masses in principle we need to only ascertain that we have obtained spectra of the most massive unevolved object in the association. In a strictly coeval population with uniform reddening, this will be equivalent to knowing the spectral type of the visually brightest member. However, given finite photometric errors, slight non-coevality, reddening which is spatially variable across a cluster, the presence of other evolved supergiants (either members or field interlopers), and the need to demonstrate coevality, our initial aim was to obtain spectra for the six or seven visually brightest stars in each of these associations. Still, this is far fewer than what would be needed to construct the IMF.
Some of these associations had extensive CCD photometry and modern spectroscopy in the literature, and for these we constructed H-R diagrams and obtained a few additional spectra where warranted. In other cases, we already had existing unpublished CCD photometry (and in some cases even spectroscopy) that had been aimed at determining the IMF; the complete data for these associations, and the IMF analysis, will be published separately elsewhere. For the most part, though, we began with published photographic photometry, using this list to select the appropriate (brightest and bluest) stars for spectroscopy, and subsequently obtained new CCD UBV data in order to better correct for reddening. In all cases we examined the preliminary H-R diagrams and then obtained spectra of the few remaining interesting stars, as needed.
## 3 New Data
We list in Table 1 the source of the data we used, be they new or from the literature, or both. For the new data, we identify the year in which it was obtained.
For most of the associations (LMC) we began with the photographic iris photometry of Lucke (1972) or older sources, and obtained spectra of the brightest and bluest stars during a run on the CTIO 1.5-m telescope during 1996 Oct 27-31. Grating 58 was used in second order with a CuSO<sub>4</sub> blocking filter, yielding wavelength coverage from $`\lambda 3750`$ to $`\lambda 5070`$ with approximately 3Å (2.8 pixels) resolution. The Loral chip was formated to 500 $`\times `$ 1200 (15–$`\mu `$m) pixels. The slit was opened to 1.5 arcsec (85$`\mu `$m) and oriented EW, except for crowded regions, where the slit angle was adjusted and/or the slit narrowed. A typical S/N of 100 per 3Å spectral resolution element was achieved in a 5 min exposure at $`V=12`$.
On the night following this run (i.e., 1996 Nov 1) we obtained UBV images of any OB associations without previous CCD data, using the Tektronix 2048$`\times 2048`$ CCD imager on the CTIO 0.9-m. The field-of-view (FOV) was 13.5 arcmin by 13.5 arcmin, quite ample for the typical 3 arcminute diameter OB associations in our sample. Exposure times were usually 100 sec in $`U`$ and 50 sec in each of $`B`$ and $`V`$. The night was mostly photometric, although the alert observing assistant reported seeing a single cloud pass by part way through the night; later we will argue that this affected the U photometry of two regions but nothing else. Standard stars were observed at the beginning, middle, and end of the night, and reduced satisfactorily (0.01 mag rms residuals in $`U`$, $`B`$, and $`V`$ in the fits to the solutions). Nevertheless, we treat the data as potentially non-photometric, comparing the derived reddening-free index $`Q=(UB)0.72\times (BV)`$ with that expected on the basis of spectral type as a check, as described in Section 4. As we discussed above, our photometric requirements are in any event modest, given our extensive spectroscopy.
About half of the OB associations in our sample had previously been imaged with an RCA CCD on the CTIO 0.9-m in 1985 October by two of the present authors (PM and KDE). The full details of these data are given in Massey et al. (1989a). Although the FOV was only $`2.5\times 4.0`$ arcmin in size, overlapping frames were taken when needed in order to include the whole of an OB association. The photometric integrity of these 1985 data is very high, as standard star observations were obtained over 10 photometric nights and used for precise determinations of zero-points and color-terms.
Similarly, some of the stars have previously unpublished spectroscopy obtained as part of our program to determine IMFs in the Clouds. Data obtained in 1989-1992 (Table 1) were taken on the CTIO 4-m telescope with the RC spectrograph. The details of these data were given by Massey et al. (1995b); here we will simply note that they were of comparable spectral resolution (3Å), and covered at least the wavelength region from Si $`\lambda 4089`$ through He II $`\lambda 4686`$. The S/N were typically 75 per 3Å spectral resolution element.
After our preliminary HRDs were constructed, we had two observing opportunities to obtain additional spectra where warranted. On 1999 Jan 3-7 we used the CTIO 4-m for significantly higher resolution and better S/N data. Grating KPGLD was used in second order with a CuSO<sub>4</sub> filter resulting in a resolution of 1Å (2.5 pixels) and a wavelength coverage of 3730Å to 4960Å using the Loral $`1024\times 3100`$ (15 $`\mu `$m) CCD. The S/N obtained was typically 160 per 1Å resolution element. We obtained one final observation for this project on 1999 Oct 21 using the CTIO 1.5-m.
### 3.1 Analysis
#### 3.1.1 Spectroscopy
We classified the spectra with reference to the Walborn & Fitzpatrick (1990) spectral atlas of O and B stars. Based upon our internal consistency and previous experience we expect that the spectral subtypes are determined to an accuracy of one subclass and one luminosity class (e.g., supergiant vs giant), except for the earliest O-type stars, for which there is little or no ambiguity in subclass. (See discussion in Massey et al. 1995a, 1995b.)
There is no metallicity dependence in classifying hot stars as to spectral subclass, as the primary spectral type (effective temperature) indicators are the relative strengths of different ionization states of the same ion; e.g., He I vs. He II for the O-type stars, and Si IV vs. Si III for the early B-type stars; however, it is our experience that the luminosity indicators are metallicity dependent, even for the O-type stars. This makes physical sense—in fact, it would be hard to see how this would fail to be the case—as the O-type luminosity indicators are primarily indicators of the strength of the stellar wind (i.e., He II emission vs. He II absorption). The B-type luminosity indicators rely upon how strong certain metal lines are relative to, say, He, and again we expect this to have a metallicity dependence. We therefore always checked the “MK” luminosity class with that expected on the basis of the absolute magnitudes, as described below; we note cases where we have adjusted the luminosity class based upon the absolute magnitudes.
All told, we classified slightly over 200 stars. We include our classification, as well as those from the literature, in the catalog we describe in Section 3.2. Here we will illustrate and comment on just a few of the more interesting spectra.
##### 3.1.1.1 R 85.
We propose that the luminous star R 85 in LH 41 be considered an LBV. Based upon their characteristic of its photometric variability, van Genderen et al. (1998) state that the star is “undoubtedly an active LBV.” We show in Fig. 6 some of the spectral changes that have taken place in recent years; we agree with van Genderen et al.’s characterization. Feast, Thackeray, & Wesselink (1960) classify the star as “B5 Iae”, and note the presence of H$`\beta `$ emission, H$`\gamma `$ and H$`\delta `$ absorption, as well as its photometric variability. Our 1996 spectra did not appear totally consistent with this description, as Mg II $`\lambda 4481`$ was present but there was little or no He I $`\lambda 4471`$; for a B5 star the latter should be somewhat stronger. We took a very high signal-to-noise spectrum with the CTIO 4-m in January 1999, and were surprised by the rapid and strong changes present since 1996; the newer spectrum shows the star to be hotter (based upon He I to Mg II) with much stronger lines. Dr. B. Bohannan was kind enough to make available a photographic spectrogram he obtained in 1985 on the Yale 1-m, along with a sensitometer exposure; there is very good agreement between his exposure, and what we obtained 11 years later. The recent change in the spectrum of R 85 suggests that further monitoring would be of interest. The photometry listed in Table 2 comes from the 1 Nov 1996 observation; e.g., $`V=10.53`$, $`BV=0.16`$, and $`UB=0.81`$. In the 1985 data (28 Nov) the star was slightly brighter: $`V=10.44`$, $`BV=0.12`$, and $`UB=0.71`$.
##### 3.1.1.2 Newly Identified O3 Stars.
As part of this investigation we came across a number of previously unrecognized O3 stars, stars whose effective temperatures are at the extreme of the spectral sequence of luminous stars. We show examples in Figs. 6 and 6.
First, let us consider the O3 supergiants (O3 If\*) and giants \[O3 III(f\*)\]. These evolved stars are still in the temperature regime covered by the O3 classification, and thus all such stars must be extremely massive. Walborn et al. (1999) classify the star LH90$`\beta `$-13 as O4 If+ on the basis of an FOS spectrum obtained with HST, but our higher signal-to-noise spectrum (with higher resolution) reveals N V $`\lambda \lambda 4603,19`$ absorption; this, combined with the lack of He I makes this an O3 star (Fig. 6). The star ST5-31 in LH 101 was classified as O3 If\* by Testor & Niemela (1998); our spectrum is in good agreement with that. We consider the star W16-8 in LH 64 an O3 III(f\*) owing to the relative weakness of He II $`\lambda 4686`$, despite the extremely strong N IV $`\lambda 4058`$ emission and very strong N V $`\lambda 4603,19`$, usually indicative of high luminosity; the absolute magnitude we derive in the next section is $`M_V=5.4`$, consistent with this classification, and reminding us that slight abundance anomalies can mask as luminosity effects in early-type stars. A detailed atmospheric analysis of this star is in progress in collaboration with Rolf Kudritzki.
Among the O3 dwarfs (Fig. 6) we include ST2-22 (in LH 90). This star was previously recognized as an O3, but called a giant by Testor et al. (1993). The lack of emission at He II $`\lambda 4686`$, and the weakness of N IV $`\lambda `$ 4058, suggest a lower luminosity class. We classify W28-23 in LH 81 as an O3 V((f)). The star ST5-27 in LH 101 was called an O4 V by Testor & Niemela (1998). The spectrum of this star is strongly contaminated by nebular emission lines. We tentatively adopt an O3 V((f)) spectral type, but our data are not inconsistent with the O4 V((f)) designation; we do not show the spectrum as the nebular lines makes casual comparisons difficult. Another star in LH 81, W28-5, appears to be intermediate between O3((f)) and O4 V((f)): the strength of He I $`\lambda 4471`$ relative to He II $`\lambda 4542`$ would argues that the star is a little bit later than O3, but there is N V $`\lambda 4602,19`$ present on our high signal-to-noise spectra, and this has usually been taken as characteristic of O3s.
The presence of He I $`\lambda 4471`$ is easy to discern on the O3 stars in Fig. 6 because of the extraordinarily high S/N (160 per 1Å resolution element). The O3 class was introduced by Walborn (1971) to describe four stars in Carina which showed no He I $`\lambda 4471`$ on well-widened IIa-O emulsion spectrograms obtained at modest resolution (2Å). When finer-grain plates were used at higher resolution, He I $`\lambda 4471`$was detected with equivalent widths of 120-250 mÅ by Kudritzki (1980) and Simon et al. (1984) for three of the Carina stars. Here we find that He I $`\lambda 4471`$ lines have equivalent widths of 75 mÅ in W28-23, and 105 mÅ in ST2-22, significantly smaller than that measured for the stars which first defined the class. Yet modern spectroscopy makes it possible to readily detect these lines.
##### 3.1.1.3 Other O-type Stars.
There are clearly other exceptions to the premise that N V $`\lambda 4603,19`$ absorption is indicative of a luminous O3 star. In Fig. 6 we show the spectrum of ST5-52, a star in LH 101 classified by Testor & Niemela (1998) as O3 V. However, the strength of He I suggests a considerably later O5.5 type. It is easy to infer the basis for the Testor & Niemela classification of this star: our spectrum shows both NIV $`\lambda 4058`$ emission and N V $`\lambda 4603,19`$ absorption, typically assumed to be only characteristic of luminous O3 stars. One possibility is that this star is a spectrum binary, consisting of an O3 III(f\*) plus a later O-type companion, which contributes the He I. However, we propose instead that this is a “nitrogen enhanced” star, and classify it as ON5.5V((f)). We prefer this latter explanation because we have identified another LMC star, not connected with the present study, whose He I to He II ratios are consistent with an O5 type, but which also shows N IV emission and N V absorption. Detailed atmospheric analysis is underway for both stars, pending HST data.
The star LH58-496 was classified as “O3-4 V” by Garmany, Massey, & Parker (1994). Our high S/N spectrum (Fig. 6) obtained with the CTIO 4-m shows a somewhat later spectral type, O5V((f)). In Fig. 6 we also show two other early-type dwarfs, an O5 V((f)) star and an O4 V((f)) star.
We illustrate a few newly discovered luminous O-type supergiants in Fig. 6. Examples shown here include supergiants from O4 through O8.
##### 3.1.1.4 A Reconsideration of Br 58 as a WR star, and A Newly Discovered WR Star.
The star Br 58 in LH 90, has long been recognized as a WN Wolf-Rayet star. Testor et al. (1993) classify it as WN6-7, while earlier work has classified it as WN5-6 (Conti & Massey 1989). We illustrate its spectrum in Fig. 6 from a new high-dispersion, high S/N observation. We note that our ground-based spectrum shows strong N V $`\lambda 4603,19`$ absorption; this, plus the considerable strength of its absorption line spectrum, would tempt us to reclassify this as an extreme O3 If\* star, i.e., O3If\*/WN6. (See Fig. 3 in Massey & Hunter 1998.) These stars are believed to be young, H-burning hot stars whose very high luminosities result in sufficiently strong stellar winds to mimic the strong emission characteristic of a WR.
The star Sk$`69^{}194`$=W28-10 in LH 81 is a newly discovered WR star, of type B0 I+WN. The spectroscopic discovery of another WR star in the LMC is not surprising, particularly given the weakness of the emission in this object. (The equivalent width of He II $`\lambda 4686`$ is $`2`$Å, compared to typical $`30`$Å for a very weak-lined WN star; presumably this is due to the continuum being dominated by the B0 I component.) We question below whether all B0 I+WN are truly binaries.
#### 3.1.2 Photometry
UBV photometry is needed only (a) to determine accurate $`M_V`$ values for the stars with spectra, and (b) to check that we obtained spectra for all of the likely “most massive unevolved star” candidates. In order to accomplish (a) we typically needed $`V`$ and B-V data for half a dozen stars or so in each association, and to accomplish (b) we also required U-B, in order to construct a reddening-free index. Nevertheless, with modern techniques it proved just as easy to measure photometry for all stars on a frame, typically several thousand stars. At least we could then be assured that the brightest stars were well-measured, in the sense that their photometry was not contaminated by resolved neighbors.
We did this by fitting point-spread-functions (PSFs) using DAOPHOT implemented under IRAF. The 1996 CCD frames were measured by E.W., while the 1985 data were measured by P.M. The method used is similar to that described by Massey et al. (1989a) and we will give only a brief overview here. Automatic star-finding algorithms were used to identify stellar sources down to the “plate-limit” (typically 4$`\sigma `$ above the noise). Aperture photometry through a small digital aperture (with a diameter corresponding roughly to the full-width at half-maximum of the stellar images) were then run in order to determine the local sky values for each star (determined from the modal value in an annulus surrounding each star) and to determine the instrumental magnitude to assign to the PSF stars. For each frame isolated, well-exposed stars were chosen to define the PSF. This PSF was then simultaneously fit to all of the stars whose brightnesses could possibly overlap. In regions of nebulosity, the sky value was also fit separately; otherwise, an average sky value was adopted for all the stars in a given fitting exercise. A frame in which the fitted PSFs had been subtracted was then examined to see how well the PSF matched and to look (by eye) for any stars that had been buried in the brightness of other stars. In addition, the $`U`$, $`B`$, and $`V`$ frames were blinked along with the fitted xy centers to make sure there was consistency. Missing stars were added back into the star list and a final run was made on each of the three colors. Aperture corrections were then determined for each frame in order to correct the instrument zero-point (based upon the small digital aperture) to the large apertures used to measure the standard stars. These instrumental magnitudes were then transformed to the standard system. In the case of the 1985 RCA CCD data there were often overlapping frames involved in covering a region, and the final photometry was combined to produce a single star list, with stars with multiple entries averaged.
One region, Lucke-Hodge 41, was common to both data sets, and thus served as an end-to-end independent check on the final, transformed photometry. If we consider the twenty brightest stars (in $`V`$) we find a mean difference (new minus old data set) of $`+0.015`$ mag in $`V`$, $`+0.011`$ mag in $`UB`$, and $`+0.014`$ in $`BV`$, with sample standard deviations of 0.06 mag, 0.02 mag, and 0.04 mag, respectively. If two outliers are removed from the $`V`$ data, and one from the $`UB`$ data, the mean differences become +0.002 mag and $`+0.001`$ mag, respectively with standard deviations of 0.03 mag and 0.04 mag. This agreement is excellent, and suggests that no systematic differences exist between the two data sets over the magnitude and color ranges of interest.
### 3.2 The Catalog
We list in Table 2 the brightest stars in each of the 14 associations for which we have new photometry; existing and new spectral types are also given. We include all stars of magnitudes $`V=15`$ or brighter; in several cases we extended this to fainter magnitudes to include additional stars with spectral types or, in the case of NGC 602c, to include at least 10 stars. For two of the associations (LH 58 and LH 101) we reply upon cited studies (cf., Table 1) but have a few new spectral types; we include these in Table 2. (For three addition associations, NGC 346, LH 9, and LH 47, we reply purely on the cited works in Table 1.)
In listing the stars we make use of published names where available finding charts exist, although the celestial coordinates given in Table 2 should be of sufficient accuracy to remove the need for finding charts. For the LMC, we have kept with the star numbering given in the finding charts of Lucke (1972), with additional stars given designations of 1000+ so as to avoid confusion. The exceptions are those associations with modern CCD studies, where we have kept with the numbering scheme employed by the authors. In a few cases the associations contained stars that were saturated on our CCDs (typically $`V<10`$); we include photometry of these stars from the literature. We describe below details related to each association, making reference to the results obtained in subsequent sections.
#### 3.2.1 Descriptions of Individual Associations
NGC 346: We rely on the CCD photometry and spectroscopy of Massey, Parker, & Garmany (1989b). The imaging data had their source in the same observing run as the 1985 imaging used for many of the other associations studied here. Four of the brightest stars were also subjected to detailed analysis by Kudritzki et al. (1989). Reanalysis of these stars by Puls et al. (1996) was used in the spectral type to effective temperature calibration of Vacca, Garmany, & Shull (1996), which we adopt in the next section; we note here that despite the different methodology involved, the masses determined by Puls et al. for these stars are in good agreement with those we compute in the following sections. The visually brightest star is HD 5980, the WN3+abs Wolf-Rayet that underwent an LBV-like outburst. The second brightest star is the O7If star Sk 80. More than a magnitude fainter visually are the very early O-type stars first found by Walborn (1978), Walborn & Blades (1986), and Niemela, Marraco, & Cabanne (1986).
Hodge 53: Our photometry here is a comprehensive mosaic of several CCD frames and extensive spectroscopy obtained with the goal of determining the IMF. However, the the region is not condensed, and there are several stars of type A-F and later, some of which are apparently foreground dwarfs or giants, and others which are SMC supergiants. Our spectrum of AV 331 shows it to be an SMC member of type A2 I, based both on its radial velocity, appearance of the hydrogen lines, and the strength of Fe II $`\lambda 4233`$ (see Jaschek & Jascheck 1990, Fig 10.2). However, our spectrum of AV 339a shows it to be an F2 foreground star, probably a dwarf, based both on its radial velocity and lack of luminosity-sensitive Sr II $`\lambda 4077`$. A fainter star, h53-144, is an A8 foreground dwarf. We lack spectra for the other yellow stars, and so we cannot comment further on their membership. Our spectroscopy has also identified a double-lined spectroscopic binary (O4 V+O6.5 V) which is among the most bolometric luminous members. When we construct the HRD, we will consider that each of the two components contributes equally to the visual flux, consistent with the appearance of our double-lined spectrum, and the expected $`M_V`$s of stars of these spectral types. The visually brightest member is the WR binary AV 332=Sk 108=R 31=AB 6 (WN3+O6.5) with a 6.54 day orbit (Moffat 1982, 1988; Hutchings et al. 1984; Hutchings, Bianchi, & Morris 1993). Hutchings et al. (1984) argue convincingly that the O-type companion dominates the visual flux by a factor of 10 to 1 (making it of luminosity class “I”), and that its location in the HRD suggests an initial mass of $`7080_{}`$, consistent too with its Keplerian mass. Our analysis will yield a very similar value. The other WR member, AV 336a=AB 7, is quite a bit fainter. The WR component is likely a WN3 (Moffat 1988), although all that is certain is that it is earlier than WN7 (Conti, Massey, & Garmany 1989). An O-type absorption spectrum is also seen. Recent work by Niemela (1999) suggests a 19.6 day period.
NGC 602c: NGC 602 is located in the wing of the SMC; the region was studied by Westerlund (1964), who identified three sub-components. Components “a” and “b” are adjacent and are immersed in nebulosity known as N90 (Henize 1956); component “a” is also known as Lindsey 105 (Lindsey 1958). Here we are concerned with the third component, “c”, which is an isolated condensation with little nebular emission. It was designated as a separate association both by Lindsey (1958) and Hodge (1985), and is known as “Lindsey 107”, and “Hodge 69”. (See Plate 5 and Figure 1 in Westerlund 1964.) We obtained new CCD photometry of NGC 602c. Its visually brightest star is the WR star AB 8, the only WC star known in the SMC. It has enhanced oxygen, and was classified by Conti et al. (1989) as “WO4 + abs”. (Crowther, De Marco, & Barlow 1998 instead call the WR component “WO3”.) A new spectrum of the star obtained as part of the present program suggests that the absorption spectrum is O4 V. Moffat, Niemela, & Marraco (1990) present an orbit for this system with a period of 16.644 days. They propose spectral types of WO4+O4 V, with which we concur, although Kingsburgh, Barlow, & Storey (1995) suggest a somewhat later type for the O star.
LH5: Our photometry and spectroscopy are the first modern study of this association. The visually brightest star is Sk $`69^{}30`$, a G-type supergiant according to Feast et al. (1960), with the next brightest star an O9 I. The WR star, Br4, was described as “WN2” by Conti & Massey (1989), as no N lines are visible, similar to the WN2 Galactic star HD 6327. Like that star, Br 4 has a faint absolute magnitude. We will find in subsequent sections that the star has a normal bolometric luminosity, and that its faintness is presumably due to a very high temperature, which shifts its light into the unobserved UV. In constructing our HRD we find that the G5 Ia star Sk $`69^{}30`$ is coeval with the rest of the massive stars.
LH9: This association was studied in detail by Parker et al. (1992), using the same 1985 imaging data and calibration that we employ here for many of the other associations. The central object, HD 32228, was clearly an unresolved cluster of many early-type stars, with a composite WC+O spectral type. The region was recently examined using HST by Walborn et al. (1999), and we adopt their photometry and spectroscopy here, ignoring the region outside of the central 30 arcsec covered by the PC frame of WFPC2. Although they were able to spectroscopically observe the WC component separately from its close neighbors for the first time, their spectral classification of WC4 is based upon only a spectrum in the blue, which lacks the crucial classification lines O V $`\lambda 5592`$, C III $`\lambda 5696`$ and C IV $`\lambda 5812`$ (e.g., Smith 1968a; van der Hucht et al. 1981). Walborn (1977) had earlier classified the WR star as WC5, but this was also based upon a blue spectrogram. Smith (1968b) called the star WC5, but this was before the earlier WC4 subclass was introduced. Breysacher et al. (1999) cite a speckle study by Schertl et al. (1995) for the spectral type, but no spectrum was actually taken as part of that study. We adopt WC4 as the spectral type, but note here that the type is uncertain. The visually brightest stars in the LH 9 association are late-O supergiants (O9 I and O8.5 I).
LH12: Ours is the the first modern study of this association. It contains the WC4 star Br 10. The visually brightest stars are B-type supergiants, although our study has revealed a very early O-type star, with type O4 V(f). To the extent that the association is coeval, the B-type supergiants evolved from stars of spectral subtype O4 V or even earlier.
LH31: This association contains two Wolf-Rayet stars, Br 16 classified by Conti & Massey (1989) as WN2.5. A second WR star has been recently discovered by Morgan & Good (1985), who classify the star as WC5+O. This star is BAT99-20 in the catalog of Breysacher et al. (1999), whose finding chart puts the star centrally located in the association boundary shown by Lucke (1972). Nebulosity prevented Lucke from photographic photometry of any by the brightest few stars. The visually brightest stars include a B1 III, an O6 I(f), and two yellow stars. One brighter of these, which we call LH31-1002, is apparently an LMC F2 supergiant, based both upon our measured radial velocity and strong Sr II $`\lambda 4077`$ (see Jaschek & Jaschek 1990). The other is clearly a late F-type foreground dwarf, based upon its radial velocity and its lack of Sr II.
LH39: The cluster was examined by Schild (1987), and again by Heydari-Malayeri et al. (1997). We obtained new photometry and a few additional spectral types. The association contains one of the rare Ofpe/WN9 stars, Br 18=Sk$`69^{}79`$. Ardeberg et al. (1972) list the star Sk$`69^{}80`$ as having a spectral type of F2 Ia; however, Schild (1987) suggests a type of B8: I. Our photometry is consistent with something intermediate between these two, and we will use its photometry to place it in the HR diagram. (The radial velocity of Ardeberg et al. does confirm it is an LMC member.) We will find that two A supergiants classified by Schild (1987) appear to be much older than the rest of the cluster. We have independent spectroscopy for one of these, LH39-22, and confirm Schild’s type.
LH41: This association contains S Doradus, the prototype LBV, and the visually brightest star in the cluster. The second brightest star, R 85=Sk-69$`69^{}92`$ we propose as an LBV, based upon its spectral and photometric variability, as discussed earlier in Section 3.1.1.1. The third brightest star is the Wolf-Rayet star Br 21, classified by Conti & Massey (1989) as B1Ia + WN3. The star LH41-4 is of M-type, but we lack the radial velocity information that would ascertain whether this is an M supergiant or foreground dwarf. There are two lower luminosity but bona fide A-type supergiants, and an F5 supergiant. The latter has been confirmed based upon our radial velocity and the strength of Sr II $`\lambda 4077`$. (It is also an excellent match to the F5Iab star HD 9973 shown in the Jacoby, Hunter, & Christian 1984 atlas.) Ours is the first modern study of this association.
LH43: The visually brightest star is an early M-type, but again we lack the proper radial velocity information to ascertain whether this is an LMC member or not. The second brightest star is a newly discovered O4 If star. The WR star Br23 is classified WN3.
LH47: This association was studied by Oey & Massey (1995) and Will, Bomans & Dieball (1997). We adopt the photometry and spectroscopy of the former, who obtained spectral types for all the brighter components, primarily of early to mid O-type. Oey & Massey (1995) suggest that there are two ages for the stars in the LH47/48 region: stars interior to the DEM 152 superbubble have an older age than stars in rim of the bubble. The WR star and other massive stars of interest are on the exterior, and we will restrict our analysis to those. In agreement with Will et al. we find no difference between the photometric $`Q`$ and that expected on the basis of spectroscopy; we cannot comment on their assertion that field-to-field differences exist in the individual $`BV`$ and $`UB`$ colors at the 0.15 mag level, other than to note our value for the reddening appears to be reasonable.
LH58: This association was recently studied by Garmany et al. (1994). It contains three WR stars, Br 32 (WC4+abs), Br 33 (WN3+abs), and Br 34 (B3I+WN3). The latter is the visually brightest star. We did obtain a spectrum of the earliest-type star in the association, reclassifying it from O3-4 V to O5.5 V((f)), as described earlier (Section 3.1.1.3). We note that LH58-473 as B0.5V must be a giant based upon its M<sub>V</sub>.
LH64: This association was studied by Westerlund (1961) as well as by Lucke (1972). Ours is the first modern study. The three visually brightest stars have colors characteristic of mid-to-late type stars, presumably foreground, although spectroscopy is needed to determine if they are supergiants. The WR star Br 39 was not classified by Conti & Massey (1989), but was called WN3 by Breysacher (1981).
LH81: Also studied by Westerlund (1961) and Lucke (1972), ours is the first CCD study of this interesting region. It contains three WR stars: the WC4 star Br 50 (classified by Conti & Massey (1989), the WN4+OB star Br 53 (classified by Breysacher 1981), and Sk$`69^{}`$194, discovered as a WR star here (B0I+WN). The visually brightest star is a foreground G dwarf. We identify two very early-type stars in the association, W28-23, a O3 V((f)) star, and W28-5. As discussed in Section 3.1.1.2, we classify the latter as O4 V((f)) based upon its He I to He II strengths, but our very high S/N spectrum shows the definite presence of N V $`\lambda 4603,19`$ absorption lines, previously associated only with O3 stars. Possibly an intermediate type (O3.5) would be warranted, but we leave that until we have been able to complete a detailed analysis of this star.
LH85: We identify the star LH85-10 as a newly discovered B\[e\]. Our study is the first since Westerlund (1961) and Lucke (1972). The association also contains the WR star Br 63, classified as WN4.5 (Breysacher et al. 1999). Westerlund (1961) treated this association and the neighboring LH 89 as one unit; we treat them separately here, following Lucke (1972), although the ages and cut-off masses we derive will prove to be essentially the same. The earliest spectral type we find in LH 85 is B0.5.
LH89: A section of LH89 was included in the study by Schild & Testor (1992) of stars in the general 30 Doradus region (their “zone 3”), in addition to the Westerlund (1961) and Lucke (1972) studies. We have used their spectral types as a supplement to our own, but use our own CCD photometry. The association contains Br6 (WN4) and Br 64=BE 381, the archetype of Ofpe/WN9 stars. The visually brightest stars are three tenth magnitude stars of intermediate color; radial velocities of the two brightest demonstrate that they are LMC members (Ardenberg et al. 1972). Our spectrum of the third shows it is a foreground F8 dwarf, based both on its radial velocity and the weakness of high-luminosity features in the spectrum, emphasizing once again the need for spectroscopy in determining membership of even bright stars in the Clouds. We will find that the two confirmed A-F supergiants turn out to be coeval with the rest of the association members.
LH90: Photometry of the LH 90 region was published by Schild & Testor (1992), who refer to the region as “Zone 2”, and provide a finding chart in their Figure 3. (Only stars 2-33, 2-34, and and 2-45 fall outside the association boundary shown by Lucke 1972.) There are three clumps of stars, designated as “clusters” $`\alpha `$, $`\beta `$, and $`\delta `$ by Loret & Testor (1984). The region was re-examined by Testor et al. (1993), who provided new photometric and spectroscopic data on knots $`\alpha `$ and $`\beta `$. Clusters $`\beta `$ and $`\delta `$ were also studied by Heydari-Malayeri et al. (1993). Recently, Walborn et al. (1999) were largely successful in further unraveling the $`\beta `$ knot of stars using WFPC-1 images and FOS spectroscopy with HST. (They refer to $`\beta `$ alternatively as “NGC 2044 West” and “HDE 269828”.) To this, we add our own UBV photometry and spectroscopy. We note that a comparison of the high resolution image of Testor et al. (his Fig. 1b) with that of Walborn et al. (Fig. 5) suggest that ground-based work actually did a remarkably good job of resolving multiple components in cluster $`\beta `$. The stars designated “TSWR2” and “TSWR1” are multiple, but the others are actually well resolved with 1” resolution. The components found independently by our PSF-fitting are an exact match to those identified by Testor et al. The most interesting star is the one Testor et al. identify as “6” in cluster $`\beta `$; this is the star labeled “9” by Heydari-Malayeri et al., and split into two components (“9A” and “9B”) by Walborn et al., although 9B is 1.5 mag fainter than 9A and hence the composite spectrum we obtained from the ground is a good representation of star 9A. We have noted earlier (Section 3.1.1.2) that the star $`\beta 13`$ is probably better considered an O3 If\* star rather than the O4If+ used by Walborn et al.
In our analysis of this region we will make use of our new ground-based data, but defer to the HST data of Walborn et al. for stars for the the group of stars called “TSWR1” (or $`\beta `$-6) by Testor et al., which is the star identified as “5” by Heydari-Malayeri et al., split into multiple components by Walborn et al. (1999). Our ground-based (composite) spectrum would have resulted in a “B0I+WN” designation, but the HST work clearly shows that these are separate stars, in accord with Testor et al.’s finding. One wonders if other “BI+WN” systems might not be similarly resolved. We also note the need for a high-resolution study of the $`\delta `$ knot in this interesting region.
In addition to the WN4 component of “TSWR1”, the association contains many other WRs: Br 56 (WN6), Br 57 (WN7), Br 58 (WN5-6), and Br 65 (WN7), all of fairly late type for the LMC, plus the WC4 star Br 62. The classifications are from Conti & Massey (1989), except for that of Br 65, which is from Breysacher (1981). Earlier (Section 3.1.1.4) we suggest that Br 58 may be better classified as O3If\*/WN6.
In analyzing this cluster in Section 4.3, we find that the $`\beta `$ subclump is no more coeval than the association as a whole, as witness the fact that both a B0 I star of modest luminosity cohabits with an O3 star of high luminosity. There is a significant range of ages.
LH101: This region has recent CCD photometry and spectroscopy by Testor & Niemela (1998). To this, we obtained our own spectra for three of the stars, as discussed in Section 3.1.1. We find that ST5-27 is an O3V((f)), as indicated both by the lack of He I and the weak presence of N V $`\lambda 4609,19`$ absorption; the star was classified as O4 V by Testor & Niemela. We confirm that their star ST5-31 is indeed an O3If. And, we reclassify ST5-52 as an ON5.5V((f)) star, rather than O3 V (Section 3.1.1.3). The association contains Br 91, another of the rare Ofpe/WN9 objects.
LH104: This association was also studied by Testor & Niemela (1998). We have obtained new CCD photometry, as well as additional spectroscopy. The association contains three WRs, all of which are spectrum binaries as described by Testor & Niemela: Br 94 (WC5+O7), Br 95 (WN3+O7), and Br 95a (WC5+O6). The visually brightest star is the B\[e\] star, S 134 (Zickgraf 1993). We note that one of the visually brighter stars is an M star, confirmed by Testor & Niemela as a supergiant on the basis of its radial velocity; this agrees with the conclusion of Massey & Johnson (1998) that WRs and M supergiants are sometimes found in the same associations, contrary to the prevailing wisdom.
## 4 Construction of HRDs: Coevality and Uncovering the Most Massive Stars
In order to identify the most massive stars, we construct “physical” H-R diagrams ($`\mathrm{log}T_{\mathrm{eff}}`$ vs. $`M_{bol}`$) for comparison with the theoretical evolutionary tracks. These tracks will allow us to test for coevality, and determine the masses for the highest mass unevolved (H-burning) stars in these associations. First, we must correct the observed photometry for reddening, and second to convert the data (spectral types and photometry) to effective temperatures and bolometric magnitudes. Next we will construct the HRDs and uncover the masses of the most massive stars.
### 4.1 Corrections for Reddening and Testing the Reddening-free Index $`Q`$
Our first step in constructing HRDs is to determine the reddening corrections for each region. For stars with spectral types, we adopt the intrinsic colors of FitzGerald (1970) as a function of spectral type and compute the color excess $`E(BV)`$ directly. Occasionally even a star with a spectral classification has a reddening which differs substantially from the other members in a region, and so we’ve chosen to constrain the reddening to the range indicated by the majority of stars for which there are spectral types. We list in Table 3 the average color excess $`\overline{E(BV)}`$ and ranges of $`E(BV)`$ we adopt for each of the 19 associations. (For consistency, we re-derived reddenings even for the associations with values already in the literature.)
Although we obtained spectral types for most of the bright stars in each association, there are some stars for which we have only photometry. Rather than de-redden these using $`\overline{E(BV)}`$ we employed a relationship between $`Q`$ and $`(BV)_o`$ to de-redden each star individually, using the star’s photometry and $`\overline{E(BV)}`$ as a gauge of whether the star’s intrinsic colors were sufficiently blue for this method to work. We found that for stars with $`Q<0.2`$ for $`(BV)_o(BV)\overline{E(BV)}=0.06`$ we could de-redden star by star; for stars with intrinsic colors redder than this amount, we adopted the average reddening. We did further constrain the reddening to the range determined by the majority of stars with spectral types in a region.
Since our earlier work (Massey et al. 1989b, 1995b) it has become clear that the intrinsic colors as a function of spectral type or effective temperatures are not extremely well know, particularly for the early B supergiants, and we have therefore computed new relationships based $`Q`$ and $`(BV)_o`$ (and the intrinsic colors and effective temperatures) using the Kurucz (1992) ATLAS9 models, using a metallicity of 0.8 times solar, a compromise between SMC, LMC, and (local) Galactic abundances. We find
$$(BV)_o=0.005+0.317\times Q$$
regardless of luminosity class.
Construction of the reddening-free index $`Q`$ for the stars with spectral type allows an independent check upon the accuracy of the photometry: is there good agreement between the observed $`Q`$ and that $`Q`$ expected on the basis of the intrinsic colors for that spectral type? We determine if there is a statistically significant shift in $`Q`$ for all the stars for which we have spectral types in each association. In general we find deviations in $`Q`$ within 1$`\sigma `$ of 0.0. The only exceptions for our new photometry are LH 43, for which we adopt a shift $`\mathrm{\Delta }Q=0.13`$, and LH 64, for which we adopt a shift $`\mathrm{\Delta }Q=0.15`$ (i.e., in both cases the photometric $`Q`$ must be made more negative to agree with the expectations of the spectroscopy). The two regions were imaged within a few minutes of each other during the 1996 night at about the same time that the observing assistant reported seeing an isolated cloud. Interestingly, the reddening values we found for these two regions are each quite reasonable, suggesting that it might have been only $`U`$ which was affected in the two fields. Inspection of the observing logs confirms that the U exposure of LH 43 was observed back-to-back with the U exposure of LH 64. The next regions observed, LH85/89, appears to have no significant photometric problems. We see no problems with any of the 1985 data, either published or new in this paper. We do find a shift of $`\mathrm{\Delta }Q=0.11`$ for the LH 101 photometry published by Testor & Niemela (1998). Although the large scatter (0.08 mag) makes this result marginal in significance, and nearly all the stars of interest to us have spectral types, we still apply this correction to their photometry.
The WFPC2 photometry of LH 9 (“HD 32228”) by Walborn et al. (1999) also shows a systematic shift in $`Q`$, with $`\mathrm{\Delta }Q=0.07\pm 0.01`$(s.d.m.) mag. Presumably this shift is an artifact of their reduction procedure. This shift is larger than any of the ground-based UBV data reported here, other than the cases noted above, and so it is unlikely due to any problems with the spectral-class to $`Q`$ relationship we adopt. We did not apply any correction to their data as we used only the stars with spectral types in constructing the HRD, although this could have some minor effect on the absolute magnitudes (0.2 mag) and hence masses we determine if the problem is in $`BV`$ rather than in $`UB`$.
### 4.2 Conversion to $`\mathrm{log}T_{\mathrm{eff}}`$ and Bolometric Luminosity
The final step in constructing the HRDs is to use the data to determine the effective temperature and bolometric luminosity of each star.
For stars with spectral types, we begin by adopting the spectral type to effective temperature scale given by Vacca et al. (1996) for O-type stars, based as it is on the results of modern hot-star models. This will yield results that are somewhat hotter and, thus, somewhat more luminous and massive than the older effective temperature scale of Chlebowski & Garmany (1991), say, or that of Conti (1973). For the early B stars we were faced with a dilemma. As discussed by Massey et al. (1995a) there is a discontinuity in the effective temperature scales of hot stars corresponding to roughly where the modern work of Conti (1973) ended (i.e., O9.5) and earlier works took over. In order to smooth the transition, we have adopted the effective temperatures of B0.5-B1 dwarfs and giants as given in Table 3-4 of Conti (1988), as those are in excellent agreement both with what we expect on the basis of the intrinsic colors from the model atmospheres, and with the spectral analysis of Kilian (1992). For B1.5 and B2 dwarfs and giants, we compromised between the latter two. For the B-type supergiants, we made use of the effective temperatures suggested by Conti (1988), the recent spectroscopic analysis of two early B supergiants by McErlean, Lennon & Dufton (1998), a comparison of the intrinsic colors listed by FitzGerald (1970) with those of the Kurucz model atmospheres, and the effective temperature scale given by Humprheys & McElroy (1984). In the past we have relied exclusively on the latter; we note here though that this disagrees with the more recent analysis by 0.1 dex from B1 I through B5 I. It is clear that a consistent effective temperature scale that extends from O through the B-type stars is currently lacking, and the compromise we use here is only a stop-gap until a comprehensive study can be done.
For stars with photometry alone, we rely upon a relationship between the reddening-free parameter $`Q`$ and $`\mathrm{log}T_{\mathrm{eff}}`$ determined from the Kurucz models; this relationship is given in Table 4, and is appropriate for intrinsically blue stars \[($`Q<0.6`$ and either $`(BV)_o<0.00`$ or $`(UB)_o<0.6`$\]. For redder stars, we use a relationship between $`(BV)_o`$ and $`\mathrm{log}T_{\mathrm{eff}}`$ also given in Table 4, based upon the Kurucz models. The latter relationship need not be of high accuracy, as the BC becomes a less steep function of $`\mathrm{log}T_{\mathrm{eff}}`$.
The bolometric correction (BC) is a function primarily of effective temperature with little dependence on $`\mathrm{log}g`$; we adopt the approximation $`BC=27.666.84\times \mathrm{log}T_{\mathrm{eff}}`$ appropriate to hot stars ($`\mathrm{log}T_{\mathrm{eff}}>4.2`$) given by Vacca et al. (1996). For the cooler supergiants we find discrepancies between the BCs listed by Humphreys & McElroy (1984) and the corresponding effective temperatures when compared to the Kurucz models; we adopt the relationship given in Table 4 based upon a fit of the BCs with $`\mathrm{log}T_{\mathrm{eff}}`$ based upon the Kurucz models.
We show the resulting HRDs in Fig 6. In these figures, we have indicated the stars with spectral types by filled circles, and those stars with only photometry with open circles. Crosses represent stars with only photometry whose placement in the HRD are uncertain for one reason or another: either their transformations failed because of unrealistic colors, resulting in superfluously high effective temperatures and locations to the left of the ZAMS, or else their colors are too red to allow us to determine their reddening using $`Q`$, or the derived reddening falls outside the range we adopted on the basis of our spectroscopy. We also mark with an asterisk stars with spectral types but whose location is uncertain, such as the components of double-lined binaries. We include in these diagrams the evolutionary tracks of Schaerer et al. (1993) computed at $`z=0.008`$ (appropriate for the LMC), and the tracks of Schaller et al. (1992) at $`z=0.001`$, similar to the $`z=0.002`$ of the SMC.
We also show isochrones corresponding to ages of 2, 4, 6, 8, and 10 Myr (dashed curves), which we computed using a program kindly provided by Georges Meynet.
### 4.3 Identification of the Most Massive Stars, and the Limits of Coevality
Using the results of our calculations in the previous section, we can now identify the mass of the highest mass unevolved (H-burning) star in each association. We list the derived quantities ($`\mathrm{log}T_{\mathrm{eff}}`$, $`M_{\mathrm{bol}}`$, mass, age) for the highest mass stars in Table 5.
For associations that are strictly coeval, we expect that the stars in the HRD will follow a single isochrone, and in that case the highest mass would correspond to a “turn-off” mass and we could be confident that any evolved members of these associations were descended from stars with masses greater than this value. Alas, the HRDs of Fig. 6 do not for the most part yield such an unambiguous picture. In all cases there is some spread across isochrones. If real, such spreads would tell us that the massive stars formed over some period of time.
How significant are these age spreads? We can answer this quantitatively by considering the errors associated with the placement of stars in the HRD. Let us first consider the systematic errors. In Fig. 6(a) we show the location of the spectral type calibration data in the HRD. The huge gap among the supergiants (upper-most string of points) corresponds to the difference in the adopted effective temperature of a B5 I and a B8 I star, which is a realistic uncertainty in spectral classification. Smaller gaps likewise correspond to differences of a single spectral type. We have adopted an absolute magnitude corresponding to each type; of course, our stars, with $`M_V`$ computed from the photometry, will fall both above and below the points shown. It is instructive to see the systematic deviation of these stars from the ZAMS as one approaches cooler temperatures among the dwarfs. By log T$`{}_{\mathrm{eff}}{}^{}=4.2`$ the locations of the dwarfs are nearly coincident with the terminal main-sequence, as indicated by the first switch-back in the tracks. In this region the isochrones are tightly spaced, and a large error in the age spread would result if we compared the ages of a high mass luminosity class “V” stars with one of lower mass; for this reason we should exclude stars below 20$`_{}`$ unless they are of high visual luminosity, such as an A-type supergiant.
We note that this progression away from the ZAMS is intrinsic to the spectral type to $`\mathrm{log}T_{\mathrm{eff}}`$ calibration we’ve adopted and/or the absolute visual magnitude scale we’ve used for the purposes of this illustration. Transformations to effective temperatures on the basis of colors are usually often based on the use of spectral types as an intermediate step, rather than going directly from model atmosphere colors to effective temperatures. In these cases, the apparent presence of stars to the right of the ZAMS might be misconstrued as evidence of pre-main-sequence objects. We emphasize the need for spectroscopic followups to establish the authenticity of such discoveries.
Next, let us consider the random errors caused by misclassifying stars by a single spectral type and/or major luminosity class; i.e., calling a star an “O8 III” when in fact it is an “O9 I”. (The absolute visual magnitudes of these two subclasses overlap, and so our photometry would pose no warning.) We would overestimate the star’s luminosity by 0.1 mag simply by assuming a slightly too negative $`(BV)_o`$, which will lead to too large a value for $`A_V`$. More significant, however, is the fact that we will overestimate the star’s effective temperature by 0.05 dex, and thus overestimate the star’s bolometric correction by 0.4 mag, for a net error of 0.5 mag. The age we calculate might be 3.80 Myr (6.58 dex) if the actual age were 5.25 Myr (6.72 dex). We expect misclassification by a single spectral subtype to be common. The size of the errors we make will depend of course upon the spectral type. We show in Fig. 6(b) the errors associated with misclassification of a star by one spectral type and/or luminosity class. (We have not included in this figure the modest addition error caused by the change in reddening adopted; this will increase these errors.)
Given this discussion, we can ask the question: what fraction of stars of 20$`_{}`$ and above, and lower-mass supergiants, are in fact consistent with some median age for the association? We assume here that our error in spectral sub-typing is only 1 subtype, except for uncertain cases. We compute the youngest and oldest ages of each star associated with such a misclassification; if the cluster’s median age falls within this range, we consider that the star is coeval with the rest of the cluster. We use only the stars for which there are spectral information, as the errors in the HRD are much greater for stars with only photometry. (Compare Figures 1c and 1d of Massey et al. 1995b.) We list the fraction of stars that we find to be coeval in Table 6, along with the median ages of the clusters.
Even for the clusters that have a large percentage of stars whose ages are within 1$`\sigma `$ of the median cluster age, we might well ask the question if the ages of the highest mass stars are in accord with this value. After all, we know that in some clusters intermediate mass stars form over some period of time (several million years), with the highest mass stars forming over a shorter time, e.g., NGC 6611 (Hillenbrand et al. 1993) and R136 (Massey & Hunter 1998). We include the median age of the three highest mass stars in Table 6.
Inspection of the HRDs in Fig. 6, and of the numbers in Table 6, suggests that there is a natural division, and that some of these associations are highly coeval while the coevality of the others are more questionable. If the match between the median cluster age and the age of the 3 highest mass stars is good ($`<0.2`$ dex, comparable to the individual errors), and a large percentage of stars ($`>80\%`$) lie within 1$`\sigma `$ of the median cluster age, we consider that degree of coevality is high. Clusters that fail to meet one or the other criterion we consider the degree of coevality questionable. We consider the coevality high in 11 of our clusters, and questionable in four. We regard the other five associations as non-coeval. This could be evidence that massive stars have formed over a prolonged period, possibly with several subgroups of different ages contributing, but it may also be simply due to line-of-sight contamination within the Magellanic Clouds.
The age structure of the LH 47/48 was discussed by Oey & Massey (1995); as mentioned earlier, we restrict ourselves here to the stars on the periphery of the associated superbubble, and confirm that these stars at least form a coeval unit. LH 90 is a very interesting association located near 30 Doradus, and its age structure was explicitly discussed by Testor et al. (1993), who found “at least” two distinct age groups (3-4 Myr and 7-8 Myr). They attempted to assign membership of the evolved stars to one or the other of these populations based, not upon spatial locale, but on the basis of bolometric luminosity, which then assumes an answer about the progenitor masses a priori. They found that the $`\alpha `$ clump itself was not coeval. We have separately examined the $`\beta `$ sub-cluster using the improved data obtained by Walborn et al. (1999) and find that the same age spread apparent in the cluster as a whole is also apparent in this subclump; the $`\beta `$ cluster contains both a B0 I star of modest luminosity and a high luminosity O3 If\* star. We are, therefore, forced to abandon this very interesting region with its large number of WR stars.
We can perform one other “reasonability test” of whether the turn-off masses are relevant for the evolved objects. What is the spatial separation between the three highest mass stars (which typically define the turn-off) and the evolved objects? We computed the projected distances, and include the median of these three values in Table 7, which we discuss in the next section. (We note cases where the turnoff is actually due to the binary companion.) Here we find that the median separation is 25 pc. As this is the median, there is always some massive star nearer the evolved object than the numbers shown here. This is consistent with the notion expressed in Section 1.1 that coeval massive stars may have originated in the same place, as drifts of this order are just what we expect over 3 Myr.
We can now proceed with some confidence to assign progenitor masses to the evolved stellar content of the coeval regions.
## 5 The Progenitor Masses and BCs
### 5.1 Progenitor Masses
In Table 7 we present the main results of this investigation: what are the progenitor masses of various evolved massive stars? We enclose in parenthesis values derived from clusters whose coevality is in question, and exclude the WR stars from the associations which are non-coeval.
What can be conclude from these values? First, we find that the masses of the progenitors of WRs in the SMC are higher than those of the LMC. The data are admittedly sparse, and this conclusion rests to some extent on what mass we assign to the progenitor of AB7: the three stars with the highest mass in Hodge 53 are all components of spectroscopic binaries. We can be fairly certain that the progenitor mass of AV 332 was greater than that of its companion (i.e., $`>80_{}`$), although this supposes that binary evolution itself did not play an important role in this system.
Turning to the WRs in the LMC, we find that there is a considerable range of progenitor masses for the WNEs, with minimum masses of 30$`_{}`$ through 60$`_{}`$. If the more questionable cases were included this would increase the mass range. It appears that stars covering a range of masses pass through a WNE stage, at least at LMC metallicities.
Both of the Ofpe/WN9 stars come from associations with very low lower limits— in fact, among the lowest in our sample. There is a third Ofpe/WN9 star, one located in LH 101, which also contains evolved stars of similarly low mass (as well as higher mass evolved stars). We might conclude then that the Ofpe/WN9 stars in fact are not extremely high-mass stars at all, as their association with (other) LBVs has led others to speculate. Our conclusion that Ofpe/WN9 stars are actually “low-mass” (30$`_{}`$) in origin is not new with us: St-Louis et al. (1998) examined five LMC associations containing Ofpe/WN9 stars, including LH 89 and LH 101, and suggested much the same, although coevality was a concern for 3 of her 5 clusters. Schild (1987) had earlier studied LH 39, and also noticed the relative high age and low mass for this cluster containing an Ofpe/WN9 star. Using the WR standard atmosphere model, Crowther et al. (1995a) derive bolometric luminosities for Br 18 (R 84) and BE 381 that suggest (present) masses of 25 $`_{}`$ and 15 $`_{}`$ respectively.
Three BI + WN3 stars appear in our sample. Stars with this (composite?) type are among the brightest stars when M 33 was imaged at $`\lambda 1500`$ with the Ultraviolet Imaging Telescope (Massey et al. 1996). To our knowledge, no BI + WN3 star has ever been demonstrated to have a spectroscopic orbit. We note with some interest the relatively high minimum masses for the progenitors suggested by our study here, and we believe that only radial velocity studies can resolve the nature of these objects.
The WCs come from high mass stars, but, interestingly, not significantly higher than do the WNs. Naively this would suggest that most massive stars of mass 45-50 and above go through both a WN and a WC stage. Similarly the WC star in the SMC, AB8, has a high minimum mass ($`>70_{}`$), not too different from the WNs in the SMC.
For the LBVs in the LMC and SMC we find extremely high minimum masses—among the highest of any stars in our study. This is in accord with the prevailing notion that they are among the highest mass stars, and owe their photometric outbursts and dramatic spectral changes to instabilities inherent to high luminosity. The two B\[e\] stars in our sample have substantially different masses, in accord with the suggestion B\[e\] stars come from a large range of luminosity (Gummersbach et al. 1995).
Although the cluster turn-offs provide only lower limits to the masses of the progenitors of the evolved stars, the mass functions of these and other OB associations we’ve studied are generally well populated (cf. Massey 1995a, 1995b). Thus these cluster turn-offs should provide substantial clues to the actual masses of the progenitors.
### 5.2 The Bolometric Corrections
We next turn to computing the BCs for these evolved stars, using the observed $`M_V`$ of the star, and the $`M_{bol}`$ of the cluster turn-off stars. Previous efforts to do this (cf. Humphreys et al. 1985) relied on the fact that little change occurs in the bolometric luminosity of a massive star as it evolves, a fact simply traced to the fact that the core mass remains relatively unaffected during main-sequence evolution. Here we propose to do somewhat better, by using the evolutionary models to make a modest correction for evolution.
Smith (1968b) introduced a narrow-band photometric system to reduce the effect of WR emission lines on photometry; her “v” filter is centered at $`\lambda 5160`$ has has a zero-point tied to the system of spectrophotometric standards. For a lightly reddened star with no emission, broad-band Johnson V and Smith’s v are equivalent. ($`Vv=0.020.36\times (bv)`$ according to Conti & Smith 1972; a typical $`bv`$ value for a MC WR star is -0.1 mag, e.g. Table VI of Smith 1968b. See also Turner 1982.) We therefore use the “v” mags listed by Breysacher et al. (1999) when available to compute $`M_V`$, using the average reddenings we find for each association. We list these values in Table 7.
We can make two assumptions for computing the BCs. The first of these is to assume that the bolometric luminosity of the WR star is the same as that of the cluster turn-off. The second is to attempt to make a correction for the luminosity evolution that the models predict. The difficulty with the latter is that what the evolutionary models predict is a very sensitive function of how mass-loss is treated, and, as we emphasized earlier in this paper, the episodic shedding of mass during the LBV phase can play an appreciable role and is difficult to model. The Geneva models do not produce WR stars when standard mass-loss rates are applied except at the very highest masses, and for this reason mass-loss rates twice that of the observed values have been assumed in some of the model calculations (e.g., Meynet et al. 1994). From the end of core H-burning (similar to the stage of the highest mass stars near the cluster turn-off) to the end of the WR phase, the evolution amounts to -1.1 mag to +0.5 mag at LMC metallicities, and +0.1 mag to +0.2 mag at SMC metallicities in the sense of $`M_{\mathrm{bol}}`$ at the end of core H-burning minus $`M_{\mathrm{bol}}`$ at the end of stellar models. We include the BCs in Table 7 computed both ways, using the $`M_{\mathrm{bol}}`$ corresponding to the end of core-H burning (i.e., the terminal age main-sequence, or TAMS) and corresponding to the adopted mass of the cluster turn-off.
We see that the BCs for the WNE stars are indeed very negative, approximately $`6`$ mag, whether evolution is taken into account or not. This is in good accord with similar analysis of Galactic clusters by Humphreys et al. (1985) and Smith et al. (1994), although this is considerably more negative than that of even the earliest O-type stars ($`5`$ mag). However, recent applications of the “standard WR model” applied to “weak-lined” WNE stars by Crowther et al. (1995c) have found similar values for the BCs, giving us confidence both in our method, and providing yet another indication that the models provide a solid basis for interpreting the spectra of WR stars. There is a large range present for the BCs of WNE stars shown in Table 7, with perhaps some trend with spectral subclass; i.e., more negative with earlier type. It will be interesting to see if additional atmosphere analysis produces similar results when applied to WN2 stars.
The Ofpe/WN9 stars have far more modest BCs ($`2`$ to $`4`$ mag); analysis by Crowther et al. (1995a) of Br 18 (R 84) BE 381 using the “standard WR model” derives BCs of $`2.6`$ and $`2.7`$ mag, also in good agreement with what we find.
Turning to the WCs, we find BCs of order $`5.5`$ mag. This is a little more negative than what Humphreys et al. (1995) and Smith et al. (1994) found, although none of the WCs in their samples were as early as those studied here.
The BCs for S Dor and R 85 are very modest ($`2`$ mag). Crowther (1997) computes a similar BC for the LBV R 127, although we note that this star is another Ofpe/WN9, or was until its outburst. We have used our own photometry obtained of HD 5980 obtained in 1985 (Massey et al. 1989b) to compute its absolute visual magnitude; given the complicated nature of this (multiple) star, it is unclear what to make if its value. The bolometric luminosity of S 134 computed by Zickgraf et al. (1986) is $`10`$, in excellent agreement with the assumptions here.
## 6 Conclusions, Discussion, and Summary
Our photometric and spectroscopic investigation of 19 OB associations in the Magellanic Clouds has found that most of the massive stars have formed within a short time ($`<`$1 Myr) in about half of the regions in our sample. Their degree of coevality is similar to that found by Hillenbrand et al. (1993) for NGC 6611, i.e., that the data are consistent with all of the massive stars “having been born on a particular Tuesday.” In other regions, star-formation of the massive stars may have proceeded over a longer time, as suggested by the presence of evolved stars of 15-20$`_{}`$ (suggesting ages of 10 Myr) along with unevolved stars of high mass (60 $`_{}`$) with ages of only 2 Myr. In some cases such apparent non-coevality may be due to chance line-of-sight coincidences within the Clouds, or to drift of lower mass stars into the space occupied by a truly coeval OB association, but in other cases, such as the $`\beta `$ subcluster of LH 90, one is forced to conclude that star-formation itself was not very coeval, but proceeded over several million years.
The turn-off masses of the coeval associations have provided considerable insight into the evolution of massive stars. We find that only the highest mass stars ($`>70_{}`$) become WRs in the SMC. The numbers are admittedly sparse, and an additional complication is the fact that most SMC WRs show the presence of absorption lines. Are these absorption lines indicative of a weak stellar wind (as evidenced by the weakness of the WR emission lines) or are these all due to binary companions? Conti et al. (1989) discuss this without reaching any conclusions, and we note here that the issue of the binary frequency of the SMC WR stars requires further investigation. Possibly a strong stellar wind due to very high luminosity and binary-induced mass-loss is needed to become a WR star in the low metallicity of the SMC.
In the LMC the mass limit for becoming a WR star would appear to be a great deal lower, possibly 30$`_{}`$. Stars with a large range of initial masses (30-60 $`_{}`$), and possibly all massive stars with a mass above 30$`_{}`$ go through a WNE stage in the LMC. Most WR stars in the LMC are of early WN type; this is not true at the higher metallicity of the Milky Way, where WN3 and WN4 stars are relatively rare. This is consistent with recent theoretical work of Crowther (2000), who finds that varying only the abundance in synthetic WN models (holding all other physical parameters consist) changes the spectral subtype, with WNEs characteristic of low abundances, and WNLs characteristic of higher abundances. Thus, it may be the excitation classes are related neither to the masses nor to stellar temperatures.
The true LBVs occurs in clusters with very high turn-off masses ($`85_{}`$), both in the LMC and the SMC. This is very similar to the turn-off mass in the Trumpler 14/16 complex with which the Galactic LBV $`\eta `$ Car is associated (Massey & Johnson 1993). This supports the standard picture, that LBVs are an important, if short-lived, phase in the evolution of the most massive stars, at least at the metallicities that characterize the Magellanic Clouds and the Milky Way. We note with interest the important study by King, Gallagher, & Walterbos (2000), who find that some LBV stars in M 31 appear to be found in relative isolation, leading them to question whether these are all high mass stars, at least at the higher metallicity of M 31.
The Ofpe/WN9 stars, some of which go through some sort of outburst, cannot be “true” LBVs, if the nature of the latter is tied to extremely high bolometric luminosities. We find that the Ofpe/WN9 stars have the lowest masses of any WRs, with the progenitors possibly as low as 25$`_{}`$. Similarly, the connection of the B\[e\] stars to LBVs seems tenuous on the basis of mass or bolometric luminosities.
We know that the relative number of WC and WN stars change drastically throughout the Local Group, in a manner well-correlated with metallicity (Massey & Johnson 1998). One obvious interpretation of this is that it is much harder to lose enough mass to become a WC star in a low-metallicity environment; i.e., only the most luminous and massive stars have sufficiently high mass-loss rates to achieve this. And, similarly, the limit for WN stars should be higher in lower metallicity systems. As long as the bar is somewhat lower for achieving WN status compared to WC status, then the IMF assures that the WC/WN ratio will change. Thus our finding here that WCs and WNs come from similar mass ranges (although higher in the SMC than in the LMC), suggest that an alternative explanation is needed. Instead, it may be that it is the relative lifetimes of the WC and WN stages which are different at different masses; i.e., at very high masses the WC stage is shorter compared to the length of the WN stage than at lower masses. Or, it could be that the metallicity itself affects the relative lifetimes of the WC and WN stages. We note that we found luminous red supergiants (RSGs) cohabiting with both WNs and WCs in many OB associations in more distant galaxies of the Local Group (Massey & Johnson 1998; see for example their Figs. 14-16). While we were unable to evaluate the degree of coevality of these associations, the statistics suggest that these stars have similar progenitor mass at a given metallicity, and that variations in the relative number of RSGs to WRs are due primarily to changes in the relative lifetimes due to the effect of metallicity on the mass-loss rates (Azzopardi, Lequeux, & Maeder 1988).
We conclude that the BCs of WNE stars are quite substantial, $`6`$ mag. This value is in very good accord with that determined from weak-lined WNE stars using the WR “standard model” of Hillier (1987, 1990) by Crowther et al. (1995c). The earliest-type WN star known (of type WN2) is included in our sample, and our data suggest an even more striking BC ($`<7.5`$ mag); a full analysis of Br 4 via the standard model would be of great interest. For the Ofpe/WN9 stars we find BCs of $`2`$ to $`4`$ mag, again in good agreement with the atmospheric analysis of several such stars by Crowther et al. (1995a). We find here that the BCs of WC4 stars are typically about $`5.5`$ mag.
In the next paper, we will extend this study to the higher metallicities found in our own Milky Way galaxy.
We are grateful to Nichole King for correspondence on the issue of LBVs and their native environments, as well as useful comments on the manuscript. Deidre Hunter was also kind enough to provide a critical reading of the paper. We thank Bruce Elmegreen for correspondence and helpful preprints concerning coevality in extended regions. Comments by an anonymous referee resulted in improved discussion. Classification of some of the older spectra were done in collaboration with C. D. Garmany. Bruce Bohannan kindly allowed us to use his photographic spectrum of R 85 in this work. The participation of one of the authors (E.W.) was made possible through the Research Experiences for Undergraduate Program, which was supported by the National Science Foundation under Grant No. 9423921. P.M. acknowledges the excellent support provided by the CTIO TELOPS group.
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# Superspace Representations of SU(2,2/N) Superalgebras and Multiplet Shortening
## 1 Introduction
The study of superconformal algebras has recently attracted renewed interest for their dual role in the $`\mathrm{AdS}_{d+1}/\mathrm{CFT}_d`$ correspondence , connected to the near-horizon geometry of $`d1`$-branes.
A special role is played by $`3`$-branes since they are related to superconformal invariant quantum Yang-Mills theories. These theories are the only ones exhibiting conformal symmetry both at weak and strong coupling and, in any case, admitting, unlike other types of branes, Yang-Mills fields in the conformal regime.
The bulk and boundary operators in this correspondence are classified by highest weight UIR’s of $`SU(2,2/N)`$ algebras where $`N=1,2`$ and 4 in the known examples, since supergravity or superstring theory can admit at most 32 ($`8N`$) supersymmetries. Nevertheless, in the study of superconformal algebras and their representations different values of $`N`$ are of interest because they help one to exhibit some general features of short representations, corresponding to conformal operators with protected dimension, but more importantly, because these algebras may be relevant for some generalizations of the known schemes in which more than 32 supersymmetries may be required .
Recently it has been shown that a known generalization of ordinary superspace, called “harmonic superspace” -, is particularly suitable to build up, in a rather simple and general manner, all possible composite operators of superconformal invariant gauge theories with $`N>1`$ extended supersymmetry.
Other approaches, like ordinary superspace or the oscillator construction - of highest weight representations, although in principle possible, are much more complicated to deal with and the complete analysis of all possible shortenings would be unnecessarily difficult.
In fact, the structure of harmonic superspace is powerful enough to allow us to extend the analysis of Ref. to all $`SU(2,2/N)`$ superalgebras with arbitrary $`N`$,<sup>1</sup><sup>1</sup>1For a thorough treatment of the action of superconformal groups on harmonic superspaces see Refs. . although no dynamical theory is known for $`N>4`$. This report contains results obtained in Ref. .
From a mathematical point of view harmonic superspace is an enlarged space where superfields are defined on “flag manifolds”
$$=\frac{SU(N)}{S\left(U(n_1)\times \mathrm{}\times U(n_p)\right)},\left(\underset{k=1}{\overset{p}{}}n_k=N\right).$$
(1)
To study the general case of multiplets it is important to use the choice $`n_k=1`$ ($`k=1\mathrm{}N`$), i.e. where we quotient the group $`SU(N)`$ by its maximal torus. Then the above manifold is the largest flag manifold with complex dimension
$`N(N1)/2`$.
The ultrashort UIR’s of $`SU(2,2/N)`$ superalgebras described by analytic harmonic superfields depend only on half of the odd coordinates (Grassmann or G-analyticity):
$$W^{12\mathrm{}k}=W^{12\mathrm{}k}(\theta _{k+1},\theta _{k+2},\mathrm{},\theta _N,\overline{\theta }^1,\overline{\theta }^2,\mathrm{},\overline{\theta }^k).$$
(2)
In addition, they are annihilated by all the “step-up” generators $`E_a`$ in the Cartan decomposition of the Lie algebra of $`SU(N)`$. In other words, these superfields correspond to highest weight states of $`SU(N)`$:
$$E_a|\mathrm{HW}=0.$$
(3)
In harmonic superspace this irreducibility condition corresponds to harmonic (or H-) analyticity. The crucial point is that the $`SU(2,2/N)`$ algebra acting on such states defines a “quasi-primary” superconformal field denoted by
$$𝒟(\mathrm{},J_1,J_2;r;a_1,\mathrm{},a_{N1})$$
(4)
where $`\mathrm{},J_1,J_2`$ are the conformal dimension and spin of the state, $`r`$ is the $`U(1)`$ $`R`$ charge and $`a_1,\mathrm{},a_{N1}`$ are the $`SU(N)`$ Dynkin labels. We assign the $`R`$ charge $`r_\theta =\frac{1}{2}(1\frac{4}{N})`$ to the Grassmann coordinates in order to be consistent with the convention that chiral superfields $`\mathrm{\Phi }(\theta )`$ have $`l=r`$ for any $`N`$. This is also the charge which naturally appears in the definition of the
$`SU(2,2/N)`$ superalgebra . <sup>2</sup><sup>2</sup>2Note that for $`N=4`$, $`r_\theta =0`$ and the $`r`$ quantum number becomes a “central charge” . In this case the analysis of section 2 refers to the $`PSU(2,2/4)`$ algebra for $`r=0`$ and to the $`PU(2,2/4)`$ algebra for $`r0`$.
The G- and H-analytic superfields (2) have their lowest (scalar) component belonging to the rank $`k`$ antisymmetric representation of $`SU(N)`$ ($`k=1\mathrm{}[\frac{N}{2}]`$), have $`R`$ charge $`r_k=\frac{2k}{N}1`$ and will be shown to describe “ultrashort” representations of the $`SU(2,2/N)`$ superalgebra. If the algebra is interpreted as acting on $`\mathrm{AdS}_5`$, these are the “supersingleton” representations . For $`k=0`$ the superfield is actually “chiral” and in this case the highest weight state may carry a spin label $`(J_L,0)`$ with $`\mathrm{}=1+J_L`$. The chiral superfield is the supersingleton representation when the top spin is $`J_L=\frac{N}{2}`$. For all other analytic superfields $`(k>1)`$ the supersingleton will have top spin $`J_L=\frac{N}{2}\frac{k}{2}`$.
It should be pointed out that the same massless multiplets can be described in terms of ordinary but constrained superfields . The reason why we prefer the harmonic superspace version is the fact that the superfields (2) are unconstrained analytic objects. Analyticity is a property which is preserved by multiplication. This will allow us to tensor the above massless UIR’s in a very simple way and thus obtain series of short multiplets of $`SU(2,2/N)`$. We observe that from the $`\mathrm{AdS}_5`$ point of view, tensoring more than two supersingleton reps produces “massive bulk” reps, while tensoring only two of them produces “massless bulk” reps . The latter are the “supercurrent” multiplets discussed in Ref. .
## 2 Unitarity bounds and shortening of UIR’s of $`SU(2,2/N)`$
The unitarity bounds of highest weight UIR’s of $`SU(2,2/N)`$ have been derived in Refs. . They correspond to some bounds on the highest weight state (4). Let us define the quantities
$$m_1=\underset{k=1}{\overset{N1}{}}a_k,m=\underset{k=1}{\overset{N1}{}}(Nk)a_k$$
(5)
and
$`X(J,r,{\displaystyle \frac{2m}{N}})`$ $`=`$ $`2+2Jr+{\displaystyle \frac{2m}{N}},`$
$`Y(r,{\displaystyle \frac{2m}{N}})`$ $`=`$ $`r+{\displaystyle \frac{2m}{N}}.`$ (6)
Then we have $`(J_1=J_L,J_2=J_R)`$:
$$\mathrm{A})\mathrm{}X(J_2,r,\frac{2m}{N})X(J_1,r,2m_1\frac{2m}{N})$$
(7)
(or $`J_1J_2`$, $`rr`$, $`\frac{2m}{N}2m_1\frac{2m}{N}`$);
$$\mathrm{B})\mathrm{}=Y(r,\frac{2m}{N})X(J_1,r,2m_1\frac{2m}{N})$$
(8)
(or $`J_1J_2`$, $`rr`$, $`\frac{2m}{N}2m_1\frac{2m}{N}`$);
$$\mathrm{C})\mathrm{}=m_1,r=\frac{2m}{N}m_1,J_1=J_2=0.$$
(9)
The massless UIR’s correspond to B) for $`a_k=0`$, $`\mathrm{}=r=1+J_L`$ and to C) for $`\mathrm{}=m_1=1`$, $`r_k=\frac{2k}{N}1`$, $`1k[\frac{N}{2}]`$. Note that the two series overlap for $`J_L=0`$ in B) and $`k=0`$ in C).
The short multiplets that we shall build in section 4 by tensoring massless multiplets from the C) series in the case of $`N=2n`$ for $`k=n`$ ($`r=0`$) will belong to the shortenings in B) and C) obtained for $`J_1=0`$ and $`r=0`$:
$`\mathrm{B})`$ $`\mathrm{}={\displaystyle \frac{2m}{N}},{\displaystyle \frac{2m}{N}}m_11`$ (10)
$`\mathrm{C})`$ $`\mathrm{}=m_1,{\displaystyle \frac{2m}{N}}=m_1`$
## 3 Massless superconformal multiplets
### 3.1 Grassmann analytic superfields
We consider superfields
$$W^{i_1\mathrm{}i_k}(x_{\alpha \dot{\alpha }},\theta _i^\alpha ,\overline{\theta }^{\dot{\alpha }i})$$
with $`k=1,\mathrm{},n`$ (where $`n=[\frac{N}{2}]`$) totally antisymmetrized indices in the fundamental representation of $`SU(N)`$. These superfields satisfy the following constraints:
$`D_\alpha ^{(j}W^{i_1)i_2\mathrm{}i_k}=0,`$ (11)
$`\overline{D}_{\dot{\alpha }\{j}W^{i_1\}i_2\mathrm{}i_k}=0`$ (12)
where $`()`$ means symmetrization and $`\{\}`$ means the traceless part. The spinor derivatives algebra is
$$\{D_\alpha ^i,\overline{D}_{\dot{\alpha }j}\}=i\delta _j^i_{\alpha \dot{\alpha }}$$
(13)
with $`_{\alpha \dot{\alpha }}=\sigma _{\alpha \dot{\alpha }}^\mu _\mu `$. In the cases $`N=2,3,4`$ these constraints define the on-shell $`N=2`$ matter (hyper)multiplet and the $`N=3,4`$ on-shell super-Yang-Mills multiplets . Their generalization to arbitrary $`N`$ has been given in Refs. where it has also been shown that they describe on-shell massless multiplets.
Our aim in this section is to rewrite the constraints (11), (12) in harmonic superspace where they will take the simple form of analyticity conditions. Using this fact we will then be able to construct tensor products of the corresponding multiplets in a very straightforward and easy way (section 4).
The main purpose of introducing harmonics is to be able to covariantly project all the $`SU(N)`$ indices in (11), (12) onto a set of $`U(1)`$ charges. To this end we choose the harmonic coset $`SU(N)/(U(1))^{N1}`$ described in terms of harmonic variables $`u_i^I`$ and their conjugates $`u_I^i=(u_i^I)^{}`$. <sup>3</sup><sup>3</sup>3The harmonic notation used here differs from the original one of Refs. . It is similar to the one introduced in Ref. for the case $`N=3`$ and in Refs. for general $`N`$. They form an $`SU(N)`$ matrix where $`i`$ is an index in the fundamental representation of $`SU(N)`$ and $`I=1,\mathrm{},N`$ is a collection of the $`N1`$ $`U(1)`$ charges corresponding to the projections of the second index (the harmonic $`u_I^i`$ carries charges opposite to those of $`u_i^I`$). They satisfy the following $`SU(N)`$ defining conditions:
$`u_i^Iu_J^i=\delta _J^I,`$ (14)
$`uSU(N):`$ $`u_i^Iu_I^j=\delta _i^j,`$ (16)
$`\epsilon ^{i_1\mathrm{}i_N}u_{i_1}^1\mathrm{}u_{i_N}^N=1.`$
Now, let us use these harmonic variables to split all the $`SU(N)`$ indices in the constraints (11), (12) into independent $`(U(1))^{N1}`$ projections. For example, the projection
$$W^{12\mathrm{}k}=W^{i_1i_2\mathrm{}i_k}u_{i_1}^1u_{i_2}^2\mathrm{}u_{i_k}^k$$
(17)
satisfies the constraints
$`D_\alpha ^1W^{12\mathrm{}k}=D_\alpha ^2W^{12\mathrm{}k}=\mathrm{}=`$
$`D_\alpha ^kW^{12\mathrm{}k}=0,`$
(18)
$`\overline{D}_{\dot{\alpha }k+1}W^{12\mathrm{}k}=\overline{D}_{\dot{\alpha }k+2}W^{12\mathrm{}k}=\mathrm{}=`$
$`\overline{D}_{\dot{\alpha }N}W^{12\mathrm{}k}=0`$ (19)
where $`D_\alpha ^I=D_\alpha ^iu_i^I`$ and $`\overline{D}_{\dot{\alpha }I}=\overline{D}_{\dot{\alpha }i}u_I^i`$. The first of them, eq. (18), is a corollary of the commuting nature of the harmonics variables, and the second one, eq. (LABEL:4), of the unitarity condition (14). The main achievement in rewriting the constraints (11), (12) in this new form is that they can be explicitly solved by going to an appropriate G-analytic basis in superspace:
$`x_A^{\alpha \dot{\alpha }}=x^{\alpha \dot{\alpha }}+i(\theta _1^\alpha \overline{\theta }^{1\dot{\alpha }}+\mathrm{}+\theta _k^\alpha \overline{\theta }^{k\dot{\alpha }}`$
$`\theta _{k+1}^\alpha \overline{\theta }^{k+1\dot{\alpha }}\mathrm{}\theta _N^\alpha \overline{\theta }^{N\dot{\alpha }}),`$
$`\theta _I^\alpha =\theta _i^\alpha u_I^i,\overline{\theta }^{\dot{\alpha }I}=\overline{\theta }^{\dot{\alpha }i}u_i^I.`$ (20)
In this basis $`W^{12\mathrm{}k}`$ becomes an unconstrained function of $`k`$ $`\overline{\theta }`$’s and $`Nk`$ $`\theta `$’s:
$$W^{12\mathrm{}k}=W^{12\mathrm{}k}(x_A,\theta _{k+1},\mathrm{},\theta _N,\overline{\theta }^1,\mathrm{},\overline{\theta }^k,u).$$
(21)
Altogether it depends on half the number of the odd variables of $`N`$-extended superspace and for this reason we call it Grassmann (or G-) analytic. We recall that the notion of Grassmann analyticity was first introduced in Ref. , still in the context of ordinary superspace. In $`N=2`$ harmonic superspace this notion became $`SU(2)`$ covariant. The generalization to $`N=3`$ was given in Ref. and later on to general $`N`$ in Refs. (under the name of “$`(N,p,q)`$ superspace”).
The massless conformal multiplets describe the ordinary massless UIR’s of the super Poincaré group obtained earlier by the Wigner method of induced representations (see, for instance, Ref. ). The self-conjugate $`N=8`$ multiplet was obtained by the oscillator method in Ref. .
### 3.2 Harmonic analyticity as $`SU(N)`$ <br>irreducibility
It is important to realize that a G-analytic superfield is an $`SU(N)`$ covariant object only because it depends on the harmonic variables. In order to recover the original harmonic-independent but constrained superfield $`W^{i_1i_2\mathrm{}i_k}(x,\theta ,\overline{\theta })`$ (11), (12) we need to impose differential conditions involving the harmonic variables. The harmonic derivatives are made out of the operators
$$_J^I=u_i^I\frac{}{u_i^J}u_J^i\frac{}{u_I^i}$$
(22)
which respect the defining relations (14), (16). These derivatives act on the harmonics as follows:
$$_J^Iu_i^K=\delta _J^Ku_i^I,_J^Iu_K^i=\delta _K^Iu_J^i.$$
(23)
The diagonal ones $`_I^I`$ count the $`U(1)`$ charges,
$$_I^Iu_i^I=u_i^I,_I^Iu_I^i=u_I^i.$$
(24)
The relation (16) implies that the charge operators $`_I^I`$ are not independent,
$$\underset{I=1}{\overset{N}{}}_I^I=0$$
(25)
(this reflects the fact that we are considering $`SU(N)`$ and not $`U(N)`$).
A basic assumption in our approach to the harmonic coset $`SU(N)/U(1)^{N1}`$ is that any harmonic function is homogeneous under the action of $`U(1)^{N1}`$, i.e., it is an eigenfunction of the charge operators $`_I^I`$,
$`_I^If_{L_1\mathrm{}L_r}^{K_1\mathrm{}K_q}(u)=(\delta _I^{K_1}+\mathrm{}+\delta _I^{K_q}\delta _{L_1}^I\mathrm{}`$
$`\delta _{L_r}^I)f^{K_1\mathrm{}K_q}_{L_1\mathrm{}L_r}(u)`$ (26)
(note that the charges $`K_1\mathrm{}K_q;L_1\mathrm{}L_r`$ are not necessarily all different). Thus it effectively depends on the $`(N^21)(N1)=N(N1)`$ real coordinates of the coset $`SU(N)/U(1)^{N1}`$. Then the actual harmonic derivatives on the coset are the $`N(N1)/2`$ complex derivatives $`_J^I`$, $`I<J`$ (or their conjugates $`_J^I`$, $`I>J`$).
The set of $`N^21`$ derivatives $`_J^I`$ (taking into account the linear dependence (25)) form the algebra of $`SU(N)`$:
$$[_J^I,_L^K]=\delta _J^K_L^I\delta _L^I_J^K.$$
(27)
The Cartan decomposition of this algebra $`L^++L^0+L^{}`$ is given by the sets
$`L^+`$ $`=`$ $`\{_J^I,I<J\},L^0=\{_I^I,{\displaystyle \underset{I=1}{\overset{N}{}}}_I^I=0\},`$
$`L^{}`$ $`=`$ $`\{_J^I,I>J\}.`$ (28)
It becomes clear that imposing the harmonic conditions
$$_J^If_{L_1\mathrm{}L_r}^{K_1\mathrm{}K_q}(u)=0,I<J$$
(29)
on a harmonic function with a given set of charges $`K_1\mathrm{}K_q;L_1\mathrm{}L_r`$ defines the highest weight of an $`SU(N)`$ irrep. In other words, the harmonic expansion of such a function contains only one irrep which is determined by the combination of charges $`K_1\mathrm{}K_q;L_1\mathrm{}L_r`$. In fact, not all of the derivatives $`_J^I,I<J`$ are independent, as follows from the algebra (27). The independent set consists of the $`N1`$ derivatives
$$_2^{\mathrm{\hspace{0.17em}1}},_3^{\mathrm{\hspace{0.17em}2}},\mathrm{},_N^{N1}$$
(30)
corresponding to the simple roots of $`SU(N)`$. Then the $`SU(N)`$ defining constraint (29) is equivalent to
$$(_2^{\mathrm{\hspace{0.17em}1}},_3^{\mathrm{\hspace{0.17em}2}},\mathrm{},_N^{N1})f_{L_1\mathrm{}L_r}^{K_1\mathrm{}K_q}(u)=0.$$
(31)
The coset $`SU(N)/U(1)^{N1}`$ can be parametrized by $`N(N1)/2`$ complex coordinates. In this case the constraints (29) take the form of covariant (in the sense of Cartan) Cauchy-Riemann analyticity conditions. For this reason we call the set of constraints (29) (or the equivalent set (31)) harmonic (H-)analyticity conditions. The above argument shows that H-analyticity is equivalent to defining a highest weight of $`SU(N)`$, i.e. it is the $`SU(N)`$ irreducibility condition on the harmonic functions.
As an example, take $`N=2`$ and the function $`f^1(u)`$ subject to the constraint
$$_2^{\mathrm{\hspace{0.17em}1}}f^1(u)=0f^1(u)=f^iu_i^1.$$
(32)
So, the harmonic function is reduced to a doublet of $`SU(2)`$. Similarly, for $`N=4`$ the function $`f^{12}(u)`$ is reduced to the $`\underset{¯}{6}`$ of $`SU(4)`$. Indeed, the constraints $`_3^{\mathrm{\hspace{0.17em}2}}f^{12}(u)=_4^{\mathrm{\hspace{0.17em}3}}f^{12}(u)=0`$ ensure that $`f^{12}(u)`$ depends on $`u^1,u^2`$ only, $`f^{12}(u)=f^{ij}u_i^1u_j^2`$. Then the constraint $`_2^{\mathrm{\hspace{0.17em}1}}f^{12}(u)=f^{ij}u_i^1u_j^1`$ = 0 implies $`f^{ij}=f^{ji}`$.
In the G-analytic basis (20) the harmonic derivatives become covariant $`D_J^I`$. In particular, the derivatives
$`D_J^I`$ $`=`$ $`_J^Ii\theta _J^\alpha \overline{\theta }^{I\dot{\alpha }}_{\alpha \dot{\alpha }}\theta _J^I+\overline{\theta }^I\overline{}_J,`$
$`I`$ $`=`$ $`1,\mathrm{},k,J=k+1,\mathrm{},N`$ (33)
acquire space-time derivative terms. The $`SU(N)`$ commutation relations among the $`D_J^I`$ are not affected by the change of basis. The same is true for the commutation relations of the $`D_J^I`$ with the spinor derivatives:
$$[D_J^I,D_\alpha ^K]=\delta _J^KD_\alpha ^I,[D_J^I,\overline{D}_{\dot{\alpha }K}]=\delta _K^I\overline{D}_{\dot{\alpha }J}.$$
(34)
Using these relations one can see that the H-analyticity conditions
$$D_J^IW^{12\mathrm{}k}=0,I<J$$
(35)
or the equivalent set
$$(D_2^{\mathrm{\hspace{0.17em}1}},D_3^{\mathrm{\hspace{0.17em}2}},\mathrm{},D_N^{N1})W^{12\mathrm{}k}=0$$
(36)
are compatible with the G-analyticity ones (LABEL:4).
### 3.3 Analyticity and massless multiplets: <br>“Singletons”
The constraints of H-analyticity (35) combined with those of G-analyticity (LABEL:4) have important implications for the components of the superfield. First of all, they make each component an irrep of $`SU(N)`$. Take, for example, the first component
$$\varphi ^{12\mathrm{}k}(x,u)=W^{12\mathrm{}k}|_0$$
(37)
where $`|_0`$ means $`\theta =\overline{\theta }=0`$. The constraints $`_{I+1}^I\varphi ^{12\mathrm{}k}(x,u)=0`$, $`I=k,\mathrm{},N`$ imply that $`\varphi ^{12\mathrm{}k}(x,u)`$ takes the form
$$\varphi ^{12\mathrm{}k}(x,u)=\varphi ^{i_1i_2\mathrm{}i_k}(x)u_{i_1}^1u_{i_2}^2\mathrm{}u_{i_k}^k.$$
This is a rank $`k`$ tensor without any symmetry, i.e. a reducible representation of $`SU(N)`$. Further, the constraint, e.g.,
$$_2^1\varphi ^{123\mathrm{}k}(x,u)=\varphi ^{113\mathrm{}k}(x,u)=$$
$$\varphi ^{i_1i_2i_3\mathrm{}i_k}u_{i_1}^1u_{i_2}^1u_{i_3}^3\mathrm{}u_{i_k}^n=0$$
removes the symmetric part in the first two indices. Similarly, the remaining constraints (35) remove all the symmetrizations and we find the totally antisymmetric rank $`k`$ irrep of $`SU(N)`$.
Another example are the spinor components
$`\chi _\alpha ^{12\mathrm{}kk+1}(x,u)=D_\alpha ^{k+1}W^{12\mathrm{}k}|_0,`$
$`\overline{\psi }_{\dot{\alpha }}^{23\mathrm{}k}(x,u)=\overline{D}_{1\dot{\alpha }}W^{12\mathrm{}k}|_0.`$ (38)
The same harmonic argument shows that these are harmonic projections of the totally antisymmetric components $`\chi _\alpha ^{[i_1i_2\mathrm{}i_{k+1}]}(x)`$ and $`\overline{\psi }_{\dot{\alpha }}^{[i_2i_3\mathrm{}i_k]}(x)`$.
Further important constraints occur at the level of 2 or more $`\theta `$’s:
$`D^{I\alpha }D_\alpha ^JW^{12\mathrm{}k}`$ $`=`$ $`0,I,J=k+1,\mathrm{},N,`$
$`\overline{D}_{I\dot{\alpha }}\overline{D}_J^{\dot{\alpha }}W^{12\mathrm{}k}`$ $`=`$ $`0,I,J=1,\mathrm{},k.`$ (40)
The easiest way to see this is to hit the defining constraint (11) with $`D^{k\alpha }`$ and then project with harmonics.
The constraints (LABEL:lin2), (40) imply that the components of the type
$`\chi _{(\alpha _1\mathrm{}\alpha _p)}^{1\mathrm{}k+p}=D_{\alpha _1}^{k+1}\mathrm{}D_{\alpha _p}^{k+p}W^{12\mathrm{}k}|_0,pNk`$
(41)
$`\overline{\psi }_{(\dot{\alpha }_1\mathrm{}\dot{\alpha }_p)}^{p+1\mathrm{}k}=\overline{D}_{1\dot{\alpha }_1}\mathrm{}\overline{D}_{p\dot{\alpha }_p}W^{12\mathrm{}k}|_0,pk`$
(42)
are totally symmetric in their spinor indices, i.e. they carry spin $`(p/2,0)`$ or $`(0,p/2)`$, correspondingly. Among them one finds the
$`\text{top spin }(\frac{N}{2}\frac{k}{2},0)\text{:}\chi _{(\alpha _1\mathrm{}\alpha _{Nk})}=`$
$`D_{\alpha _1}^{k+1}\mathrm{}D_{\alpha _{Nk}}^NW^{12\mathrm{}k}|_0`$ (43)
which is also an $`SU(N)`$ singlet. Note that in the case $`N=2n`$, $`k=n`$ the top spin occurs both as $`(n/2,0)`$ and $`(0,n/2)`$ (we call this a “self-conjugate” multiplet). Moreover, if $`N=4n`$ and $`k=2n`$ one can impose a reality condition on the superfield $`W^{12\mathrm{}2n}`$ which implies, in particular, that
$$\chi _{(\alpha _1\mathrm{}\alpha _{2n})}=(\psi _{(\dot{\alpha }_1\mathrm{}\dot{\alpha }_{2n})})^{}.$$
(44)
Next, one can show that all the components of the type (41), (42) satisfy massless field equations. Indeed, from the constraint (LABEL:lin2) and from G-analyticity it follows that
$`0`$ $`=`$ $`\overline{D}_{k+1\dot{\beta }}D^{k+1\alpha _1}D_{\alpha _1}^{k+1}\mathrm{}D_{\alpha _p}^{k+p}W^{12\mathrm{}k}`$ (45)
$`=`$ $`2i_{\dot{\beta }}^{\alpha _1}D_{\alpha _1}^{k+1}\mathrm{}D_{\alpha _p}^{k+p}W^{12\mathrm{}k}`$
$``$ $`_{\dot{\beta }}^{\alpha _1}\chi _{(\alpha _1\mathrm{}\alpha _p)}^{1\mathrm{}k+p}=0`$
and similarly for $`\overline{\psi }_{(\dot{\alpha }_1\mathrm{}\dot{\alpha }_p)}^{p+1\mathrm{}k}`$. The leading scalar component (37) satisfies the d’Alembert equation:
$`0`$ $`=`$ $`(D^1)^2(\overline{D}_1)^2W^{12\mathrm{}k}=4\mathrm{}W^{12\mathrm{}k}`$ (46)
$``$ $`\mathrm{}\varphi ^{12\mathrm{}k}=0.`$
Finally, all the components of mixed type,
$`f_{\dot{\alpha }_1\mathrm{}\dot{\alpha }_p\alpha _1\mathrm{}\alpha _q}^{p+1\mathrm{}k+q}`$
$`=\overline{D}_{1\dot{\alpha }_1}\mathrm{}\overline{D}_{p\dot{\alpha }_p}D_{\alpha _1}^{k+1}\mathrm{}D_{\alpha _q}^{k+q}W^{12\mathrm{}k}|_0,`$
$`pk,qNk`$ (47)
are expressed in terms of the space-time derivatives of lower components. Indeed,
$`D_{k+q}^1f_{\dot{\alpha }_1\mathrm{}\dot{\alpha }_p\alpha _1\mathrm{}\alpha _q}^{p+1\mathrm{}k+q}`$
$`=\overline{D}_{k+q\dot{\alpha }_1}\overline{D}_{2\dot{\alpha }_2}\mathrm{}D_{\alpha _q}^{k+q}W^{12\mathrm{}k}|_0`$
$`=(1)^{p+q1}i_{\dot{\alpha }_1\alpha _q}\overline{D}_{2\dot{\alpha }_2}\mathrm{}D_{\alpha _{q1}}^{k+q1}W^{12\mathrm{}k}|_0`$
$`f_{\dot{\alpha }_1\mathrm{}\dot{\alpha }_p\alpha _1\mathrm{}\alpha _q}^{p+1\mathrm{}k+q\mathrm{1\hspace{0.33em}1}}`$
$`=(1)^{p+q1}i_{\dot{\alpha }_1\alpha _q}g_{\dot{\alpha }_2\mathrm{}\dot{\alpha }_p\alpha _1\mathrm{}\alpha _{q1}}^{1p+1\mathrm{}k+q1}`$ (48)
To summarize, the superfield $`W^{12\mathrm{}k}`$ subject to the constraints of G- and H-analyticity has the following component content (the derivative terms are not shown):
$`W^{12\mathrm{}k}=\varphi ^{12\mathrm{}k}`$
$`+\overline{\theta }_{\dot{\alpha }}^1\overline{\psi }^{\dot{\alpha }\mathrm{\hspace{0.33em}23}\mathrm{}k}+\mathrm{}+\overline{\theta }_{\dot{\alpha }}^k\overline{\psi }^{\dot{\alpha }\mathrm{\hspace{0.33em}12}\mathrm{}k1}`$
$`+\theta _{k+1}^\alpha \chi _\alpha ^{1\mathrm{}kk+1}+\mathrm{}+\theta _N^\alpha \chi _\alpha ^{1\mathrm{}kN}`$
$`+\overline{\theta }_{\dot{\alpha }}^1\overline{\theta }_{\dot{\beta }}^2\overline{\psi }^{(\dot{\alpha }\dot{\beta })\mathrm{\hspace{0.33em}3}\mathrm{}k}+\mathrm{}+\overline{\theta }_{\dot{\alpha }}^{k1}\overline{\theta }_{\dot{\beta }}^k\overline{\psi }^{(\dot{\alpha }\dot{\beta })\mathrm{\hspace{0.33em}1}\mathrm{}k2}`$
$`+\theta _{k+1}^\alpha \theta _{k+2}^\beta \chi _{(\alpha \beta )}^{1\mathrm{}kk+1k+2}+`$
$`\mathrm{}+\theta _{N1}^\alpha \theta _N^\beta \chi _{(\alpha \beta )}^{1\mathrm{}kN1N}\mathrm{}`$
$`+\overline{\theta }_{\dot{\alpha }_1}^1\mathrm{}\overline{\theta }_{\dot{\alpha }_k}^k\overline{\psi }^{(\dot{\alpha }_1\mathrm{}\dot{\alpha }_k)}`$
$`+\theta _{k+1}^{\alpha _1}\mathrm{}\theta _N^{\alpha _{Nk}}\chi _{(\alpha _1\mathrm{}\alpha _{Nk})}`$ (49)
where all the fields belong to totally antisymmetric irreps of $`SU(N)`$ and satisfy the massless field equations
$`\mathrm{}\varphi ^{[i_1\mathrm{}i_k]}=0,`$
$`^{\beta \dot{\alpha }_1}\overline{\psi }_{(\dot{\alpha }_1\mathrm{}\dot{\alpha }_p)}^{[i_1\mathrm{}i_{kp}]}=0,1pk`$ (50)
$`^{\alpha _1\dot{\beta }}\chi _{(\alpha _1\mathrm{}\alpha _p)}^{[i_1\mathrm{}i_p]}=0,1pNk`$
This is the content of an $`N`$-extended superconformal multiplet of the C) series of section 2. It is characterized by the $`SU(N)`$ irrep of the first component (described by the Young tableau $`m_1=\mathrm{}=m_k=1,m_{k+1}=\mathrm{}=m_{N1}=0`$), by its $`R`$ charge
$$r_k=\frac{2k}{N}1$$
(51)
and conformal dimension $`\mathrm{}=1`$ and by the top spin $`J_{\mathrm{top}}=(\frac{N}{2}\frac{k}{2},0)`$.
### 3.4 Chiral superfields
The G-analytic superfields considered above contain at least one $`\overline{\theta }`$. The case of “extreme” G-analyticity will be the absence of any $`\overline{\theta }`$’s. These are the well-known chiral superfields satisfying the constraint
$$\overline{D}_{i\dot{\alpha }}W=0W=W(x_L^{\alpha \dot{\alpha }},\theta _i^\alpha )$$
(52)
where
$$x_L^{\alpha \dot{\alpha }}=x^{\alpha \dot{\alpha }}i\theta _i^\alpha \overline{\theta }^{i\dot{\alpha }}.$$
(53)
Note that in this case we do not need harmonic variables, since G-analyticity involves a subset of odd coordinates forming an entire irrep of $`SU(N)`$, and not a set of $`U(1)`$ projections. Consequently, in order to put such a superfield on shell, we cannot use H-analyticity but need to impose a new type of constraint:
$$D^{\alpha i}D_\alpha ^jW=0.$$
(54)
The resulting components are multispinors of the same chirality (cf. eq. (49)):
$`W`$ $`=`$ $`\varphi +\theta _i^\alpha \chi _\alpha ^i`$
$`+`$ $`\mathrm{}+\theta _{i_1}^{\alpha _1}\mathrm{}\theta _{i_n}^{\alpha _n}\chi _{(\alpha _1\mathrm{}\alpha _n)}^{[i_1\mathrm{}i_n]}+\mathrm{}+(\theta )^{2N}\chi `$
satisfying massless field equations. The tops spin is $`(\frac{N}{2},0)`$.
The chiral superfields above are scalar, but there exist conformally covariant chiral superfields with an arbitrary $`(J_L,0)`$ index of the highest weight: $`W_{\alpha _1\mathrm{}\alpha _{2J_L}}`$. In this case the masslessness condition is $`D^{\alpha _1i}W_{\alpha _1\mathrm{}\alpha _{2J_L}}=0`$.
## 4 Short superconformal multiplets: bulk “ massless” and “massive” states
In this section we shall concentrate on the case $`N=2n`$ for reasons of simplicity. The analytic superfield $`W^{12\mathrm{}n}(\theta _{n+1},\mathrm{},\theta _{2n},\overline{\theta }^1,\mathrm{},\overline{\theta }^n)`$ describes a superconformal multiplet characterized by the Young tableau $`m_1=\mathrm{}=m_n=1,`$ $`m_{n+1}=\mathrm{}=m_{2n1}=0`$ of its first component (a Lorentz scalar), by its dimension $`\mathrm{}=1`$ and $`R`$ charge $`r=0`$ (see (51)). Now we shall use this multiplet as a building block for constructing other “short” superconformal multiplets.
The building block $`W^{12\mathrm{}n}`$ can be equivalently rewritten by choosing different harmonic projections of its $`SU(N)`$ indices and, consequently, different sets of G-analyticity constraints. This amounts to superfields of the type
$$W^{I_1I_2\mathrm{}I_n}(\theta _{J_1},\mathrm{},\theta _{J_n},\overline{\theta }^{I_1},\mathrm{},\overline{\theta }^{I_n})$$
(56)
where $`I_1,\mathrm{},I_n`$ and $`J_1,\mathrm{},J_n`$ are two complementary sets of $`n`$ indices. Each of these superfields depends on $`2N=4n`$ Grassmann variables, i.e. half of the total number of $`4N=8n`$. This is the minimal size of a G-analytic superspace, so we can say that the $`W`$’s are the “shortest” superfields (superconformal multiplets). Another characteristic of these $`W`$’s is the absence of $`R`$ charges.
The idea now is to start multiplying different species of the $`W`$’s of the type (56) in order to obtain composite objects depending on various numbers of odd variables. The sets $`I_1,\mathrm{},I_n`$ can be chosen in $`(2n)!/(n!)^2`$ different ways. However, we do not need consider all of them. The following choice of $`W`$’s and of the order of multiplication covers all possible intermediate types of G-analyticity:
$$A(p_1,p_2,\mathrm{},p_{2n1})=$$
$`[W^{1\mathrm{}n}(\theta _{n+1\mathrm{}2n}\overline{\theta }^{1\mathrm{}n})]^{p_1+\mathrm{}+p_{2n1}}`$
$`\times [W^{1\mathrm{}n1n+1}(\theta _{\underset{¯}{n}n+2\mathrm{}2n}`$
$`\times \overline{\theta }^{1\mathrm{}n1\underset{¯}{n+1}})]^{p_2+\mathrm{}+p_{2n1}}`$
$`\times [W^{1\mathrm{}n1n+2}(\theta _{nn+1n+3\mathrm{}2n}`$
$`\times \overline{\theta }^{1\mathrm{}n1\underset{¯}{n+2}})]^{p_3+\mathrm{}+p_{2n1}}`$
$`\mathrm{}`$
$`\times [W^{1\mathrm{}n\mathrm{1\hspace{0.33em}2}n1}(\theta _{n\mathrm{}2n\mathrm{2\hspace{0.33em}2}n}`$
$`\times \overline{\theta }^{1\mathrm{}n1\underset{¯}{2n1}})]^{p_n+\mathrm{}+p_{2n1}}`$
$`\times [W^{1\mathrm{}n2nn+1}(\theta _{\underset{¯}{n1}n+2\mathrm{}2n}`$
$`\times \overline{\theta }^{1\mathrm{}n2nn+1})]^{p_{n+1}+\mathrm{}+p_{2n1}}`$
$`\times [W^{1\mathrm{}n3n1nn+1}(\theta _{\underset{¯}{n2}n+2\mathrm{}2n}`$
$`\times \overline{\theta }^{1\mathrm{}n3n1nn+1})]^{p_{n+2}+\mathrm{}+p_{2n1}}`$
$`\mathrm{}`$
$`\times [W^{13\mathrm{}n+1}(\theta _{\underset{¯}{2}n+2\mathrm{}2n}\overline{\theta }^{13\mathrm{}n+1})]^{p_{2n2}+p_{2n1}}`$
$`\times [W^{23\mathrm{}n+1}(\theta _{\underset{¯}{1}n+2\mathrm{}2n}\overline{\theta }^{23\mathrm{}n+1})]^{p_{2n1}}.`$ (57)
The power $`_{r=k}^{2n1}p_r`$ of the $`k`$-th $`W`$ is chosen in such a way that each new $`p_r`$ corresponds to bringing in a new type of $`W`$. As a result, at each step a new $`\theta `$ or $`\overline{\theta }`$ appears (they are underlined in (57)), thus adding new odd dimensions to the G-analytic superspace. The only exception of this rule is the second step at which both a new $`\theta `$ and a new $`\overline{\theta }`$ appear. So, the series (57) covers all possible subspaces with $`4n,4n+4,4n+6,\mathrm{},8n2`$ odd coordinates (notice once again the missing subspace with $`4n+2`$ odd coordinates). In this sense we can say that the G-analytic superfield $`A(p_1,p_2,\mathrm{},p_{2n1})`$ realizes a “short” superconformal multiplet.
The superfield $`A(p_1,p_2,\mathrm{},p_{2n1})`$ should be submitted to the same H-analyticity constraints as one would impose on $`W^{1\mathrm{}n}`$ alone,
$`D_{I+1}^IA(p_1,p_2,\mathrm{},p_{2n1})=0,`$
$`I=1,2,\mathrm{},2n1.`$ (58)
This is clearly compatible with G-analyticity since the conditions on a generic $`A(p_1,p_2,\mathrm{},p_{2n1})`$ form a subset of these on $`W^{1\mathrm{}n}`$. As before, H-analyticity makes $`A(p_1,p_2,\mathrm{},p_{2n1})`$ irreducible under $`SU(N)`$. Here is the structure of Young tableau which corresponds to the first (scalar) component of this superfield (and characterizes the supermultiplet as a whole):
The top row is filled with indices projected with $`u_i^1`$ (hence the symmetrization among them), the second row \- with $`u_i^2`$, etc. The harmonic conditions (4) remove all the symmetrizations among indices belonging to different projections (rows). By counting the number of occurrences of the projection $`1`$ in (57), we easily find the relation
$$m_1=\mathrm{}p_{2n1}$$
(59)
where $`\mathrm{}`$ is the total number of $`W`$’s (equal to the dimension of the superfield $`A`$, since $`\mathrm{}_W=1`$). Another simple counting shows the relation
$$\underset{k=1}{\overset{2n1}{}}m_k=n\mathrm{}=\frac{N}{2}\mathrm{}.$$
(60)
If the last $`W`$ in (4) is not present there is an additional relation among the Young tableau labels:
$$p_{2n1}=0m_1=\frac{2}{N}\underset{k=1}{\overset{2n1}{}}m_k.$$
(61)
Finally, introducing the Dynkin labels $`[a_1,\mathrm{},`$ $`a_{2n1}]`$ where $`a_1=m_{2n1}`$ and $`a_k=m_{2nk+1}m_{2nk}`$ for $`k2`$, we find
$`a_1={\displaystyle \underset{k=n}{\overset{2n1}{}}}p_k,`$
$`a_2=p_{n1},\mathrm{},a_{n2}=p_3,`$
$`a_{n1}=p_2+{\displaystyle \underset{k=n+1}{\overset{2n1}{}}}(kn)p_k,`$
$`a_n=p_1,`$ (62)
$`a_{n+1}=(n2){\displaystyle \underset{k=n+1}{\overset{2n1}{}}}p_k+{\displaystyle \underset{k=2}{\overset{n}{}}}(k1)p_k,`$
$`a_{n+2}=p_{n+1},\mathrm{},a_{2n1}=p_{2n2}.`$
## 5 Conclusion
In this paper we studied representations of four-dimensional superconformal algebras with an arbitrary number of supersymmetries.
This analysis also provides the classification of short multiplets of superalgebras on AdS<sub>5</sub> and in particular “massless” and “massive” fields in anti-de Sitter geometries, in terms of boundary “composite” operatprs
## Acknowledgements
This work has been supported in part by the European Commission TMR programme ERBFMRX-CT96-0045 (Laboratori Nazionali di Frascati, INFN) and by DOE grant DE-FG03-91ER40662, Task C.
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# Wetting-induced effective interaction potential between spherical particles
## I Introduction
In view of understanding a particular phenomenon in condensed matter, theory is supposed to identify the corresponding relevant degrees of freedom and to provide the effective interaction between them by, approximately, integrating out the remaining ones so that one is left with a manageable model. It is a major challenge to determine the effective interactions because that requires to calculate the partition function of the whole system under the constraint of a fixed configuration of the relevant degrees of freedom. The benefit for carrying out this constrained calculation, which in general is more difficult than the original full problem, is twofold. First, there is a gain in transparency by describing the system in terms of relevant degrees of freedom. Secondly, it is typically less risky to apply approximations for the partial trace because they only concern the less relevant degrees of freedom.
The determination of the phase behavior and of the structural properties of multi-component fluids represents a case study for this general approach. If the composing particles of the mixture are of comparable size and shape their degrees of freedom have to be treated on equal footing. The well developed machinery of liquid state theory offers various techniques to cope with this problem. However, these techniques fail to yield reliable results if, e.g., one component is much larger than the others; in this case numerical simulations become inefficient and integral theories lose their accuracy. Colloidal suspensions are a paradigmatic case for such highly asymmetric solutions. For their description these difficulties can be overcome by resorting to the general scheme laid out at the beginning with the positions of the colloidal particles as the relevant degrees of freedom. Accordingly the degrees of freedom of the small solvent particles are to be integrated out for a fixed configuration of the colloidal particles which we assume to be smooth, monodisperse spheres. At sufficiently low concentrations of the suspended particles this leads to an effective pair potential between them. In many cases the effective potential resembles the bare one, i.e., the one in the absence of the solvent, but with modified, effective interaction parameters which depend on the thermodynamic variables of the system such as pressure and temperature.
The effective pair potential acquires additional new features if the solvent is enriched with particles of medium size such as, e.g., polymers. If the colloidal particles come close to each other the depletion zones around them, generated by the finite size of the medium particles, overlap leading to an entropically driven attraction of the colloidal particles . Correlation effects can modify the form and the range of these depletion forces considerably . These effective potentials have indeed turned out to be successful in describing the phase behavior of colloidal suspensions .
Qualitatively new aspects arise if the solvent particles exhibit a strong cooperative behavior of their own such as a phase transition which proliferates to the effective potential between the large particles. If the solvent undergoes a continuous phase transition, thermal Casimir forces between the large particles are induced due to the geometrical constraint they pose for the critical fluctuations . Such forces are long-ranged and have a strong influence on the phase behavior of the colloidal particles . If the solvent is thermodynamically close to a first-order phase transition, wetting phenomena can occur at the surfaces of the dissolved particles (see Ref. and references therein). If the bulk phase of the solvent is the vapor phase of a one-component fluid, the surfaces of the large spheres can be covered by a liquidlike wetting film. This situation corresponds to aerosol particles floating in a vapor. If the bulk phase of the solvent is the A-rich liquid phase of a binary liquid mixture composed of (small) A and B molecules, the dissolved colloidal particles can be coated by the B-rich liquid phase of the mixture. If the wet spheres approach each other, at a critical distance the two wetting films snap to a bridgelike structure. This morphological phase transition is expected to yield a nonanalytic form of the effective interaction potential between the large spheres. This nonanalyticity demonstrates that cooperative phenomena among those degrees of freedom which are integrated out can leave clearly visible fingerprints on the effective interaction between the remaining relevant degrees of freedom. The study of this kind of profileration is not only of theoretical interest in its own right but seems to play an important (albeit not exclusive ) role for the experimentally observed flocculation of colloidal particles dissolved in a binary liquid mixture close to its demixing transition into an A-rich and a B-rich liquid phase . This observation has triggered numerous theoretical efforts devoted to various possible explanations of it. Since they are reviewed in Sec. I of Ref. and more recently in Ref. the interested reader is referred to there and we refrain from repeating this discussion here.
In our present analysis of this problem we apply density functional theory which offers two advantages. First, this technique is particularly well suited to calculate, as required here, free energies under constraints. Secondly, it allows one to keep track of the basic molecular interaction potentials of the system. We focus our interest on thermodynamic states of the solvent which are sufficiently far away from its critical point so that the emerging liquid-vapor interfaces of the wetting films exhibit only a small width. Therefore we can apply the so-called sharp-kink approximation which considers only steplike variations of the solvent density distribution and thus leaves the interface position as the main statistical variable. This approximation has turned out to be surprisingly accurate for the description of wetting phenomena . Our analysis extends and goes beyond previous efforts which are based on a similar interface displacement model grounded on a phenomenological ansatz. Whereas Refs. and are aimed at mapping out the phase diagram in terms of interaction parameters for the bridging transition mentioned above, we focus on the effective interaction potentials between the wet spheres, which are not presented in Refs. and , and on their microscopic origin. Inter alia, this allows us to compare the effective interaction potential between the colloidal particles with the bare one, i.e., in the absence of the solvent, and thus to comment on the quantitative relevance of the solvent-mediated interaction. Moreover, we present the phase diagram of the system in terms of the thermodynamic variables temperature and chemical potential which is also not contained in Refs. and .
In Sec. II we describe the implementation of a simple version of density functional theory for the present problem. For reasons of simplicity we confine our analysis to liquid-vapor coexistence of a one-component solvent; the generalization to a binary solvent is straightforward. In Sec. III we present some examples for the numerically calculated wetting film morphologies and discuss a phase diagram for the aforementioned morphological transition, and in Sec. IV we analyze the effective wetting-induced interaction potential between the spheres as a function of the distance between the spheres and the undersaturation. The experimental relevance of our model calculations is discussed in Sec. V and Sec. VI summarizes our main results. The Appendix contains some technical details.
## II Density functional theory
### A Model
We consider two identical, homogeneous, and smooth spherical particles of radius $`R`$ whose centers of mass are separated by a distance $`D`$ (see Fig. 1). They are immersed in a fluid of particles of number density $`\rho (𝐫)`$ which interact via a Lennard-Jones potential
$$\varphi (r)=4ϵ\left(\left(\frac{\sigma }{r}\right)^{12}\left(\frac{\sigma }{r}\right)^6\right).$$
(1)
The system is symmetric with respect to a rotation around the axis which connects the centers of mass of the spheres (Fig. 1) and with respect to a reflection at a plane in the middle between the spheres that is perpendicular to the symmetry axis. Since we work in a grand canonical ensemble and the fluid particles are subject to the external potential exerted by the spheres, the equilibrium number density profile of the fluid particles exhibits these symmetries, too. Therefore we describe the system in cylindrical coordinates $`(r_{},\varphi ,z)`$, with the $`z`$ axis being the symmetry axis of the system. The two centers of mass of the spheres are located at $`(r=0,z=\pm D/2)`$ such that the spheres occupy the volumes $`𝒮_\pm =\{𝐫(r_{},\varphi ,z)=(x,y,z)=(r_{}\mathrm{cos}\varphi ,r_{}\mathrm{sin}\varphi ,z)^3|\pm D/2Rz\pm D/2+R,\sqrt{r_{}^2+(zD/2)^2}R\}`$. The external potential exerted by both spheres on each individual fluid particle is
$$v_{tot}(r_{},z;R)=v(\sqrt{r_{}^2+(zD/2)^2};R)+v(\sqrt{r_{}^2+(z+D/2)^2};R)$$
(2)
where (see Eq. (A.4) in Ref. )
$`v(r;R)`$ $`=`$ $`{\displaystyle \frac{9}{8}}u_9\left({\displaystyle \frac{1}{r(r+R)^8}}{\displaystyle \frac{1}{r(rR)^8}}\right)u_9\left({\displaystyle \frac{1}{(r+R)^9}}{\displaystyle \frac{1}{(rR)^9}}\right)`$ (4)
$`{\displaystyle \frac{3}{2}}u_3\left({\displaystyle \frac{1}{r(r+R)^2}}{\displaystyle \frac{1}{r(rR)^2}}\right)+u_3\left({\displaystyle \frac{1}{(r+R)^3}}{\displaystyle \frac{1}{(rR)^3}}\right)`$
is the interaction potential between a single sphere of radius $`R`$ and a fluid particle at a distance $`r>R`$ from the center of mass of the sphere. In a continuum description, $`v(r;R)`$ follows from an integration of the Lennard-Jones potential
$$\varphi _{sf}(r)=4ϵ_{sf}\left(\left(\frac{\sigma _{sf}}{r}\right)^{12}\left(\frac{\sigma _{sf}}{r}\right)^6\right)$$
(5)
between a molecule of the spherical *s*ubstrate and a *f*luid particle. The subscript $`sf`$ denotes the parameters of the dispersion interaction between a particle in the fluid and a particle in the spheres. One has $`u_3=\frac{2\pi }{3}ϵ_{sf}\rho _s\sigma _{sf}^6`$ and $`u_9=\frac{4\pi }{45}ϵ_{sf}\rho _s\sigma _{sf}^{12}`$ where $`\rho _s`$ is the number density of the particles forming the spheres. (Many colloidal particles exhibit an even more complicated substrate potential because they are coated by a material different from their core so that they are no longer radially homogeneous as assumed for Eq. (4).)
Within our density functional approach the equilibrium particle number density distribution of the inhomogeneous fluid surrounding the spheres in a grand canonical ensemble minimizes the functional
$`\mathrm{\Omega }([\rho (𝐫)];T,\mu )`$ $`=`$ $`{\displaystyle \underset{𝒱_f}{}}d^3r\left(f_{HS}(\rho (𝐫),T)+\left(v_{tot}(𝐫)\mu \right)\rho (𝐫)\right)`$ (7)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{𝒱_f}{}}{\displaystyle \underset{𝒱_f}{}}d^3rd^3r^{}w(|𝐫𝐫^{}|)\rho (𝐫)\rho (𝐫^{}).`$
$`𝒱_f=𝒱(𝒮_+𝒮_{})`$ is the volume accessible for the fluid particles, $`𝒱`$ is the total volume of the system; $`𝒱^3`$ in the thermodynamic limit. Equation (7) does not include the bare interaction potential $`\mathrm{\Phi }(D;R)`$ (see, c.f., Sec. V) between the solid spheres, separated by vacuum, generated by the dispersion forces between the molecules forming the two spheres. $`f_{HS}(\rho ,T)`$ is the free energy density of a hard-sphere fluid of number density $`\rho `$ at temperature $`T`$. In Eq. (7), the hard-sphere reference fluid is treated in local density approximation. We apply the Weeks-Chandler-Andersen procedure to split up $`\varphi (r)`$ into an attractive part $`\varphi _{att}(r)`$ and a repulsive part $`\varphi _{rep}(r)`$. The latter gives rise to an effective, temperature dependent hard-sphere diameter
$$d(T)=\underset{0}{\overset{2^{1/6}\sigma }{}}𝑑r\left(1\mathrm{exp}\left(\frac{\varphi _{rep}(r)}{k_BT}\right)\right)$$
(8)
which is inserted into the Carnahan-Starling expression
$$f_{HS}(\rho ,T)=k_BT\rho \left(\mathrm{ln}(\rho \lambda ^3)1+\frac{4\eta 3\eta ^2}{(1\eta )^2}\right)$$
(9)
for the free energy density $`f_{HS}`$ of the *h*ard-*s*phere fluid, where $`\eta =\frac{\pi }{6}\rho (d(T))^3`$ is the dimensionless packing fraction and $`\lambda `$ is the thermal de Broglie wavelength. We approximate the attractive part of the interaction $`\varphi _{att}(r)`$ by
$$w(r)=\frac{4w_0\sigma ^3}{\pi ^2}(r^2+\sigma ^2)^3$$
(10)
with
$$w_0=_^3d^3rw(r)=_^3d^3r\varphi _{att}(r)=\frac{32}{9}\sqrt{2}\pi ϵ\sigma ^3$$
(11)
in order to simplify subsequent analytical calculations. The double integral in Eq. (7) takes into account this attractive interaction within mean-field approximation.
In the bulk the particle density $`\rho _\gamma `$ (where $`\gamma =l,g`$ denotes the liquid and vapor phase, respectively) is spatially constant, leading to (see Eq. (7))
$$\mathrm{\Omega }_b(\rho _\gamma ,T,\mu )=f_{HS}(\rho _\gamma ,T)+\frac{1}{2}w_0\rho _\gamma ^2\mu \rho _\gamma $$
(12)
for the grand canonical free energy density of the *b*ulk fluid. Minimization of $`\mathrm{\Omega }_b`$ with respect to $`\rho _\gamma `$ yields the equilibrium densities. The line $`\mu =\mu _0(T)`$ of bulk liquid-vapor coexistence and the two bulk densities $`\rho _l`$ and $`\rho _g`$ at coexistence follow from
$$\frac{\mathrm{\Omega }_b}{\rho }|_{\rho =\rho _g}=\frac{\mathrm{\Omega }_b}{\rho }|_{\rho =\rho _l}=0\text{and}\mathrm{\Omega }_b(\rho _g)=\mathrm{\Omega }_b(\rho _l).$$
(13)
For $`\mu \mu _0`$, i.e., off coexistence, only the liquid or the vapor phase is stable. In this case the density of the metastable phase corresponds to the second local minimum of $`\mathrm{\Omega }_b`$.
### B General expressions for the contributions to the effective interaction potential
Henceforth we consider the case that the substrate is sufficiently attractive so that the liquid phase is preferentially adsorbed. Therefore, if in the bulk the vapor phase is stable ($`\mu \mu _0`$), the fluid density is significantly increased in the vicinity of both spheres. In the spirit of the so-called sharp-kink approximation (see Sec. I and Ref. ) we assume that a thin film of constant density $`\rho _l`$ but with locally varying thickness is adsorbed at the surfaces of the spheres, separating the spheres from the bulk vapor phase of density $`\rho _g`$. This wetting film encapsulating both spheres is characterized by a function $`h(z)`$:
$$\rho (𝐫)=\rho (r_{},\varphi ,z)=\mathrm{\Theta }(r_{}(R+d_s))\left(\mathrm{\Theta }(h(z)r_{})\rho _l+\mathrm{\Theta }(r_{}h(z))\rho _g\right)$$
(14)
where $`\mathrm{\Theta }`$ denotes the Heaviside step function. The length $`d_s`$ takes into account the excluded volume at the surfaces of the spheres which the centers of the fluid particles cannot penetrate due to repulsive forces. The profile $`h(z)`$ as given by Eq. (14) can describe both a configuration in which the wetting films surrounding each sphere are connected by a liquid bridge as well as the configuration in which both single spheres are surrounded by disjunct wetting layers. In the latter configuration there is a region around $`z=0`$ with $`h(z)=0`$.
Inserting $`\rho (r_{},\varphi ,z)`$ from Eq. (14) into the functional $`\mathrm{\Omega }`$ in Eq. (7) leads to a decomposition of $`\mathrm{\Omega }=\text{Vol}(𝒱_f)\mathrm{\Omega }_b(\rho _g)+\mathrm{\Omega }_S`$ into a bulk and subdominant contributions. The bulk contribution is $`\text{Vol}(𝒱_f)\mathrm{\Omega }_b(\rho _g)`$ (with $`\mathrm{\Omega }_b`$ given by Eq. (12)) and corresponds to the vapor phase which is stable in the bulk. The subdominant contribution is
$$\mathrm{\Omega }_S[h]=\mathrm{\Omega }_{sl}+\mathrm{\Omega }_{ex}[h]+\mathrm{\Omega }_{ei}[h]+\mathrm{\Omega }_{lg}[h]$$
(15)
where only $`\mathrm{\Omega }_{sl}`$ is independent of $`h(z)`$ and all the other three contributions are functionals of $`h(z)`$. Since we have not found an indication for spontaneous symmetry breaking, in the following we discuss only symmetric configurations with $`h(z)=h(z)`$.
$$\mathrm{\Omega }_{ex}[h(z)]=\text{Vol}()\left(\mathrm{\Omega }_b(\rho _l)\mathrm{\Omega }_b(\rho _g)\right)$$
(16)
with
$$\text{Vol}()=2\pi \underset{0}{\overset{L_z}{}}𝑑zh^2(z)\frac{8\pi }{3}R^3$$
(17)
is an *ex*cess contribution which takes into account that the volume $`=𝒦(𝒮_{}𝒮_+)`$ is filled with the metastable liquid instead of the vapor phase; $`𝒦=\{𝐫(r_{},\varphi ,z)^3|r_{}h(z)\}`$ is the volume enclosed by the liquid-vapor interface. (The excluded volume due to $`d_s`$ enters into $`\mathrm{\Omega }_{sl}`$ (see, c.f., Eq. (23)).) This free energy contribution $`\mathrm{\Omega }_{ex}`$ vanishes at two-phase coexistence $`\mu =\mu _0(T)`$ (compare Eq. (13)). $`2L_z`$ is the extension of the total volume of the system $`𝒱`$ in $`z`$ direction; $`L_z\mathrm{}`$ in the thermodynamic limit and $`h(z>z_{max})=0`$ with $`z_{max}L_z`$.
$$\mathrm{\Omega }_{ei}[h(z)]=2\mathrm{\Delta }\rho \underset{𝒱_{}𝒦_{}}{}d^3r\left(\rho _l(t(𝐫,𝒮_{})+t(𝐫,𝒮_+))v_{tot}(𝐫)\right)$$
(18)
can be interpreted as the integrated *e*ffective *i*nteraction between the spheres and the liquid-vapor interface described by $`h(z)`$; $`\mathrm{\Delta }\rho =\rho _l\rho _g`$. $`𝒱_{}`$ is that part of the volume $`𝒱`$ with $`z<0`$ (we note again that $`𝒱_{}_{}^3`$ in the thermodynamic limit which is always considered here), analogously $`𝒦_{}`$ is the part of the set $`𝒦`$ with $`z<0`$. In Eq. (18) we have introduced the interaction potential
$$t(𝐫;)=\underset{}{}d^3r^{}w(|𝐫𝐫^{}|)$$
(19)
between a fluid particle at $`𝐫`$ and a region $``$ (with $`𝐫`$) homogeneously filled with the same fluid particles (analogous to the function $`t(z)`$ introduced in Refs. and in the case of a planar substrate). $`v_{tot}`$ is the total interaction potential between a fluid particle and both spheres (see Eq. (2)). Finally,
$$\mathrm{\Omega }_{lg}[h(z)]=(\mathrm{\Delta }\rho )^2\underset{𝒱_{}𝒦_{}}{}d^3r\left(t(𝐫;𝒦_{})+t(𝐫;𝒦_+)\right)$$
(20)
is the free energy contribution from the free *l*iquid-*g*as interface. It is a *nonlocal* functional of $`h(z)`$ in contrast to $`\mathrm{\Omega }_{ex}`$ and $`\mathrm{\Omega }_{ei}`$ whose dependence on $`h(z)`$ enters only via the integration volume $`𝒦_{}`$. The *local* approximation thereof, which is provided by the gradient expansion of Eq. (20), is
$$\mathrm{\Omega }_{lg}^{loc}=4\pi \sigma _{lg}^{(p)}\underset{0}{\overset{L_z}{}}𝑑zh(z)\sqrt{1+\left(\frac{dh}{dz}\right)^2}.$$
(21)
In Eq. (21)
$$\sigma _{lg}^{(p)}=\frac{1}{2}(\mathrm{\Delta }\rho )^2\underset{0}{\overset{\mathrm{}}{}}𝑑z\underset{z}{\overset{\mathrm{}}{}}𝑑z^{}\underset{^2}{}d^2r_{}w\left(\sqrt{r_{}^2+z^2}\right)$$
(22)
is the interfacial tension of a planar, free liquid-vapor interface in sharp-kink approximation. We note that, strictly speaking, the surface tension of a curved liquid-vapor interface depends on the local radius of curvature (see Fig. 2 in Ref. and the references therein concerning the Tolman length). This curvature dependence is omitted in the local model presented here. However, for spheres of radius $`R20\sigma `$ as considered henceforth the curvature correction is less than $`1\%`$. Similar arguments hold for the deviation of the actual liquidlike density in the wetting film from the bulk value $`\rho _l`$.
For our choice of interaction potentials (Eq. (1) and (5)) a tedious calculation leads to explicit expressions for the contributions $`\mathrm{\Omega }_{ei}`$ and $`\mathrm{\Omega }_{lg}`$ which are given in the Appendix. The remaining contribution
$$\mathrm{\Omega }_{sl}=\rho _l\underset{𝒱_{}𝒮_{}}{}d^3r\left(\rho _l(t(𝐫,𝒮_{})+t(𝐫,𝒮_+))2v_{tot}(𝐫)\right)\mathrm{\Omega }_b(\rho _l)\frac{8\pi }{3}((R+d_s)^3R^3),$$
(23)
which is independent of $`h(z)`$, is the *s*phere-*l*iquid interfacial free energy corresponding to the interface between the spheres and the liquid phase. The last term in Eq. (23) takes into account the excluded volumes at the surfaces of the spheres. In the limit of large separations $`D`$ one has
$$\mathrm{\Omega }_{sl}(D\mathrm{})2\mathrm{\Omega }_{sl}^{(1)}D^6$$
(24)
with the sphere-liquid interfacial free energy $`\mathrm{\Omega }_{sl}^{(1)}`$ of a single sphere immersed in the liquid phase. The leading power law $`D^6`$ in Eq. (24) can be inferred from the following consideration: if present, the second sphere displaces a spherical volume from the homogeneous liquid phase so that the free energy of the interaction of the first sphere with the bulk liquid is reduced by the interaction free energy of that sphere with the displaced spherical liquid volume. This latter interaction decays as $`D^6`$ for large separations $`D`$, at which the dispersion interaction between two spherical objects resembles the dispersion interaction between two pointlike particles. (Here, as before, we have not yet taken into account the bare interaction potential $`\mathrm{\Phi }(D;R)`$ between the two solid spheres; but see, c.f., Sec. V.)
Up to the bulk contribution the grand canonical potential of the system is the minimum of $`\mathrm{\Omega }_S[h(z)]`$ with respect to the profile $`h(z)`$:
$$\mathrm{\Omega }_S=\mathrm{\Omega }_S(D;R)=\underset{\{h(z)\}}{\mathrm{min}}(\mathrm{\Omega }[h(z)]).$$
(25)
Thus the equilibrium interface morphology $`h(z)`$ minimizes $`\mathrm{\Omega }_S[h(z)]`$ which includes the contributions $`\mathrm{\Omega }_{ex}[h(z)]`$, $`\mathrm{\Omega }_{ei}[h(z)]`$, $`\mathrm{\Omega }_{lg}[h(z)]`$, and $`\mathrm{\Omega }_{sl}`$. The functional used in Refs. and (Eq. (1) in both references) is, albeit formulated in another coordinate system and using a more phenomenological ansatz for the basic interaction potentials, essentially identical with the sum $`(\mathrm{\Omega }_{lg}^{loc}+\mathrm{\Omega }_{ex}+\mathrm{\Omega }_{ei})[h(z)]`$. However, this model description does neither incorporate the bare dispersion interaction of the two spheres (c.f., Sec. V) nor the free energy contribution $`\mathrm{\Omega }_{sl}`$ which describes the sphere-liquid interfacial free energy. We emphasize that the consideration of the contribution $`\mathrm{\Omega }_{sl}`$ – which does not depend on $`h(z)`$ – is not essential for the determination of the equilibrium wetting film morphology and hence it is not relevant for the thermodynamic phase diagram of thin-thick and bridging transitions (Fig. 2 in Ref. ) for a *fixed* separation $`D`$ between the spheres. But the term $`\mathrm{\Omega }_{sl}`$ is crucial to the *shape* of the effective, wetting-induced interaction potential between the spheres, i.e., its dependence on $`D`$ (see Eq. (24)).
## III Morphology of the wetting layers
### A Interface profiles
The actual wetting layer morphology $`h(z)`$ follows from numerical minimization of the functional $`\mathrm{\Omega }_S[h(z)]`$ (Eq. (15)) for a given temperature $`T`$ and undersaturation $`\mathrm{\Delta }\mu =\mu _0(T)\mu `$, with the contributions $`\mathrm{\Omega }_{ex}`$ (Eq. (16)), $`\mathrm{\Omega }_{ei}`$ (Eq. (18)), $`\mathrm{\Omega }_{sl}`$ (Eq. (23)), and $`\mathrm{\Omega }_{lg}`$ (Eq. (20) within the nonlocal and Eq. (21) for the local theory). Within a range of parameters $`(T,\mathrm{\Delta }\mu )`$ the numerical minimization yields two different solutions for $`h(z)`$, one with a liquid bridge and one without, depending on the initial function $`h(z)`$ used in the iteration scheme for the minimization. For small separations $`a2R`$ only the solution which exhibits a liquid bridge is stable whereas for large separations $`a2R`$ only the solution without bridge minimizes $`\mathrm{\Omega }_S`$. For large distances $`D`$ the minimization consistently yields twice the result known for a single individual sphere enclosed by a wetting film (compare Ref. ). This observation amounts to a useful check of the numerical procedure.
As a first example, in Fig. 2 we present the numerical results for a wetting layer enclosing two spheres of radius $`R=20\sigma `$. For our particular choice of interaction parameters, at coexistence $`\mathrm{\Delta }\mu =0`$ the wetting film on each of the single spheres alone exhibits a first-order *t*hin-*t*hick transition (which is the remnant of the first-order wetting transition on the corresponding planar substrate, see Fig. 8(a) in Ref. ) at $`T_{tt}^{}=k_BT/ϵ1.271`$ (which corresponds to $`T_{tt}/T_c0.9`$ where $`T_c`$ is the critical temperature of gas-liquid coexistence in the bulk). The planar substrate, i.e., a single sphere in the limit $`R\mathrm{}`$, exhibits a genuine first-order wetting transition (with the film thickness jumping to a macroscopic value) at $`T_w^{}1.053`$ ($`T_w/T_c0.75`$, $`T_{tt}/T_w1.21`$). Figure 2(a) depicts a typical solution with a bridge, here for a separation $`a=D2R=10\sigma `$ ($`D=50\sigma `$) and the thermodynamic parameters $`T^{}=1.3>T_{tt}^{}`$ and $`\mathrm{\Delta }\mu =0`$, i.e., at liquid-vapor coexistence. The solution without a bridge for the same choice of parameters is shown in Fig. 2(b). The latter solution has a higher free energy than the former one. Therefore the solution with bridge is thermodynamically stable whereas the solution without bridge is metastable. For the solution without a bridge the distortion of the liquidlike layer around one sphere due to the presence of the other sphere is not visible. Finally, Fig. 2(c) displays the wetting film morphology for the stable state with bridge at the temperature $`T^{}=1.2`$, i.e., below the thin-thick transition temperature $`T_{tt}^{}`$. (We note that the thin-thick transition temperature $`T_{tt}`$ for each sphere is slightly shifted by the presence of the second sphere. However, as already pointed out in Ref. , this effect is negligibly small.) In any case, the difference between the nonlocal and the local theory is very small. This latter result is in accordance with the findings for the comparison between the nonlocal and the local description of the three-phase contact line on a homogeneous substrate and of the wetting layer morphology on a chemically structured substrate (compare Ref. ). For this reason, henceforth we only consider the local theory.
Figure 3 shows another pertinent example. Here we study the wetting layer morphology for two larger spheres of radius $`R=50\sigma `$ as a function of the undersaturation $`\mathrm{\Delta }\mu `$ along the isotherm $`T^{}=1.2`$. The interaction potential parameters are the same as for the previous first example and the separation of the surfaces $`a`$ is $`20\sigma `$ ($`D=120\sigma `$). At coexistence each single sphere exhibits a first-order thin-thick transition at $`T_{tt}^{}1.191`$ (i.e., $`T_{tt}/T_c0.84`$ and $`T_{tt}/T_w1.13`$). In analogy to the prewetting line on a homogeneous substrate there is a line of thin-thick transitions $`(T,\mathrm{\Delta }\mu _{tt}(T))`$ which intersects the liquid-vapor coexistence line at $`(T=T_{tt},\mathrm{\Delta }\mu =0)`$ (compare with, c.f., Fig. 4 and Fig. 8(a) in Ref. ). At the temperature $`T^{}=1.2>T_{tt}^{}`$ considered here the thin-thick transition occurs at $`\mathrm{\Delta }\mu _{tt}^{}=\mathrm{\Delta }\mu _{tt}/ϵ0.0103`$. Upon reducing the undersaturation along the isotherm, starting at, e.g., $`\mathrm{\Delta }\mu ^{}=0.05`$, first the configuration with thin films and without bridge is stable (Fig. 3(a)). For $`\mathrm{\Delta }\mu \mathrm{\Delta }\mu _{bt}`$ (*b*ridging *t*ransition) with $`\mathrm{\Delta }\mu _{bt}^{}0.0235>\mathrm{\Delta }\mu _{tt}^{}(T)`$ the solution with bridge becomes stable, but the layers enclosing the spheres still remain thin (Fig. 3(b)). Upon further reduction of $`\mathrm{\Delta }\mu `$, at $`\mathrm{\Delta }\mu _{tt}(T)`$ the second transition from a solution with bridge and thin films to a solution with bridge and thick films (Fig. 3(c)) takes place. (As before, concerning the value of $`T_{tt}^{}`$ at coexistence, also the value $`\mathrm{\Delta }\mu _{tt}^{}(T)`$ is practically unchanged by the presence of the second sphere – even for the bridge configuration.) We note that for this choice of parameters and in the case of a solution with bridge and *thin* films (Fig. 3(b)) the profile $`h(z)`$ exhibits *six* turning points instead of only two as for the case of a solution with bridge and *thick* films (Fig. 3(c)). This rich curvature behavior is caused by the details of the effective interaction potential between the spherical substrate surfaces and the liquid-vapor interface (see Sec. 2.3 in Ref. ), similar to the curvature behavior of the liquid-vapor interface when it meets a homogeneous, planar substrate forming a three-phase contact line (compare Ref. ). These features may also occur for a bridge configuration with thin films at coexistence and $`T<T_{tt}`$.
### B Phase diagram
The example presented in the previous paragraph shows that besides the gas-liquid coexistence curve $`\mathrm{\Delta }\mu =0`$ the $`T`$-$`\mathrm{\Delta }\mu `$ phase diagram of the system contains two distinct lines of first-order phase transitions: a line of thin-thick transitions $`(T,\mathrm{\Delta }\mu _{tt}(T))`$ on the single spheres (which is the remnant of the line of prewetting transitions on the corresponding flat substrate and which is, as stated above, practically unshifted by the presence of the second sphere) and a second, *independent* line of bridging transitions $`(T,\mathrm{\Delta }\mu _{bt}(T))`$. If one crosses the latter along an isotherm $`T=T_0`$ approaching coexistence $`(T_0,\mathrm{\Delta }\mu 0)`$, at $`\mathrm{\Delta }\mu =\mathrm{\Delta }\mu _{bt}(T_0)`$ a transition from the configuration without bridge $`(\mathrm{\Delta }\mu >\mathrm{\Delta }\mu _{bt}(T_0))`$ to a configuration with bridge $`(\mathrm{\Delta }\mu <\mathrm{\Delta }\mu _{bt}(T_0))`$ occurs. The derivative $`\mathrm{\Omega }_S/\mathrm{\Delta }\mu `$ is discontinuous at $`\mathrm{\Delta }\mu _{bt}`$, indicating that the bridging transition is first order. Figure 4 shows the $`T`$-$`\mathrm{\Delta }\mu `$ phase diagram for the two spheres with $`R=20\sigma `$ for $`D=50\sigma `$ ($`a=10\sigma `$). The line of thin-thick transitions intersects the liquid-vapor coexistence line at $`T_{tt}^{}1.271`$ with a finite, negative slope (compare Fig. 8(a) in Ref. ). It extends into the vapor phase region ($`\mathrm{\Delta }\mu >0`$) of the phase diagram and ends at a critical point. The line of bridging transitions intersects the coexistence line also with a finite, negative slope. On the other end, within our sharp-kink interface model, it happens to be cut off at that metastability line in the phase diagram at which the second minimum of the bulk free energy at high fluid density (Eq. (12)) ceases to exist so that for larger undersaturations the liquid phase is not even metastable. Within a more sophisticated approach, e.g., by seeking the full minimal density distributions of Eq. (7), the line of bridging transitions is expected to end in a critical point, too. (Concerning the effect of fluctuations on these mean field predictions see the following paragraph.) The line of bridging transitions is entirely located in the region where the liquidlike films on the spheres are thin. Moreover, the effect of the presence of the liquid bridge on the line of thin-thick transitions is negligibly small. In Fig. 4 the relative location of the bridging transition line and the thin-thick transition line corresponds to our specific choice of the interaction potential parameters as well as the chosen size of and distance between the spheres. Changing these parameters will lead to shifts of these lines and, possibly, to different topologies of the phase diagram. Here we refrain from exhaustingly presenting all possibilities which can occur according to Refs. and .
Since the liquid volume enclosed by the interface $`h(z)`$ is quasi-zerodimensional, fluctuation effects destroy the sharp first-order phase transition (see Refs. and ). In Sec. 4 of Ref. it has been extensively discussed how finite size effects smear out the thin-thick transition such that the thickness increases sharply but continuously within a range $`\delta \mu `$ around $`\mathrm{\Delta }\mu _{tt}(T)`$; these results apply analogously to the present case. Using similar approximations we obtain a range $`\delta \mu `$ between $`\delta \mu ^{}0.004`$ for $`T^{}=1.16`$ and $`\delta \mu ^{}0.02`$ for $`T^{}=1.26`$ over which the *bridging* transitions shown in Fig. 4 are smeared out around $`\mathrm{\Delta }\mu _{bt}(T)`$. Thus close to $`\mathrm{\Delta }\mu =0`$ the quasi-first-order thin-thick transitions are clearly visible. However, for larger values of $`\mathrm{\Delta }\mu `$ they become progressively smeared out such that their critical points predicted by mean field theory are erased by fluctuations.
## IV Effective film-induced interaction potential
### A Shape of the effective potential, metastability, and asymptotic behavior
In the following we change our point of view: we vary the distance $`D`$ between the centers of mass of the spheres instead of the thermodynamic parameters $`T`$ and $`\mathrm{\Delta }\mu `$. Figure 5 shows the grand canonical potential $`\mathrm{\Omega }_S`$ corresponding to the wetting layer morphologies for the case $`R=20\sigma `$ and $`T^{}=1.2`$ (Fig. 2(c)) as a function of the separation $`a=D2R`$ for several values of $`\mathrm{\Delta }\mu `$. $`\mathrm{\Omega }_S`$ is the minimum of $`\mathrm{\Omega }_S[h(z)]`$ (Eq. (15)) for the given set of parameters $`T`$, $`\mathrm{\Delta }\mu `$, and $`D=2R+a`$. For each value of $`\mathrm{\Delta }\mu `$ there are two branches of the free energy, one corresponding to the solution without bridge, which for the case $`R=20\sigma `$ considered here exists only for $`a0.15R`$, and the other corresponding to the solution with bridge which exists up to $`a0.65R`$ and $`a0.6R`$ for $`\mathrm{\Delta }\mu ^{}=0`$ and $`\mathrm{\Delta }\mu ^{}=0.01`$, respectively. At a certain value $`D=D_{bt}`$ or, equivalently, $`a=a_{bt}`$, which are functions of $`\mathrm{\Delta }\mu `$, a first-order phase transition occurs with discontinuous derivative $`\mathrm{\Omega }_S/D`$ between the solutions with and without bridge. The main effect of increasing the undersaturation $`\mathrm{\Delta }\mu `$ is that the free-energy curves are rigidly shifted upwards. This shift is approximately proportional to $`\mathrm{\Delta }\mu `$ and larger in the case of the solution with bridge, resulting in the dependence of $`D_{bt}`$ on $`\mathrm{\Delta }\mu `$. The values of $`\mathrm{\Omega }_S`$ shown in Fig. 5 are obtained within the local theory. The nonlocal theory yields the same functional dependence $`\mathrm{\Omega }_S(D)`$ but with a slight and rigid shift of the free-energy curves, relative to the results of the local theory, of the order of $`0.1\%`$ and of the same sign and size for both the solutions with and without bridge. Finite-size effects again destroy the sharp first-order bridging transition; we obtain a range $`\delta D0.1\sigma `$ (corresponging to $`\delta D0.005R`$) over which the bridging transitions shown in Fig. 5 are smeared out.
The thermodynamic states which are located on the metastable branches of the free energy curves survive during an average lifetime $`\tau \tau _0\mathrm{exp}(\mathrm{\Delta }\mathrm{\Omega }_S/k_BT)`$ where $`\mathrm{\Delta }\mathrm{\Omega }_S`$ is the height of the energy barrier that separates the metastable from the stable branch and $`\tau _0`$ is a characteristic microscopic time scale for the dynamics associated with the transition from a metastable to a stable wetting layer configuration. The energy barrier is highest in the vicinity of the bridging transition and vanishes near the ends of the metastable branches. An estimation of the energy barrier height yields, e.g., $`\mathrm{\Delta }\mathrm{\Omega }_S75ϵ`$ for $`\mathrm{\Delta }\mu =0`$ and $`D=50\sigma `$ ($`a=0.5R`$), and with $`k_BTϵ`$ it follows that $`\mathrm{exp}(\mathrm{\Delta }\mathrm{\Omega }_S/k_BT)10^{32}`$, i.e., the metastable unbridged state for $`a=0.5R`$ near the bridging transition remains stable practically forever. However, at, e.g., $`a=0.2R`$ one has $`\mathrm{exp}(\mathrm{\Delta }\mathrm{\Omega }_S/k_BT)10^{11}`$ so that with $`\tau _01`$ps$`\mathrm{}1`$ns one may observe a decay of the metastable states near the ends of the metastable branches within seconds or minutes. Thus the change of the morphology of the wetting films is expected to exhibit pronounced hysteresis effects as function of $`D`$.
Obviously, in the limit of large separation $`D\mathrm{}`$ (in which only the configuration without a bridge is stable) the grand canonical potential $`\mathrm{\Omega }_S(D)`$ approaches the limiting value $`2\mathrm{\Omega }_S^{(1)}`$ corresponding to the free energy of two individual spheres, each surrounded by a wetting layer. It is convenient to separate this constant contribution $`2\mathrm{\Omega }_S^{(1)}`$ from the grand canonical potential $`\mathrm{\Omega }_S`$ of the system and thus to define an *e*xcess free energy $`\mathrm{\Omega }_E(D)=\mathrm{\Omega }_S(D)2\mathrm{\Omega }_S^{(1)}`$ which contains all contributions from the wetting-layer induced interaction between the two spheres. In the limit $`D\mathrm{}`$, i.e., in the absence of a liquid bridge, this excess free energy $`\mathrm{\Omega }_E(D)`$ decays as $`D^6`$ (see Eq. (24) and, c.f., Sec. V). We note that for the example shown in Fig. 5 the coefficient of this leading order is *positive*, i.e., the effective potential in the absence of a liquid bridge is *repulsive*. This is owed to the choice $`T<T_{tt}`$ for this example: the spheres disfavor the adsorption of thick liquid films and the presence of the second sphere with its surrounding liquidlike layer leads to an additional cost in free energy which diminishes for increasing $`D`$. For the choice $`T>T_{tt}`$, i.e., if the spheres favor the adsorption of liquid (e.g., for $`T^{}=1.3`$ as in Figs. 2(a) and (b)) the coefficient of $`D^6`$ is *negative* and the effective interaction is *attractive*. However, in the presence of a liquid bridge, i.e., for sufficiently small values of $`D`$, the effective potential shows the same qualitative behavior as in Fig. 5 also for the case of thick wetting layers ($`T^{}=1.3>T_{tt}^{}`$) as well as for the larger spheres ($`R=50\sigma `$) with thin or thick films.
### B Effective interaction potential for large spheres
In this subsection we consider the limiting case that the sphere radius $`R`$ is much larger than the diameter $`\sigma `$ of the solvent particles and that the separations $`a`$ between the surfaces of the spheres are proportional to $`R`$: $`R\sigma `$, $`\sigma aR`$. For such large separations as compared to $`\sigma `$ the contributions $`\mathrm{\Omega }_{sl}`$ (Eq. (23)), $`\mathrm{\Omega }_{ei}`$ (Eq. (18)), and $`\mathrm{\Phi }`$ (c.f., Eq. (31)) become vanishingly small relative to the contributions $`\mathrm{\Omega }_{lg}`$ (Eq. (20)) and $`\mathrm{\Omega }_{ex}`$ (Eq. (16)). For the case described above $`\mathrm{\Omega }_{lg}`$ and $`\mathrm{\Omega }_{ex}`$ scale proportional to the surface area of the spheres, i.e., $`R^2`$, whereas for $`a/\sigma \mathrm{}`$, $`R/\sigma \mathrm{}`$, $`a/R`$ finite, $`\mathrm{\Phi }(D;R)`$ remains finite $`ϵ_{ss}\sigma _{ss}^6\rho _s^2`$ with a proportionality constant of the order $`1`$. Analogously, in the same limit $`\mathrm{\Omega }_{ei}2\mathrm{\Omega }_{ei}^{(1)}`$ (Eq. (18)) and $`\mathrm{\Omega }_{sl}2\mathrm{\Omega }_{sl}^{(1)}`$ (Eq. (23)) are determined by finite terms $`\mathrm{\Delta }\rho \rho _lϵ\sigma ^6`$ and $`\mathrm{\Delta }\rho \rho _sϵ_{sf}\sigma _{sf}^6`$ and of terms $`\rho _l^2ϵ\sigma ^6`$ and $`\rho _l\rho _sϵ_{sf}\sigma _{sf}^6`$, respectively, each with a proportionality constant of the order $`1`$. Therefore measured in units of $`8\pi R^2`$ the unbridged branch of $`\mathrm{\Omega }_E`$ in Fig. 5(b) vanishes in the limit $`R\mathrm{}`$. Moreover, on this scale the excluded volume at small $`a`$ disappears from the figure, too, because $`d_s/R0`$.
Fig. 6 shows the excess effective interaction potential $`\mathrm{\Omega }_E`$ in the limit of large spheres for the case $`\mathrm{\Delta }\mu =0`$, i.e., at two-phase coexistence in the solvent. In this limit and for $`\mathrm{\Delta }\mu =0`$, $`\mathrm{\Omega }_{lg}`$ is the only relevant contribution to $`\mathrm{\Omega }_S`$ because $`\mathrm{\Omega }_{ex}(\mathrm{\Delta }\mu =0)=0`$. Accordingly, in this case the bridging transition is determined by the equality of the surface areas of the liquid-vapor interfaces for the unbridged and bridged configuration. From this condition and from dimensional analysis it follows that for large spheres $`D_{bt}(\mathrm{\Delta }\mu =0)`$ is determined by the equation
$$8\pi (R+l_0)^2\sigma _{lg}^{(p)}=8\pi (R+l_0)^2\sigma _{lg}^{(p)}f\left(\frac{D_{bt}}{R+l_0}\right)$$
(26)
where $`f`$ is, for dimensional reasons, a universal function of $`D/(R+l_0)`$ alone which describes the surface area of the bridged configuration; $`l_0`$ is the equilibrium wetting layer thickness on a single sphere. Since the line of bridging transitions lies below the line of thin-thick transitions, $`l_0`$ remains microscopicly small at the bridging transition (Fig. 4). Therefore one has
$$D_{bt}(\mathrm{\Delta }\mu =0)=\lambda (R+l_0)$$
(27)
with a universal number
$$\lambda 2.32$$
(28)
determined by $`f(\lambda )=1`$ (compare Fig. 6). If one applies this reasoning to Fig. 5 one finds $`\lambda 2.39`$. Therefore even for $`R=20\sigma `$ this macroscopic approximation leads to a surprisingly small error of only $`3\%`$ for $`D_{bt}(\mathrm{\Delta }\mu =0)`$. Accordingly, in Fig. 5 the full curves corresponding to $`\mathrm{\Delta }\mu =0`$ closely resemble the ones in Fig. 6 describing the case of large spheres. The only differences appear for small separations $`a`$ where the bridged branch linearly extends down to its minimum value $`\mathrm{\Omega }_E/(8\pi (R/\sigma )^2)0.0227ϵ`$ at $`a/R=0`$ (Fig. 6). Only in this range of separations the effect of the contributions $`\mathrm{\Omega }_{ei}`$ and $`\mathrm{\Omega }_{sl}`$ becomes significant, leading to the deeper minimum visible in Fig. 5. Thus for $`\mathrm{\Delta }\mu =0`$ and large $`R`$ the dependence of the effective interaction potential on $`R`$ for the bridged configuration is captured by the indicated rescaling of the axes in Fig. 5(b). However, our numerical analysis shows that the smallness of the deviations between the macroscopic description valid for $`R\sigma `$ and the actual results for $`R=20\sigma `$ is somewhat fortuitous. Whereas the dependence of $`D_{bt}(\mathrm{\Delta }\mu =0)`$ on $`R`$ is indeed weak, the shape of the potential (for $`\mathrm{\Delta }\mu =0`$) reduces to that shown in Fig. 6 only for $`R`$ larger than several hundred $`\sigma `$ and, surprisingly, for $`R`$ up to $`20\mathrm{}30\sigma `$ with the deviations being maximal for $`R100\sigma `$.
Off coexistence $`\mathrm{\Delta }\mathrm{\Omega }_b=\mathrm{\Omega }_b(\rho _l)\mathrm{\Omega }_b(\rho _g)\mathrm{\Delta }\mu \mathrm{\Delta }\rho `$ is positive so that Eq. (26) has to be augmented correspondingly:
$$8\pi (R+l_0)^2+\frac{8\pi \mathrm{\Delta }\mathrm{\Omega }_b}{3\sigma _{lg}^{(p)}}\left((R+l_0)^3R^3\right)=𝒜+\frac{\mathrm{\Delta }\mathrm{\Omega }_b}{\sigma _{lg}^{(p)}}\text{Vol}()$$
(29)
where $`𝒜`$ and $`\text{Vol}()`$ (Eq. (17)) are the area of the liquid-vapor interface and the volume of the liquid, respectively, for the bridged configuration. They are obtained by inserting into Eqs. (21) and (17) that profile $`h(z)`$ which solves the differential equation determining the minimum of $`\mathrm{\Omega }_{lg}[h]+\mathrm{\Omega }_{ex}[h]`$ together with the appropriate boundary conditions. By splitting off a factor $`(R+l_0)^2`$ from $`𝒜`$ and $`(R+l_0)^3`$ from $`\text{Vol}()`$ dimensional analysis shows that up to terms $`l_0/R`$ the critical distance for the bridging transition is given by a universal scaling function $`\mathrm{\Lambda }`$:
$$D_{bt}(\mathrm{\Delta }\mu )=\mathrm{\Lambda }\left(\frac{\mathrm{\Delta }\rho \mathrm{\Delta }\mu R}{\sigma _{lg}^{(p)}}\right)R$$
(30)
with $`\mathrm{\Lambda }(0)=\lambda `$. Thus off coexistence the critical bridging transition depends, apart from an explicit factor $`R`$, on $`R`$ and $`\mathrm{\Delta }\mu `$ via the scaling variable $`\mathrm{\Delta }\rho \mathrm{\Delta }\mu R/\sigma _{lg}^{(p)}`$. This property is shared by the whole bridged branch of the effective interaction potential. Thus increasing $`R`$ for fixed undersaturation $`\mathrm{\Delta }\mu `$ has the same effect as increasing $`\mathrm{\Delta }\mu `$ for fixed $`R`$. From Fig. 5(b), in which the unbridged branch will disappear in the limit $`R\sigma `$, one infers that the range and the depth of $`\mathrm{\Omega }_E`$ decrease for increasing $`R`$ at fixed undersaturation $`\mathrm{\Delta }\mu `$. The behavior of $`D_{bt}`$ and of the bridged branch of the effective interaction potential off coexistence and for $`R\mathrm{}`$ is determined by the behavior of the scaling function $`\mathrm{\Lambda }(x)`$ in the limit $`x\mathrm{}`$. Our numerical data indicate that $`\mathrm{\Lambda }(x\mathrm{})<2`$ so that due to the geometric constraint $`D2R`$ there is no bridging transition and the bridged branch of the effective potential vanishes for any value of $`\mathrm{\Delta }\mu `$ in the limit $`R\mathrm{}`$. The cost in free energy due to the excess contribution $`\mathrm{\Omega }_{ex}`$ suppresses the formation of a liquidlike bridge in the case of macroscopicly large spheres. In turn, this means that for any finite value of $`\mathrm{\Delta }\mu `$ there is a large but finite critical radius $`R_c`$ for which the critical separation $`a_{bt}`$ for the bridging transition attains the value $`a_{bt}=0`$, such that for $`R>R_c`$ there is no bridging transition. The determination of $`R_c`$ requires to analyze the full dependence of $`\mathrm{\Lambda }`$ on the scaling variable $`x`$. This, however, implies such a large numerical effort that it is beyond the scope of the present paper.
## V Discussion
### A Total interaction potential
The bare dispersion interaction between the two spheres is not included in Eq. (7). According to Hamaker this contribution is given by
$$\mathrm{\Phi }(D;R)=\frac{A_{ss}}{12}\left(\frac{4R^2}{(D2R)(D+2R)}+\frac{4R^2}{D^2}+2\mathrm{ln}\left(\frac{(D2R)(D+2R)}{D^2}\right)\right)$$
(31)
as the dispersion interaction between two identical spheres of radius $`R`$ at center-of-mass distance $`D`$. In the limit $`a/R1`$, where $`a=D2R`$ (see Fig. 1(a)) is the smallest separation between the surfaces of the spheres, Eq. (31) reduces to
$$\mathrm{\Phi }(D=2R+a;Ra)\frac{A_{ss}}{12}\frac{R}{a}$$
(32)
which corresponds to the Derjaguin approximation whereas $`\mathrm{\Phi }(DR;R)=16A_{ss}R^6/(9D^6)`$. Thus except for the $`D`$-independent bulk contribution $`\text{Vol}(𝒱_f)\mathrm{\Omega }_b(\rho _g)`$ the total grand canonical potential of the system is
$$\mathrm{\Omega }_{tot}(D;R)=\mathrm{\Omega }_S(D;R)+\mathrm{\Phi }(D;R)$$
(33)
where $`\mathrm{\Omega }_S(D;R)`$ is given by the minimum value $`\mathrm{min}_{\{h(z)\}}(\mathrm{\Omega }_S[h(z)])`$ for given $`D`$ and $`R`$ (Eqs. (15) and (25)); in analogy to $`\mathrm{\Omega }_E`$ we define the excess total free energy $`\mathrm{\Omega }_{E,tot}=\mathrm{\Omega }_{tot}2\mathrm{\Omega }_S^{(1)}`$. $`A_{ss}`$ is the Hamaker constant appertaining to the bare dispersion interaction between the particles in the spheres. In the case of pairwise additivity of the molecular interactions and in the absence of retardation effects one has $`A_{ss}=4\pi ^2ϵ_{ss}\sigma _{ss}^6\rho _s^2`$ if the interaction potential between two individual molecules in the spheres is given by a Lennard-Jones potential (Eq. (5)) with the parameters $`ϵ_{ss}`$ and $`\sigma _{ss}`$. Typically $`A_{ss}`$ is of the order of $`10^{19}`$J or, equivalently, $`10\mathrm{}100ϵ`$. If the vacuum between the spheres is replaced by a medium of condensed matter the interaction between the spheres is screened . In our present model this medium is the bulk vapor phase modified by the presence of the liquidlike films adsorbed on the spheres and the screening effect is described microscopicly by the functional $`\mathrm{\Omega }[\rho (𝐫)]`$.
In Refs. and this additional screening effect, due to spherical shells of adsorbed, homogeneous layers surrounding spherical particles, on the dispersion interaction between the latter immersed in another homogeneous medium has been calculated macroscopicly. Beyond molecular scales these results should closely correspond to the configuration without liquid bridge discussed herein because the deviation of the spherical shape of one wetting layer due to the presence of the second sphere is very small. Indeed, the interaction energy calculated in Refs. and is practically the same as the sum of the $`D`$-dependent contributions in $`\mathrm{\Omega }_{ei}`$ (Eq. (18)) and $`\mathrm{\Omega }_{sl}`$ (Eq. (23)) for the configurations without bridge – for these configurations $`\mathrm{\Omega }_{lg}`$ and $`\mathrm{\Omega }_{ex}`$ do not contribute to the dependence of $`\mathrm{\Omega }_S`$ on $`D`$ – and the direct dispersion interaction $`\mathrm{\Phi }(D;R)`$. In Ref. the total dispersion interaction is shown to be always attractive if the Hamaker constants $`A_{ij}`$ corresponding to the interaction between any two media $`i`$ and $`j`$ are chosen such that $`A_{ij}=\sqrt{A_{ii}A_{jj}}`$. Although the effective interaction induced by the wetting layers shown in Figs. 5(b) and 7 for the configuration without bridge is repulsive, we note that the sum $`\mathrm{\Omega }_{tot}`$ of this interaction and of the bare dispersion potential $`\mathrm{\Phi }(D;R)`$ is also *attractive* if we choose the Hamaker constant in Eq. (31) accordingly, i.e., $`A_{ss}=A_{sf}^2/A`$ (Fig. 7). Therefore our results are consistent with those obtained in Ref. . Since only effective interactions between finite volumes enter into the total excess interaction potential $`\mathrm{\Omega }_{E,tot}`$ and these effective interactions decay as $`D^6`$ in the limit of large separations $`D`$, the same holds for $`\mathrm{\Omega }_{E,tot}`$.
Figure 5 shows that as soon as the wetting films snap to a liquidlike bridge, whether it is stable or metastable, there is an attractive wetting-layer-induced force $`\mathrm{\Omega }_E/D`$ that pulls the spheres together. From Fig. 5 one can infer that this attractive force is of the order of $`40ϵ/\sigma `$ in the range between $`a4\sigma `$ (i.e., $`0.2R`$ for $`R=20\sigma `$ discussed in this figure) and $`a10\sigma `$ ($`0.5R`$) where the effective potential varies almost linearly. At the small separation $`a_{min}2.5\sigma `$ the effective potential $`\mathrm{\Omega }_E`$ induced by the bridgelike wetting layer is minimal and the wetting-induced force is zero. Finally, at smaller separations the interaction is repulsive leading to a stabilization of the spheres at $`D=D_{min}=2R+a_{min}`$. Within the range $`aR`$ the bare, direct dispersion interaction between the spheres (Eqs. (31) and (32)) gives rise to a force $`F_{bare}(a)A_{ss}R/12a^2`$. The estimate $`A_{ss}4\pi ^2ϵ_{sf}^2\sigma _{sf}^{12}\rho _s^2/ϵ\sigma ^6400ϵ`$ for the case of pairwise additive interactions without retardation follows from the ansatz $`A_{sf}=\sqrt{AA_{ss}}`$, so that the bare dispersion force in our example with $`R=20\sigma `$ is $`F_{bare}(a)670ϵ\sigma /a^2`$. Therefore in the range where the bridge-induced force is almost constant ($`4\sigma a10\sigma `$) the direct, bare dispersion force decays from approximately $`40ϵ/\sigma `$ (which is of the same order of magnitude as the bridge-induced force) to approximately $`6ϵ/\sigma `$, whereas for smaller separations it becomes the dominant force.
### B Relevance for force microscopy
Our model calculations can be tested experimentally by force microscopy. This can be done by suitably fixing one sphere in the fluid and by attaching the second one to the tip of a force microscope. Alternatively, both spheres can be positioned by optical tweezers and the force law can be inferred by monitoring optically their dynamics after switching off the tweezers. At separations between the spheres which are comparable with the diameter $`\sigma `$ of the solvent particles the actual effective interaction potential will exhibit an additional oscillatory contribution due to packing effects which decays exponentially on the scale $`\sigma `$ . In order to obtain these oscillations one would have to resort to density functional theories which are more sophisticated than the one in Eq. (7). This, in turn, would make it much more difficult to obtain the bridgelike configuration, to map out the complete phase diagram, and to obtain results for large spheres. According to Subsec. IV B, for $`R\sigma `$ and at two-phase coexistence $`\mathrm{\Delta }\mu =0`$ the bridging transition occurs at distances $`a`$ which are proportional to $`R`$. In this case, due to $`R\sigma `$, the effective interaction potential will be practically unaffected by this oscillatory contribution for the vast portion $`\sigma aa_{bt}`$ of the range of the effective interaction potential.
### C Relevance for charge stabilized colloidal suspensions
Whereas the kinds of experiments considered in the previous subsection are focused on two individual spherical particles, we discussed in the Introduction that the effective interaction potential enters into the collective behavior of colloidal suspensions such that the bridging transition may trigger flocculation. If colloidal suspensions would be governed by dispersion forces alone, most of them would flocculate even in the absence of the wetting-induced forces discussed here because the dispersion forces generate the so-called primary minimum in the effective interaction potential close to contact. Since this minimum is much deeper than $`k_BT`$ the colloidal particles would simply stick together permanently. This effect, which is undesired for many applications, can be avoided by endowing the particles with electrical charges which adds a screened Coulomb repulsion between the charged particles. As a result, such charge stabilized colloidal suspensions are characterized by effective interaction potentials in which a substantial energy barrier separates the aforementioned primary minimum from a second, much more shallow minimum at larger distances. Since this potential barrier is typically large compared with $`k_BT`$ the phase behavior of the colloidal particles is practically independent of the primary minimum formed by the dispersion forces and determined by the shape of the potential *outside* the barrier. As demonstrated by Figs. 5 and 7 the range of the wetting-induced forces is about $`0.55R`$, in good agreement with $`D_{bt}2.32(R+l_0)`$ (see Eqs. (27) and (28)). On the other hand the position (and height) of the aforementioned energy barrier depends sensitively on the size of the total charge on the spheres, the amount of salt in the solvent, and the dispersion forces and can be varied over a wide range. With a high salt concentration the barrier position can be as small as a few tens of nm. Thus under such circumstances the wetting-induced interaction potentials would be relevant even for colloidal particles whose radii are only a few tens of nm.
### D Relevance for stericly stabilized colloidal suspensions
There is another class of colloidal suspensions for which the wetting-induced forces can be of practical importance. By coating the colloidal particles with polymers and by matching the refractive indices of the colloidal particles and the bulk fluid (in our case study the vapor phase or, more realisticly in the present context, the A-rich liquid phase of a binary liquid mixture acting as the solvent) the colloidal particles behave effectively like hard spheres (see, e.g., Refs. and ). Through this index matching the sum of the bare interaction potential $`\mathrm{\Phi }(D;R)`$ and the effective interaction potential $`\mathrm{\Omega }_{sg}`$, which would arise if the spheres were immersed in the homogeneous and unperturbed bulk solvent, vanishes. Within our model $`\mathrm{\Omega }_{sg}`$ is given by the expression in Eq. (23) with $`\rho _l`$ replaced by $`\rho _g`$, which is the density of the bulk phase. Since the index matching works for the bulk phase, it does not work for the wetting phase. As a consequence the wetting-induced forces appear against a background effective potential of hard spheres. Therefore for this class of colloidal suspensions the wetting phenomena discussed here are expected to have a pronounced effect on their phase behavior. Within our model, for *i*ndex-*m*atched suspensions the *total* effective interaction potential is given by
$$\mathrm{\Omega }_{tot,im}(D;R)=\mathrm{\Omega }_{tot}(D;R)(\mathrm{\Phi }(D;R)+\mathrm{\Omega }_{sg}(D;R))=\mathrm{\Omega }_S(D;R)\mathrm{\Omega }_{sg}(D;R)$$
(34)
and in analogy to $`\mathrm{\Omega }_E`$ and $`\mathrm{\Omega }_{E,tot}`$ we define
$$\mathrm{\Omega }_{E,im}(D;R)=\mathrm{\Omega }_{tot,im}(D;R)2\mathrm{\Omega }_{im}^{(1)}(R)$$
(35)
with $`\mathrm{\Omega }_{E,im}(D\mathrm{};R)=0`$ for the unbridged solutions. Figure 8 displays $`\mathrm{\Omega }_{tot,im}`$ and $`\mathrm{\Omega }_{E,im}`$ as function of $`a=D2R`$ for the same system as in Figs. 5 and 7. $`\mathrm{\Omega }_{sg}`$ is about $`30\%`$ smaller than $`\mathrm{\Omega }_S`$ for the unbridged solution and also approaches its asymptote $`2\mathrm{\Omega }_{sg}^{(1)}`$ from above. As before (see the discussion of Fig. 7 above), the resulting total effective interaction between spheres in an index-matched bulk fluid for the state with liquid bridge is still attractive and of the same order of magnitude as the bare dispersion interaction between the spheres, i.e., in the absence of the solvent.
## VI Summary
We have obtained the following main results:
1. Based on microscopic interaction potentials and within a simple version of density functional theory (Eqs. (7)–(11)) we have calculated the grand canonical potential of a system of two spheres immersed in a bulk fluid phase (Fig. 1). The microscopic interactions are chosen such that the spheres prefer the adsorption of a second fluid phase which is thermodynamically close to the bulk fluid phase. Accordingly, a single sphere immersed in the fluid is covered by a homogeneous wetting layer of this second phase of thickness $`l_0`$. These thin wetting layers covering the spheres lead to an effective wetting-induced interaction potential $`\mathrm{\Omega }_S(D)`$ between the spheres. We have systematically determined the dependence of $`\mathrm{\Omega }_S`$ on the distance $`D`$ between the spheres in terms of the morphology $`h(z)`$ of the wetting film enclosing the spheres (Eqs. (15)–(23)). We find that the shape of the effective interaction potential $`\mathrm{\Omega }_S(D)`$ depends, inter alia, on the effective interaction of two spheres immersed in the homogeneous *wetting* phase (Eq. (23)). This contribution, which is independent of $`h(z)`$, is not incorporated in previous phenomenological models for this system .
2. The equilibrium interfacial profiles of the wetting layers are determined numerically by minimizing the free energy functional $`\mathrm{\Omega }_S[h(z)]`$ in Eqs. (15)–(23). We have calculated the rich structure of these equilibrium profiles (Fig. 3) for spheres of radii $`R=20\sigma `$ (Fig. 2) and $`R=50\sigma `$ (Fig. 3) where $`\sigma `$ denotes the diameter of the solvent particles. As function of distance $`D`$, temperature $`T`$, and undersaturation $`\mathrm{\Delta }\mu `$ the system undergoes a first-order “bridging transition” between the two configurations shown in Fig. 1. For a fixed distance $`D`$ we have mapped out the phase diagram of bridging transitions in the $`T`$-$`\mathrm{\Delta }\mu `$ plane (Fig. 4). It turns out that the bridging transition differs from and to a large extent is independent of the thin-thick transition of the wetting layer on each single sphere which is a remnant of the prewetting transition on the corresponding flat substrate. Thus one has to distinguish between the prewetting line for a first-order wetting transition on a planar substrate, the thin-thick transition line for wetting on a single sphere, and the bridging transition line for two spheres (Fig. 4). At two-phase coexistence $`\mathrm{\Delta }\mu =0`$ and for $`R\sigma `$ the bridging transition is determined by the equality of the surface areas of the interfaces in the bridged and the unbridged configuration, leading to a universal ratio $`D_{bt}(\mathrm{\Delta }\mu =0)/(R+l_0)2.32`$ for the critical distance $`D_{bt}(\mathrm{\Delta }\mu =0)`$ of the bridging transition at coexistence (Fig. 6 and Subsec. IV B). Off coexistence $`D_{bt}(\mathrm{\Delta }\mu ,R)`$ is described by a universal scaling function (Eq. (30)).
3. At large distances and depending on the temperature relative to the thin-thick transition temperature on a single sphere the wetting-induced effective interaction potential can be either attractive or repulsive; in both cases it decays $`D^6`$ for large $`D`$. The bridging transition leads to a strong break in slope of the effective interaction potential at $`D=D_{bt}`$. This is the fingerprint of a cooperative phenomenon among the fluid particles whose degrees of freedom have been integrated out (see Sec. I). Metastable branches of the effective potential give rise to pronounced hysteresis effects (Fig. 5).
4. In the case that a bridge of the wetting phase connects the spheres (i.e., $`D<D_{bt}`$) there is an attractive wetting-induced interaction (Fig. 5) that pulls the spheres together. Within a wide range of separations $`a=D2R`$ of the spherical surfaces this force is of the same order of magnitude as the bare dispersion interaction potential $`\mathrm{\Phi }`$ (Eq. (31)) between the spheres. This bare interaction of two spheres (corresponding to the case that they are separated by vacuum) has to be added to the effective potential $`\mathrm{\Omega }_S`$ to yield the total interaction potential $`\mathrm{\Omega }_{tot}`$ between the spheres which is attractive at large distances (Eqs. (31), (33), and Fig. 7).
5. The wetting-induced force between spherical particles is experimentally accessible *directly* through suitable force microscopy techniques (Subsec. V B). Moreover, in Subsec. V D we argue that this force influences the phase behavior of *stericly* stabilized, index-matched colloidal suspensions. The total effective interaction potential for such a case is shown in Fig. 8; it is repulsive at large distances. The phase behavior of *charge* stabilized colloidal suspensions (Subsec. V C) is only affected by the wetting-induced interaction potential if the screening length of the Coulomb repulsion in the solvent is smaller than $`a_{bt}=D_{bt}2R0.32R`$. Depending on the size of the charges, the salt concentration of the solvent, and the underlying dispersion forces this criterion may be fulfilled even for colloidal particles whose radii are only a few tens of nm.
###### Acknowledgements.
We gratefully acknowledge financial support by the German Science Foundation within the special research initiative *Wetting and Structure Formation at Interfaces*.
## Contributions to the free energy
Our choice of interaction potentials $`\varphi (r)`$ (Eq. (1)) and $`\varphi _s(r)`$ (Eq. (5)) leads to the following expressions for the contributions to the free energy $`\mathrm{\Omega }_S`$ (with the thermodynamic limit already carried out):
$$\mathrm{\Omega }_{ei}[h(z)]=2\mathrm{\Delta }\rho \left(\rho _l\underset{0}{\overset{\mathrm{}}{}}𝑑z(g_+(z)+g_{}(z))\underset{0}{\overset{\mathrm{}}{}}𝑑z(f_+(z)+f_{}(z))\right)$$
(36)
with
$`g_\pm (z)`$ $`=`$ $`2w_0\sigma ^2{\displaystyle \frac{R_1}{\sigma }}w_0\sigma ^2\left({\displaystyle \frac{h^2(z)}{\sigma ^2}}+\left({\displaystyle \frac{z}{\sigma }}\pm {\displaystyle \frac{D}{2\sigma }}\right)^2{\displaystyle \frac{R_1^2}{\sigma ^2}}+1\right)`$ (39)
$`\times (\mathrm{arctan}(\sqrt{{\displaystyle \frac{h^2(z)}{\sigma ^2}}+\left({\displaystyle \frac{z}{\sigma }}\pm {\displaystyle \frac{D}{2\sigma }}\right)^2}+{\displaystyle \frac{R_1}{\sigma }})`$
$`\mathrm{arctan}(\sqrt{{\displaystyle \frac{h^2(z)}{\sigma ^2}}+\left({\displaystyle \frac{z}{\sigma }}\pm {\displaystyle \frac{D}{2\sigma }}\right)^2}{\displaystyle \frac{R_1}{\sigma }}))`$
where $`R_1=R+d_s`$ and
$`f_\pm (z)`$ $`=`$ $`{\displaystyle \frac{\pi u_9}{4}}\left({\displaystyle \frac{1}{7}}\left({\displaystyle \frac{1}{(k_\pm +R)^7}}{\displaystyle \frac{1}{(k_\pm R)^7}}\right)+R\left({\displaystyle \frac{1}{(k_\pm +R)^8}}+{\displaystyle \frac{1}{(k_\pm R)^8}}\right)\right)`$ (41)
$`\pi u_3\left({\displaystyle \frac{1}{k_\pm +R}}{\displaystyle \frac{1}{k_\pm R}}+R\left({\displaystyle \frac{1}{(k_\pm +R)^2}}+{\displaystyle \frac{1}{(k_\pm R)^2}}\right)\right)`$
with $`k_\pm =\sqrt{h^2(z)+(z\pm D/2)^2}`$. The contribution $`\mathrm{\Omega }_{lg}`$ is given by
$`\mathrm{\Omega }_{lg}[h(z)]`$ $`=`$ $`w_0\sigma ^3(\mathrm{\Delta }\rho )^2{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}dz{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}dz^{}({\displaystyle \frac{1}{y_{}^4}}({\displaystyle \frac{qy_{}^6+y_{}^4(2q^2+p^2)+3p^2qy_{}^2+p^4}{(p^2+2y_{}^2q+y_{}^4)^{3/2}}}p)`$ (43)
$`+{\displaystyle \frac{1}{y_+^4}}({\displaystyle \frac{qy_+^6+y_+^4(2q^2+p^2)+3p^2qy_+^2+p^4}{(p^2+2y_+^2q+y_+^4)^{3/2}}}p))`$
where the abbreviations $`y_\pm `$, $`p`$, and $`q`$ are defined by
$$y_\pm ^2=\sigma ^2+(z\pm z^{})^2,$$
(44)
$$p=h^2(z)h^2(z^{}),$$
(45)
and
$$q=h^2(z)+h^2(z^{}).$$
(46)
The double integral in Eq. (43) demonstrates the nonlocal functional dependence of $`\mathrm{\Omega }_{lg}`$ on $`h(z)`$.
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# 1 Introduction
## 1 Introduction
During the past couple of years microlensing proved to be a new and potentially very powerful tool of modern astrophysics. Originally proposed by Paczyński (1986, 1991) as a method of searching for dark matter in the Galaxy it has been used for such different applications as searching for planets, determination of parameters of stellar atmospheres, studies of Galactic structure and many others.
After the original reports on discovery of the first cases of microlensing in September 1993 (toward the LMC: MACHO survey – Alcock et al. 1993, EROS survey – Aubourg et al. 1993; toward the Galactic bulge: OGLE survey – Udalski et al. 1993) much observing work was done to convince the astronomical community on the potentials of the newly discovered class of events. Soon more cases of classical microlensing in the Galactic bulge were announced (Udalski et al. 1994a, Alcock et al. 1995a), first cases of ”exotic microlensing” like binary microlensing (Udalski et al. 1994b) or events with parallax effect (Alcock et al. 1995b) were found. Also first estimates of the observed optical depth to the Galactic bulge were published (Udalski et al. 1994c, Alcock et al. 1995a, 1997a) indicating that it is much larger than that predicted from modeling. Many theoretical interpretations of these intriguing results followed (see review by Paczyński 1996). First interpretation of results of observations toward the LMC was also published (Alcock et al. 1996a).
Another important step was development and implementation of the so called alert systems (the Early Warning System, EWS, of the OGLE survey – Udalski et al. 1994d and MACHO Alert system – Alcock et al. 1996b) which allow to detect the microlensing phenomena when an event is in progress. This step changed the observing strategy of microlensing surveys. The ability of detection in real time made it possible to perform follow-up observations, both photometric and spectroscopic, of many events. New kind of microlensing studies, ”follow-up” projects concentrated on high time resolution observations of events discovered by survey projects, were formed (e.g., PLANET – Albrow et al. 1998; MPS – Rhie et al. 1999). From 1995 on, the vast majority of microlensing events have been detected by the alert systems.Information on microlensing events in progress can be found:
OGLE-EWS: http://www.astrouw.edu.pl/~ogle
MACHO (1995–1999): http://darkstar.astro.washington.edu/
EROS: http://www-dapnia.cea.fr/Spp/Experiences/EROS/alertes.html
PLANET: http://www.astro.rug.nl/~planet/
MPS: http://bustard.phys.nd.edu/MPS/
The microlensing field of astrophysics matured rapidly and entered the second phase of extensive observations to largely increase statistic of collected microlensing events. Both EROS and OGLE surveys considerably increased their observing capabilities in 1996. First microlensing events were discovered in other lines of sight – toward the SMC (MACHO – Alcock et al. 1997b) and in the Galactic disk (EROS – Derue et al. 1999) even as far as $`70^{}`$ from the Galactic center (OGLE – Mao 1999). Many interesting results were published by follow-up teams (PLANET – Albrow et al. 2000a,b, MPS – Bennett et al. 1999a). Up to now one can estimate the total number of registered microlensing events to be about 500.
In this paper we present the catalog of microlensing event candidates detected during the second phase of the Optical Gravitational Lensing Experiment (OGLE-II) in three observing seasons 1997–1999. Although during the OGLE-II phase observations are conducted in many lines of sight only two events were discovered in the directions other than the Galactic bulge so far.1999-CAR-01 and 1999-LMC-01 see http://www.astrouw.edu.pl/~ogle for more information on these events. Therefore we decided to limit our catalog to the Galactic bulge events only. Presented microlensing events were extracted from the OGLE-II photometric databases with a technique similar to that applied in the EWS alert system. Great care was paid to achieve the highest possible completeness of the catalog. The Catalog comprises 214 cases of microlensing toward the Galactic bulge.
The main goal of the paper is to provide the astronomical microlensing community with the large set of photometric data of microlensing events for further analysis. Photometry of all objects is available from the OGLE Internet archive (see Section 5). Beside of the typical single mass microlensing cases the Catalog includes several cases of ”exotic microlensing”, like binary microlensing events etc.
The sample of presented events is already large enough that in spite of possible incompleteness the distribution of microlensing parameters can be studied. In particular we present the first preliminary distribution of the rate of microlensing in different parts of the Galactic bulge. Studies of such a distribution can provide important constraints on the origin of the Galactic bulge microlensing (bar vs. disk) and in general on the structure of the Galaxy.
## 2 Observational Data
All observations presented in this paper were carried out during the second phase of the OGLE experiment with the 1.3-m Warsaw telescope at the Las Campanas Observatory, Chile, which is operated by the Carnegie Institution of Washington. The telescope was equipped with a ”first generation” camera with a SITe $`2048\times 2048`$ back-illuminated CCD detector working in drift-scan mode. The pixel size was 24 $`\mu `$m giving the scale of 0.417 arcsec/pixel. Observations of the Galactic bulge were performed in the ”medium” speed reading mode of CCD detector with the gain 7.1 e<sup>-</sup>/ADU and readout noise of about 6.3 e<sup>-</sup>. Details of the instrumentation setup can be found in Udalski, Kubiak and Szymański (1997).
Regular observations of the Galactic bulge started on March 23, 1997. The observing season of the Galactic bulge lasts typically from mid February up to the end of October. In this paper we present microlensing events detected during the seasons 1997–1999.
Forty nine driftscan fields were observed, each covering $`14.2\times 57`$ arcmins on the sky, giving total area of about 11 square degrees. Forty seven of these fields were monitored frequently for detection of microlensing events. The two remaining fields (BUL$`\mathrm{\_}`$SC45 and BUL$`\mathrm{\_}`$SC46) were observed only from time to time, mostly for maintaining phasing of variable stars discovered during the OGLE-I phase of the OGLE survey. Fields BUL$`\mathrm{\_}`$SC47, BUL$`\mathrm{\_}`$SC48 and BUL$`\mathrm{\_}`$SC49 were added to the list of targets in the second observing season, so observations of these fields span seasons 1998 and 1999 only.
Table 1 lists the equatorial and Galactic coordinates of the center of each field and its acronym. Also the number of collected frames is provided. Fig. 1 shows schematically location of the Galactic fields in ($`l`$, $`b`$) and (RA, DEC) coordinate systems. Most fields are located along the regions of smaller interstellar extinction, i.e., large stellar density at $`b3\text{.}^{}5`$ and spanning large range of the Galactic longitude – from $`l=11^{}`$ to $`l=+11^{}`$. A few fields are located on the other side of the Galactic equator at positive $`b`$. In general the interstellar extinction there is much larger making selection of dense stellar fields suitable for microlensing search much more difficult.
Observations of the Galactic bulge fields were obtained in the standard BVI-bands. The effective exposure time was 87, 124 and 162 seconds for the I, V and B-band, respectively. The vast majority of observations were done through the I-band filter while images on only several epochs were collected in the BV-bands. The instrumental system closely resembles the standard BVI. All photometric data of the Galactic bulge will be ultimately calibrated and released in the form of photometric maps (similar to the SMC data: Udalski et al. 1998). In this paper we present only I-band observations. The data are only preliminarily calibrated but uncertainty of the zero points of photometry should not be larger than $`\pm 0.05`$ mag. The median seeing of the entire set of several thousand collected frames of the Galactic bulge is equal to 1.<sup>′′</sup>29.
## 3 Search Algorithm
The Galactic bulge photometric data collected during the OGLE-II survey are naturally divided into separate seasons. We used the I-band photometric data from each season to create databases of non-variable stars in every field. Then we compared brightness of selected non-variable stars with that in other seasons, looking for objects which vary in brightness. Databases of non-variable stars were created for each season, i.e., 1997, 1998 and 1999 separately, and the search for variable objects was performed independently with all possible combinations of ”non-variable” and ”variable” seasons. In this manner we were able to detect inter-seasonal events and verify most candidates by comparison of results in a given season obtained with databases of non-variable stars from two other seasons.
The star entered the database of non-variable objects in a given season when the standard deviation of all observations collected during that season was smaller than the maximum standard deviation of non-variable stars at given magnitude. The latter was determined by analysis of the distribution of the brightness standard deviation of all observed stars in the field. We limited the databases to objects brighter than $`I=19.5`$ mag. The total number of non-variable stars searched for microlensing phenomena in our fields was more than 20.5 millions (about $`410^9`$ photometric measurements).
Comparison of number of stars which were included into databases of non-variable stars in different seasons provided information on the actual number of real stellar objects. They were similar to within 5% indicating that only a small fraction of stars in these databases are spurious objects like artifacts in saturated regions of bright stars etc.
Candidates for variable objects were selected by comparison of their photometry with the mean magnitude in the template ”non-variable” season. When five consecutive observations were brighter than $`3\sigma `$ limit over the ”non-variable” magnitude, such an object was marked for further analysis. The typical $`3\sigma `$ threshold was equal to 0.06 mag for a star of $`I<16`$ mag, 0.20 mag for $`I=18`$ mag and 0.60 mag for the limiting $`I=19.5`$ mag. At this stage we decided to visually inspect light curves of all such candidates to avoid removing non-standard events like for instance caustic crossing binary events. All candidates which displayed more than one episode of brightening (non consistent with microlensing type), turned out to be evident long-term variable stars and those which evidently exhibited spurious variability were removed from the list of candidates. In the next step images taken at the epochs of normal (constant) brightness and near the maximum of brightening were inspected to verify whether the brightening is real. In many cases the brightening was spurious due to proximity of large amplitude variable star etc. Such objects were also removed. We did not checked achromaticity of the brightening. Because of sparse coverage in the bands other than I it would be practically impossible to perform that test for most of our candidates. On the other hand the experience of the OGLE-I phase as well as experience of other teams indicate that in the case of stellar population in the Galactic bulge (contrary to the Magellanic Clouds) there are practically no objects which would mimic microlensing event light curve. Thus, the achromaticity is not that important to extract the real microlensing events from the background of variable stars. Moreover, in the very dense stellar fields of the Galactic bulge blending can cause quite severe deviations from achromaticity making the latter condition not very strong.
We were left with a list of 214 candidates for microlensing events which passed all our tests. They are shown and analyzed in the next Sections.
## 4 Catalog of Microlensing Events
The Catalog of Microlensing Events in the Galactic Bulge consists of two parts: the table with basic data of the microlensing event and coordinates of the lensed object and the atlas presenting photometry of the event.
Table 2 lists all microlensing event candidates presented in the Catalog. The first column contains the ID of the lensed object in regular OGLE convention i.e., consisting of the field name and database number. In the next column we provide cross identification with the EWS name when the microlensing event was detected in real time. The two next columns contain equatorial coordinates (J2000) of the lensed star followed by parameters of the microlensing event. The latter were calculated assuming point mass microlensing and include the moment of maximum brightness (HJD–2450000), $`T_0`$, time of the Einstein radius crossing, $`t_0`$, magnification at maximum brightness, $`A_{\mathrm{max}}`$, and I-band baseline magnitude. The final column of Table 2 includes remarks on non-standard events. In particular we note there candidates for binary microlensing. When such a candidate is a typical caustic crossing binary event we do not provide single point mass microlensing fit parameters because they are meaningless in such cases and modeling of light curve requires much more complicated approach. In weak cases of possible binary microlensing single mass parameters are provided. Abbreviations in the remark column have the following meaning: bcc – caustic crossing binary microlensing; bin – possible binary microlensing; dls – binary microlensing or double source star; ble – severely blended event; var – possible variable star.
The atlas of light curves of microlensing events in the Galactic bulge detected during the seasons 1997-1999 is presented in Appendix. We do not provide there finding charts, but they are available electronically from the OGLE Internet archive. Also photometric data of all presented events can be retrieved from there.
## 5 Discussion
Two hundred fourteen microlensing event candidates were detected in the OGLE-II fields during the observing seasons 1997–1999. Significant number of events from 1998–1999 seasons, i.e., when the EWS alert system was implemented, were detected in real time. Many of these events were then followed up by other groups for detailed study (e.g., 1998-BUL-14, Albrow et al. 2000a). Because part of the OGLE-II fields overlap with those observed by the MACHO team several events were discovered by both teams. Cross-identification of those discovered by alert systems can be found on the WWW alert pages of both teams while for the remaining objects (in particular all 1997 season events) it can be done using accurate equatorial coordinates provided in this paper and on the MACHO WWW alert page.
The sample of OGLE-II microlensing events contains several cases of evident lensing caused by binary object. The sample of characteristic caustic crossing events consists of 14 events. In five additional cases it is possible that the observed light variations are also caused by binary lens. In two cases, BUL$`\mathrm{\_}`$SC39 259656 and BUL$`\mathrm{\_}`$SC40 434222, two separate microlensing like episodes were detected. These cases might be explained either as due to binary microlensing without caustic crossing or double lensed star. In the latter case the time scales of both episodes should be equal what can discriminate between those two possibilities. Time scales of BUL$`\mathrm{\_}`$SC39 259656 episodes are equal to about 15 and 27 days suggesting that this case is a binary microlensing. The second brightening episode of BUL$`\mathrm{\_}`$SC40 434222 started at the end of 1999 season but observations collected up to the moment of writing this paper confirm that it can be of microlensing origin (see OGLE EWS WWW page). In this case the time scales are within errors the same ($`t_0130`$ days) suggesting double source star. However, the second episode is still in the rising part of its light curve.
The rate of the binary caustic crossing microlensing toward the Galactic bulge observed in our sample of all microlensing events is equal to $`6.5`$% while that of all binary microlensing $`9.3`$%. It is worth noting that this is much more than reported by the MACHO team (Alcock et al. 1999). On the other hand Mao and Paczyński (1991) in their classical paper on binary microlensing predicted about 10% rate of strong binary microlensing events. The agreement between the predicted and observed rates is very good.
Among many interesting cases of microlensing presented in the Catalog we draw attention to the object, BUL$`\mathrm{\_}`$SC5 244353. Unfortunately, only falling part of its light curve was covered during the entire observation period. While the shape of the light curve strongly resembles that of microlensing event, it cannot be excluded that the object is a variable star. If, however, its brightening was caused by microlensing, the time-scale of the event was very long. The event would closely resemble those reported by Bennett et al. (1999b) which are supposed to be caused by black hole lenses.
The statistic of collected events toward the Galactic bulge is becoming significantly larger after each observing season. While an accurate analysis requires additional information like efficiency of detection, the number of presented microlensing events is already so large that interesting correlations can be shown.
Efficiency of detection of microlensing events depends at least on two factors: time sampling of the microlensing light curve and brightness of the lensed star. Typical sampling of one/two observations per night limits detection to events with sufficiently long time-scale. For fainter lensed stars the $`3\sigma `$ detection threshold is larger, therefore only larger magnification events can be triggered. In Fig. 2 we plot distribution of the I-band baseline magnitude of lensed stars from our sample. The bins are 0.3 mag wide. Dotted line presents the distribution of all events while the bold solid line the one of subsample of events with magnification $`A_{\mathrm{max}}>1.5`$. The shape of both distributions indicate that the entire sample of microlensing events is reasonably complete up to $`I18.0`$ mag while subsample of events with $`A_{\mathrm{max}}>1.5`$ up to $`I18.8`$ mag.
Selection effect due to brightness of the lensed star can also be assessed from Fig. 3 which shows magnification at the maximum,$`A_{\mathrm{max}}`$, and minimum impact parameter, $`u_{\mathrm{min}}`$, plotted against the I-band baseline magnitude. It is clearly seen from Fig. 3 that our sample is quite complete for stars of $`I<18`$ mag and $`A_{\mathrm{max}}>1.3`$, i.e., $`u<1.0`$. As can be expected for fainter stars only higher magnification events were triggered. Nevertheless at $`I=19`$ mag the limit of reasonable completeness is still $`A_{\mathrm{max}}1.5`$ ($`u_{\mathrm{min}}0.8`$).
Fig. 4 presents distribution of maximum magnification of our entire sample. This parameter ranges from as low as 1.1 up to about 50 (in a few cases the single point mass model predicts extremely large magnifications but because of lack of observations at the very maximum these values are not reliable. In these cases we provide lower limit of magnification resulting from the brightest observation). Below $`A_{\mathrm{max}}=1.3`$ incompleteness of our sample is larger – only for brighter stars so low magnification events could be detected. However, the completeness becomes much higher for events with $`A_{\mathrm{max}}>1.5`$. The number of events is gradually falling from $`A_{\mathrm{max}}1.7`$ to $`A_{\mathrm{max}}=7`$ with a long tail of single events with higher $`A_{\mathrm{max}}`$.
Fig. 5 shows the distribution of the Einstein radius crossing time, $`t_0`$, for our sample of microlensing events. This parameter is also likely to be affected by incompleteness resulting from the sampling of the light curve. Although we have not performed yet detailed analysis of the dependence of detection efficiency on the event time-scale for our entire data set, preliminary tests performed on 1997 databases of constant stars in similar manner as in Udalski et al. (1994c) indicate that the region of reasonable efficiency is extended toward shorter time scales as compared to the OGLE-I phase. For events with $`t_0>8`$ days efficiency of detection becomes relatively flat, so we may expect that the distribution of $`t_0`$ is also relatively complete for events longer than that limit. The distribution of $`t_0`$ peaks at $`t_017`$ days with a long tail of longer time-scale events. One should also note a small excess of events with $`t_050`$ days. It was marginally seen in the MACHO data (Alcock et al. 1997a). The time-scale is plotted as a function of the I-band baseline magnitude in Fig. 6. No evident correlation is seen in this plot.
The spatial distribution of events with different time-scales is presented in Fig. 7. The total sample was divided into three sub-samples of short ($`t_0<20`$ days), medium ($`20<t_0<40`$ days) and long time-scale ($`t_0>40`$ days) events. Location on the sky of these three samples in the Galactic coordinates ($`l,b`$) is plotted with different symbols in Fig. 7. If the different time-scale events were to be caused by different populations of lensing objects one could expect some differences in the distribution of our three sub-samples on the sky. However, this does not seem to be true – the distribution of all our sub-samples is rather similar. One should be, however, aware of still small statistic of events in the regions at larger $`|l|`$.
Table 3 lists the average number of stars searched for microlensing events in each of the Galactic bulge fields. In the second column the number of detected events in each field is provided. Because the number of searched stars is different by a factor of more than four in our fields, we normalized the observed number of microlensing events in each field to one million stars. Normalized number of events is listed in the last column of Table 3.
Finally, because we typically observe two slightly overlapping driftscan fields in a given part of the Galactic bulge (see Fig. 1), we averaged the normalized numbers of events in such adjacent fields to increase statistic. Table 4 lists our 24 lines of sight in the Galactic bulge with their Galactic coordinates and the average number of observed microlensing events per one million stars during seasons 1997–1999 (the total span of presented observations is about 940 days, i.e., 2.58 years).
Fig. 8 presents the number of events per one million stars observed in 24 lines of sight in the Galactic bulge in years 1997–1999. The number at given ($`l`$,$`b`$) indicates the normalized number of events observed in a given direction. Of course, one should be aware that the presented numbers are not the true optical depth. Accurate estimate of the optical depth would require precise determination of detection efficiency, assessment of the blending effect and contribution of the Galactic disk stars to the stars searched for microlensing events. On the other hand one can expect that these factors are in the first approximation similar in so uniformly observed fields. Also the spatial distribution of the time-scale of events is similar. Therefore the normalized number of events is likely a crude approximation of the optical depth.
As one can expect the number of events is a strong function of the Galactic latitude. It falls by almost an order of magnitude when $`b`$ changes from $`b=1\text{.}^{}3`$ to $`b=6^{}`$ (at $`l0^{}`$). It also changes with the Galactic longitude. The numbers of events observed in the fields at $`b3\text{.}^{}5`$ indicate that there is a clear dependence on $`l`$ – the number of events falls by a factor of 2–4 at $`|l|10^{}`$ as compared to $`l=0^{}`$. There is a noticable asymmetry with larger number of events at negative $`l`$. The clear dependence of the number of microlensing events on $`l`$ strongly suggests that the majority of microlensing events are caused by lenses located in the Galactic bar, inclined to the line of sight toward the Galactic center, rather than in the Galactic disk.
At positive $`b`$ the number of microlensing events in the line of sight located at $`l=0\text{.}^{}4,b=3^{}`$ is consistent with that observed at negative $`b`$ but one can notice possible excess of events at $`l=5\text{.}^{}3,b=2\text{.}^{}7`$.
Looking at the numbers presented in Fig. 8 and Table 4 one should be aware of the possible bias resulting from somewhat different time sampling of some fields. While the numbers of observations in most fields are similar there are six fields with about 40% larger number of analyzed frames (see Table 1). These were the fields located closest to the Galactic center and they were observed with frequency of 3–5 observations per night during a part of the 1997 season for ultra-short time events. No such events (time scale $`t_0<2`$ days) were, however, found. Nevertheless, the number of detected microlensing events can be larger in more frequently observed fields resulting in somewhat overestimated ratios between the number of microlensing events close to the Galactic center and in other directions.
We estimated the possible magnitude of this effect by comparison of the median time scale, $`t_0`$, of events from the least-, medium- and most-frequently observed fields. One could expect that for the most frequently sampled fields, the efficiency of detection is much higher for short time scale events and the median $`t_0`$ of that subsample should be significantly shorter than that of the remaining subsamples if the bias is strong. However, the median $`t_0`$ is equal to 17.0, 21.5 and 19.7 days for the most-, medium- and least-frequently sampled fields, respectively. The differences are small indicating that the bias is also small and can be in the first approximation neglected. It would be certainly much stronger if much larger part of events had very short time scales.
It is interesting to compare our empirical results with the modeling predictions. We used the map of the optical depth for the Galactic bar calculated by Stanek et al. (1997) to derive the optical depth in our lines of sight. Contribution of the Galactic disk in these directions was taken from Kiraga (1994) and added to the contribution of the Galactic bar. Fig. 9 shows the relation between the model optical depth and the observed number of microlensing events per one million stars in our lines of sight. The correlation between both values is clearly seen. The slope of the relation is equal to $`0.86\pm 0.11`$ giving a crude calibration between the observed numbers and the optical depth. However, one should also note the potential problem – the linear relation does not cross the (0,0) point. This may indicate that the ordinate in Fig. 9 is not that simply related with the optical depth and the observed normalized number of events is a function not only of the optical depth but also of an additional factor like, for instance, spatial dependent efficiency of detection. Another possibility is that the model underestimate significantly the optical depth close to the Galactic center. In general, however, Figs. 8 and 9 indicate that models of the bar might provide a reasonable approximation of the Galactic bulge structure (see also for example Evans 1994, Zhao and Mao 1996, Grenacher et al. 1999) and that the microlensing will be a very powerful tool in further constraining the Galactic bulge properties when significantly larger statistic of events is collected.
While Fig. 8 can provide a first outlook on the distribution of microlensing events in the Galactic bulge one should be aware that this is only the first approximation. The statistic of microlensing events, in particular in the fields with larger $`|l|`$, is still small. Observations with the same set-up will be continued during the next observing season (2000) providing about 60–70 new cases of microlensing events. After that a large instrumental upgrade to the OGLE-III phase is planned by implementation of a new mosaic CCD camera. This will increase the number of discovered events by a factor of 3–5 leading to fast increase of the statistic of Galactic bulge events. Also larger area of the Galactic bulge will be monitored. It is also planned to reanalyze the OGLE-II photometric data with new techniques like image subtraction method. It would allow to get rid of blending effect uncertainties. By limiting to well defined population of the Galactic bulge stars like for instance red clump giants it will be possible to provide much more accurate information on the optical depth distribution in the Galactic bulge in the near future.
The Catalog of Microlensing Events in the Galactic Bulge and all photometric data presented in this paper are available now to the astronomical community from the OGLE Internet archive:
http://www.astrouw.edu.pl/~ogle
ftp://sirius.astrouw.edu.pl/ogle/ogle2/microlensing/gb/
or its US mirror
http://www.astro.princeton.edu/~ogle
ftp://astro.princeton.edu/ogle/ogle2/microlensing/gb/
Acknowledgements. We would like to thank Prof. Bohdan Paczyński for many discussions and help at all stages of the OGLE project. The paper was partly supported by the Polish KBN grants 2P03D00814 to A. Udalski, 2P03D00916 to M. Szymański, and 2P03D00717 to K. Żebruń. Partial support for the OGLE project was provided with the NSF grant AST-9820314 to B. Paczyński.
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# Einstein manifolds of non-negative sectional curvature and entropy
## 1. Introduction
Let $`(M^n,g)`$ be a complete smooth connected Riemannian manifold. One of our main concerns here will be the study of Einstein metrics $`g`$ of non-negative sectional curvature. Unless $`(M,g)`$ is flat, the Ricci curvature $`r`$ must be positive and we might as well rescale the metric so that $`r=g`$. Therefore all the sectional curvatures satisfy $`0K(P)1`$. Myers’ theorem asserts that $`M`$ is compact, has finite fundamental group and its diameter is $`\sqrt{n1}\pi `$. By passing to the universal covering we shall assume in the sequel that $`M`$ is in fact closed and simply connected.
When $`M`$ is four-dimensional, strong topological restrictions to the existence of such metrics were first given by N. Hitchin . He proved that the Euler characteristic $`\chi `$ and the signature $`\tau `$ of $`M`$ must satisfy
$$\chi \left(\frac{3}{2}\right)^{3/2}|\tau |.$$
This inequality was recently improved by M.J. Gursky and C. LeBrun who showed that in fact one has $`9\chi >15/4|\tau |`$. This result together with Freedman’s classification theorem tells us that there are at most 12 homeomorphism types (which can be explicitly listed) of simply connected closed Einstein manifolds of non-negative sectional curvature. If one assumes further that $`M`$ has positive definite intersection form, then Gursky and LeBrun proved that $`(M,g)`$ must, up to isometry, be $`P^2`$ equipped with a constant multiple of the Fubini-Study metric. All these results rely heavily on the four-dimensional Gauss-Bonnet theorem.
Even though it is believed that simply connected closed Einstein manifolds with non-negative sectional curvature are very rare, there are basically no known topological obstructions to the existence of such metrics in dimensions $`5`$, except, of course, Gromov’s celebrated result in which gives an obstruction to the existence of metrics just with non-negative sectional curvature. Gromov’s result asserts that if $`M`$ admits such a metric, then the sum of the Betti numbers $`b_i`$ must verify:
$$\text{dim}H_{}(M,)=\underset{i=0}{\overset{n}{}}b_iC(n),$$
where
$$C(n):=((n+1)J(n))^{100^n},J(n):=2^{M(n)},M(n)=8^n10^{n^2+4n}.$$
Clearly, the bound $`C(n)`$ is quite unrealistic, and Gromov himself conjectured that $`C(n)`$ should be just $`2^n`$.
We now state our first main result. Let $`R`$ be the radius of convergence of the series
$$\underset{i2}{}\text{dim}(\pi _i(M))t^i.$$
Theorem A. Let $`(M^n,g)`$ be a smooth simply connected closed manifold of positive Ricci curvature. Suppose that we normalize $`g`$ so that $`rg`$ and let $`k:=\mathrm{max}_PK(P)`$. Then
$$\mathrm{log}R\frac{\pi \sqrt{n1}}{2}\left((n1)\sqrt{k}\frac{1}{\sqrt{k}}\right).$$
In particular, if $`(M^n,g)`$ is a simply connected closed Einstein manifold of non-negative sectional curvature, then
$$\mathrm{log}R\frac{\pi \sqrt{n1}(n2)}{2}.$$
We remark that Theorem A also holds if we just assume that $`M`$ has finite fundamental group. Also note that the second inequality in the theorem (as well as all the corollaries below) can also be obtained just assuming the existence of a metric $`g`$ with $`rg`$ and $`K(P)1`$ for all $`P`$.
A simply connected closed manifold $`M`$ is said to be rationally elliptic if the rational homotopy $`\pi _{}(M)`$ is finite dimensional, i.e. there exists a positive integer $`i_0`$ such that for all $`ii_0`$, $`\pi _i(M)=0`$. The manifold $`M`$ is said to be rationally hyperbolic if it is not rationally elliptic (cf. and references therein).
It was proved by Y. Félix and S. Halperin that if $`M`$ is rationally hyperbolic, then the integers $`\rho _i=_{ji}\text{dim}\pi _j(M)`$ grow exponentially in $`i`$ (i.e. there exist $`C>1`$ and a positive integer $`k`$ such that if $`i>k`$ then $`\rho _iC^i`$). Hence if $`M`$ is rationally elliptic we have $`R=+\mathrm{}`$ and if $`M`$ is rationally hyperbolic $`R<1`$. It follows that Theorem A is meaningful only if $`M`$ is rationally hyperbolic.
The “generic” manifold is rationally hyperbolic; rational ellipticity is a severely restrictive condition. Examples of rationally elliptic manifolds are simply connected homogeneous spaces , manifolds that admit a codimension one compact action , Dupin hypersurfaces and any known manifold that admits a Riemannian metric of non-negative sectional curvature. A conjecture attributed to R. Bott states that any compact simply connected manifold that admits a metric of non-negative sectional curvature must be rationally elliptic (cf. ). It is known that if $`M`$ is rationally elliptic then $`\text{dim}H_{}(M,)2^n`$ and hence Bott’s conjecture implies Gromov’s conjecture on the optimal bound for the sum of the Betti numbers of $`M`$.
A manifold $`M`$ is said to be formal if there exists a morphism of differential graded algebras from the minimal model of $`M`$ to $`(H^{}(M,),0)`$ that induces an isomorphism in cohomology. The interest for this class of spaces lies in the fact that for them all the rational homotopy invariants of $`M`$ can be obtained from $`H^{}(M,)`$. In , it is shown that compact simply connected Kähler manifolds are formal. Also, any manifold with dimension $`6`$ is formal and, more generally, any $`p1`$-connected manifold of dimension $`4p2`$ is formal .
In Section 2 we shall explain how the results of Y. Félix and J.C. Thomas in combined with Theorem A yield the following two corollaries:
Corollary 1. Let $`(M^n,g)`$ be a smooth simply connected closed Einstein manifold of non-negative sectional curvature. If $`M`$ is formal we have:
$$\text{dim}H_{}(M,)\left[1+\mathrm{exp}\left(\frac{\pi \sqrt{n1}(n2)}{2}\right)\right]^n.$$
Corollary 2. Let $`(M^n,g)`$ be a smooth simply connected closed Einstein manifold of non-negative sectional curvature. Suppose that $`M`$ is formal and $`(p1)`$-connected. Then
$$b_p\frac{n}{p}\mathrm{exp}\left(\frac{p\pi \sqrt{n1}(n2)}{2}\right).$$
From Corollary 2, we obtain right away:
Corollary 3. Let $`(M^5,g)`$ be a smooth simply connected closed Einstein manifold of non-negative sectional curvature. Then
$$b_2\frac{5}{2}\mathrm{exp}\left(6\pi \right).$$
Corollary 3 implies, for instance, that the connected sum of $`k`$ copies of $`S^3\times S^2`$ cannot admit an Einstein metric of non-negative sectional curvature for $`k>\frac{5}{2}\mathrm{exp}\left(6\pi \right)`$.
Certainly the bounds in Corollaries 1 and 3 are far better than Gromov’s bound. On the other hand our bounds are not as good as those obtained by Gursky and LeBrun in for four manifolds. The best lower bound for $`\mathrm{log}R`$ in dimension four has been obtained by I. Babenko in . He proved that
$$1/R\frac{b_2+\sqrt{b_2^24}}{2}.$$
Combining this bound with Theorem A we obtain:
$$\frac{b_2+\sqrt{b_2^24}}{2}\mathrm{exp}(\pi \sqrt{3}).$$
This implies that $`b_2230`$ while Gursky and LeBrun’s result implies that $`b_27`$.
The proof of Theorem A is based on ideas first introduced by M. Berger and R. Bott in and further developed by N. Grossman in . Given $`x`$ and $`y`$ in $`M`$ and $`T>0`$, define $`n_T(x,y)`$ as the number of geodesic arcs joining $`x`$ and $`y`$ with length $`T`$. For each $`T>0`$, the counting function $`n_T(x,y)`$ is finite and locally constant on an open full measure subset of $`M\times M`$, and integrable on $`M\times M`$. If $`g`$ has positive Ricci curvature, i.e., $`r\delta g`$ with $`\delta >0`$, then we shall see in Section 2 that the Morse theory of the loop space yields
(1)
$$\frac{\sqrt{\delta }\mathrm{log}R}{\pi \sqrt{n1}}\underset{T+\mathrm{}}{lim\; sup}\frac{1}{T}\mathrm{log}_Mn_T(x,y)𝑑y,$$
for any point $`xM`$. G.P. Paternain explained in how Yomdin’s theorem can be used to prove that
(2)
$$\underset{T+\mathrm{}}{lim\; sup}\frac{1}{T}\mathrm{log}_Mn_T(x,y)𝑑yh_{top}(g),$$
where $`h_{top}(g)`$ denotes the topological entropy of the geodesic flow of $`g`$. Combining (1) and (2) yields:
(3)
$$\frac{\sqrt{\delta }\mathrm{log}R}{\pi \sqrt{n1}}h_{top}(g).$$
Inequality (3) is also pointed out by I. Babenko in \[1, Theorem 4\]. While Berger, Bott and Grossman estimated the average counting function by using Rauch’s comparision theorem we shall use dynamical ideas to estimate $`h_{top}(g)`$ from above. As we shall see below our bound is better than those of Berger, Bott and Grossman and does not involve any lower bound on the sectional curvature. In Section 3 we shall prove the following result which will imply Theorem A and has independent interest.
Theorem B. Let $`(M^n,g)`$ be a closed Riemannian manifold and let $`K_{max}`$ be a positive upper bound for the sectional curvature. Then
$$h_{top}(g)\frac{n1}{2}\sqrt{K_{max}}\frac{\mathrm{min}_{vSM}r(v)}{2\sqrt{K_{max}}},$$
where $`SM`$ is the unit sphere bundle of $`M`$ and $`r(v)`$ is the Ricci curvature in the direction of $`vSM`$.
Let $`k`$ be a positive number such that $`|K(P)|k`$ for all 2-planes $`P`$. Then, clearly $`r(n1)kg`$ and hence Theorem B gives
$$h_{top}(g)\frac{n1}{2}\sqrt{k}+\frac{n1}{2}\sqrt{k}=(n1)\sqrt{k}.$$
The latter inequality, which is certainly weaker than that of Theorem B, was first proved by A. Manning in . The upper bound for $`h_{top}(g)`$ in Theorem B is probably the best that one can obtain in terms of the $`C^0`$-norm of the curvature tensor. Note that the bound becomes sharp for all the space forms (Manning’s bound is not sharp for positively curved space forms). Also observe that there is no hope to obtain an upper bound for $`h_{top}(g)`$ purely in terms of bounds of the Ricci curvature. Indeed, the $`K3`$ surfaces admit Ricci flat metrics, but $`h_{top}(g)>0`$ for any $`C^{\mathrm{}}`$ metric $`g`$ since a $`K3`$ surface is rationally hyperbolic (any simply connected four-manifold with second Betti number strictly bigger than two is rationally hyperbolic).
We conclude this introduction by comparing the upper bound in Theorem B with the bounds obtained by Berger and Bott in and Grossman in . Grossman’s bound gives Berger and Bott’s bound if one assumes that the manifold has positive sectional curvature. The first important difference between their bounds and ours is that ours does not involve any lower bound on the sectional curvatures; this is certainly an advantage if one is interested in upper bounds for $`\mathrm{log}R`$ as in Theorem A. But even in the case of manifolds of non-negative sectional curvature our bound is sharper as we now explain. To make the comparison easier suppose that our Riemannian metric $`g`$ has $`0K(P)1`$ for all 2-planes $`P`$. In this case our Theorem B gives
$$h_{top}(g)\frac{1}{2}\left(n1\underset{vSM}{\mathrm{min}}r(v)\right).$$
Under the same hypotheses on $`g`$ the inequality in Proposition 5.1 in yields
$$\underset{T+\mathrm{}}{lim\; sup}\frac{1}{T}\mathrm{log}_Mn_T(x,y)𝑑y\frac{2(n1)}{\pi }\mathrm{log}\left(2+\frac{\pi }{2}\right).$$
which is certainly weaker than our bound even ignoring the term involving the Ricci curvature since $`\frac{2}{\pi }\mathrm{log}\left(2+\frac{\pi }{2}\right)`$ is aproximately $`0.8103`$.
## 2. Proof of Theorem A and Corollaries 1 and 2
### 2.1. Proof of Theorem A
As in the introduction, given $`x`$ and $`y`$ in $`M`$ and $`T>0`$, let $`n_T(x,y)`$ be the number of geodesic arcs joining $`x`$ and $`y`$ with length $`T`$. For each $`T>0`$, the counting function $`n_T(x,y)`$ is finite and locally constant on an open full measure subset of $`M\times M`$, and integrable on $`M\times M`$ .
Suppose now that the points $`x`$ and $`y`$ are not conjugate along any geodesic connecting them. Let $`\mathrm{\Omega }(x,y)`$ be the space of piecewise smooth paths $`\alpha :[0,1]M`$ with $`\alpha (0)=x`$ and $`\alpha (1)=y`$. Given a non-negative integer $`q`$, let $`i_q(x,y)`$ be the number of geodesic arcs from $`x`$ to $`y`$ with index $`q`$. The Morse inequalities applied to the energy functional on $`\mathrm{\Omega }(x,y)`$ give right away :
(4)
$$\underset{i=0}{\overset{q}{}}b_i(\mathrm{\Omega }(x,y))i_q(x,y).$$
It follows from the proof of Myers’ theorem that if the Ricci curvarture $`r`$ satisfies $`r\delta g`$ for $`\delta >0`$, then any geodesic arc between $`x`$ and $`y`$ with length at least $`\pi \sqrt{\frac{n1}{\delta }}q`$ has index at least $`q`$. Hence for any non-negative integer $`q`$ we have:
(5)
$$i_q(x,y)n_{\pi \sqrt{\frac{n1}{\delta }}q}(x,y).$$
Combining (4) and (5) we obtain
$$\underset{i=0}{\overset{q}{}}b_i(\mathrm{\Omega }(x,y))n_{\pi \sqrt{\frac{n1}{\delta }}q}(x,y),$$
and integrating this inequality with respect to $`y`$ gives:
$$\underset{i=0}{\overset{q}{}}b_i(\mathrm{\Omega }(x,y))\frac{1}{\text{Vol}(M)}_Mn_{\pi \sqrt{\frac{n1}{\delta }}q}(x,y)𝑑y,$$
and hence
(6)
$$\underset{q+\mathrm{}}{lim\; sup}\frac{1}{q}\mathrm{log}\underset{i=0}{\overset{q}{}}b_i(\mathrm{\Omega }(x,y))\pi \sqrt{\frac{n1}{\delta }}\underset{T+\mathrm{}}{lim\; sup}\frac{1}{T}\mathrm{log}_Mn_T(x,y)𝑑y.$$
This argument and the last inequality are taken from .
Now let $`R_\mathrm{\Omega }`$ be the radius of convergence of the Poincaré series:
$$\underset{i0}{}b_i(\mathrm{\Omega }(x,y),)t^i.$$
Since this series always has infinitely many non-zero coefficients we clearly have:
$$\mathrm{log}R_\mathrm{\Omega }=\underset{q+\mathrm{}}{lim\; sup}\frac{1}{q}\mathrm{log}\underset{i=0}{\overset{q}{}}b_i(\mathrm{\Omega }(x,y),).$$
On the other hand Babenko showed in that if $`M`$ is rationally hyperbolic, then $`R=R_\mathrm{\Omega }`$ where $`R`$ is the radius of convergence of:
$$\underset{i2}{}\text{dim}(\pi _i(M))t^i.$$
Using (6) we obtain inequality (1) in the introduction, namely:
(7)
$$\frac{\sqrt{\delta }\mathrm{log}R}{\pi \sqrt{n1}}\underset{T+\mathrm{}}{lim\; sup}\frac{1}{T}\mathrm{log}_Mn_T(x,y)𝑑y.$$
In , G.P. Paternain explained how Yomdin’s theorem can be used to prove that (see also ):
(8)
$$\underset{T+\mathrm{}}{lim\; sup}\frac{1}{T}\mathrm{log}_Mn_T(x,y)𝑑yh_{top}(g),$$
where $`h_{top}(g)`$ denotes the topological entropy of the geodesic flow of $`g`$. Combining (7) and (8) yields:
(9)
$$\frac{\sqrt{\delta }\mathrm{log}R}{\pi \sqrt{n1}}h_{top}(g).$$
Theorem A is an immediate consequence of inequality (9), Theorem B and the following observation: if $`M`$ admits an Einstein metric $`g`$ with non-negative sectional curvature, then we can rescale $`g`$ so that $`r=g`$ and $`0K(P)1`$.
### 2.2. Proof of Corollary 1
We will use the following result of Y. Félix and J.C. Thomas :
###### Theorem 2.1.
Suppose that $`M`$ is formal and rationally hyperbolic and let $`P_M`$ be the Poincaré polynomial of $`H_{}(M,)`$. Write $`P_M(t)=_{i=1}^n(tz_i)`$. Then
$$R\underset{1in}{\mathrm{min}}|z_i|.$$
###### Corollary 2.2.
Let $`M^n`$ be a closed formal rationally hyperbolic manifold. Then
$$\text{dim}H_{}(M,)\left(1+\frac{1}{R}\right)^n.$$
###### Proof.
By Poincaré duality if $`z`$ is a root of $`P_M(t)`$ then $`1/z`$ is also a root of $`P_M(t)`$. Hence if we write
(10)
$$P_M(t)=\underset{i=1}{\overset{n}{}}(tz_i)$$
it follows that
(11)
$$\underset{1in}{\mathrm{min}}|z_i|=\frac{1}{\mathrm{max}_{1in}|z_i|}.$$
Let us set for brevity $`A:=\mathrm{max}_{1in}|z_i|`$. From (10) we get:
$$b_i({}_{i}{}^{n})A^i,$$
hence
$$\text{dim}H_{}(M,)\underset{i=0}{\overset{n}{}}({}_{i}{}^{n})A^i=(1+A)^n.$$
By Theorem 2.1 and (11) we have $`A1/R`$ which yields
$$\text{dim}H_{}(M,)\left(1+\frac{1}{R}\right)^n$$
as desired.
Corollary 1 follows from Theorem A and Corollary 2.2 if $`M`$ is rationally hyperbolic. If $`M`$ is rationally elliptic we always have $`\text{dim}H_{}(M,)2^n`$.
###### Remark 2.3.
It seems interesting to observe that from inequality (9) and Corollary 2.2 it follows that if $`(M^n,g)`$ is a closed simply connected formal manifold with $`r\delta g`$, $`\delta >0`$, then
$$\text{dim}H_{}(M,)\left[1+\mathrm{exp}\left(\pi h_{top}(g)\sqrt{\frac{n1}{\delta }}\right)\right]^n.$$
This suggests that the following smooth invariant of a manifold that admits a metric of positive Ricci curvature might be of interest:
$$h_r(M):=\underset{\{g:r(g)g\}}{inf}h_{top}(g).$$
From the last inequality, if $`M`$ is formal and admits a metric of positive Ricci curvature we have
$$\text{dim}H_{}(M,)\left[1+\mathrm{exp}\left(\pi \sqrt{n1}h_r(M)\right)\right]^n.$$
### 2.3. Proof of Corollary 2
In , Félix and Thomas point out the following corollary of Theorem 2.1:
###### Corollary 2.4.
If $`M^n`$ is a closed $`(p1)`$-connected formal manifold then
$$R\left(\frac{n}{pb_p}\right)^{1/p}.$$
Corollary 2 follows right away from the Corollary 2.4 and Theorem A.
## 3. Proof of Theorem B
Let $`SMTM`$ be the unit sphere bundle of $`M`$. We will consider the geodesic flow $`\mathrm{\Phi }_t`$ of $`g`$ acting on $`SM`$. Given any $`\theta SM`$ we will denote by $`c_\theta `$ the geodesic with initial condition $`\theta `$. The expression $`\theta =(x,v)`$ will mean that $`vT_xM`$.
Recall that the metric $`g`$ on $`M`$ induces a metric on $`TM`$ (the Sasaki metric) and for any $`\theta TM`$ an orthogonal decomposition of $`T_\theta TM`$ into horizontal and vertical parts: $`T_\theta TM=H_\theta V_\theta `$ (the vertical space is of course the tangent space of the fiber). Recall also that the differential of $`\mathrm{\Phi }_t`$ has a nice expression in geometric terms; given $`\xi =(w_1,w_2)T_\theta TM`$, $`d_\theta \mathrm{\Phi }_t(\xi )=(J_\xi (t),\dot{J}_\xi (t))`$, where $`J_\xi `$ is the Jacobi field along $`c_\theta `$ with initial conditions $`J_\xi (0)=w_1`$ and $`\dot{J}_\xi (0)=w_2`$.
If $`\theta =(x,v)`$ we will denote by $`S(\theta )`$ the orthogonal complement of $`(v,0)`$ in $`T_\theta SM`$. Equivalently, $`S(\theta )`$ is the orthogonal complement of the subspace spanned by $`(v,0)`$ and $`(0,v)`$ in $`T_\theta M`$. It is easy to see that the subspaces $`S(\theta )`$ are invariant through the differential of the geodesic flow.
The proof of Theorem B will be based on Przytycki’s inequality for the topological entropy of the geodesic flow which we will describe now.
Given two real vector spaces with inner product $`V`$ and $`W`$ of the same dimension and a linear transformation $`f:VW`$ the expansion of $`f`$, $`\text{ex}(f)`$, is the supremum over all non-trivial subspaces of $`V`$ of the absolute of the determinant of $`f|_V`$. An important (and trivial) property of the expansion is that given two linear maps $`f`$ and $`g`$ we have $`\text{ex}(fg)\text{ex}(f)\text{ex}(g)`$. This will be used below.
Przytycki’s inequality is the following:
$$h_{top}(g)\underset{t\mathrm{}}{lim\; inf}\frac{1}{t}\mathrm{log}_{SM}\text{ex}(d_\theta \mathrm{\Phi }_t)𝑑\theta .$$
Here we consider $`d_\theta \mathrm{\Phi }_t`$ as a map between $`S(\theta )`$ and $`S(\dot{c}_\theta (t))`$.
###### Remark 3.1.
We only need to use Przytycki’s inequality although it is actually true that for a $`C^{\mathrm{}}`$ Riemanniann metric one has Mañé’s formula :
$$h_{top}(g)=\underset{t\mathrm{}}{lim}\frac{1}{t}\mathrm{log}_{SM}\text{ex}(d_\theta \mathrm{\Phi }_t)𝑑\theta =\underset{T\mathrm{}}{lim}\frac{1}{T}\mathrm{log}_{M\times M}n_T(x,y)𝑑x𝑑y.$$
Moreover, O.S. Kozlovski has shown that Przytycki’s inequality is an equality for arbitrary $`C^{\mathrm{}}`$ maps.
We will find our upper bound for the topological entropy by studying $`\text{ex}(d_\theta \mathrm{\Phi }_t)`$ for small values of $`t`$. Given any $`\delta >0`$ we have:
$$h_{top}(g)\underset{i\mathrm{}}{lim\; inf}\frac{1}{\delta i}\mathrm{log}_{SM}\text{ex}(d_\theta ((\mathrm{\Phi }_\delta )^i))𝑑\theta .$$
Now suppose that for some $`\delta >0`$ and $`\alpha >0`$ we have that $`\text{ex}(d_\theta \mathrm{\Phi }_\delta )\alpha `$ for all $`\theta `$. Then it follows that
$$h_{top}(g)\underset{i\mathrm{}}{lim}\frac{1}{\delta i}\mathrm{log}(\alpha ^i\text{Vol}(SM))=\frac{\mathrm{log}(\alpha )}{\delta }.$$
We will get an estimate of the form $`\alpha =1+\beta \delta +O(\delta ^2)`$. Then $`\mathrm{log}(\alpha )=\beta \delta +O(\delta ^2)`$ and since the last inequality holds for any $`\delta >0`$ we get
(12)
$$h_{top}(g)\beta .$$
We will prove Theorem B by estimating $`\beta `$. Consider the polar decomposition of $`d_\theta \mathrm{\Phi }_t`$: $`d_\theta \mathrm{\Phi }_t=O_t(\theta )L_t(\theta )`$ where $`O_t(\theta ):S(\theta )S(\dot{c}_\theta (t))`$ is a linear isometry and $`L_t(\theta )`$ is a symmetric positive endomorphism of $`S(\theta )`$. Then $`ex(d\mathrm{\Phi }_t)=ex(L_t)`$. Of course, $`L_t`$ can be given explicitly: $`L_t=((d\mathrm{\Phi }_t)^{}(d\mathrm{\Phi }_t))^{1/2}`$.
Consider the map $`:S(\theta )S(\theta )`$ which in the decomposition into horizontal and vertical parts is given by $`(w_1,w_2)=(w_2,R(v,w_1)v)`$, where $`R`$ is the curvature tensor and $`\theta =(x,v)SM`$.
###### Lemma 3.2.
$`((d\mathrm{\Phi }_\delta )^{}(d\mathrm{\Phi }_\delta ))^{1/2}=Id+\frac{\delta }{2}(+^{})+O(\delta ^2).`$
###### Proof.
Given any $`\theta =(x,v)SM`$ let $`c_\theta `$ be the geodesic with initial condition $`\theta `$. Let $`T_t`$ be the parallel transport along $`c_\theta `$ from $`T_{c_\theta (0)}`$ to $`T_{c_\theta (t)}`$. Let $`e_1,e_2,\mathrm{},e_{n1}`$ be an orthonormal basis of $`\{v\}^{}T_xM`$ by eigenvectors of the symmetric transformation $`uR(v,u)v`$. Let $`E_i`$ be the parallel vector field along $`c_\theta `$ with initial condition $`e_i`$. Given $`\xi T_\theta SM`$ we can write $`J_\xi (t)=_{i=1}^{n1}a_i(t)E_i(t)`$ for some smooth functions $`a_i`$. Then,
$$J_\xi (\delta )=\underset{i=1}{\overset{n1}{}}(a_i(0)+\delta a_i^{}(0))T_\delta E_i(0)+O(\delta ^2)$$
and since $`\dot{J}_\xi (t)=_{i=1}^{n1}a_i^{}(t)E_i(t)`$, we have
$$\dot{J}_\xi (\delta )=\underset{i=1}{\overset{n1}{}}(a_i^{}(0)+\delta a_i^{\prime \prime }(0))T_\delta E_i(0)+O(\delta ^2).$$
Now, from the Jacobi equation, we get that $`_{i=1}^{n1}a_i^{\prime \prime }(0)e_i=R(v,J_\xi (0))v`$. Therefore
$`(J_\xi (\delta ),`$ $`\dot{J}_\xi (\delta ))=`$
$`(T_\delta \left({\displaystyle \underset{i=1}{\overset{n1}{}}}(a_i(0)+\delta a_i^{}(0))e_i\right),T_\delta \left({\displaystyle \underset{i=1}{\overset{n1}{}}}a_i^{}(0)e_i\delta R(v,J_\xi (0))v\right))+O(\delta ^2),`$
i.e.
$$d\mathrm{\Phi }_\delta =(T_\delta ,T_\delta )(Id+\delta )+O(\delta ^2).$$
Since $`T_\delta `$ is orthogonal we get
$$(d\mathrm{\Phi }_\delta )^{}(d\mathrm{\Phi }_\delta )=Id+\delta (+^{})+O(\delta ^2)$$
and therefore,
$$((d\mathrm{\Phi }_\delta )^{}(d\mathrm{\Phi }_\delta ))^{1/2}=Id+(\delta /2)(+^{})+O(\delta ^2).$$
From the lemma we obtain that $`\text{ex}(d\mathrm{\Phi }_\delta )=\text{ex}\left(Id+\frac{\delta }{2}(+^{})\right)+O(\delta ^2)`$; and we are left to compute $`\text{ex}(Id+\frac{\delta }{2}(+^{}))`$. This is a positive definite symmetric endomorphism of $`S(\theta )`$ and therefore the expansion is the product of the eigenvalues which are greater than or equal to one, provided that there exists at least one eigenvalue greater than or equal to one.
We can compute the eigenvalues explicitly using the orthonormal basis of $`S(\theta )`$ given by:
$$\{(e_1,e_1),\mathrm{},(e_{n1},e_{n1}),(e_1,e_1),\mathrm{},(e_{n1},e_{n1})\}.$$
Suppose that $`R(v,e_i)v=\lambda _ie_i`$. Then it is easy to check that
$$\left(Id+(\delta /2)(+^{})\right)(e_i,e_i)=\left(1+(\delta /2)(1\lambda _i)\right)(e_i,e_i)$$
and
$$\left(Id+(\delta /2)(+^{})\right)(e_i,e_i)=\left(1+(\delta /2)(\lambda _i1)\right)(e_i,e_i).$$
If $`\lambda _i1`$ for all $`i`$ we get that
$$\text{ex}(Id+(\delta /2)(+^{}))=\underset{i=1}{\overset{n1}{}}(1+(\delta /2)(1\lambda _i))=1+(\delta /2)(n1r(v))+O(\delta ^2),$$
where $`r(v)`$ denotes, as before, the Ricci curvature in the direction of $`v`$.
Therefore if the sectional curvature of $`g`$ is bounded above by 1 we get, from the previous discussion, that
$$h_{top}(g)\frac{1}{2}\left(n1\underset{vSM}{\mathrm{min}}r(v)\right).$$
For a general $`g`$, let $`k`$ be the maximum of all the sectional curvatures. If $`k>0`$ the metric $`g_k=kg`$ has sectional curvature bounded above by one and from the previous observation we get that
$$h_{top}(g_k)\frac{1}{2}\left(n1\frac{\mathrm{min}_{vSM_g}r_g(v)}{k}\right).$$
But $`h_{top}(g_k)=(1/\sqrt{k})h_{top}(g)`$ and therefore we get
$$h_{top}(g)\frac{1}{2}\left(\sqrt{k}(n1)\frac{\mathrm{min}_{vSM}r(v)}{\sqrt{k}}\right).$$
This finishes the proof of Theorem B.
###### Remark 3.3.
If the sectional curvature of $`g`$ is non-positive (i.e. if $`k0`$), we get that for any positive number $`\rho `$,
$$h_{top}(g_\rho )\frac{1}{2}\left(n1\frac{\mathrm{min}_{vSM_g}r_g(v)}{\rho }\right)$$
and therefore,
$$h_{top}(g)\frac{1}{2}\left(\sqrt{\rho }(n1)\frac{\mathrm{min}_{vSM}r(v)}{\sqrt{\rho }}\right).$$
The sharpest inequality is obtained by taking $`\rho =\frac{\mathrm{min}_{vSM}r(v)}{n1}`$ and hence we have:
(13)
$$h_{top}(g)\sqrt{(n1)\underset{vSM}{\mathrm{min}}r(v)}.$$
The Bishop comparison theorem implies right away that if $`r(n1)g`$ then $`\lambda `$, the exponential growth rate of volume of balls in the universal covering, satisfies
$$\lambda n1.$$
But when the sectional curvature is non-positive, Manning proved that $`h_{top}(g)=\lambda `$ and we recover inequality (13).
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# 1 Left - Dipole cross section for different colour states. Right - Evolution of the colour probabilities 𝑃_𝑆, 𝑃_𝑂, 𝑃_𝐷, along 𝑧 and with 𝑏⃗=0. Also shown is the profile of a Pb nucleus.
A NEW MECHANISM FOR
J/$`\psi `$ SUPPRESSION IN NUCLEAR COLLISIONS
ALBERTO POLLERI
Institut für Theoretische Physik der Universität,
Philosophenweg 19, D-69120 Heidelberg, Germany.
e-mail: polleri@tphys.uni-heidelberg.de
## Abstract
We present new results based on an improved version of the Glauber model, used to describe the nuclear absorption of J/$`\psi `$ in reactions between nuclei at high energy. The excitation of the colour degrees of freedom of nucleons, due to collisions mediated by one-gluon exchange, is taken into account. It is found that the proposed mechanism leads to larger nuclear absorption of J/$`\psi `$ than previously considered.
The origin of the anomalous behaviour of the J/$`\psi `$ cross section measured in Pb+Pb collisions is still not understood and competing interpretations have been proposed . In order to resolve this important and controversial issue, more detailed calculations are demanded. Providing further insight in this difficult problem, preliminary results, obtained with a new formulation of the Glauber approach to the early stage of a nuclear collision, are hereby presented.
Preparing the ground for the discussion, it is useful to recall the Glauber model expression for the cross section in a A+B collision at impact parameter $`\stackrel{}{b}`$. Denoting a generic colour singlet $`c\overline{c}`$ state by $`\psi `$, one has
$$\frac{d^2\sigma _\psi ^{AB}}{d^2\stackrel{}{b}}(\stackrel{}{b})=\sigma _\psi ^{NN}d^2\stackrel{}{s}𝑑z_A𝑑z_B\rho _A(z_A,\stackrel{}{s})\rho _B(z_B,\stackrel{}{b}\stackrel{}{s})S^{abs}(z_A,z_B,\stackrel{}{b},\stackrel{}{s}),$$
(1)
where $`\rho _A`$ and $`\rho _B`$ are the densities of the colliding nuclei and
$$S^{abs}(z_A,z_B,\stackrel{}{b},\stackrel{}{s})=\mathrm{exp}\left[\sigma _{\psi N}^{abs}(_{\mathrm{}}^{z_A}𝑑z_A^{}\rho _A(z_A^{},\stackrel{}{s})+_{z_B}^+\mathrm{}𝑑z_B^{}\rho _B(z_B^{},\stackrel{}{b}\stackrel{}{s}))\right]$$
(2)
is the nuclear suppression factor, which contains the effective absorption cross section $`\sigma _{\psi N}^{abs}`$ for $`\psi `$-nucleon scattering.
Provided that $`\sigma _{\psi N}^{abs}=7.3`$ mb, the above expressions can account for p+A and S+Pb data, but not for the Pb+Pb measurements. On the other hand, the usual version of the Glauber model must be improved, since in eqs. (1) and (2) one assumes that the nucleons which absorb the produced $`\psi `$ are in their ground state. This cannot be correct. In fact, before encountering the produced meson, they scatter several times while the nuclei stream through each other. In doing so they leave the ground state. A first attempt to address this problem was recently made, by including in the calculation the effect of the cloud of prompt gluons around nucleons . On the other hand, nucleon themselves were still treated as non-interacting. Since the centre of mass energy of each NN collision is $`17`$ GeV at the SPS, it is reasonable to ascribe the main contribution to the inelastic cross sections to one-gluon exchange processes, as done, for example, in the Dual Parton Model . This has an important consequence: two colour-singlet nucleons jump into octet states after the elementary collision and therefore become coloured. In same way one also allows the possibility that some ‘nucleons’ are in a decuplet state. From now on the quotes will be dropped when referring to a coloured nucleon.
Having realized that nucleons become coloured, soon after enough rescattering has taken place, it is necessary to establish whether this fact has any effect at all on $`\psi `$ absorption. In other words one must quantify the differences in the inelastic cross sections for $`\psi `$N scattering due to the various colour states of the nucleon. This can be achieved within the Low-Nussinov model , first evaluating the elastic cross section for a meson scattering off a nucleon, and then making use of the optical theorem, which relates the forward elastic amplitude to the total cross section. The model consists in the exchange of two gluons characterised by a phenomenological mass $`\mu _G140`$ MeV, which effectively mimics confinement. Neglecting the elastic contribution compared to the inelastic one, one can show that
$$\sigma _{\psi N}^{abs}=d^2\rho \left|\mathrm{\Phi }_\psi (\rho )\right|^2\sigma _a(\rho ),$$
(3)
where $`\mathrm{\Phi }_\psi (\rho )`$ is the transverse part of the $`\psi `$ meson wave function, while
$$\sigma _a(\rho )=\frac{16}{3}\alpha _s^2d^2k\frac{1}{(k^2+\mu _G^2)^2}\left[1aF_N(3k^2)\right]\left(1e^{i\stackrel{}{k}\stackrel{}{\rho }}\right)$$
(4)
is the colour dipole-nucleon cross section. The function $`F_N`$ is the two-quark form factor of the nucleon, usually identified with the charge form factor. The coefficient $`a`$ in front of the form factor specifies the colour state of the nucleon. For the singlet state one has $`a=1`$ and the expression reduces to the usual one, exhibiting the partial cancellation of amplitudes, implying $`\sigma (\rho )0`$ for $`\rho 0`$ (Colour transparency). A detailed calculation shows that $`a=1/4`$ for an octet nucleon while $`a=1/2`$ for the decuplet state.
The dipole cross section can be computed analytically in terms of modified Bessel functions. One can fix the value of the strong coupling in order to reproduce $`\pi `$N cross section of $`23`$ mb, therefore taking $`\alpha _s0.65`$. Successively, the $`\psi `$N cross sections can be obtained using eqs. (3) and (4) without additional parameter fixing. The results of the calculations are illustrated in the left side of Figure 1, which shows the dipole cross section for the three colour states of the nucleon. It largely increases when the dipole scatters off a coloured object. This new result is important and could not have been guessed prior to calculation. One must now evaluate the corresponding absorption cross sections. To do so it is useful to first rewrite eq. (4) in the form
$$\sigma _a(\rho )=\sigma (\rho )\left[\mathrm{\hspace{0.17em}1}+(1a)\mathrm{\Delta }(\rho )\right],$$
(5)
where the function $`\mathrm{\Delta }(\rho )`$ was found to be approximately constant and amounts to $`0.46`$. This shows that the colour cross sections have a simple $`\rho `$-dependence and scale with respect to the singlet one. Then, using eq. (3), one obtains the absorption cross sections for $`\psi `$ as
$$\sigma _{\psi N_j}^{abs}=\sigma _{\psi N}^{abs}\times \{\begin{array}{c}\text{1}\text{Singlet}\\ \text{1.35}\hfill & \text{Octet}\hfill \\ \text{1.7}\hfill & \text{Decuplet}\hfill \end{array}.$$
(6)
A large increase with respect to the singlet value is found, suggesting that the newly proposed absorption mechanism is indeed important. Using a Gaussian parametrisation for the $`\psi `$ wave function, such that the root mean squared radius is $`0.2`$ fm, one obtains $`\sigma _{\psi N}^{abs}6`$ mb. On the other hand, at this preliminary stage of the calculation, the effect of the feeding into J/$`\psi `$ from $`\psi ^{}`$ and $`\chi _c`$ states is neglected. To compare with the conventional Glauber model, it is therefore preferable to use the effective value used to reproduce the p+A data, therefore setting $`\sigma _{\psi N}^{abs}=7.3`$ mb and scaling the colour cross section accordingly by means of eq. (6).
To correct eqs. (1) and (2), one must now understand how the colour state of a nucleon evolves while passing through a nucleus. This is a non-trivial problem which can be formulated by means of a master equation. Its content is to express the evolution of the colour probabilities $`P_S`$, $`P_O`$ and $`P_D`$, due to subsequent scatterings. The framework is again that of the Glauber model of multiple collisions. What one finds is the solution
$`P_S(z,\stackrel{}{b})`$ $`=`$ $`{\displaystyle \frac{1}{27}}+{\displaystyle \frac{20}{27}}F(z,\stackrel{}{b})+{\displaystyle \frac{6}{27}}G(z,\stackrel{}{b}),`$ (7)
$`P_O(z,\stackrel{}{b})`$ $`=`$ $`{\displaystyle \frac{16}{27}}{\displaystyle \frac{40}{27}}F(z,\stackrel{}{b})+{\displaystyle \frac{24}{27}}G(z,\stackrel{}{b}),`$ (8)
$`P_D(z,\stackrel{}{b})`$ $`=`$ $`{\displaystyle \frac{10}{27}}+{\displaystyle \frac{20}{27}}F(z,\stackrel{}{b}){\displaystyle \frac{30}{27}}G(z,\stackrel{}{b}),`$ (9)
where
$$F(z,\stackrel{}{b})=\mathrm{exp}\left(X(z,\stackrel{}{b})\right),G(z,\stackrel{}{b})=\mathrm{exp}\left(\frac{2}{3}X(z,\stackrel{}{b})\right)$$
(10)
and
$$X(z,\stackrel{}{b})=\frac{9}{8}\sigma _{NN}^{in}_{\mathrm{}}^z𝑑z^{}\rho (z^{},\stackrel{}{b}),\text{with}\sigma _{NN}^{in}=30\text{mb}.$$
(11)
One notices several properties of the found probabilities. First of all, the limit $`z\mathrm{}`$ implies $`X0`$ and $`F,G1`$. This means that $`P_S1`$ and $`P_O,P_D0`$. In other words the nucleon, before entering the nucleus, is in singlet state as it should be. Moreover, $`P_S+P_O+P_D=1`$ for any $`z`$, therefore probability is always conserved. Finally, if $`z+\mathrm{}`$ and if the nucleus is large enough, one has $`X1`$ and $`F,G1`$. This implies that if enough scatterings take place, the colour probabilities reach the statistical limit, given by the first coefficients of eqs. (7), (8) and (9). The $`z`$-dependence of the colour probabilities is illustrated in the right side of Figure 1, together with the longitudinal profile of a Pb nucleus at $`\stackrel{}{b}=0`$. Soon after the nucleon has penetrated the nuclear profile, a process that involves $`4`$ fm, the statistical limit is reached. About $`2/3`$ of the nucleons are in octet state while $`1/3`$ are in decuplet. The amount of singlet is negligible.
It is now possible to describe how to perform the calculation of the $`\psi `$ cross section in A+B collisions. One must modify eqs. (1) and (2) to account for the previously discussed colour absorption cross sections and probabilities. It is necessary to replace eq. (2) with $`S^{abs}=\mathrm{exp}\left[(f_A+f_B)\right]`$, where
$`f_A(z_A,z_B,\stackrel{}{b},\stackrel{}{s})`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{z_A}}𝑑z_A^{}\mathrm{\Sigma }_B(\stackrel{~}{z}_B,\stackrel{}{b}\stackrel{}{s})\rho _A(z_A^{},\stackrel{}{s}),`$ (12)
$`f_B(z_A,z_B,\stackrel{}{b},\stackrel{}{s})`$ $`=`$ $`{\displaystyle _{z_B}^+\mathrm{}}𝑑z_B^{}\mathrm{\Sigma }_A(\stackrel{~}{z}_A,\stackrel{}{s})\rho _B(z_B^{},\stackrel{}{b}\stackrel{}{s})`$ (13)
The effective cross sections $`\mathrm{\Sigma }`$ take into account the colour cross sections and probabilities, as previously discussed. They are
$`\mathrm{\Sigma }_A(\stackrel{~}{z}_A,\stackrel{}{s})`$ $`=`$ $`{\displaystyle \underset{j=S,O,D}{}}\sigma _{\psi N_j}^{in}P_j^A(\stackrel{~}{z}_A,\stackrel{}{s}),`$ (14)
$`\mathrm{\Sigma }_B(\stackrel{~}{z}_B,\stackrel{}{b}\stackrel{}{s})`$ $`=`$ $`{\displaystyle \underset{j=S,O,D}{}}\sigma _{\psi N_j}^{in}P_j^B(\stackrel{~}{z}_B,\stackrel{}{b}\stackrel{}{s})`$ (15)
In the above expressions the values $`\stackrel{~}{z}`$ correspond to the $`\psi `$ absorption points as in the illustration shown in the left side of Figure 2. With a simple geometrical construction one obtains
$$\stackrel{~}{z}_A=z_A(z_B^{}z_B)\frac{1v_\psi }{1+v_\psi }\text{and}\stackrel{~}{z}_B=z_B(z_A^{}z_A)\frac{1+v_\psi }{1v_\psi }.$$
(16)
The velocity $`v_\psi `$ of the meson is related to the measured value of $`x_F=p_\psi /p_{max}0.15`$ and to the centre of mass energy $`\sqrt{s_{NN}}`$ of the NN collisions. One has $`v_\psi =\sqrt{x_F^2s_{NN}/(4m_\psi ^2+x_F^2s_{NN})}0.4`$.
The $`\psi `$ cross section expressed by eq. (1) can now be evaluated. Its ratio with the Drell-Yan cross section is obtained in a conventional manner and is converted into a $`E_T`$-dependent function by fixing the scale to the number of participants, in order to describe the minimum bias data as measured by the NA50 experiment. No spread in the $`E_T(\stackrel{}{b})`$ correlation is taken into account for simplicity. The calculated ratio is compared with the standard Glauber approach and to the data as shown in the right side of Figure 2. The result exhibits a stronger suppression when compared to the usual Glauber calculation. Although several improvements are under investigation, the effect is clear and cannot be neglected in future work. Among the aforementioned improvements is the inclusion of a retardation effect for the so far sudden switch of colour in NN collisions. This becomes relevant at large impact parameters, where there are only few collisions, implying that the present result overestimates the suppression at small transverse energy. Another improvement which works in the same direction consists in accounting for important formation time effects . All this is presently under careful study.
I warmly thank Jörg Hüfner and Boris Kopeliovich for the very pleasant and stimulating collaboration.
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# Macroscopic Quantum Coherence in a Magnetic Nanoparticle Above the Surface of a Superconductor
## Abstract
We study macroscopic quantum tunneling of the magnetic moment in a single-domain particle placed above the surface of a superconductor. Such a setup allows one to manipulate the height of the energy barrier, preserving the degeneracy of the ground state. The tunneling amplitude and the effect of the dissipation in the superconductor are computed.
Tunneling of the magnetic moment in nanoparticles and molecular clusters has been intensively studied theoretically and experimentally in the last decade . The interest in this problem is two-fold. Firstly, magnetic tunneling reveals itself at a quasiclassical level, that is, in situations where all three components of the magnetic moment, $`𝐌`$, can be rather accurately determined by a macroscopic measurement. The interaction of $`𝐌`$ with microscopic degrees of freedom makes this problem one of tunneling with dissipation . Secondly, tunneling of the magnetic moment changes the magnetic properties of small magnets, with potential implications for the data- storage technology. It also adds nanomagnets to the list of candidates for qubits - the elements of quantum computers.
In zero magnetic field the magnetic state of a classical magnet is degenerate with respect to $`𝐌𝐌`$ due to time-reversal symmetry. In nanoparticles the $`|>`$ and $`|>`$ minima of the energy are separated by a barrier, $`U`$, due to the magnetic anisotropy. The thermal rate of switching between the two classical states is proportional to $`\mathrm{exp}(U/T)`$. At high temperatures, when the thermal rate is high, the particle is in the superparamagnetic regime. At low temperature, as far as the thermal rate is concerned, the magnetic moment should freeze along one of the anisotropy directions. In particles of sufficiently small size, however, even at $`T=0`$ the magnetic moment can switch due to quantum tunneling. The quantum switching rate, $`\mathrm{\Gamma }`$, scales with the total spin $`S`$ of the nanoparticle according to $`\mathrm{ln}\mathrm{\Gamma }S`$. If the switching time, $`\mathrm{\Gamma }^1`$, is small compared to the measuring time, the nanoparticle remains superparamagnetic in the limit of $`T=0`$. If, in addition to that, the interaction of $`𝐌`$ with microscopic degrees of freedom (phonons, itinerant electrons, nuclear spins, etc.) is small, the nanoparticle, during a certain decoherence time, can exist in a coherent quantum superposition of two classical states, $`\mathrm{\Psi }_o=|>+|>`$. In that state the probability that the moment of the particle has a certain orientation oscillates in time as $`\mathrm{cos}(2\mathrm{\Gamma }t)`$.
For a magnetic particle of considerable size to be in the quantum superparamagnetic (not necessarily coherent) regime, the energy barrier between the $`|>`$ and $`|>`$ states must be made sufficiently small. One way to decrease the barrier is to use an external magnetic field. This method has been used in experiments on magnetic particles and molecular clusters performed to date . It has a clear drawback if one attempts to create a coherent superposition of the $`|>`$ and $`|>`$ states. Namely, the external field, unless it is applied exactly parpendicular to the anisotropy axis, removes the degeneracy between the $`|>`$ and $`|>`$ states. In this Letter, we suggest a method of controlling the barrier without breaking the degeneracy. This method is illustrated in Fig. 1. The nanoparticle is placed above the surface of a superconductor at a variable distance controlled by, e.g., a piezoelectric layer or holder. The current induced in the superconductor creates the magnetic image of the nanoparticle. The interaction between the nanoparticle and the superconductor is then equivalent to the dipole interaction between the nanoparticle and its image. The reduction of the barrier is similar to the one from the external magnetic field applied opposite to $`𝐌`$. However, contrary to the situation with the external field, the system (the nanoparticle plus the superconductor) is now degenerate with respect to $`𝐌𝐌`$. It should be emphasized that such a degeneracy is a very general property of the system that is independent of the shape of the particle and the landscape of the superconducting surface. It is rooted in the time-reversal symmetry of the system in the absence of the field. The tunneling rate for the situation shown in Fig. 1 will be computed below.
A high tunneling rate does not automatically provide the coherent superposition of quasiclassical states. Different mechanisms of decoherence due to interactions inside the nanoparticle have been worked out in recent years . At low temperature, in the absence of nuclear spins and itinerant electrons, the effect of dissipation on the tunneling rate can be very small . Decoherence is a more subtle issue. Generally speaking, macroscopic quantum coherence (MQC), that is $`\mathrm{cos}(2\mathrm{\Gamma }t)`$ oscillations of $`𝐌`$, occur only if the decohering interactions are small compared to the tunnel splitting of the ground state, $`\mathrm{}\mathrm{\Gamma }`$. We will show that this condition can be satisfied at least as far as the interaction between the nanoparticle and the superconductor is concerned. This question is of interest also in a more general context of measuring the rotation of a mesoscopic spin with the help of a superconducting device. Indeed, most of the experiments on individual nanoparticles used SQUIDs. Although the problem studied here is different from those experimental situations, some of the ideas should also apply in those cases.
Let us consider tunneling in a nanoparticle above the flat surface of a superconductor , as shown in Fig. 1, in the absence of dissipation. The external magnetic field will be considered zero throughout this paper. To make the classical electrodynamics of the problem less cumbersome, we will make certain simplifying assumptions about the superconductor, the shape of the particle and its magnetic anisotropy. None of them is important and any generalization can be studied along the same lines. First, we shall assume that the superconductor is in the Meissner regime; that is, the magnetic field at the surface of the superconductor does not exceed $`H_{c1}`$. We shall also assume that the characteristic geometrical dimensions of the problem, the size of the particle and its distance to the surface, are large compared to the penetration depth $`\lambda _L`$. In that case, it is a good approximation to take the field of the particle at the surface of the superconductor to be parallel to that surface. The effect of the superconductor on the particle is then equivalent to the magnetic dipole interaction with the image shown in Fig. 1.
Next we assume that the particle is of ellipsoidal shape (that is, uniformly magnetized) with crystal fields either small or dominated by the single-ion anisotropy. The total energy of the magnetic anisotropy of such a particle must be quadratic in the magnetization ,
$$E_{an}=\frac{2\pi }{V}N_{ik}M_iM_k,$$
(1)
where $`𝐌`$ is the total magnetic moment of the particle, $`V`$ is its volume, and the tensor $`N_{ik}`$ includes both the demagnetizing effect (that is, shape anisotropy) and the magnetocrystalline anisotropy. In a ferromagnetic particle, $`M`$ is proportional to $`V`$, while $`N_{ik}`$ is independent of $`V`$. The factor $`V^1`$ in Eq. (1) is, therefore, needed to provide the correct linear scaling of $`E_{an}`$ with $`V`$. We shall assume that the principal axes of $`N_{ik}`$ coincide with the coordinate axes in Fig. 1; with $`Z`$ being the easy magnetization direction and $`N_{xx}>N_{yy}>N_{zz}`$. The magnitude of $`𝐌`$ is assumed, as usual, to be formed by a strong exchange interaction and, thus, independent of the orientation. That is, $`M_x^2+M_y^2+M_z^2=M^2=const.`$ The energy of the magnetic anisotropy then becomes
$$E_{an}=\frac{1}{V}(\beta _xM_x^2\beta _zM_z^2),$$
(2)
where $`\beta _x,\beta _z`$ are positive dimensionless coefficients of order unity.
Eq. (2) describes a magnet having a $`YZ`$ easy magnetization plane with $`Z`$ being the easy axis in that plane. The two degenerate minima of Eq. (2) correspond to $`𝐌`$ looking along and opposite to the $`Z`$ axis. In the Meissner state of the superconductor, the magnetic field of the particle induces superconducting currents whose field is equivalent to the field of the image shown in Fig. 1. As the particle moves closer to the superconductor, the interaction between the particle and its image increases and the barrier between the two equilibrium orientations of $`𝐌`$ decreases. The magnetic moment of the particle, $`𝐌`$, and the moment of the image, $`𝐦`$, are related through
$$M_x=m_x,M_y=m_y,M_z=m_z.$$
(3)
With reasonable accuracy the energy of the magnetic dipole interaction between the particle and its image is given by
$$E_{int}=\frac{[𝐌𝐦3(𝐧𝐌)(𝐧𝐦)]}{(2d)^3},$$
(4)
where $`d`$ is the distance from the (center of the) particle to the surface of the superconductor. With the help of relations (4), one obtains
$$E_{int}=\frac{M_z^2}{(2d)^3}.$$
(5)
The total energy of the system, $`E=E_{an}+E_{int}`$ then becomes
$$E=\frac{1}{V}(\beta _zϵM_z^2+\beta _xM_x^2),$$
(6)
where we have introduced $`ϵ=1V/\beta _z(2d)^3`$.
According to Eq. (6) the energy barrier between the degenerate $`|>`$ and $`|>`$ states of the particle is given by $`U=\beta _zϵM_o^2V`$, where $`M_o=M/V`$ is the volume magnetization of the particle. Our main idea is to manipulate $`d`$ in such a way that $`ϵ`$ and, consequently, $`U`$ become small enough to provide a significant tunneling rate. Note that most deviations from the simplifying assumptions made above will renormalize $`\beta _{x,z}`$ and $`d`$ in Eq. (6) but will not change the form of the total energy. It is convenient to introduce the total dimensionless spin of the particle, $`𝐒=𝐌/\mathrm{}\gamma `$ ($`\gamma `$ being the gyromagnetic ratio) and two characteristic frequencies:
$$\omega _{}=\beta _zϵ\gamma M_o,\omega _{}=\beta _x\gamma M_o.$$
(7)
The total energy can be then written as
$$E=\frac{1}{S}[\mathrm{}\omega _{}S_z^2+\mathrm{}\omega _{}S_x^2].$$
(8)
The corresponding tunneling problem has been studied by a number of authors . In the limit of $`\omega _{}<<\omega _{}`$, i.e. at $`ϵ<<1`$, the tunneling rate at $`T=0`$ in the absence of dissipation is given by $`\mathrm{\Gamma }_o=A_o\mathrm{exp}(B_o)`$ with
$`A_o`$ $`=`$ $`{\displaystyle \frac{16}{\sqrt{\pi }}}S^{1/2}\omega _{}^{3/4}\omega _{}^{1/4}`$ (9)
$`B_o`$ $`=`$ $`2S\left({\displaystyle \frac{\omega _{}}{\omega _{}}}\right)^{1/2}.`$ (10)
So far, we have neglected non-dissipative terms in the total energy that come from the superconductor. One such term is the kinetic energy of the Cooper pairs,
$$E_{sc1}=d^3r\frac{n_sm𝐯_s^2}{2},$$
(11)
where $`n_s`$ is the concentration of superconducting electrons, $`m`$ is their mass, and $`𝐯_s=𝐣_s/en_s`$ is their drift velocity expressed in terms of the superconducting current $`j_s`$. Eq. (10) can be written as
$$E_{sc1}=\frac{4\pi \lambda _L^2}{c^2}d^3r𝐣_s^2,$$
(12)
where $`\lambda _L=mc^2/4\pi e^2n_s`$ is the London penetration depth. The superconducting current is concentrated near the surface, resulting in the surface current
$$𝐠_s=𝑑z𝐣_s=\frac{c}{4\pi }𝐧\times 𝐁(𝐫),$$
(13)
where $`𝐧`$ is the unit vector in the $`Z`$ direction and $`𝐁(𝐫)`$ is the sum of the dipole fields of the magnetic particle and its image at $`z=0`$. A somewhat tedious but straightforward calculation then gives
$$E_{sc1}=\frac{\lambda _L}{8\pi }d^2r𝐁^2(𝐫)=const+\frac{3\lambda _L}{16d^4}M_z^2.$$
(14)
The contribution of the kinetic energy of the superconducting electrons to the total energy is small if $`d>>\lambda _L`$. Even at $`d\lambda _L`$, however, it reduces to the renormalization of the $`d`$ dependence of the interaction between the magnetic particle and the superconductor. Since we are interested in $`d`$ close to the critical value at which the barrier becomes zero, the form of the Hamiltonian, Eq. (8), and the expressions, Eq. (9), for the tunneling rate remain unaffected by that renormalization.
The next non-dissipative term in the energy is the inertia coming from the energy of the electric field,
$$E_{sc2}=d^3r\frac{𝐄^2}{8\pi }=\frac{1}{8\pi c^2}d^3r\left(\frac{d𝐀}{dt}\right)^2.$$
(15)
where $`𝐀=(4\pi \lambda _L^2/c)𝐣_s`$ is the vector potential in the superconductor. Then, similarly to the previous case, one obtains
$$E_{sc2}=\frac{\lambda _L^3}{16\pi c^2}d^2r\left(\frac{d𝐁}{dt}\right)^2=\frac{3\lambda _L^3}{32c^2d^4}(\dot{𝐌}^2+\dot{M}_z^2).$$
(16)
To estimate the effect of this inertia term on tunneling, it should be compared, at the instanton frequency $`\omega _{inst}=(\omega _{}\omega _{})^{1/2}`$, with $`E`$ of Eq. (8). The ratio of the energies is
$$E_{sc2}/E(\lambda _L/d)^3(l/d)(l/\lambda _{inst})^2,$$
(17)
where we have introduced $`l=V^{1/3}`$ and $`\lambda _{inst}=c/\omega _{inst}`$. Even if the experimental values of $`l`$,$`d`$ and $`\lambda _L`$ do not differ in order of magnitude, $`\lambda _{inst}`$ can hardly be less than 1 cm, making the ratio of energies in Eq. (16) negligible for nanoparticles used in tunneling experiments. We may then conclude that Eq. (9) gives a good estimate of the tunneling rate in the absence of dissipation.
Spin tunneling with dissipation due to the interaction of $`𝐌`$ with microscopic degrees of freedom inside the nanoparticle has been intensively studied . Interactions with phonons, magnons, nuclear spins, etc. have been considered. Here we will study the mechanism of dissipation specific to our problem: the interaction of the magnetic moment with normal quasiparticles in the superconductor. This analysis may also be relevant to experiments on spin tunneling performed using SQUIDs.
We begin with the derivation of the energy dissipation in the superconductor due to the rotation of $`𝐌`$. With the help of the relations $`𝐄=c^1𝐀/t`$ and $`𝐣=\sigma 𝐄`$, one obtains
$$Q=d^3r𝐣𝐄=\frac{1}{c^2}d^3r\sigma (t)\dot{𝐀}^2.$$
(18)
where $`\sigma `$ is the conductivity due to quasiparticles. Its Fourier transform can be approximated as $`\sigma _\omega =e^2n_q/m(\nu +i\omega )`$, where $`n_q`$ is the quasiparticle concentration and $`\nu `$ is their scattering rate. The latter is typically 2-3 orders of magnitude greater than the instanton frequency for the tunneling of $`𝐌`$. Consequently, the time dependence of $`\sigma `$ in Eq. (17) can be neglected and one can use $`\sigma =e^2n_q(T)/m\nu `$. Equation (17) then becomes similar to Eq. (14) and the same argument gives the surface integral that is proportional to $`(\dot{𝐌}^2+\dot{M}_z^2)[\dot{\theta }^2+(\dot{\theta }^2+\dot{\varphi }^2)\mathrm{sin}^2\theta ]`$. Due to the smallness of $`ϵ`$, the hard-axis anisotropy in Eq. (8) is small compared to the easy-axis anisotropy. Under this condition, quasiclassical trajectories of $`𝐌`$ must be close to the easy plane. This means $`\varphi \pi /2`$, while $`\theta `$ for the tunneling trajectory changes from $`0`$ to $`\pi `$. A more rigorous analysis shows that $`\dot{\varphi }^2ϵ\dot{\theta }^2`$. Thus, with good accuracy,
$$Q=\frac{\sigma \lambda _L^3}{2c^2}d^2r\dot{𝐁}^2=\frac{3\lambda _LM^2}{16\nu d^4}\left(\frac{n_q}{n_s}\right)\dot{\theta }^2(1+\mathrm{sin}^2\theta ).$$
(19)
If Eq. (18) were quadratic in $`\dot{\theta }`$, it could be interpreted as linear dissipation with a friction coefficient $`\eta =3\lambda _LM^2n_q/16\nu d^4n_s`$. This would allow the Caldeira-Leggett approach to tunneling with dissipation. The dissipation in the rotation of $`𝐌`$ due to its interaction with quasiparticles is nonlinear in $`\theta `$, however. Nevertheless, since we only want to obtain an estimate of the effect of dissipation on tunneling, and because the $`\mathrm{sin}^2\theta `$ term in Eq. (18) can hardly change this effect significantly, we shall go ahead and estimate the effective Caldeira-Leggett action as
$$I_{CL}=\frac{\eta }{4\pi }_{\mathrm{}}^{\mathrm{}}𝑑\tau ^{}_o^{\mathrm{}}𝑑\tau \frac{[\theta (\tau )\theta (\tau ^{})]^2}{(\tau \tau ^{})^2}.$$
(20)
Note that the double-integral in Eq. (19) is dimensionless.
The measure of the effect of dissipation on tunneling is the ratio $`I_{CL}/\mathrm{}B_o`$ where $`\mathrm{}B_o`$ is the effective action in the absence of dissipation, with $`B_o`$ given by Eq. (9). After simple algebra, one obtains
$$\frac{I_{CL}}{\mathrm{}B}\frac{3\pi }{128\sqrt{ϵ}}\left(\frac{n_q}{n_s}\right)\left(\frac{\lambda _L}{d}\right)\left(\frac{l}{d}\right)^3\left(\frac{\gamma M_o}{\nu }\right).$$
(21)
In a typical experiment, one should expect $`\lambda _L<d`$ and $`l<d`$, while the ratio $`\gamma M_o/\nu `$ may hardly exceed $`10^2`$. Consequently, even at $`TT_c`$, when $`n_qn_s`$, the effect of the dissipation on the tunneling rate may become visible only at $`ϵ`$ less than $`10^5`$.
A more rigorous approach along the lines of Ref. shows that the ratio $`n_q/n_s`$ in equations (19)-(21) should be replaced by the ratio of the coherence factors. Similar to other dissipation problems due to quasiparticles , its effect on tunneling must have a maximum at $`T`$ slightly lower than $`T_c`$. Even at the maximum this effect should still be small. Coherence is a more subtle issue. It may be destroyed by a dissipative environment even if the effect of dissipation on tunneling is weak. The mechanism of dissipation discussed above goes down with temperature as $`\mathrm{exp}(\mathrm{\Delta }/T)`$, where $`\mathrm{\Delta }`$ is the superconducting gap. Consequently, preserving coherent oscillations of $`𝐌`$ between $`|>`$ and $`|>`$ states requires $`T<<\mathrm{\Delta }`$.
The work of E.M.C. has been supported by the U.S. Department of Energy through Grant No. DE-FG02-93ER45487.
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# Inflationary models with a flat potential enforced by non-abelian discrete gauge symmetries
## I Introduction
Inflation explains many basic features of our universe . It is also thought to have generated the density perturbations needed to form galaxies and all the other large scale structure in the observable universe . There are many types of inflation that are natural from the particle physics point of view.
If the energy density of our vacuum (the cosmological constant) is positive, then it will eventually give rise to inflation. Observations suggest this is just beginning now. Although this type of inflation is natural, the fact that it is just beginning now can (in the opinion of EDS) only be explained by anthropically selected fine-tuning of the cosmological constant.
If in the past the universe became trapped in a positive energy false vacuum for sufficiently long, one will get an epoch of false vacuum (old) inflation . This probably did happen, though in the unobservably distant past. A false vacuum with near Planck scale energy density could start (eternal) inflation from fairly generic initial conditions. The desirable properties (and maybe even necessity) of eternal inflation have been stressed by Linde in the context of $`\varphi ^n`$ chaotic inflationary potentials. Unfortunately, such potentials generically do not survive the inclusion of gravitational strength effects, especially for the extremely large field values needed to start eternal inflation at the Planck density. However, much the same ideas can be realized using the generic false vacuum inflation.
Thermal inflation just needs a potential $`V=V_0\frac{1}{2}m^2\varphi ^2+\mathrm{}`$ with $`mV_0^{1/4}`$, typical of supersymmetric theories. It occurs when $`\varphi `$ is held at $`\varphi =0`$ by thermal effects, and is probably needed to solve the moduli (Polonyi) problem . It also has important implications for baryogenesis and dark matter .
Rolling scalar field inflation just needs a potential $`V=V_0\frac{1}{2}m^2\varphi ^2+\mathrm{}`$ with $`mV_0^{1/2}/M_{\mathrm{Pl}}`$ where $`M_{\mathrm{Pl}}=1/\sqrt{8\pi G}2.4\times 10^{18}\mathrm{GeV}`$, typical of moduli potentials. It occurs as the inflaton $`\varphi `$ rolls off the maximum of the potential. This may also have happened.
However, observations constrain the density perturbations to be approximately scale-invariant. Therefore, the natural way to produce these is with an approximately scale-invariant inflation. The only known scale-invariant inflation is a limit of rolling scalar field inflation called slow-roll inflation . It requires the stronger condition $`mV_0^{1/2}/M_{\mathrm{Pl}}`$, or more generally
$$\left(\frac{V^{}}{V}\right)^2\frac{1}{M_{\mathrm{Pl}}^2}$$
(1)
and
$$\left|\frac{V^{\prime \prime }}{V}\right|\frac{1}{M_{\mathrm{Pl}}^2}$$
(2)
The first condition suggests we should be near a maximum, or other extremum, of the potential. The second is non-trivial . For example, many models of inflation are built ignoring gravitational strength interactions, and so are implicitly setting $`M_{\mathrm{Pl}}=\mathrm{}`$. Clearly one cannot achieve the second condition in this context. In supergravity, the potential is composed of two parts, the $`F`$-term and the $`D`$-term. If the inflationary potential energy is dominated by the $`F`$-term then one can show that
$$\frac{V^{\prime \prime }}{V}=\frac{1}{M_{\mathrm{Pl}}^2}+\text{model dependent terms}$$
(3)
Unless the model dependent terms cancel the first term, the second slow roll condition, Eq. (2) above, is violated. Thus to build a model of slow-roll inflation one must be able to control the gravitational strength corrections.
There have been various attempts at achieving slow-roll inflation naturally, which are summarized below. For extensive references on inflationary models, see, for example, .
Special forms for the Kahler potential : The $`F`$-term part of the potential is determined by the superpotential $`W`$ and the Kahler potential $`K`$. The Kahler potential contains most of the terms which make slow-roll inflation difficult. Choosing a special form for the Kahler potential combined with some other conditions can allow one to cancel off the model independent gravitational strength corrections that generically destroy slow-roll inflation. Kahler potentials of the required form arise in large radius, weak coupling limits of string theory or in models with some effective extended supersymmetry.
D-term domination of the inflationary potential energy The first $`D`$-term model of inflation was given in but the model and the motivation were different.: Naively simple, but in order to obtain the COBE normalisation one must stabilize a modulus at a very large value without the aid of $`F`$-term supersymmetry breaking.
Flattening the inflaton’s potential with quantum corrections : This is completely natural but is being tested by observations and may not succeed.
Cancellation mechanism : Here the expectation value of a Nambu-Goldstone boson is used to cancel the inflaton’s mass to produce slow-roll.
In this paper we use non-abelian discrete gauge symmetries to guarantee the flatness of the inflaton’s potential. The basic idea was presented in . Here two full inflationary models utilizing this idea are constructed, a hybrid model and a mutated hybrid model. The inflationary mechanism requires the inclusion of higher order terms in the superpotential (and Kahler potential and supersymmetric loop corrections), and quantitative calculation of the properties of the exit. As the hybrid model can have a very flat potential, it can have a low energy scale, but this also brings with it the possibility of large fluctuations during the exit which provides a stringent constraint. This inflationary mechanism has the advantage that one can work in the low energy effective field theory, without needing to know the detailed high energy theory.
In Section II we briefly review the construction of a low energy effective supergravity theory; see textbooks, for example Ref. , and references therein for further information. Readers familiar with low energy effective supergravity model building can skip this section. In Section III we describe our basic idea. In Sections IV and V we give examples of models implementing this idea. In Section VI we give our conclusions. In the Appendix we list useful properties of the non-abelian discrete group $`\mathrm{\Delta }`$(96) that we use to build the models of Sections IV and V.
## II Review of low energy effective supergravity model building
The scalar potential of a supergravity theory is specified by its full Kahler potential $`K(\varphi _i,\overline{\varphi }_i)`$, superpotential $`W(\varphi _i)`$ and $`D`$-terms. We use discrete gauge symmetries which have no associated gauge fields and hence no $`D`$-terms. In the full supergravity theory, the scalar potential is
$$V(\varphi _i)=e^K\left[\left(W_{\varphi _i}+WK_{\varphi _i}\right)K_{\varphi _i\overline{\varphi }_j}^1\left(\overline{W}_{\overline{\varphi }_j}+\overline{W}K_{\overline{\varphi }_j}\right)3|W|^2\right]+D\text{-terms}$$
(4)
where the $`\{\varphi _i\}`$ include fields in all sectors, hidden and not hidden. As there are no $`D`$-terms in our case we will not discuss them further.
Writing $`\varphi `$ for the set of fields $`\{\varphi _i\}`$, only the combination
$$G(\varphi ,\overline{\varphi })=K+\mathrm{ln}|W|^2$$
(5)
is physically relevant, and so the freedom to make a Kahler transformation remains
$$K(\varphi ,\overline{\varphi })K(\varphi ,\overline{\varphi })F(\varphi )\overline{F}(\overline{\varphi })$$
(6)
$$W(\varphi )e^{F(\varphi )}W(\varphi )$$
(7)
Thus the Kahler potential can be chosen to be independent of holomorphic and anti-holomorphic terms. The kinetic term is
$$K_{\varphi _i\overline{\varphi }_j}_\mu \varphi _i^\mu \overline{\varphi }_j$$
(8)
and so, for a standard kinetic term, the leading term in the Kahler potential will be
$$K=\overline{\varphi }_i\varphi _i+\text{higher order terms}$$
(9)
The superpotential $`W`$ consists of all holomorphic terms allowed by the symmetries. This is the expression in terms of all the fields in the theory.
For the effective field theory, only some of the fields are dynamical, and the rest are integrated out. The symmetries will dictate allowed terms for the dynamical fields. To leading order (in the flat space, i.e. the $`M_{\mathrm{Pl}}\mathrm{}`$ limit), the potential for the dynamic fields (written here as $`\varphi _i`$ as well) in the effective field theory is given by (see Eq. 7.5 of Ref. )
$`V_{\mathrm{low}\mathrm{energy},\mathrm{leading}}(\varphi _i)`$ $`=`$ $`V_0+{\displaystyle \underset{j}{}}\left|{\displaystyle \frac{W}{\varphi _j}}\right|^2+m_j^2\left|\varphi _j\right|^2+\mu \stackrel{~}{W}+\text{c.c.}`$ (10)
$`=`$ $`{\displaystyle \underset{j}{}}\left|{\displaystyle \frac{W}{\varphi _j}}\right|^2+V_{\mathrm{susy}\text{}}`$ (11)
(Ref. discusses the presence of a possible cosmological constant which in our case is the vacuum energy during inflation and is taken nonzero in order for inflation to occur.) The combination $`\stackrel{~}{W}`$ is an expansion in the dynamical fields obeying the symmetries and holomorphicity just as $`W`$ does, but the coefficients of the various allowed terms are different from those in $`W`$ (the low energy coefficients are induced by integrating out terms with the heavier fields and thus do not have any fixed relation in the low energy theory). The coefficient $`\mu `$ comes from $`W_{\varphi _{\mathrm{hid}}}`$ and so is naively the size of $`|F_{\mathrm{hid}}|=|W_{\varphi _{\mathrm{hid}}}+K_{\varphi _{\mathrm{hid}}}W|`$. The masses $`m_j^2`$ can be positive or negative with a magnitude $`\left|m_j^2\right|\text{ }>|F_{\mathrm{hid}}|^2\text{ }>V_0`$. See and references therein for more discussion.
Before going on to include higher order terms in $`M_{\mathrm{Pl}}`$, it is useful to compare the scale of supersymmetry breaking and the potential energy scale. Having supersymmetry breaking in the full theory means that $`|F_{\varphi _i}|=|W_{\varphi _i}+WK_{\varphi _i}|0`$ for some $`\varphi _i`$. We will for simplicity take this to be for one field $`\varphi _i`$ and so write $`|F|`$ for the scale of supersymmetry breaking. This field is in the hidden sector by construction. One has in the full theory during inflation
$$V=V_0>0|F|^2|W_{\varphi _{hid}}+WK_{\varphi _{hid}}|^2>3|W|^2$$
(12)
(the Kahler potential merely multiplies this out front). With a vacuum energy, $`V=V_0>0`$, as is the case during inflation, one has in the full formula that $`|F|\text{ }>|W|`$ (for a positive potential) as well as $`|F|^2\text{ }>V_0`$ (as it is decreased by $`3|W|^2`$). Note that $`V_0`$ does not have to be the dominant source of susy breaking in our models (i.e. inflation can occur below the scale of susy breaking), so $`V_0`$ does not necessarily characterize the scale of the susy breaking parameters. This large degree of freedom and uncertainty makes it difficult to relate the scale of supersymmetry breaking and the vacuum energy in a model independent way.
To include the higher order possible terms in the theory, we have contributions from the visible and hidden sector. The leading terms from supergravity are determined by leading terms in the expansion of the superpotential $`W`$ and $`V_{\mathrm{susy}\text{}}`$, both of which will be given for each model. The higher order terms nonzero around the background will be discussed as they appear with the exception of the leading term in $`V_0K`$, which has already been included in the masses. These higher order terms include terms of the form $`|W|^2`$, $`W_\varphi K_{\overline{\varphi }}\overline{W}+c.c.`$ and $`|K_\varphi W|^2`$. By construction $`F_{\mathrm{hid}}`$ has a constant leading term (for supersymmetry breaking) plus, in terms of the low energy effective theory fields, some expansion in the low energy fields. The other terms will also be expansions in the low energy fields, and all of these must obey the symmetries (as the symmetries do not mix the hidden and visible sectors by construction). The leading terms are thus
$`V`$ $`=`$ $`{\displaystyle \underset{i}{}}\left|W_{\varphi _i}\right|^2+|F|^2\times \text{(arbitrary real function of the fields)}`$ (14)
$`+W\times \text{(antiholomorphic function)}+\text{c.c.}`$
which can each be multiplied by some arbitrary real function (e.g. from the expansion of the Kahler potential). The higher order terms are usually irrelevant because they are only small corrections to existing terms.
It should be kept in mind that this effective field theory is likely to correspond to a different region of field space than we are in now. In particular, although there are expected values for vacuum energy and supersymmetry breaking today, these numbers correspond to expansion of the full high energy theory around the background we have today. In this earlier time, the vacuum energy and supersymmetry breaking terms will correspond to an expansion around a different point in field space, and thus may have very different values. As the slow roll conditions are violated, inflation ends, and the field values begin to change by large amounts. This takes the theory outside the realm of the validity of the effective field theory under consideration, which is an expansion around a particular point in field space. Thus, although there will be contributions to the potential that will be seen to subtract from the vacuum energy $`V_0`$, the ultimate value of the cosmological constant after this inflationary era is not determined within this theory. That is, although the cosmological constant is decreasing when inflation ends, the subsequent behavior of the fields causing this cancellation (and thus some specified value of $`V_0`$ after inflation ends) requires knowledge of the potential for these fields outside of the regime of the effective field theory. The very important question of the observed cosmological constant today is a major outstanding problem in theoretical particle physics and outside the issues of concern in this paper. This value is a combination of the value reached after the inflationary stage described here (determined by the minimum of the potential of the full theory at the end of inflation, unspecified in an effective field theory) plus any other contributions which arise as the universe evolves to its current day temperature and field configuration. More (although less comprehensive) information is needed about the full theory as well in order to specify what happens next, in particular details of (p)reheating (which involves fields which were not dynamical during inflation and thus did not appear in the effective field theory), or which fields are dynamical in the next era, etc.
## III The idea
One of the better early attempts to naturally achieve a flat inflaton potential was Natural Inflation . The inflaton was the pseudo-Nambu-Goldstone boson $`\theta `$ of an approximate U(1) global symmetry. The potential was of the form
$$V=ϵf(\theta )$$
(15)
where $`ϵ0`$ in the limit of exact symmetry. Thus the inflaton’s mass
$$V^{\prime \prime }ϵ$$
(16)
can be made arbitrarily small. However, in this model one can not use the U(1) global symmetry to enforce
$$\left|\frac{V^{\prime \prime }}{V}\right|\frac{1}{M_{\mathrm{Pl}}^2}$$
(17)
because $`V`$ also vanishes in the limit where the symmetry is exact.
This problem can be solved by adding a constant to the potential
$$V=V_0+ϵf(\theta )$$
(18)
in which case one could in principle make $`|V^{\prime \prime }/V|`$ arbitrarily small.
However, one must now find a way to end inflation. Inflation can end if there is some critical value of the inflaton, $`\theta =\theta _\mathrm{c}`$, at which the potential destabilizes. For example, this could happen due to a term in the potential of the form $`[\lambda ^2\varphi _0^2\mathrm{sin}^2(n\theta )m^2]\psi ^2`$ with $`\lambda \varphi _0>m`$. For $`\theta <\theta _\mathrm{c}=(1/n)\mathrm{sin}^1(m/\lambda \varphi _0)`$, the potential is unstable to $`\psi \mathrm{}`$ and the runaway of $`\psi `$ can cancel out the vacuum energy $`V_0`$. See Refs. for general discussion of the hybrid inflation mechanism. This critical value must violate the U(1) symmetry, as a particular value of $`\theta `$ is singled out. However, special values of $`\theta `$ can be consistent with a discrete subgroup of the U(1) symmetry being unbroken, $`𝐙_{2n}`$ in the above example.
Furthermore, if this discrete subgroup is gauged, it can be regarded as fundamental, with the approximate U(1) global symmetry arising as a consequence. For example, if one had fields $`\varphi _+`$ and $`\varphi _{}`$ with charges $`+1`$ and $`1`$ respectively under a $`𝐙_4`$ symmetry, then the lowest dimension (and thus dominant) invariants, $`\varphi _+\varphi _{}`$, $`\left|\varphi _+\right|^2`$, and $`\left|\varphi _{}\right|^2`$, are invariant under the extended global U(1) symmetry, while terms which explicitly break the U(1), such as $`\varphi _+^4`$, are of higher order. The exact discrete $`𝐙_4`$ symmetry thus gives rise to an approximate U(1) symmetry in the region of field space in which $`\left|\varphi _+\right|`$ and $`\left|\varphi _{}\right|`$ are small.
In order to realize the couplings necessary for the hybrid inflation mechanism, for example $`\lambda \varphi ^2\psi ^2`$, it is more natural to use a non-abelian discrete symmetry. The inflaton then corresponds to the pseudo-Nambu-Goldstone bosons, $`\mathrm{\Phi }_a/|\mathrm{\Phi }|`$, of the approximate non-abelian continuous symmetry, and the hybrid exit is implemented when the magnitude of one of the components of a representation of the symmetry reaches some critical value, for example $`\left|\mathrm{\Phi }_1\right|=\mathrm{\Phi }_\mathrm{c}`$, rather than when the phase of a field reaches some critical value, which would be the case if one were to use an abelian discrete symmetry.
For a (discrete) gauge theory to be consistent it must be anomaly free . However, only the linear anomaly conditions survive for discrete abelian gauge symmetries . For the same reasons we expect only linear anomaly conditions to survive for non-abelian discrete gauge symmetries. However, there are no linear anomaly conditions for non-abelian gauge symmetries. Therefore, by this argument, non-abelian discrete gauge symmetries should be automatically anomaly free. Of course, any other gauge symmetries in the model will have to satisfy the usual anomaly conditions.
In order to have our pseudo-Nambu-Goldstone bosons, we need a potential which spontaneously breaks the extended continuous symmetry, fixing $`|\mathrm{\Phi }|\left(_a\left|\mathrm{\Phi }_a\right|^2\right)^{1/2}`$ at some value $`\mathrm{\Phi }_0>0`$. Our non-abelian discrete gauge symmetry is then non-linearly realized on the pseudo-Nambu-Goldstone bosons which include the inflaton. This protects the inflaton mass from large corrections. In this paper, we assume a hidden sector breaks supersymmetry. This generates supersymmetry breaking terms, including a vacuum energy $`V_0`$ and masses for the scalars, in our effective potential. We then use the renormalization group running of the supersymmetry breaking mass term for $`\mathrm{\Phi }`$ to generate a potential for $`\mathrm{\Phi }`$ with non-trivial minimum $`|\mathrm{\Phi }|=\mathrm{\Phi }_0`$ . The renormalization is induced (to leading order) by low dimension couplings symmetric under the extended continuous symmetry. Thus the renormalization group masses and the potential will be symmetric under the extended continuous symmetry.
However, this potential could be obtained in several other ways. One particularly interesting possibility would be to generate the potential from strong coupling dynamics symmetric under the extended continuous symmetry, allowing the inflaton to be intimately connected with the strong coupling dynamics that presumably also generates the vacuum energy that drives the inflation. Normally it would be difficult to control the inflaton’s mass in such a context, but here it is protected by the discrete gauge symmetries.
In this paper we use the non-abelian discrete symmetry $`\mathrm{\Delta }`$(96) $``$ SU(3) described in the Appendix. However, many other choices for the non-abelian discrete symmetry are possible; for example, one could use non-abelian discrete subgroups of SU(2) which would lead to more minimal models. We use $`\mathrm{\Delta }`$(96) simply for ease of model building.
To build a model one makes a suitable choice of gauge group and representations. The symmetries strongly constrain the allowed terms in the superpotential and Kahler potential. The resulting effective field theory is determined by the gauge symmetries, the representations, the couplings, and the supersymmetry breaking parameters.
The supersymmetry breaking parameters are in principle determined by the supersymmetry breaking and how it is mediated to the relevant fields. If supersymmetry is broken at a scale $`F`$, then one expects to induce a vacuum energy $`V_0|F|^2`$, scalar mass squareds of order $`|F|^2/M^2`$, where $`M`$ parametrizes the strength of the mediation and should be somewhere in the range $`|F|^{1/2}MM_{\mathrm{Pl}}`$, etc. However, in our vacuum today supersymmetry is known to be broken at a scale $`|F|^2\text{ }>\mathrm{TeV}^4`$ but the vacuum energy is known to be $`\mathrm{\Lambda }10^{59}\mathrm{TeV}^4`$. This fine tuning of the cosmological constant will be transferred to $`V_0`$ if vacuum supersymmetry breaking is the dominant supersymmetry breaking in our effective field theory, allowing $`V_0|F|^2`$. Note that $`V_0`$ itself breaks supersymmetry and so we can not have $`V_0>|F|^2`$.
However, as we do not wish to restrict ourselves to one particular model of supersymmetry breaking, we do not specify $`F`$ or $`M`$ but rather just consider the supersymmetry breaking parameters that appear in our low energy effective theory and apply the appropriate constraints to them. Our treatment is very similar to that of the Minimal Supersymmetric Standard Model. Specifically, we have no constraint on $`V_0`$ and require the scalar mass squareds to be greater than the greater of $`V_0/M_{\mathrm{Pl}}^2`$ and $`\mathrm{TeV}^4/M_{\mathrm{Pl}}^2`$. The mechanism discussed in this paper is directed at protecting the mass of the inflaton from this relatively large value.
Other parameters in our effective field theories include couplings, which can be both dimensionless and dimensionful (for example in the first model below $`\lambda `$ is dimensionless, while $`\sigma `$ has dimensions of inverse mass). The dimensionless parameters should be of order one, unless protected by some symmetry, but the dimensionful couplings go as some inverse mass parameter in the theory, related to the fields that have been integrated out and are thus not specified. These dimensionful parameters can have a very wide range. In the units we adopt, $`M_{\mathrm{Pl}}1`$, a coupling $`\sigma 1/M`$ for $`M<M_{\mathrm{Pl}}`$ will obey $`\sigma >1`$ and can be very large. However, the mass scales integrated out should be larger than the values of the dynamical fields, i.e. the dynamical fields should have values small enough for the effective field theory expansion to remain valid, i.e. for given couplings, successively higher order terms should get smaller and smaller. Note that this requires the field values to be $`M_{\mathrm{Pl}}`$.
With a given lagrangian in hand, we then impose the constraints coming from inflation. The effective field theory needs to provide a potential flat enough for slow-roll inflation to occur (flatness), a way for it to end (exit) and viable predictions for the primordial flucutuation amplitude and tilt, including the absence of any large spikes in the spectrum on observable wavelengths . Inflation ends soon after violation of slow-roll. These constraints will determine allowed ranges for the parameters in the models.
## IV A hybrid model
We choose the gauge symmetries and fields shown in Table I. The non-abelian discrete symmetry $`\mathrm{\Delta }`$(96) is described in the Appendix. The model is anomaly free.
For this choice of symmetries and fields, the most general superpotential is
$$W=\lambda \mathrm{\Phi }\mathrm{{\rm Y}}\mathrm{\Xi }+\frac{\sigma }{2}\underset{a=1}{\overset{3}{}}\mathrm{\Phi }_a^2\mathrm{\Psi }_a^2+\frac{\rho }{2}\underset{a=1}{\overset{3}{}}\mathrm{\Psi }_a^2\mathrm{{\rm Y}}_a\mathrm{\Xi }_a$$
(19)
plus dimension 6 and higher terms. Here, and throughout most of the rest of the paper, we have set $`M_{\mathrm{Pl}}=1`$ (not $`M_{\mathrm{Pl}}=\mathrm{}`$!). Some other sector breaks supersymmetry, and in our low energy effective field theory gives rise to the following general supersymmetry breaking terms:
$`V_{\mathrm{susy}\text{}}`$ $`=`$ $`V_0+\stackrel{~}{m}_\mathrm{\Phi }^2\left|\mathrm{\Phi }\right|^2m_\mathrm{\Psi }^2\left|\mathrm{\Psi }\right|^2+m_\mathrm{{\rm Y}}^2\left|\mathrm{{\rm Y}}\right|^2+m_\mathrm{\Xi }^2\left|\mathrm{\Xi }\right|^2`$ (21)
$`\left(\mu _\lambda \mathrm{\Phi }\mathrm{{\rm Y}}\mathrm{\Xi }+\mu _\sigma {\displaystyle \underset{a=1}{\overset{3}{}}}\mathrm{\Phi }_a^2\mathrm{\Psi }_a^2+\mu _\rho {\displaystyle \underset{a=1}{\overset{3}{}}}\mathrm{\Psi }_a^2\mathrm{{\rm Y}}_a\mathrm{\Xi }_a+\text{c.c.}\right)`$
plus dimension 6 and higher terms. Here $`\stackrel{~}{m}_\mathrm{\Phi }^2\left(\left|\mathrm{\Phi }\right|\right)`$ is the SU(3) symmetric renormalization group running mass squared of $`\mathrm{\Phi }`$ induced by the SU(3) symmetric coupling $`\lambda \mathrm{\Phi }\mathrm{{\rm Y}}\mathrm{\Xi }`$ in the superpotential. We assume that $`\stackrel{~}{m}_\mathrm{\Phi }^2\left(\left|\mathrm{\Phi }\right|\right)\left|\mathrm{\Phi }\right|^2`$ has a minimum at $`\left|\mathrm{\Phi }\right|=\mathrm{\Phi }_0`$. We also assume that $`m_\mathrm{\Psi }^2>0`$, $`m_\mathrm{{\rm Y}}^2>0`$, and $`m_\mathrm{\Xi }^2>0`$. As mentioned earlier, generically the masses squared have magnitude greater than or equal to $`V_0`$ due to supergravity corrections. Recall that $`V_0`$ is the vacuum energy at this time which is not necessarily equal to the scale of supersymmetry breaking.
We consider the minimum in field space corresponding to the background with $`\mathrm{{\rm Y}}=\mathrm{\Xi }=0`$. The symmetries guarantee that this background is an extremum and one can verify explicitly that it is stable if $`\left|\mathrm{\Phi }\right|>\left|\mu _\lambda /\lambda ^2\right|`$ or $`\left|\lambda \right|^2\left(m_\mathrm{{\rm Y}}^2+m_\mathrm{\Xi }^2\right)>\left|\mu _\lambda \right|^2`$. We assume that $`\mathrm{\Phi }`$ is located in the neighborhood of $`|\mathrm{\Phi }|=\mathrm{\Phi }_0`$ and replace the term $`\stackrel{~}{m}_\mathrm{\Phi }^2\left(\left|\mathrm{\Phi }\right|\right)\left|\mathrm{\Phi }\right|^2`$ by $`m_\mathrm{\Phi }^2\left(\left|\mathrm{\Phi }\right|\mathrm{\Phi }_0\right)^2`$. The leading terms are now
$$W=\frac{\sigma }{2}\underset{a=1}{\overset{3}{}}\mathrm{\Phi }_a^2\mathrm{\Psi }_a^2$$
(22)
and
$$V_{\mathrm{susy}\text{}}=V_0+m_\mathrm{\Phi }^2\left(\left|\mathrm{\Phi }\right|\mathrm{\Phi }_0\right)^2m_\mathrm{\Psi }^2\left|\mathrm{\Psi }\right|^2\left(\mu _\sigma \underset{a=1}{\overset{3}{}}\mathrm{\Phi }_a^2\mathrm{\Psi }_a^2+\text{c.c.}\right)$$
(23)
Note that the $`D`$-term is zero. The term $`m_\mathrm{\Phi }^2\left(\left|\mathrm{\Phi }\right|\mathrm{\Phi }_0\right)^2`$ will constrain $`\mathrm{\Phi }`$ to lie on the sphere $`\left|\mathrm{\Phi }\right|=\mathrm{\Phi }_0`$. The lowest order terms in the potential are then
$$V=V_0+\underset{a=1}{\overset{3}{}}\left[\left|\sigma \right|^2\left|\mathrm{\Phi }_a\right|^4\left|\mathrm{\Psi }_a\right|^2+\left|\sigma \right|^2\left|\mathrm{\Phi }_a\right|^2\left|\mathrm{\Psi }_a\right|^4\left(\mu _\sigma \mathrm{\Phi }_a^2\mathrm{\Psi }_a^2+\text{c.c.}\right)m_\mathrm{\Psi }^2\left|\mathrm{\Psi }_a\right|^2\right]$$
(24)
with the constraint $`\left|\mathrm{\Phi }\right|=\mathrm{\Phi }_0`$.
This is a hybrid inflation type potential. When
$$\left|\mathrm{\Phi }_a\right|>\mathrm{\Phi }_\mathrm{c}\sqrt{\frac{\alpha m_\mathrm{\Psi }}{\left|\sigma \right|}},a=1,2,3$$
(25)
$`\mathrm{\Psi }`$ is constrained to zero, leaving the potential
$$V=V_0$$
(26)
with the constraint $`\left|\mathrm{\Phi }\right|=\mathrm{\Phi }_0`$. When one of the $`\left|\mathrm{\Phi }_a\right|`$ drops below $`\mathrm{\Phi }_\mathrm{c}`$, the potential becomes unstable to $`\left|\mathrm{\Psi }_a\right|\mathrm{}`$. This may cause inflation to rapidly end, see Section IV A, or there could be more inflation as $`\left|\mathrm{\Psi }_a\right|\mathrm{}`$, see Section IV B. We have assumed
$$\mathrm{\Phi }_0>\sqrt{3}\mathrm{\Phi }_\mathrm{c}$$
(27)
The constant $`\alpha `$ is given by
$$\alpha =\sqrt{1+\left(\frac{\left|\mu _\sigma \right|}{\left|\sigma \right|m_\mathrm{\Psi }}\right)^2}+\frac{\left|\mu _\sigma \right|}{\left|\sigma \right|m_\mathrm{\Psi }}$$
(28)
We expect $`\left|\mu _\sigma \right|\left|\sigma \right|m_\mathrm{\Psi }`$ so that $`\alpha 1`$.
The potential, Eq. (26), is flat with respect to the Nambu-Goldstone bosons $`\mathrm{\Phi }_a/\left|\mathrm{\Phi }\right|`$. The gauge symmetries have forbidden any terms which might produce a large mass for the inflaton. However, the higher dimension terms in the Kahler potential and superpotential that are allowed by our gauge symmetries but which violate the approximate global symmetry, and that we have neglected up to now, will generate a gentle slope. The relevant higher dimension invariants are $`\mathrm{\Phi }_1^2\mathrm{\Phi }_2^2\mathrm{\Phi }_3^2`$, $`_a\left|\mathrm{\Phi }_a\right|^4`$, and $`_{ab}\left|\mathrm{\Phi }_a\right|^2\left|\mathrm{\Phi }_b\right|^2`$, which generate the terms
$$W=\mathrm{}+\frac{1}{2}\nu \mathrm{\Phi }_1^2\mathrm{\Phi }_2^2\mathrm{\Phi }_3^2$$
(29)
and
$$V_{\mathrm{susy}\text{}}=\mathrm{}+m_1^2\underset{a}{}\left|\mathrm{\Phi }_a\right|^4+m_2^2\underset{ab}{}\left|\mathrm{\Phi }_a\right|^2\left|\mathrm{\Phi }_b\right|^2\left(\mu _\nu \mathrm{\Phi }_1^2\mathrm{\Phi }_2^2\mathrm{\Phi }_3^2+\text{c.c.}\right)$$
(30)
Now for $`\left|\mathrm{\Phi }\right|=\mathrm{\Phi }_0`$ we have
$$\underset{a}{}\left|\mathrm{\Phi }_a\right|^4=\mathrm{\Phi }_0^42\underset{ab}{}\left|\mathrm{\Phi }_a\right|^2\left|\mathrm{\Phi }_b\right|^2$$
(31)
and so
$`V`$ $`=`$ $`V_0+m_1^2\mathrm{\Phi }_0^4+m_K^2{\displaystyle \underset{ab}{}}\left|\mathrm{\Phi }_a\right|^2\left|\mathrm{\Phi }_b\right|^2\left(\mu _\nu \mathrm{\Phi }_1^2\mathrm{\Phi }_2^2\mathrm{\Phi }_3^2+\text{c.c.}\right)`$ (33)
$`+\left|\nu \right|^2\left|\mathrm{\Phi }_1\right|^2\left|\mathrm{\Phi }_2\right|^2\left|\mathrm{\Phi }_3\right|^2{\displaystyle \underset{ab}{}}\left|\mathrm{\Phi }_a\right|^2\left|\mathrm{\Phi }_b\right|^2`$
where $`m_K^2m_2^22m_1^2`$.
We assume the terms derived from the non-holomorphic invariants dominate over the ones derived from the holomorphic invariant. This can be ensured either by adding extra symmetry to the model, which could set $`\nu =0`$, or just by being in the appropriate region of parameter space ($`m_K^2\mu _\nu \mathrm{\Phi }_0^2`$). We also require $`m_1^2\mathrm{\Phi }_0^4V_0`$. In order for the non-holomorphic term to drive the inflaton towards the hybrid exit to inflation, we require $`m_K^2>0`$. Then
$$V=V_0+m_K^2\underset{ab}{}\left|\mathrm{\Phi }_a\right|^2\left|\mathrm{\Phi }_b\right|^2$$
(34)
with the constraint $`\left|\mathrm{\Phi }\right|=\mathrm{\Phi }_0`$. For simplicity, we assume<sup>§</sup><sup>§</sup>§If instead $`\left|\mathrm{\Phi }_1\right|^2\left|\mathrm{\Phi }_2\right|^2\left|\mathrm{\Phi }_3\right|^2`$, the dynamics of $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$ do not decouple. $`\left|\mathrm{\Phi }_1\right|^2,\left|\mathrm{\Phi }_2\right|^2\left|\mathrm{\Phi }_3\right|^2`$. Then
$$V=V_0+m_K^2\mathrm{\Phi }_0^2\underset{a=1}{\overset{2}{}}\left|\mathrm{\Phi }_a\right|^2$$
(35)
Quantum corrections will also generate a small slope
$`V_{1\mathrm{l}\mathrm{o}\mathrm{o}\mathrm{p}}`$ $`=`$ $`{\displaystyle \frac{1}{64\pi ^2}}\text{Str}^4\mathrm{ln}{\displaystyle \frac{^2}{\mathrm{\Lambda }^2}}`$ (36)
$`=`$ $`{\displaystyle \frac{1}{64\pi ^2}}{\displaystyle \underset{a=1}{\overset{3}{}}}\{\left[\right|\sigma |^2|\mathrm{\Phi }_a|^4m_\mathrm{\Psi }^2+2|\mu _\sigma \left|\right|\mathrm{\Phi }_a|^2]^2\mathrm{ln}{\displaystyle \frac{\left|\sigma \right|^2\left|\mathrm{\Phi }_a\right|^4m_\mathrm{\Psi }^2+2\left|\mu _\sigma \right|\left|\mathrm{\Phi }_a\right|^2}{\mathrm{\Lambda }^2}}`$ (39)
$`+\left[\left|\sigma \right|^2\left|\mathrm{\Phi }_a\right|^4m_\mathrm{\Psi }^22\left|\mu _\sigma \right|\left|\mathrm{\Phi }_a\right|^2\right]^2\mathrm{ln}{\displaystyle \frac{\left|\sigma \right|^2\left|\mathrm{\Phi }_a\right|^4m_\mathrm{\Psi }^22\left|\mu _\sigma \right|\left|\mathrm{\Phi }_a\right|^2}{\mathrm{\Lambda }^2}}`$
$`2\left[\right|\sigma |^2\left|\mathrm{\Phi }_a|^4\right]^2\mathrm{ln}{\displaystyle \frac{\left|\sigma \right|^2\left|\mathrm{\Phi }_a\right|^4}{\mathrm{\Lambda }^2}}\}`$
$`=`$ $`{\displaystyle \frac{\left|\sigma \right|^4}{64\pi ^2}}{\displaystyle \underset{a=1}{\overset{3}{}}}\{\left[\left(\right|\mathrm{\Phi }_a|^2\alpha ^2\mathrm{\Phi }_\mathrm{c}^2)\left(\right|\mathrm{\Phi }_a|^2+\mathrm{\Phi }_\mathrm{c}^2)\right]^2\mathrm{ln}{\displaystyle \frac{\left|\sigma \right|^2\left(\left|\mathrm{\Phi }_a\right|^2\alpha ^2\mathrm{\Phi }_\mathrm{c}^2\right)\left(\left|\mathrm{\Phi }_a\right|^2+\mathrm{\Phi }_\mathrm{c}^2\right)}{\mathrm{\Lambda }^2}}`$ (42)
$`+\left[\left(\left|\mathrm{\Phi }_a\right|^2\mathrm{\Phi }_\mathrm{c}^2\right)\left(\left|\mathrm{\Phi }_a\right|^2+\alpha ^2\mathrm{\Phi }_\mathrm{c}^2\right)\right]^2\mathrm{ln}{\displaystyle \frac{\left|\sigma \right|^2\left(\left|\mathrm{\Phi }_a\right|^2\mathrm{\Phi }_\mathrm{c}^2\right)\left(\left|\mathrm{\Phi }_a\right|^2+\alpha ^2\mathrm{\Phi }_\mathrm{c}^2\right)}{\mathrm{\Lambda }^2}}`$
$`2\left|\mathrm{\Phi }_a|^8\mathrm{ln}{\displaystyle \frac{\left|\sigma \right|^2\left|\mathrm{\Phi }_a\right|^4}{\mathrm{\Lambda }^2}}\right\}`$
For $`\left|\mathrm{\Phi }\right|^2=\mathrm{\Phi }_0^2\mathrm{\Phi }_\mathrm{c}^2`$ and $`\left|\mathrm{\Phi }_1\right|^2,\left|\mathrm{\Phi }_2\right|^2\left|\mathrm{\Phi }_3\right|^2`$, this gives
$`V_{1\mathrm{l}\mathrm{o}\mathrm{o}\mathrm{p}}={\displaystyle \frac{\left(4\alpha ^2\alpha ^2\right)}{16\pi ^2}}\mathrm{ln}\left({\displaystyle \frac{\left|\sigma \right|^2\mathrm{\Phi }_0^4}{\mathrm{\Lambda }^2}}\right)\left|\sigma \right|^2m_\psi ^2\mathrm{\Phi }_0^2{\displaystyle \underset{a=1}{\overset{2}{}}}\left|\mathrm{\Phi }_a\right|^2`$ (43)
This can be absorbed into Eq. (35) if
$$m_K^2\left|\sigma \right|^2m_\mathrm{\Psi }^2$$
(44)
We assume $`\left|\mathrm{\Phi }_1\right|\left|\mathrm{\Phi }_2\right|`$ so that $`\left|\mathrm{\Phi }_1\right|`$ controls the end of inflation and so is the relevant degree of freedom. Defining $`\varphi =\sqrt{2}\left|\mathrm{\Phi }_1\right|`$, $`\psi =\sqrt{2}\left|\mathrm{\Psi }_1\right|`$, and $`\varphi _\mathrm{c}=\sqrt{2}\mathrm{\Phi }_\mathrm{c}`$, and reintroducing the hybrid exit terms, Eq. (24), (with phases relaxed and irrelevant terms dropped), we get our effective model of inflation
$$V=V_0+\frac{1}{2}m_K^2\mathrm{\Phi }_0^2\varphi ^2+\frac{1}{2}\left(\frac{1}{4}\left|\sigma \right|^2\varphi ^4\left|\mu _\sigma \right|\varphi ^2m_\mathrm{\Psi }^2\right)\psi ^2$$
(45)
There are two possibilities for when astronomically observable scales could leave the horizon during inflation; either at $`\varphi >\varphi _\mathrm{c}`$ or at $`\varphi <\varphi _\mathrm{c}`$. The former requires a quick hybrid exit in order to avoid possible problems with a spike in the density perturbation spectrum at $`\varphi =\varphi _\mathrm{c}`$ . The latter occurs in the opposite limit of a slow exit.
### A Fast exit
Here astronomically observable scales leave the horizon when $`\varphi >\varphi _\mathrm{c}`$. The slow roll conditions are satisfied if $`m_K^2\mathrm{\Phi }_0^2V_0`$. The number of $`e`$-folds until $`\varphi =\varphi _\mathrm{c}`$ is
$$N=_t^{t_\mathrm{c}}H𝑑t_{\varphi _\mathrm{c}}^\varphi \frac{V}{V^{}}𝑑\varphi =\frac{V_0}{m_K^2\mathrm{\Phi }_0^2}\mathrm{ln}\frac{\varphi }{\varphi _\mathrm{c}}$$
(46)
The COBE normalization gives
$$\frac{V^{3/2}}{V^{}}=\frac{V_0^{3/2}}{m_K^2\mathrm{\Phi }_0^2\varphi }=\frac{V_0^{3/2}}{m_K^2\mathrm{\Phi }_0^2\varphi _\mathrm{c}}\mathrm{exp}\left(\frac{m_K^2\mathrm{\Phi }_0^2N}{V_0}\right)=6\times 10^4$$
(47)
Substituting in for $`\varphi _c=\sqrt{2}\mathrm{\Phi }_c`$ and using Eq. (25), this can be rewritten
$$V_0^{1/4}=10^3\left(\frac{\alpha }{\left|\sigma \right|}\right)^{1/2}\left(\frac{m_\mathrm{\Psi }^2}{V_0}\right)^{1/4}\left(\frac{m_K^2\mathrm{\Phi }_0^2}{V_0}\right)\mathrm{exp}\left(\frac{m_K^2\mathrm{\Phi }_0^2N}{V_0}\right)$$
(48)
The spectral index is
$$n=1+2\frac{V^{\prime \prime }}{V}=1+\frac{2m_K^2\mathrm{\Phi }_0^2}{V_0}$$
(49)
A quick hybrid exit avoids problems at $`\varphi \varphi _c`$, caused by $`\psi `$’s fluctuations leading to too large a spike in the density perturbation spectrum, by making the time at which inflation ends effectively controlled by $`\varphi `$’s classical motion rather than by $`\psi `$’s stochastic fluctuations. The rough idea is that $`\psi `$’s effective mass squared goes from $`H^2V_0`$ to $`H^2V_0`$ in a time-scale short compared with the Hubble time so that the stochastic fluctuations in $`\psi `$, which do actually cause the end of inflation, do not lead to large fluctuations in the number of e-folds of expansion, and so do not lead to large density perturbations. In terms of parameters this means
$$\frac{dM_\psi ^2}{dN}|_{\varphi =\varphi _\mathrm{c}}=\frac{dM_\psi ^2}{d\varphi }|_{\varphi =\varphi _\mathrm{c}}\frac{d\varphi }{dN}|_{\varphi =\varphi _\mathrm{c}}=\frac{2\left(\alpha ^2+1\right)m_\mathrm{\Psi }^2m_K^2\mathrm{\Phi }_0^2}{V_0}V_0$$
(50)
where $`M_\psi ^2=\frac{1}{4}\left|\sigma \right|^2\varphi ^4\left|\mu _\sigma \right|\varphi ^2m_\mathrm{\Psi }^2`$ is the effective mass of $`\psi `$.
This constraint, when combined with the others mentioned above, severely restricts the parameter space. However, pushing things to the limit, one can still come up with interesting numbers. For example, taking $`\mathrm{\Phi }_0=10^{3.5}\left|\sigma \right|^1`$, $`m_\mathrm{\Psi }=10^8\left|\sigma \right|^1`$, and $`m_K=10^8`$ gives $`V_0^{1/4}=10^5\left|\sigma \right|^{1/2}`$ and $`n=1.002`$. Taking $`\left|\sigma \right|=10^8`$ would then give $`m_\mathrm{\Psi }=10^{16}200\mathrm{GeV}`$ and $`V_0^{1/4}=10^92\times 10^9\mathrm{GeV}`$.
### B Slow exit
When $`\varphi \varphi _\mathrm{c}`$, $`\psi `$’s mass is partially canceled This is similar to the scenario of Ref. in which the expectation value of a Nambu-Goldstone boson is used to cancel off the mass of the inflaton. Our scenario has very different parameters, which leads to different terms dominating the potential when observable scales leave the horizon during inflation. allowing $`\psi `$ to slow-roll in addition to $`\varphi `$. Here astronomically observable scales leave the horizon when $`\varphi <\varphi _\mathrm{c}`$. Define $`\phi =\varphi _\mathrm{c}\varphi `$. Then
$$V=V_0m_K^2\mathrm{\Phi }_0^2\varphi _\mathrm{c}\phi \frac{\left(\alpha ^2+1\right)m_\psi ^2}{\varphi _\mathrm{c}}\phi \psi ^2+𝒪\left(\frac{\phi ^2}{\varphi _\mathrm{c}^2}\right)$$
(51)
Note that when $`\phi `$ becomes of order $`\varphi _\mathrm{c}`$, $`\psi `$’s mass is no longer suppressed and inflation ends rapidly, if it has not already ended. Thus $`\phi \varphi _\mathrm{c}`$ will be a good approximation during inflation. The slow-roll equations of motion are
$$\frac{d\phi }{dN}=\frac{m_K^2\mathrm{\Phi }_0^2\varphi _\mathrm{c}}{V_0}\frac{\left(\alpha ^2+1\right)m_\psi ^2}{V_0\varphi _\mathrm{c}}\psi ^2$$
(52)
$$\frac{d\psi }{dN}=2\frac{\left(\alpha ^2+1\right)m_\psi ^2}{V_0\varphi _\mathrm{c}}\phi \psi $$
(53)
where
$$N=_t^{t_\mathrm{e}}H𝑑t$$
(54)
is the number of $`e`$-folds until the end of inflation. Once
$$\psi ^2\frac{m_K^2\mathrm{\Phi }_0^2\varphi _\mathrm{c}^2}{\left(\alpha ^2+1\right)m_\psi ^2}$$
(55)
one can solve this system of equations to give
$$\frac{1}{2}\psi ^2=\phi ^2+A^2$$
(56)
where $`A`$ is a constant. Substituting this into Eq. (52) and integrating gives
$$N=\frac{V_0\varphi _\mathrm{c}}{2\left(\alpha ^2+1\right)Am_\psi ^2}\left[\mathrm{tan}^1\frac{A}{\phi }\mathrm{tan}^1\frac{A}{\phi _\mathrm{e}}\right]$$
(57)
Therefore, once $`\phi `$ and $`\psi `$ have rolled to values much greater than $`A`$, we have
$$\phi \frac{1}{\sqrt{2}}\psi \frac{V_0\varphi _c}{2(\alpha ^2+1)m_\psi ^2N}$$
(58)
Therefore, in terms of $`N`$, the condition Eq. (55) translates to
$$2\left(\alpha ^2+1\right)N^2m_\psi ^2m_K^2\mathrm{\Phi }_0^2V_0^2$$
(59)
i.e. the limit opposite to that of Eq. (50) of the previous section.
Because both $`\phi `$ and $`\psi `$ are slow-rolling, we need to use the method of Ref. to calculate the density perturbations.Note that the dangerous spike in the density perturbations produced at $`\varphi \varphi _\mathrm{c}`$, i.e. $`\phi 0`$, is inflated to unobservably large scales by the inflation that occurs at $`\phi >0`$. Our direct calculation shows that the density perturbations are acceptable on observable scales. The physics behind this method is very intuitive. Stochastic fluctuations in the scalar fields lead to perturbations in the number of $`e`$-folds of expansion. Perturbations in the number of $`e`$-folds of expansion then induce curvature perturbations. Finally, once these curvature perturbations re-enter the horizon after inflation, they are naturally reinterpreted as density perturbations. Now from Eq. (57)
$$N=\frac{V_0\varphi _\mathrm{c}}{2\left(\alpha ^2+1\right)m_\psi ^2\phi }\left[1\frac{A^2}{3\phi ^2}+𝒪\left(\frac{A^4}{\phi ^4}\right)+𝒪\left(\frac{\phi }{\phi _\mathrm{e}}\right)\right].$$
(60)
To calculate the change in $`N`$ as $`\phi `$ and $`\psi `$ are changed, one also needs to use from Eq. (56) that
$$\frac{A}{\phi }=\frac{\phi }{A}$$
(61)
and
$$\frac{A}{\psi }=\frac{\psi }{2A}$$
(62)
that is, one also needs to take into account fluctuations between trajectories characterized by a given value of $`A`$, as well as fluctuations along a given trajectory. Therefore, including this,
$$\frac{N}{\phi }=\frac{2\left(\alpha ^2+1\right)m_\psi ^2N^2}{3V_0\varphi _\mathrm{c}}$$
(63)
$$\frac{N}{\psi }=\frac{2\sqrt{2}\left(\alpha ^2+1\right)m_\psi ^2N^2}{3V_0\varphi _\mathrm{c}}$$
(64)
The COBE normalisation is
$$\frac{H}{2\pi }\sqrt{\left(\frac{N}{\phi }\right)^2+\left(\frac{N}{\psi }\right)^2}=6\times 10^5$$
(65)
Therefore
$$V_0^{1/2}\sqrt{\left(\frac{N}{\phi }\right)^2+\left(\frac{N}{\psi }\right)^2}=\frac{2\left(\alpha ^2+1\right)m_\psi ^2N^2}{V_0^{1/2}\varphi _\mathrm{c}}=6\times 10^4$$
(66)
and so
$$V_0^{1/4}=10^7\left|\sigma \right|^{1/2}\left(\frac{2\sqrt{\alpha }}{\alpha ^2+1}\right)\left(\frac{45}{N}\right)^2\left(\frac{V_0}{m_\psi ^2}\right)^{3/4}$$
(67)
The spectral index is
$$n=1\frac{4}{N}$$
(68)
This is the same as one would get if one had a potential of the form $`V=V_0a\varphi ^3`$, for example Ref. . However, the two models can in principle be distinguished by the fact that our model does not satisfy the single component inflaton consistency condition $`n_T=bT/S`$. Here $`n_T`$ is the spectral index of the gravitational waves, $`b`$ is a constant that depends on conventions, and $`S`$ and $`T`$ are the amplitudes of the scalar perturbations and the gravitational waves, respectively . Instead we have
$$n_T=3b\frac{T}{S}$$
(69)
In practice, though, this will be impossibly difficult to measure.
An interesting feature of this model is that it can easily produce inflation at very low scales; for instance, to take an extreme example, one can get $`V_0^{1/4}=10^{14}20\mathrm{TeV}`$ with $`m_\psi 10^{24}`$ and $`\sigma 10^2`$. This would, for example, be a low enough scale to replace thermal inflation . It would also make embedding the model in the MSSM, or modest extensions thereof, plausible. However, the low scale of inflation means that less inflation is needed and so observable scales leave the horizon at relatively small values of $`N`$. This, combined with the relatively large factor of $`4`$ in Eq. (68), results in a spectral index $`n`$ which is too small to agree with observations. One can get a more viable spectral index, i.e. $`n`$ closer to $`1`$, by raising the scale of inflation; for instance taking $`V_0^{1/4}10^8`$. Other parameters are then constrained by Eqs. (27), (44), (59) and (67).
## V A mutated hybrid model
To get a mutated hybrid inflation model, one can instead take the symmetries and field content shown in Table II.
The most general superpotential is
$$W=\lambda \mathrm{\Phi }\mathrm{{\rm Y}}\mathrm{\Xi }+\frac{\sigma }{3}\underset{a=1}{\overset{3}{}}\mathrm{\Phi }_a\mathrm{\Psi }_a^3+\frac{\rho }{3}\underset{a=1}{\overset{3}{}}\mathrm{\Phi }_a\mathrm{\Psi }_a\mathrm{\Omega }_a\mathrm{\Gamma }_a$$
(70)
plus higher dimension terms, and the most general supersymmetry breaking terms are
$`V_{\mathrm{susy}\text{}}`$ $`=`$ $`V_0+\stackrel{~}{m}_\mathrm{\Phi }^2\left|\mathrm{\Phi }\right|^2m_\mathrm{\Psi }^2\left|\mathrm{\Psi }\right|^2+m_\mathrm{{\rm Y}}^2\left|\mathrm{{\rm Y}}\right|^2+m_\mathrm{\Xi }^2\left|\mathrm{\Xi }\right|^2+m_\mathrm{\Omega }^2\left|\mathrm{\Omega }\right|^2+m_\mathrm{\Gamma }^2\left|\mathrm{\Gamma }\right|^2`$ (72)
$`\left(\mu _\lambda \mathrm{\Phi }\mathrm{{\rm Y}}\mathrm{\Xi }+\mu _\sigma {\displaystyle \underset{a=1}{\overset{3}{}}}\mathrm{\Phi }_a\mathrm{\Psi }_a^3+\mu _\rho {\displaystyle \underset{a=1}{\overset{3}{}}}\mathrm{\Phi }_a\mathrm{\Psi }_a\mathrm{\Omega }_a\mathrm{\Gamma }_a+\text{c.c.}\right)`$
plus higher dimension terms. $`\mathrm{\Phi }`$’s mass squared acquires a $`\mathrm{\Phi }`$ dependence from the renormalization group running induced by the coupling $`\lambda \mathrm{\Phi }\mathrm{{\rm Y}}\mathrm{\Xi }`$ in the superpotential. Since this coupling is SU(3) symmetric, the $`\mathrm{\Phi }`$ dependence induced by it will also be SU(3) symmetric, i.e. $`\stackrel{~}{m}_\mathrm{\Phi }^2=\stackrel{~}{m}_\mathrm{\Phi }^2\left(\left|\mathrm{\Phi }\right|\right)`$. We assume $`\stackrel{~}{m}_\mathrm{\Phi }^2\left(\left|\mathrm{\Phi }\right|\right)\left|\mathrm{\Phi }\right|^2`$ has a minimum at $`\left|\mathrm{\Phi }\right|=\mathrm{\Phi }_0`$. The higher dimension, SU(3) asymmetric couplings will induce a small SU(3) asymmetric $`\mathrm{\Phi }`$ dependence in the potential. These small quantum corrections will be considered later.
The potential is minimized for $`\mathrm{{\rm Y}}=\mathrm{\Xi }=\mathrm{\Omega }=\mathrm{\Gamma }=0`$. We assume that $`\mathrm{\Phi }`$ is located in the neighborhood of $`|\mathrm{\Phi }|=\mathrm{\Phi }_0`$ and replace $`\stackrel{~}{m}_\mathrm{\Phi }^2\left(\left|\mathrm{\Phi }\right|\right)\left|\mathrm{\Phi }\right|^2`$ by $`m_\mathrm{\Phi }^2\left(\left|\mathrm{\Phi }\right|\mathrm{\Phi }_0\right)^2`$.
In this background, the model simplifies to
$$W=\frac{\sigma }{3}\underset{a=1}{\overset{3}{}}\mathrm{\Phi }_a\mathrm{\Psi }_a^3$$
(73)
and
$`V`$ $`=`$ $`V_0+m_\mathrm{\Phi }^2\left(\left|\mathrm{\Phi }\right|\mathrm{\Phi }_0\right)^2m_\mathrm{\Psi }^2\left|\mathrm{\Psi }\right|^2\left(\mu _\sigma {\displaystyle \underset{a=1}{\overset{3}{}}}\mathrm{\Phi }_a\mathrm{\Psi }_a^3+\text{c.c.}\right)`$ (75)
$`+\left|\sigma \right|^2{\displaystyle \underset{a=1}{\overset{3}{}}}\left|\mathrm{\Phi }_a\right|^2\left|\mathrm{\Psi }_a\right|^4+{\displaystyle \frac{1}{9}}\left|\sigma \right|^2{\displaystyle \underset{a=1}{\overset{3}{}}}\left|\mathrm{\Psi }_a\right|^6`$
This is a mutated hybrid inflation type potential. During inflation $`\mathrm{\Phi }`$ constrains $`\mathrm{\Psi }`$ to small but non-zero values
$$\left|\mathrm{\Psi }_a\right|=\frac{\beta m_\mathrm{\Psi }}{\sqrt{2}\left|\sigma \right|\left|\mathrm{\Phi }_a\right|}$$
(76)
where
$$\beta =\sqrt{1+\left(\frac{3\left|\mu _\sigma \right|}{2\sqrt{2}\left|\sigma \right|m_\mathrm{\Psi }}\right)^2}+\frac{3\left|\mu _\sigma \right|}{2\sqrt{2}\left|\sigma \right|m_\mathrm{\Psi }}$$
(77)
and we have neglected the $`\left|\mathrm{\Psi }_a\right|^6`$ term. The effective potential for $`\mathrm{\Phi }`$ is therefore
$$V=V_0+m_\mathrm{\Phi }^2\left(\left|\mathrm{\Phi }\right|\mathrm{\Phi }_0\right)^2\underset{a=1}{\overset{3}{}}\frac{\beta ^2\left(\beta ^2+2\right)m_\mathrm{\Psi }^4}{12\left|\sigma \right|^2\left|\mathrm{\Phi }_a\right|^2}$$
(78)
In the limit $`\left|\mathrm{\Phi }_1\right|^2\left|\mathrm{\Phi }_2\right|^2+\left|\mathrm{\Phi }_3\right|^2`$ this simplifies to
$$V=V_0\frac{\beta ^2\left(\beta ^2+2\right)m_\mathrm{\Psi }^4}{12\left|\sigma \right|^2\left|\mathrm{\Phi }_1\right|^2}$$
(79)
which is a mutated hybrid inflation potential . During inflation $`|\mathrm{\Phi }_1|`$, or more precisely the field corresponding to the trajectory Eq. (76), rolls to smaller values and eventually rolls fast enough to end inflation.
Mutated hybrid inflation has a spectral index
$$n=1\frac{3}{2N}0.97$$
(80)
and the COBE normalisation gives
$$V_0^{1/4}=\frac{10^5}{\sqrt{\left|\sigma \right|}}\left(\frac{50}{N}\right)^{3/4}\frac{m_\mathrm{\Psi }}{V_0^{1/2}}$$
(81)
## VI Discussion and Conclusions
We have discussed a mechanism to obtain potentials flat enough for slow-roll inflation in the presence of supergravity corrections, and given a hybrid and mutated hybrid example. Our context has been that of a low energy effective field theory. Discrete gauge symmetries are used to guarantee that Planck scale effects do not destroy the flatness of the potential, which is determined by the choice of gauge symmetries, representations, and signs of the supersymmetry breaking masses. Constraints on the viable models we considered were related to the mutated or hybrid exits. The exit had to be approached via the slow roll potential and additionally not generate fluctuations inconsistent with observation. As this is a only a first attempt at building models implementing this mechanism, we hope and believe it is likely that more elegant versions are possible.
One attractive feature of this way of obtaining inflation is that in principle, the inflationary scales for the hybrid models can be very low. In the specific case we looked at, the spectral index becomes unviably small as the scale of inflation is lowered, but we do not have any reason to expect this to be a generic limitation for these sorts of models. Inflation at very low scales has several advantages. For example, it might obviate the need for a round of thermal inflation , as mentioned above, to solve the moduli problem. In addition, due to the low energy scales involved, the model might have a simple relation to phenomenological particle theory models such as the minimal supersymmetric standard model. One might also be able to make some correspondence with the discrete gauge symmetries used here to obtain flatness and the discrete symmetries in various parts of the standard model and its supersymmetric extensions, for example those used for fermion masses, to suppress flavour changing neutral currents, or in certain grand unified theories.
It should be stressed that this model is in the context of an effective field theory. As a result, certain properties of the more complete theory cannot be deduced from the effective theory alone, as they are more model dependent than the inflationary mechanism and its exit described here. These include the details of (pre)heating and the value of the cosmological constant today.
On a related note, we have not discussed constraints from gravitino production in the cases where these models have a higher inflationary scale. This is primarily because, aside from the low reheating temperature case mentioned above, a short era of low scale inflation is needed to dilute the moduli, and will serve to dilute the gravitinos as well. In addition, the amount of gravitino production is strongly model dependent, and thus our effective field theory does not necessarily contain enough information to predict it. Future directions include implementing this idea for different gauge groups, and embedding an effective theory with this mechanism into a more complete model.
## Appendix
$`\mathrm{\Delta }`$(96) is the discrete subgroup of SU(3) with elements
$$X_{mn}A_{mn}X_{00}$$
(82)
where
$$A_{mn}\left(\begin{array}{ccc}i^m& 0& 0\\ 0& i^n& 0\\ 0& 0& i^{mn}\end{array}\right)$$
(83)
and
$`X_{00}`$ $``$ $`\{\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right),\left(\begin{array}{ccc}0& 1& 0\\ 0& 0& 1\\ 1& 0& 0\end{array}\right),\left(\begin{array}{ccc}0& 0& 1\\ 1& 0& 0\\ 0& 1& 0\end{array}\right),`$ (103)
$`\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& 1\end{array}\right),\left(\begin{array}{ccc}1& 0& 0\\ 0& 0& 1\\ 0& 1& 0\end{array}\right),\left(\begin{array}{ccc}0& 0& 1\\ 0& 1& 0\\ 1& 0& 0\end{array}\right)\}`$
It can be generated by
$$\{\left(\begin{array}{ccc}0& 1& 0\\ 0& 0& 1\\ 1& 0& 0\end{array}\right),\left(\begin{array}{ccc}0& i& 0\\ i& 0& 0\\ 0& 0& 1\end{array}\right)\}$$
(104)
Let $`\mathrm{\Phi }_a`$, $`\mathrm{\Psi }_a`$, $`\mathrm{{\rm Y}}_a`$, $`\mathrm{\Xi }_a`$, $`\mathrm{\Omega }_a`$, and $`\mathrm{\Gamma }_a`$ transform as fundamental representations of $`\mathrm{\Delta }`$(96), where $`a=1,2,3`$ labels the components of the representation.
The holomorphic invariants of $`\mathrm{\Delta }`$(96) are
$$\mathrm{\Phi }\mathrm{\Psi }\mathrm{{\rm Y}}\underset{a,b,c}{}ϵ_{abc}\mathrm{\Phi }_a\mathrm{\Psi }_b\mathrm{{\rm Y}}_c$$
(105)
$$\underset{a}{}\mathrm{\Phi }_a\mathrm{\Psi }_a\mathrm{{\rm Y}}_a\mathrm{\Xi }_a$$
(106)
$$\underset{abca}{}\mathrm{\Phi }_a\mathrm{\Psi }_a\mathrm{{\rm Y}}_b\mathrm{\Xi }_b\mathrm{\Omega }_c\mathrm{\Gamma }_c$$
(107)
plus dimension 7 and higher invariants.
Non-holomorphic invariants are
$$\mathrm{\Phi }^{}\mathrm{\Psi }\underset{a}{}\mathrm{\Phi }_a^{}\mathrm{\Psi }_a$$
(108)
$$\underset{a}{}\mathrm{\Phi }_a^{}\mathrm{\Psi }_a^{}\mathrm{{\rm Y}}_a\mathrm{\Xi }_a$$
(109)
$$\underset{ab}{}\mathrm{\Phi }_a^{}\mathrm{\Psi }_b^{}\mathrm{{\rm Y}}_a\mathrm{\Xi }_b$$
(110)
plus dimension 5 and higher invariants.
Note that the lowest dimension holomorphic and non-holomorphic invariants, Eqs. (105) and (108), are symmetric under the full continuous SU(3) group.
### Acknowledgements
This work was supported by the DOE and the NASA grant NAG 5-7092 at Fermilab and by Grant No. 1999-2-111-002-5 from the interdisciplinary Research Program of the KOSEF. JDC was supported in part by NSF-PHY-9800978 and NSF-PHY-9896019. EDS thanks M. White and R. Leigh at UIUC for hospitality while this work was begun, JDC thanks the Aspen Center for Physics for hospitality while this work was in progress, and we both thank the Santa Fe 99 Workshop on Structure Formation and Dark Matter for hospitality while this work was mostly completed. JDC is grateful to Martin White for discussions and Jochen Weller for a helpful question, and we both thank M. Dine for helpful suggestions on the draft.
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# Depolarisation of spherical–membrane quantum well: Gap renormalisation for closed–shell fullerenes
## I Introduction
One–electron approximation (owing to its cheapness) is often used for carbon nanoclusters though at an expense of uncontrolled incorrectness. It is worth to be aware when the simple approach fails because a true many–body calculation seems to be too complicated to apply it for any new cluster appearing in the large fullerene family. While a complete account for Coulomb interaction has to include all renormalization effects, a summing of some of diagrams can lead to the depolarization lost. A model estimation will be performed below which captures some physics, usually covered only by the sophisticated many–body theory. We stress that our term is a counterpart of the standard vertex renormalization (i.e., electron–hole attraction), as increasing the one–electron gap.
In the paper we go to reveal the depolarization correction which follows from an interaction of an electron with an electromagnetic field created by all other valence electrons of the closed–shell cluster. Therefore, the electron interaction is treated selfconsistently within the approach. This continues our consideration of C<sub>60</sub> in frame of a simple quantum mechanical model, namely, the model of the spherical–membrane quantum well (SMQW) , which has fruitful analogies with a standard quantum well model in the theory of low–dimensional structures.
The group of full rotations, SO(3), was shown to be useful in the fullerene physics to label the one–electron states, to simplify a theory of an electron–electron interaction on the sphere and surface plasmons, as well as to facilitate a computation of a high–harmonic generation spectrum in the fullerene. An essential simplification is achieved (also in the depolarization calculation) using the spherical symmetry because of the angular momentum subspaces can be separated readily in many cases. For a linear dipole response, the SO(3) symmetry cancels a lot of matrix elements and allows one to get analytically the solution for the selfconsistent RPA response function of C<sub>60</sub>. A peak of a collective excitation shows up in this spectrum, resulting from fast coherent oscillations of a total electron density of valence states.
This surface density oscillation can be thought as a confined electrical field mode or the surface plasmon (which was studied before). A zero–point oscillation of an electromagnetic vacuum is well known to manifest itself as a Casimir force between close surfaces of a polarizable substance, as a van–der–Waals interaction, as a standard Lamb shift in a hydrogen–like atom. We will consider in the paper the shift of electron levels in the field fluctuations of the confined modes (connected with the nanocluster), which effect is much stronger than of the zero–point oscillations of the free electromagnetic field.
Below we will show that, within the spherical–symmetry of SMQW model, the one–electron level shift due to the interaction with the zero–point oscillations of the electric field of the local collective mode (or the depolarization) results in a strong renormalization of a gap of the closed–shell fullerene. The relative gap increment is independent of the cluster size.
## II Perturbation theory for energy level shift
We put forward a semiclassical theory of an energy level shift (LS) for an arbitrary shell object in Ref., keeping the simplicity of the one–electron calculation and outlining the depolarization. The method follows the book. For completeness we give here the bases of the computation. The frequency of the (zero–point) oscillations of the external field is much higher than the inverse period of the electron orbit. Therefore, the adiabatic approximation has to be used. It means that the fast (field) variables can be integrated out in a motion equation for the (slow) electron. A simple model for the depolarization states that the LS results from short fast deflections of the electron from its original orbit in the random high–frequency field of the electromagnetic wave. The energy correction is given by the second order perturbation theory (see the diagram in Fig.1) and reads as:
$$\overline{\mathrm{\Delta }H}=H(r+\delta )H(r)=H\stackrel{}{\delta }+\frac{1}{2}^2H\stackrel{}{\delta }\stackrel{}{\delta }+\mathrm{}=\frac{1}{4}^2H\overline{\delta ^2}+o(\overline{\delta ^2})$$
(1)
where $`H(r)`$ is the unperturbed (one–electron) Hamiltonian and $`H(r+\delta )`$ is the Hamiltonian with account for the electron deflection $`\delta `$. The energy difference is expanded in series on $`\delta `$, then averaged over the fluctuations and a first nonzero contribution is taken. The small parameter of the perturbation theory will be proved in the end of the section. It is the ratio of the deflection to a characteristic length of the potential which estimates Laplasian, $`\overline{\delta ^2}^2\overline{\delta ^2}/R_C^21`$. We stress that the spherical symmetry allows us to limit the calculation to the subspace of the fixed angular momentum as well as to get the eigen modes of the confined field. The angular momentum plays the role of the simple momentum for a space invariant system. It conserves for definite types of diagrams (for example, the bubble diagrams). That means that the spherical plasmon modes do not mix.
The expression for the mean square of the deflection, $`\overline{\delta ^2}`$, caused by the zero–point fluctuation of the confined electromagnetic modes, was deduced in Ref.. A contribution from the free electromagnetic field, which will not be considered below, was also calculated and shown to be miserable comparing with the confined field. We will not repeat the derivation completely but refer to the expression (5) of that article, which gives the deflection from the spherical modes, $`|L`$, as follows:
$$\overline{\delta ^2}=\underset{L=1}{\overset{L_c}{}}(2L+1)\overline{\delta _L^2}\frac{\pi 2^{3/2}}{3}\left(\frac{R}{Na_B}\right)^{3/2}a_B^2\left(L_\mathrm{c}+\frac{1}{2}\right)^{3/2},$$
(2)
where $`R`$ is a nanocluster radius, $`a_B=\mathrm{}^2/me^2`$ is the Bohr radius (atomic length unit), $`L_\mathrm{c}`$ is the maximum allowed angular momentum of the plasmon state. A number of atoms of the cluster, $`N`$, reads as follows:
$$N=\frac{4\pi R^2}{3\sqrt{3}b^2/4},$$
(3)
where $`b1.4`$ Å is the carbon–carbon distance in the graphite–like lattice of the nanocluster. Using this definition we are able to evaluate the mean square deflection. For the infinitely large cluster ($`R,N\mathrm{}`$) the (infinitely large) angular momentum can be related to the (finite) 2D wave–number of the surface excitation via: $`\widehat{L}\widehat{k}R`$. The maximum wave–number $`k_{\mathrm{max}}\pi /\sqrt{3}b`$ lies on the ”Brilluene zone” boundary. Substituting the corresponding angular momentum value $`L_\mathrm{c}+1/2\pi R/(b\sqrt{3})`$ and the number of atoms into the expression (2) we get the semiclassical value of the electron deflection. In the units of the atomic length it reads as:
$$\frac{\overline{\delta ^2}}{a_B^2}=\frac{3^{5/4}}{2^{9/2}\sqrt{\pi }}\left(\frac{b}{a_B}\right)^{3/2}\left(\frac{b}{R}\right)^{3/2}\left(L_\mathrm{c}+\frac{1}{2}\right)^{3/2}=\frac{3^{1/2}\pi }{2^{9/2}}\left(\frac{b}{a_B}\right)^{3/2}1.03.$$
(4)
Though the estimation is semiquantitative, the deflection seems to be of the order of the atomic unit, which proves the expansion (1). It follows from Eq.(4) that the perturbation theory works as long as the potential changes on the scale larger than the atomic one.
The first result of the model is that the mean square deflection of the electron in the SMQW does not depend on the radius, neither on the number of atoms. It is ocularly because of the density of the valence electrons is constant (precisely, it grows slightly with $`N`$ reflecting the fact that the hexagonal carbon lattice of the spherical cluster includes 12 pentagons those lessen the density, which becomes insignificant for the large enough cluster). The independence of the deflection on the number of atoms follows from the extreme quantization both of the electron and the field mode.
## III Level ordering in SMQW and angular momentum dependent shift
Suppose that the one–electron model works for some cluster C<sub>N</sub>. To make a numerical estimation we will think about C<sub>60</sub>, which spectrum was well studied experimentally. The result does not depend essentially on the one–electron model chosen, therefore, in order to present a manifestation of the depolarization, the simplest SMQW model will be used for the bare level ordering. Then the one–electron Hamiltonian reads as :
$$H_o=E_n+\frac{\mathrm{}^2}{2mR^2}\widehat{L}^2,$$
(5)
where $`E_n`$ is the energy of a lowest level of $`n`$th radial series; an orbital quantization energy $`\mathrm{}\omega _o\mathrm{}^2/mR^2`$ defines the SO(3) level spacing between states $`|n,LM`$ which are the eigenstates of the angular momentum operator and are $`2L+1`$ degenerated. We will refer below to the single series with $`n=1`$ corresponding to $`\pi `$ electron system of the nanocluster, therefore the radial index will be omitted. It will be convenient to substitute the classical value $`L+1/2`$ for the angular momentum operator eigenvalues, which is correct for the large enough momentum.
Let us rewrite the orbital energy of the $`L`$th electron state in the following form:
$$E_L^{(o)}=\frac{\mathrm{}\omega _o}{2}\left(L+\frac{1}{2}\right)^2\frac{8\pi }{3\sqrt{3}}\left(\frac{a_B}{b}\right)^2\frac{\left(L+\frac{1}{2}\right)^2}{N}E_B=E_{\mathrm{max}}^{(o)}\frac{\left(L+\frac{1}{2}\right)^2}{N},$$
(6)
where $`E_B=e^2/a_B`$ is the atomic energy unit. The meaning of the energy $`E_{\mathrm{max}}^{(o)}=E_B(a_B/b)^28\pi /3\sqrt{3}16.8`$ eV will become clear in the end of the section. The expression (6) is derived using the surface carbon density appearing in the denominator of Eq.(3).
The electrons move within a very thin spherical shell layer (spherical membrane) which is approximated by a delta–function $`\delta (rR)`$. Hence, we use 2D–Laplasian operator in Eq.(1), which is nothing more than its angular part in the radial co–ordinate system $`\widehat{L}^2/R^2`$. Evidently, one has $`^2H_o\frac{\mathrm{}^2}{2mR^2}\frac{1}{R^2}\left(L+\frac{1}{2}\right)^4`$. Finally, Eq.(1) yields the SMQW level shift as follows:
$$\delta E_L=\frac{\pi ^3\sqrt{2}}{9\sqrt{3}}\left(\frac{a_B}{b}\right)^{5/2}\frac{\left(L+\frac{1}{2}\right)^4}{N^2}E_B,$$
(7)
here the last fraction is dimensionless. As it will be shown, it is actually independent of the nanocluster size for some characteristic $`L`$, e.g. for $`L_F`$ — Fermi momentum dividing the empty levels from the occupied ones.
The second result of the model is that the universal law for the level shift is independent on the specific nanocluster (excluding the trivial dependence in $`N`$ which also drops as will be explained in the end of the section). The depolarisation heightens the one–electron energy as:
$$E_L=E_L^{(o)}\left(1+\kappa \frac{\widehat{L}^2}{N}\right),$$
(8)
where $`\kappa 0.36`$ is the numerical coefficient depending only on the carbon atom density: $`\kappa =\sqrt{a_B/b}\pi ^2/(2^{2.5}3)`$. In the contraction limit $`LkR`$, the rest term goes to the squared wave–number $`\widehat{L}^2/N(kb)^2`$ with the accuracy of some factor. The lowest level does not shift because of the zero value of the angular momentum (in contrast to the standard Lamb shift which originates from the interaction with the charge of the hydrogen–like nuclei core. This charge is concentrated in the co–ordinate origin, therefore, the maximum Lamb shift is for the lowest, $`s`$, state). Note that the $`L=0`$ term is absent in Eq.(2) owing to no monopole plasmon exists.
Let us evaluate the maximum depolarization LS. It occurs for the maximum angular momentum $`L_{\mathrm{max}}+1/2`$ which is derived from the sum of the electron states. Because the total number of (double degenerate due to the spin) $`\pi `$states equals the number of carbon atoms, $`N=\underset{L=0}{\overset{L_{\mathrm{max}}}{}}(2L+1)=(L_{\mathrm{max}}+1)^2`$, we substitute for $`L_{\mathrm{max}}+1/2\sqrt{N}`$. Now the meaning of the energy $`E_{\mathrm{max}}^{(o)}`$, entered Eq.(6), is clear. It is the upper limit for the bare energy (see spectra in Fig.2). From Eq.(8) the maximum depolarization LS reads as:
$$\frac{\delta E_L}{E_L^{(o)}}\frac{\pi ^2}{12\sqrt{2}}\sqrt{\frac{a_B}{b}}\frac{(L_{\mathrm{max}}+1/2)^2}{N}=\kappa 0.36.$$
(9)
As we claimed before, the LS for the fixed (upper) level does not depend neither on the state label $`L`$ nor on the cluster size $`N`$. Let us now reflect on the size scaling of the shift of the state $`|L`$. The only state with $`L`$ increasing as $`\sqrt{N}`$ has a physical sense because of such momentum, scaling with the cluster size, remains in the same point of the ”Brilluene zone”. The LS of this fixed state constitutes the fixed percentage (of the bare energy) which is equal for any cluster size. The depolarisation for the different $`L`$ varies from 0 for the lowest state of a closed–shell cluster, to $`\kappa `$ for the upper one. The universality for our scaling law means that for an arbitrary cluster size and an arbitrary state momentum the relative LS falls into the same straight line as shown in Fig.2. Clearly, Eq.(9) proves that the perturbation theory is applicable as its correction is still less than unity for the highest possible level which has the maximum shift.
## IV Gap increase within SMQW–depolarization model
Within the closed–shell model the optical gap occurs between the levels $`|L_F`$ and $`|L_F+1`$ (see Fig.3) with the value:
$$E_g^{(o)}=\frac{\mathrm{}\omega _o}{2}\left[(L_F+1)(L_F+2)L_F(L_F+1)\right]=\mathrm{}\omega _o(L_F+1).$$
(10)
The gap value does depend on the cluster size, decreasing to the zero as $`N`$ going to infinity in order to approach the gapless graphite. For the buckminsterfullerene C<sub>60</sub> ($`R3.6`$ Å) the orbital energy quantum is $`\mathrm{}\omega _o0.3`$ eV, and the Fermi momentum is about $`45`$ (the uncertainty is due to the exact number of $`\pi `$ electrons is more than 50 for $`L_F=4`$ and less than 72 for $`L_F=5`$). Then the estimation for the one–electron gap, $`1.51.8`$ eV, is in a reasonable agreement with the experimental value about $`1.8`$ eV. We note that the cluster radius has to be a fitting parameter due to the 2D approximation lying in the base of the SMQW model.
The gap should increase owing to the zero–point oscillations. It is because of the higher level shifts faster. The energy difference between $`L_F`$ and $`L_F+1`$ levels reads as:
$$E_g=\mathrm{}\omega _o(L_F+1)\left(1+2\frac{(L_F+1)^2}{N}\kappa \right),$$
(11)
where the parameter $`\kappa 0.36`$ is the same as before. Within the closed–shell approximation, similarly to what done to get Eq.(9), the Fermi momentum follows from the condition: $`N=2\underset{L=0}{\overset{L_F}{}}(2L+1)=2(L_F+1)^2`$, because of the number of the occupied states is one half of the total number of states. Thus Eq.(11) becomes extremely simple and contains no fitting parameters. The gap correction is universal for any closed–shell spherical cluster and amounts about 40 % to the bare value:
$$E_g=E_g^{(o)}(1+\kappa )1.36E_g^{(o)}.$$
(12)
## V Conclusions
In the paper we deduced the semiclassical theory of the depolarization level shift (LS) in the electronic spherical–shell system for the fullerene nanoclusters. The LS materializes by the adiabatic interaction of the charge carrier with the local field of the fast zero–point oscillations of the plasmon modes confined to the cluster spherical surface. The analytical expression was derived for the depolarization LS. It is shown that the perturbation theory is applicable even for the highest (unoccupied) electron level which has the largest shift. The LS depends on the cluster size as well as on the angular momentum of the one–electron state. As a function of the bare one–electron energy all levels shifts, even for different clusters, collapse onto a straight line (Fig.2). Though, for the scaling, solely the fixed state with $`L^2/N=`$const (for example, the state at the Fermi level) has a physical sense. The relative energy correction of this fixed state is the same for a cluster of any size. The more the number of atoms, the higher the Fermi momentum, while the ratio $`L_F^2/N`$, contained in the expression for the depolarization, remains a constant number about $`1/2`$. Therefore, the Fermi level correction is independent of the size. This universal law for the LS which is non–equal for the states above and below the Fermi level (see Fig.3) gives a universal rise to the one–electron gap. As the result the renormalized gap in any closed–shell spherical cluster is wider by 1.36 times.
Acknowledgments. This work was partially supported by RFBR grants no. 96-15-96348 and 99-02-18170.
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# Determinantal Random Point Fields
## 1 Definition and General Properties of Determinantal Random Point Fields
Let $`E`$ be a one-particle space and $`X`$ a space of finite or countable configurations of particles in $`E`$. In general $`E`$ can be a separable Hausdorff space, however for our purposes it is enough to consider
$$E=\underset{j=1}{\overset{m}{}}E_j,\text{ where }E_j^d\text{ (or }^d)$$
(1.1)
If it is not mentioned specifically otherwise we always assume below $`E=^d`$ with the understanding that all resuls can be easily generalized to (1.1). We assume that each configuration $`\xi =(x_i),x_iE,i^1`$ (or $`_+^1`$ if $`d>1`$), is locally finite, that is for every compact $`KE\mathrm{\#}_K(\xi )=\mathrm{\#}(x_iK)`$ is finite. The particles in $`\xi `$ are ordered in some natural way, e.g., $`x_ix_{i+1}`$ for $`d=1`$, and if $`d>1`$ then either $`x_i=x_{i+1}`$, or
$$\begin{array}{c}\hfill |x_i|=\left(\underset{j=1}{\overset{d}{}}(x_i^{(j)})^2\right)^{\frac{1}{2}}<|x_{i+1}|=\left(\underset{j=1}{\overset{d}{}}(x_{i+1}^{(j)})^2\right)^{\frac{1}{2}}\end{array}$$
(1.2)
where $`x_i=(x_i^{(1)},\mathrm{},x_i^{(d)})`$, or $`|x_i|=|x_{i+1}|\text{ and there exists }1rd`$ such that $`x_i^{(j)}x_{i+1}^{(j)},1jr1`$, and $`x_i^{(r)}<x_{i+1}^{(r)}`$
To define a $`\sigma `$-algebra of measurable subsets of $`X`$ we first construct the so-called cylinder sets. Let $`BE`$ be any bounded Borel set and $`n0`$. We call $`C_n^B=\{\xi X:\mathrm{\#}_B(\xi )=n\}`$ a cylinder set. We define $``$ as a $`\sigma `$-algebra generated by all cylinder sets (i.e., $``$ is a minimal $`\sigma `$-algebra that contains all $`C_n^B`$).
Definition 1. A random point field is a triplet $`(X,,P)`$ where $`P`$ is a probability measure on $`(X,)`$.
This definition raises a natural question, namely how one can construct such probability measures. The corresponding theory was developed by Lenard in \[L1–L3\] where a general case of $`E`$ locally compact Hausdorff space satisying the second axiom of countability was studied. If $`E=^d`$ or $`^d`$ one can proceed quite naively by employing Kolmogorov’s fundamental theorem from the theory of stochastic processes (\[K\]). Let $`t`$ and $`s`$ be two vectors from $`E`$ with rational coordinates $`t=(t^{(1)},\mathrm{},t^{(d)}),s=(s^{(1)},\mathrm{},s^{(d)})`$. We denote an open rectangle $`\{x=(x^{(1)},\mathrm{}x^{(d)})E:x^{(j)}=t^{(j)}+\theta _j(s^{(j)}t^{(j)}),0<\theta _j<1,j=1,\mathrm{},d\}`$ by $`_{t,s}`$. Let us denote the family of finite unions of open, closed or semi-closed rectangles with rational $`t,s`$ by $``$. Suppose we are able to construct a joint distribution of non-negative integer-valued random variables $`\eta _D,D`$ (that we later identify with $`\mathrm{\#}_D`$) such that the following finite-additivity condition holds
$$\eta _D=\underset{i=1}{\overset{n}{}}\eta _{D_i}\text{ (a.e.)}$$
(1.3)
if $`D=_{i=1}^nD_i,D,D_i,i=1,\mathrm{}n`$. One immediately can replace (1.3) then by $`\sigma `$-additivity property
$$\eta _D=\underset{i=1}{\overset{\mathrm{}}{}}\eta _{D_i}\text{ (a.e.)},$$
(1.4)
$`D={\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}D_i`$, $`D,D_i,i=1,\mathrm{}`$, (of course the fact that $`\eta _D`$ takes only non-negative integers is essential here!).
It is then easy to see that the joint distribution of random variables $`\mathrm{\#}_D=\eta _D`$, $`D`$ with (1.3) (or (1.4) for that matter) uniquely defines a probability distribution on $`(X,)`$.
Since in many cases it is convenient to define the distribution of random variables through their moments the following definition appears natural:
Definition 2. Locally integrable function $`\rho _k:E^k_+^1`$ is called the $`k`$-point correlation function of the random point field $`(X,,P)`$ if for any disjoint bounded Borel subset $`A_1,\mathrm{},A_m`$ of $`E`$ and $`k_i_+^1`$, $`i=1,\mathrm{}m`$, $`{\displaystyle \underset{i=1}{\overset{m}{}}}k_i=k`$ the following identity holds:
$$𝔼\underset{i=1}{\overset{m}{}}\frac{(\mathrm{\#}_{A_i})!}{(\mathrm{\#}_{A_i}k_i)!}=_{A_1^{k_1}\times \mathrm{}\times A_m^{k_m}}\rho _k(x_1,\mathrm{},x_k)𝑑x_1\mathrm{}𝑑x_k$$
(1.5)
where by $`𝔼`$ we denote the mathematical expectation with respect to $`P`$. In particular $`\rho _1(x)`$ is the density of particles, since
$$𝔼\mathrm{\#}_A=_A\rho _1(x)𝑑x$$
for any bounded Borel $`AE`$. In general $`\rho _k(x_1,\mathrm{},x_k)`$ has the following probabilistic interpretation: let $`[x_1,x_i+dx_i]`$, $`i=1,\mathrm{},k`$ be infinitesimally small boxes around $`x_i`$, then $`\rho _k(x_1,x_2,\mathrm{},x_k)dx_1\mathrm{}dx_k`$ is the probability to find a particle in each of these boxes. The problem of existence and uniqueness of a random point field defined by its correlation functions was studied in \[L1–L3\]. Not very surprisingly, Lenard’s papers revealed many similarities to the classical moment problem (\[A\], \[S2\]). In particular the random point field is uniquely defined by its correlation functions if the distribution of random variables $`\{\mathrm{\#}_A\}`$ is uniquely determined by its moments. The sufficient condition for the uniqueness derived in \[L1\] reads
$$\underset{k=0}{\overset{\mathrm{}}{}}(\frac{1}{(k+j)!}_{A^{k+j}}\rho _{k+j}(x_1,\mathrm{},x_{k+j})𝑑x_1,\mathrm{}dx_{k+j})^{\frac{1}{k}}=\mathrm{}$$
(1.6)
for any bounded Borel $`AE`$ and any integer $`j0`$, however we invite the reader to check that the divergence of the series with $`j=0`$, namely
$$\underset{k=0}{\overset{\mathrm{}}{}}(\frac{1}{k!}_{A^k}\rho _k(x_1,\mathrm{},x_k)𝑑x_1,\mathrm{}dx_k)^{\frac{1}{k}}=\mathrm{}$$
(1.6’)
implies (1.6) for any $`j0`$. In \[L2\], \[L3\] Lenard obtained the necessary and sufficient condition for the existence of a random point field with the prescribed correlation functions.
Theorem 1. (Lenard)
Locally integrable functions $`\rho _k:E^k^1`$, $`k=1,2,`$ are the correlation functions of some random point field if and only if the Symmetry and Positivity Conditions below are satisfied.
1. Symmetry Condition
$`\rho _k`$ is invariant under the action of the symmetric group $`S_k`$, i.e.,
$$\rho _k(x_{\sigma (1)},\mathrm{},x_{\sigma (k)})=\rho _k(x_1,\mathrm{},x_k)$$
(1.7)
for any $`\sigma S_k`$.
2. Positivity Condition
For any finite set of measurable bounded functions $`\phi _k:E^k^1`$, $`k=0,1,\mathrm{},N`$ with compact support, such that
$$\phi _0+\underset{k=1}{\overset{N}{}}\underset{i_1\mathrm{}i_k}{}\phi _k(x_{i_1},\mathrm{},x_{i_k})0$$
(1.8)
for all $`\xi =(x_i)X`$, the next inequality must be valid:
$$\phi _0+\underset{k=1}{\overset{N}{}}_{E^k}\phi _k(x_1,\mathrm{},x_k)\rho _k(x_1,\mathrm{},x_k)𝑑x_1\mathrm{}𝑑x_k0.$$
(1.9)
The necessary part of the theorem is quite easy since both conditions have an obvious probabilistic interpretation. In particular the Positivity Condition means that the mathematical expectations of a certain class of non-negative random variables must be non-negative. The sufficient part is more elaborate and relies on an analogue of the Riesz Representation Theorem and the Riesz–Krein Extension Theorem (a close relative of the Hahn–Banach Theorem). It should be noted that Lenard established his results in a general setting when $`E`$ is locally compact Hausdorff space with the second axiom of countability.
One can obtain a slightly weaker (but still hopelessly ineffective!) variant of the Positivity Condition by approximating $`\phi _k`$ from above by step functions. Let $`𝒫_k`$ be the class of polynomials in $`k`$ variables that take non-negative values on non-negative integers. Since the polynomials $`\{{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \underset{j=0}{\overset{m_i1}{}}}(x_ij),m_i0\}`$ form a linear basis in the vector space of all polynomials in $`k`$ variables, we can represent any $`q(x_1,\mathrm{},x_k)𝒫_k`$ as
$$q(x_1,\mathrm{}x_k)=\underset{m_1,\mathrm{},m_k0}{}a_{m_1,\mathrm{},m_k}\underset{i=1}{\overset{k}{}}\underset{j=0}{\overset{m_i1}{}}(x_ij)$$
(1.10)
Positivity Condition: For any $`q𝒫_k,k1`$, any bounded Borel sets $`A_1,\mathrm{},A_kE`$, the following condition must be satisfied:
$$a_{0,\mathrm{},0}+\underset{m1}{}\underset{m_1+\mathrm{}+m_k=m}{}a_{m_1,\mathrm{},m_k}_{_{i=1}^kA_i^{m_i}}\rho _m(x_1,\mathrm{},x_m)𝑑x_1\mathrm{}𝑑x_m0.$$
(1.11)
Indeed, the l.h.s. at (1.11) is equal to
$$\begin{array}{cc}\hfill 𝔼q(\mathrm{\#}_{A_1},\mathrm{},\mathrm{\#}_{A_k})& =𝔼[a_{0,\mathrm{}0}+\underset{m1}{}\underset{m_1+\mathrm{}+m_k=m}{}a_{m_1,\mathrm{},m_k}\underset{i_1\mathrm{}i_m}{}\hfill \\ & \chi _{A_1^{m_1}\times \mathrm{}\times A_k^{m_k}}(x_{i_1},\mathrm{},x_{i_m})]\hfill \end{array}$$
(1.12)
One can notice that in a sense the Positivity Condition is similar to the condition on the moments of the integer-valued nonnegative random variable.
In our paper we will study a special class of random point fields introduced by Macchi in \[Ma\] (see also \[DVJ\]). We start with an integral operator $`K:L^2(^d)L^2(^d)`$ that we assume to be non-negative and locally trace class. The last condition means that for any compact $`B^d`$ the operator $`K\chi _B`$ is trace class, where $`\chi _B(x)`$ is an indicator of $`B`$. Therefore we have
$$K0,\mathrm{Tr}(\chi _BK\chi _B)<+\mathrm{}$$
(1.13)
The kernel of $`K`$ is defined up to a set of measure zero in $`^d\times ^d`$. For our purposes it is convenient to choose it in such a way that for any bounded measurable $`B`$ and any positive integer $`n`$
$$\mathrm{Tr}((\chi _BK\chi _B))=_BK(x,x)𝑑x$$
(1.14)
It appears that one can indeed achieve this. We start with
Lemma 1. (\[S3\], \[AvSS\] Remark 3.4) Let $`K`$ be trace class on $`L^2(^d)`$. Then its integral kernel may be chosen so that the function $`M(x,y)K(x,x+y)`$ is a continuous function of $`y`$ with values in $`L^1(^d)`$. Furthermore if $`m(y)=M(x,y)𝑑x`$, then $`\mathrm{Tr}K=m(0)=K(x,x)𝑑x`$.
Proof. We give the proof only when $`K`$ is non-negative. The general case is quite similar. Let $`\{\lambda _j\}_{j1}`$ is the set of non-zero eigenvalues of $`K`$ and $`\{\phi _j\}_{j1}`$ is the set of the corresponding eigenfunctions. The canonical form of $`K`$ (as a selfadjoint compact operator) is
$$K=\underset{j1}{}\lambda _j(\phi _j,)\phi _j$$
(1.15)
Fix $`y^d`$ and consider $`M(x,y)={\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\lambda _j\phi _j(x)\overline{\phi _j(x+y)}`$ as a function of $`x`$. Since $`\phi _j()\overline{\phi _j(+y)}_1=_^d|\phi _j(x)\overline{\phi _j(x+y)}|𝑑x\phi _j_2\phi _j_2=1`$, the series defining $`M(,y)`$ converges in $`L^1(^d)`$ for any $`y`$ and $`M(,y)_1{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\lambda _j=\mathrm{Tr}K<+\mathrm{}`$. If we now consider $`K(x,y)M(x,yx)`$, it is well defined for a.e. $`(x,y)^d\times ^d`$ and gives a kernel for $`K`$. The $`L^1`$-continuity of $`M(,y)`$ follows from $`{\displaystyle \underset{j1}{}}\lambda _j(\phi _j()\overline{\phi _j(+y_1)}\phi _j()\overline{\phi _j(+y_2)})_1{\displaystyle \underset{j=1}{\overset{N}{}}}\lambda _j\phi _j_2\phi _j(+y_1)\phi _j(+y_2)_2+{\displaystyle \underset{jN}{}}\lambda _j`$. Choosing $`N`$ sufficiently large one can make $`{\displaystyle \underset{jN}{}}\lambda _j<\frac{ϵ}{2}`$. Choosing $`y_1`$ sufficiently close to $`y_2`$ so that for each $`1jN\phi _j()\phi _j(+y_2y_1)_2\frac{ϵ}{2_{j=1}^N\lambda _j}`$ we have the first term less than $`\frac{ϵ}{2}`$ as well. $`\mathrm{}`$
With the help of Lemma 1 we derive
Lemma 2. Let $`K`$ be a non-negative locally trace class operator on $`L^2(^d)`$. Then its integral kernel can be chosen in such a way that for any bounded measurable $`B^d`$ the function
$$M_B(x,y)=(K\chi _B)(x,x+y)$$
is a continuous function of $`y`$ with values in $`L^1(B)`$. Furthermore,
$$\mathrm{Tr}(\chi _BK\chi _B)=_BK(x,x)𝑑x$$
Proof. Let $`K_n=\chi _{[n,n]^d}K\chi _{[n,n]^d}`$. By Lemma 1 one can choose the kernel $`K_n(x,y)`$, such that $`K_n(,+y)`$ is continuous in $`L^1([n,n]^d)`$-norm. We denote $`M_n(x,y)=K_n(x,x+y)`$. Since $`K_{n+1}(x,y)=K_n(x,y)`$ for almost all $`(x,y)[n,n]^d\times [n,n]^d`$ we conclude that for almost all $`|y|nM_{n+1}(x,y)=M_n(x,y)`$ for a.e. $`|x|n|y|`$. The $`L^1`$-continuity of $`M_{n+1}(,y),M_n(,y)`$ allows to replace “almost all $`|y|n`$” by “all $`|y|n`$”. Therefore for any $`y`$ the values of $`M_n(x,y)`$ eventually agree for a.e. $`x`$. We denote this value by $`M(x,y)`$. The function $`M(,y)`$ inherits local $`L^1`$-continuity from $`\{M_n(,y)\}`$ and (1.14) follows. It may be worthwhile to note that for any positive integer $`k`$ and bounded measurable $`B_1^d,\mathrm{},B_k^d`$ the kernel
$$M_B^{(k)}(x,y)=\left(\underset{k\text{times}}{(K\chi _{B_1})\mathrm{}(K\chi _{B_k})}\right)(x,x+y)$$
(1.16)
is also a continuous function of y with the values in $`L^1(^d)`$ and
$$Tr(K\chi _{B_1})\mathrm{}(K\chi _{B_k})=_{B_1\times \mathrm{}\times B_k}K(x_1,x_2)\mathrm{}K(x_k,x_1)𝑑x_1\mathrm{}𝑑x_k,$$
(1.17)
in particular
$$Tr(K\chi _B)\mathrm{}(K\chi _B)=_{B^k}K(x_1,x_2)\mathrm{}K(x_k,x_1)𝑑x_1\mathrm{}𝑑x_k.$$
(1.18)
Indeed, for any two versions of the integral kernel and $`k>1`$ the expressions in the last two integrals coincide up to a set of measure zero. Since we have already proved that there exists the kernel of
$$K\chi _{B_1}\mathrm{}K\chi _{B_k}$$
(1.19)
satisfying the $`L^1`$\- continuity condition from Lemma 2, the same condition is satisfied by any variant of the kernel.
$`\mathrm{}`$
Definition 3. A random point field in $`E`$ is called determinantal (or fermion) if its $`n`$-point correlation functions are given by
$$\rho _n(x_1,\mathrm{},x_n)=det\left(K(x_i,x_j)\right)_{1in}$$
(1.20)
In the case $`E=_{j=1}^ME_j,E_j^d`$ the definition takes the following form: Let $`K`$ be a trace class operator on $`L^2\underset{m\text{ times}}{(^d)\mathrm{}L^2(^d)}`$. Then $`K`$ has a matrix valued kernel $`(K_{rs}(x,y))_{1r,sm},x,y^d`$.
Definition 3’. A random point field in $`E`$ is called determinantal (or fermion) if its $`n`$-point correlation functions are given by
$$\rho _n(x_{11},x_{12},\mathrm{}x_{1i_1},\mathrm{},x_{m1},x_{m2},\mathrm{},x_{mi_m})=det(K_{rs}(x_{ri},x_{sj}))_{\underset{1ji_s,s=1,\mathrm{},m}{1ii_r,r=1,\mathrm{},m}},$$
(1.21)
where $`n=i_1+i_2+\mathrm{}+i_m,x_{ri}E_r,1rm,1ii_r`$.
Remark 1. If the kernel is Hermitian-symmetric then the non-negativity of $`n`$-point correlation functions implies that the kernel $`K(x,y)`$ is non-negative definite and therefore indeed $`K`$ must be a non-negative operator. It should be noted however that there exist determinantal random point fields corresponding to non-Hermitian kernels (see the remark after (1.36) and the examples in the sections 2.2 and 2.5).
Remark 2. The condition (1.13) is satisfied for all continuous non-negative definite kernels (see \[GK\], section III.10 or \[RS\], vol. III, section XI.4). In general situation when $`K(x,x)`$ is locally integrable, non-negative definiteness of $`K(x,y)`$ implies that $`K_B`$ is a Hilbert–Schmidt operator and one can use a theorem of Gohberg–Krein (\[GK\], section III.10, theorem 10.1) that claims that a non-negative Hilbert-Schmidt operator A is trace class iff
$$\overline{\underset{h0}{lim}}\frac{1}{(2h)^{2d}}\underset{j=1}{\overset{d}{}}[2h|x^jy^j|]_+A(x,y)dxdy<\mathrm{}$$
(1.22)
where $`t_+=\mathrm{max}(t,0),x=(x^1,\mathrm{},x^d),y=(y^1,\mathrm{}y^d)`$, and TrA is then given by (1.22).
An interesting generalization of determinantal random point fields, so called immanantal random point fields (processes) was introduced by Diaconis and Evans in \[DE\].
The classical formula of Fredholm (see \[S1\], Chapter 3) claims that a trace class operator with a continuous (in a usual sense) kernel satisfies
$$\mathrm{Tr}\left(^n(A)\right)=\frac{1}{n!}det\left(A(x_i,x_j)\right)_{1i,jn}dx_1,\mathrm{}dx_n$$
(1.23)
In general the kernel $`K(x,y)`$ may not be continuous, however (1.18) and the Lidskii theorem (see e.g., \[RS\], volume IV, section XIII.17 or \[S1\], Theorem 3.7) imply
$$_{B^n}K(x_1,x_2)\mathrm{}K(x_n,x_1)𝑑x_1\mathrm{}𝑑x_n=\underset{j=1}{\overset{\mathrm{}}{}}\lambda _j^n(K_B),$$
(1.24)
$$\mathrm{Tr}\left(^n(K_B)\right)=\underset{j_1<\mathrm{}<j_n}{}\lambda _{j_1}(K_B)\mathrm{}\lambda _{j_n}(K_B)$$
(1.25)
Combining (1.24) and (1.25) one arrives at
$$\mathrm{Tr}\left(^n(K_B)\right)=\frac{1}{n!}_{B^n}det\left(K(x_i,x_j)\right)_{1i,jn}dx_1,\mathrm{}dx_n$$
(1.26)
It follows then from (1.17) that
$$\begin{array}{cc}& \mathrm{Tr}\left((K\chi _{B_1})\mathrm{}(K\chi _{B_n})\right)=\hfill \\ & \frac{1}{n!}det\left(K(x_i,x_j)\chi _{B_j}(x_j)\right)_{1i,jn}dx_1\mathrm{}dx_n\hfill \end{array}$$
(1.27)
Definition 4. Let the kernel $`K`$ as in Lemma 2. We say that it defines a determinantal random point field $`(X,B,P)`$ if (1.21) holds.
Theorem 2. Let $`(X,B,P)`$ be a determinantal random point field with the kernel $`K`$. For any finite number of disjoint bounded Borel sets $`B_jE,j=1,\mathrm{},n`$, the generating function of the probability distribution of $`\mathrm{\#}_{B_j}=\mathrm{\#}\{x_iB_j\}`$ is given by
$$𝔼\underset{j=1}{\overset{n}{}}z_j^{\mathrm{\#}_{B_j}}=det(\mathrm{Id}+\chi _B\underset{j=1}{\overset{n}{}}(z_j1)K\chi _{B_j})$$
(1.28)
Remark 3. (1.28) is the equality of two entire functions. The r.h.s. of (1.28) is well defined as a Fredholm determinant of a trace class operator (see e.g., \[RS\], volume IV, section XIII.17 or \[S1\], section 3).
Recall that by definition
$$𝔼\underset{j=1}{\overset{n}{}}z^{\mathrm{\#}_{B_j}}=\underset{k_1,\mathrm{},k_n=0}{\overset{\mathrm{}}{}}P(\mathrm{\#}_{B_j}=k_j,j=1,\mathrm{},n)\underset{j=1}{\overset{n}{}}z_j^{k_j}$$
(1.29)
and
$$\begin{array}{cc}& det\left(\mathrm{Id}+\chi _B\underset{j=1}{\overset{n}{}}(z_j1)K\chi _{B_j}\right)=1+\underset{m=1}{\overset{\mathrm{}}{}}\underset{j_1,\mathrm{}j_m=1}{\overset{n}{}}\underset{\mathrm{}=1}{\overset{m}{}}\hfill \\ & (z_j_{\mathrm{}}1)\mathrm{Tr}(\chi _BK\chi _{B_{j_1}}\mathrm{}\chi _BK\chi _{B_{j_m}})\hfill \end{array}$$
(1.30)
Proof of Theorem 2. The Taylor expansion of the generating function near $`(z_1,\mathrm{},z_n)=(1,\mathrm{},1)`$ is given by
$$𝔼\underset{j=1}{\overset{n}{}}z_j^{\mathrm{\#}_{B_j}}=1+\underset{m=1}{\overset{\mathrm{}}{}}\underset{m_1+\mathrm{}+m_n=m}{}𝔼\underset{j=1}{\overset{n}{}}\frac{(\mathrm{\#}_{B_j})!}{(\mathrm{\#}_{B_j}m_j)!(m_j)!}\underset{j=1}{\overset{n}{}}(z_j1)^{m_j}$$
(1.31)
The radius of convergence of (1.30) is infinite since
$$\mathrm{Tr}(K\chi _{B_{j_1}}\mathrm{}K\chi _{B_{j_m}}\frac{1}{m!}\mathrm{Tr}(K\chi _B)^m,\text{ where }B=\underset{j=1}{\overset{n}{}}B_j.$$
(1.32)
Therefore, it is enough to show that the coefficients in the series (1.30), (1.31) coincide. The case $`n=1`$ follows then from (1.5), (1.21), (1.26). Using (1.27) instead of (1.26) we prove the case $`n1`$ as well.
Remark 4. Theorem 2 is well known in the Theory of Random Point Fields (see \[DVJ\], p. 140, exercise 5.4.9) and in the Random Matrix Theory (see \[TW1\]).
As we already mentioned above, if an operator $`K`$ defines a determinantal random point field it must be non-negative because of the non-negativity of the correlation functions. It follows from Theorem 2, formula (1.28) that $`K`$ must also be bounded from above by the identity operator, i.e., $`K1`$. Indeed, suppose $`K>1`$. Then there exists a bounded Borel $`BE`$ such that $`K_B>1+\frac{K1}{2}>1`$. Let $`\lambda _1(K_B)\lambda _2(K_B)\lambda _3(K_B)\mathrm{}`$ be the eigenvalues of $`K_B`$ and choose $`0<z_0<1`$ so that $`1+(z_01)\lambda _1(K_B)=0`$. Then $`𝔼z_0^{\mathrm{\#}_B}=_{k=1}^{\mathrm{}}P(\mathrm{\#}_B=k)z_0^k=det(\mathrm{Id}+(z_01)K_B)=`$(by Theorem XIII.106 from \[RS\])$`=_{j1}(1+(z_01)\lambda _j(K_B))=0`$. Therefore $`P(\mathrm{\#}_B=k)=0`$ for any $`k`$, a contradiction. On the other side assume $`0K1`$ and let (1.28) define what we hope to be the distribution of non-negative integer-valued random variables $`\{\mathrm{\#}_B\}`$.
Lemma 3. Let $`0K1`$ and $`K`$ be a locally trace class operator. Then (1.28) defines the distribution of non-negative integer-valued random variables $`\{\mathrm{\#}_B\}`$ with the additional property that for $`B=_{i=1}^nB_i`$
$$\mathrm{\#}_B=\underset{i=1}{\overset{n}{}}\mathrm{\#}_{B_i}(a.e.).$$
(1.33)
$`\mathrm{}`$
We need to show three things: first, that (1.28) defines some finite-dimensional distributions; second, that the finite-dimensional distributions satisfy the additivity property (1.33); and third, that the finite-dimensional distributions are consistent and therefore we can apply the Kolmogorov’s Fundamental Theorem to prove the existence of the distribution of $`\{\mathrm{\#}_B\}`$. Since the Fredholm determinant in (1.28) is 1 when $`z_i=1,i=1,\mathrm{},n`$, the first statement would follow from the non-negativity of the Taylor coefficients of the Fredholm determinant at $`z_i=0,i=1,\mathrm{},n`$. Consider $`0z_i1,i=1,\mathrm{},n`$ and assume for a moment $`K<1`$ (the case $`K=1`$ would be treated later by a limiting argument). Let $`B=_{i=1}^nB_i`$. Then $`K_B<1`$ and $`(\mathrm{Id}K_B)^1`$ is a bounded linear operator such that $`(\mathrm{Id}K_B)^1\mathrm{Id}=K_B(\mathrm{Id}K_B)^1`$ is trace class. Applying Theorem XIII.p105 from \[RS\], vol. IV we obtain
$$\begin{array}{cc}& det(\mathrm{Id}+\chi _B\underset{j=1}{\overset{n}{}}(z_j1)K\chi _{B_j})=det((\mathrm{Id}K_B)(\mathrm{Id}+\underset{j=1}{\overset{n}{}}z_j\hfill \\ & (\mathrm{Id}K_B)^1\chi _BK\chi _{B_j}))=det(\mathrm{Id}K_B)det(\mathrm{Id}+\underset{j=1}{\overset{n}{}}z_j\hfill \\ & (\mathrm{Id}K_B)^1\chi _BK\chi _{B_j})=det(\mathrm{Id}K_B)^{\mathrm{}}_{k=1}\mathrm{Tr}(^k(\underset{j=1}{\overset{n}{}}z_j\hfill \\ & (\mathrm{Id}K_B)^1\chi _BK\chi _{B_j}))=_{k_1,\mathrm{},k_n0}\frac{(k_1+\mathrm{}+k_n)!}{k_1!\mathrm{}k_n!}^n_{j=1}\hfill \\ & z_j^{k_j}det(\mathrm{Id}K_B)\mathrm{Tr}(_{j=1}^n(^{k_j}(\chi _{B_j}(\mathrm{Id}K_B)^1\chi _BK\chi _{B_j})))\hfill \end{array}$$
(1.34)
One can see from (1.34) that Taylor coefficients are, up to some positive factors, the traces of the exterior products of the non-negative operators, and, therefore, non-negative. We conclude that (1.28) defines some finite-dimensional distributions. Since $`𝔼_{i=1}^nz^{\mathrm{\#}_{B_i}}=det(\mathrm{Id}+\chi _B_{i=1}^n(z1)K\chi _{B_i})=det(\mathrm{Id}+(z1)K_B)=𝔼z^{\mathrm{\#}_B}`$, we conclude that $`\mathrm{\#}_B=_{i=1}^n\mathrm{\#}_{B_i}`$ (a.e.). The formula (1.28) defined the finite dimensional distributions of $`\mathrm{\#}_{B_i}`$ for disjoint compact sets. In the case of non-empty self-intersections one represents $`B_i`$ as $`C_{k_i}`$, where $`\{C_k\}`$ are disjoint sets, defines distributions of $`\mathrm{\#}_{C_k}`$ and then uses the additivity property (1.33) to define the distributions of $`\mathrm{\#}_{B_i}`$. To prove the consistency of the finite-dimensional distributions we note that (1.33) allows us to check it only for the disjoint $`B_1,\mathrm{},B_{n+1}`$. But then it trivially follows from $`det(\mathrm{Id}+\chi _B_{j=1}^n(z_j1)K\chi _{B_j}+\chi _B(11)K\chi _{B_{n+1}})=det(\mathrm{Id}+\chi _B_{j=1}^n(z_j1)K\chi _{B_j})`$. The case $`K<1`$ is proven. Now let $`K=1`$. Denote by $`K^{(ϵ)}:=K(1ϵ),ϵ>0`$ and $`\mathrm{\#}_B^{(ϵ)}`$ the random variables corresponding to the kernel $`K^{(ϵ)}`$. Since $`K^{(ϵ)}<1`$ the arguments above establish the result of Lemma 3 for $`K^{(ϵ)}`$. It is an easy exercise to see that $`𝔼_{i=1}^nz_i^{\mathrm{\#}_{B_i}^{(ϵ)}}=1+_{m=1}^{\mathrm{}}_{m_1+\mathrm{}+m_n=m}𝔼_{j=1}^n\frac{(\mathrm{\#}_{B_j^{(ϵ)}})!}{(\mathrm{\#}_{B_j^{(ϵ)}}m_j)!(m_j)!}_{j=1}^n(z_11)^{m_j}`$ uniformly converges with all derivatives to $`𝔼_{i=1}^nz_i^{\mathrm{\#}_{B_i}}`$ on compact sets as $`ϵ0`$.
Lemma 3 is proven.
The results above prove
Theorem 3. Hermitian locally trace class operator $`K`$ on $`L^2(E)`$ defines a determinantal random point field if and only if $`0K1`$. If the corresponding random point field exists it is unique.
$`\mathrm{}`$ The necessary and sufficient condition for the existence of the field has been already established. The uniqueness result easily follows from the general criterion (1.6’) since $`\frac{1}{k!}_{A^k}\rho _k(x_1,\mathrm{},x_k)𝑑x_1\mathrm{}𝑑x_k=\mathrm{Tr}(^k(K_A))\frac{\mathrm{Tr}(K_A)^k}{k!}\frac{1}{k!}`$
Consider arbitrary bounded Borel set $`BE`$. Then $`\mathrm{Tr}(K_B)=𝔼\mathrm{\#}_B<\mathrm{}`$ and the number of particles in $`B`$ is finite with probability 1. Let us write $`X=_{0k<\mathrm{}}C_k^B`$, where as before $`C_k^B=\{\xi X:\mathrm{\#}_B(\xi )=k\}`$. We choose a kernel for $`\chi _BK\chi _B`$ in such a way (see Lemma 1) that $`(\chi _BK\chi _B)(x,x+y)=_{i=1}^{\mathrm{}}\lambda _i(B)\phi _i(x)\overline{\phi _i(x+y)}`$ is a continuous function of $`y`$ in $`L^1(B)`$ norm. Assume for a moment that $`K_B<1`$. Then
$$L_B(x,x+y)=\underset{i=1}{\overset{\mathrm{}}{}}\frac{\lambda _i(B)}{1\lambda _i(B)}\phi _i(x)\overline{\phi _i(x+y)}$$
(1.35)
is also a continuous function of $`y`$ in $`L^1(B)`$ norm and is a kernel of $`L_B=(\mathrm{Id}K_B)^1K_B`$. Taking $`B_j`$ in (1.34) infinitesimally small one concludes that for each $`C_k^B`$ the distribution of $`k`$ particles $`x_1x_2\mathrm{}x_k`$ in $`B`$ has a density with respect to the Lebesgue measure. Denoting this density by $`p_k(x_1,\mathrm{},x_k)`$ we obtain
$$p_k(x_1,\mathrm{},x_k)=det(\mathrm{Id}K_B)det\left(L_B(x_i,x_j)\right)_{1i,jk}$$
(1.36)
(It should be noted that (1.36) may be nonnegative even for non-Hermitian kernel $`K`$, it is easy to see that such $`K`$ still has nonnegative minors). It follows from the definition of $`k`$-point correlation functions that
$$\rho _k(x_1,\mathrm{},x_k)=\underset{j=1}{\overset{\mathrm{}}{}}\frac{1}{j!}_{B^j}p_{k+j}(x_1,\mathrm{},x_k,x_{k+1},\mathrm{},x_{k+j})𝑑x_{k+1}\mathrm{}𝑑x_{k+j}$$
(1.37)
The system of equations can be inversed :
$$p_k(x_1,\mathrm{},x_k)=\underset{j=0}{\overset{\mathrm{}}{}}\frac{(1)^j}{j!}_{B^j}\rho _{k+j}(x_1,\mathrm{},x_k,x_{k+1},\mathrm{},x_{k+j})𝑑x_{k+1}\mathrm{}𝑑x_{k+j}$$
(1.38)
Functions $`p_k(x_1,\mathrm{},x_k)`$ are called Janossy probability densities (see \[DVJ\], p. 122) or exclusion probability densities (see \[Ma\]). It is easy to check that
$$\underset{j=0}{\overset{\mathrm{}}{}}\frac{1}{j!}_{B_j}p_j(x_1,\mathrm{},x_j)𝑑x_1\mathrm{}𝑑x_j=1$$
(1.39)
The r.h.s. of (1.36) still makes sense when $`K_B=\lambda _1(B)=1`$ (and therefore $`p_k(x_1,\mathrm{},x_k)`$ are properly defined in this case too). Indeed, $`det(\mathrm{Id}K_B)=_{j=1}^{\mathrm{}}(1\lambda _j(B))`$ as a function of $`\lambda _1`$ has a zero of order 1 at $`\lambda _1=1`$. We claim that $`det(L(x_i,x_j))_{1i,jk}`$ has a pole at $`\lambda _1=1`$ also of order 1. To see this we write $`L=\stackrel{~}{L}+\stackrel{}{L}`$, where $`\stackrel{~}{L}_{i,j}=\frac{\lambda _1(B)}{1\lambda _1(B)}\phi _1(x_i)\overline{\phi _1(x_j)}`$, $`\stackrel{}{L}=_\mathrm{}2\frac{\lambda _{\mathrm{}}(B)}{1\lambda _{\mathrm{}}(B)}\phi _{\mathrm{}}(x_i)\overline{\phi _{\mathrm{}}(x_j)}.`$ Then $`det(L(x_i,x_j))_{1i,jk}=^k(L(x_i,x_j)_{1i,jk})`$, and we use the fact that rank $`(\stackrel{~}{L})=1`$. If 1 is a multiple eigenvalue of $`K\mathrm{}_B`$, say $`\lambda _1(B)=\lambda _2(B)=\mathrm{}=\lambda _m(B)=1>\lambda _{m+1}(B)`$ one defines $`\stackrel{~}{L}_{i,j}=_{\mathrm{}=1}^m\frac{\lambda _{\mathrm{}}(B)}{1\lambda _{\mathrm{}}(B)}\phi _{\mathrm{}}(x_i)\overline{\phi _{\mathrm{}}(x_j)}`$ and proceeds in a similar manner.
Remark 5. Following Macchi, we call a random point field regular if for any Borel $`BE`$ satisfying $`\mathrm{\#}_B<\mathrm{}`$ (P-a.e.), the generating function $`𝔼z^{\mathrm{\#}_B}`$ is entire. It follows from our results (see also Theorem 4 below) that any determinantal random point field is regular.
Remark 6. In \[Ma\] (Theorem 12, p. 113) (see also \[DVJ\], p. 138) Macchi essentially claimed that a necessary and sufficient condition on the integral operator $`K`$, locally trace class, to define a regular fermion (=determinantal in our notations) random point field is $`0K<1`$. As one can see from Theorem 3 above this condition is sufficient, but not necessary (as we established in Theorem 3, the necessary and sufficient condition is $`0K1)`$. For completeness it should be noted that Macchi studied the case of continuous $`K(x,y)`$ with Tr$`K<\mathrm{}`$.
Remark 7. Formula (1.36) was established in \[Ma\], p. 113 (see also \[DVJ\], p. 138 and \[TW1\], p. 820).
We finish §1 with a few more results of general nature about determinantal random point fields.
Theorem 4
1. The probability of the event that the number of all particles is finite is either 0 or 1, depending on whether $`\mathrm{Tr}K`$ is finite or infinite.
2. The number of particles is less or equal to $`n`$ with probability 1 if and only if $`K`$ is a finite rank operator with rank $`(K)n`$.
3. The number of particles is $`n`$ with probability 1 if and only if $`K`$ is an orthogonal projector with rank $`(K)=n`$.
4. For any determinantal random point field with probability 1 no two particles coincide.
5. To obtain results of the theorem for $`BE`$ one has to replace $`K`$ by $`K_B`$ .
Proof of Theorem 4.
1. One direction is obvious. Indeed, if $`\mathrm{Tr}K=𝔼\mathrm{\#}_E<+\mathrm{}`$, then $`\mathrm{\#}_E<+\mathrm{}`$ with probability 1. Let us now assume $`\mathrm{Tr}K=+\mathrm{}`$. Consider a monotone absorbing family of compact sets $`\{B_j\}_{j=1}^{\mathrm{}}`$ (i.e., $`B_iB_{i+1}`$ and $`_{i=1}^{\mathrm{}}B_i=E)`$. Then $`\mathrm{Tr}K_{B_j}\underset{j\mathrm{}}{}+\mathrm{}`$. Fix arbitrary large $`N`$. By the construction of $`\{B_j\}`$ we have $`P(\mathrm{\#}_EN)=lim_j\mathrm{}P(\mathrm{\#}_{B_j}N)`$. But $`P(\mathrm{\#}_{B_j}N)2^N𝔼2^{\mathrm{\#}_{B_j}}=2^Ndet(\mathrm{Id}\frac{1}{2}K_{B_j})2^Ne^{\frac{1}{2}\mathrm{Tr}(K_{B_j})}\underset{j\mathrm{}}{}0`$.
2. If rank $`(K)=n`$, then writing $`K(x,y)=_{i=1}^n\lambda _i\phi _i(x)\overline{\phi _i(y)}`$ (a.e.), and $`\rho _n(x_1,\mathrm{},x_n)=det(K(x_i,x_j))_{1i,jn}`$ we observe that $`\rho _m(x_1,\mathrm{}x_m)=0`$ (a.e.) for any $`m>n`$. Therefore $`𝔼\mathrm{\#}_E(\mathrm{\#}_E1)\mathrm{}(\mathrm{\#}_En)=\rho _{n+1}(x_1,\mathrm{}x_{n+1})𝑑x_1\mathrm{}𝑑x_{n+1}=0`$ which implies $`\mathrm{\#}_En`$ with probability 1.
In the opposite direction, if $`\mathrm{\#}_En`$ (a.e.) we have $`_{B^{n+1}}\rho _{n+1}(x_1,\mathrm{},x_{n+1})𝑑x_1\mathrm{}𝑑x_{n+1}=0`$ for any bounded Borel $`BE`$, therefore $`\mathrm{Tr}(^{n+1}(K_B))=0`$. Since $`K0`$ we obtain rank $`(K_B)n`$ for arbitrary compact $`B`$, which implies rank $`(K)n`$.
3. follows from b) and the formula Var$`(\mathrm{\#}_E)=\mathrm{Tr}(KK^2)=_{i=1}^n\lambda _i(1\lambda _i)`$.
4. Let $`B_n=[n,n]^d`$. It is enough to show that for any $`n`$ with probability 1 no two particles in $`B_n`$ coincide. Let $`ϵ`$ be arbitrary small. Then $`P\{ij:x_i=x_jB_n\}P\{ij:|x_ix_j|<ϵ,x_iB_n,x_jB_n\}_{B_n}(_{|xy|<ϵ}\rho _2(x,y)𝑑x)𝑑y`$. Since $`\rho _2(x,y)`$ is locally integrable, the last integral can be made arbitrary small by letting $`ϵ0`$. $`\mathrm{}`$
The next result gives a criterion for the weak convergence of determinantal random point fields.
Theorem 5. Let $`P`$ and $`P_n,n=1,2,`$ be probability measures on $`(X,B)`$ corresponding to the determinantal random point fields defined by the Hermitian kernels $`K`$ and $`K_n`$. Let $`K_n`$ converge to $`K`$ in the weak operator topology and $`\mathrm{Tr}(\chi _BK_n\chi _B)\underset{n\mathrm{}}{}\mathrm{Tr}(\chi _BK\chi _B)`$ for any bounded Borel $`BE`$. Then the probability measures $`P_n`$ converge to $`P`$ weakly on the cylinder sets.
Proof of Theorem 5. It follows from \[S1\], Theorem 2.20, p. 40 that the assumptions of the theorem imply
$$\mathrm{Tr}|(K_nK)_B|=(K_nK)_B_1\underset{n\mathrm{}}{}0.$$
(1.40)
As a consequence of (1.40) we have
$$\mathrm{Tr}(K_n\chi _{B_1}\mathrm{}K_n\chi _{B_m})\underset{n\mathrm{}}{}\mathrm{Tr}(K\chi _{B_1}\mathrm{}K\chi _{B_m})$$
(1.41)
for any compact $`B_1,\mathrm{},B_m`$.
Thus using (1.26), (1.27) one can see that the joint moments of $`\{\mathrm{\#}_B\}`$ with respect to $`P_n`$ converge to the joint moments with respect to $`P`$. Since the moments of $`\mathrm{\#}_B`$ in the case of the determinantal random points define the distribution of $`\mathrm{\#}_B`$ uniquely one can see (exercise) that $`P_n\stackrel{𝑊}{}P`$. $`\mathrm{}`$
The rest of the notes is organized as follows. Section 2 is devoted to the various examples of determinantal random point fields arising in Quantum Mechanics, Statistical Mechanics, Random Matrix Theory, Representation Theory, Probability Theory (Renewal Process, 2D Random Growth Models). In §3 we discuss ergodic properties of the translation invariant determinantal random point fields. We also point out a special role played by the sine kernel $`K(x,y)=\frac{\mathrm{sin}\pi (xy)}{\pi (xy)}`$. In §4 we discuss the Central Limit Theorem for the counting measure and the Functional Central Limit Theorem for the empirical distribution function of spacings.
It is a great pleasure to thank Ya. Sinai for the encouragement to write this paper, B. Simon for the explaination of the result of Lemma 1 , G. Olshanski for many valuable remarks, and A. Borodin, B. Khoruzhenko, R. Killip, and Yu. Kondratiev for useful conversations.
## 2 Examples of Determinantal Random Point Fields
### 2.1 Fermion Gas
Let $`H=\frac{d^2}{dx^2}+V(x)`$ be a Schrödinger operator with discrete spectrum acting on $`L^2(E)`$. Let $`\{\phi _{\mathrm{}}\}_{\mathrm{}=0}^{\mathrm{}}`$ be an orthonormal basis of the eigenfunctions, $`H\phi _{\mathrm{}}=\lambda _{\mathrm{}}\phi _{\mathrm{}},\lambda _0<\lambda _1\lambda _2\mathrm{}`$. Consider the $`n^{th}`$ exterior power of $`H`$, $`^n(H):^n(L^2(E))^n(L^2(E))`$, where $`^n(L^2(E))=A_nL^2(E^n))`$ is the space of square-integrable antisymmetric functions of $`n`$ variables and $`^n(H)=_{i=1}^n(\frac{d^2}{dx_i^2}+V(x_i))`$. In Quantum Mechanics $`^n(H)`$ describes the Fermi gas with $`n`$ particles. The ground state of the Fermi gas is given by
$$\begin{array}{cc}& \psi (x_1,\mathrm{},x_n)=\hfill \\ & \frac{1}{\sqrt{n!}}\underset{\sigma S_n}{}(1)^\sigma \underset{i=1}{\overset{n}{}}\phi _{i1}(x_{\sigma (i)})=\frac{1}{\sqrt{n!}}det(\phi _{i1}(x_j))_{1i,jn}\hfill \end{array}$$
(2.1)
It could be noted that $`\psi (x_1,\mathrm{},x_n)`$ coincides up to a sign $`ϵ(x_1,\mathrm{},x_n)`$ with the ground state of $`_{i=1}^n(\frac{d^2}{dx_i^2}+V(x_i))`$ acting on $`S_nL^2(E^n)`$ with the boundary conditions $`\psi |_{x_i=x_j}=0`$. According to the postulate of Quantum Mechanics the absolute value squared of the ground state defines the probability distribution of $`n`$ particles. We write
$$\begin{array}{cc}& p(x_1,\mathrm{},x_n)=|\psi (x_1,\mathrm{},x_n)|^2=\frac{1}{n!}det\left(\phi _{i1}(x_j)\right)_{1i,jn}\hfill \\ & det\left(\overline{\phi _{j1}(x_i)}\right)_{1i,jn}=\frac{1}{n!}det\left(K_n(x_i,x_j)\right)_{1i,jn},\hfill \end{array}$$
(2.2)
where $`K_n(x,y)=_{i=0}^{n1}\phi _{i1}(x)\overline{\phi _{i1}(y)}`$ is the kernel of the orthogonal projector onto the subspace spanned by the first $`n`$ eigenfunctions of $`H`$. We claim that (2.2) defines a determinantal random point field. Indeed, the $`k`$-point correlation functions are given by
$$\begin{array}{cc}& \rho _k^{(n)}(x_1,\mathrm{},x_n)=\frac{n!}{(nk)!}p_n(x_1,\mathrm{},x_n)𝑑x_{k+1}\mathrm{}𝑑x_n=\hfill \\ & det\left(K_n(x_1,x_j)\right)_{1i,jk}\hfill \end{array}$$
(2.3)
The last equality in (2.3) follows from the general lemma well known in Random Matrix Theory.
Lemma 4 \[Me\], p. 89). Let $`(E,d\mu )`$ be a measurable space and a kernel $`K:E^2^1`$ satisfy
$$_EK(x,y)K(y,z)𝑑\mu (y)=K(x,z)$$
(2.4)
$$_EK(x,x)𝑑\mu (x)=\mathrm{const}$$
(2.5)
Then
$$\begin{array}{cc}& _Edet\left(K(x_i,x_j)\right)_{1i,jn}d\mu (x_n)=\hfill \\ & (\mathrm{const}n+1)det\left(K(x_i,x_j)\right)_{1i,jn1}.\hfill \end{array}$$
(2.6)
We shall consider in more detail two special cases of $`H`$. The first case is the harmonic oscillator
1. $`H=\frac{d^2}{dx^2}+x^2,E=^1`$. Then
$$\phi _{\mathrm{}}(x)=\frac{(1)^{\mathrm{}}}{\pi ^{\frac{1}{4}}(2^{\mathrm{}}\mathrm{}!)^{\frac{1}{2}}}\mathrm{exp}\left(\frac{x^2}{2}\right)\frac{d^{\mathrm{}}}{dx^{\mathrm{}}}\left(\mathrm{exp}(x^2)\right)$$
(2.7)
are known as Weber-Hermite functions. To pass to the thermodynamic limit $`n\mathrm{}`$ we make a proper rescaling
$$x_i=\frac{\pi }{(2n)^{\frac{1}{2}}}y_i,i=1,\mathrm{},n.$$
(2.8)
Then the Christoffel-Darboux formula and the Plancherel-Rotach asymptotics of the Hermite polynomials (\[E\]) imply that
$$\begin{array}{cc}& K_n(x_1,x_2)=\underset{\mathrm{}=0}{\overset{n1}{}}\phi _{\mathrm{}}(x_1)\phi _{\mathrm{}}(x_2)=\hfill \\ & \left(\frac{n}{2}\right)^{\frac{1}{2}}\left[\frac{\phi _n(x_1)\phi _{n1}(x_2)\phi _n(x_2)\phi _{n1}(x_1)}{x_1x_2}\right]\hfill \end{array}$$
has a limit as $`n+\mathrm{}`$
$$K_n(x_1,x_2)\underset{n\mathrm{}}{}K(y_1,y_2)=\frac{\mathrm{sin}\pi (y_1y_2)}{\pi (y_1y_2)}$$
(2.9)
The convergence of kernels implies the convergence of $`k`$-point correlation functions, which in turn implies the weak convergence of the distribution
$$\left(\frac{\pi }{(2n)^{\frac{1}{2}}}\right)^np_n(\frac{\pi }{(2n)^{\frac{1}{2}}}y_1,\mathrm{},\frac{\pi }{(2n)^{\frac{1}{2}}}y_n)dy_1\mathrm{}dy_n$$
to the translation-invariant determinantal random point field with the “sine kernel” $`K(y_1,y_2)=\frac{\mathrm{sin}\pi (y_1y_2)}{\pi (y_1y_2)}`$.
2. For another example let $`E=S^1=\{z=e^{i\theta },0\theta <2\pi \},H=\frac{d^2}{d\theta ^2}`$. Then
$$\begin{array}{cc}& \phi _{\mathrm{}}(\theta )=\frac{1}{\sqrt{2\pi }}e^{i\mathrm{}\theta },\hfill \\ & p_n(\theta _1,\mathrm{},\theta _n)=\frac{1}{n!}det\left(\underset{\mathrm{}=0}{\overset{n1}{}}\frac{1}{2\pi }e^{i\mathrm{}(\theta _j\theta _k)}\right)_{1j,kn}=\hfill \\ & \frac{1}{n!}det(K_n(\theta _i,\theta _j)_{1i,jn}\hfill \end{array}$$
(2.10)
where
$$K_n(\theta _1,\theta _2)=\frac{1}{2\pi }\frac{\mathrm{sin}\left(\frac{n}{2}(\theta _2\theta _1)\right)}{\mathrm{sin}\left(\frac{\theta _2\theta _1}{2}\right)}$$
(2.11)
After rescaling $`\frac{n}{2\pi }\theta _i=y_i,i=1,\mathrm{},n`$ the rescaled correlation functions have the same limit as in (2.9), in particular
$$\underset{n\mathrm{}}{lim}\frac{2\pi }{n}K_n(\frac{2\pi }{n}y_1,\frac{2\pi }{n}y_2)=\frac{\mathrm{sin}\pi (y_2y_1)}{\pi (y_2y_1)}.$$
For more information we refer the reader to \[D1\]-\[D3\], \[L4\]-\[L5\], \[Sp\].
### 2.2 Coulomb Gas at $`\beta =2`$
Examples a), b) from §2.1 can be reinterpreted as the equilibrium distribution of $`n`$ unit charges confined to the one-dimensional line (ex. 2.1a)) or the unit circle (ex. 2.1b)) repelling each other according to the Coulomb law of two-dimensional electrostatics. Writing the potential energy as $`H(z_1,\mathrm{},z_n)=_{1i<jn}\mathrm{log}|z_iz_j|+_{i=1}^nV(z_i)`$, where $`V`$ is an external potential, we note that the Boltzmann factor $`\frac{1}{Z}\mathrm{exp}(\beta H(z_1,\mathrm{},z_n)),\beta =2`$, is exactly $`p_n(z_1,\mathrm{}z_n)`$ in 2.1a) with $`V(z)=\frac{1}{2}z^2`$, and $`p_n(\theta _1,\mathrm{},\theta _n)`$ in §2.1b) with $`V(z)=0,z_j=e^{i\theta _j},j=1,\mathrm{},n`$.
The one-component Coulomb gas in two dimensions (a.k.a. a two-dimen- sional one-component plasma) was studied in a number of papers including \[Gin\], \[Ja1\], \[Ja2\], \[AL\], \[DFGIL\], \[FJ1\]. This subject is closely related to the theory of non-Hermitian Gaussian random matrices (to be discussed in §2.3d). The two-component two-dimensional Coulomb gas (i.e. a system of positively and negatively charged particles) was studied in \[Ga\], \[CJ1\]-\[CJ3\], \[AF\], \[FJ2\]. Let us start with a neutral system of $`n`$ positive and $`n`$ negative particles. After denoting the complex coordinates by $`u_j`$ and $`v_j,j=1,\mathrm{},n`$, we write the Boltzmann factor at $`\beta =2`$ as
$$\begin{array}{cc}& \mathrm{exp}\left(2\underset{1i<jn}{}(\mathrm{log}|u_iu_j|+\mathrm{log}|v_iv_j|2\mathrm{log}|u_iv_j|)\right)=\hfill \\ & \frac{_{1i<jn}|u_iu_j|^2|v_iv_j|^2}{_{i,j}|u_iv_j|^2}=\left|det\left(\frac{1}{u_iv_j}\right)_{1i,jn}\right|^2.\hfill \end{array}$$
Discretizing the model one allows the positive particles to occupy only the sites of the sublattice $`\gamma ^2`$ and the negative particles to occupy only the sites of the sublattice $`\gamma (^2+(\frac{1}{2},\frac{1}{2}))`$. The grand canonical ensemble is defined by the partition function (let $`\gamma =1`$)
$$\begin{array}{cc}& Z=1+\underset{u,v}{}\lambda _+(u)\lambda _{}(v)\frac{1}{|uv|^2}+\left(\frac{1}{2!}\right)^2\hfill \\ & \underset{u_1,u_2,v_1,v_2}{}\lambda _+(u_1)\lambda _+(u_2)\lambda _+(v_1)\lambda _+(v_2)\left|det\left(\frac{1}{u_iv_j}\right)_{1i,j2}\right|^2+\mathrm{},\hfill \end{array}$$
where $`\lambda _+(u)=e^{V(u)},\lambda _{}(u)=e^{V(u)}`$, are fugacities and $`V`$ is an external potential. One can rewrite the last formula as
$$\begin{array}{cc}\hfill Z=& det(\mathrm{Id}+(\lambda _+\frac{1+\sigma _z}{2}+\lambda _{}\frac{1\sigma _z}{2})\hfill \\ & (\frac{\sigma _x+i\sigma _y}{2}\frac{1}{zz^{}}+\frac{\sigma _xi\sigma _y}{2}\frac{1}{\overline{z}\overline{z^{}}})),\hfill \end{array}$$
where $`\sigma _x,\sigma _y,\sigma _z`$ are $`2\times 2`$ Pauli matrices. In particular we see that the grand canonical ensemble is a discrete fermion random point field (the appearance of matrix-valued kernel reflects the fact that $`E=^2(^2+(\frac{1}{2},\frac{1}{2})))`$. Passing to the continuous limit $`(\gamma =0)`$ one can see that two- and higher order correlation functions have a limit, and the limiting kernel $`K`$ can be expressed in terms of the Green function of a differential Dirac operator, namely
$$\begin{array}{cc}\hfill K=& (m_+\frac{1+\sigma _z}{2}+m_{}\frac{1\sigma _z}{2})\hfill \\ & \left(\sigma _x_x+\sigma _y_y+m_+\frac{1+\sigma _z}{2}+m_{}\frac{1\sigma _z}{2}\right)^1,\hfill \end{array}$$
where $`m_+,m_{}`$ are rescaled fugacities. In the special case $`m_+=m_{}`$ const (i.e., $`V0`$), $`K=\left(\begin{array}{cc}K++,& K+\\ K+,& K\end{array}\right)`$ can be expressed in terms of modified Bessel function (for the details see e.g., \[CJ3\]).
### 2.3 Random Matrix Models
a) Unitary Invariant Ensembles of Hermitian Randon Matrices
The probability distribution in §2.1a) (formulas (2.2), (2.7)) allows yet another interpretation. It is well known in Random Matrix Theory as the distribution of the eigenvalues in the Gaussian Unitary Ensemble (G.U.E.). We recall the definition of G.U.E. Consider the space of $`n\times n`$ Hermitian matrices $`\{A=(A_{ij})_{1i,jn},\mathrm{Re}(A_{ij})=\mathrm{Re}(A_{ji}),\mathrm{Im}(A_{ij})=\mathrm{Im}(A_{ji})\}`$. A G.U.E. random matrix is defined by its probability distribution
$$P(dA)=\mathrm{const}_n\mathrm{exp}(\mathrm{Tr}A^2)dA,$$
(2.12)
where $`dA`$ is a flat (Lebesgue) measure, i.e., $`dA=_{i<j}d\mathrm{Re}(A_{ij})d\mathrm{Im}(A_{ij})_{k=1}^ndA_{kk}`$. The definition of G.U.E. is equivalent to the requirement that $`\{\mathrm{Re}(A_{ij}),\mathrm{Im}(A_{ij}),1i<jn,A_{kk},1kn\}`$ are mutually independent and Re$`(A_{ij})N(0,\frac{1}{4}),\mathrm{Im}(A_{ij})N(0,\frac{1}{4}),A_{kk}N(0,\frac{1}{2})`$. The eigenvalues of a random Hermitian matrix are real random variables. For the derivation of their joint distribution we refer the reader to \[De\], sections 5.3-5.4 and \[Me\], chapters 3, 5. It appears that the density of the joint distribution with respect to the Lebesgue measure is given exactly by (2.2), (2.7).
We remark that the distribution of a G.U.E. random matrix is invariant under the unitary transformation $`AUAU^1,UU(n)`$. A natural generalization of (2.12) that preserves the unitary invariance is
$$P(dA)=\mathrm{const}_n\mathrm{exp}(2\mathrm{Tr}V(A))dA$$
(2.13)
where $`V(x)`$ can be, for example, a polynomial of even degree with a positive leading coefficients (see \[De\], section 5). The derivation of the formula for the joint distribution of the eigenvalues is very similar to the G.U.E. case. The density $`p_n(\lambda _1,\mathrm{},\lambda _n)`$ is given by (2.2), where $`\{\phi _{\mathrm{}}(x)e^{V(x)}\}_{\mathrm{}=0}^{n1}`$ are the first $`n`$ orthonormal polynomials with respect to the weight $`\mathrm{exp}(2V(x))`$. $`K_n(x,y)`$ is then again a kernel of a projector and therefore satisfies the conditions of Lemma 4.
b) Random Unitary Matrices
Let us consider the group of $`n\times n`$ unitary matrices $`U(n)`$. There exists a unique translation invariant probability measure on $`U(n)`$ (see \[We\]). It is called the Haar measure, we will denote it by $`\mu _{Haar}`$. The probability density of the induced distribution of the eigenvalues is given by
$$p_n(\theta _1,\mathrm{},\theta _n)=(2\pi )^n\frac{1}{n!}\underset{1k<\mathrm{}n}{}|e^{i\theta _k}e^{i\theta _{\mathrm{}}}|^2,$$
which coincides with (2.10)-(2.11) (see \[Me\], ch. 9-10, \[D1\]-\[D3\]). In the last formula we used the notations
$$\lambda _1=e^{i\theta _1},\mathrm{},\lambda _n=e^{i\theta _n}.$$
If one starts with the probability measure $`\mathrm{const}_ne^{\mathrm{Tr}V(U)}d\mu _{Haar}(U)`$ on the unitary group instead of the Haar measure, and replaces the monomials $`\frac{1}{\sqrt{2\pi }}e^{i\mathrm{}\theta }`$ by $`\psi _{\mathrm{}}(\theta )e^{{\scriptscriptstyle \frac{1}{2}}V(\theta )}`$, where $`\{\psi _{\mathrm{}}\}_{\mathrm{}=0}^{n1}`$ are the first $`n`$ orthonormal polynomials in $`e^{i\theta }`$ with respect to the weight $`e^{V(\theta )}d\theta `$, one stills arrives at the formula (2.10) for the $`k`$-point correlation functions.
c) Random Orthogonal and Symplectic Matrices
The distribution of the eigenvalues of a random orthogonal or symplectic matrix (with respect to the Haar measure) also has a form of a determinantial random point field with a fixed number of particles. For the convenience of the reader we draw below the chart of the kernels appearing in the ensembles of random matrices from the Classical Compact Groups.
| | $`K_n(x,y)`$ |
| --- | --- |
| $`U(n)`$ | $`\frac{1}{2\pi }\frac{\mathrm{sin}\left(\frac{n}{2}(xy)\right)}{\mathrm{sin}\left(\frac{xy}{2}\right)};E=[0,2\pi ]`$ |
| $`SO(2n)`$ | $`\frac{1}{2\pi }\left(\frac{\mathrm{sin}\left(\frac{2n1}{2}(xy)\right)}{\mathrm{sin}\left(\frac{xy}{2}\right)}+\frac{\mathrm{sin}\left(\frac{2n1}{2}(x+y)\right)}{\mathrm{sin}\left(\frac{x+y}{2}\right)}\right);E=[0,\pi ]`$ |
| $`SO(2n+1)`$ | $`\frac{1}{2\pi }\left(\frac{\mathrm{sin}(n(xy))}{\mathrm{sin}\left(\frac{xy}{2}\right)}\frac{\mathrm{sin}(n(x+y))}{\mathrm{sin}\left(\frac{x+y}{2}\right)}\right);E=[0,\pi ]`$ |
| $`Sp(n)`$ | $`\frac{1}{2\pi }\left(\frac{\mathrm{sin}(\frac{2n+1}{2}(xy))}{\mathrm{sin}\left(\frac{xy}{2}\right)}\frac{\mathrm{sin}(\frac{2n+1}{2}(x+y))}{\mathrm{sin}\left(\frac{x+y}{2}\right)}\right);E=[0,\pi ]`$ |
| $`0_{}(2n+2)`$ | the same as for $`Sp(n)`$ |
For additional information we refer the reader to \[Jo1\], \[DS\], \[KS\], \[So1\], \[So2\], \[So3\].
d) Complex Non-Hermitian Gaussian Random Matrices
In \[Gin\] Ginibre considered the ensemble of complex non-Hermitian random $`n\times n`$ matrices where all $`2n^2`$ parameters $`\{\mathrm{Re}A_{ij},\mathrm{Im}A_{ij},1i,jn\}`$ are independent Gaussian random variables with zero mean and variance $`\frac{1}{2}`$. The joint probability distribution of the matrix elements is then given by the formula
$$\begin{array}{cc}& P(dA)=\mathrm{const}_n\mathrm{exp}(\mathrm{Tr}(A^{}A))dA,\hfill \\ & dA=\underset{1j,kn}{}d\mathrm{Re}A_{jk}d\mathrm{Im}A_{jk}.\hfill \end{array}$$
(2.14)
The equivalent definition of (2.14) is that $`A=\stackrel{~}{A}+i\stackrel{}{A}`$, where $`\stackrel{~}{A}`$ and $`\stackrel{}{A}`$ are two independent G.U.E. matrices. The eigenvalues $`\lambda _1,\mathrm{},\lambda _n`$ are complex random variables. It was shown that their distribution is given by the determinantal random point field in $`^2`$ with a fixed number of particles $`(\mathrm{\#}=n)`$ and the correlation functions
$$\rho _k^{(n)}(z_1,\mathrm{},z_k)=det\left(K_n(z_j,\overline{z_m})\right)_{1j,mn}$$
(2.15)
where $`K_n(z_1,\overline{z_2})=\frac{1}{\pi }\mathrm{exp}(\frac{|z_1|^2}{2}\frac{|z_2|^2}{2})_{\mathrm{}=0}^{n1}\frac{z_1^{\mathrm{}}\overline{z_2}^{\mathrm{}}}{\mathrm{}!}`$. We mention in passing that $`K_n(z_1,\overline{z_2})`$ converges to the kernel
$$K(z_1,\overline{z_2})=\frac{1}{\pi }\mathrm{exp}(\frac{|z_1|^2}{2}\frac{|z_2|^2}{2}+z_1\overline{z_2})$$
(2.16)
which defines the limiting random point field. A generalization of (2.14) was studied in \[Gir1\], \[Gir2\], \[SCSS\], \[FKS1\], \[FKS2\]. Let $`A=\stackrel{~}{A}+iv\stackrel{}{A}`$, where $`\stackrel{~}{A}`$ and $`\stackrel{}{A}`$ are, as above, two independent G.U.E. matrices and $`v`$ is a real parameter (it is enough to consider $`0v1`$). Let us introduce a new parameter $`\tau =\frac{1v^2}{1+v^2}`$. The distribution of the matrix elements is given by
$$P(dA)=\mathrm{const}_n\mathrm{exp}\left(\frac{1}{1\tau ^2}\mathrm{Tr}(A^{}A\tau \mathrm{Re}(A^2))\right)dA$$
(2.17)
It induces the distribution of the eigenvalues
$$\begin{array}{cc}& p_n(z_1,\mathrm{},z_n)\underset{j=1}{\overset{n}{}}dz_jd\overline{z_j}=\mathrm{const}_n\mathrm{exp}[\frac{1}{1\tau ^2}\underset{j=1}{\overset{n}{}}\hfill \\ & (|z_j|^2\frac{\tau }{2}(z_j^2+\overline{z_j}^2))]_{j<k}|z_jz_k|^2^n_{j=1}dz_jd\overline{z_j}\hfill \end{array}$$
(2.18)
It should be noted that the expression (2.18) also appeared in the papers by DiFranceso et al. (\[DFGIL\]) and Forrester-Jancovici (\[FJ1\]) as the Boltzmann factor of the two-dimensional one-component plasma. For the calculation of the correlation functions we refer the reader to \[DFGIL\],\[FJ1\],\[FKS1\],\[FKS2\]. The crucial role there is played by the orthonormal polynomials in the complex plane with the weight
$$w^2(z)=\mathrm{exp}\left[\frac{1}{1\tau ^2}\left(|z|^2\frac{\tau }{2}(z^2+\overline{z}^2)\right)\right]$$
(2.19)
Such orthonormal polynomials can be expressed in terms of the Hermite polynomials,
$$\psi _{\mathrm{}}(z)=\frac{\tau ^\frac{\mathrm{}}{2}}{\pi ^{\frac{1}{2}}(\mathrm{}!)^{\frac{1}{2}}(1\tau ^2)^{\frac{1}{4}}}H_{\mathrm{}}\left(\frac{z}{\sqrt{\tau }}\right),\mathrm{}=0,1,\mathrm{}$$
(2.20)
where $`_{n=0}^{\mathrm{}}H_n(z)\frac{t^n}{n!}=\mathrm{exp}(zt\frac{t^2}{2})`$. We remark that if $`\tau =0`$ (Ginibre case) $`\psi _{\mathrm{}}(z)=\frac{1}{\pi ^{\frac{1}{2}}(\mathrm{}!)^{\frac{1}{2}}}z^{\mathrm{}}`$. Formula for the correlation functions ($`\tau `$ arbitrary) generalzies (2.15):
$$\begin{array}{c}\hfill \rho _k^{(n)}=det\left(K_n(z_i,\overline{z_j})\right)_{1i,jk},\\ \hfill K_n(z_1,\overline{z_2})=w(z_1)w(\overline{z_2})\underset{\mathrm{}=0}{\overset{n1}{}}\psi _{\mathrm{}}(z_1)\psi _{\mathrm{}}(\overline{z_2})\end{array}$$
(2.21)
In the limit $`n\mathrm{}K_n(z,\overline{z_2})`$ converges to
$$\begin{array}{cc}& K(z_1,\overline{z_2})=\underset{n\mathrm{}}{lim}K_n(z_1,\overline{z_2})=\hfill \\ & \frac{1}{\pi (1\tau ^2)}\mathrm{exp}\left(\frac{1}{1\tau ^2}\left(\frac{|z_1|^2}{2}+\frac{|z_2|^2}{2}z_1\overline{z_2}\right)\right)\hfill \end{array}$$
(2.22)
We reamrk that the last formula differs from (2.16) only by the trivial rescaling $`zz\sqrt{1\tau ^2}`$. A special regime, called the regime of weak non-Hermiticity, was discovered for the model (2.17) by Fyodorov, Khoruzhenko and Sommers in \[FKS1\], \[FKS2\]. Let
$`\mathrm{Re}(z_1)`$ $`=`$ $`n^{\frac{1}{2}}x+n^{\frac{1}{2}}x_1,`$
$`\mathrm{Re}(z_2)`$ $`=`$ $`n^{\frac{1}{2}}x+n^{\frac{1}{2}}x_2,`$
$`\mathrm{Im}(z_1)`$ $`=`$ $`n^{\frac{1}{2}}y_1,`$
$`\mathrm{Im}(z_2)`$ $`=`$ $`n^{\frac{1}{2}}y_2,`$
Assume the parameters $`x,x_1,x_2,y_1,y_2`$ fixed and take the limit $`n\mathrm{}`$ in such a way that $`lim_n\mathrm{}n(1\tau )=\frac{\alpha ^2}{2}`$. Then
$$\begin{array}{cc}& \underset{n\mathrm{}}{lim}\frac{1}{n}K_n(z_1,z_2)=\frac{1}{\pi \alpha }\hfill \\ & \mathrm{exp}\left[\frac{y_1^2+y_2^2}{\alpha ^2}+ix\frac{(y_1y_2)}{2}\right]g_\alpha \left(\frac{y_1+y_2}{2}i\frac{(x_1x_2)}{2}\right),\hfill \end{array}$$
(2.23)
where
$$g_\alpha (y)=_{\sqrt{1\frac{x^2}{4}}}^{\sqrt{1\frac{x^2}{4}}}\frac{du}{\sqrt{2\pi }}\mathrm{exp}\left[\frac{\alpha ^2u^2}{2}2uy\right].$$
(2.24)
(if $`x>2`$ the limit in (2.23) is equal to zero). The formulas (2.23), (2.24) define determinantal random point field in $`^2`$, different from (2.16).
e) Positive Hermitian Random Matrices
Following Bronk \[Br\] we define the Laguerre ensemble of positive Hermitian $`n\times n`$ matrices. Any positive Hermitian matrix $`M`$ can be written as $`M=A^{}A`$, where $`A`$ is some complex matrix. The probability distribution of a random matrix $`M`$ is given by
$$\mathrm{const}_n\mathrm{exp}(\mathrm{Tr}A^{}A)[det(A^{}A)]^\alpha dA,$$
(2.25)
where $`dA`$ is defined as in (2.14) and $`\alpha >1`$ (the values of $`\alpha `$ of special interest are $`\pm \frac{1}{2},0`$). The induced probability distribution of the (positive) eigenvalues is given by
$$\mathrm{const}_n\mathrm{exp}\left(\underset{i=1}{\overset{n}{}}\lambda _i\right)\underset{i=1}{\overset{n}{}}\lambda _i^\alpha \underset{1i<jn}{}(\lambda _i\lambda _j)^2d\lambda _1\mathrm{}\lambda _n$$
(2.26)
Employing the associated Laguerre polynominals
$$L_m^\alpha (x)\frac{1}{n!}e^xx^\alpha \frac{d^m}{dx^m}(e^xx^{m+\alpha }),m=0,1,\mathrm{}$$
one can rewrite (2.26) as
$$\frac{1}{n!}det\left(K_n(x_i,x_j)\right)_{1i,jn}$$
(2.27)
where
$$K_n(x,y)=\underset{\mathrm{}=0}{\overset{n1}{}}\phi _{\mathrm{}}^{(\alpha )}(x)\phi _{\mathrm{}}^{(\alpha )}(y),$$
(2.28)
and $`\{\phi _{\mathrm{}}^{(\alpha )}(x)=(\mathrm{\Gamma }(\alpha +1)\left(\genfrac{}{}{0pt}{}{n+\alpha }{n}\right))^{{\scriptscriptstyle \frac{1}{2}}}L_{\mathrm{}}^\alpha (x)\}_{\mathrm{}=0}^{\mathrm{}}`$ is the orthonormal basis with respect to the weight $`e^xx^\alpha `$ on the positive semiaxis. Once again Lemma 4 allows us to calculate explicitly $`k`$-point correlation functions and show that they are given by determinants of $`k\times k`$ matrices with the kernel (2.28).
f) Hermitian Matrices Coupled in a Chain
Let $`A_1,\mathrm{},A_p`$ be complex hermitian random $`n\times n`$ matrices with the joint probability density
$$\begin{array}{cc}& \mathrm{const}_n\mathrm{exp}[\mathrm{Tr}(\frac{1}{2}V_1(A_1)+V_2(A_2)+\mathrm{}+V_{p1}(A_{p1})+\frac{1}{2}V_p(A_p)\hfill \\ & +c_1A_1A_2+c_2A_2A_3+\mathrm{}+c_{p1}A_{p1}A_p)]\hfill \end{array}$$
(2.29)
We denote the eigenvalues of $`A_j`$ (all real) by $`\stackrel{~}{\lambda _j}=(\lambda _{j1},\mathrm{},\lambda _{jn}),j=1,\mathrm{},p`$. The induced probability densit of the eigenvalues is then equal to
$$\begin{array}{cc}& P_n(\stackrel{~}{\lambda _1},\mathrm{},\stackrel{~}{\lambda _p})=\mathrm{const}_n\left[\underset{1r<sn}{}(\lambda _{1r}\lambda _{1s})(\lambda _{pr}\lambda _{ps})\right]\hfill \\ & \left[\underset{k=1}{\overset{p1}{}}det\left[w_k(\lambda _{kr},\lambda _{k+1s})\right]_{r,s=1,\mathrm{},n}\right]\hfill \end{array}$$
(2.30)
where
$$w_k(x,y)=\mathrm{exp}\left(\frac{1}{2}V_k(x)\frac{1}{2}V_{k+1}(y)+c_kxy\right)$$
(2.31)
Eynard and Mehta \[EM\] established that the correlation functions of this model $`\rho _{k_1,\mathrm{},k_p}(\lambda _{11},\mathrm{},\lambda _{1k_1};\mathrm{};\lambda _{p1},\mathrm{},\lambda _{pk_p})=_{j=1}^p\frac{n!}{(nk_j)!}p_n(\stackrel{~}{\lambda _1},\mathrm{},\stackrel{~}{\lambda _p})_{j=1}^p_{r_j=k_j+1}^nd\lambda _{jr_j}`$ can be written as a $`k\times k`$ determinant with $`k=k_1+\mathrm{}+k_p`$
$$det\left[K_{ij}(\lambda _{ir},\lambda _{js})\right]_{r=1,\mathrm{},k_i;s=1,\mathrm{}k_j;i,j=1,\mathrm{},p}$$
(2.32)
For the exact formulas for the kernels $`K_{ij}(x,y)`$ we refer the reader to \[EM\] (see also \[AM\]). We remark that (2.32) defines a determinantal random point field with one-particle space $`E`$ being the union of $`p`$ copies of $`^1`$.
g) Universality in Random Matrix Models. Airy, Bessel and sine Random Point Fields.
We start with a general class of kernels of the form
$$K(x,y)=\frac{\phi (x)\psi (y)\phi (y)\psi (x)}{xy}$$
(2.33)
where
$$\begin{array}{cc}& m(x)\phi ^{}(x)=A(x)\phi (x)+B(x)\psi (x)\hfill \\ & m(x)\psi ^{}(x)=C(x)\phi (x)A(x)\psi (x)\hfill \end{array}$$
(2.34)
and $`m(x),A(x),B(x),C(x)`$ are polynomials. It was shown by Tracy and Widom (\[TW2\]) that Fredholm determinants of integral operators with kernels (2.33)-(2.34) restricted to a finite union of intervals satisfy certain partial differential equations. Airy, Bessel and sine kernels are the special cases of (2.33), (2.34). To define sine kernel we set $`\phi (x)\frac{1}{\pi }\mathrm{sin}(\pi x),\psi (x)\phi ^{}(x)(m(x)1,A(x)0,B(x)1,C(x)\pi ^2)`$. For the Airy kernel $`\phi (x)A_i(x),\psi (x)\phi ^{}(x)(m(x)1,A(x)0,B(x)1,C(x)x)`$. For the Bessel kernel $`\phi (x)J_\alpha (\sqrt{x}),\psi (x)x\phi ^{}(x)(m(x)x,A(x)0,B(x)1,C(x)\frac{1}{4}(x\alpha ^2))`$. Here $`A_i(x)`$ is the Airy function and $`J_\alpha (x)`$ is the Bessel function of order $`\alpha `$ (see \[E\]). Writing down these kernels explicitly we have (see \[TW2\], \[TW3\], \[TW4\])
$$K_{\mathrm{sine}}(x,y)=\frac{\mathrm{sin}\pi (xy)}{\pi (xy)},$$
(2.35)
$$\begin{array}{cc}\hfill K_{\mathrm{Airy}}(x,y)& =\frac{𝒜_i(x)𝒜_i^{}(y)𝒜_i(y)𝒜_i^{}(x)}{xy}\hfill \\ & =_0^{\mathrm{}}A_i(x+t)𝒜_i(y+t)𝑑t\hfill \end{array}$$
(2.36)
$$\begin{array}{cc}\hfill K_{\mathrm{Bessel}}(x,y)& =\frac{J_\alpha (\sqrt{x})\sqrt{y}J_\alpha ^{}(\sqrt{y})\sqrt{x}J_\alpha ^{}(\sqrt{x})J_\alpha (\sqrt{y})}{2(xy)}\hfill \\ & =\frac{\sqrt{x}J_{\alpha +1}(\sqrt{x})J_\alpha (\sqrt{y})J_\alpha (\sqrt{x})\sqrt{y}J_{\alpha +1}(\sqrt{y})}{2(xy)}\hfill \end{array}$$
(2.37)
As we already mentioned, sine kernel appears as a scaling limit in the bulk of the spectrum in G.U.E. (\[Me\], Chapter 5). In its turn, Airy kernel appears as a scaling limit at the edge of the spectrum in G.U.E. and at the (soft) right edge of the spectrum in the Laguerre ensemble, while Bessel kernel appears as a scaling limit at the (hard) left edge in the Laguerre ensemble (\[F\], \[TW3\], \[TW4\]). Universality conjecture in Random Matrix Theory asserts that such limits should be universal for a wide class of Hermitian random matrices. Recently this conjecture was proven for unitary invariant ensembles (2.13) in the bulk of the spectrum (\[PS\], \[BI\], \[DKMLVZ\], \[De\]) and for some classes of Wigner matrices in the bulk of the spectrum \[Jo4\] and at the edge \[So4\].
In the next subsection we completely characterize determinantal random point fields in $`^1(^1)`$ with independent identically distributed spacings.
### 2.4 Determinantal random point fields with i.i.d. spacings. Renewal processes
We start with some basic facts from the theory of renewal processes (see e.g., \[Fe\], \[DVJ\]). Let $`\{\tau _k\}_{k=1}^{\mathrm{}}`$ be independent identically distributed non-negative random variables and $`\tau _0`$ another non-negative random variable independent from $`\{\tau _k\}_{k=1}^{\mathrm{}}`$ (in general the distribution of $`\tau _0`$ will be different). We define
$$x_k=\underset{j=0}{\overset{k}{}}\tau _j.$$
(2.38)
This gives us a random configuration $`\{x_k\}_{k=0}^{\mathrm{}}`$ in $`_+^1`$. In probability theory a random sequence $`\{x_k\}_{k=0}^{\mathrm{}}`$ is known as delayed renewal process. We assume that the distribution of random variables $`\tau _k,k1`$ has density $`f(x)`$, called interval distribution density, and a finite mathematical expectation $`E\tau _1=_0^{\mathrm{}}xf(x)𝑑x`$. The renewal density is defined than as
$$\begin{array}{cc}\hfill u(x)=& \underset{k=1}{\overset{\mathrm{}}{}}f^k(x)=f(x)+_0^xf(xy)f(y)𝑑y+\hfill \\ & _0^x_0^{xy_2}f(xy_1y_2)f(y_1)f(y_2)𝑑y_1𝑑y_2+\mathrm{}\hfill \end{array}$$
(2.39)
One can express higher order correlation functions of the renewal process through its one-point correlation function and the renewal density. Indeed (see \[DVJ\], p. 136) for $`t_1t_2\mathrm{}t_k`$ and $`k>1`$ the following formula takes place
$$\rho _k(t_1,\mathrm{},t_k)=\rho _1(t_1)u(t_2t_1)u(t_3t_2)\mathrm{}u(t_kt_{k1})$$
(2.40)
It follows immediately from the above definitions that a random point field in $`_+^1`$ has i.i.d. nearest spacings iff it is a renewal process (2.38). To make this process translation-invariant the probability density of $`\tau _0`$ must be given by
$$\frac{1}{E\tau _1}_x^+\mathrm{}f(t)𝑑t\text{([DVJ], p. 72, [Fe], section XI.3)}$$
(2.41)
Then one-point correlation function is identically constant, $`\rho _1(x)\rho >0`$, so (2.40) implies that the distribution of the process is uniquely defined by the renewal density (in particular one can obtain $`\rho `$ from $`u(x)`$ since $`\rho =(E\tau _1)^1`$ and the Laplace transforms of $`f`$ and $`u`$ are simply related). Macchi (\[Ma\]) considered a special class of translation-invariant renewal processes with the interval distribution density
$$f(x)=2\rho (12\rho \alpha )^{\frac{1}{2}}e^{\frac{x}{\alpha }}\mathrm{sinh}\left((12\rho \alpha )^{\frac{1}{2}}\left(\frac{x}{\alpha }\right)\right),$$
(2.42)
where
$$2\rho \alpha 1,\rho >0,\alpha >0,$$
(2.43)
and showed that it is a determinantal random point field with the kernel
$$K(x,y)=\rho \mathrm{exp}(|xy|/\alpha )$$
(2.44)
(restrictions (2.43) are exactly $`0<K\mathrm{Id}`$).
In the next theorem we classify all delayed renewal processes that are also determinantal random point fields in $`_+^1`$.
Theorem 6. Determinantal random point field in $`_+^1`$ with Hermitian kernel has i.i.d. spacings if and only if its kernel satisfies the following two conditions in addition to $`0K\mathrm{Id}`$ locally trace class :
a) for almost all $`x_1x_2x_3`$
$$K(x_1,x_2)K(x_2,x_3)=K(x_1,x_3)K(x_2,x_2),$$
(2.45)
b) for almost all $`x_1x_2`$ the function
$$K(x_2,x_2)\frac{K(x_1,x_2)K(x_2,x_1)}{K(x_2,x_1)}$$
(2.46)
depends only on the difference $`x_2x_1`$. If a determinantal random point field is both translation-invariant and with i.i.d. spacings, it is given by (2.42)-(2.44).
Remark 8. Of course a translation-invariant d.r.p.f. in $`_+^1`$ can be extended in a unique way to the translation-invariant d.r.p.f. in $`^1`$.
Proof of Theorem 6. First we prove the “only if” part of the theorem. Suppose that a determinantal random point field with a kernel $`K(x,y)`$ is also a delayed renewal process. From (2.40), $`k=2,3`$, we obtain the formula for the renewal density
$$u(yx)=K(y,y)\frac{K(x,y)K(y,x)}{K(x,x)},yx,$$
(2.47)
and the expression for $`\rho _3(x_1,x_2,x_3),x_1x_2x_3`$:
$$\begin{array}{cc}& \rho _3(x_1,x_2,x_3)=K(x_1,x_1)u(x_2x_1)u(x_3x_2)\hfill \\ & =K(x_1,x_1)\left(K(x_2,x_2)\frac{K(x_1,x_2)K(x_2,x_1)}{K(x_1,x_1)}\right)\hfill \\ & \left(K(x_3,x_3)\frac{K(x_2,x_3)K(x_3,x_2)}{K(x_2,x_2)}\right)\hfill \end{array}$$
(2.48)
Since with probability 1 there are no particles outside $`A=\{x:K(x,x)>0\}`$, we can always consider random point field restricted to $`A`$.
Comparing
$$\rho _3(x_1,x_2,x_3)=det(K(x_i,x_j))_{1i,j3}$$
(2.49)
with (2.48), we have
$$\begin{array}{cc}& K(x_1,x_2)K(x_2,x_1)K(x_2,x_3)K(x_3,x_2)\frac{1}{K(x_2,x_2)}=\hfill \\ & K(x_1,x_3)K(x_3,x_1)K(x_2,x_2)+K(x_1,x_2)K(x_2,x_3)\hfill \\ & K(x_3,x_1)+K(x_1,x_3)K(x_3,x_2)K(x_2,x_1)\hfill \end{array}$$
which is equivalent to
$$\begin{array}{cc}& \frac{1}{K(x_2,x_2)}\left(K(x_1,x_2)K(x_2,x_3)K(x_2,x_2)K(x_1,x_3)\right)\hfill \\ & (K(x_3,x_2)K(x_2,x_1)K(x_3,x_1)K(x_2,x_2))=0\hfill \end{array}$$
The third factor in the last equality is a complete conjugate of the second factor, and we obtain (2.45). Condition b) of the theorem has been already established in (2.47). For the translation-invariant d.r.p.f. the kernel $`K(x,y)`$ depends only on the difference, therefore $`K(x,y)=\rho e^{|xy|/\alpha }e^{i\beta (xy)}`$, and the unitary equivalent kernel $`e^{i\beta x}K(x,y)e^{i\beta y}`$ coincides with (2.44). Now we turn to the proof of the “if” part of the theorem. Once we are given the kernel satisfying (2.45) and (2.46) the candidate for the renewal density must obey $`u(x_2x_1)=K(x_2,x_2)\frac{K(x_1,x_2)K(x_2,x_1)}{K(x_1,x_1)}`$ for almost all $`x_1x_2`$. Let $`x_1x_2\mathrm{}x_k`$. Our goal is to deduce the algebraic identity
$$\begin{array}{cc}& det\left(K(x_i,x_i)\right)_{1i,jk}=K(x_1,x_1)\underset{i=1}{\overset{k1}{}}(K(x_{i+1},x_{i+1})\hfill \\ & \frac{K(x_i,x_{i+1})K(x_{i+1},x_i)}{K(x_i,x_i)})\hfill \end{array}$$
(2.50)
from the basic identities between the commuting variables $`K(x_i,x_j),\overline{K(x_i,x_j)}`$ satisfying $`K(x_i,x_j)K(x_j,x_{\mathrm{}})=K(x_i,x_{\mathrm{}})K(x_j,x_j),1ij\mathrm{}k,K(x_i,x_j)=\overline{K(x_j,x_i)}`$. Let us introduce $`a(x)=K(x,x)K(0,x)^1,b(x)=K(0,x)^1`$. Then for $`ij`$
$$K(x_i,x_j)=a(x_i)b(x_j)^1,$$
$$K(x_j,x_i)=\overline{a(x_i)}\overline{b(x_j)}^1.$$
This allows us to write the determinant as
$$\begin{array}{cc}& \left|\begin{array}{cccc}a(x_1)b(x_1)^1,& a(x_1)b(x_2)^1,& \mathrm{}& ,a(x_1)b(x_n)^1\\ \overline{a(x_1)}\overline{b(x_2)}^1,& a(x_2)b(x_2)^1,& \mathrm{}& ,a(x_2)b(x_n)^1\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \overline{a(x_1)}\overline{b(x_n)}^1,& \overline{a(x_2)}\overline{b(x_n)}^1,& \mathrm{}& ,a(x_n)b(x_n)^1\end{array}\right|\hfill \\ & =a(x_1)b(x_1)^1\underset{i=1}{\overset{n1}{}}(b(x_{i+1})^1\overline{b(x_{i+1})}^1\hfill \\ & (a(x_{i+1})\overline{b(x_{i+1})}\overline{a(x_i)}b(x_i))),\hfill \end{array}$$
(2.51)
which is exactly the r.h.s. of (2.50). Once we established $`\rho _k(x_1,\mathrm{},x_k)=\rho _1(x_1)_{i=1}^{n1}u(x_{i+1}x_i),x_1x_2\mathrm{}x_k`$, the rest of the proof is quite easy. Let $`p_k(x_1,\mathrm{},x_k)`$ be Janossy densities, i.e. the probability density of the event to have particles at $`x_1,\mathrm{},x_k`$, and no particles in between. We recall that $`p_k(x_1,\mathrm{},x_k)=_{j=1}^{\mathrm{}}\frac{(1)^j}{j!}\rho _{k+j}(x_1,\mathrm{},x_k;y_{k+1},\mathrm{},y_{k+j})𝑑y_{k+1}\mathrm{}𝑑y_{k+j}`$, where the integration in the $`j`$th term is over $`\underset{j\text{ times}}{(x_1,x_k)\times \mathrm{}\times (x_1,x_k)}`$. We claim that
$$p_k(x_1,\mathrm{},x_k)=\rho _1(x_1)\underset{i=1}{\overset{k1}{}}f(x_{i+1}x_i)$$
(2.52)
where the interval distribution density $`f`$ and the renewal distribution density $`u`$ are related via the convolution equation
$$u=f+uf$$
(2.53)
Theorem 6 is proven. $`\mathrm{}`$
Remark 9. The analogue of Theorem 6 is valid in the discrete case and the proof is the same. One has to replace (2.42) by the solution of the discrete convolution equation (2.53) with $`u(n)=1\rho e^{2\beta n},K(n_1,n_2)=\rho e^{\beta |n_1n_2|},0<\rho 1,\beta >0`$, so that
$$\widehat{f}(t)=\underset{n=0}{\overset{\mathrm{}}{}}f(n)e^{int}=\frac{(1\rho )(e^{2\beta }\rho )e^{it}}{(2\rho )(2e^{2\beta }\rho +1)e^{it}+e^{2\beta }e^{2it}}$$
(2.54)
Remark 10. One can consider a generalization of Theorem 6 to the case when the multiplicative identity (2.45) still holds, but the renewal density
$$u(x_1,x_2)=K(x_2,x_2)\frac{K(x_1,x_2)K(x_2,x_1)}{K(x_1,x_1)}$$
no longer depends only on the difference of $`x_1`$ and $`x_2`$. Such processes have independent but not necessarily identically distributed spacings since the spacings distribution $`f(x,y)dy`$ depends on the position $`x`$ of the left particle. Thus
$$\begin{array}{cc}& u(x_1,x_2)=f(x_1,x_2)+_{x_1}^{x_2}f(x_1,y_1)f(y_1,x_2)𝑑y_1\hfill \\ & +_{x_1}^{x_2}_{x_1}^{y_2}f(x_1,y_1)f(y_1,y_2)f(y_2,x_2)𝑑y_1𝑑y_2+\mathrm{}\hfill \end{array}$$
(2.55)
where $`f(x,y)`$ is a one-parameter family of probability densities, such that
$$\text{supp }f(x,)[x,+\mathrm{}],f0,f(x,y)𝑑y=1.$$
We recall the inversion formula for the equation (2.55):
$$\begin{array}{cc}& f(x_1,x_2)=u(x_1,x_2)_{x_1}^{x_2}u(x_1,y_1)u(y_1,x_2)𝑑y_1\hfill \\ & +_{x_1}^{x_2}_{x_1}^{y_2}u(x_1,y_1)u(y_1,y_2)u(y_2,x_2)𝑑y_1𝑑y_2\mathrm{}\hfill \end{array}$$
(2.56)
Writing $`K(x,y)=a(x)b(y)^1,xy`$, where $`a(x)=\frac{K(x,x)}{K(0,x)},b(y)=\frac{1}{K(0,y)}`$, and $`u(x,y)=\frac{1}{|b(y)|^2}(a(y)\overline{b(y)}\overline{a(x)}b(x))`$, one can in principle characterize through (2.56) the class of corresponding interval densities $`u(x,y)`$.
### 2.5 Plancherel Measure on Partitions and its Generali- <br>zations–$`z`$-Measures and Schur Measures
By a partition of $`n=1,2,\mathrm{}`$ we understand a collection of non-negative integers $`\lambda =(\lambda _1,\mathrm{},\lambda _m)`$ such that $`\lambda _1+\mathrm{}+\lambda _m=n`$ and $`\lambda _1\lambda _2\mathrm{}\lambda _m`$. We shall use a notation $`\mathrm{Par}(n)`$ for the set of all partitions of $`n`$. For the basic facts about partitions we refer the reader to \[St\], \[Fu\], \[Mac\], \[Sa\]. In particular recall that each partition $`\lambda `$ of $`n`$ (denoted $`\lambda n`$) can be identified with a Young diagram with $`|\lambda |=n`$ boxes. A partition $`\lambda ^{}`$ corresponds to the transposed diagram. Let $`d`$ be the number of diagonal boxes in $`\lambda `$ (i.e., the number of diagonal boxes in the Young diagram corresponding to $`\lambda `$). We define the Frobenius coordinates of $`\lambda `$ as $`(p_1,\mathrm{},p_d|q_1,\mathrm{},q_k)`$, where $`p_j=\lambda _jj,q_j=\lambda _j^{}j,j=1,\mathrm{},d`$. The importance of partitions in Representation Theory can be most easily understood from the fact that the elements of $`\mathrm{Par}(n)`$ label the irreducible representations of the symmetric group $`S_n`$ (see e.g., \[Sa\], \[Fu\]). The Plancherel measure $`M_n`$ on the set $`\mathrm{Par}(n)`$ of all partitions of $`n`$ is given by
$$M_n(\lambda )=\frac{(dim\lambda )^2}{n!},$$
(2.57)
where $`dim\lambda `$ is the dimension of the corresponding representation of $`S_n`$. The dimension $`dim\lambda `$ can be expressed in terms of the Frobenius coordinates via a determinantal formula
$$\frac{dim\lambda }{n!}=det\left[\frac{1}{(p_i+q_j+1)p_i!q_i!}\right]_{1i,jd}$$
(2.58)
where $`|\lambda |=n`$ (\[Ol\], Proposition 2.6, formula (2.7)). Let $`\mathrm{Par}=_{n=0}^{\mathrm{}}\mathrm{Par}(n)`$. Consider the measure $`M^\theta `$ on Par, which in analogy with statistical mechanics can be called the grand canonical ensemble:
$$\begin{array}{cc}& M^\theta (\lambda )=e^\theta \frac{\theta ^n}{n!}M_n(\lambda ),\hfill \\ & \text{ if }\lambda \mathrm{Par}(n),n=0,1,2,\mathrm{},0\theta <\mathrm{}.\hfill \end{array}$$
(2.59)
$`M^\theta `$ is also called the poissonization of the measures $`M_n`$. It follows from (2.59) that $`|\lambda |`$ is distributed by the Poisson law with the mean $`\theta `$, and $`M^\theta (||\lambda |=n)=M_n`$. In the Frobenius coordinates measures $`M^\theta ,M_n`$ can be viewed as random point fields on the lattice $`^1`$. Recently Borodin, Okounkov and Olshanski (\[BOO\]) and, independently, Johansson \[Jo3\]) proved that $`M^\theta `$ is a determinantal random point field (to be exact in \[Jo3\] only the restriction of $`M^\theta `$ to the first half of the Frobenius coordinates $`(p_1,\mathrm{},p_{d(\lambda )})`$ was studied, and as a result, only the part of (2.60) corresponding to $`xy>0`$ was obtained). To formulate the results of \[BOO\], \[Jo3\] we define the modified Frobenius coordinates of $`\lambda `$ by
$$\mathrm{Fr}(\lambda ):=\{p_1+\frac{1}{2},\mathrm{},p_d+\frac{1}{2},q_1\frac{1}{2},\mathrm{},q_d\frac{1}{2}\}.$$
Let $`\rho _k^\theta (x_1,\mathrm{},x_k)`$ be the $`k`$-point correlation function of $`M^\theta `$ in the modified Frobenius coordinates, where
$$\{x_1,\mathrm{},x_k\}^1+\frac{1}{2}.$$
Then
$$\rho _k^\theta (x_1,\mathrm{},x_k)=det[K(x_i,x_j)]_{1i,jk},$$
where $`K`$ is a so-called discrete Bessel kernel,
$$K(x,y)=\{\begin{array}{cc}\sqrt{\theta }\frac{J_{|x|\frac{1}{2}}(2\sqrt{\theta })J_{|y|+\frac{1}{2}}(2\sqrt{\theta })J_{|x|+\frac{1}{2}}(2\sqrt{\theta })J_{|y|\frac{1}{2}}(2\sqrt{\theta })}{|x||y|},\text{ if }xy>0,\hfill & \\ \sqrt{\theta }\frac{J_{|x|\frac{1}{2}}(2\sqrt{\theta })J_{|y|\frac{1}{2}}(2\sqrt{\theta })J_{|x|+\frac{1}{2}}(2\sqrt{\theta })J_{|y|+\frac{1}{2}}(2\sqrt{\theta })}{xy},\text{ if }xy<0,\hfill & \end{array}$$
(2.60)
where $`J_x()`$ is the Bessel function of order $`x`$. We note that the kernel $`K(x,y)`$ is not Hermitian symmetric, however the restriction of this kernel to the positive and negative semi-axis is Hermitian. (2.60) can be seen as a limiting case of a more general theorem obtained by Borodin and Olshanski for the so-called $`z`$-measures (see Theorem 3.3 of \[BO1\], also \[BO2\], \[BO3\], \[KOV\] and references therein). Let $`z,z^{}`$ be complex numbers such that either
$$\begin{array}{cc}& z^{}=\overline{z}\hfill \\ & \text{or}\hfill \\ & [z]<\mathrm{min}(z,z^{})\mathrm{max}(z,z^{})<[z]+1,\hfill \end{array}$$
(2.61)
where $`z,z^{}`$ real and \[ \] denotes the integer part. Let $`(x)_j=x(x+1)\mathrm{}(x+j1),(x)_0=1`$. Below we introduce a 2-parametric family of probability measures $`M_{z,z^{}}^{(n)}`$ on $`\mathrm{Par}(n)`$. These measures take their origin in harmonic analysis on the infinite symmetric group (\[KOV\], \[Ol\]). By definition
$$\begin{array}{cc}& M_{z,z^{}}^{(n)}(\lambda )=\frac{(zz^{})^{d(\lambda )}}{(zz^{})_n}\underset{i=1}{\overset{d(\lambda )}{}}(z+1)_{p_i}(z^{}+1)_{p_i}\hfill \\ & (z+1)_{q_i}(z^{}+1)_{q_i}\frac{dim^2\lambda }{|\lambda |!}\hfill \end{array}$$
(2.62)
The conditions on $`z,z^{}`$ stated above are equivalent to the requirement that $`(z)_j(z^{})_j`$ and $`(z)_j(z^{})_j`$ are positive for any $`j=1,2,\mathrm{}`$ We note that $`M_{z,z^{}}^{(n)}`$ converges to the Plancherel measure $`M_n`$ if $`z,z^{}\mathrm{}`$. The measure $`M_{z,z^{}}^{(n)}`$ is called the $`n`$-th level $`z`$-measure. Consider now the negative binomial distribution on the non-negative integers
$$(1\xi )^{zz^{}}\frac{(zz^{})_n}{n!}\xi ^n,n=0,1,\mathrm{}$$
where $`\xi `$ is an additional parameter, $`0<\xi <1`$. The corresponding mixture of the $`n`$-level $`z`$-measures defines measure $`M_{z,z^{},\xi }`$ on Par. We remark that $`M_{z,z^{},\xi }`$ degenerates into $`M^\theta `$ if $`z,z^{}\mathrm{},\xi 0`$ in such a way that $`zz^{}\xi \theta `$. It was shown in \[BO1\] that in the modified Frobenius coordinates $`M_{z,z^{},\xi }`$ is a determinantal random point field on $`^1+\frac{1}{2}`$. The corresponding kernel can be expressed in terms of the Gauss hypergeometric function and is called the hypergeometric kernel. It appears that a number of familiar kernels can be obtained in terms of the hypergeometric kernel, in particular Hermite kernel ((2.2), (2.3), (2.7)), Laguerre kernel ((2.2), (2.28)), Meixner kernel ((2.67) below), Charlier kernel. For the hierarchy of the degenerations of the hypergeometric kernel we refer the reader to \[BO2\] §9. Recently Okounkov \[Ok1\] showed that the measures $`M_{z,z^{},\xi }`$ are the special case of an infinite parameter family of probability measures on Par, called the Schur measures, and defined as
$$M(\lambda )=\frac{1}{z}s_\lambda (x)s_\lambda (y),$$
(2.63)
where $`s_\lambda `$ are the Schur functions (for the definition of the Schur functions see \[St\] or \[Mac\]), $`x=(x_1,x_2,\mathrm{})`$ and $`y=(y_1,y_2,\mathrm{})`$ are parameters such that
$$Z=\underset{\lambda \mathrm{Par}}{}s_\lambda (x)s_\lambda (y)=\underset{i,j}{}(1x_iy_j)^1$$
(2.64)
is finite and $`\{x_i\}_{i=1}^{\mathrm{}}=\overline{\{y_i\}_{i=1}^{\mathrm{}}}`$. Measures $`M_{z,z^{},\xi }`$ formally correspond to $`_{i=1}^{\mathrm{}}x_i^m=\xi ^{\frac{m}{2}}z,_{i=1}^{\mathrm{}}y_i^m=\xi ^{\frac{m}{2}}z^{},m=1,2\mathrm{}`$ .To be precise one should consider the Newton power sums as real parameters and express the Schur functions as polynomials in the power sums. By now the reader probably would not be very surprised to learn that the Schur measures also can be considered as determinantal random point fields (\[Ok1\], Theorems 1,2)!
### 2.6 Two-Dimensional Random Growth Model
As our last example we consider the following two-dimensional random growth model (\[Jo2\]). Let $`\{a_{ij}\}_{i,j1}`$ be a family of independent identically distributed random variables with a geometric law
$$p(a_{ij}=k)=pq^k,k=0,1,2,\mathrm{}$$
(2.65)
where $`0<q<1,p=1q`$. One may think about (2.65) as the distribution of the first success time in a series of Bernoulli trials. We define
$$G(M,N)=\underset{\pi }{\mathrm{max}}\underset{(i,j)\pi }{}a_{ij},$$
(2.66)
where the maximum in (2.66) is considered over all up/right paths $`\pi `$ from (1,1) to $`(M,N)`$, in other words over $`\pi =\{(i_1,j_1)=(1,1),(i_2,j_2),(i_3,j_3),\mathrm{},(i_{M+N1},j_{M+N1})=(M,N)\}`$, such that $`(i_{k+1},j_{k+1})(i_k,j_k)\{(0,1),(1,0)\}`$. We mention in passing that distribution of random variables $`\{G(M,N)\}`$ can be interpreted in terms of randomly growing Young diagrams and totally asymmetric exclusion process with discrete time (for the details see \[Jo2\]). Without loss of generality we may assume $`MN1`$. To state explicitly the connection to the determinantal random point fields we introduce the discrete weight $`w_K^q(x)=\left(\genfrac{}{}{0pt}{}{x+K1}{x}\right)q^x,K=MN+1`$, on non-negative integers $`x=0,1,2\mathrm{}`$ The normalized orthogonal polynomials $`\{M_n(x)\}_{n0}`$ with respect to the weight $`w_K^q`$ are proportional to the classical Meixner polynomials (\[Ch\]). The kernel
$$K_{M,N}(x,y)=\underset{j=0}{\overset{N1}{}}M_j(x)M_j(y)\left(w_K^q(x)w_K^q(y)\right)^{\frac{1}{2}}$$
(2.67)
satisfies the conditions of Lemma 4 with respect to the counting measure on non-negative integers. Therefore
$$P_N(x_1,\mathrm{},x_N)=\frac{1}{N!}det\left(K_{M,N}(x_i,x_j)\right)_{1i,jN}$$
(2.68)
defines a discrete determinantal random point field. It was shown by Johansson that the distribution of the random variable $`G(M,N)`$ coincides with the distribution of the right-most particle in (2.68). After appropriate rescaling in the limit $`N\mathrm{},M\mathrm{},\frac{M}{N}\mathrm{const}`$, this distribution converges to the distribution of the right-most particle in the Airy random point field (2.36). Additional information on the subject of the last two subsections can be found in the recent papers/preprints \[AD\], \[BDJ1\], \[BDJ2\], \[BR1\], \[BR2\], \[BR3\], \[Bor\], \[ITW\], \[Ku\], \[Ok2\], \[PS1\], \[PS2\], \[TW5\], \[TW6\].
## 3 Translation Invariant Determinantal Random Point Fields
As before $`(X,,P)`$ denotes a random point field with a one-particle space $`E`$, hence $`X`$ is a space of locally finite configurations of particles in $`E`$, $``$ is a Borel $`\sigma `$-algebra of measurarble subsets of $`X`$ and $`P`$ is a probability measure on $`(X,B)`$. Throughout this section we always assume $`E=^d`$ or $`^d`$. We define a continuous action $`\{T^t\}_{tE}`$ of $`E`$ on $`X`$ in a natural way:
$$T^t:XX,(T^t\xi )_i=(\xi )_i+t.$$
(3.1)
Definition 5. Random point field $`(X,,P)`$ is called translation invariant if for any $`A`$, any $`tE`$
$$P(T^tA)=P(A).$$
The translation invariance of a random point field implies the invariance of $`k`$-point correlation functions:
$$\begin{array}{cc}& \rho _k(x_1+t,\mathrm{},x_k+t)=\rho _k(x_1,\mathrm{},x_k),\text{a.e. }\hfill \\ & k=1,2,\mathrm{},tE.\hfill \end{array}$$
(3.2)
Conversely, if $`\{\rho _k\}`$ are invariant under $`\{T^t\}`$, then there exists a corresponding random point field which is translation invariant (\[L3\]). In particular, if the translation invariant correlation functions define $`P`$ uniquely then the random point field is translation invariant. In the case of a determinantal random point field this implies the following criterion: a determinantal random point field is translation invariant if and only if the kernel $`K`$ is translation invariant, i.e., $`K(x,y)=K(xy,0)=:K(xy)`$. In this section we restrict our attention to the translation invariant determinantal random point fields. We are interested in the ergodic properties of the dynamical system $`(X,B,P,\{T^t\})`$. For the convenience of the reader recall some basic definitions of Ergodic Theory (\[CFS\]).
1. A dynamical system is said to be ergodic if the measure $`P(A)`$ of any invariant set $`A`$ equals 0 or 1.
2. A dynamical system has the mixing property of multiplicity $`r1`$ if for any functions $`f_0,f_1,\mathrm{},f_rL^{r+1}(X,,P)`$ we have
$$\underset{t_1,\mathrm{},t_r\mathrm{}}{lim}_Xf_0(\xi )f_1(T^{t_1}\xi ):\mathrm{}:f_r(T^{t_1+\mathrm{}+t_r}\xi )dF=\underset{i=0}{\overset{r}{}}_Xf_i(\xi )𝑑P$$
(3.3)
3. A dynamical system has an absolute continuous spectra if for any $`fL^2(X,B,P)`$ orthogonal to constants
$$_Xf(\xi )\overline{f(T^t\xi )}𝑑P=e^{i(t\lambda )}h_f(\lambda )𝑑\lambda ;$$
(3.4)
where the integration at the r.h.s. of (3.4) is over $`^d`$ in the continuous case and over $`[0,2\pi ]^d`$ in the discrete case, and $`h_f(\lambda )d\lambda `$ is a finite measure absolutely continuous with respect to the Lebesgue measure. One can interpret (3.4) in the following way. We define a $`d`$-parameter group of unitary operators $`\{U^t\}_{tE}`$ on $`L^2(X,B,P)`$ as
$$(U^tf)(\xi )=f(T^t\xi ).$$
Usually such family of unitary operators is called adjoint to the dynamical system. It is easy to see that $`\{U^t\}`$ commute. Since $`L^2(X,,P)`$ is separable and $`(U^t\psi ,\phi )`$ is a measurable function of $`t`$ for any $`\psi ,\phi L^2(X,B,P)`$ one can apply the von Neumann theorem (\[RS\], vol. 1, Theorem VIII.9) to conclude that $`U^t`$ is strongly continuous. In the case $`E=^d`$ one has $`h_f(\lambda )d\lambda =d(f,Q_\lambda f)`$, where $`dQ_\lambda `$ is a projection-valued measure, $`Q_\lambda =Q_{(\mathrm{},\lambda _1)\times \mathrm{}\times (\mathrm{},\lambda _d)}=_{j=1}^d\chi _{(\mathrm{},\lambda _j)}(A_j),\{A_j\}_{j=1}^d`$ are the generators of the one-parameter groups $`U^{(0,\mathrm{},t_j,0,\mathrm{}0)}`$ and $`\chi _{(\mathrm{},t)}`$ is the indicator of $`(\mathrm{},t)`$ (\[RS\], vol. I, Theorem VIII.12). In the discrete case $`E=^ddQ_\lambda `$ is a projection-valued measure on a $`d`$-dimensional torus,
$$Q_{[1,e^{i\lambda _1}]\times \mathrm{}\times [1,e^{i\lambda _d}]}=\underset{j=1}{\overset{d}{}}\chi _{[1,e^{i\lambda _j}]}(U_j),U_j=U^{(0,\mathrm{},t_j=1,\mathrm{}0)}.$$
Theorem 7. Let $`(X,B,P)`$ be a translation invariant determinantal random point field. Then the dynamical system $`(X,B,P,\{T^t\})`$ is ergodic, has the mixing property of any multiplicity and its spectra is absolutely continuous.
Remark 11. Recall that the absolute continuity of the spectra implies the mixing property of multiplicity 1, which in turn implies ergodicity (\[CFS\]).
Proof of Theorem 7. We note that the linear combinations of
$$\begin{array}{cc}& f(\xi )=\underset{j=1}{\overset{N}{}}S_{g_j}(\xi ),\hfill \\ & N1,S_g(\xi )=\underset{i}{}g(x_i),g_jC_0^{\mathrm{}}(^d),j=1,\mathrm{}N\hfill \end{array}$$
(3.5)
are dense in $`L^2(X,B,P)`$. Therefore it is enough to establish (3.3), (3.4) for the functions of such form. We start with the lemma calculating the mathematical expectation of (3.5).
Lemma 5.
a)
$$\begin{array}{cc}& 𝔼_P\underset{j=1}{\overset{N}{}}S_{g_j}(\xi )=\underset{m=1}{\overset{N}{}}\underset{\stackrel{\text{over partitions}}{_{\mathrm{}=1}^mC_{\mathrm{}}=\{1,\mathrm{},N\}}}{}\underset{\mathrm{}=1}{\overset{m}{}}[\underset{k_{\mathrm{}}=1}{\overset{\mathrm{\#}(C_{\mathrm{}})}{}}\underset{\stackrel{\text{over partitions}}{_{i=1}^k_{\mathrm{}}B_\mathrm{}i=C_{\mathrm{}}}}{}\{\underset{\sigma S^k_{\mathrm{}}}{}\hfill \\ & \frac{(1)^\sigma }{k_{\mathrm{}}}\underset{i=1}{\overset{k_{\mathrm{}}}{}}g_{B_{\mathrm{}\sigma (i)}}(x_i)K(x_{i+1}x_i)dx_1\mathrm{}dx_k_{\mathrm{}}\}]\hfill \end{array}$$
(3.6)
where $`g_{B_\mathrm{}i}(x)=_{jB_\mathrm{}i}g_j(x)`$.
b) $`𝔼_{j=1}^{N_1+\mathrm{}+N_{r+1}}S_{g_j}(\xi )_{s=1}^{r+1}\left(𝔼_{N_1+\mathrm{}+N_{s1}+1}^{N_1+\mathrm{}+N_s}S_{g_j}(\xi )\right)=(`$similar expression to (3.6), with the only difference that partitions
$$\underset{\mathrm{}=1}{\overset{m}{}}C_{\mathrm{}}=\{1,2,\mathrm{},\underset{s=1}{\overset{r+1}{}}N_s\}$$
(3.7)
satisfy (\*)), where
(\*) There exists at least one element $`C_{\mathrm{}}`$ of the partition such that the intersections of $`C_{\mathrm{}}`$ with at least two of the following sets $`\{1,\mathrm{}N_1\},\mathrm{},\{N_1+\mathrm{}+N_{s1}+1,\mathrm{},N_1+\mathrm{}+N_s\},\mathrm{},\{N_1+\mathrm{}+N_r+1,\mathrm{},N_1+\mathrm{}+N_{r+1}\}`$ are non-empty.
Proof of Lemma 5. The proof of part a) is rather straightforward and quite similar to the one given at the beginning of §2 in \[So3\] (see formulas (2.1)-(2.7) from the reference). The proof of part b) follows from a). $`\mathrm{}`$
To derive the mixing property (3.3) we replace $`g_j()`$ for $`N_1+\mathrm{}+N_{s1}+1jN_1+\mathrm{}+N_s,s=1,\mathrm{},r+1`$, in (3.7) by $`g_j(+t_1+\mathrm{}+t_{s1})`$. Fix a partition $`_{\mathrm{}=1}^mC_{\mathrm{}}=\{1,2,\mathrm{},N_1+\mathrm{}+N_{r+1}\}`$. Since $`\{g_j\}`$ are bounded functions with compact support, each of $`m`$ factors at the r.h.s. of (3.7) is bounded. We claim that the $`\mathrm{}^{th}`$ factor (corresponding to $`C_{\mathrm{}}`$, where $`\mathrm{}`$ is the same index as in (\*)) goes to zero. To see this we fix an arbitrary partition of $`C_{\mathrm{}},_{i=1}^k_{\mathrm{}}B_\mathrm{}i=C_{\mathrm{}}`$. By assumption, $`C_{\mathrm{}}`$ contains indices $`1u<vN_1+\mathrm{}+N_{r+1}`$, such that $`u`$ and $`v`$ belong to different subsets $`\{1,\mathrm{},N_1\},\mathrm{}\{N_1+\mathrm{}+N_{s1}+1,\mathrm{},N_1+\mathrm{}+N_s\},\mathrm{},\{N_1+\mathrm{}+N_r+1,\mathrm{},N_1+\mathrm{}+N_{r+1}\}`$. We claim that
$$\underset{i=1}{\overset{k_{\mathrm{}}}{}}g_{B_{\mathrm{}\sigma (i)}}(x_i)K(x_{i+1}x_i)dx_1\mathrm{}dx_k_{\mathrm{}}$$
(3.8)
goes to zero as min$`\{t_s,1sr\}\mathrm{}`$. Indeed if min$`\{t_s,1sr\}`$ is sufficiently large, the indices $`u,v`$ belong to different $`B_\mathrm{}i`$’s or the corresponding $`g_{B_\mathrm{}i}`$ is zero (the supports of the factors in $`g_{B_\mathrm{}i}`$ will not intersect). Once $`u`$ and $`v`$ belong to different $`B_\mathrm{}i`$’s the argument in $`K(x_{i+1}x_i)`$ for some $`i`$ is greater than min$`\{t_s,1sr\}`$. Since the Fourier transform of $`K(x),\widehat{K}(t)=e^{ixt}K(x)𝑑x`$ is a non-negative integrable function (bounded from above by 1), applying the Riemann-Lebesgue lemma we obtain that $`K(x_{i+1}x_i)`$ goes to zero. The other terms in (3.8) are bounded and the integration is over a bounded set, therefore (3.8) goes to zero and the proof of the mixing property follows.
To establish the absolute continuity of the spectrum we apply (3.7) when $`r=2,N_1=N_2=N,g_{N+j}(x)=\overline{g_j(x+t)},j=1,\mathrm{},N,f(\xi )=_{j=1}^NS_{g_j}(\xi ),\overline{f(T^t\xi )}=_{j=1}^NS_{\overline{g_j}}(T^t\xi )=_{j=N+1}^{2N}S_{\overline{g_j}}(\xi )`$. We have
$$\begin{array}{cc}& 𝔼(f(\xi )𝔼f)(\overline{f(T^t\xi )}\overline{𝔼f})=\underset{m=1}{\overset{2N}{}}\underset{\stackrel{\text{over partitions}}{_{\mathrm{}=1}^mC_{\mathrm{}}=\{1,\mathrm{},2N\}}}{^{}}\underset{\mathrm{}=1}{\overset{m}{}}[\underset{k_{\mathrm{}}=1}{\overset{\mathrm{\#}(C_{\mathrm{}})}{}}\hfill \\ & \underset{\stackrel{\text{over partitions}}{_{i=1}^k_{\mathrm{}}B_\mathrm{}i=C_{\mathrm{}}}}{}\{\underset{\sigma S^k_{\mathrm{}}}{}\frac{(1)^\sigma }{k_{\mathrm{}}}\underset{i=1}{\overset{k_{\mathrm{}}}{}}g_{B_{\mathrm{}\sigma (i)}}(x_i)K(x_{i+1}x_i)dx_1\mathrm{}dx_k_{\mathrm{}}\}],\hfill \end{array}$$
(3.9)
where we assume that $`x_{k_{\mathrm{}}+1}=x_1`$ in the integral, and the sume in $`^{}`$ is over partitions $`\{C_1,\mathrm{},C_m\}`$ such that for at least one element $`C_{\mathrm{}}`$ of the partition both $`C_{\mathrm{}}\{1,2,\mathrm{},N\}`$ and $`C_{\mathrm{}}\{N+1,\mathrm{},2N\}`$ are non-empty (we denoted above this property by (\*)). The terms in the product $`_{\mathrm{}=1}^m`$ corresponding to those $`\mathrm{}`$ that do not satisfy (\*) are constants as functions of $`t`$. Fix now $`\mathrm{}`$ satisfying (\*). We claim that
$$\begin{array}{cc}& \underset{i=1}{\overset{k_{\mathrm{}}}{}}g_{B_{\mathrm{}\sigma (i)}}(x_i)K(x_{i+1}x_i)dx_1\mathrm{}dx_k_{\mathrm{}}=\hfill \\ & \left(\frac{1}{2\pi }\right)^k_{\mathrm{}}\underset{i=1}{\overset{k_{\mathrm{}}}{}}\widehat{g}_{B_{\mathrm{}\sigma (i)}}(y_{i+1}y_i)\widehat{K}(y_{i+1})dy_1\mathrm{}dy_k_{\mathrm{}}\hfill \end{array}$$
(3.10)
can be written as $`e^{i(t\lambda )}h(\lambda )𝑑\lambda `$, where $`h(\lambda )`$ is an integrable function. The check is rather straightforward and we leave the details to the reader. We infer that (3.9) is a linear combination of the products of the Fourier transforms of integrable functions. Since the product of the Fourier transforms is the Fourier transform of the convolution the proof of the absolute continuity of the spectrum follows. Theorem 7 is proven. $`\mathrm{}`$
One can without difficulty calculate the spectral density of the centralized linear statistics
$$S_g(\xi )𝔼S_g=\underset{i}{}g(x_i)𝔼\underset{i}{}g(x_i).$$
Namely
$$\begin{array}{cc}& 𝔼(S_g𝔼S_g)(\overline{S_g(T^t)}𝔼\overline{S_g})=e^{i(t\lambda )}(K(0)\widehat{|K|^2}(\lambda ))\hfill \\ & \frac{1}{2\pi }|\widehat{g}(\lambda )|^2d\lambda ,\text{ and }h_{S_g}(\lambda )=(K(0)\widehat{|K|^2}(\lambda ))\frac{1}{2\pi }|\widehat{g}(\lambda )|^2\hfill \end{array}$$
(3.11)
We conclude that
$$\mu (d\lambda )=(K(0)\widehat{|K|^2}(\lambda ))d\lambda $$
(3.12)
is the spectral measure of the restriction of $`\{U^t\}`$ to the subspace of the centralized linear statistics. Since $`0\widehat{K}(\lambda )1,K(0)=\frac{1}{2\pi }\widehat{K}(\lambda )𝑑\lambda `$, we see that
$$0\frac{d\mu }{d\lambda }=K(0)\widehat{|K|^2}(\lambda )=K(0)\frac{1}{2\pi }\widehat{K}(y)\widehat{K}(y\lambda )𝑑yK(0).$$
We note that $`\frac{d\mu }{d\lambda }>0`$ for $`\lambda 0`$, and $`\frac{d\mu }{d\lambda }(0)=0`$ if and only if $`\widehat{K}(\lambda )`$ is an indicator. In particular the spectral measure $`\mu `$ is equivalent to the Lebesgue measure.
Before we formulate the next lemma recall that by $`\mathrm{\#}_{[L,L]^d}(\xi )`$ we denote the number of particles in $`[L,L]^d`$.
Lemma 6.
$$\mathrm{Var}(\mathrm{\#}_{[L,L]^d})=\mathrm{Vol}([L,L]^d)\left(\frac{d\mu }{d\lambda }(0)+\overline{o}(1)\right)\text{ as }L\mathrm{}.$$
(3.13)
Proof of Lemma 6. The probabilitists are well familiar with the analogue of this result in the Theory of Random Processes: let $`\{\eta _n\}`$ be $`L^2`$-stationary random sequence and $`h(\lambda )`$ its spectral density, $`𝔼\eta _n\overline{\eta _m}=b(nm)=\frac{1}{2\pi }_0^{2\pi }e^{i\lambda (nm)}h(\lambda )𝑑\lambda `$, then Var$`(\eta _n+\mathrm{}+\eta _n)=(h(0)+\overline{o}(1))n`$ (\[IL\], section XVIII.2). To prove the lemma we write
$$\begin{array}{cc}& \mathrm{Var}(\mathrm{\#}_{[L,L]^d})=_{[L,L]^d}_{[L,L]^d}\rho _2(x,y)\rho _1(x)\rho _1(y)dxdy\hfill \\ & +_{[L,L]^d}\rho _1(x)𝑑x=_{[L,L]^d}_{[L,L]^d}|K|^2(xy)𝑑x𝑑y\hfill \\ & +K(0)\mathrm{Vol}([L,L]^d)=(K(0)_^d|K|^2(x)𝑑x+\overline{o}(1))\mathrm{Vol}([L,L]^d)=\hfill \\ & (K(0)\widehat{|K|^2}(0)+\overline{o}(1))\mathrm{Vol}([L,L]^d).\hfill \end{array}$$
$`\mathrm{}`$
The subleading terms in (3.13) also depend on the behavior of $`\frac{d\mu }{d\lambda }`$ near the origin. For example, let $`\widehat{K}(\lambda )`$ be an indicator, $`\widehat{K}(\lambda )=\chi _B(\lambda ),B^d`$. As we have seen above this is equivalent to $`\frac{d\mu }{d\lambda }(0)=0`$. For simplicity we will assume $`d=1`$. If $`B`$ is a union of $`m`$ disjoint intervals
$$\begin{array}{cc}& \frac{d\mu }{d\lambda }(\lambda )=K(0)\frac{1}{2\pi }\widehat{K}(y)\widehat{K}(y\lambda )𝑑y=\hfill \\ & \frac{1}{2\pi }\left[\mathrm{length}(B)\mathrm{length}(B(B+\lambda ))\right]=\hfill \\ & \frac{m}{2\pi }|\lambda |(1+\overline{o}(1)),\lambda 0,\hfill \end{array}$$
(3.14)
and after more careful evaluation of the asymptotics of $`_L^L_L^L|K|^2(xy)𝑑x𝑑y=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}\widehat{|K|^2}(\lambda )(\frac{2\mathrm{sin}(L\lambda )}{\lambda })^2𝑑\lambda `$ we arrive at
$$\mathrm{Var}(\mathrm{\#}_{[L,L]^d})=\frac{m}{\pi ^2}\mathrm{log}L(1+\overline{o}(1))$$
(3.15)
Choosing $`m=1`$, $`\widehat{K}(\lambda )=X_{[\pi ,\pi ]}(\lambda )`$ one obtains the sine kernel $`K(xy)=\frac{\mathrm{sin}\pi (xy)}{\pi (xy)}`$. A special role played by the sine kernel can be highlighted by the fact that $`\frac{1}{\pi ^2}\mathrm{log}L`$ rate of the growth for Var$`(\mathrm{\#}_{[L,L]})`$ is the slowest among all translation-invariant kernels $`K(xy)`$ corresponding to projectors, $`\widehat{K}=\chi _B`$, for which $`inf(B),sup(B)`$ are the density points of $`B`$ (if $`K`$ is not a projector it follows from Lemma 6 that the rate of the growth of the variance is linear).
As an example consider $`B=_{n1}[n,n+\frac{1}{n^\gamma }],\gamma >1`$, than one has $`\frac{d\mu }{d\lambda }|\lambda |^{1{\scriptscriptstyle \frac{1}{\gamma }}}`$ and Var$`(\mathrm{\#}_{[L,L]})L^{{\scriptscriptstyle \frac{1}{\gamma }}}`$. More generally, $`\frac{d\mu }{d\lambda }|\lambda |^\alpha ,0<\alpha <1`$, implies Var$`(\mathrm{\#}_{[L,L]})L^{1\alpha }`$.
## 4 Central Limit Theorem for Counting Function and Empirical Distribution Function of Spacings
In \[CL\] Costin and Lebowitz proved the Central Limit Theorem for $`\mathrm{\#}_{[L,L]}`$ in the case of the sine kernel. The article also contains a remark on p. 71, due to Widom, that the result holds for a larger class of Random Matrix models. In its general form this theorem appeared in \[So2\].
Theorem 8. Let $`E`$ be as in (1.1), $`\{0<K_t1\}`$ a family of locally trace class operators in $`L^2(E),\{(X,,P_t)\}`$ a family of the corresponding determinantal random point fields in $`E`$, and $`\{I_t\}`$ a family of measurable subsets in $`E`$ such that
$$\mathrm{Var}_t\mathrm{\#}_{I_t}=\mathrm{Tr}(K_t\chi _{I_t}(K_t\chi _{I_t})^2)\mathrm{}\text{ }\text{as }t\mathrm{}.$$
(4.1)
Then the distribution of the normalized number of particles in $`I_t`$ (with respect to $`P_t`$) converges to the normal law, i.e.,
$$\frac{\mathrm{\#}_t𝔼\mathrm{\#}_{I_t}}{\sqrt{\mathrm{Var}_t\mathrm{\#}_t}}\stackrel{𝑤}{}N(0,1)$$
Remark 12. It was shown in \[So2\] that the condition (4.1) from Theorem 8 (the growth of the variance) is satisfied for the Airy kernel ($`K_tK`$ from (2.36), $`I_t`$ expanding), the Bessel kernel ($`K_tK`$ from (2.37), $`I_t`$ expanding) and for the families of kernels $`\{K_n\}`$ corresponding to random matrices from the Classical Compact Groups (§2.3b), §2.3c)). In all these cases Var$`{}_{t}{}^{}\mathrm{\#}_{I_t}^{}`$ growth logarithmically with respect to $`𝔼_t\mathrm{\#}_{I_t}`$.
Remark 13. To construct an example of the kernel $`0K\mathrm{Id}`$ such that $`E\mathrm{\#}_{[n,n]}=\mathrm{Tr}K\chi _{[n,n]}\mathrm{}`$ as $`n\mathrm{}`$, but Var$`\mathrm{\#}_{[n,n]}=\mathrm{Tr}(K\chi _{[n,n]}(K\chi _{[n,n]})^2)`$ stays bounded, consider $`\{\phi _n(x)\}_{n=\mathrm{}}^{\mathrm{}}`$ satisfying
a) supp $`\phi _n(n,n+1)`$,
b) $`\phi _n_{L^2}=1`$.
Then $`K(x,y)=_{n=\mathrm{}}^{\mathrm{}}(1\frac{1}{n^2+1})\phi _n(x)\overline{\phi _n(y)}`$ is the desired kernel. Indeed, $`𝔼\mathrm{\#}_{[n,n]}=_{k=n}^n(1\frac{1}{k^2+1})\stackrel{n\mathrm{}}{}\mathrm{},\mathrm{Var}\mathrm{\#}_{[n,n]}=_{k=n}^n(1\frac{1}{k^2+1})\frac{1}{k^2+1}_{\mathrm{}}^{\mathrm{}}(1\frac{1}{k^2+1})\frac{1}{k^2+1}<\mathrm{}`$. From the other side if $`0K\mathrm{Id}`$ is compact, locally trace class and $`\mathrm{Tr}K\chi _{[n,n]}+\mathrm{}`$, then $`\mathrm{Tr}K\chi _{[n,n]}(KX_{[n,n]})^2+\mathrm{}`$.
The result of Theorem 8 can be generalized to a finite number of intervals. Namely, if $`I_t^{(1)},\mathrm{},I_t^{(m)}`$ are disjoint subsets such that Cov$`{}_{t}{}^{}(\mathrm{\#}_{I_t^{(k)}},\mathrm{\#}_{I_t^{(j)}})/V_tb_{ij}`$ as $`t\mathrm{},1i,jm`$, where $`V_t`$ is some function of $`t`$ growing to infinity, then the distribution of $`((\mathrm{\#}_{I_t^{(k)}}𝔼_t\mathrm{\#}_{I_t^{(k)}}/V_t^{{\scriptscriptstyle \frac{1}{2}}})_{k=1,\mathrm{},m}`$ converges to the $`m`$-dimensional centralized normal vector with the covariance matrix $`(b_{ij})_{1i,jm}`$ (see \[So2\]).
Finally, we turn our attention to the problem of the global distribution of spacings. Let $`E=^d`$ or $`^d,\{B_j\}_{j=1}^k`$ be some bounded measurable subsets of $`E`$, and $`\{n_j\}_{j=1}^k`$ be some non-negative integers. We will be interested in the counting statistics of the following type
$$\eta _L(B_1,\mathrm{},B_k;n_1,\mathrm{},n_k):=\mathrm{\#}(x_i[L,L]^d:\mathrm{\#}_{x_i+B_j}=n_j,j=1,\mathrm{},k)$$
(4.2)
We can assume without loss of generality that $`\{B_j\}`$ are disjoint and do not include the origin. If $`d=1,k=1,B_1=(0,s]`$, then $`\eta _L((0,s]),0)`$ is the number of the nearest spacings in $`[L,L]`$ greater than $`s:\eta _L((0,s],0)=\mathrm{\#}\{x_i[L,L]:x_{i+1}x_i>s\}`$, and $`\eta _L((0,s]),n)`$ is the number of $`n`$-spacings greater than $`s:\eta _L((0,s],n)=\mathrm{\#}\{x_i[L,L]:x_{i+n+1}x_i>s\}`$. In \[So1\] we proved the convergence in law of the process $`\frac{\eta _L((0,s],0)𝔼\eta _L((0,s],0)}{L^{{\scriptscriptstyle \frac{1}{2}}}}`$ to the limiting Gaussian process in the case $`K(x,y)=\frac{\mathrm{sin}\pi (xy)}{\pi (xy)}`$. Recall that the convergence in law (Functional Central Limit Theorem) implies not only the convergence of the finite-dimensional distributions, but also the convergence of functionals continuous in the appropriate (e.g. locally uniform) topology on the space of sample paths. The proof of the Central Limit Theorem for the finite-dimensional distributions of $`\eta _L((0,s],0)`$ can be extended essentially word by word to the case of arbitrary, not necessarily translation invariant, kernel $`K(x,y)`$ and dimension $`d1`$, assuming the conditions (4.33), (4.34), (4.35) are satisfied. One can also replace $`(0,s]`$ by an arbitrary measurable bounded $`BE`$. For the convenience of the reader we sketch the main ideas of the proof of the finite-dimensional Central Limit Theorem below. Let us fix $`B_1,\mathrm{},B_k;n_1,\mathrm{},n_k`$. We construct a new (called modified) random point field such that $`\eta _L(B_1,\mathrm{}B_k;n_1,\mathrm{},n_k)`$ is equal to the number of all particles of the modified random point field in $`[L,L]^d`$. Namely we keep only those particles of the original r.p.f. for which
$$\mathrm{\#}_{x_i+B_j}=n_j,j=1,\mathrm{},k,$$
(4.3)
and throw away the particles for which (4.3) is violated. The modified r.p.f. in general will no longer be a d.r.p.f. What is important is that its correlation functions and cluster functions (see Definition 6 below) can be expressed in terms of the correlation functions of the original determinantal r.p.f. Let us denote by $`\rho _{\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}};B_1,\mathrm{},B_k;n_1,\mathrm{},n_k)`$ the $`\mathrm{}`$-point correlation function of the modified r.p.f. Suppose that
$$x_ix_j+B_p,1ij\mathrm{},1pk.$$
(4.4)
Then by the inclusion-exclusion principle
$$\begin{array}{cc}& \rho _{\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}};B_1,\mathrm{},B_k;n_1,\mathrm{},n_k)=\underset{m=0}{\overset{\mathrm{}}{}}\frac{(1)^m}{m!}\hfill \\ & \underset{\mathrm{}\text{ times}}{_{(x_1+B_1)^{n_1}\times \mathrm{}\times (x_1+B_k)^{n_k}}_{(x_{\mathrm{}}+B_1)^{n_1}\times \mathrm{}\times (x_{\mathrm{}}+B_k)^{n_k}}}\hfill \\ & _{((x_1+_{j=1}^kB_j){\scriptscriptstyle \mathrm{}(x_{\mathrm{}}+_{j=1}^kB_j)})^m}\rho _{\mathrm{}+\mathrm{}n+m}(x_1,\mathrm{},x_{\mathrm{}};\hfill \\ & x_{11},\mathrm{},x_{1n},x_{21},\mathrm{},x_{2n},\mathrm{},x_\mathrm{}1,\mathrm{},x_\mathrm{}n,y_1,\mathrm{}y_m)dy_1\mathrm{}dy_m\hfill \\ & dx_\mathrm{}1\mathrm{}dx_\mathrm{}n\mathrm{}dx_{11}\mathrm{}dx_{1n},\hfill \\ & n=n_1+\mathrm{}+n_k.\hfill \end{array}$$
(4.5)
If (4.4) is violated then the formula is quite similar, the only difference is that the exponent $`n_j`$ in $`(x_i+B_j)^{n_j}=\underset{n_j\text{ times}}{(x_i+B_j)\times \mathrm{}\times (x_i+B_j)}`$, $`1i\mathrm{},1jk`$, has to be replaced by $`n_j\mathrm{\#}(1rik:x_rx_i+B_j`$). While formulas (4.5) appear to be cumbersome and lengthy, they are nevertheless quite useful for calculating the asymptotics of the moments of $`\eta _L(B_1,\mathrm{},B_k;n_1,\mathrm{}n_k)`$. (Of course the assumption that the correlation functions of the original r.p.f. are the determinants is the key here.) Recall the definition of the cluster functions.
Definition 6. The $`\mathrm{}`$-point cluster functions $`r_{\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}}),\mathrm{}=1,2,\mathrm{},`$ of a random point field are defined by the formula
$$r_{\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}})=\underset{G}{}(1)^{m1}(m1)!\underset{j=1}{\overset{m}{}}\rho _{|G_j|}(\overline{x}(G_j))$$
(4.6)
where the sum is over all partitions $`G`$ of $`[\mathrm{}]=\{1,2,\mathrm{},\mathrm{}\}`$ into subsets $`G_1,\mathrm{},G_m,m=1,\mathrm{},\mathrm{}`$, and $`\overline{x}(G_j)=\{x_i:iG_j\},|G_j|=\mathrm{\#}(G_j)`$.
The cluster functions are also known in the Statistical Mechanics as the truncated correlated funciton and the Ursell functions. Sometimes in the literature the r.h.s. of (4.6) defines $`(1)^\mathrm{}1r_{\mathrm{}}`$. Correlation functions can be obtained from cluster functions by the inversion formula
$$\rho _{\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}})=\underset{G}{}\underset{j=1}{\overset{m}{}}r_{|G_j|}(\overline{x}(G_j)).$$
(4.7)
((4.6) is just the Möbius inversion formula to (4.7).) The integrals of$`r_{\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}})`$ over $`\underset{\mathrm{}\text{ times}}{[L,L]^d\times \mathrm{}\times [L,L]^d}=[L,L]^\mathrm{}d`$ are closely related to the cumulants $`C_j(L)`$ of the number of particles in $`[L,L]^d`$ :
$`V_1(L)=_{[L,L]^d}r_1(x_1)𝑑x_1=C_1(L)=𝔼\mathrm{\#}_{[L,L]^d},`$
$`V_2(L):=_{[L,L]^d}_{[L,L]^d}r_2(x_1,x_2)𝑑x_1𝑑x_2=C_2(L)C_1(L)=\mathrm{Var}\mathrm{\#}_{[L,L]^d}𝔼\mathrm{\#}_{[L,L]^d},`$
$`V_3(L):=_{[L,L]^d}_{[L,L]^d}_{[L,L]^d}r_3(x_1,x_2,x_3)𝑑x_1𝑑x_2𝑑x_3=C_3(L)3C_2(L)+2C_1(L)`$.
In general,
$$\underset{n=1}{\overset{\mathrm{}}{}}\frac{C_n(L)}{n!}z^n=\underset{n=1}{\overset{\mathrm{}}{}}\frac{V_n(L)}{n!}(e^z1)^n$$
(4.8)
(see \[CL\], \[So1\]). For the determinantal random point fields
$$r_{\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}})=(1)^\mathrm{}1\underset{\text{cyclic }\sigma S_{\mathrm{}}}{}K(x_1,x_2)K(x_2,x_3)\mathrm{}K(x_{\mathrm{}},x_1),$$
(4.9)
where the sum in (4.9) is over all cyclic permutations, and the term written in the body of the sum corresponds to $`\sigma =(\mathrm{1\; 2\; 3}\mathrm{}\mathrm{})`$. One can also rewrite (4.9) as
$$\begin{array}{cc}& r_{\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}})=(1)^\mathrm{}1\frac{1}{\mathrm{}}\underset{\sigma s_{\mathrm{}}}{}K(x_{\sigma (1)},x_{\sigma (2)})\hfill \\ & K(x_{\sigma (2)},x_{\sigma (3)})\mathrm{}K(x_{\sigma (\mathrm{})},x_{\sigma (1)}).\hfill \end{array}$$
(4.10)
We note that the difference between (4.9) and the formula for $`\mathrm{}`$-point correlation
$$\rho _{\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}})=\underset{\sigma s_{\mathrm{}}}{}(1)^\sigma K(x_1,x_{\sigma (1)})K(x_2,x_{\sigma (2)})\mathrm{}K(x,x_{\sigma (\mathrm{})})$$
(4.11)
is that the summation in (4.9) is only over cyclic permutations. It appears that a relation between $`\rho _{\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}};B_1,\mathrm{},B_k;n_1,\mathrm{},n_k)`$ and $`r_{\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}};B_1,\mathrm{},B_k;n_1,\mathrm{},n_k)`$ (at least when (3.19) is satisfied) is of a similar nature.
Lemma 7. Let (4.4) be satisfied. Then
$$\begin{array}{cc}& r_{\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}};B_1,\mathrm{},B_k;n_1,\mathrm{},n_k)=\underset{m=0}{\overset{\mathrm{}}{}}\frac{(1)^m}{m!}\hfill \\ & \underset{\mathrm{}\text{ times}}{_{(x_1+B_1)^{n_1}\times \mathrm{}\times (x_1+B_k)^{n_k}}\mathrm{}_{(x_{\mathrm{}}+B_1)^{n_1}\times \mathrm{}\times (x_{\mathrm{}}+B_k)^{n_k}}}\hfill \\ & _{((x_1+_{j=1}^kB_j){\scriptscriptstyle }\mathrm{}{\scriptscriptstyle }(x_{\mathrm{}}+_{j=1}^kB_j)^m}\rho _{\mathrm{}+\mathrm{}n+m,\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}};x_{11},\mathrm{},x_{1n},\hfill \\ & x_{21},\mathrm{},x_{2n}\mathrm{},x_\mathrm{}1,\mathrm{},x_\mathrm{}n,y_1,\mathrm{},y_m)dy_1\mathrm{}dy_m\hfill \\ & dx_\mathrm{}1\mathrm{}dx_\mathrm{}n\mathrm{}dx_{11}\mathrm{}dx_{1n},\hfill \end{array}$$
(4.12)
where $`\rho _{\mathrm{}+\mathrm{}n+m,\mathrm{}}`$ is defined below in (4.13).
To define $`\rho _{\mathrm{}+\mathrm{}n+m,\mathrm{}}`$ recall that $`\rho _{\mathrm{}+\mathrm{}n+m}(x_1,\mathrm{},y_m)=_{\sigma S_{\mathrm{}+\mathrm{}n+m}}(1)^\sigma K(x_1,\sigma (x_1))\mathrm{}K(y_m,\sigma (y_m))`$, where $`\sigma `$ is a permutation on the set of variables $`(x_1,\mathrm{},x_{\mathrm{}},x_{11},\mathrm{},x_\mathrm{}n,y_1,\mathrm{}y_m)`$. We write
$$\begin{array}{cc}& \rho _{\mathrm{}+\mathrm{}n+m,\mathrm{}}(x_1,\mathrm{},y_m)=\underset{\sigma S_{\mathrm{}+\mathrm{}n+m}}{^{}}(1)^\sigma K(x_1,\sigma (x_1))\mathrm{}K(y_m,\sigma (y_m)),\hfill \end{array}$$
(4.13)
where the summation in $`^{}`$ is over the permutations $`\sigma `$ satisfying the following property:
Let $`\tau `$ be a multivalued map defined on $`\{1,\mathrm{},\mathrm{}\}`$ with the values in $`\{1,\mathrm{},\mathrm{}\}`$:
$$\begin{array}{cc}& \tau (i)=\{j:\sigma (\{x_i,x_{i1},\mathrm{},x_{in}\}(\{y_1,\mathrm{},y_m\}(x_i+\underset{p=1}{\overset{k}{}}B_p)))\hfill \\ & (\{x_j,x_{j1},\mathrm{},x_{jn}\}(\{y_1,\mathrm{},y_m\}(x_j+\underset{p=1}{\overset{k}{}}B_p)))\mathrm{}\};\hfill \end{array}$$
(4.14)
then for any $`1i,j\mathrm{}`$ there exists $`N=N(i,j)`$ such that
$$\tau ^N(i)j.$$
(4.15)
Remark 14. The proof of Lemma 7 in the case $`d=1,K(x,y)=\frac{\mathrm{sin}\pi (xy)}{\pi (xy)},B_1=(0,s],n_1=0`$, was given in §3 of \[So1\]. In the general case the argument is absolutely the same. As a corollary of Lemma 7 we obtain
Lemma 8. Let
$$|K(x,y)|\psi (xy),$$
(4.16)
and (4.4) hold for the $`\mathrm{}`$-tuple $`(x_1,\mathrm{},x_{\mathrm{}})`$. Then for any $`\delta >0`$ the following estimate takes place:
$$\begin{array}{cc}& |r_{\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}};B_1,\mathrm{},B_k;n_1,\mathrm{},n_k|\mathrm{const}(\mathrm{},\delta )\hfill \\ & \underset{\text{cyclic }\sigma S_{\mathrm{}}}{}\left(\psi (x_2x_1)\psi (x_3x_2)\mathrm{}\psi (x_1x_{\mathrm{}})\right)^{1\delta }\hfill \end{array}$$
(4.17)
$`\mathrm{}`$ For the proof of Lemma 8 we refer the reader to \[So1\] §3. The key element of the proof is the upper bound on the absolute value of the $`m^{th}`$ term in (4.12) by
$$\begin{array}{cc}& \mathrm{const}_1(n,\mathrm{})\frac{1}{m!}\mathrm{const}_2^m\mathrm{min}\{\mathrm{const}_3(n,\mathrm{});(\mathrm{}+\mathrm{}n+m)!\hfill \\ & \underset{\text{cyclic }\sigma S_{\mathrm{}}}{}(\psi (x_2x_1)\psi (x_3x_2)\mathrm{}\psi (x_1x_{\mathrm{}}))\}.\hfill \end{array}$$
If $`\psi ^{1\delta }L^2(E)`$ for some $`0<\delta <1`$, then $`_{[L,L]^d}\mathrm{}_{[L,L]^d}\psi (x_2x_1)^{1\delta }\mathrm{}\psi (x_1x_{\mathrm{}})^{1\delta }𝑑x_1\mathrm{}𝑑x_{\mathrm{}}\mathrm{const}(\psi )_{[L,L]^d}\psi (xy)^{22\delta }𝑑x𝑑y=O(L^d)`$, therefore by Lemma 8
$$\begin{array}{cc}& _{[L,L]^d\mathrm{}(4.4)}r_{\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}};B_1,\mathrm{},B_k;n_1,\mathrm{}n_k)\hfill \\ & dx_1\mathrm{}dx_{\mathrm{}}=O(L^d),\mathrm{}=1,2,\mathrm{}\hfill \end{array}$$
(4.18)
In particular
$$\begin{array}{cc}& 𝔼\eta _L(B_1,\mathrm{},B_k;n_1,\mathrm{},n_k)=V_1(L)=\hfill \\ & _{[L,L]^d}r_1(x;B_1,\mathrm{},B_k;n_1,\mathrm{},n_k)𝑑x=O(L^d)\hfill \end{array}$$
(4.19)
Suppose that one could show
$$\begin{array}{cc}& \mathrm{Var}\eta _L(B_1,\mathrm{},B_k;n_1,\mathrm{}n_k)=V_1(L)+V_2(L)=\hfill \\ & _{[L,L]^d}r_1(x;B_1,\mathrm{}B_k;n_1,\mathrm{},n_k)𝑑x+_{[L,L]^d}_{[L,L]^d}\hfill \\ & r_2(x_1,x_2;B_1,\mathrm{},B_k;n_1,\mathrm{}n_k)dx_1dx_2=\mathrm{const}L^d(1+\overline{o}(1)),\hfill \end{array}$$
(4.20)
$$\begin{array}{cc}& _{[L,L]^\mathrm{}d(4.4)}r_{\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}};B_1,\mathrm{},B_k;n_1,\mathrm{}n_k)𝑑x_1\mathrm{}𝑑x_{\mathrm{}}\hfill \\ & =\overline{o}\left(L^{\frac{\mathrm{}d}{2}}\right),\mathrm{}>2.\hfill \end{array}$$
(4.21)
Since the $`\mathrm{}^{th}`$ cumulant of $`\eta _L`$ is a linear combination of $`V_i(L),i=1,2,\mathrm{}\mathrm{}`$ (see (4.8)), the estimates (4.18)–(4.21) would imly that the $`\mathrm{}^{th}`$ cumulant of $`\eta _L`$ is $`\mathrm{const}L(1+\overline{o}(1))`$ for $`\mathrm{}=2`$ and grows slower than $`L^{{\scriptscriptstyle \frac{\mathrm{}d}{2}}}`$ for $`\mathrm{}>2`$. This in turn would imply that while the second cumulant of $`\frac{\eta _L𝔼\eta _L}{\sqrt{\mathrm{Var}\eta _L}}`$ is 1, all the other cumulants of $`\frac{\eta _L𝔼\eta _L}{\sqrt{\mathrm{Var}\eta _L}}`$ go to zero as $`L+\mathrm{}`$. The last statement is equivalent to the statement that the moments of $`\frac{\eta _L𝔼\eta _L}{\sqrt{\mathrm{Var}\eta _L}}`$ converge to the moments of the normal distribution, and in particular
$$\frac{\eta _L𝔼\eta _L}{\sqrt{\mathrm{Var}\eta _L}}\stackrel{𝑤}{}N(0,1).$$
Of course the devil is in the details. It turns out that there is no nice extension of the formulas (4.12), (4.13) to the case when (4.4) is not satisfied. Below we show how one can overcome these difficulties in the case of $`\eta _L(B;0)`$ (i.e. $`k=1,n_1=0`$). We introduce the centralized $`\mathrm{}`$-point correlation functions by the formula
$$\rho _{\mathrm{}}^{(c)}(x_1,\mathrm{},x_{\mathrm{}})=\underset{G}{^{}}\underset{j=1}{\overset{m}{}}r_{|G_j|}(\overline{x}(G_j),$$
(4.22)
where $`^{}`$ is the sum over all partitions $`G=\{G_1,\mathrm{},G_m\},m=1,2,\mathrm{}`$ of $`\{1,\mathrm{},\mathrm{}\}`$ into two- and more element subsets (i.e. $`|G_j|>1,j=1,\mathrm{},m)`$. It follows from (4.7), (4.22) that
$$\begin{array}{cc}& \rho _{\mathrm{}}^{(c)}(x_1,\mathrm{},x_{\mathrm{}})=\rho _{\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}})+\underset{p=1}{\overset{\mathrm{}}{}}(1)^p\underset{1i_1<\mathrm{}<i_p\mathrm{}}{}\underset{s=1}{\overset{p}{}}\hfill \\ & \rho _1(x_{i_s})\rho _\mathrm{}p\left((x_1,\mathrm{},x_{\mathrm{}})(x_{i_1},\mathrm{},x_{i_p})\right)=\rho _{\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}})\hfill \\ & \underset{p=1}{\overset{\mathrm{}}{}}\underset{1i_1<\mathrm{}<i_p\mathrm{}}{}\underset{s=1}{\overset{p}{}}\rho _1(x_{i_s})\rho _\mathrm{}p^{(c)}\left((x_1,\mathrm{},x_{\mathrm{}})(x_{i_1},\mathrm{},x_{i_p})\right).\hfill \end{array}$$
(4.23)
Let us denote by $`M_{(\mathrm{})}^{(c)}(L)`$ the integral of the centralized $`\mathrm{}`$-point correlation function of the modified random point field over $`[L,L]^\mathrm{}d`$,
$$M_{(\mathrm{})}^{(c)}(L)=_{[L,L]^d}\mathrm{}_{[L,L]^d}\rho _{\mathrm{}}^{(c)}(x_1,\mathrm{},x_{\mathrm{}};B_1;0)𝑑x_1\mathrm{}𝑑x_{\mathrm{}}.$$
(4.24)
We have
$$\begin{array}{cc}& \underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\frac{t^{\mathrm{}}}{\mathrm{}!}𝔼(\eta _L𝔼\eta _L)^{\mathrm{}}=e^{t𝔼\eta _L}\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\frac{t^{\mathrm{}}}{\mathrm{}!}𝔼\eta _L^{\mathrm{}}=e^{t𝔼\eta _L}\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\frac{(e^t1)^{\mathrm{}}}{\mathrm{}!}\hfill \\ & 𝔼\eta _L(\eta _L1)\mathrm{}(\eta _L\mathrm{}+1)=e^{t𝔼\eta _L}e^{(e^t1)𝔼\eta _L}\hfill \\ & \underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\frac{(e^t1)^{\mathrm{}}}{\mathrm{}!}M_{(\mathrm{})}^{(c)}(L)\hfill \end{array}$$
(4.25)
If we can show that
$$M_{\mathrm{}}^{(c)}(L)=\{\begin{array}{cc}(2n1)!!\mathrm{const}_1^nL^{nd}(1+\overline{o}(1))\text{ for }\mathrm{}=2n,\hfill & \\ \overline{o}\left(L^{\frac{\mathrm{}d}{2}}\right)\text{ for }\mathrm{}=2n+1,\hfill & \end{array}$$
(4.26)
and
$$𝔼\eta _L=\mathrm{const}_2L^d(1+\overline{o}(1)),$$
(4.27)
then (4.25) implies
$$𝔼(\eta _L𝔼\eta _L)^{\mathrm{}}=\{\begin{array}{cc}(2n1)!!(\mathrm{const}_1+\mathrm{const}_2)^nL^{nd}(1+\overline{o}(1))\text{ for }\mathrm{}=2n,\hfill & \\ \overline{o}\left(L^{\frac{\mathrm{}d}{2}}\right)\text{ for }\mathrm{}=2n+1,\hfill & \end{array}$$
(4.28)
and
$$\frac{\eta _L𝔼\eta _L}{L^{\frac{d}{2}}}\stackrel{𝑤}{}N(0,\mathrm{const}_1+\mathrm{const}_2)$$
One can in principle calculate $`M_{\mathrm{}}^{(c)}(L)`$ from (4.12), (4.13). Indeed, if
$$x_ix_jB,$$
(4.29)
(we remark that (4.29) is exactly (4.4) written in the case $`k=1,n_1=0`$), then the expression for $`\rho _{\mathrm{}}^{(c)}(x_1,\mathrm{},x_{\mathrm{}};B;0)`$ can be obtained from (4.22), (4.12), (4.13). Otherwise $`\rho _{\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}};B;0)=0`$, and (4.23) implies
$$\begin{array}{cc}& \rho _{\mathrm{}}^{(c)}(x_1,\mathrm{},x_{\mathrm{}};B;0)=\underset{p=1}{\overset{\mathrm{}}{}}(1)^p\underset{1i_1<\mathrm{}<i_p\mathrm{}}{}\underset{s=1}{\overset{p}{}}\hfill \\ & r_1(x_{i_s};B;0)\rho _\mathrm{}p\left((x_1,\mathrm{},x_{\mathrm{}})(x_{i_1},\mathrm{},x_{i_p})\right).\hfill \end{array}$$
(4.30)
If for an $`(\mathrm{}p)`$-tuple $`(x_1,\mathrm{},x_{\mathrm{}})(x_{i_1},\mathrm{},x_{i_p})`$ the condition (4.29) is not satisfied, then the corresponding term $`\rho _\mathrm{}p((x_1,\mathrm{},x_{\mathrm{}})(x_{i_1},\mathrm{},x_{i_p}))`$ in (4.20) is zero. If (4.29) is satisfied for $`(x_1,\mathrm{},x_{\mathrm{}})(x_{i_1},\mathrm{},x_{i_p})`$, then we iterate (4.23) again
$$\begin{array}{cc}& \rho _\mathrm{}p\left((x_1,\mathrm{},x_{\mathrm{}})(x_{i_1},\mathrm{},x_{i_p})\right)=\hfill \\ & \rho _\mathrm{}p^{(c)}\left((x_1,\mathrm{},x_{\mathrm{}})\right)(x_{i_1},\mathrm{},x_{i_p}))+\mathrm{}\hfill \end{array}$$
We claim
Lemma 9. Let the condition (4.29) be not satisfied for the $`\mathrm{}`$-tuple $`(x_1,\mathrm{},x_{\mathrm{}})`$. Then
$$\rho _{\mathrm{}}^{(c)}(x_1,\mathrm{},x_{\mathrm{}};B;0)=\underset{\mathrm{}D\{1,\mathrm{},\mathrm{}\}}{}C_D\underset{iD}{}r_1(x_i;B;0)\rho _{|D|}^{(c)}(\overline{x}(D)),$$
(4.31)
where
$$C_D=\underset{\underset{\text{(4.29) is satisfied for }\overline{x}(A)}{AD,}}{}(1)^{|A|}.$$
(4.32)
In particular $`C_D=0`$ if (4.29) is not satisfied for $`\overline{x}(D)`$ or if there exists $`1i\mathrm{},iD`$, such that for any $`1j\mathrm{}x_ix_jB(B)`$.
Proof easily follows from the above arguments.
Theorem 9. Let $`(X,B,P)`$ be a determinantal random point field with the kernel
$$|K(x,y)|\psi (xy),$$
(4.33)
where $`\psi `$ is a bounded non-negative function such that $`\psi (\mathrm{log}(\frac{\psi +1}{\psi }))^nL^2(E)`$ for any $`n>0`$. Let for $`\eta _L(B;0)=\mathrm{\#}(x_i[L,L]^d:\mathrm{\#}(x_i+B)=0)`$ we have
$$\mathrm{Var}\eta _L(B;0)=\sigma ^2L^d(1+\overline{0}(1))$$
(4.34)
Then the Central Limit Theorem holds:
$$\frac{\eta _L(B,0)𝔼\eta _L(B;0)}{L^{\frac{d}{2}}}\stackrel{𝑤}{}N(0,\sigma ^2).$$
Remark 15. If $`\mathrm{Cov}(\eta _L(B_i;0),\eta _L(B_j;0))=b_{ij}L^d(1+\overline{0}(1)),1i,jp`$, then
$$\left(\frac{\eta _L(B_i;0)𝔼\eta _L(B_i;0)}{L^{\frac{d}{2}}}\right)_{1ip}\stackrel{𝑤}{}N(0,(b_{ij})_{1i,jp}).$$
(4.35)
Recall that
$$\begin{array}{cc}& \mathrm{Cov}(\eta _L(B_i;0);\eta _L(B_j;0))=𝔼(\eta _L(B_i;0)𝔼\eta _L(B_i,0))(\eta _L(B_j;0)\hfill \\ & 𝔼\eta _L(B_j;0))=_{\stackrel{[L,L]^{2d}}{\{x_1x_2B_i(B_j)\}}}(\underset{m=0}{\overset{\mathrm{}}{}}\frac{(1)^m}{m!}_{((x_1+B_i)(x_2+B_j))^m}\hfill \\ & \rho _{2+m,2}(x_1,x_2;y_1,\mathrm{},y_m)dy_1\mathrm{}dy_m)dx_1dx_2_{[L,L]^d}r_1(x_1;B_i;0)\hfill \\ & _{(x_1+B_i)(x_1B_j)}r_1(x_2;B_j;0)𝑑x_2𝑑x_1+_{[L,L]^d}r_1(x;B_1B_2,0)𝑑x\hfill \end{array}$$
Remark 16. Lemma 8 suggests slightly more restrictive condition on $`\psi `$, namely $`\psi ^{1\delta }L^2(E)`$ for some $`0<\delta <1`$. However, looking at the proof of Lemma 6 one immediately realizes that it is possible to replace $`\psi ^{1\delta }`$ in (4.17) by $`\psi (\mathrm{log}(\frac{\psi +1}{\psi }))^n`$ with $`n>3\mathrm{}`$.
Proof of Theorem 9. It follows from (4.25)-(4.28) that it is enough to show
$$\begin{array}{cc}& _{[L,L]^{2nd}}\rho _{2n}^{(c)}(x_1,\mathrm{},x_{2n};B;0)𝑑x_1\mathrm{}𝑑x_{2n}=\hfill \\ & (2n1)!!(_{\stackrel{[L,L]^{2d}}{\{xyB(B)\}}}r_2(x,y;B;0)dxdy_{[L,L]^d}\hfill \\ & r_1(x;B;0)_{(x+B)(xB)}r_1(y;B;0)dydx)^n+\overline{o}(L^{nd}),\hfill \\ & n=1,2,\mathrm{},\hfill \end{array}$$
(4.36)
$$\begin{array}{cc}& _{[L,L]^{(2n+1)d}}\rho _{2n+1}^{(c)}(x_1,\mathrm{},x_{2n+1};B;0)𝑑x_1\mathrm{}𝑑x_{2n+1}=\overline{o}\left(L^{\frac{2n+1}{2}d}\right),\hfill \\ & n=1,2,\mathrm{}\hfill \end{array}$$
(4.37)
Lemma 10.
$$\begin{array}{cc}& _{[L,L]^{2nd}(4.29)}\rho _{2n}^{(c)}(x_1,\mathrm{},x_{2n};B;0)𝑑x_1\mathrm{}𝑑x_{2n}=\hfill \\ & (2n1)!!\left(_{\stackrel{[L,L]^{2d}}{\{xyB(B)\}}}r_2(x,y;B;0)𝑑x𝑑y\right)^n+\overline{o}(L^{nd}),\hfill \end{array}$$
(4.38)
$$_{[L,L]^{(2n+1)d}(4.29)}\rho _{2n+1}^{(c)}(x_1,\mathrm{},x_{2n+1};B;0)𝑑x_1\mathrm{}𝑑x_{2n+1}=\overline{o}(L^{\frac{2n+1}{2}d}).$$
(4.39)
Recall that all $`r_{\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}};B;0)`$ are bounded functions (see (4.17)). Let us rewrite (4.22) as
$$\rho _{\mathrm{}}^{(c)}(x_1,\mathrm{},x_{\mathrm{}})=\underset{G}{^{}}\underset{j=1}{\overset{m}{}}r_{|G_j|}(\overline{x}(G_j))+\underset{G}{^{}}\underset{j=1}{\overset{m}{}}r_{|G_j|}(\overline{x}(G_j)),$$
where $`^{}`$ is the sum over all partitions of $`\{1,\mathrm{},\mathrm{}\}`$ into pairs, and $`^{\prime \prime }`$ is the sum over all other two- and more element partitions. Let $`\mathrm{}`$ be even, $`\mathrm{}=2n`$. Integrating $`_G^{}`$ over $`[L,L]^{2nd}(4.29)`$ we obtain exactly the r.h.s. of (4.38) (there are $`(2n1)!!`$ partitions of $`\{1,\mathrm{},2n\}`$ into two-element sets). It follows from (4.17) and the estimate below Lemma 8 that $`_{[L,L]^\mathrm{}d}|r_{\mathrm{}}(x_1,\mathrm{},x_{\mathrm{}};B;0)|𝑑x_1\mathrm{}𝑑x_{\mathrm{}}=\underset{¯}{O}(L^d)`$. Therefore the integral of $`_G^{\prime \prime }`$ over $`[L,L]^{2nd}(4.29)`$ is $`\overline{o}(L^{nd})`$. The formula (4.39) can be proven in the same way. $`\mathrm{}`$
To estimate
$$_{[L,L]^{2nd}(4.29)}\rho _{2n}^{(c)}(x_1,\mathrm{},x_{2n};B;0)𝑑x_1\mathrm{}𝑑x_{2n}$$
(4.40)
we introduce the equivalence relation on $`\{x_1,\mathrm{},x_{2n}\}`$ by calling $`x_i,x_j`$ “neighbors” if there exists a sequence of indices $`1i_0,i_1,\mathrm{},i_u2n,1u2n`$, such that $`i_0=i,i_u=j`$, and $`x_{i_{s+1}}x_{i_s}B(B),s=0,\mathrm{},u1`$. We claim that the contributions of order $`O(L^{nd})`$ appear in (4.40) only from such sets of $`(x_1,\mathrm{}x_{2n})`$ where each equivalence class of “neighbors” has either one or two indices. Consider for example the case when we have $`k`$ two-element classes $`\{x_1,x_2\},\mathrm{},\{x_{2k1},x_{2k}\}`$ and $`2n2k`$ one-element equivalence classes $`\{x_{2k+1}\},\mathrm{},\{x_{2n}\}`$. Similarly to the calculations on pp. 596-597 of \[So1\] we verify that the integral of $`\rho _{2n}^{(c)}(x_1,\mathrm{},x_{2n};B;0)`$ over the subset of $`[L,L]^{2nd}`$ corresponding to the above partition is equal to
$$\begin{array}{cc}& (2n2k1)!!\left(_{[L,L]^d}r_1(x;B;0)_{(x+B)(xB)}r_1(y;B;0)𝑑y𝑑x\right)^k\hfill \\ & \left(_{\underset{\{xyB(B)\}}{[L,L]^{2n}}}r_2(x,y;B;0)𝑑x𝑑y\right)^{nk}+\overline{o}(L^{nd}).\hfill \end{array}$$
(4.41)
After the summation over all partitions into one- and two-element equivalence classes of “neighbors” (we remark that (4.38) corresponds to the partition into singletons), we obtain exactly (4.36). It follows from Lemma 7 and (4.17) that all other partitions into the equivalence classes give negligible contributions. (4.37) can be proven in a similar fashion.
The conditions of Theorem 9 are very unrestrictive in the case of translation invariant kernels. The covariance function of the limiting Gaussian process w. $`lim\frac{\eta _L((0,\overline{s}],0)𝔼\eta _L((0,\overline{s}];0)}{L^{d/2}}`$ is then given by the $`d`$-dimensional analogues of the formulas (37), (38), (26) from \[So1\] (of course one has to replace $`\frac{\mathrm{sin}\pi (xy)}{\pi (xy)}`$ by $`K(xy)`$). Here and below we denote by $`(0,\overline{s}]`$ the rectangle $`(0,s_1]\times \mathrm{}\times (0,s_d],\overline{s}=(s_1,\mathrm{},s_d)`$. In particular, if $`K(x)`$ is continuously differentiable the limiting Gaussian process is Hölder-continuous with any exponent less than $`\frac{1}{2}`$. Among other characteristics of the modified random point field (with respect to $`B=(0,\overline{s}],n=0)`$ one may be interested in the spectral measure of the restriction of the group $`\{U^t\}`$ to the subspace of the centralized linear statistics. We shall denote the spectral measure by $`\mu ^{(s)}(d\lambda )`$. Recall that the spectral measure $`\mu ^{(0)}(d\lambda )=\mu (d\lambda )`$ of the original determinantal random point field is given by (3.12). In particular, for the sine kernel
$$\frac{d\mu }{d\lambda }=\{\begin{array}{cc}\frac{|\lambda |}{2\pi },\hfill & |\lambda |2\pi ,\hfill \\ 1,\hfill & |\lambda |>2\pi .\hfill \end{array}$$
After lengthy, but rather straightforward calculations one can obtain that in the case of the sine kernel :
$$\frac{d\mu ^{(s)}}{d\lambda }=\frac{\pi ^2s^3}{9}+\frac{|\lambda |}{2\pi }\left(1\frac{4}{3}\pi ^2s^3\right)+O(s^4)+O(|\lambda |s^4)+O(|\lambda ^2|s^2)$$
(4.42)
We note that $`\frac{d\mu ^{(s)}}{d\lambda }(0)0`$ if $`s0`$, $`s`$ small, which is consistent with Var $`\eta _L((0,s];0)L`$. For the proof of the Functional Central Limit Theorem we refer the reader to pp. 577, 598–600 of \[So1\]. Suppose that
$$L^d\frac{}{s}\eta _L((0,\overline{s}];0)),L^d\frac{}{s}\mathrm{Cov}(\eta _L((0,\overline{s}];0);\eta _L((0,\overline{t}];0))$$
(4.43)
are uniformly bounded in $`L,\overline{s},\overline{t}`$, where $`\overline{s},\overline{t}`$ belong to compact subsets of $`_+^d(_+^d)`$. By smoothing with a $`C^{\mathrm{}}`$ approximate $`\delta `$-function one can construct a continuous approximation $`\stackrel{~}{\eta }_L((0,\overline{s}];0)`$ such that $`|\stackrel{~}{\eta }_L((0,\overline{s}];0)\eta _L((0,\overline{s}];0)|1`$. As a result
$$\frac{\stackrel{~}{\eta }_L((0,\overline{s}];0)𝔼\stackrel{~}{\eta }_L((0,\overline{s}];0)}{L^{\frac{d}{2}}}$$
is a random continuous function in $`\overline{s}`$, and
$$\left|\frac{\stackrel{~}{\eta }_L((0,\overline{s}];0)𝔼\stackrel{~}{\eta }_L((0,\overline{s}];0)}{L^{\frac{d}{2}}}\frac{\eta _L((0,\overline{s}];0)𝔼\eta _L((0,\overline{s}];0)}{L^{\frac{d}{2}}}\right|\frac{2}{L^{\frac{d}{2}}}$$
(4.44)
The distribution of the random process $`\frac{\stackrel{~}{\eta }_L((0,\overline{s}];0)𝔼\stackrel{~}{\eta }_L((0,\overline{s}];0)}{L^{\frac{d}{2}}}`$ defines a probability measure on $`C([0,\mathrm{})^d)`$. By the convergence in law of random processes we mean the weak convergence of the induced probability measures on $`C([0,\mathrm{})^d)`$ (see \[B\], in general one can consider different spaces of sample paths, e.g. the space of càdlàg functions, instead of the space of continuous functions).
Theorem 10. Let the condition (4.33), (4.34), (4.35), (4.43) be satisfied. Then the random process
$$\frac{\stackrel{~}{\eta }_L((0,\overline{s}];0)𝔼\stackrel{~}{\eta }_L((0,\overline{S}];0)}{L^{\frac{1}{2}}}$$
converges in law to the limiting Gaussian process.
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# DYSON-SCHWINGER EQUATION APPROACH TO THE QCD DECONFINEMENT TRANSITION AND J/𝜓 DISSOCIATION
## 1 Introduction
Recent results of the NA50 collaboration at CERN SPS on anomalous J/$`\psi `$ suppression in Pb-Pb collisions at 158 A GeV $`^{\mathrm{?},\mathrm{?}}`$ have renewed the quest for a proper description of charmonium production in heavy-ion collisions. It has been emphasized that a threshold effect in the E<sub>T</sub> dependence of J/$`\psi `$ supression like the one observed by NA50 could be a signal for quark-gluon plasma formation $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$, see Fig. 1. However, up to now there is no theory for the hadroproduction of charmonium and the present approaches have to be considered as parametrizations of the data. We provide here a reparametrization of this threshold effect where we extract effective comover cross sections from the difference of the observed data points to the extrapolated nuclear absorption within a generalized Glauber model $`^\mathrm{?}`$, see Fig. 1, lower panel.
This comover cross section shows an onset at the transverse energy of $`E_T40`$ GeV which we interpret as a signal for a critical phenomenon in dense matter due to the chiral/deconfinement transition. Within the present paper, we want to investigate the charmonium dissociation cross section within a DSE model of QCD at high temperature and density $`^\mathrm{?}`$ and thus go beyond previous estimates of this quantity within a nonrelativistic potential model $`^\mathrm{?}`$. This is particularly important since pions as quasi Goldstone bosons are involved in the dissociation process and a nonrelativistic description of these states fails to describe their wave functions accurately. The DSE approach will be applied for the description of the temperature dependence of heavy meson masses, which can lead to a dramatic increase in the J/$`\psi `$ dissociation cross section, as has been demonstrated previously $`^\mathrm{?}`$. This behaviour could contribute to an interpretation of the still puzzling experimental situation.
## 2 Dyson-Schwinger equation approach to light and heavy mesons
The nonperturbative features of low energy QCD such as dynamical chiral symmetry breaking and confinement can be successfully modeled by the truncation of the full DSE for the renormalized dressed quark propagator $`S(p)`$, see Fig. 2, to the rainbow level $`^{\mathrm{?},\mathrm{?}}`$.
This approach has been systematically extended to the case of finite temperatures $`^\mathrm{?}`$ and chemical potentials $`^\mathrm{?}`$ within the imaginary time (Matsubara) formalism, where it gives insight into the interplay between chiral symmetry restoration and deconfinement found in Lattice QCD simulations $`^\mathrm{?}`$ and predicts the phase transition parameters in the quark matter phase diagram $`^\mathrm{?}`$.
Masses and decay properties of light mesons ($`\pi ,\rho `$) have been studied using the finite-temperature Bethe-Salpeter equations (BSEs) in the corresponding channels $`^{\mathrm{?},\mathrm{?}}`$. The DSE approach has also proven to be successful for the description of heavy meson observables $`^\mathrm{?}`$ which play a crucial rôle in the present work.
The approach we present in the following will employ a confining, separable model $`^\mathrm{?}`$ recently extended to finite temperatures in the low-mass sector $`^{\mathrm{?},\mathrm{?}}`$. In a rank-1 truncation the model gluon propagator reads $`D(p,q)=D_0\phi (p^2)\phi (q^2)`$, and the renormalized quark propagator takes the form
$`S(p)=i\gamma _\mu p_\mu \sigma _V(p^2)+\sigma _S(p^2)`$ (1)
with $`\sigma _V(p^2)=[p^2+m^2(p^2)]^1,\sigma _S(p^2)=m(p^2)\sigma _V(p^2)`$ and the dynamical quark mass function is given by
$`m(p^2)`$ $`=`$ $`m_0+\mathrm{\Delta }m(T,\mu )\phi (p^2).`$ (2)
The solution of the quark DSE reduces to that of the gap equation
$`\mathrm{\Delta }m(T,\mu )={\displaystyle \frac{16D_0}{3}}{\displaystyle \frac{d^4q}{(2\pi )^4}\phi (q^2)\sigma _S(q^2)},`$ (3)
where a nonvanishing mass gap $`\mathrm{\Delta }m(T,\mu )`$ in the chiral limit $`m_0=0`$ signals spontaneous chiral symmetry breaking in the $`T,\mu `$ -plane of the phase diagram. We choose a Gaussian formfactor $`\phi (p^2)=\mathrm{exp}(p^2/\mathrm{\Lambda }^2)`$ and obtain a fit of $`\pi `$\- and $`\rho `$\- meson properties with $`D_0=227.8`$ GeV<sup>-2</sup>, $`\mathrm{\Lambda }=0.60`$ GeV, $`m_0=7.2`$ MeV. The temperature dependence of the solution of the mass gap equation (3) results in a chiral restoration transition temperature $`T_c=150`$ MeV, see upper panel of Fig. 6. The model is confining for all $`\mathrm{\Delta }m(T,\mu )>\mathrm{\Lambda }/\sqrt{2\mathrm{e}}`$ since it has then no solution of $`p^2+m^2(p^2)=0`$ for real $`p^2`$ (no quasiparticle poles). For the heavy quark propagators the limit $`m(p^2)=m_Q=\mathrm{const}`$ will be used where $`m_Q`$ is the current mass of the heavy quark $`^\mathrm{?}`$.
The BSE in ladder approximation for the mesonic bound states has the form
$`\lambda (P^2)\mathrm{\Gamma }_M(q,P)={\displaystyle \frac{4}{3}}{\displaystyle \frac{d^4p}{(2\pi )^4}D(q,p)\gamma _\mu S_{q_1}(p_+)\mathrm{\Gamma }_M(p,P)S_{q_2}(p_{})\gamma _\mu }`$ (4)
with $`p_+=p+\eta P,p_{}=p(1\eta )P`$, where $`\eta =1/2`$ for $`m_{q_1}=m_{q_2}`$ and $`\eta =1`$ for $`m_{q_1}m_{q_2}`$. Here we will use the one-covariant meson vertex functions $`\mathrm{\Gamma }_M=\gamma _MN_M\phi _M(p^2)`$, where for the pseudoscalar mesons $`\gamma _\pi =\gamma _D=\gamma _{\overline{D}}=i\gamma _5`$ and for the vector mesons $`\gamma _\rho =\gamma _D^{}=\gamma _{\mathrm{J}/\psi }=\gamma _\mu `$ with the proper normalization constants $`N_M`$. For the present study we will use the exploratory ansatz $`\phi _M(p^2)=\phi (p^2)`$. The condition $`\lambda (P^2=M_M^2)=1`$ leads to the meson mass formulae
$`1={\displaystyle \frac{8}{3}}D_0{\displaystyle \frac{d^4p}{(2\pi )^4}\phi ^2(p^2)\left[V_M^{(V)}(p,P)\sigma _V^+\sigma _V^{}+V_M^{(S)}(p,P)\sigma _S^+\sigma _S^{}\right]}|_{P=iM_M},`$ (5)
where the functions $`V_M^{(V,S)}(p,P)`$ can be found elsewhere $`^\mathrm{?}`$. With $`m_Q=m_c=1.844`$ GeV the resulting masses for the heavy mesons (experimental value in brackets) are $`M_D=1.869`$ GeV $`((1.8693\pm 0.0005)`$ GeV), $`M_D^{}=2.006`$ GeV $`((2.0067\pm 0.0005)`$ GeV) and $`M_{J/\psi }=3.459`$ GeV $`((3.09688\pm 0.00004)`$ GeV).
## 3 Triangle diagram for the $`D^{}D\pi `$ decay
We consider this vector-pseudoscalar-pseudoscalar decay as a test for our relativistic quark model. In the light meson sector this process describes the $`\rho \pi \pi `$ decay $`^\mathrm{?}`$, which is an important ingredient of recent studies of dilepton production in heavy-ion collisions due to vector meson dominance. It has also been studied in the present model at finite temperature recently $`^\mathrm{?}`$. Here we calculate the $`D^{}D\pi `$ decay process as a quark loop diagram, see Fig. 3.
The decay width of this process is given by
$`\mathrm{\Gamma }_{D^{}D\pi }={\displaystyle \frac{g_{D^{}D\pi }^2}{192\pi M_D^{}^5}}\lambda ^3(M_D^{}^2,M_D^2,M_\pi ^2),`$ (6)
with the kinematic factor $`\lambda (s,M_1^2,M_2^2)=\sqrt{[s(M_1+M_2)^2][s(M_1M_2)^2]}`$. The $`D^{}D\pi `$ coupling constant $`g_{D^{}D\pi }=\frac{1}{2}(AB)`$ is given by the functions $`A,B`$, defined from the transition amplitude of this process
$`T_{D^{}D\pi }^\sigma `$ $`=`$ $`AP_\pi ^\sigma +BP_D^\sigma `$ (7)
$`=`$ $`N_c{\displaystyle }{\displaystyle \frac{d^4p}{(2\pi )^4}}(p^2)\text{tr}_D\{[\gamma _\sigma ][i(p_1\gamma )\sigma _V(p_1^2)+\sigma _S(p_1^2)][i\gamma _5]`$
$`[i(p_2\gamma )\sigma _V(p_2^2)+\sigma _S(p_2^2)][i\gamma _5][i(p_3\gamma )\sigma _V(p_3^2)+\sigma _S(p_3^2)]\},`$
where $`(p^2)=N_D^{}N_\pi N_D\phi ^3(p^2)`$ and $`p_1=pP_\pi /2`$, $`p_2=p+P_\pi /2`$, $`p_3=p+P_\pi /2P_D`$. Our result for the coupling constant $`g_{D^{}D\pi }=10.7`$ agrees well with the experimental value $`g_{D^{}D\pi }^{exp}=10.0\pm 1.3`$ $`^\mathrm{?}`$. The next type of diagram to be calculated is the box diagram that occurs at fourth order in the meson field expansion.
## 4 $`J/\psi `$ dissociation cross section as a quark exchange process
The $`J/\psi `$ dissociation cross section in a hot and dense medium of correlated quarks (mesons) has been studied within a nonrelativistic potential model $`^\mathrm{?}`$ for quark exchange (string-flip) processes. This calculation gives large energy dependent cross sections with a threshold enhancement (peak value about $`6`$ mb) and an exponential tail. It is, however, questionable whether the inadequate treatment of the pions in these models spoils the results. The pions which are Goldstone bosons of the broken chiral symmetry should be described within a chiral Lagrangian model. The corresponding reformulation of the charmonium dissociation process given previously $`^\mathrm{?}`$ has recently been extended and improved $`^\mathrm{?}`$. These approaches disregard the quark substructure effects which give rise to the nonlocality of the formfactors and which are expected to become apparent at finite temperatures and densities close to the QCD phase transition.
Therefore, we consider here the formulation of the $`J/\psi `$ dissociation process within a relativistic quark model at finite temperature. As a generic new mechanism on the quark level that we will study in detail is the anomalous process $`J/\psi +\pi D+\overline{D}`$ is described by the two types of diagram shown in Fig. 4.
One type represents the quark exchange process (box diagram) and the other one describes the $`D^{}`$ exchange (ladder resummed box diagram). In the language of the effective meson Lagrangian model, the box diagram is a contact term which is found to dominate over the heavy-meson exchange. Therefore we will focus on the evaluation of the cross section for the anomalous box diagram as a function of the center of mass energy $`\sqrt{s}=\sqrt{(P_{\mathrm{J}/\psi }+P_\pi )^2}`$
$`\sigma (s)={\displaystyle _{t_{min}}^{t_{max}}}𝑑t{\displaystyle \frac{|T_{J/\psi \pi D\overline{D}}|^2}{16\pi \lambda (s,M_{\mathrm{J}/\psi }^2,M_\pi ^2)}}`$ (8)
with the transition amplitude for this process being $`T_{J/\psi \pi D\overline{D}}`$ and the kinematic function $`\lambda `$ as defined above. In our calculations we put $`M_\pi =0`$, so that $`\lambda =[sM_{J/\psi }^2],s4M_D^2`$, and the physical region for this process is determined by
$`t_{\mathrm{min},\mathrm{max}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(sM_{J/\psi }^2)[(1{\displaystyle \frac{4M_D^2}{s}})^{\frac{1}{2}}1]+M_D^2.`$ (9)
In the DSE quark model the transition amplitude is given by
$`T_{J/\psi \pi D\overline{D}}`$ $`=`$ $`ϵ^\mu N_c{\displaystyle \frac{d^4p}{(2\pi )^4}\frac{1}{[p_2^2+m_Q^2][p_3^2+m_Q^2]}}`$ (10)
$`\times \text{tr}_D\{[i(p_1\gamma )\sigma _V+\sigma _S]\mathrm{\Gamma }_{\overline{D}}[i(p_2\gamma )+m_Q]\mathrm{\Gamma }_{J/\psi }^\mu `$
$`[i(p_3\gamma )+m_Q]\mathrm{\Gamma }_D[i(p_4\gamma )\sigma _V+\sigma _S]\mathrm{\Gamma }_\pi \},`$
where $`p_1=p+P_\pi /2,p_2=p+P_\pi /2P_{\overline{D}},p_3=pP_\pi /2+P_D,p_4=pP_\pi /2`$. Our result for the s dependence of the cross section $`\sigma _{J/\psi \pi D\overline{D}}^{\mathrm{box}}`$ is shown in Fig. 5. We see a critical enhancement just above the reaction threshold. Then the cross section drops down for higher values of s. This qualitative behavior is similar to the result of a calculation within a non-relativistic potential model $`^\mathrm{?}`$ but the absolute value of the cross section without in-medium effects ($`n=0`$) is much smaller since we have considered the anomalous process first which vanishes in the non-relativistic treatment.
Now we investigate the influence of a hot and dense medium on charmonium dissociation. A first effect is the modification of the $`D`$-meson mass $`M_D(T)`$ due to chiral symmetry restoration. Similar to the temperature dependence of the light quark mass the $`D`$-mass is nearly constant for lower temperatures and becomes smaller in the vicinity of the critical temperature, see Fig. 6. Secondly, this effect of chiral symmetry restoration enhances the cross section for the breakup reaction $`J/\psi +\pi D+\overline{D}`$. The enhancement of the breakup cross section has a critical temperature/density dependence at the chiral restoration transition.
## 5 J/$`\psi `$ kinetics at chiral symmetry restoration
Aiming at a qualitative discussion of possible observable effects of the chiral restoration transition to be seen in the J/$`\psi `$ production cross section in ultrarelativistic heavy-ion collisions we consider a relaxation time approximation for the evolution of the J/$`\psi `$ distribution in a hot and dense pion gas with the pion distribution function $`f_\pi (\stackrel{}{p};T,\mu _\pi )=\{\mathrm{exp}[(\sqrt{|\stackrel{}{p}|^2+M_\pi ^2}\mu _\pi )/T]1\}^1`$. This distribution function increases strongly near and above the critical temperature. This leads to a critical enhancement of thermally averaged charmonium dissociation cross section $`\sigma _{\mathrm{J}/\psi \pi D\overline{D}}v_{\mathrm{rel}}`$ and to a critical drop of the relaxation time $`\tau `$
$`\tau ^1=\sigma _{J/\psi \pi D\overline{D}}v_{\mathrm{rel}}n_\pi (T)={\displaystyle \frac{d\stackrel{}{p_\pi }}{(2\pi )^3}f_\pi (\stackrel{}{p}_\pi ,T)\sigma (\stackrel{}{p}_\pi ,\stackrel{}{p}_\psi )v_{\mathrm{rel}}\frac{P_\psi P_\pi }{E_\psi E_\pi }}`$ (11)
of the J/$`\psi `$ distribution in a comoving dense pion (parton) gas.
The result is shown in Fig. 6 lower panel. The critical enhancement of the thermally averaged dissociation cross section at the QCD phase transition temperature may lead to the result that in a heavy-ion collision above a critical energy density ($`E_T`$) threshold this additional absorption process is switched on and results in an enhanced J/$`\psi `$ suppression. This pattern may serve as an explanation for the puzzling observation of anomalous J/$`\psi `$ suppression in the NA50 experiment at CERN SpS, see Fig. 1.
## 6 Conclusions
A recently developed dynamical approach to the chiral symmetry restoration and deconfinement transition is generalized for heavy mesons and quarkonia; in-medium modifications of masses and decay constants are studied. We apply this model to calculate the cross section for charmonium dissociation by hadron impact and find that contributions from the quark exchange dominate over those from D resonance propagation processes. At the QCD phase transition the reaction rates for the J/$`\psi `$ breakup are strongly enhanced.
The approach will be further developed in order to include the $`\rho `$-meson which catalyzes exothermic J/$`\psi `$ breakup and a treatment of the normal process J/$`\psi +\pi D\overline{D^{}},D^{}\overline{D}`$, which should be larger than the anomalous one considered here. On this basis, a detailed quantitative calculation of J/$`\psi `$ production in the NA50 experiment can be attacked and predictions for the RHIC and LHC colliders can be made. A particularly new aspect is the occurrence of $`D\overline{D}`$ annihilation into $`J/\psi +\pi (\rho )`$ (gain process) under these conditions which modifies the charmonium kinetics $`^{\mathrm{?},\mathrm{?}}`$ and may lead to a saturation or even enhancement of the J/$`\psi `$ abundance.
## Acknowledgments
The authors thank P. Braun-Munzinger, J. Hüfner, B. Müller, C.D. Roberts, G. Röpke and S.M. Schmidt for stimulating discussions. MII and YLK acknowledge the hospitality of the Physics Department at the University of Rostock where part of this work has been done; GB was supported by a stipend from the Max-Planck-Gesellschaft. This work has been supported by the Heisenberg-Landau programme, the Deutscher Akademischer Austauschdienst (DAAD), the Russian Fund for Fundamental Research, contract number 97-01-01040, the National Science Foundation under Grant No. INT-9603385 and by the DFG-Graduiertenkolleg “Stark korrelierte Vielteilchensysteme” at the University of Rostock.
## References
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# Localization and delocalization in dirty superconducting wires
The discovery of the $`d`$-wave nature of the order parameter in high $`T_c`$ materials has renewed interest in unconventional superconductors with low energy quasiparticles near the Fermi energy $`\epsilon _F`$. An important question is how disorder affects the quasiparticle dynamics and the corresponding low-temperature properties of the superconductor. In Ref. it is predicted that, on energy scales $`|\epsilon \epsilon _F|`$ less than the inverse mean free time and on length scales beyond the mean free path, weak impurity scattering leads to a finite density of states (DoS) and to a diffusive dynamics of quasiparticles. In normal metals, it has been known for a long time that quantum interference imposes corrections to this picture, in the form of weak localization, and eventually, for dimensions $`2`$, exponential (Anderson) localization. The analogous question for low energy quasiparticles in unconventional superconductors has been considered only recently .
The crucial distinction between quasiparticles in a superconductor and in a normal metal, is that the former are described by a Hamiltonian of Bogoliubov-de Gennes (BdG) type. Such a Hamiltonian has an additional particle-hole grading, accompanied by a discrete particle-hole symmetry, which is absent in the Hamiltonian for (electron-like) quasiparticles in a normal metal. Symmetry plays a crucial role in the problem of Anderson localization. A classification of the symmetry classes for BdG Hamiltonians, depending on the presence or absence of time-reversal (TR) and spin-rotation (SR) symmetry, has been given by Altland and Zirnbauer . The four possibilities are denoted C, CI, D, and DIII, see table I. Ref. addressed the “zero-dimensional” (0D) case of chaotic quantum dots with superconducting leads. The higher dimensional realizations of the BdG symmetry classes, relevant for the question of localization, were studied in Refs. , mainly by field-theoretical methods involving construction and analysis of non-linear sigma models with appropriate symmetries.
In this Letter we study localization in the BdG symmetry classes for the geometry of a quantum wire, i.e., in quasi-one-dimension (quasi-1D). For this purpose, we use the Fokker-Planck (FP) approach , which is complementary to the non-linear sigma model of Refs. . Using the classification scheme of Ref. , we obtain four FP equations that control quasiparticle transport at the Fermi level in a dirty superconducting wire. Our findings are remarkable: While for classes C and CI the mean and typical values of the quasiparticle conductance $`g`$ decay exponentially with the length $`L`$ of the wire for large $`L`$, the situation in classes D and DIII is quite different. There the mean $`g`$ decays only algebraically to zero for large $`L`$ and $`\mathrm{ln}g`$ is not self-averaging, indicating a very broad distribution of the conductance and the absence of the exponential localization of the quasiparticle states at $`\epsilon _F`$. (The absence of exponential localization for class D has been announced independently in Ref. .)
It should be stressed that the BdG Hamiltonians do not conserve charge. Instead, the conserved densities are those of the energy (in all four classes) and spin (when the SR symmetry is present). Thus, the transport properties (the conductance $`g`$) studied in this Letter refer to transport of heat and spin.
We now proceed with a detailed statement of our results and their derivation. The model that we consider is that of a disordered quantum wire, with a Hamiltonian of the BdG form. We distinguish gradings corresponding to spin up/down, particle/hole, left/right movers. Denoting these with Pauli matrices $`\sigma `$, $`\gamma `$, and $`\tau `$, respectively, we write our model Hamiltonian as
$$=𝒦+𝒱,𝒦=iv_F_x\sigma _0\gamma _0\tau _3𝟙_{},$$
(1)
where $`\sigma _0`$ is the $`2\times 2`$ unit matrix in the spin grading etc. The kinetic energy $`𝒦`$ describes the propagation of right and left moving quasiparticles in $`N`$ channels at the Fermi level. The “potential” $`𝒱(x)`$ is an $`8N\times 8N`$ matrix that accounts both for the presence of disorder and of superconducting correlations. In particle/hole ($`\gamma `$) grading it reads
$$𝒱=\left(\begin{array}{cc}v& \mathrm{\Delta }\\ \mathrm{\Delta }^{}& v^\mathrm{T}\end{array}\right),$$
(2)
where $`v`$ ($`\mathrm{\Delta }`$) is a hermitian (antisymmetric) $`4N\times 4N`$ matrix, representing the impurity potential (superconducting order parameter). The form (2) of the potential $`𝒱`$ ensures that the Hamiltonian $``$ obeys particle-hole symmetry, $`=\gamma _1^\mathrm{T}\gamma _1`$ . In addition, $``$ (and hence $`𝒱)`$ may obey TR invariance $`=𝒯^{}𝒯^1`$, with $`𝒯=i\tau _1\sigma _2`$, and/or SR invariance $`=\gamma _2^\mathrm{T}\gamma _2`$.
Spatial fluctuations of the order parameter $`\mathrm{\Delta }`$ and the potential $`v`$ are taken into account by assuming that $`𝒱`$ is a Gaussian random variable with vanishing mean, i.e. that its probability functional $`P[𝒱]`$ is of the form
$$P[𝒱]\mathrm{exp}\left[\frac{\gamma \mathrm{}}{4c}_0^L𝑑x\mathrm{tr}𝒱^2(x)\right],$$
(3)
where $`\mathrm{}`$ is the mean free path, $`\gamma `$ is a numerical constant to be defined below, and $`c=1`$ ($`2`$) for class C/D (CI/DIII). The transport properties of the Hamiltonian $``$ describe the transport of spin and heat by quasiparticles in a disordered superconducting quantum wire.
Before we continue with the analysis of our model (13), some remarks about its validity and relevance are in place. One key property of the model is that, apart from corrections at very low energies due to quasiparticle localization or the appearance of a critical state, the DoS of the Hamiltonian $``$ near the Fermi level $`\epsilon _F`$ is nonzero and finite. This is related to the fact that the statistical average of the order parameter $`\mathrm{\Delta }`$ is zero in our model, cf. Eq. (3). For a dirty superconductor, such behavior is plausible if the order parameter is unconventional, as in $`d`$-wave superconductors, or when it breaks TR symmetry, as is believed to be the case for, e.g., the ruthenates , vortex lines in a (conventional) superconductor , or a normal metal wire with magnetic impurities that is weakly connected to a superconducting substrate. In all these cases the disorder leads to the existence of low-energy quasiparticle states . The Hamiltonian (1) then describes diffusion and localization of these “disorder-facilitated” quasiparticles. An altogether different scenario is that of a wire made out of an unconventional superconductor with very weak disorder. If boundary conditions are suitably chosen, one or several propagating modes can exist at $`\epsilon _F`$, whose localization properties are described by Eq. (1). In any case, one should view $``$ as an effective or coarse grained Hamiltonian, whose validity is restricted to length scales beyond the microscopic mean free path $`\mathrm{}`$. It is universal in the sense that its form is determined solely by the symmetry, and the distribution (3) provides for the existence of the diffusive regime with a finite DoS at the proper energy scale. (Note that the restrictions to the validity of our model are not different from those of related field theoretic descriptions appearing in the literature .)
We describe transport properties of the model (13) through its $`8N\times 8N`$ transfer matrix $``$ that encodes the $`x`$-dependence of an $`8N`$-component quasiparticle wavefunction $`\psi `$ satisfying the Schrödinger equation $`\psi =\epsilon \psi `$ at $`\epsilon =0`$, $`\psi (x+L)=(x+L,x)\psi (x).`$ Formally, $``$ is related to the Hamiltonian (1) as
$$(x+L;x)=\mathrm{T}_y\mathrm{exp}\left[i_x^{x+L}𝑑y\tau _3𝒱(y)\right],$$
(4)
where $`\mathrm{T}_y`$ denotes the path ordering operator for the $`y`$-integration along the wire. From Eq. (4) one finds that flux conservation (i.e., Hermiticity of $``$) and particle-hole symmetry imply that $`^{}\tau _3=\tau _3`$ and $`\gamma _1\gamma _1=^{}`$, respectively. Further, TR invariance requires $`𝒯𝒯^1=^{}`$, while SR invariance is obeyed if $`\gamma _2\gamma _2=^{}`$. The transfer matrix $``$ obeys the multiplicative rule $`(z,x)=(z,y)(y,x)`$ for $`x<y<z`$ and hence is an element of a certain Lie group $``$. The appropriate Lie groups for the four symmetry classes are listed in Table I. We note that the actual transfer matrix group is an $`8N`$-dimensional representation of the Lie group $``$, where $``$ also allows a lower dimensional (irreducible) representation for the classes C, CI, and DIII . Elements of $``$ are conveniently parameterized in terms of their polar decomposition, which, in an irreducible representation, takes the form
$`\left(\begin{array}{cc}V_1& 0\\ 0& V_2\end{array}\right)\left(\begin{array}{cc}\mathrm{cosh}X& \mathrm{sinh}X\\ \mathrm{sinh}X& \mathrm{cosh}X\end{array}\right)\left(\begin{array}{cc}V_3& 0\\ 0& V_4\end{array}\right)`$ $`(\text{C, D}),`$
$`V_1\left(\begin{array}{cc}\mathrm{cosh}X& i\mathrm{sinh}X\\ i\mathrm{sinh}X& \mathrm{cosh}X\end{array}\right)V_2`$ $`(\text{CI, DIII}).`$
Here $`V_i\mathrm{O}(4N)`$ \[Sp$`(N)`$\] for classes D/DIII \[C/CI\], for all $`i=1,2,3,4`$, and $`X`$ is a diagonal matrix with positive entries $`x_j`$. (By Kramers’ degeneracy, the elements of $`X`$ occur in pairs in class C.) The $`x_j`$ serve as radial coordinates on the Lie Group $``$. One verifies that the eigenvalues of the true $`8N\times 8N`$ transfer matrix $`^{}`$ occur in $`d`$-fold degenerate inverse pairs $`\mathrm{exp}(\pm 2x_j)`$, where the degeneracy $`d`$ is listed in Table I. Hence the number of independent $`x_j`$’s is $`4N/d`$. Finally, we note that the $`x_j`$ are related to the conductance $`g`$ through
$$g=d\underset{j=1}{\overset{4N/d}{}}\mathrm{cosh}^2x_j.$$
(5)
Our aim is to find the probability distribution of the $`x_j`$ for a transfer matrix corresponding to the model (13). Increasing the length $`L`$ of the wire by a small increment $`\delta L`$ amounts to multiplication of its transfer matrix $`(L)=(x+L,x)`$ by a transfer matrix $`^{}=(x+L+\delta L,x+L)`$. Since $`^{}`$ is close to the unit matrix, random, and statistically independent from $`(L)`$, we find that as a function of $`L`$, $`(L)`$ performs a random trajectory on its Lie group $``$. Actually, we do not need to know the full trajectory on $``$ if we are only interested in the conductance $`g`$. It is sufficient to know the trajectory of the radial coordinates $`x_j`$ of $`(L)`$ after dividing out a maximal compact subgroup $`𝒢`$ of $``$ corresponding to the angular degrees freedom of $``$ that leave the product $`^{}`$ invariant, or, in other words, to know the trajectory of the $`x_j`$ in the symmetric space $`/𝒢`$ . The subgroups $`𝒢`$ are listed in Table I. Starting from the microscopic model (1-3), one can show that the trajectory obeyed by the $`x_j`$ is a Brownian motion on the coset space $`/𝒢`$ described by the joint probability distribution $`P(x_1,\mathrm{},x_{4N/d};L)`$. The $`L`$-evolution of $`P`$ is described by a FP equation, which follows either from a direct calculation starting from Eq. (1), or from the general theory of symmetric spaces . In both cases we find
$`{\displaystyle \frac{P}{L}}`$ $`=`$ $`{\displaystyle \frac{1}{2\gamma \mathrm{}}}{\displaystyle \underset{j=1}{\overset{4N/d}{}}}{\displaystyle \frac{}{x_j}}\left[J\left({\displaystyle \frac{}{x_j}}J^1P\right)\right],`$ (6)
$`J`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{4N/d}{}}}|\mathrm{sinh}2x_j|^{m_l}{\displaystyle \underset{k>j}{\overset{4N/d}{}}}{\displaystyle \underset{\pm }{}}|\mathrm{sinh}(x_j\pm x_k)|^{m_o},`$ (7)
where the numbers $`m_l`$ and $`m_o`$ are the long and ordinary root multiplicities in the symmetric spaces $`/𝒢`$, see Table I, and $`\gamma =(4Nm_o/d)+1m_o+m_l`$. The FP equation (6) is supplemented with the boundary condition $`P/x_j=(P/J)J/x_j`$ at $`x_j=0`$. The initial condition $`=1`$ for $`L=0`$ corresponds to $`P(x_1,\mathrm{},x_{4N/d};0)=_j\delta (x_j)`$.
The FP equation (6) is the fundamental equation that governs quasiparticle transport and localization in quantum wires of the symmetry classes C, CI, D, and DIII.
In the localized regime $`LN\mathrm{}`$, typically all $`x_j`$ and their spacings are much bigger than unity, and the conductance is governed by the smallest coordinate $`x_1`$. That coordinate has a Gaussian distribution, with mean $`m_lL/\gamma \mathrm{}`$ and variance $`L/\gamma \mathrm{}`$. For classes C and CI this implies that $`g`$ is exponentially small, with
$$\mathrm{ln}g=\frac{2m_lL}{\gamma \mathrm{}},\text{var}\mathrm{ln}g=\frac{4L}{\gamma \mathrm{}},$$
(8)
with $`m_l=2`$ for class CI and $`m_l=3`$ for class C. Exponential localization for class C in quasi-1D was previously obtained by Bundschuh et al. , using the non-linear sigma model. For class D and DIII however, $`m_l=0`$, so that there is no exponential localization. Instead, $`g`$ has a very broad distribution (broader than log-normal), with an algebraic decay of the mean and the variance and an $`L^{1/2}`$-dependence of $`\mathrm{ln}g`$,
$`g=d\sqrt{{\displaystyle \frac{2\gamma \mathrm{}}{\pi L}}},\text{var}g={\displaystyle \frac{2d}{3}}g,`$ (9)
$`\mathrm{ln}g=4\sqrt{{\displaystyle \frac{L}{2\pi \gamma \mathrm{}}}},\text{var}\mathrm{ln}g={\displaystyle \frac{4(\pi 2)L}{\pi \gamma \mathrm{}}}.`$ (10)
Hence in classes D and DIII, quasiparticle states are not localized at the Fermi level. Since they are neither truly extended (typically $`g1`$ in class D and DIII), we label them critical, following terminology from the case of quantum wires with off-diagonal disorder, where similar behavior is found at the center of the band .
The effect of disorder is much less pronounced in the diffusive regime $`\mathrm{}LN\mathrm{}`$. Here, the conductance has only small fluctuations around its mean. Following the method of moments , we find $`g`$ from the FP equation by construction of evolution equations for the moments of $`g_a=d_j\mathrm{cosh}^{2a}x_j`$, $`a=1,2,\mathrm{}`$,
$`{\displaystyle \frac{\gamma \mathrm{}}{a}}{\displaystyle \frac{g_a}{L}}`$ $`=`$ $`{\displaystyle \frac{m_o}{d}}{\displaystyle \underset{n=1}{\overset{a1}{}}}g_{an}g_n{\displaystyle \frac{m_o}{d}}{\displaystyle \underset{n=1}{\overset{a}{}}}g_{an+1}g_n`$ (13)
$`+(am_o2a1+m_l)g_{a+1}`$
$`+(2aam_o+m_o2m_l)g_a.`$
In the diffusive regime one may replace the average of a product by the product of the averages, and hence one finds for $`\mathrm{}LN\mathrm{}`$
$`g={\displaystyle \frac{4N\mathrm{}}{L+\mathrm{}}}+{\displaystyle \frac{d(m_o2m_l)}{3m_o}}+𝒪(\mathrm{}/L,L/N\mathrm{}).`$ (14)
The first term in Eq. (14) is the Drude conductance, while the second term is the first quantum interference correction to the average conductance. For classes C, CI, D, and DIII it takes the values $`2/3`$, $`4/3`$, $`1/3`$, and $`2/3`$, respectively. The weak localization correction for class C was obtained earlier in Ref. . For the classes D and DIII, the correction is positive, i.e., quantum interference enhances the conductance relative to the classical Drude-like leading behavior (see also Ref. ). This is similar to the phenomenon of anti-localization in the standard symplectic symmetry class, though, as pointed out by Bocquet et al., here it is a precursor of the breakdown of exponential localization, while in the standard symplectic class localization takes over in higher order quantum corrections.
In the presence of TR symmetry, Eq. (6) is soluble. This is in contrast to the case of the FP equations for the standard and chiral symmetry classes, where only the case of broken TR symmetry was exactly solvable . The first step is a map of Eq. (6) onto a Schrödinger equation in imaginary time for the wave function $`\mathrm{\Psi }(\{x_j\};s)=\mathrm{exp}\left[\frac{1}{2}\mathrm{ln}J(\{x_j\})\right]P(\{x_j\};s)`$ for $`4N/d`$ fermions in one dimension with coordinates on the half-line $`x>0`$. They interact through a two-body potential proportional to $`m_o2`$ in the presence of a one-body potential proportional to $`(m_l2)m_l`$. Hence, for classes DIII and CI these fermions are free and they only differ by the boundary condition obeyed by their wave functions $`\mathrm{\Psi }`$ at the origin. We thus find the solutions
$`P`$ $``$ $`{\displaystyle \underset{j}{}}(x_j\mathrm{sinh}2x_j)^{m_l/2}e^{\gamma x_j^2\mathrm{}/2L}`$ (16)
$`\times {\displaystyle \underset{j<k}{}}\left(x_k^2x_j^2\right)\left(\mathrm{sinh}^2x_k\mathrm{sinh}^2x_j\right).`$
Using the method of bi-orthogonal functions it is then possible to calculate the average conductance $`g`$ for all $`N`$ and $`L`$. Here we report the result for the limit of large $`N`$, leaving the results for finite $`N`$ for a future publication,
$`g`$ $`=`$ $`{\displaystyle \frac{1}{s}}{\displaystyle \frac{4}{3}}+4{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\pi ^2n^2/4s}\left({\displaystyle \frac{1}{s}}+{\displaystyle \frac{2}{\pi ^2n^2}}\right)\text{CI},`$ (17)
$`g`$ $`=`$ $`{\displaystyle \frac{1}{s}}+{\displaystyle \frac{2}{3}}4{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\pi ^2n^2/2s}{\displaystyle \frac{1}{\pi ^2n^2}}\text{DIII},`$ (18)
where $`s=L/(4N\mathrm{})`$. Note the agreement with Eq. (14) in the diffusive regime $`s1`$. In the localized regime $`s1`$, Eq. (18) may be resummed, and the asymptotic result (10) is reproduced for class DIII, while for class CI one finds $`g=8(\pi s)^{1/2}\mathrm{exp}(4s)`$. Quite remarkably, the exact results (18) for $`g`$ in classes CI and DIII are related to the average conductance $`g_{\mathrm{ch}}`$ in the chiral unitary ensemble for odd channel number as $`g(s/2)_{\mathrm{CI}}+2g(s)_{\mathrm{DIII}}=4g(s)_{\mathrm{ch}}`$.
The absence of localization in wires of classes D and DIII may have important implications for higher dimensions, provided our results can be extended beyond 1D, and provided they are not restricted to the regime of weak disorder. With respect to the latter restriction, we can point to the close formal similarity of the delocalization for the FP equations of class D/DIII and the corresponding FP equation for the chiral symmetry classes with odd $`N`$, where it is understood that the absence of localization holds both for weak and strong disorder . Thus arguing that quasiparticle states at the Fermi level remain delocalized for arbitrary disorder strength and dimensionality in the D-classes, our result suggests a possible resolution of a controversy in the literature surrounding 2D disordered superconductors of class D . While all Refs. assumed existence of two localized phases, distinguished by the quantized value of the Hall conductivity $`\sigma _{xy}`$, and a metallic phase, the proposed global phase diagrams and transitions between the phases differ considerably. We suggest that the solution might simply lie in the absence of localized phases for classes D and DIII in any dimension $`1`$.
In conclusion, we considered quasiparticle transport and localization in disordered quasi-1D superconducting wires at the Fermi level for the four Bogoliubov-de Gennes symmetry classes C, CI, D, and DIII. We obtained and solved the Fokker-Planck equations for the probability of the radial coordinates of the transfer matrix. While quasiparticle states are localized in classes C/CI, localization is absent if spin-rotation symmetry is broken (classes D/DIII).
We thank A. Altland, L. Balents, M. P. A. Fisher, N. Read, T. Senthil, and M. Sigrist for valuable discussions. PWB gratefully acknowledges that this problem was suggested to him in an earlier stage by A. Altland. Close to completion of this work, we learned that J. T. Chalker and coworkers obtained independently similar results for class D, see also Ref. . This work was supported by a Grant-in-Aid for Scientific Research from Japan Society for the Promotion of Science No. 11740199 (AF), and by the NSF under grant No. DMR-9528578 (IAG).
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# Blocking measures for asymmetric exclusion processes via coupling
## 1 Introduction
We consider the exclusion process $`\eta _t`$ on $`\{0,1\}^{}`$ with generator $`L`$ given by
$$Lf(\eta )=\underset{x}{}\underset{y}{}p(x,y;\eta )[f(\eta ^{x,y})f(\eta )].$$
(1.1)
Here $`f`$ is a continuous function on $`\{0,1\}^{}`$ (with the product topology) and for $`x,y,z`$,
$$\eta ^{x,y}(z)=\{\begin{array}{cc}\eta (x),\hfill & \text{if }z=y\text{,}\hfill \\ \eta (y),\hfill & \text{if }z=x\text{,}\hfill \\ \eta (z),\hfill & \text{otherwise.}\hfill \end{array}$$
(1.2)
The jump rate of particles from $`x`$ to $`y`$ in configuration $`\eta `$, $`p(x,y;\eta )`$, is a continuous function of $`\eta `$ which is zero unless $`\eta (x)=1\eta (y)=1`$.
Let
$$𝒳_n=\{\eta \{0,1\}^Z:\underset{xn}{}\eta (x)=\underset{x>n}{}(1\eta (x))<\mathrm{}\}$$
(1.3)
and
$$𝒳=_n𝒳_n.$$
(1.4)
The set $`𝒳`$ is countable; we will call elements of $`𝒳`$ blocking configurations, and call probability measures supported on $`𝒳`$ blocking measures. Our interest is to find sufficient conditions on the rates $`p`$ for the existence of blocking measures which are invariant for the process $`\eta _t`$.
We will construct the process on blocking configurations directly. For the construction we will use two conditions on the rates, which we assume throughout the paper; these could be somewhat weakened at the price of increasing the complexity of the exposition. (Liggett gives conditions on the rates which assure the existence of the process started from an arbitrary initial condition.) First, we take the rates to be uniformly bounded; we can set the upper bound equal to one by a time-scale change, and thus assume that
$$0p(x,y;\eta )1.$$
(1.5)
Second, we assume that the total rate for exiting any configuration of $`𝒳`$ is finite:
$$\text{For any }\zeta 𝒳,\underset{x}{}\underset{y}{}\overline{p}(x,y;\zeta )<\mathrm{}.$$
(1.6)
This condition follows from (1.5) if there is an upper bound on the range of jumps.
Various special cases are of interest. The rates are simple when they are independent of the configuration except for the exclusion condition, so that
$$p(x,y;\eta )=c(x,y)\eta (x)(1\eta (y)),$$
(1.7)
and are translation invariant when
$$p(x,y;\eta )=p(0,yx;\tau _x\eta ),$$
(1.8)
where $`\tau _z`$ is the operator of translation by $`z`$. When both of these conditions are satisfied the rates can be written in the form
$$p(x,y;\eta )=a(yx)\eta (x)(1\eta (y)).$$
(1.9)
Liggett exhibits invariant blocking measures in the case of simple, translation invariant rates with jumps restricted to length 1: $`a(z)=0`$ for $`|z|>1`$ and $`a(1)>a(1)`$. A trivial extension of his result is the following: if for some $`\alpha <1`$ the rates have the form (1.7) with
$$c(x,y)=\alpha ^{xy}c(y,x)\text{for all }x<y,$$
(1.10)
then the product measure $`\mu `$ with marginals
$$\mu (\eta (x)=1)=\frac{1}{1+\alpha ^x}$$
(1.11)
is reversible for the process $`\eta _t`$. This is a special case of a more general construction which we describe in the appendix.
For a more general set of rates, one might expect that blocking measures exist when the process has a sufficiently strong positive drift, for example in the translation invariant simple case (that is, for rates satisfying (1.9)) when
$$\underset{\{y\}}{}ya(y)>0,$$
(1.12)
(positive mean drift for the underlying random walk). Proving that (1.12) or a similar condition implies the existence of blocking measures seems quite difficult; this is one of the open problems of . When the rates $`p(x,y;\eta )`$ depend on the configuration $`\eta `$ at sites other than $`x`$ and $`y`$, it is not even clear what necessary and/or sufficient condition to conjecture. We do not deal directly with conditions like (1.12), but give a different sort of sufficient condition, showing that when the rates of two processes are appropriately related, existence of a blocking measure for one implies existence for the second.
Note that if $`\mu `$ is any invariant blocking measure for $`\eta _t`$ then $`\mu (𝒳_n)0`$ for some $`n`$; since each $`𝒳_n`$ is a closed set for the process, the conditional measure $`\mu _n=\mu (|𝒳_n)`$ is then also an invariant blocking measure. Thus, if we permit ourselves a translation of the entire system, there is no loss of generality in treating existence of a blocking measure on $`𝒳`$ as equivalent to the existence of a blocking measure on $`𝒳_0`$. We remark that if the rates are simple and translation invariant (see (1.9)) then $`𝒳_n`$ is irreducible whenever there is a positive rate for some forward and some backward jump, and the greatest common divisor of $`\{x0:a(x)>0\}`$ is 1, so that under these condition each $`\mu _n`$ is unique and extremal in the class of invariant blocking measures.
We now compare the process $`\eta _t`$ with a second process $`\overline{\eta }_t`$ for which the generator $`\overline{L}`$ is constructed as in (1.1) but with rates $`\overline{p}(x,y;\eta )`$. Our main result, presented in Section 4, gives conditions on the rates $`p`$ and $`\overline{p}`$ under which the existence of a blocking invariant measure for the process $`\overline{\eta }_t`$ implies the existence of such a measure for $`\eta _t`$. In the case in which the rates are simple and translation invariant, it takes the following form:
###### Theorem 1.1
Suppose that $`p(x,y;\eta )=a(yx)\eta (x)(1\eta (y))`$ and that
$`a(x)\overline{a}(y),`$ for $`0<xy,`$ (1.13)
$`a(y)\overline{a}(x),`$ for $`yx<0.`$ (1.14)
Then if $`\overline{\eta }_t`$ has a blocking invariant measure, so does $`\eta _t`$.
For example, we may take the weights $`\overline{p}`$ to have the form (1.10), with $`c(x,y)=a(yx)`$ for $`x<y`$, as in (1.9), so that the requisite blocking measure is given by (1.11).
We remark that establishing the existence of invariant blocking measures is a special case, and perhaps a first step toward the general case, of the problem of establishing the existence of invariant shock measures: measures on $`\{0,1\}^{}`$ which have distinct asymptotic limits to the right and left of the origin and which are time invariant in some appropriate sense, usually for the process as seen from a suitable random viewpoint. Such measures are related to the shock solutions of the Burgers equation, which describes the process in the hydrodynamical limit. The left and right asymptotic measures will be time invariant for the process in the usual sense, so that invariant shock measures appear in systems that have more than one translation invariant state. Given two such asymptotic measures, the shock measure describes one ultimate fate of the system when it starts with one of these on each side of the origin (another is the so called rarefaction fan). The blocking measures are the simplest shock measures: conceptually, because they are invariant when seen from a fixed viewpoint, and technically, because they have support on a countable state space.
In the case of simple exclusion the extremal time and translation invariant measures are the one parameter family of homogeneous product measures indexed by density. In nearest neighbor asymmetric simple exclusion, existence of invariant shock measures has been established for the process as seen from a “second class particle”, (). The approach of and was closely based on the known blocking measures for this process, the product measures (1.11). In other approaches for the problem of describing shock measures are proposed.
The paper is organized as follows. In Section 2 we construct $`\eta _t`$ on $`𝒳`$ using Poisson processes (the Harris graphical construction); the construction is made in such a way as to facilitate an appropriate coupling of two such process. We describe in Section 3 the key idea for the proof of our results: the introduction of a certain partial order $``$ on the space $`𝒳_0`$ of blocking configurations with the property that, under the coupling, the conditions of Theorem 1.1 (or the more general conditions to be given later) imply that if the initial configurations $`\eta _0`$ and $`\overline{\eta }_0`$ satisfy $`\eta _0\overline{\eta }_0`$, then this ordering is preserved by the dynamics: $`\eta _t\overline{\eta }_t`$ for all $`t0`$. In Section 4 we state and prove our general result, of which Theorem 1.1 is an immediate corollary. In Section 5 we give some applications, and in the appendix discuss the construction of a class of possible comparison processes $`\overline{\eta }_t`$.
## 2 Construction of the process
We exhibit now a special construction of the process in $`𝒳`$. The construction requires that the rates $`p(x,y;\zeta )`$ satisfy conditions (1.5) and (1.6) of the introduction.
For a configuration $`\eta 𝒳`$ we define ordered positions of the particles and empty sites by
$`x_0(\eta )=\mathrm{min}\{x:\eta (x)=1\},`$ (2.1)
$`x_k(\eta )=\mathrm{min}\{x>x_{k1}(\eta ):\eta (x)=1\},`$ (2.2)
$`y_0(\eta )=\mathrm{max}\{x:\eta (x)=0\},`$ (2.3)
$`y_k(\eta )=\mathrm{max}\{x<y_{k1}(\eta ):\eta (x)=0\}.`$ (2.4)
For each pair $`(i,j)`$ with $`i,j0`$ let
$$\mathrm{\Theta }^{i,j}=:\{((T_n^{i,j},U_n^{i,j}),(R_m^{i,j},V_m^{i,j})):n,m1\}$$
(2.5)
be a process with the following properties:
* Both $`(T_n^{i,j}T_{n1}^{i,j})_{n1}`$ and $`(R_m^{i,j}R_{m1}^{i,j})_{m1}`$, where by convention $`T_0^{i,j}=R_0^{i,j}=0`$, are families of independent exponentially distributed random variables of mean one. In other words, $`(T_n^{i,j})`$ and $`(R_m^{i,j})`$ are Poisson processes of rate $`1`$ for all $`i,j`$.
* Both $`(U_n^{i,j})_{n1}`$ and $`(V_m^{i,j})_{m1}`$ are families of independent random variables, uniformly distributed in $`[0,1]`$.
* All four of these families of variables are mutually independent.
We also assume that $`\{\mathrm{\Theta }^{i,j}:i,j0\}`$ is a family of mutually independent processes. The times $`T_n^{i,j}`$ and $`R_m^{i,j}`$ will be called Poisson events and the associated random variables $`U_n^{i,j}`$ and $`V_m^{i,j}`$ will be called marks.
We now construct the process $`\eta _t`$ as a function of the marked Poisson processes and the initial configuration $`\eta _0𝒳`$. Set $`\tau _0=0`$ and suppose inductively that we have defined times $`\tau _0,\mathrm{},\tau _{n1}`$ and configurations $`\eta _{\tau _0},\mathrm{},\eta _{\tau _{n1}}`$. Define
$$\tau _n=\mathrm{min}\{\underset{i,j,k}{inf}\{T_k^{i,j}>\tau _{n1}:U_k^{i,j}<A_+(\eta _{\tau _{n1}},i,j)\},\underset{i,j,k}{inf}\{R_k^{i,j}>\tau _{n1}:V_k^{i,j}<A_{}(\eta _{\tau _{n1}},i,j)\}\},$$
(2.6)
where for $`i,j0`$,
$`A_+(\eta ,i,j)`$ $`=`$ $`p(x_i(\eta ),y_j(\eta );\eta )\mathbf{\hspace{0.17em}1}\{y_j(\eta )>x_i(\eta )\},`$ (2.7)
$`A_{}(\eta ,i,j)`$ $`=`$ $`p(x_i(\eta ),y_j(\eta );\eta )\mathbf{\hspace{0.17em}1}\{y_j(\eta )<x_i(\eta )\}.`$ (2.8)
Here $`\mathrm{𝟏}S`$ denotes the characteristic function of the set $`S`$. If $`(I_n,J_n)`$ is the pair $`(i,j)`$ such that $`T_k^{i,j}`$ or $`R_k^{i,j}`$ realizes the infimum $`\tau _n`$ for some $`k`$, set
$`X_n`$ $`=`$ $`x_{I_n}(\eta _{\tau _{n1}}),`$ (2.9)
$`Y_n`$ $`=`$ $`y_{J_n}(\eta _{\tau _{n1}}),`$ (2.10)
and define
$`\eta _{\tau _n}`$ $`=`$ $`(\eta _{\tau _{n1}})^{X_n,Y_n}.`$ (2.11)
This completes the induction step. To finish the construction after all $`\tau _n`$ and $`\eta _{\tau _n}`$ are defined, set
$$\eta _t=\underset{n0}{}\eta _{\tau _n}\mathrm{𝟏}\{\tau _nt<\tau _{n+1}\}\text{for all }t0.$$
(2.12)
It is important to notice that after each jump the particles and holes are effectively relabeled according to (2.1)–(2.4), so that for all times $`t`$,
$$x_i(\eta _t)x_{i+1}(\eta _t)\text{ and }y_j(\eta _t)y_{j+1}(\eta _t),i,j0.$$
(2.13)
The construction may be described in words as follows. We use independent times ($`T_n^{i,j}`$ and $`R_m^{i,j}`$, respectively) for jumps to the right and jumps to the left; this is not necessary for the construction here but ensures that the coupling we define later preserves a certain partial order on configurations. The instant $`\tau _n`$ is the first time after $`\tau _{n1}`$ at which a jump is performed, and is the minimum of the first scheduled jump times to the right and to the left. The first scheduled jump time to the right is the first $`T_k^{i,j}`$ for which the corresponding uniform random variable $`U_k^{i,j}`$ is smaller than the threshold $`A_+`$, defined by (2.7) to ensure that the jump is indeed to the right and occurs at the correct rate (here we use the condition (1.5) that $`p(x,y;\eta )1`$). Similarly, the first scheduled jump time to the left is the first $`R_k^{i,j}`$ for which the corresponding uniform random variable $`V_k^{i,j}`$ is smaller than the threshold $`A_{}`$ defined by (2.8). The configuration at time $`\tau _n`$ is then the one obtained by interchanging the hole and the particle whose indexes $`i,j`$ correspond to the $`R_k^{i,j}`$ or $`T_k^{i,j}`$ that realizes the time $`\tau _n`$.
To see that the above is well defined for initial configurations in $`𝒳`$ it suffices to see that, for any initial $`\eta _0𝒳`$, $`\tau _n`$ is with probability one a strictly increasing sequence of (finite) times. The conditional distribution of $`\tau _n\tau _{n1}`$ given the past up to $`\tau _{n1}`$ is
$`(\tau _n\tau _{n1}>s|\eta _{\tau _{n1}})`$ $`=`$ $`\mathrm{exp}\left\{{\displaystyle \underset{x,y}{}}p(x,y;\eta _{\tau _{n1}})\right\},`$ (2.14)
by (2.6) (it is the minimum of independent random variables with exponential distribution and inverse-mean $`p(x,y;\eta _{\tau _{n1}})`$). Since $`\eta _{\tau _{n1}}`$ is obtained by doing at most $`n1`$ modifications to the initial configuration $`\eta _0`$, it belongs to $`𝒳`$. By condition (1.6) of the introduction, the conditional law (2.14) is that of a non-degenerate exponential random variable.
It is tedious but easy to show that the process $`\eta _t`$ so constructed in $`𝒳`$ has generator $`L`$ (restricted to $`𝒳`$). We remark that the above construction works also if the process restricted to $`𝒳`$ has explosions, that is, if $`lim_n\mathrm{}\tau _n<\mathrm{}`$.
We give now a graphical interpretation of this construction, and of the coupling of the processes to be introduced later. For simplicity assume $`𝒳=𝒳_0`$. To each configuration $`\eta 𝒳_0`$ associate an interface $`\mathrm{\Phi }\eta `$ corresponding to the integrated profile of $`\eta `$. Here $`\mathrm{\Phi }:𝒳_0_+^{}`$ is defined by either of two equivalent expressions:
$`(\mathrm{\Phi }\eta )(x)`$ $`=`$ $`x+2{\displaystyle \underset{yx}{}}\eta (y)`$ (2.15)
$`=`$ $`x+2{\displaystyle \underset{y>x}{}}(1\eta (y))`$ (2.16)
Note that $`\mathrm{\Phi }\eta `$ increases by one when a particle is present at $`x`$ or decreases by one when no particle is present at $`x`$, so that in particular, $`|\mathrm{\Phi }\eta (x)\mathrm{\Phi }\eta (x+1)|=1`$. The graph $`\{(x,(\mathrm{\Phi }\eta )(x))x\}`$ is a subset of the lattice $`_{\mathrm{even}}^2=\{(x,y)^2x+y\text{ is even}\}`$. The Heaviside configuration $`\eta ^H`$, given by $`\eta ^H(x)=\mathrm{𝟏}\{x1\}`$, gives rise to the interface $`\mathrm{\Phi }\eta ^H(x)=|x|`$.
The interface picture yields a geometric interpretation of the construction of the process $`\eta _t`$. Index the squares (plaquettes) of the lattice $`_{\mathrm{even}}^2`$ as $`\{S_{i,j}i,j\}`$ as shown in Figure 1 ($`S_{i,j}=\{(x,y)2i<x+y<2i+2,2j<yx<2j+2\}`$); with this convention, the interface $`\mathrm{\Phi }\eta `$ lies above $`S_{i,j}`$ ($`i,j0`$) if and only if $`x_i(\eta )<y_j(\eta )`$. Now think of the marked processes $`(T_n^{i,j},U_n^{i,j})`$ and $`(R_m^{i,j},V_m^{i,j})`$ as associated with $`S_{i,j}`$. When at the Poisson event $`T_n^{i,j}`$ the corresponding uniform variable $`U_n^{i,j}`$ is less than $`p(x_i(\eta ),y_j(\eta );\eta )`$, then, if the interface $`\mathrm{\Phi }\eta `$ lies above $`S_{i,j}`$, we update the interface by decreasing its height by two units in the interval $`(x_i,y_j]`$. Similarly, when at time $`R_m^{i,j}`$ the corresponding mark satisfies $`V_m^{i,j}<p(x_i(\eta ),y_j(\eta );\eta )`$ and the interface lies below $`S_{i,j}`$, we increase by two units the height of the interface in the interval $`(x_i,y_j]`$. All of this is shown in Figure 1.
## 3 An order relation on configurations
For configurations $`\eta `$ and $`\overline{\eta }𝒳_0`$, we say that
$$\eta \overline{\eta }\text{ if and only if for all }i,j0\text{}x_i(\eta )x_i(\overline{\eta })\text{ and }y_j(\eta )y_j(\overline{\eta }).$$
(3.1)
It is easy to see that this is a partial order which corresponds to the natural order on interfaces:
$$\eta \overline{\eta }\text{ if and only if }(\mathrm{\Phi }\eta )(x)(\mathrm{\Phi }\overline{\eta })(x)\text{ for all }x.$$
(3.2)
Under this ordering, the Heaviside configuration $`\eta ^H`$ precedes every other configuration: $`\eta ^H\eta `$ for any $`\eta 𝒳_0`$. From (2.15) and (2.16) it follows that if $`\eta \overline{\eta }`$ then for all $`z`$,
$`(\mathrm{\Phi }\overline{\eta })(z)(\mathrm{\Phi }\eta )(z)`$ $`=`$ $`2{\displaystyle \underset{i}{}}\mathrm{𝟏}\{x_i(\overline{\eta })z<x_i(\eta )\}`$ (3.3)
$`=`$ $`2{\displaystyle \underset{j}{}}\mathrm{𝟏}\{y_j(\eta )z<y_j(\overline{\eta })\},`$ (3.4)
and for all $`x,y`$ such that $`\eta (x)=1`$ and $`\eta (y)=0`$ and all $`z`$,
$$(\mathrm{\Phi }\eta ^{x,y})(z)=(\mathrm{\Phi }\eta )(z)2\mathbf{\hspace{0.17em}1}\{xz<y\}+2\mathbf{\hspace{0.17em}1}\{yz<x\}.$$
(3.5)
The following lemma says essentially that if we have two configurations which are ordered by $``$ then they will remain ordered after either (i) a jump in both configurations, in the same direction, of the $`i^{\mathrm{th}}`$ particle to the $`j^{\mathrm{th}}`$ hole, or (ii) certain jumps in only one of the configurations.
###### Lemma 3.1
Assume $`\eta \overline{\eta }`$, fix $`i`$ and $`j`$, and let $`x=x_i(\eta )`$, $`y=y_j(\eta )`$, $`\overline{x}=x_i(\overline{\eta })`$, and $`\overline{y}=y_j(\overline{\eta })`$. Then jumps preserve ordering in the following cases:
If $`\overline{x}x<y\overline{y}`$, then $`\eta ^{x,y}\overline{\eta }`$. (3.6)
$`\text{If }y\overline{y}<\overline{x}x\text{, then }\eta \overline{\eta }^{\overline{x},\overline{y}}.`$ (3.7)
If $`\overline{x}x<y\overline{y}`$, then $`\eta ^{x,y}\overline{\eta }^{\overline{x},\overline{y}}`$ (3.8)
If $`y\overline{y}<\overline{x}x`$, then $`\eta ^{x,y}\overline{\eta }^{\overline{x},\overline{y}}`$. (3.9)
$`\text{If }x>y\text{ and }\overline{x}<\overline{y}\text{, then }\eta \overline{\eta }^{\overline{x},\overline{y}}\text{ and }\eta ^{x,y}\overline{\eta }.`$ (3.10)
Before giving a formal proof of this lemma, we describe its graphical interpretation. The interface $`\mathrm{\Phi }\eta `$ lies below $`\mathrm{\Phi }\overline{\eta }`$. In cases (3.6) and (3.8) the square $`S_{i,j}`$ lies below both interfaces, so that for either interface a jump of the $`i^{\mathrm{th}}`$ particle to the $`j^{\mathrm{th}}`$ hole—briefly, an $`(i,j)`$ jump—lowers the interface; (3.6) and (3.8) assert respectively that the order is preserved by either a jump in the lower interface only, or a jump for both interfaces. Similarly, in cases (3.7) and (3.9) $`S_{i,j}`$ lies above both interfaces, an $`(i,j)`$ jump raises either interface, and the order is preserved by such a jump in either the upper interface alone or in both. Finally, in case (3.10) $`S_{i,j}`$ lies between the two interfaces, an $`(i,j)`$ jump for the lower interface raises it and for the upper interface lowers it, and (3.10) asserts that such a jump for either interface alone preserves the order. These properties are easy to check in the graphical representation.
Proof of Lemma 3.1: Statements (3.6) and (3.7) follow immediately from (3.5). Under the hypothesis of (3.8) $`x<y`$ and $`\overline{x}<\overline{y}`$. Hence, by (3.5),
$`(\mathrm{\Phi }\eta ^{x,y})(z)`$ $`=`$ $`(\mathrm{\Phi }\eta )(z)2\mathbf{\hspace{0.17em}1}\{xz<y\};`$ (3.11)
an analogous identity holds for $`\overline{\eta }`$. Since $`\eta \overline{\eta }`$ and $`\overline{x}x<y\overline{y}`$, by (3.3) and (3.4),
$`(\mathrm{\Phi }\eta )(z)(\mathrm{\Phi }\overline{\eta })(z)2\mathbf{\hspace{0.17em}1}\{\overline{x}z<x\}2\mathbf{\hspace{0.17em}1}\{yz<\overline{y}\}.`$ (3.12)
Subtracting $`2\mathbf{\hspace{0.17em}1}\{x<zy\}`$ in both members of the above inequality we get
$`(\mathrm{\Phi }\eta )(z)2\mathbf{\hspace{0.17em}1}\{xz<y\}`$ $``$ $`(\mathrm{\Phi }\overline{\eta })(z)2\mathbf{\hspace{0.17em}1}\{\overline{x}z<\overline{y}\},`$ (3.13)
which by (3.5) is the same as $`(\mathrm{\Phi }\eta )(z)<(\mathrm{\Phi }\overline{\eta }^{\overline{x},\overline{y}})(z)`$. In this way we get $`\eta ^{x,y}\overline{\eta }^{\overline{x},\overline{y}}`$ and (3.8) is proven. Display (3.9) is verified analogously.
By (3.3),
$`(\mathrm{\Phi }\overline{\eta })(z)(\mathrm{\Phi }\eta )(z)`$ $``$ $`2\mathbf{\hspace{0.17em}1}\{yz<\overline{y}\},`$ (3.14)
$`(\mathrm{\Phi }\overline{\eta })(z)(\mathrm{\Phi }\eta )(z)`$ $``$ $`2\mathbf{\hspace{0.17em}1}\{\overline{x}z<x\}.`$ (3.15)
Under the hypothesis of (3.10), this implies that
$`(\mathrm{\Phi }\overline{\eta })(z)(\mathrm{\Phi }\eta )(z)`$ $``$ $`2\mathbf{\hspace{0.17em}1}\{\mathrm{min}\{y,\overline{x}\}z<\mathrm{min}\{\overline{y},x\}\}.`$ (3.16)
Applying (3.5), we get (3.10).
## 4 Statement and proof of main result
Now we consider two processes $`\eta _t`$ and $`\overline{\eta }_t`$ with rates $`p`$ and $`\overline{p}`$, respectively, as discussed in the introduction. Our main result is:
###### Theorem 4.1
Suppose that whenever $`\eta \overline{\eta }`$ and $`\eta (x)=\overline{\eta }(\overline{x})=1`$, $`\eta (y)=\overline{\eta }(\overline{y})=0`$,
$`p(x,y;\eta )\overline{p}(\overline{x},\overline{y};\overline{\eta }),`$ if $`\overline{x}x<y\overline{y},`$ (4.1)
$`p(x,y;\eta )\overline{p}(\overline{x},\overline{y};\overline{\eta }),`$ if $`y\overline{y}<\overline{x}x.`$ (4.2)
Then if $`\overline{\eta }_t`$ restricted to $`𝒳`$ has a blocking invariant measure, so does $`\eta _t`$.
Theorem 1.1 is an immediate corollary of Theorem 4.1.
We construct simultaneously the two processes $`\eta _t`$ and $`\overline{\eta }_t`$ using the same marked Poisson processes $`((T_n^{i,j},U_n^{i,j}),(R_m^{i,j},V_m^{i,j}))`$. This joint construction is called coupling and is the key to the proof.
###### Lemma 4.2
Assume that $`\eta _t`$ and $`\overline{\eta }_t`$ are processes with rates $`p`$ and $`\overline{p}`$ satisfying (4.1)—(4.2). Under the coupling, if $`\eta _0\overline{\eta }_0`$ are both configurations of $`𝒳`$, then for all $`t0`$, $`\eta _t\overline{\eta }_t`$.
Proof. This is a mark-by-mark proof. Set $`\theta _0=0`$ and let $`\theta _1<\theta _2<\mathrm{}`$ be the instants at which there is a jump for at least one of the processes $`\eta _t,\overline{\eta }_t`$. Assume inductively that $`\eta _{\theta _{n1}}\overline{\eta }_{\theta _{n1}}`$, so that if $`(x_i,y_j)`$ and $`(\overline{x}_i,\overline{y}_j)`$ are the sites and holes of $`\eta _{\theta _{n1}}`$ and $`\overline{\eta }_{\theta _{n1}}`$, respectively, at time $`\theta _{n1}`$, then
$$x_i\overline{x}_i,\text{and}y_j\overline{y}_j,i,j0.$$
(4.3)
Let $`\tau _n`$ and $`\overline{\tau }_n`$ be the times defined as in (2.6) for the processes $`\eta _t`$ and $`\overline{\eta }_t`$, so that
$$\theta _n=\mathrm{min}\{\mathrm{min}\{\tau _k>\theta _{n1}\},\mathrm{min}\{\overline{\tau }_k>\theta _{n1}\}\}.$$
(4.4)
Let $`(I,J,K)`$ be the indices which realize the infimum (2.6) defining the time $`\theta _n`$, so that $`\theta _n\{T_K^{I,J},R_K^{I,J}\}`$. Let $`U\{U_K^{I,J},V_K^{I,J}\}`$ be the uniform random variable related with the indexes realizing the infimum, and let $`\sigma =\pm `$ indicate the direction of the jump at $`\theta _n`$: $`\sigma =+`$ if $`\theta _n=T_K^{I,J}`$ and $`U=U_K^{I,J}`$, $`\sigma =`$ if $`\theta _n=R_K^{I,J}`$ and $`U=V_K^{I,J}`$. Let
$`X=x_I,\overline{X}=\overline{x}_I,Y=y_J,\overline{Y}=\overline{y}_J;`$ (4.5)
$`\xi =\eta _{\theta _{n1}}\overline{\xi }=\overline{\eta }_{\theta _{n1}};`$ (4.6)
$`B=A_\sigma (\xi ,I,J),\overline{B}=A_\sigma (\overline{\xi },I,J).`$ (4.7)
Since (4.3) implies that $`\overline{X}X`$ and $`Y\overline{Y}`$, there are three possibilities:
1. $`\overline{X}X<Y\overline{Y}`$. By hypothesis (4.1), $`\overline{B}B`$. Hence there are two possibilities:
(a) $`U<\overline{B}B`$. In this case $`\eta _{\theta _n}=\xi ^{X,Y}`$ and $`\overline{\eta }_{\theta _n}=\overline{\xi }^{\overline{X},\overline{Y}}`$. By (3.8), $`\eta _{\theta _n}\overline{\eta }_{\theta _n}`$.
(b) $`\overline{B}U<B`$. In this case $`\eta _{\theta _n}=\xi ^{X,Y}`$ and $`\overline{\eta }_{\theta _n}=\overline{\xi }`$. By (3.6), $`\eta _{\theta _n}\overline{\eta }_{\theta _n}`$.
2. $`Y\overline{Y}<\overline{X}X`$. By hypothesis (4.2), $`B\overline{B}`$. Hence there are two possibilities:
(a) $`U<B\overline{B}`$. In this case $`\eta _{\theta _n}=\xi ^{X,Y}`$ and $`\overline{\eta }_{\theta _n}=\overline{\xi }^{\overline{X},\overline{Y}}`$. By (3.9), $`\eta _{\theta _n}\overline{\eta }_{\theta _n}`$.
(b) $`BU<\overline{B}`$. In this case $`\eta _{\theta _n}=\xi ^{X,Y}`$ and $`\overline{\eta }_{\theta _n}=\overline{\xi }`$. By (3.7), $`\eta _{\theta _n}\overline{\eta }_{\theta _n}`$.
3. $`X>Y`$ and $`\overline{X}<\overline{Y}`$. There are two possibilities:
(a) $`\sigma =+`$ and $`0=BU=U_K^{I,J}<\overline{B}`$. In this case $`\eta _{\theta _n}=\xi `$ and $`\overline{\eta }_{\theta _n}=\overline{\xi }^{\overline{X},\overline{Y}}`$. By (3.10), $`\eta _{\theta _n}\overline{\eta }_{\theta _n}`$.
(b) $`\sigma =`$ and $`0=\overline{B}U=V_K^{I,J}<B`$. In this case $`\eta _{\theta _n}=\xi ^{X,Y}`$ and $`\overline{\eta }_{\theta _n}=\overline{\xi }`$. Again by (3.10), $`\eta _{\theta _n}\overline{\eta }_{\theta _n}`$.
Notice that if we had used the same Poisson process for both forward and backward jumps then in the situation of case 3 above jumps could have occurred simultaneously in $`\eta `$ and $`\overline{\eta }`$, in opposite directions, which could destroy the ordering.
We remark that explosions are not excluded in Lemma 4.2.
Proof of Theorem 4.1: As remarked in the introduction, it suffices to show that if $`\overline{\eta }_t`$ has an invariant measure in $`𝒳_0`$, so does $`\eta _t`$. By restricting to a subset of $`𝒳^{}𝒳_0`$ (if necessary) we may assume that $`\overline{\eta }_t`$ is ergodic with invariant measure $`\overline{\mu }`$ having support $`𝒳^{}`$. This excludes explosions for the process $`\overline{\eta }_t`$ starting with configurations in $`𝒳^{}`$.
Start the coupled process with any two configurations $`\zeta \overline{\zeta }`$, with $`\overline{\zeta }𝒳^{}`$ and $`\zeta 𝒳`$. We know that:
1. $`\eta _t\overline{\eta }_t`$, by Lemma 4.2;
2. No explosions occur for $`\eta _t`$ (by an argument similar to the one in the proof of Lemma 4.2);
3. Since $`\overline{\eta }_t`$ is a continuous time ergodic Markov process in a countable state space, it converges in distribution to its unique invariant measure $`\overline{\mu }`$.
Hence any weak Cesaro-limit $`\mu `$ of the distribution of $`\eta _t`$ is coupled with $`\overline{\mu }`$ in such a way that, calling $`\nu `$ the coupled measure with marginals $`\mu `$ and $`\overline{\mu }`$, $`\nu `$ satisfies
$$\nu ((\eta ,\overline{\eta }):\eta \overline{\eta })=1.$$
(4.8)
This in particular implies that $`\mu (𝒳)=1`$. Since $`\mu `$ is a Cesaro-limit, $`\mu `$ is invariant for $`\eta _t`$. This implies the theorem.
## 5 Applications
To apply Theorem 4.1 one needs a suitable comparison process $`\overline{\eta }`$ which is known to have an invariant blocking measure. Obvious candidates are processes satisfying (1.10), for which the product measures (1.11) are invariant; in this section we draw some simple conclusions from this comparison. In the appendix we discuss briefly the existence of other possible comparison processes: those which satisfy detailed balance with respect to a Gibbs measure obtained from a suitable potential (Hamiltonian).
###### Theorem 5.1
Suppose that the exclusion process $`\eta _t`$ has simple, translation invariant rates $`p(x,y;\eta )=a(yx)\eta (x)(1\eta (y))`$ which for some $`\alpha `$ with $`0\alpha <1`$ satisfy
$$a(x)\alpha ^x\underset{0<yx}{inf}a(y)$$
(5.1)
for all $`x>0`$. Then $`\eta _t`$ has an invariant blocking measure.
Proof. The process with rates $`\overline{p}(x,y;\eta )=\overline{a}(yx)\eta (x)(1\eta (y))`$, where for $`x>0`$,
$$\overline{a}(x)=\underset{0<yx}{inf}a(y)\text{and}\overline{a}(x)=\alpha ^x\overline{a}(x),$$
(5.2)
has an invariant measure of the form (1.11). Thus the process $`\eta _t`$ has an invariant blocking measure by Theorem 1.1.
As a second example, consider a process with symmetric “disorder,” in which translation invariant, asymmetric, nearest neighbor rates are perturbed by arbitrary, bounded, symmetric nearest neighbor rates. Specifically, take $`p(x,y;\eta )=(c_0(x,y)+c_1(x,y))\eta (x)(1\eta (y))`$, where $`c_1(x,y)=c_2(x,y)=0`$ if $`|xy|>1`$ and
$$c_0(x,x+1)=K,c_0(x+1,x)=0,c_1(x,x+1)=c_1(x+1,x)=h(x),$$
(5.3)
with $`K>0`$ and $`h:_+`$ an arbitrary bounded function. It follows from Theorem 4.1 that this process has a blocking measure. A suitable comparison process has rates $`\overline{p}(x,y;\eta )=\overline{c}(x,y)\eta (x)(1\eta (y))`$ with $`\overline{c}(x,x+1)=c(x,x+1)`$, $`\overline{c}(x+1,x)=\alpha c(x,x+1)`$, and $`\overline{c}(x,y)=0`$ if $`|xy|>1`$, where $`\alpha =M/(M+K)`$ with $`M`$ an upper bound on $`h(x)`$; these rates satisfy (1.10) and hence have a blocking measure as given in (1.11). We single out this rather trivial example because in this case it is easy to see that the product measures with constant density are invariant measures, since if $`\mu `$ is such a measure then $`L_1^{}\mu =L_2^{}\mu =0`$ and hence $`L^{}\mu =0`$, where $`L_i^{}`$ is the adjoint of the generator for the process with rates $`c_i`$.
## Acknowledgments
JLL was supported in part by NSF Grant DMR-9813268. PAF and JLL thank DIMACS and its supporting agencies, the NSF under contract STC-91-19999 and the N. J. Commission on Science and Technology. This work started while PAF. was visiting Rutgers University with support of DIMACS. PAF was supported in part by FAPESP.
## Appendix
The remark that processes satisfying (1.10) have invariant product blocking measures of the form (1.11) can be generalized to processes which satisfy detailed balance with respect to a Gibbs measure obtained from a suitable potential (Hamiltonian). The latter is specified by a collection of real numbers $`\{J_R\}`$ indexed by finite subsets $`R`$ of $``$ and satisfying $`_{Rx}|J_R|<\mathrm{}`$ for each $`x`$. We show that if these coupling constants are chosen appropriately, then blocking Gibbs measures for this potential arise as the limit of finite volume measures.
Let $`T_N=[N+1,N]`$ and $`Y_N=\{0,1\}^{T_N}`$. For $`\eta Y_N`$ let $`\eta ^{}𝒳`$ be the configuration which agrees with $`\eta `$ in $`T_N`$ and with $`\eta ^H`$ outside $`T_N`$. The energy of the configuration $`\eta `$ is
$$H_N(\eta )=\underset{\{RRT_N\mathrm{}\}}{}J_R\chi _R(\eta ^{}),$$
(A.1)
where $`\chi _R(\zeta )=_{xR}(2\zeta (x)1)`$; the variables $`2\zeta (x)1`$ are spins which take values $`\pm 1`$. The corresponding finite-volume Gibbs measure $`\nu _N`$ on $`Y_N`$ is defined by
$$\nu _N(\{\eta \})=Z_N^1\mathrm{exp}\left(H_N(\eta )\right)$$
(A.2)
for $`\eta Y_N`$, with $`Z_N=_{\zeta Y_N}\mathrm{exp}(H(\zeta ))`$ a normalization constant; $`\nu _N`$ defines a measure on $`\{0,1\}^{}`$ by setting $`\nu _N(A)=\nu _N(\{\eta \{0,1\}^{}\eta ^{}A\}`$. Now let us assume for simplicity that all $`n`$-body terms in the potential, for $`n2`$, are translation invariant, i.e., that $`J_{R+k}=J_R`$ for $`k`$ and $`|R|2`$ (this assumption could easily be relaxed), and let $`K=_{Rx,|R|2}|J_R|`$.
###### Theorem A.1
Suppose that the one particle potential $`J_{\{x\}}`$ approaches $`\mathrm{}`$ as $`x`$ approaches $`\pm \mathrm{}`$, respectively, sufficiently fast that
$$\underset{x1}{}\mathrm{exp}\left(2J_{\{x\}}\right)<\mathrm{}\text{and}\underset{x0}{}\mathrm{exp}\left(2J_{\{x\}}\right)<\mathrm{}.$$
(A.3)
Then $`\nu =lim_N\mathrm{}\nu _N`$ exists and is a blocking measure. Moreover, if the rates $`p(x,y;\eta )`$ satisfy the detailed balance condition
$$p(x,y;\eta )e^{_{\{RxR\mathrm{or}yR\}}J_R\chi _R(\eta )}=p(y,x;\eta ^{x,y})e^{_{\{RxR\mathrm{or}yR\}}J_R\chi _R(\eta ^{x,y})}$$
(A.4)
then $`\nu `$ is reversible for the process with rates $`p`$.
Proof. We want to compare the measures $`\nu _N`$ and $`\nu _M`$, where $`N<M`$. For $`\eta Y_N`$ we let $`\eta ^{}Y_M`$ be the configuration which agrees with $`\eta `$ in $`T_N`$ and with $`\eta ^H`$ in $`T_MT_N`$, and for $`\zeta Y_M`$ we let $`\widehat{\zeta }Y_N`$ be the restriction of $`\zeta `$ to $`T_N`$; thus $`\widehat{\zeta }^{}(\widehat{\zeta })^{}T_M`$. Now fix $`\zeta Y_M`$, let $`S=\{x\zeta (x)\widehat{\zeta }^{}(x)\}`$, and set $`S_+=S\{x1\}`$, $`S_{}=S\{x0\}`$. Then
$`H_M(\zeta )`$ $`=`$ $`H_M(\widehat{\zeta }^{})2{\displaystyle \underset{xS_+}{}}J_{\{x\}}+2{\displaystyle \underset{xS_{}}{}}J_{\{x\}}+{\displaystyle \underset{\genfrac{}{}{0pt}{}{RT_M\mathrm{}}{|R|2}}{}}J_R[\chi _R(\zeta ^{})\chi _R(\widehat{\zeta }^{}{}_{}{}^{})]`$ (A.5)
$``$ $`H_M(\widehat{\zeta }^{})2{\displaystyle \underset{xS_+}{}}(J_{\{x\}}+K)+2{\displaystyle \underset{xS_{}}{}}(J_{\{x\}}K).`$
Thus if $`\eta Y_N`$,
$$e^{H_M(\eta ^{})}\underset{\widehat{\zeta }=\eta }{}e^{H_M(\zeta )}e^{H_M(\eta ^{})}\underset{x=N+1}{\overset{M}{}}\left(1+e^{2(J_{\{x\}}+K)}\right)\underset{x=N}{\overset{M1}{}}\left(1+e^{2(J_{\{x\}}+K)}\right).$$
(A.6)
Since the infinite products $`_{x1}(1+e^{2(J_{\{x\}}+K)})`$ and $`_{x0}(1+e^{2(J_{\{x\}}+K)})`$ converge by (A.3), we have for any $`ϵ>0`$,
$$e^{H_M(\eta ^{})}\underset{\widehat{\zeta }=\eta }{}e^{H_M(\zeta )}e^{H_M(\eta ^{})}(1+ϵ),$$
(A.7)
when $`N`$ is sufficiently large, uniformly in $`M`$.
Now suppose that $`A\{0,1\}^{}`$ is such that $`\mathrm{𝟏}A(\eta )`$ depends on $`\eta `$ only through the variables $`\eta (x)`$ for a finite number of sites—say for $`xT_L`$. Since for $`\eta Y_N`$, $`H_N(\eta )H_M(\eta ^{})`$ is independent of $`\eta `$,
$$\nu _N(A)=\frac{\underset{\eta Y_N,\eta ^{}A}{}e^{H_N(\eta )}}{_{\eta Y_N}e^{H_N(\eta )}}=\frac{\underset{\eta Y_N,\eta ^{}A}{}e^{H_M(\eta ^{})}}{_{\eta Y_N}e^{H_M(\eta ^{})}},$$
(A.8)
and with (A.7) this implies that if $`NL`$,
$$(1+ϵ)^1\nu _M(A)\nu _N(A)(1+ϵ)\nu _M(A).$$
(A.9)
Hence $`lim_N\mathrm{}\nu _N(A)`$ exists, so that $`\nu `$ exists. Similarly, if $`B\{0,1\}^{}`$ is the event that $`\eta (x)=\eta ^H(x)`$ for $`xT_N`$ then $`\nu _M(B)=Z_M^1_{\eta Y_N}\mathrm{exp}(H_M(\eta ^{}))(1+ϵ)^1`$ by (A.7), so that $`\nu `$ is a blocking measure.
The measure $`\nu `$ is reversible for the process with rates $`p`$ if for any continuous $`f`$ defined on $`\{0,1\}^{}`$ and any $`x,y`$,
$$p(x,y;\eta )[f(\eta ^{x,y})f(\eta )]𝑑\nu =0;$$
(A.10)
see the proof of the analogous result for stochastic Ising models in . But this integral may be calculated to arbitrary accuracy by replacing $`\nu `$ with $`\nu _N`$ for suitably large $`N`$ (here continuity of $`p`$ in $`\eta `$ is needed), and the fact that the integral with respect to $`\nu _N`$ vanishes is an immediate consequence of (A.4).
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# A Model for the 3He(𝑑⃗, p)4He Reaction at Intermediate Energies))footnote ) Dedicated to Prof. Shinsho Oryu on the occasion of his 60th birthday
## 1 Introduction
The measurement of the $`\stackrel{}{{}_{}{}^{3}He}`$ ($`\stackrel{}{d}`$, p)$`{}_{}{}^{4}He`$ reaction at RIKEN aimed at an investigation of the high momentum components of the deuteron wave function and the d-state admixture linked to them. High precision data resulted for the polarization observables $`A_y`$, $`A_{yy},A_{xx},C_{y,y}`$ and $`C_{x,x}`$. Out of them the linear combination $`C_{}=1+\frac{1}{4}(A_{yy}+A_{xx})+\frac{3}{4}(C_{y,y}+C_{x,x})`$ has been formed. The Dubna and Saturne groups also obtained the polarization correlation coefficient $`C_{}`$ built in this case from the measurements of $`T_{20}`$ and $`\kappa _0`$ in d + p backward scattering and from the inclusive deuteron breakup process. The polarization correlation coefficient $`C_{}`$ at forward angles of the outgoing proton is directly related to the ratio of deuteron wavefunction components if one uses the plane wave impulse approximation (PWIA) :
$`C_{}(PWIA){\displaystyle \frac{9}{4}}{\displaystyle \frac{w^2(k_{pn})}{u^2(k_{pn})+w^2(k_{pn})}}`$ (1.1)
Here $`u`$ and $`w`$ are the S-, D-wave components of the deuteron wavefunction, and $`k_{pn}`$ a kinematically fixed relative momentum of the $`pn`$ pair. These PWIA calculations are very poor in relation to the data. This is shown in table I for $`C_{}`$. There we also exhibit the different D-state probabilities for the modern realistic NN potentials, CD-Bonn, AV18 and Nijmegen I,II and 93. Clearly one needs a better calculation for the analysis of the $`{}_{}{}^{3}\stackrel{}{H}e`$($`\stackrel{}{d}`$,p)<sup>4</sup>He reaction.
A theoretical analysis has been reported by the SUT group based on a <sup>3</sup>He-n-p and d-d-p three-cluster model. However, the evaluations performed up to now in this model lead only to a tiny deviation from the PWIA calculations just mentioned. Recently, the Hosei group analyzed $`T_{20}`$ and $`\kappa _0`$ with the <sup>3</sup>He-n-p cluster model by an analogy between <sup>3</sup>He and the proton (T=1/2, S=1/2). They conclude that PWIA describes the global features of the experimental data.
In this letter we would like to introduce again a 3N model, which when evaluated correctly leads to a great similarity of various polarization observables to the ones found in the reaction $`{}_{}{}^{3}\stackrel{}{H}e`$($`\stackrel{}{d}`$,p)<sup>4</sup>He.
## 2 Model
For the <sup>3</sup>He($`\stackrel{}{d}`$,p)<sup>4</sup>He reaction we assume a model which is based on a three-body reaction process. This is shown in Fig. 1. The wavefunctions for <sup>3</sup>He and <sup>4</sup>He take on maximal values if the momenta of the subclusters are zero in their respective rest systems. These are for <sup>3</sup>He the momenta of p and d and for <sup>4</sup>He the momenta of the two deuterons. This means that for the moving nuclei the subcluster momenta should be equal. Therefore to form the $`\alpha `$-particle with highest probability in the picture of Fig. 1 one has to assume that the two deuterons, $`d^{}`$ and $`\stackrel{~}{d}`$, have equal momenta. Likewise for <sup>3</sup>He one has to assume that the proton and deuteron, $`\stackrel{~}{p}`$ and $`\stackrel{~}{d}`$, have equal momenta. This turns out to be kinematically inconsistent. Therefore we make a choice and assume that only the two deuterons forming the $`\alpha `$ particle have equal momenta. We justify this choice by the larger binding energy of the $`\alpha `$ particle.
It is easy to see that our basic assumption
$`\stackrel{}{k}_{\stackrel{~}{d}}=\stackrel{}{k}_d^{}`$ (2.1)
fixes the kinematics uniquely. It follows by simple kinematical arguments that
$`\stackrel{}{k}_{\stackrel{~}{p}}^{lab}={\displaystyle \frac{1}{2}}\stackrel{}{k}_p^{cm}{\displaystyle \frac{2}{5}}\stackrel{}{k}_d^{lab}=\stackrel{}{k}_{\stackrel{~}{d}}^{lab}`$ (2.2)
Here the superscripts $`lab`$ and $`cm`$ denote the laboratory and 5-body cm systems, respectively. Further the total momentum of the picked up proton and the incoming deuteron in the lab system is
$`\stackrel{}{K}={\displaystyle \frac{1}{2}}\stackrel{}{k}_p^{cm}+{\displaystyle \frac{3}{5}}\stackrel{}{k}_d^{lab}`$ (2.3)
Also we get the momentum of the picked up proton in the 3-body center of mass system (3CM) as
$`\stackrel{}{k}_{\stackrel{~}{p}}^{3CM}={\displaystyle \frac{1}{3}}\stackrel{}{k}_p^{cm}{\displaystyle \frac{3}{5}}\stackrel{}{k}_d^{lab}`$ (2.4)
and the 3CM energy as
$`E_{3CM}={\displaystyle \frac{3}{4m}}(\stackrel{}{k}_{\stackrel{~}{p}}^{3CM})^2`$ (2.5)
We show in Fig. 2 the relevant kinematics for the cm and the 3CM systems. From the relation
$`\stackrel{}{k}_{\stackrel{~}{p}}^{3CM}={\displaystyle \frac{2}{5}}\stackrel{}{k}_p^{3CM}{\displaystyle \frac{3}{5}}\stackrel{}{k}_d^{lab}`$ (2.6)
it follows under our condition, that the angles shown in Fig. 2 are related as
$`\theta ^{3CM}=\theta _p^{3CM}\theta _{\stackrel{~}{p}}^{3CM}`$ (2.7)
(note that $`\theta _p\theta _p^{3CM}=\theta _p^{cm}`$ ). The dependence of $`E_{3CM}`$ on $`\theta _p^{cm}`$ is illustrated in Fig. 4 for 3 deuteron energies. The scattering angle $`\theta ^{3CM}`$ is shown against $`\theta _p^{cm}`$ in Fig. 4 again for the same 3 deuteron energies.
Our claim is now that $`𝒪(E_d,\theta _p^{cm})𝒪_{pd}(E_{3CM},\theta ^{3CM})`$ where $`𝒪_{pd}`$ are the elastic pd deuteron polarization observables and $`𝒪`$ the ones for the reaction $`{}_{}{}^{3}\stackrel{}{H}e`$($`\stackrel{}{d}`$, p)<sup>4</sup>He.
Before calculating these 3N observables we introduce one more approximation. Looking at Fig.4 we see that $`E_{3CM}`$ varies with $`\theta _p^{cm}`$ and consequently for each $`\theta _p^{cm}`$ one would have to solve the 3N Faddeev equation. We avoided that for that qualitative investigation and have chosen available Faddeev results at three energies which lie in the three energy bands for $`0<\theta _p^{cm}<40^{}`$. They are $`E_{3CM}`$= 66.7, 100, 133 MeV corresponding to $`E_d`$= 140, 200, 270 MeV, respectively.
## 3 Results
As NN potential we used AV18 in the Faddeev calculations. The operator $`U`$ for elastic scattering has the form (see, for instance, )
$`U=PG_0^1+PT`$ (3.1)
where $`G_0`$, $`P`$ and $`T`$ are the free 3N propagator, permutation operators and a partial 3N break-up operator, which is determined by a Faddeev equation. The first term, the famous nucleon exchange term, is essentially related to the PWIA mentioned in introduction. In order to see the importance of solving the Faddeev equation correctly and not just replacing $`U`$ by $`PG_0^1`$ we compare the corresponding predictions for $`A_{yy}`$, $`A_{xx}`$ and $`A_{xz}`$ in Figs. 6-7. We see large differences especially above about 15 degrees. Trivially $`A_y`$ is identically zero using only the real term $`PG_o^1`$.
The predictions of the full Faddeev solution are shown in Figs. 9-11 at $`E_{3CM}`$=66.7, 100 and 133 MeV, respectively. This should be compared to recent data. We see a behavior qualitatively similar to those data, especially for $`A_y`$. For the $`A_y`$ data the minima shift to smaller $`\theta _p^{cm}`$ value with increasing energy like in Fig. 9. Also for $`A_{yy}`$ the qualitative behavior is similar in our model and the data, especially at the highest energy. For $`A_{xx}`$ the shapes are again very similar. In Fig. 9 and 11 we include one data point from . This shows that our absolute values are too high. For $`A_{xz}`$ shown in Fig. 11 there are not yet data.
## 4 Summary and Outlook
We assumed that the reaction $`{}_{}{}^{3}\stackrel{}{H}e`$($`\stackrel{}{d}`$,p)<sup>4</sup>He at forward angles is mainly driven by elastic pd scattering. In this model the deuteron picks up a proton from <sup>3</sup>He, scatters elastically and combines then again with the spectator nucleons to an $`\alpha `$ particle. Our main assumption is that the momentum of the scattered deuteron equals the spectator momentum of the deuteron in <sup>3</sup>He. This leads to a high probability to form the final $`\alpha `$-particle. The resulting spin-observables are in astonishingly good qualitative agreement with the data. Important thereby is, that the elastic pd amplitude is a full solution of the 3N Faddeev equation and not only a simple PWIA expression. This model should be generalized by the mechanism that also a neutron from <sup>3</sup>He can be picked up. In this case one has to use the nd break-up amplitude. Since the polarization of <sup>3</sup>He is carried by more than 90 % by the neutron this second mechanism is of course mandatory for a description of $`C_{x,x}`$ and $`C_{y,y}`$. The proton pick-up alone is too poor for those spin correlation observables. Also we neglected the momentum distributions of the proton in <sup>3</sup>He and of the deuteron in the $`\alpha `$ particle. As an additional improvement the spin of the deuteron should be properly rotated for the deuteron polarization observables.
Based on the promising qualitative results achieved it appears worthwhile to improve and enrich the model along the lines mentioned.
## Acknowledgements
This letter is dedicated to Prof. Shinsho Oryu on the occasion of his 60th’s birthday. Authors would like to thank Prof. Hideyuki Sakai, Dr. Tomohiro Uesaka, Mr. Yositeru Satou and Ms. Kimiko Sekiguchi for fruitful discussions in RIKEN. This work was supported by the Deutsche Forschungsgemeinschaft. The numerical calculations have been performed on the CRAY T90 of the John von Neumann Institute for Computing in Jülich, Germany.
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# THERMAL INSTABILITY AND THE FORMATION OF CLUMPY GAS CLOUDS
## 1 Introduction
The interstellar medium and cold gas clouds are characterized by a clumpy substructure and a turbulent velocity field (Larson 1981, Blitz 1993). As molecular clouds are the sites of star formation, their formation, internal structure and dynamics determines the rate of star formation and the properties of young stars, such as their mass function or binarity. The understanding of the origin of cold clouds and their internal substructure is therefore of fundamental importance for a consistent theory of star formation and galactic evolution.
In the nearby clouds, the dispersion velocity inferred from molecular line width is often larger than the gas sound speed inferred from the line transition temperatures (Solomon et al 1987). MHD turbulence may be responsible for the stirring of these clouds (Arons & Max 1975). This conjecture is supported by the polarization maps and direct measurements of field strength in some star forming regions (Myers & Goodman 1988, Crutcher et al. 1993). Recent simulations of MHD turbulence, however, suggest that it dissipates rapidly (Gammie & Ostriker 1996, MacLow et al 1998, MacLow 1999, Ostriker et al. 1999). One possible source of energy supply is winds and outflows from young stellar objects (Franco & Cox 1983, McKee 1989). But in regions where star formation is inactive, clumpy structure with velocity dispersion is also observed. Thus, the origin and energy supply of clumpy cloud structure remains an outstanding issue.
On small scales, magnetic field pressure is important in regulating infall and collapse of protostellar clouds and the formation of low-mass stars (Mouschovias & Spitzer 1976, Nakano 1979, Shu 1993). For clouds with sub-critical masses, gravitational contraction is proceeded by ambipolar diffusion which for typical cloud densities operates on a timescale $`\tau _B10^{78}`$ yr (Lizano & Shu 1989, Mouschovias 1991). In regions with intense star formation activities such as the central region of Orion, $`\tau _B`$ for individual dense clumps is comparable to the typical age of the young stellar objects. But, the spread in stellar ages ($`\mathrm{\Delta }\tau _{}10^6`$ yrs) appears to be considerably shorter than $`\tau _B`$ (Carpenter et al. 1997, Hillenbrand, 1997). This coeval star formation history requires either a coordinated trigger mechanism for star formation within initially magnetically supported clumps or subcritical collapse, fragmentation and star formation of a larger molecular cloud region in which the magnetic field plays a weak role.
A rapid and coordinated episode of star formation can also be inferred in globular clusters (Brown et al. 1991, 1995, Murray & Lin 1992, Lin & Murray 1992). In some metal deficient clusters such as M92, the total amount of heavy elements corresponds to the yield of a few supernovae. If star formation has proceeded over a duration $`\mathrm{\Delta }\tau _{}`$ comparable to the expected life span ($``$ a few $`10^6`$ yrs) of massive stars, a significant metallicity spread would be expected, in contrast to the observations (e.g. Kraft 1979). At least in these systems, $`\mathrm{\Delta }\tau _{}<\tau _B`$ and star formation may have proceeded through supercritical collapse. The dynamical timescale of most clusters at their half mass radius is $`\tau _d10^6`$ yr. Any energy dissipation associated with the episode of star formation would imply an even longer dynamical timescale in the proto cluster cloud prior to that event. We infer that $`\mathrm{\Delta }\tau _{}`$ was comparable to or shorter than the dynamical timescale of the proto cluster clouds indicating a rapid fragmentation and star formation episode.
In this paper, we focus on the rapid emergence of clumpy structure during the formation and collapse of a thermally unstable supercritical cloud. This process is relevant to the formation of stellar clusters as well as galaxies. We assume that the clouds condense out of a diffuse hot medium as a result of thermal instability. Large condensations with cooling timescales $`\tau _c>\tau _d`$ are thermally stable because they can adjust through contraction such that their radiative losses may be compensated by the release of their gravitational energy. Runaway cooling of the gas through thermal instability however occurs in clouds with $`\tau _c<\tau _d`$. In order to form clumps within an initially almost homogeneous cloud, internal density fluctuations must grow rapidly on a timescale short compared to the mean dynamical timescale of the entire cloud. Small scale density fluctuations would begin to dominate if either the growth timescale or the limiting amplitude is a decreasing function of the perturbations’ length scale. One possible fragmentation mechanism is gravitational collapse. The reduction in the cloud’s temperature reduces its Jeans’ mass, leading to the onset of gravitational instability and collapse. However, for a non rotating, cold, homogeneous gaseous region, gravitational instability alone cannot induce fragmentation because the growth rate is essentially independent of length scale such that the growth timescale for the density contrast is comparable to the dynamical timescale of the whole cloud (Hunter 1962). This has also been shown by numerical collapse simulations of initially gravitationally unstable perturbed gas clouds (e.g. Burkert & Bodenheimer 1993, 1996, Burkert, Bate & Bodenheimer 1997). If the initial density perturbations $`\delta _0`$ are linear ($`\delta _0<1`$), fragmentation is suppressed until the gas cloud has collapsed into either a disk or a dense filamentary substructure.
We propose that clumpyness in clouds arises naturally from their formation through a cooling instability which acts on timescales that can be much shorter than the dynamical timescale of the cloud. In a pioneering paper, Field (1965) derived a criterion for a cooling gas to be unstable to the growth of thermal condensations. He showed that thermal instability can lead to the rapid growth of density perturbations from infinitesimal $`\delta _0`$ to nonlinear amplitudes on a cooling timescale $`\tau _c`$ which for typical conditions in the interstellar medium is short compared to the dynamical timescale. If $`\tau _c`$ increases with decreasing density any small density difference would induce a temperature difference between the cooler perturbed region and the warmer background. Across the interface between the two-phase medium, differential cooling leads to a pressure gradient which induces a gas flow from the lower-density background towards the higher-density perturbed region. The density enhancement in the cooler region further reduces its cooling timescale compared to that of the background where $`\tau _c`$ increases. A more detailed investigation of the growth of condensations in cooling regions has been presented by Schwarz et al. (1972) who included also the effects of ionization and recombination and by Balbus (1986) who examined the effect of magnetic fields. The classical model of the interstellar medium where heating balances cooling was presented by Field et al. (1969). A recent progress report on the theory of thermal instability is given by Balbus (1995).
Although thermal instability proceeds faster than the collapse of the cloud, its growth rate is determined by the local cooling rate. During the initial linear evolution, variations in the initial over density (or under temperature) might lead only to a weak dependence of the growth timescale on the wavelength. In this paper we show however that there exist two important transitions which are very sensitively determined by the wavelengths of perturbations. 1) The growth of a perturbation is limited by its transition from isobaric to isochoric cooling, when the cooling time scale is reduced below the sound crossing time scale across the wavelength of the perturbation. This transition occurs at a lower temperature, with correspondingly larger over density, for perturbations with smaller wavelengths. 2) For those perturbation which can become nonlinear before the isobaric to isochoric transition, advection overtakes the pressure gradient in promoting the compression and growth of the perturbed region at an accelerated rate. The fluctuations which can first reach nonlinearity would dominate the growth of all perturbations with longer wavelengths and homogenize disturbances with smaller wavelengths. Thus, they determine the characteristic size and mass of the cold dense clumps which would emerge from the cooling of an initially nearly homogeneous cloud. Thermal conduction could in general erase these fluctuations, suppressing the instability. Weak, tangled magnetic fields would however be efficient enough in reducing the conductive flux, allowing the medium to break up into cold clumps on the characteristic length scale.
We study the cooling and fragmentation of gas using simplified power-law cooling functions. Since we are primarily interested in supercritical clouds, we neglect the effect of magnetic fields. Note that even a weak magnetic field could have an important destabilizing influence in thermal instability (Loewenstein 1990, Balbus 1995). In §2, we obtain approximate analytic solutions which describe the evolution of a linear density perturbation in the isobaric and nearly isochoric regime. We show that the growth of over density in a thermally unstable fluctuation is limited by a transition from isobaric to isochoric evolution and that the limiting amplitude is a decreasing function of the length scale. We verify our analytic approximations with numerical, hydrodynamical calculations which are also used in §3 to study the transition into the non-linear regime. In §4 we investigate the cooling of interacting perturbations and determine the critical length scale of clumps that emerge through thermal instability. The importance of thermal conduction is investigated in §5. In §6 we discuss the affect of heating processes and the formation of a stable 2-phase medium. Finally, we summarize our results and discuss their implications in §7.
## 2 The Initial Evolution of Thermal Instability
The dynamical evolution of the gas is described by the hydrodynamical equations
$$\frac{\rho }{t}+\underset{k=1}{\overset{3}{}}\frac{\rho U_k}{x_k}=0$$
(1)
$$\frac{U_j}{t}+\underset{k=1}{\overset{3}{}}U_k\frac{U_j}{x_k}+\frac{R_g}{\mu \rho }\frac{}{x_j}\left(\rho T\right)=0$$
(2)
$$\frac{T}{t}+\underset{k=1}{\overset{3}{}}U_k\frac{T}{x_k}+(\mathrm{\Gamma }1)T\underset{k=1}{\overset{3}{}}\frac{U_k}{x_k}=\frac{\rho \mathrm{\Lambda }}{C_v}$$
(3)
where j=1,2,3 is the coordinate index, $`C_v=R_g/\mu (\mathrm{\Gamma }1)`$ is the heat capacity, $`R_g,\mu `$ and $`\mathrm{\Gamma }`$ are the gas constant, mean molecular weight and adiabatic index, respectively.
In the unperturbed state, the gas remains at rest ($`U_j=0`$) and its density attains a constant value, $`\rho _0`$. The time dependent energy equation (3) gives
$$C_v\frac{T_0}{t}=\rho _0\mathrm{\Lambda }$$
(4)
where the cooling rate $`\mathrm{\Lambda }=\mathrm{\Lambda }_0T_0^\beta `$. The power index is determined by the detailed atomic processes. Since we are primarily interested in the physical evolution of thermal instability, we adopt a simple constant $`\beta `$ prescription. The cooling would be thermally unstable (with $`\tau _c=T_0/T_0/t`$ as an increasing function of $`T_0`$) in the isochoric region if $`\beta <1`$ and in the isobaric region if $`\beta <2`$. In the absence of external heating, the unperturbed gas temperature $`T_0`$ can be expressed as a function of the dimensionless time variable $`\tau t/\tau _c(0)`$ such that
$$T_0(t)=T_0(0)\left(1(1\beta )\tau \right)^{\frac{1}{1\beta }}$$
(5)
where $`T_0(0)`$ and $`\tau _c(0)C_v/\rho _0\mathrm{\Lambda }_0T_0(0)^{\beta 1}`$ are the initial (at $`t=0`$) temperature and cooling timescale, respectively.
### 2.1 The Perturbed Quantities
The evolution of the perturbed density ($`\rho _1=\rho \rho _0`$), temperature ($`T_1=TT_0`$), and velocities ($`U_j`$) are derived from the linearization of the equations (1) to (3):
$$\frac{}{t}\frac{\rho _1}{\rho _0}=\underset{j=1}{\overset{3}{}}\frac{U_j}{x_j},$$
(6)
$$\frac{U_j}{t}=\frac{R_gT_0}{\mu }\frac{}{x_j}\left(\frac{T_1}{T_0}+\frac{\rho _1}{\rho _0}\right),$$
(7)
and
$$\frac{}{t}\frac{T_1}{T_0}=(\mathrm{\Gamma }1)\underset{j=1}{\overset{3}{}}\frac{U_j}{x_j}\frac{1}{\tau _c}\left(\frac{\rho _1}{\rho _0}+(\beta 1)\frac{T_1}{T_0}\right)$$
(8)
where
$$\tau _c(t)\frac{T_0}{dT_0/dt}=\tau _c(0)(1\beta )t$$
is the characteristic cooling timescale at the instant of time t.
Since the perturbation equations are linear in $`x_j`$, we adopt a local approximation in which the positional dependence of all the perturbed quantities is proportional to exp(i$`k_jx_j`$) where $`k_j`$ is the wave number in the $`j`$th direction. Substituting a dimensionless velocity variable $`V_j=ik_j\tau _c(0)U_j`$, the perturbed equations reduce to
$$\frac{}{\tau }\frac{\rho _1}{\rho _0}=\underset{j=1}{\overset{3}{}}V_j,$$
(9)
$$\frac{V_j}{\tau }=K_j^2\left(1(1\beta )\tau \right)^{\frac{1}{1\beta }}\left(\frac{P_1}{P_0}\right),$$
(10)
where $`\frac{P_1}{P_0}=\frac{T_1}{T_0}+\frac{\rho _1}{\rho _0}`$ is the perturbed pressure, $`K_j\tau _c(0)k_j\sqrt{R_gT_0(0)/\mu }`$ is the ratio of the initial cooling to sound crossing timescale over a characteristic wavelength $`2\pi /k_j`$, and
$$\frac{}{\tau }\frac{T_1}{T_0}=(\mathrm{\Gamma }1)\underset{j=1}{\overset{3}{}}V_j\frac{\tau _c(0)}{\tau _c}\left(\frac{\rho _1}{\rho _0}+(\beta 1)\frac{T_1}{T_0}\right).$$
(11)
For a perfect gas, the unperturbed pressure $`P_0=R_g\rho _0T_0/\mu `$ decreases at the same rate everywhere. We find from Eqs (9) and (11) that the amplitude of the perturbed pressure is
$$\frac{}{\tau }\frac{P_1}{P_0}=\frac{}{\tau }\left(\frac{\rho _1}{\rho _0}+\frac{T_1}{T_0}\right)=\mathrm{\Gamma }\underset{j=1}{\overset{3}{}}V_j\frac{1}{1(1\beta )\tau }\left((2\beta )\frac{\rho _1}{\rho _0}(1\beta )\frac{P_1}{P_0}\right).$$
(12)
### 2.2 The initially isochoric regime with K $``$ 1
For computational simplicity, we now consider a 1-D limit treatment in which the initial (at $`\tau =0`$) amplitude of $`\rho _1`$ equals to a finite value $`\rho _a`$ with that of $`V_1`$ and $`P_1`$ equal to zero. These conditions correspond to an initially almost homogeneous, hot region of gas in pressure equilibrium. To third order in $`\tau `$ the eqs (9), (10), and (12) give the following approximate solution
$$\frac{\rho _1}{\rho _0}\frac{\rho _a}{\rho _0}\left(1+\frac{K^2}{6}\left(2\beta \right)\tau ^3\right)$$
(13)
$$V\frac{(\beta 2)K^2}{6}\frac{\rho _a}{\rho _0}\left(3\tau ^2+2\beta \tau ^3\right)$$
(14)
$$P(\beta 2)\frac{\rho _a}{\rho _0}\left(\tau ^2+\left(1\beta \right)\tau ^3+\left(\frac{4}{3}\frac{7}{3}\beta +\beta ^2\frac{\mathrm{\Gamma }K^2}{6}\right)\tau ^3\right)$$
(15)
Figure 1 compares this solution with a numerical integration of the complete non-linear hydrodynamical equations (1) to (3) for K=1 and K=0.5. We use a 1-dimensional version of the second-order Eulerian hydro code which is described in Burkert and Bodenheimer (1993). The agreement between the numerical results (solid lines) and analytical solution (dots) is excellent, even for large values of $`\tau 1`$ where the basis of the analytic approximation is no longer valid.
Due to slightly more efficient cooling within the density perturbation a small pressure gradient builds up. In an attempt to maintain pressure balance, the slightly warmer gas in the low-density regions continually compresses the more dense and cooler parts. Consequently the over-density in the perturbed region increases as the gas cools. Figure 1 and the equations (13) to (15) show that for small values of $`K1`$ the growth rate of the density enhancement depends on the size of the perturbation and increases with increasing values of K or decreasing wavelength. However, due to its long sound crossing timescale, the perturbation cannot be compressed significantly while cooling; it cools almost isochorically (Parker 1953). After one cooling timescale the gas temperature has reached its minimum value with the density enhancement still in the linear regime. Now, the pressure gradient reverses, erasing the fluctuation.
### 2.3 The initially isobaric regime: K $``$ 1
The solid lines in figure 2 show a numerical calculation of the evolution of a density perturbation with K = 200, $`\beta =0`$ and $`\mathrm{\Gamma }=5/3`$. For $`K>>1`$, perturbations can react quickly on any pressure gradients due to the short sound crossing timescale, relative to the cooling timescale. The simulations indicate a solution which consists of a fast oscillatory part and a slowly growing part. Linearizing the slowly growing part, we find from the equations (9) to (12) the following approximate solution:
$$\frac{\rho _1}{\rho _0}\frac{\rho _a}{\rho _0}\left(\frac{i\omega \mathrm{\Gamma }}{i\omega \mathrm{\Gamma }2+\beta }\right)\left(1+\frac{2\beta }{\mathrm{\Gamma }}\left(\tau +\frac{i}{\omega }e^{i\omega \tau }\right)\right),$$
(16)
$$V\frac{\beta 2}{\mathrm{\Gamma }}\frac{\rho _a}{\rho _0}\left(\frac{i\omega \mathrm{\Gamma }}{i\omega \mathrm{\Gamma }2+\beta }\right)\left(1+i\omega \tau e^{i\omega \tau }\right),$$
(17)
and
$$\frac{P_1}{P_0}\frac{(2\beta )i\omega }{\mathrm{\Gamma }K^2}\frac{\rho _a}{\rho _0}\left(\frac{i\omega \mathrm{\Gamma }}{i\omega \mathrm{\Gamma }2+\beta }\right)\left(1+(1\beta )\tau e^{i\omega \tau }\right).$$
(18)
The characteristic frequency is determined by a cubic dispersion relation
$$\omega ^3+i(1\beta )\omega ^2K^2\mathrm{\Gamma }\omega i(2\beta )K^2=0.$$
(19)
A similar relation was discussed by Balbus (1995). Note that the linearized solution is also valid for $`\tau >1/\omega `$ as long as $`\tau <<1`$. As $`K>>1`$ the two dominant real roots ($`\omega \pm \sqrt{\mathrm{\Gamma }}K`$) of the dispersion eq(19) yield oscillatory parts in $`\rho _1/\rho _0`$, $`V`$, and $`P_1/P_0`$. The upper panels of figure 2 show that for $`\tau <0.1`$ the analytical approximation is in good agreement with the numerical integration of the nonlinear hydrodynamical equations.
For all length scales, the ratio of sound propagation to cooling timescale
$$Q\tau _c(t)k\sqrt{R_gT_0(t)/\mu }=K(1(1\beta )\tau )^{\frac{32\beta }{22\beta }}$$
(20)
decreases during the subsequent evolution. Provided $`Q>>1`$, eq(18) implies that the magnitude of $`P_1/P_0`$ is much smaller than both $`V`$ and $`\rho _1/\rho _0`$. That is, the fluctuation reacts isobaric. Adopting $`P_1/P_00`$ and neglecting the oscillatory term, the equations (9) and (12) can be combined to
$$V=\frac{}{\tau }\left(\frac{\rho _1}{\rho _0}\right)=\frac{(2\beta )}{\mathrm{\Gamma }(1(1\beta )\tau )}\frac{\rho _1}{\rho _0}$$
(21)
with the solution
$$\frac{\rho _1}{\rho _0}=\frac{\rho _a}{\rho _0}\left(1(1\beta )\tau \right)^{\frac{2\beta }{(1\beta )\mathrm{\Gamma }}}$$
(22)
and
$$V=\frac{(2\beta )}{\mathrm{\Gamma }}\frac{\rho _a}{\rho _0}(1(1\beta )\tau )^{\frac{2\beta }{(1\beta )\mathrm{\Gamma }}1}.$$
(23)
In contrast to fluctuations with $`K<1`$ the evolution of isobaric fluctuations is independent of K. For $`(1\beta )^1>>\tau >\omega ^1`$, solutions for $`\rho _1/\rho _0`$ in Eqs (16) and (22) are in agreement to first order in $`\tau `$. The lower panels of figure 2 show that within a cooling time both $`\rho _1/\rho _0`$ and $`V`$ are amplified to very large values. The agreement between the numerical calculation and the analytical prediction (equation 22 and 23) is excellent. The opposite signs of $`\rho _1/\rho _0`$ and V confirm that mass is being pushed into the cool dense regions. For $`\tau 1`$, the numerically derived density enhancement falls below the predicted values as the fluctuation becomes isochoric and contributions from the perturbed pressure cannot be neglected anymore.
### 2.4 Transition to Isochoric Evolution and the Emergence of Small Scale Perturbations
Figure 3 shows the density evolution of initially isobaric fluctuations with different ratios of cooling to sound crossing times K as determined from the numerical calculations. The initially isobaric growth of the density fluctuations is independent of wavelength and K and in excellent agreement with equation (22) (dashed curve). The perturbations transform however to the isochoric solution for the epoch after $`Q`$ has declined below unity. Thereafter, gas in the perturbed region cools off faster than it can adjust to a pressure equilibrium with the surrounding region. Subsequently, the over density of the perturbed region is slowly modified by the inertial motion $`V_{trans}`$ of the gas at the time of the transition and its growth stalls. In Figure 3 the transition into the isochoric regime is indicated by the overdensity falling below the expected value shown by the dashed thick line.
Although the growth of the perturbed quantities does not explicitly depend on the wavelength $`k`$, the critical transition time when $`Q1`$
$$\tau _{trans}=\frac{1K^{\frac{2\beta 2}{32\beta }}}{1\beta }$$
(24)
is a function of $`K`$ (and $`k`$). At this transition point, the over density in the perturbed region is
$$\frac{\rho _{trans}}{\rho _0}=\frac{\rho _a}{\rho _0}K^{\frac{(42\beta )}{(32\beta )\mathrm{\Gamma }}}.$$
(25)
and the velocity is
$$V_{trans}=\frac{2\beta }{\mathrm{\Gamma }}\frac{\rho _a}{\rho _0}K^{(\frac{22\beta }{32\beta })(1+\frac{2\beta }{(1\beta )\mathrm{\Gamma }})}.$$
(26)
In the isochoric regime the amplitude of the perturbed density increases as
$$\frac{\rho _1}{\rho _0}=\frac{\rho _{trans}}{\rho _0}V_{trans}(\tau \tau _{trans})$$
(27)
which increases much less steeply than the isobaric fluctuations (equation 22).
For thermally unstable clouds, $`\beta <1`$ such that $`\rho _{trans}/\rho _0`$ is an increasing function of K or a decreasing function of the wavelength ($`\lambda `$) of the perturbations. Despite the independence of the rate of change of $`\rho _1/\rho _0`$ on $`\lambda `$, equation 24 shows that for $`\beta <1`$ the short length scale disturbances undergo isobaric to isochoric transition at a later time and therefore acquire a greater limiting amplitude than the long length scale disturbances. Thus, the short length scale disturbances would emerge to dominate the structure of the cloud unless the initial perturbation amplitude $`\rho _a/\rho _0`$ increases with $`K`$ more rapidly than $`K^{(2\beta 4)/(32\beta )\mathrm{\Gamma }}`$. This evolution is physically equivalent to the fragmentation process in which the contrast between the enhanced density in a disturbance and the average cloud density becomes most pronounced on the smallest scales.
## 3 The Transition into the Nonlinear Regime
In Fig. 3 fluctuations with very large values of K show yet another evolution: for later times the overdensity rises faster than predicted by equation (22). These fluctuations become nonlinear with $`\rho _1/\rho _0>1`$ before the transition into the isochoric regime. The critical value of K for this evolution can be estimated from equation (25) assuming $`\rho _{trans}/\rho _0=1`$:
$$K_{crit}=\left(\frac{\rho _a}{\rho _0}\right)^{\frac{(2\beta 3)\mathrm{\Gamma }}{42\beta }}$$
(28)
For $`K>K_{crit}`$ the analytical approximations discussed previously are not valid anymore and we have to investigate the evolution numerically, solving the complete non-linear hydrodynamical equations. The simulations shown in figure 3 assumed $`\beta =0`$, $`\mathrm{\Gamma }=5/3`$ and $`\rho _a/\rho _0=10^3`$. For these values the simple approximation (28) predicts $`K_{crit}=5600`$ which is roughly in agreement with the numerical results where the transition into the nonlinear regime occurs more smoothly between K=1000 and K=5000. Equation (28) somewhat overestimates $`K_{crit}`$ because nonlinear effects actually become important earlier, when the overdensity is in the range $`0.1<\rho _1/\rho _0<1`$.
Figure 4 shows the structure and evolution of a non-linear fluctuation. During the early isobaric evolution the pressure gradient (lower right panel) is negligible. A small pressure gradient builds up in the nonlinear regime, where the profiles cannot be approximated anymore by sinusoidal functions but instead become strongly peaked towards the center. Non-linear fluctuations grow fast with the density and temperature reaching their maximum and minimum values, respectively, at a time $`\tau _{crit}<1`$ which is shorter than a cooling time. Due to the fast growth in the nonlinear regime, $`\tau _{crit}`$ is roughly given by the time when $`\rho _1/\rho _0=1`$. From Eq (22), we find
$$\tau _{crit}=\frac{1}{1\beta }\left(1\left(\frac{\rho _a}{\rho _0}\right)^{\frac{(1\beta )\mathrm{\Gamma }}{2\beta }}\right).$$
(29)
At $`t=\tau _{crit}`$, the dimensionless velocity (see Eq. 23)
$$V_{crit}=V(\tau _{crit})=\frac{\beta 2}{\mathrm{\Gamma }}\left(\frac{\rho _a}{\rho _0}\right)^{\mathrm{\Gamma }\frac{1\beta }{2\beta }}$$
(30)
is much larger than unity for perturbations with small initial amplitudes such that contributions due to nonlinear advection (such as $`U_j\rho /x_j`$, $`U_jU_j/x_j`$, and $`U_jT/x_j`$) would exceed the linear contributions contained in the perturbed equations (6-8) before the over density $`\rho _1`$ has become comparable to $`\rho _0`$ (see above). Advection generally enhances the effect of compression and promotes the growth of density contrast at an accelerated rate.
Although the time of maximum compression for a fluctuation with $`K>K_{crit}`$ does not depend explicitly on the length scale, it is determined by the initial amplitude $`\rho _a/\rho _0`$ of the perturbation which may be a function of the wavelength. For an initial power-law perturbation in which $`\rho _a/\rho _0=A_0(k/k_0)^\eta `$,
$$\tau _{crit}=\frac{1}{1\beta }\left(1\left(A_0^{\frac{1}{\eta }}\frac{k}{k_0}\right)^{\frac{(1\beta )\mathrm{\Gamma }\eta }{2\beta }}\right).$$
(31)
If the amplitude of the initial perturbation is an increasing function of the wavelength (which corresponds to a negative $`\eta `$), $`\tau _{crit}`$ would be an increasing function of $`k`$ in the thermally unstable region with $`\beta <1`$. In this case, nonlinearity would be first reached on the largest length scale with $`K>K_{crit}`$. If, however, the amplitude of the initial perturbation is a decreasing function of the wavelength (i.e. $`\eta >0`$), $`\tau _{crit}`$ would be a decreasing function of $`k`$ and nonlinearity would be reached on the smallest scale first. Note that for $`1<\beta <2`$, the dependence of $`\tau _{crit}`$ on $`\eta `$ and $`k`$ is reversed.
In Figure 3, a temperature independent cooling function has been used. In order to determine the dependence on the specific form of the cooling function, additional simulations have been performed, adopting a more realistic cooling function (Dalgarno & McCray 1972) which assumes solar element abundance and collisional equilibrium ionization. Note that for temperatures T $`>10^4`$ K the cooling rate is several orders of magnitudes larger than for T $`<10^4`$ K, defining two different temperature regimes with very different cooling timescales. The simulations show that the previous results remain valid for each of these temperature regimes. Starting in the low-temperature regime, a fluctuation will become non-linear for K $`>`$ K<sub>crit</sub> and cool down to the lowest allowed temperature. The same is true for fluctuations that start in the high-temperature regime. Non-linear fluctuations in this regime do however stop cooling efficiently at T $`10^4`$ K, leading to high-density clumps with such a temperature.
## 4 Interacting Fluctuations and the Emergence of Substructure with a Critical Wave Length
Up to now we have investigated the evolution of isolated fluctuations. In reality however a cooling gaseous region consists of a superposition of fluctuations with different wavelengths and amplitudes. As we indicated in the previous section, the outcome of the thermal instability may be determined by the wavelength dependence in the initial amplitude of the perturbations.
In order to illustrate various competing effects such as isobaric to isochoric transition and the onset of nonlinear growth, we present in Figure 5 a series of models with $`\beta =0`$ and $`\mathrm{\Gamma }=5/3`$, where the initial density distribution consists of the superposition of two fluctuations with ratios of wavelengths $`\lambda _1/\lambda _2=20`$ and amplitude ratios $`\rho _{a,1}/\rho _{a,2}=2`$ which corresponds to $`\eta =0.23`$. Four values of $`\lambda _1`$ were chosen and they correspond to $`K_1`$=1,10,100 and 1000, respectively. The $`K_2`$ values for the smaller perturbation are always a factor 20 larger. Since the initial overdensity $`\rho _{a,1}/\rho _0=0.01`$, the critical value of K for nonlinear evolution is $`K_{crit}=316`$, according to equation (28). We show the density distribution after 1 $`\tau _c(0)`$. In the upper left panels of Fig. 5, the fluctuations have values of $`K_1=1`$ and $`K_2=20`$ which are small compared to $`K_{crit}`$. Their growth therefore stalls due to transition into the isochoric regime and the overdensity after a cooling time is still linear. The smaller fluctuation dominates at the end because its isochoric transition occurs later and at a higher overdensity than for the larger perturbation. In the upper right panels with $`K_1=10`$ and $`K_2=200`$ the smaller perturbation is again dominating after $`\tau =1`$ although, now, the density distribution is also affected by the underlying larger perturbation. In both cases the density within the density peaks does not decrease much with respect to its initial value. The situation is different in the lower left panel with $`K_1=100`$ and $`K_2=2000`$. Here the smaller perturbation has become nonlinear, generating small dense clumps which stand out against the larger perturbation. Up to now, the smaller perturbation was always dominating the density distribution after a cooling time. The situation is however different in the lower right panel where $`K_1>K_{crit}`$. Now, the larger perturbation becomes nonlinear and advection drives all the gas and its small scale fluctuations into one very dense, cold clump that is embedded in a hot diffuse environment, erasing smaller scale fluctuations.
The dependence of structure formation on the initial power-law perturbation index $`\eta `$ is illustrated in figure 6 which shows the initial and final density distribution of two interacting perturbations with $`\lambda _1/\lambda _2=10`$, $`K_1=2\times 10^4`$, $`\rho _{a,1}/\rho _{a,2}=10`$ ($`\eta =1`$) in the left panels and $`\rho _{a,1}/\rho _{a,2}=0.1`$ ($`\eta =1`$) in the right panels. As expected, in the case of $`\eta =1`$ the larger perturbation becomes nonlinear first, leading to one massive density peak after a cooling time. For $`\eta =1`$, the small length scale perturbations begin to dominate after a cooling time, breaking the region up into dense clumps on the smallest scale. More specifically, if the amplitude of the initial perturbation increases with increasing wavelength ($`\eta <0`$), clumps will form with length scales $`\lambda \lambda _{crit}`$. Otherwise, the sizes of the fastest growing perturbations will be $`\lambda \lambda _\kappa `$. In this case, the clump sizes should decrease with increasing magnetic field strength.
## 5 The importance of thermal conduction
During the growth of linear density perturbations in the isobaric regime the resulting temperature gradient will induce conductive heating of the fluctuations.
### 5.1 Thermal conduction in the absence of magnetic fields
Several studies (e.g. McKee & Begelman 1990, Ferrara & Shchekinov 1993) have demonstrated that thermal conduction could stabilize and even erase a density perturbation if its scale is smaller than the Field length (Field 1965)
$$\lambda _F=\left(\frac{\kappa T}{n^2\mathrm{\Lambda }}\right)^{1/2}$$
(32)
where $`\kappa `$ is the thermal conduction coefficient.
In order to include the effect of thermal conduction, the term $`\left(\kappa T\right)/\rho C_v`$ has to be added to the right-hand side of equation (3). The linearized pressure equation (12) is then
$$\frac{}{\tau }\frac{P_1}{P_0}=\mathrm{\Gamma }\underset{j=1}{\overset{3}{}}V_j\frac{1}{1(1\beta )\tau }\left((2\beta )\frac{\rho _1}{\rho _0}(1\beta )\frac{P_1}{P_0}\right)\frac{\kappa }{\rho _0C_v}\frac{T_1}{T_0}k^2\tau _c(0).$$
(33)
In the isobaric regime with $`P_1/P_00`$ and $`T_1/T_0\rho _1/\rho _0`$ thermal conduction will become important if
$$\frac{\kappa }{\rho _0C_v}k^2\tau _c(0)\frac{2\beta }{1(1\beta )\tau }$$
(34)
In the early stages of cooling ($`\tau <<1`$) fluctuations will therefore be erased by thermal conduction if their wavelenghts are
$$\lambda \lambda _\kappa =\frac{2\pi }{(2\beta )^{1/2}}\lambda _F.$$
(35)
Figure 7 shows the evolution of an initially isobaric, 1-dimensional fluctuation with a ratio of cooling- to sound crossing time K=200 and $`\beta =0`$. The 1-dimensional, non-linear hydrodynamical equations (1) - (3) are solved numerically, including thermal conduction. The solid line shows the evolution of the density contrast $`\rho _1/\rho _0`$ as predicted by the analytical model (equation 22) which is in excellent agreement with the numerical result (filled points) for $`\lambda >\lambda _\kappa `$. For $`\lambda =\lambda _\kappa `$ (upper dashed line) conductive heating is non-negligible anymore and the fluctuation grows less fast. For $`\lambda <\lambda _\kappa `$ the growth of fluctuations is suppressed by thermal conduction.
In summary, thermal conduction can play a significant role in regulating the break-up of a radiatively cooling gaseous medium. Small scale substructure can only emerge in a limited wavelength regime which is given by
$$\lambda _\kappa \lambda \lambda _{crit}$$
(36)
The growth of perturbations is completely suppressed if $`\lambda _{crit}<\lambda _\kappa `$ for all perturbations.
The dotted lines in figure 8a show $`\lambda _{crit}`$ for fluctuations with initial overdensities $`\mathrm{log}(\rho _a/\rho _0)=3,2,1`$, the solid line shows $`\lambda _\kappa `$. A cooling function assuming collisional equilibrium ionization (Spitzer 1978, Dalgarno & McCray 1972) and a realistic conduction coefficient (Ferrara & Shchekinov 1993) has been adopted. The dashed line shows the mean free path (Cowie & McKee 1977)
$$\lambda _e10^4\left(\frac{T}{K}\right)^2\left(\frac{cm^3}{n}\right)cm$$
(37)
for electron energy exchange. Note that for a given temperature the ratios $`\lambda _\kappa /\lambda _{crit}/\lambda _e`$ are independent of pressure.
The classical thermal conductivity is based on the assumption that $`\lambda _e`$ is short compared to the temperature scale height $`h_T\lambda /(T_1/T_0)>\lambda _{crit}`$. Otherwise, the heat flux q is saturated (Cowie & McKee 1977) and no longer equal to $`q=\kappa T`$. Indeed, figure 8a shows that linear fluctuations in astrophysical plasmas will in general lie in the non-saturated regime ($`\lambda _{crit}>\lambda _e`$) for initial amplitudes $`\rho _a/\rho _00.001`$.
However, we also find that in general $`\lambda _\kappa >\lambda _{crit}`$ for perturbations with amplitudes $`\mathrm{log}(\rho _a/\rho _0)2`$. This implies that a cooling instability, resulting from linear density perturbations will in general be suppressed by thermal conduction.
### 5.2 Thermal conduction, including magnetic fields
The interstellar medium is in general penetrated by magnetic fields. In most situations the electron mean free path $`\lambda _e`$ is large compared to the length scale at which the resistive destruction of the magnetic field is significant. A tangled magnetic field can then develop, concentrated on scales $`l_B`$ which are smaller than $`\lambda _e`$ (Chandran & Cowley 1998). When the gyroradius
$$a=\frac{v_Tm_ec}{eB}2.2\times 10^8\sqrt{\frac{T}{10^8K}}\left(\frac{\mu G}{B}\right)cm$$
(38)
of thermal electrons with typical velocities $`v_T=(kT/m_e)^{1/2}`$ is much smaller than $`l_B`$ or $`\lambda _e`$ the magnetic field controls the motion of individual electrons. This condition is satisfied in many astrophysical plasmas even if the magnetic field is too weak to be hydrodynamically important.
If $`a`$ is small compared to the length scale $`\lambda `$ of a fluctuation, heat is conducted according to the classcal thermal conduction equation (Spitzer 1962), however with a thermal conductivity $`\kappa _B`$ which is reduced from the classical Spitzer value $`\kappa `$ as a result of the tangled magnetic field by (Chandran & Cowley 1998)
$$\kappa _B\frac{0.1}{\mathrm{ln}(l_B/a)}\kappa .$$
(39)
If we normalize B to its value for magnetic-to-thermal energy equipartition $`B_T=(24\pi \rho R_gT)^{1/2}`$ we find
$$\mathrm{ln}\left(\frac{l_B}{a}\right)=3.1+\mathrm{ln}\left(\frac{B}{B_T}\right)+2\mathrm{ln}\left(\frac{T}{K}\right)0.5\mathrm{ln}\left(\frac{n}{cm^3}\right)+\mathrm{ln}\left(\frac{l_B}{\lambda _e}\right).$$
(40)
For weak magnetic fields ($`B0.01B_T`$) and length scales $`ł_B`$ of order the electron mean free path equation 40 leads to $`\mathrm{ln}(l_B/a)10`$, by this reducing $`\kappa _B`$ by two orders of magnitudes and the length scale of thermal conduction by one order of magnitude. As an example, the shaded area in figure 8b shows the wavelength regime $`\lambda _\kappa \lambda \lambda _{crit}`$ where linear fluctuations with $`\rho _a/\rho _0=10^2`$ could grow as a result of cooling, assuming a pressure of $`P/k_B=10^3`$ K cm<sup>-3</sup> and $`l_B=0.1\lambda _e`$. The presence of a weak magnetic field can suppress thermal conduction efficiently, allowing small scale structure with wavelengths $`\lambda 10^2`$ pc to emerge as a result of cooling.
## 6 The importance of heating and the emergence of a stable two-phase medium.
The calculations in sections 3 and 4 showed that cooling gas clouds with small initial perturbations break up on a critical wave length $`\lambda _{crit}`$ below which over density first becomes nonlinear. If the initial amplitude is a decreasing function of $`\lambda `$, the clouds would break up on the smallest length where the local radiative cooling law remains valid and fluctuations are not destroyed by conduction. But, if the initial amplitude is an increasing function of the wavelength,
$$\lambda _{crit}=\frac{2\pi \tau {}_{c}{}^{}(0)}{K_{crit}}\sqrt{\frac{R_gT_0(0)}{\mu }}.$$
(41)
small perturbations on scales $`\lambda <\lambda _{crit}`$ are erased when the gas accumulates in the center of fluctuations with $`\lambda =\lambda _{crit}`$. Perturbations with $`\lambda >\lambda _{crit}`$ do not become nonlinear but they break up into substructures with $`\lambda =\lambda _{crit}`$.
After a cooling time the dense, cold, non-linear perturbations are embedded in a warmer, diffuse environment (see Fig. 9). However, the example in figure 9 shows that the cooling timescale of the inter-clump gas remains short compared to the initial cooling timescale. This gas therefore cools to a ground state temperature $`T_{min}`$ shortly after $`\tau =1`$. Subsequently the reversed pressure gradient would remove the fluctuations unless they are gravitationally bound.
In order to maintain a stable two-phase medium a heating term must be included (Field et al. 1969). Here we assume a power law dependence of the heating rate
$$\mathrm{\Gamma }_h=\mathrm{\Gamma }_0\rho ^\gamma $$
(42)
where $`\mathrm{\Gamma }_0`$ and $`\gamma `$ are constants. In general, the size of the whole cooling region is large compared to $`\lambda _{crit}`$ such that it cannot establish pressure equilibrium with the surrounding confining medium during a cooling timescale. In this case, the region cools isochorically and breaks up into substructures on scales of $`\lambda _{crit}`$ before establishing pressure equilibrium with the environment. If the average gas density $`\rho _0`$ is smaller than a critical value
$$\rho _\mathrm{\Gamma }=\left(\frac{\mathrm{\Gamma }_0}{\mathrm{\Lambda }_0}T^\beta \right)^{\frac{1}{2\gamma }}$$
(43)
heating would dominate everywhere and the region would adjust to a thermal equilibrium state where heating is balanced against cooling. If $`\rho _0>\rho _\mathrm{\Gamma }`$ cooling dominates and the density fluctuations would grow and become non-linear as discussed in the previous sections. Eventually, after a cooling time, the region would break up into cold high-density condensations which are separated by warm gas with densities $`\rho _{min}<<\rho _0`$. If $`\rho _{min}>\rho _\mathrm{\Gamma }`$ this interclump medium would cool as shown in figure 9 and the density fluctuations would be erased. Figure 10 shows a situation with $`\rho _{min}<\rho _\mathrm{\Gamma }`$. Heating dominates in the interclump region where the gas temperature and gas pressure rise, until pressure equilibrium is established. A stable 2-phase medium has formed with cold clouds of minimum temperature embedded in a hot interclump medium with a temperature that is determined by the balance of cooling and heating.
## 7 Discussions
The discussions in this paper focussed on the emergence of small scale perturbations. We have assumed the pre-existence of small initial perturbations which is a reasonable assumption for dynamically evolving systems like the interstellar medium in galaxies or in galactic clusters. We have limited our analysis to the optically thin regime such that radiation transfer is solely due to optically thin local radiative processes. This approximation is appropriate for the collapse of supercritical clouds where the effect of a magnetic field is dominated by thermal processes. Such a situation may be particularly relevant for the formation of stellar clusters and first generation stars in galaxies. Provided that the density of the progenitor clouds is relatively small, the local cooling approximation is adequate. We also neglected the interaction and merging of clumps. These processes become important for the subsequent evolution and they will be considered in subsequent papers.
In the context of our approximations, we have shown that thermal instability can lead to the breakup of large clouds into cold, dense clumps with a characteristic length scale which is given by $`\lambda _{crit}`$ in eq. (41) or by the smallest unstable wavelength that is not erased by thermal conduction, depending on whether the amplitude of the initial perturbation is an increasing or decreasing function of wavelength. For linear perturbations with overdensities $`\rho _a/\rho _00.01`$ the critical wavelength lies in the regime of $`10^3`$ pc to $`10^1`$ pc, depending on the initial temperature. The emergence of small scale dense subcondensations is equivalent to fragmentation. As in a thermally unstable region the cooling timescale is shorter than the dynamical timescale, gravity has no time to play an important role during this fragmentation process. $`\lambda _{crit}`$ may be either smaller or bigger than the Jeans’ length. In the latter case gravity becomes important eventually. In general however, thermally induced fragmentation of clouds with small initial density fluctuations proceeds the onset of gravitational instability of their individual clumps.
In our analyses, we adopted an idealized power-law cooling function. In reality, the cooling efficiency would terminate when the main cooling agents reach their ground state or establish an equilibrium with some external heating source. The latter is necessary for the clouds to attain a two-phase medium. Interaction between these two phases may determine the pressure, density and infall rate of the cloud complex as well as the dynamical evolution and size distribution of cloudlets and sub condensations. The analysis of this interaction will be presented elsewhere.
We would like to thank the referee, Andrea Ferrara, for helpful suggestions and for pointing out the importance of thermal conduction. A. Burkert thanks the staff of UCO/Lick Observatory for the hospitality during his visits. We thank S.D. Murray for useful conversation. This research was supported by NASA through NAG5-3056, and an astrophysics theory program which supports a joint Center for Star Formation Studies at NASA-Ames Research Center, UC Berkeley, and UC Santa Cruz.
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# On the evolution of young radio-loud AGN
## 1 Introduction
### 1.1 Gigahertz Peaked Spectrum and Compact Symmetric Objects.
An important element in the investigation of the evolution of extra-galactic radio sources is the study of young counterparts of ‘old’ FRI/FRII extended objects. Two classes of compact radio source can be found in the literature as most likely representatives of this early evolutionary stage: I) Gigahertz Peaked Spectrum (GPS) sources, which are characterised by a convex-shaped radio spectrum peaking at about 1 GHz in frequency (O’Dea 1998), and II) Compact Symmetric Objects (CSO) which are characterised by their small size ($`<500`$ pc) and two-sided radio structure, e.g. having jets and lobes on both sides of a central core (Wilkinson et al. 1994). Clearly, samples of GPS sources and CSOs are selected in very different ways. GPS sources are selected on their broadband radio spectra, while CSOs are selected on their multi-frequency milli-arcsecond morphology. Therefore studies of these objects have mostly been presented separately. However, a significant overlap between the two classes of sources exists. GPS sources optically identified with galaxies are most likely to possess compact symmetric radio morphologies (Stanghellini et al. 1997a, 1999), and the large majority of CSOs exhibit a gigahertz-peaked spectrum. The large but not complete overlap between these two classes of source is most likely caused by the synchrotron self-absorbed mini-lobes, located at the extremities of most CSOs, being the main contributors to the overall radio spectrum, and producing the peak at about 1 GHz in frequency.
### 1.2 Evidence for the young nature of GPS sources and CSOs.
Since the early discovery of GPS sources, it has been speculated that these were young objects (Shklovsky 1965, Blake 1970). However, a commonly discussed alternative to them being young was that they are small due to confinement by a particularly dense and clumpy interstellar medium that impedes the outward propagation of the jets (van Breugel, Miley & Heckman 1984; O’Dea, Baum & Stanghellini 1991). This latter hypothesis now looks less likely since recent observations show that the surrounding media of peaked spectrum sources are not significantly different from large scale radio sources, and insufficiently dense to confine these sources. The most compelling evidence for youth is found in observations of the propagation velocities of the hot spots of several GPS/CSOs (Owsianik & Conway 1998; Owsianik, Conway & Polatidis 1998; Tschager et al. 1999). They all appear to have separation velocities of typically $`0.2h^1c`$, indicating a dynamical age of $`10^3`$ year and clearly showing that these are indeed young objects. Recent measurements of the high frequency breaks in Compact Steep Spectrum (CSS) sources, indicate that these somewhat larger objects have radiative ages in the range of $`10^3`$ to $`10^5`$ years (Murgia et al. 1999).
### 1.3 Current views on Radio Source Evolution
Observational constraints on the luminosity evolution of radio sources mainly come from the source density in the power - linear size $`(PD)`$ diagram (Shklovsky 1963). It was found that sources with large sizes ($`D>1`$ Mpc, eg. Schoenmakers 1999) and high radio luminosities ($`P>10^{26}`$ W/Hz at 178 MHz) are rare, suggesting that the luminosities of sources should decrease quickly with linear sizes approaching 1 Mpc. Several authors have compared the number densities of GPS and CSS sources with those of large radio sources to investigate the luminosity evolution of young radio sources (Fanti et al. 1995; Readhead et al. 1996; O’Dea & Baum 1997). Fanti et al. (1995) argue that the luminosities of CSS sources decrease by a factor of $`10`$ as they evolve to extended objects. Readhead et al. (1996) find a factor of 8 decrease in luminosity as a source expands from 500 pc to 200 kpc in overall size. Taking into account their CSO statistics, they find that the luminosity evolution in the phases CSO-MSO-LSO (MSO = Medium Symmetric Object; LSO = Large Symmetric Object), i.e. from 10 pc to 150 kpc, is consistent with a single power-law luminosity evolution. This conclusion is not supported by O’Dea & Baum (1997), who found that GPS and CSS sources must decrease in luminosity at a faster rate than the classical 3CR doubles. Blundell, Rawlings & Willott (1999) showed that any radio source evolution involving a decrease in luminosity with time would, at the highest redshifts, result in a bias towards young sources in flux density limited samples. Since this effect is only important at $`z>23`$, it is not relevant to the analysis presented in this paper.
Several GPS and CSOs (eg. 0108+388; Baum et al. 1990) exhibit low level, steep spectrum, extended emission on arcsecond scales, which seem to be relics of much older radio-activity. These objects are often classified as being intermittent or re-occurent, and therefore not as young objects. However, the components related to their gigahertz-peaked spectra and CSO morphologies are certainly young. The presence of faint relic emission only indicates that the active nucleus has been active before, and may constrain the typical timescale and frequency of such events. Based on the current knowledge of the formation of massive black-holes in the centers of galaxies, it is unlikely that the central engine itself is young, but just the radio source (Richstone et al. 1998).
It is unclear whether all young sources actually evolve to large extended objects. Some, or even the majority, may be short-lived phenomena due to a lack of significant fuel (Readhead 1994). The possible existence of these objects can have a large influence on the source statistics of young radio sources.
## 2 Samples of GPS and CSS sources
In this paper, we present a study of the evolution of young radio sources from the analysis of three samples of faint and bright GPS and CSS radio sources: The faint GPS sample from WENSS (Snellen et al. 1998a), the bright (Stanghellini et al. 1998) GPS sample, and the Fanti et al. (1990) CSS sample. As discussed in the companion paper (Snellen, Schilizzi & van Langevelde, 2000), we do not regard GPS quasars to be related to their galaxy counterparts; in the following we choose not to use the quasars for further analysis and only concentrate on the GPS and CSS galaxies. Unless stated otherwise, we will assume that all GPS galaxies evolve to large scale radio galaxies, and that all large scale radio galaxies were once GPS galaxies.
### 2.1 The faint GPS sample from WENSS
The selection of this sample has been described in detail in Snellen et al. (1998a). Candidate GPS sources were selected in two ways; i) with inverted spectra between 325 MHz and 625 MHz in WENSS (Rengelink et al. 1997), and ii) with inverted spectra between WENSS 625 MHz and Greenbank 5 GHz (Gregory & Condon 1991). The sources are located in two regions of the survey; one at $`15^h<\alpha <20^h`$ and $`58^{}<\delta <75^{}`$, which is called the mini-survey region (Rengelink et al. 1997), and the other at $`4^h00^m<\alpha <8^h30^m`$ and $`58^{}<\delta <75^{}`$. Additional observations at 1.4, 5, 8.4 and 15 GHz were carried out with the WSRT and the VLA, yielding a sample of 47 genuine GPS sources with peak frequencies ranging from 500 MHz to more than 15 GHz, and peak flux densities ranging from $`30`$ to $`900`$ mJy. This sample has been imaged in the optical and near-infrared, resulting in an identification fraction of $``$ 87 % (Snellen et al. 1998b, 1999). All the galaxies in the sample were used for the morphological evolution study. The redshifts of the majority of the objects had to be estimated from their optical magnitudes, using the well determined Hubble diagram for GPS galaxies (Snellen et al. 1996, O’Dea et al.1996). Some, assumed to be galaxies, have only faint lower limits to their magnitudes. For these a redshift of z=1.5 was assumed.
The overall angular sizes were measured from the VLBI observations as the maximum angular separation of components or the angular size for single component source (see companion paper) Their 5 GHz radio power was determined, assuming $`H_o=50`$ km sec<sup>-1</sup> Mpc<sup>-1</sup> and $`\mathrm{\Omega }_o=1`$. For a few sources, the rest-frame peak frequency was above 5 GHz. The radio power of these galaxies was corrected for the spectral turnover by extrapolating their optically thin spectrum to rest-frame 5 GHz. In the most extreme case (B0752+6355), this correction is $`<20\%`$. B0531+6121 was omitted from the sample since it does not have a genuine GPS spectrum.
For the luminosity evolution study to be discussed in section 4, it is crucial to have a good understanding of the selection effects. We therefore applied more strict constraints than in the original sample. Only the GPS sources which have inverted spectra between 325 MHz and 5 GHz, and with flux densities of $`>20`$ mJy at 325 MHz, 14 in total, were selected.
All the 26 galaxies in the sample are given in table 1. Column 1 gives the B1950 name, column 2 indicates whether the source is in the complete sub-sample or not, column 3 gives the (estimated) redshift, column 4 gives the observed peak frequency, column 5 the observed peak flux density, column 6 the rest-frame 5 GHz radio power, and column 7 the overall angular size.
### 2.2 The bright Stanghellini et al. GPS sample
A sample of radio bright GPS sources has been constructed by Stanghellini et al. (1998) from GPS candidates selected from the Kühr et al. (1981) 1 Jy catalogue, with declination $`>25^{}`$ and galactic latitude $`|b|>10^{}`$. Stanghellini et al. supplemented this data set with multi-frequency observations from the VLA, WSRT and data from the literature, and selected sources with a turnover frequency between 0.4 and 6 GHz, and an optical thin spectral index $`\alpha _{thin}<0.5`$ at high frequency. The final complete sample consists of 33 GPS sources, of which 19 are optically identified with galaxies. Four galaxies do not have a spectroscopic redshift. Their redshifts were estimated from their optical magnitudes, in the same way as for galaxies in the WENSS-sample. Their rest-frame radio power at 5 GHz has also been calculated in the same way as for the objects in the WENSS sample.
All the galaxies in the sample are given in table 2. Column 1 gives the B1950 name, column 2 the redshift, column 3 the observed peak frequency, column 4 the observed peak flux density, column 5 the rest-frame 5 GHz radio power, and column 6 the overall angular size. Column 7 gives the reference for the angular size.
### 2.3 The Fanti et al. CSS sample
The sample of CSS sources used in this paper is from Fanti et al. (1990). They constructed a sample without spectral bias by integrating the 3CR sample with sources from the Peacock and Wall sample (1982) which would be stronger than 10 Jy at 178 MHz, if corrected for low frequency absorption by extrapolation of the straight high-frequency part of the spectrum. All sources were included with projected linear size $`<15`$ kpc, (corrected) flux density at 178 MHz $`>`$10 Jy, and with Log $`P_{178}>`$ 26.5, in a well defined area of sky ($`|b|>10^{}`$, $`\delta >10^{}`$). A number of sources, which are included in the Stanghellini et al. sample are omitted from the Fanti et al. sample to avoid duplication. The remaining CSS galaxies are listed in table 3.
## 3 The spectral turnovers and the morphological evolution of young radio sources
Early measurements of the angular sizes of GPS and CSS sources using VLBI strongly suggested that their spectral turnovers are caused by synchrotron self absorption (SSA, Jones, O’Dell & Stein 1974; Hodges, Mutel & Phillips 1984; Mutel, Hodges & Phillips 1985) It was realised by Jones, O’Dell & Stein (1974) that if the optical depth due to SSA was less than unity at the spectral peak-frequency of these sources, lower magnetic fields, far from equipartition, would be present which should result in detectable self-compton radiation. Fanti et al. (1990) showed that there is a strong anti-correlation between the linear size and the turnover frequencies of CSS sources, as expected for SSA. However, more recently it was suggested by Bicknell, Dopita & O’Dea (1997) that such a correlation can also be explained by a particular model in which these sources undergo free-free absorption by ionised gas surrounding the lobes. In addition, Kuncic, Bicknell & Dopita (1998) argued that in addition to free-free absorption, induced Compton scattering will also have an important effect in forming the spectral peak. As a result, it opened up the debate again about the cause of the spectral turnovers in CSS and GPS sources.
The combination of bright and faint GPS and CSS samples as presented here, gives us a unique opportunity to carefully investigate the correlation between size and spectral peak. Not surprisingly, we confirm the anti-correlation between peak frequency $`\nu _p`$ and maximum angular size $`\theta _{max}`$ (see figure 1, top left panel). However, in addition we find a correlation between peak flux density $`S_p`$ and $`\theta _{max}`$. This is shown in the top right panel of figure 1. Note that only sources from the bright and faint GPS samples are with $`0.8<\nu _p<3`$ GHz are plotted here. This is necessary, since the peak flux densities are correlated with the peak frequencies, which would erroneously result in a correlation between peak flux density and angular size.
From SSA theory, it is expected that the angular size $`\theta `$ of a radio source is proportional to (Kellerman & Pauliny-Toth,1981):
$$\theta B^{1/4}S_p^{1/2}(1+z)^{1/4}\nu _p^{5/4}$$
(1)
where $`B`$ is the magnetic field strength and $`z`$ the redshift. Note that $`\theta `$ is only weakly dependent on both $`B`$ and $`z`$. Most remarkably, the strength and signs of the correlations between $`v_p`$, $`S_p`$ and $`\theta _{max}`$ as shown in figure 1 are exactly as expected from equation 1.
The overall angular size, $`\theta _{max}`$ (eg. the distance between the two mini-lobes), is used in the analysis above, but $`\theta `$ in equation 1 corresponds to the size of the radio components which are dominant at the peak-frequency (the mini-lobes). Therefore, these correlations have implications for the morphological evolution of these radio sources. The lower left panel of figure 1 shows the maximum angular size as function of $`S_p^{1/2}\nu _p^{5/4}`$. The solid line indicates the best linear fit. The dependence of this relation on redshift is proportional to $`(1+z)^{1/4}`$, which in any case is smaller than $`<20\%`$ and negligible for our z-range. Therefore the same relation is expected in the rest-frame of the objects. In the rest-frame, we can solve for the magnetic field $`B`$ by assuming equipartition. For this we use the equation derived by Scott & Readhead (1977) assuming an optically thin spectral index $`\alpha =1`$,
$$L=3.5\times (1(1+z)^{1/2})^{1/17}(1+z)^{1/2}S_p^{8/17}\nu _p^{33/34}$$
(2)
where $`L`$ is the equipartition component size. The projected linear size is shown as function of the equipartition component size in the lower right panel of figure 1. The dashed line indicates the dependence for which both quantities are the same. The solid line is the best linear least-squares fit, indicating a ratio of overall size to component size of $`56`$, throughout the samples of faint and bright GPS and CSS galaxies. This means that if GPS sources evolve into CSS sources, their ratio of component size to overall linear size remains constant, implying a self-similar evolution.
Note that the main difference between the lower left and right panels of figure 1 is that in the first a constant magnetic field is assumed, and in the second an equipartition magnetic field. It appears that the first correlation is slightly flatter than expected for self-similar evolution. Indeed the ratio of the component to overall angular size is on average a factor 2 smaller for the CSS sample than for the GPS samples, while these ratios are virtually the same assuming an equipartition magnetic field. This may indicate that young radio sources stay in equipartition while evolving in a self-similar way. This would require that the magnetic fields in CSS sources are typically a factor $`20`$ lower than in GPS sources. The linear correlation itself only indicates a constant ratio between the magnetic field and particle energies. This constant does not have to be equal to unity, as required for equipartition. However, for a ratio of unity in energies, the bottom right panel of Figure 1 requires a ratio of overall to component size of typically $`56`$, which is close to the result seen in VLBI observations. This means that the energy ratio is not only constant, but also close to unity, which indicates that equipartition probably holds.
Figure 1 demonstrates that the data is consistent with a combination of SSA, equipartition, and self-similar growth. It is not obvious that the same correlation should apply for free-free absorption. Although other more complicated combinations of mechanisms such as free-free absorption with induced Compton scattering (Kuncic, Bicknell & Dopita, 1998) may also fit the data, the simplest explanation by far is to assume that SSA, equipartition and self-similar source growth all individually hold. We therefore believe that SSA is indeed the cause of the spectral turnovers in GPS and CSS sources.
It may not be surprising that young radio sources evolve in a self-similar way. Leahy and Williams (1984) showed that the cocoons of FRII sources of very different physical size had similar axial ratios. More recently, Subrahmanyan, Saripalli & Hunstead (1996) found very similar ratios for sources of linear sizes above 900 Kpc, also suggesting that radio sources evolve in a self-similar way. An analytical model for radio sources with pressure confined jets developed by Kaiser & Alexander (1997) shows that the properties of the bow shock and of the surrounding gas force the sources to grow in a self-similar way, provided that the density of the surrounding gas falls off less steeply than $`1/r^2`$.
## 4 The luminosity evolution of young radio sources
The number count statistics and linear size distributions used in studies to constrain the luminosity evolution of radio sources, have all been averaged over a wide redshift range and only include the brightest objects in the sky (Fanti et al. 1995, Readhead et al. 1996, O’Dea & Baum 1997). However, in flux density limited samples, the redshift distribution of GPS galaxies is significantly different from that of large size radio galaxies (see figure 2). This suggests that the interpretation of the number count statistics is not straightforward. Note that given the expected luminosity evolution as sources evolve in size, many of the present day GPS sources will have FRI luminosities. It is therefore assumed that GPS galaxies evolve into both FRI and FRII sources.
The bias of GPS galaxies towards higher redshifts than large size radio galaxies itself provides an important clue about the luminosity evolution of radio sources. It implies that GPS galaxies are more likely to have higher radio power than extended objects in flux density limited samples. If GPS and large size radio sources are identical objects, observed at different ages, their cosmological density evolution, for example their birth rate as function of redshift, should be the same. Since their lifetimes are short compared to the Hubble time, the redshift distributions of the GPS galaxies, and the objects they evolve to, should also be the same. The bias of GPS sources towards higher redshifts and radio powers therefore implies that their luminosity function must be flatter than that of large size radio sources. We argue that the luminosity evolution of the individual objects strongly influences their collective luminosity function, and propose an evolution scenario in which GPS sources increase in luminosity and large size sources decrease in luminosity with time (see section 4.1). In the simplified case, in which source to source variations in the surrounding medium can be ignored, the luminosity of a radio source depends only on its age and jet power. Consider first the luminosity function of large size sources. It is expected that large size sources decrease in luminosity with age (see section 4.1). Therefore high luminosity sources will tend to be biased towards objects with both small ages and high jet powers. The intrinsic space density for high power jet sources will of course tend to be small. Furthermore, for a given jet power there are fewer young sources than old sources, simply because sources spend only a small fraction of their time being young. The result is a very low space density of large size sources of high power. In contrast, large size sources of low power are biased to be both old and with low jet power, both common conditions, hence the space density of large size sources with low power is much higher than that for those with high power, and the luminosity function for large size sources is steep. In contrast, the luminosity of GPS sources is expected to increase with sources age (see section 4.1). High luminosity GPS sources are therefore biased to be old and of high jet power, while low luminosity objects are biased to be young and of low jet power. Instead of reinforcing each other as in the case for large size sources, for GPS sources the age and jet power space density biases partly counteract. The result is a much less difference in the space density of low and high power GPS sources and hence a much flatter luminosity function for GPS sources.
In the next section we will show that the luminosity evolution as proposed is expected for a ram-pressure confined, self similarly evolving radio source in a surrounding medium with a King-profile density. In the inner parts of the King profile, the density of the medium is constant and the radio source builds up its luminosity (eg. Baldwin 1982), but after it grows large enough the density of the medium declines and the luminosity of the radio source decreases.
In section 4.2 we will show how the luminosity evolution of the individual sources modifies the luminosity function, and in section 4.3, the local luminosity function of GPS sources is constructed and compared with that of large size radio sources.
### 4.1 A self-similar evolution model
An important parameter in evolution models of radio sources is the density profile of the surrounding medium. In general, X-ray observations of nearby ellipticals have shown that their ISM are well fitted by a King profile distribution (Trinchieri et al. 1986):
$$\rho (r)=\rho _0\left[1+\left(\frac{r}{r_c}\right)^2\right]^{\beta /2}$$
(3)
where $`\rho `$ is the density of the medium as function of distance to the centre of the host galaxy $`r`$, $`r_c`$ is the core radius, and $`\beta `$ the slope parameter. Typical core radii in giant ellipticals are observed to be $`r_c=5001000`$ pc (Trinchieri et al. 1986), For simplicity, we treat the two regimes separately: 1) The GPS phase at $`r<r_c`$ where the density of the medium, $`\rho _{ism}`$, is assumed to be constant. 2) The large size (LS) phase at $`r>r_c`$ where $`\rho _{ism}r^\beta `$.
If the thrust of the radio jet is balanced by the ram-pressure of the surrounding medium, the growth of the radio source is equal to
$$dr/dt\left(\frac{P_J}{\rho _{ism}(r)A}\right)^{1/2}$$
(4)
where $`dr/dt`$ is the propagation velocity of the hot-spots, $`P_J`$ is the jet power, and $`A`$ is the cross-sectional area (Begelman 1996). In the previous section we showed that young radio sources seem to evolve in a self-similar way. Since it is in close agreement with the theoretical work of Kaiser & Alexander (1997), we will assume self-similar evolution, with $`Ar^2`$. Note however, that in the work by Kaiser & Alexander (1997), the cross-sectional area of the jet grows slightly more slowly with size, which is therefore not completely self-similar, but allowing the expansion of the bow-shock and cocoon to be fully self-similar (Kaiser 2000, private communications). Here, by assuming $`Ar^2`$, this will not be the case, but differences are small and for simplicity’s sake we use this anyway.
From integrating eq. 4 it follows that in the GPS phase a source grows in linear size with time as $`t^{1/2}`$, assuming that the jet-power is constant with time. The average internal density of the radio source, $`\rho _i`$, is proportional to $`P_Jt/V`$, where $`V`$ is the volume of the radio source which is proportional to $`r^3`$. Hence, $`\rho _ir^1`$, indicating that the radio emitting plasma expands proportionally to its linear size, $`r`$, and that expansion losses have to be taken into account. If the energy spectrum of the electrons is $`n(E)=n_oE^\gamma `$, $`n_o`$ varies proportionally to $`r^{4/3}`$ for $`\gamma =2`$ ($`\alpha =0.5`$, Moffet, 1977). We will assume that the radio source is in equipartition, so that $`n_oB^2`$, where $`B`$ is the magnetic field. The radio power $`L_\nu `$ at a particular frequency in the optically thin part of the spectrum scales as
$$L_\nu n_o^{7/4}VP_J^{7/8}r^{2/3}$$
(5)
for $`\gamma =2`$. Hence radio sources increase in luminosity in the GPS phase.
In the LS phase, radio sources grow as $`t^{2/(4\beta )}`$, and the density of the radio emitting plasma varies as $`\rho _ir^{\frac{\beta +2}{2}}`$, Taking into account expansion losses in a similar way as for the GPS phase, this means that $`n_or^{4\frac{\beta +2}{6}}`$, and under equipartition conditions,
$$L_\nu P_J^{7/8}r^{\frac{2}{3}\frac{7}{6}\beta }$$
(6)
Hence radio sources in the LS phase decrease in radio luminosity. Note that we do not take synchrotron losses and losses due to scattering of the CMB into account, which may influence the LS phase and cause sources to decrease faster in luminosity with time. The schematic evolution in radio power of a radio source according to this model is shown in figure 3.
It is interesting to determine what the expected evolution in peak frequency and peak flux density is for a radio source in the GPS phase and the LS phase. A source will become optically thick at a frequency where $`\kappa _\nu l1`$, where $`\kappa _\nu `$ is the absorption coefficient and $`l`$ the pathlength through the radio plasma. For synchrotron self absorption, assuming equipartition and self-similar evolution, this means that,
$$\kappa _\nu ln_oB^2\nu _p^3rn_o^2\nu _p^3r=1$$
(7)
for $`\gamma =2`$ (Moffet, 1977). The optically thin radio power, as determined above, will be frequency dependent and proportional to $`P_\nu S_p\nu _p^{1/2}`$. Therefore in the GPS phase,
$$\nu _pr^{5/9},S_pr^{17/18},S_p\nu _p^{17/10}$$
(8)
In the LS phase, assuming $`\beta =1.5`$, which is a typical value based on observations of X-ray halos (Trinchieri et al. 1986),
$$\nu _pr^{11/9},S_p\nu _p^{17/44}$$
(9)
From figure 1 we can estimate the transition between the two phases (500-1000 pc) to occur at $`\nu _p100500`$ MHz The evolutionary tracks are shown in figure 4.
### 4.2 Luminosity evolution and the luminosity function.
In this section the influence of the luminosity evolution of the individual objects on the slope of their collective luminosity function is derived. We will ignore source to source variations in the surrounding medium and use the radio-size dependent luminosity evolution as derived in the previous section.
Suppose that the comoving number density of sources with a jet power $`N(P_J)`$ is a power-law distribution
$$N(P_J)P_J^\delta $$
(10)
between $`P_{}`$ and $`P_+`$, and the sources have a flat distribution of ages below a certain maximum age, then the source density as function of age and jet power is represented by the grey scales in figure 5. The radio power of a source, $`L_\nu `$, can be parameterised as
$$L_\nu P_J^\kappa r^ϵ$$
(11)
where $`ϵ=2/3`$ and $`ϵ=13/12`$ in the GPS phase and LS phase respectively, and $`\kappa =7/8`$, as derived in the previous section. A line-integral over a solid line in figure 5 gives the total number of sources in the volume with a particular luminosity. It can be seen that in the GPS phase, these lines are approximately perpendicular to the density gradient, indicating that a change in luminosity results in only a small change in the number of sources. In the LS phase, they are parallel to the density gradient, and a change in luminosity results in a large change in the number of objects. The luminosity function $`N(L)`$ can be derived from,
$$N(L_\nu )=\frac{\delta }{\delta L_{}}\underset{L_\nu (p,r)<L_{}}{}N(p,r)𝑑p𝑑r$$
(12)
where $`p`$ is the jet power and $`r`$ is the size of the radio source. As can be seen in figure 5, the integration limits of this equation are a different function of age and jet power, depending on the luminosity . The equation should be solved separately for a high and low luminosity regime, in both the LS and the GPS phase. Since the border between the GPS and LS phase is at a constant source size, $`r_{}`$, it is better to integrate over the source size $`r`$ than over the source age $`t`$. For the low luminosity regime in the GPS phase,
$$N(L_\nu <L_{})=\underset{P_{}}{\overset{P_+}{}}\underset{0}{\overset{\frac{L_{}^{1/ϵ}}{p^{\kappa /ϵ}}}{}}p^{\delta \frac{1}{2}}r𝑑r𝑑p=L_{}^{\frac{2}{ϵ}}\underset{P_{}}{\overset{P_+}{}}p^{\delta \frac{1}{2}\frac{2\kappa }{ϵ}}𝑑p$$
(13)
where $`N(L_\nu <L_{})`$ are the total number of sources below a particular luminosity $`L_{}`$, and therefore integrating over jet power $`p`$ and age $`t`$ gives,
$$N(L_\nu )L_\nu ^{2/ϵ1}$$
(14)
For the high radio power regime in the GPS phase,
$$N(L_\nu >L_{})=\underset{\left(\frac{L_{}}{P_+^\kappa }\right)^{\frac{1}{ϵ}}}{\overset{r_{}}{}}r\underset{\left(\frac{L_{}}{r^ϵ}\right)^{\frac{1}{\kappa }}}{\overset{P_+}{}}p^{\delta \frac{1}{2}}𝑑p𝑑rL_{}^{(\delta +\frac{1}{2})\frac{1}{\kappa }}$$
(15)
with a sharp cut-off near $`L_\nu =P_+^\kappa r_{}^ϵ`$. In this regime of radio power,
$$N(L_\nu )L_\nu ^{(\delta +\frac{1}{2})\frac{1}{\kappa }1}$$
(16)
For the low luminosity regime in the LS phase,
$$N(L_\nu <L_{})=\underset{p_{}}{\overset{\left(\frac{L_{}}{r_+^ϵ}\right)^{\frac{1}{\kappa }}}{}}p^{\delta \frac{1}{2}}\underset{\left(\frac{L_{}}{p^\kappa }\right)^{\frac{1}{ϵ}}}{\overset{r_+}{}}r^{\frac{1}{2}}𝑑r𝑑pL_{}^{(\delta +\frac{1}{2})\frac{1}{\kappa }}$$
(17)
with a sharp cut-off near $`L_\nu =P_{}^\kappa r_+^ϵ`$. In this regime of radio power,
$$N(L_\nu )L_\nu ^{(\delta +\frac{1}{2})\frac{1}{\kappa }1}$$
(18)
For the high radio power regime in the LS phase,
$$N(L_\nu >L_{})=\underset{r_{}}{\overset{\left(\frac{L_{}}{P_+^\kappa }\right)^{\frac{1}{ϵ}}}{}}r^{\frac{1}{2}}\underset{\left(\frac{L_{}}{r^ϵ}\right)^{\frac{1}{\kappa }}}{\overset{P_+}{}}p^{\delta \frac{1}{2}}𝑑p𝑑rL_{}^{3/2ϵ}$$
(19)
with a cut-off near $`L_\nu =P_+^\kappa r_{}^ϵ`$. In this radio power regime,
$$N(L_\nu )L_\nu ^{3/2ϵ1}$$
(20)
As can be seen from equations 15 and 18, the slope of the luminosity function is expected to be the same in the high luminosity and low luminosity regimes for the GPS and LS phases respectively, since $`\delta `$ and $`\kappa `$ are independent of the age of the radio source. The low luminosity regime and the high luminosity regime of the GPS phase and the LS phase are expected to have a slope of $`+2`$ and $`2.4`$ respectively for the proposed evolution model.
### 4.3 The Local Luminosity Function of GPS sources.
As is shown in the previous section, the comparison of the local luminosity function (LLF) of young and old radio sources can put strong constraints on the rise and decay of their radio luminosity. One would like to compare the LLF of GPS sources with the model derived in section 4.1 & 4.2 directly. This is not possible due to the lack of local GPS sources in present samples (the low local number density of GPS sources catalysed this discussion in the first place). For example, only 2 GPS galaxies in the Stanghellini et al. sample are at $`z<0.2`$. However, since we assume that GPS sources evolve into large size sources and their lifetimes are short compared to cosmological timescales, their birth rate as function of redshift should be the same. Therefore the cosmological evolution as determined for large scale radio sources can be used to describe the cosmological evolution for GPS sources. In this way, the GPS LLF can be estimated using the GPS galaxies at all redshifts, which will be attempted in this section. This estimated GPS LLF will then be compared with what is expected from the model, as derived in section 4.2.
The LLF for powerful radio sources and its cosmological evolution, has been studied by Dunlop & Peacock (1990). We will use the pure luminosity evolution model, since it fits the available redshift and source-count data well, and it is relatively straightforward to implement. In this particular model, the overall shape of the luminosity function does not change with cosmological epoch, only the normalisation in luminosity (see fig. 6). Dunlop and Peacock (1990) parameterise an evolving two-power-law luminosity function as
$$\rho (P_\nu ,z)=\rho _o\left\{\left(\frac{P}{P_c(z)}\right)^a+\left(\frac{P}{P_c(z)}\right)^b\right\}^1$$
(21)
where $`a`$ and $`b`$ are the two power-law slopes, $`P_c(z)`$ is the evolving ‘break’ luminosity, and $`\rho _o`$ is determined by normalisation at z=0. The redshift dependence, $`P_c(z)`$, was parameterised by Dunlop & Peacock as
$$\mathrm{log}P_c(z)=a_0+a_1z+a_2z^2$$
(22)
The best-fit model parameters for pure luminosity evolution ($`\mathrm{\Omega }_0=1`$) are, $`\rho _o=6.91`$, $`a=0.69`$, $`b=2.17`$, $`a_0=25.99`$ (in W/Hz), $`a_1=1.26`$, $`a_2=0.26`$. Since Dunlop & Peacock (1990) did their analysis at 2.7 GHz, their radio powers have to be transformed to 5 GHz. Assuming a mean spectral index of $`0.75`$, we use a conversion factor of $`0.20`$ in the logarithm. This luminosity evolution parameterisation, as shown in figure 6, is used to derive the LLF of GPS sources. First the radio powers, as given in table 1 and 2, are corrected for the cosmological evolution of the luminosity function. This correction factor as function of redshift is equation 22. For example, at z=1, the luminosity function has shifted a factor 10 towards higher luminosities, and therefore 2128+048 with a radio power of $`10^{27.8}`$ W/Hz will contribute to the LLF at $`10^{26.8}`$ W/Hz. Note that this correction is independent of luminosity, and therefore the difference in radio luminosity of young and old sources does not have to be accounted for. Note however, that the increase in number density is dependent on radio luminosity due to a change in the slope of the luminosity function. The number densities increase from z=0 to z=1 by a factor of 5 and 150 for low and high luminosity sources respectively.
Figure 7 shows the corrected and uncorrected radio powers for a source with flux densities of 10 mJy, 100 mJy and 1 Jy (assuming a spectral index of $`0.5`$ at about 5 GHz). Interestingly, the corrected radio power for a source with an certain observed flux density, does not significantly change at $`0.6<z<2.0`$. Hence, although for many GPS galaxies no spectroscopic redshift has been measured, this is not likely to influence the result, since they will probably all be in this redshift range.
The next step is to correct the number of sources observed for the volume of space over which they can be observed. As can be seen from figure 7, a flux density limit in the sample of 1 Jy means that all sources with corrected luminosities greater than 25.4 $`WHz^1Sr^1`$ can be detected out to z=2, and that only source with lower radio power, and consequently at $`z<0.6`$, have to be corrected for the fact that they only could have been seen out to a certain redshift. However, a possible additional redshift limit results from the lower limit in peak frequency at 0.4 GHz in the bright Stangellini et al. sample, and in the faint sample due to the limit in 325-5000 MHz spectral index. For these sources a weight-factor is used equal to the volume of the survey (assuming a redshift limit of 2.0) divided by the maximum volume over which they could have been in the sample, which is dependent on the maximum observable redshift. The corrections above are relatively straight-forward. However, some additional, more complicated corrections have to be made for the faint GPS sample. Firstly, this sample is originally selected at 325 MHz frequency, eg. on the optically thick part of their spectrum. Furthermore, only sources with positive spectral indices between this frequency and 5 GHz were initially selected. Therefore the faint WENSS sample is more biased towards GPS sources with higher peak frequencies than the bright Stanghellini et al. sample. To correct for this we assumed that the parent distribution of peak frequencies is independent of flux density and radio power, and determined what fraction of the Stangellini et al sample would have been included in the sample if it would have been selected as for the faint WENSS sample. It turns out that 26 % of the galaxies in the bright sample have 325 MHz flux densities $`>`$ 1 Jy and positive spectral indices between 325 MHz and 5 GHz. In addition, the bright GPS sample has a limit in optically thin spectral index of $`0.5`$, while several sources in the faint sample have a flatter optically thin spectral index. Taking into account these two effects, the number densities for the faint sample are multiplied by a factor 3.2.
The resulting local luminosity function of GPS sources is shown in figure 8. Note that a luminosity bin (centred at 23.75 Watt/Hz) containing only a single source (B0830+5813) is omitted due to its large uncertainty.
We compared the resulting LLF with an LLF of a simulated radio source population of $`10^6`$ objects, with random ages, and a jet-power distribution defined as in equation 10. The ‘observed’ luminosity of a source was calculated assuming that it had evolved over its lifetime according to the luminosity evolution derived in section 4.1, out to a maximum size, $`r_+`$. At $`r<r_{}`$, the size of the source evolves as $`r=t^{1/2}P_J^{1/4}`$. To avoid a discontinuity in propagation velocity at $`r_{}`$, the source evolves from $`r_{}`$ as,
$$r(P_J,t)=\frac{\gamma }{2}\gamma ^{\frac{1}{\gamma }}(\gamma P_J^{\frac{1}{2}}t)^{\frac{1}{\gamma }}+(1\frac{\gamma }{2})$$
(23)
with $`\gamma =(4\beta )/2`$. The luminosity of a source increases at $`r<r_{}`$ as $`L=P_J^{7/8}(r/r_{})^{2/3}`$, and as $`L=P_J^{7/8}(r/r_{})^{\frac{2}{3}\frac{7}{6}\beta }`$ at $`r>r_{}`$. This results in a similar evolution for large size radio sources as derived in section 4.1, with the luminosity at $`r=r_{}`$ only dependent on $`P_J`$. It was not our aim to determine absolute values for number densities and radio powers with these simulations. The results of the simulation were scaled in such way, that the LLF obtained for large size radio sources, matched the LLF of steep spectrum sources as derived by Dunlop & peacock (1990).
Table 4 lists the important characteristics of the simulated LLF of large size and GPS sources, and their dependence on the free model parameters. The parameters $`\delta `$ and $`\beta `$, as defined in equations 10 and 3, determine the slope of the low and high luminosity part of the LLF of large size radio sources. These were chosen to be similar to the parameters $`a1`$ and $`b1`$ as derived by Dunlop & Peacock (1990), with $`\delta =1.10`$ and $`\beta =1.16`$. This value of $`\beta `$ is slightly lower than derived from X-ray observations of nearby ellipticals ($`\beta =1.52`$, Trinchieri et al. 1986). Note however, that the radio source population is dominated by objects with size $`>20`$ kpc, for which the surrounding medium is dominated by intra-cluster gas, which is expected to have a flatter density gradient. With the parameters $`\delta `$ and $`\beta `$ and the slopes of the LLF of large size radio sources fixed, the relative positions of the break luminosities could be determined. A sharp cut-off will occur near the highest luminosity, $`L_{max}`$. The number of GPS galaxies in the highest luminosity bin, as shown in figure 8, is lower than expected from the extrapolation of the LLF at lower luminosities. This can be explained if this luminosity bin is near the cut-off luminosity $`L_{max}`$. We therefore chose log $`L_{max}`$ to be 27.1 (W Hz$`{}_{}{}^{}1`$). The break luminosity of large size radio sources, is also determined by Dunlop & Peacock (1990) to be log $`L_{LS}`$ = 25.79 (corrected to 5 GHz). As can be seen from table 4, the luminosity ratio $`L_{max}/L_{ls}`$ determines the value of $`r_+/r_{}`$. This corresponds to a maximum size for a radio source of 100 kpc, assuming $`r_{}=1`$ kpc. This value is quite near the turnover seen in the linear size distribution of 3CR galaxies, as shown by O’Dea & Baum (1997). The break luminosity of GPS sources, $`L_{gps}`$, relative to $`L_{max}`$, is dependent on the range of jet-powers $`(P_+/P_{})`$. To let $`L_{gps}`$ coincide with the peak in the observed GPS LLF, a value of $`(P_+/P_{})`$=200 was used.
Although the uncertainties on the datapoints are large and several free parameters enter the simulation, figure 8 shows that the shape of the LLF of GPS sources is as expected. Note that most free parameters are determined by fitting the LLF of large size radio sources to that of Dunlop & Peacock (1990), except $`r_+`$ and $`(P_+/P_{})`$. This analysis should be regarded as an example of how future large and homogeneously defined samples of GPS sources can constrain the luminosity evolution of extragalactic radio sources.
The proposed increase in luminosity for young radio sources seems to be in contradiction to the high number counts of GPS sources with respect to large size radio sources suggesting that they should decrease in radio luminosity by a factor $`10`$ during their lifetime (Fanti et al. 1995, Readhead et al. 1996, O’Dea & Baum 1997). However, this is not the case. Flux density limited samples, as used for these analyses, only probe the most luminous objects at any redshift. As can be seen from figure 8, at high luminosities, the two luminosity function approach each other, due to the flatter slope of the luminosity function of GPS sources. This results in a relatively high number density of GPS sources in flux density limited samples.
### 4.4 Summary and Conclusions
In this paper we show that in addition to the well known correlation between spectral peak frequency and angular size (eg. Fanti et al. 1990), a correlation exists between the peak flux density and angular size of GPS & CSS sources. The strength and sign of these correlations are exactly as expected from SSA theory, assuming equipartition, and are therefore a strong indication that SSA is indeed the cause of the spectral turnovers in these objects. Furthermore, these correlations are consistent with GPS & CSS sources evolving in a self-similar way. Interestingly, the self-similar evolution scenario is better fitted by assuming an equipartition than a constant magnetic field.
In flux density limited samples, GPS galaxies are found at higher redshifts than large size radio sources. Since the lifetimes of radio sources are short compared to cosmological timescales, this can only mean that the slope of their luminosity functions are different, if GPS sources are to evolve into large size radio sources. It is shown that that the slope of a luminosity function is strongly dependent on the evolution of radio power of the individual sources. A new method is introduced to constrain the luminosity evolution of radio sources using the luminosity functions of ‘young’ and ‘old’ objects. It is shown that if GPS sources are increasing in radio power with time, it would result in a relatively flatter slope of their luminosity function compared to that of large size radio sources which decrease in radio power.
A simple model was developed in which a radio source, embedded in a King profile medium, evolves in a self similar way under the equipartition energy assumption. This model indeed results in the suggested increase in luminosity for young radio sources, and decrease in luminosity for old, extended objects. The calculated luminosity function for large size radio sources shows a break and slopes at low and high luminosity comparable to that derived by Dunlop & Peacock (1990) for steep spectrum sources. The local luminosity function (LLF) of GPS sources can not be measured directly since so few GPS sources are found at low redshift. Therefore, the knowledge of the cosmological evolution of the luminosity function of steep spectrum sources, as derived by Dunlop & Peacock (1990), is used, so that the LLF can be derived from the complete samples of bright and faint GPS sources. It is shown that this LLF is as expected for radio sources which increase in luminosity with time, which is confirmed by simulations of the young radio source population. Note however, that the bright and faint GPS samples are constructed in very different ways and that therefore large corrections had to be made. Uncertainties are still very large and ideally large samples of GPS sources should be constructed which are uniformly selected at low and high flux densities.
## 5 Acknowledgements
We like to thank the anonymous referee for carefully reading the manuscript and valuable suggestions. We also thank Christian Kaiser and Jane Dennett-Thorpe for helpful comments. This research was in part funded by the European Commission under contracts ERBFMRX-CT96-0034 (CERES) and ERBFMRX-CT96-086 (Formation and Evolution of Galaxies), and SCI\*-CT91-0718 (The Most Distant Galaxies).
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# The effect of Hund’s Coupling on one-dimensional Luttinger Liquids
## I Introduction
One-dimensional (1D) electron systems provide examples where the electron correlation effects play a dominant role. The metallic phase of 1D interacting electron system is characterized by the separated gapless charge and spin excitations, and it is called the Luttinger liquid phase. Recently there has been great interest on the two coupled Luttinger liquid systems as the model for ladder systems, zigzag system, stripe phase in cuprate superconductor, and nano tubes . Theoretically, the relevancy of inter-chain hopping $`t_{}`$, the opening of spin gap, and the superconducting instabilities have been studied.
In this paper, we investigate two 1D Luttinger liquid systems coupled by Hund’s coupling without interchain hopping $`t_{}`$. This coupled system is interesting as a first step towards 2D $`\mathrm{Sr}_2\mathrm{RuO}_4`$ . A Hubbard chain away from half-filling can be modelled as a Luttinger liquid for arbitary values of on-site repulsion $`U`$, and two coupled Hubbard chains in the weak coupling regime $`U,J_H<t`$ can be investigated using the bosonization technique ($`t`$ is a band width, $`J_H`$ is Hund’s coupling). Nagaosa and Oshikawa investigated two coupled Luttinger liquids using a semiclassical analysis under the condition $`\delta \mathrm{max}(t,t_{})<|J_H|`$, where $`\delta `$ is the density measured from half-filling. They showed that a spin gap opens up in this system and that the charge fluctuations between chains are massive. For generic values of $`t,t_{},J_H`$, their result is valid only near half-filling.
In present paper we investigate this model at arbitrary filling factors away from half filling. We focus on the influence of Hund’s coupling on the charge and spin excitations and its interplay with the repulsive Hubbard interaction. In order to preserve SU(2) invariance in renormalization group analysis we use non-Abelian bosonization. The weak coupling phase diagram can be mapped out from the analysis of the resulting perturbative R.G. equations. We find that Hund’s coupling is relevant even for infinitesimally small initial value of $`J_H`$, and the system flows into a new strong coupling fixed point. At the strong coupling fixed point, all spin excitations are gapped. As for the charge exciations, the symmetric charge excitations $`\varphi _{\rho +},\theta _{\rho +}`$ (See Eq.(9))remain gapless, while the asymmetric charge excitations $`\varphi _\rho ,\theta _\rho `$ are gapped. Thus, our result extends the results by Nagaosa and Oshikawa to general filling for the case $`t_{}=0`$. Furthermore, we have investigated the dependence of the physical properties on the value of the Hubbard repulsion, which is not addressed by Nagaosa and Oshikawa: We find that, even the system is in gapped phase for all $`0U<t`$ and $`0<J_H<t`$, the gap (i.e. the scale where the renormalized coupling constants diverge or become the order of the bare cut-off scale) strongly depends on the relative magnitude of $`U,J_H`$. If $`U>J_H`$ , the strong coupling fixed point can be reached faster than the opposite case $`J_H>U`$.
This paper is organized as follows: In section II, we specify the model and re-cast it in a bosonized form. In section III, the renormalization group equations are derived using the operator product expansion method. We present the renormalization group flows in section IV and conclude with discussions in section V.
## II Model
We consider two Hubbard chains coupled by Hund’s coupling.
$$H=t\underset{ij,l}{}c_{li\alpha }^{}c_{lj\alpha }+U\underset{i,l}{}n_{il}n_{il}\underset{i}{}J_H𝐒_{1i}𝐒_{2i}=H_0+H_U+H_H,$$
(1)
where $`\alpha `$ is a spin index and $`l=1,2`$ is a band index. $`J_H`$ is positive. Note that there is no inter-chain hopping. We assume the implicit normal ordering of all interaction terms in the Hamiltonian. For simplicity, the band dependences of hopping parameter $`t_{ij,l}`$ and Hubbard interaction $`U_l`$ are ignored, two coupled chains are considered instead of three, and only the nearest neighbor hopping is considered.
We investigate only the weak coupling regime $`U<t,J_H<t`$, in which the linearization of non-interacting electron spectrum near Fermi points and the perturbative treatment of interactions are legitimate. A generic filling away from half-filling is considered, so that the umklapp process can be neglected. Because the inter-chain hopping term is absent, the non-interacting energy spectrums of two bands are degenerate. The linearization gives the decomposition $`c_{i\alpha l}/\sqrt{a}\psi _{R\alpha l}(x)e^{ik_Fx}+\psi _{L\alpha l}(x)e^{ik_Fx}`$. $`a`$ is a lattice spacing. Substituting the decomposition into $`H_0`$ of Eq.(1), we get
$$H_0=\underset{\alpha l}{}𝑑xv_F\left(\psi _{Rl\alpha }^{}i_x\psi _{Rl\alpha }\psi _{Ll\alpha }^{}i_x\psi _{Ll\alpha }\right),$$
(2)
where $`v_F=2ta\mathrm{sin}(k_Fa)`$. To express the Hamiltonian in the bosonized form, it is convenient to introduce the (chiral) charge and spin current operators.
$$J_{R,l}=\underset{\alpha }{}\psi _{R,l,\alpha }^{}\psi _{R,l,\alpha },𝐉_{R,l}=\underset{\alpha \beta }{}\psi _{R,l,\alpha }^{}\frac{𝝈_{\alpha \beta }}{2}\psi _{R,l,\beta }.$$
(3)
The left moving currents are defined analogously. The operator product expansion allows us to express $`H_0`$ in terms of currents .
$$H_0=\underset{l}{}𝑑x\left[\frac{\pi v_F}{2}(J_{Ll}^2+J_{Rl}^2)+\frac{2\pi v_F}{3}(𝐉_{Ll}𝐉_{Ll}+𝐉_{Rl}𝐉_{Rl})\right].$$
(4)
The Hubbard part $`H_U`$ can be expressed in terms of currents alone, while the Hund’s coupling term $`H_H`$ cannot be, since
$$𝐒_l/a=𝐉_{Ll}+𝐉_{Rl}+e^{2ik_Fx}\psi _{Rl\alpha }^{}\frac{𝝈_{\alpha \beta }}{2}\psi _{Ll\beta }+e^{+2ik_Fx}\psi _{Ll\alpha }^{}\frac{𝝈_{\alpha \beta }}{2}\psi _{Rl\beta }.$$
(5)
At this point, we apply Abelian and Non-abelian bosonization to the charge and spin part, respectively. The phase field $`\varphi _{\rho l}`$ and its conjugate momentum $`\mathrm{\Pi }_{\rho l}`$ are introduced for the charge currents .
$$J_{Rl}+J_{Ll}=\sqrt{\frac{2}{\pi }}_x\varphi _{\rho l},J_{Rl}J_{Ll}=\sqrt{\frac{2}{\pi }}\mathrm{\Pi }_{\rho l}.$$
(6)
And for the spin currents SU(2) matrix field $`g_l(x,\tau )`$ is introduced .
$$J_{Ll}^a=\frac{i}{\pi }\mathrm{tr}(_zg_lg_l^1\frac{\sigma ^a}{2}),J_{Rl}^a=+\frac{i}{\pi }\mathrm{tr}(g_l^1_{\overline{z}}g_l\frac{\sigma ^a}{2}),$$
(7)
with $`z=x+iv_F\tau `$, $`\overline{z}=xiv_F\tau `$. The elements of the matrix field $`g_l`$ can also be expressed in terms of scalar field through the vertex operator construction
$$g_{\alpha \beta }\left(\begin{array}{cc}e^{i\sqrt{2\pi }\varphi _{\sigma l}}& e^{i\sqrt{2\pi }\theta _{\sigma l}}\\ e^{i\sqrt{2\pi }\theta _{\sigma l}}& e^{i\sqrt{2\pi }\varphi _{\sigma l}}\end{array}\right).$$
(8)
$`\varphi _{\sigma l}`$ is the phase field of the Abelian bosonized spin currents, and $`\theta _{\sigma l}`$ is the conjugate field. The above Abelian bosonized representation is useful in interpreting the results of R.G. equations physically. For later convenience, we define the symmetric and asymmetric charge and spin modes.
$$\varphi _{\rho \pm }=\frac{1}{\sqrt{2}}\left(\varphi _{\rho 1}\pm \varphi _{\rho 2}\right),\varphi _{\sigma \pm }=\frac{1}{\sqrt{2}}\left(\varphi _{\sigma 1}\pm \varphi _{\sigma 2}\right).$$
(9)
According to the Non-abelian bosonization rule the fermion bilinear can be represented as
$$\underset{\alpha \beta }{}\psi _{Rl\alpha }^{}\frac{𝝈_{\alpha \beta }}{2}\psi _{Ll\beta }\mathrm{tr}\left(g_l^{}(x)\frac{𝝈}{2}\right)e^{i\sqrt{2\pi }\varphi _{\rho l}(x)}.$$
(10)
As a result of the parametrizations of currents and fermion bilinear Eq.(6,7,10), it is possible to express the total Hamiltonian in terms of the scalar fields $`\varphi _{\rho l}`$ and the SU(2) matrix fields $`g_l`$. For the renormalization group analysis, the Euclidean Lagrangian formulation is more convenient.
$`S`$ $`=`$ $`{\displaystyle \frac{v_c}{2K_c}}{\displaystyle \underset{l}{}}{\displaystyle 𝑑x𝑑\tau \left[\left(\frac{\varphi _{\rho l}}{v_c\tau }\right)^2+\left(\frac{\varphi _{\rho l}}{x}\right)^2\right]}+{\displaystyle \underset{l}{}}S[g_l]_{WZW}`$ (11)
$``$ $`\lambda _1{\displaystyle \underset{l}{}}{\displaystyle 𝑑x𝑑\tau 𝐉_{Rl}𝐉_{Ll}}`$ (12)
$``$ $`{\displaystyle 𝑑x𝑑\tau \left[\lambda _2\left(𝐉_{L1}𝐉_{L2}+𝐉_{R1}𝐉_{R2}\right)+\lambda _3\left(𝐉_{L1}𝐉_{R2}+𝐉_{R1}𝐉_{L2}\right)\right]}`$ (13)
$``$ $`\lambda _4{\displaystyle }dxd\tau [\mathrm{tr}\left(g_1^{}(x){\displaystyle \frac{𝝈}{2}}\right)e^{i\sqrt{2\pi }\varphi _{\rho 1}(x)}\mathrm{tr}\left(g_2(x){\displaystyle \frac{𝝈}{2}}\right)e^{+i\sqrt{2\pi }\varphi _{\rho 2}(x)}+\mathrm{h}.\mathrm{c}]`$ (14)
$``$ $`\lambda _5{\displaystyle }dxd\tau [\mathrm{tr}\left(g_1^{}(x)\right)e^{i\sqrt{2\pi }\varphi _{\rho 1}(x)}\mathrm{tr}\left(g_2(x)\right)e^{+i\sqrt{2\pi }\varphi _{\rho 2}(x)}+\mathrm{h}.\mathrm{c}]`$ (15)
$`+`$ $`{\displaystyle \frac{v_c}{2K_2}}{\displaystyle \underset{l}{}}{\displaystyle 𝑑x𝑑\tau \mathrm{\hspace{0.17em}2}\left[\left(\frac{\varphi _{\rho 1}}{v_c\tau }\frac{\varphi _{\rho 2}}{v_c\tau }\right)+\left(\frac{\varphi _{\rho 1}}{x}\frac{\varphi _{\rho 2}}{x}\right)\right]}`$ (16)
where
$`v_c=v_F\sqrt{1+{\displaystyle \frac{Ua}{\pi v_F}}}`$ (17)
and
$`K_c=1/\sqrt{1+{\displaystyle \frac{Ua}{\pi v_F}}}.`$ (18)
$`S[g]_{WZW}`$ is the so-called Wess-Zumino-Witten action, and we don’t need its explicit form. We just note that the Fermi velocity for the matrix field $`g`$ is renormalized by the interaction $`v_Fv_s=v_F(1\frac{Ua}{2\pi v_F})`$. The initial value (bare value) of $`\lambda _1`$ is $`2Ua>0`$, and the initial value of $`\lambda _2`$, $`\lambda _3`$ and $`\lambda _4`$ is $`J_Ha>0`$. Although the initial values of $`\lambda _2`$, $`\lambda _3`$ and $`\lambda _4`$ are identical they will scale differently under the R.G. flow. The $`1/K_2`$ and $`\lambda _5`$ terms in the above action is absent in the bare action (the initial values are 0), but they are generated by the second order perturbation.
The action Eq.(11) defines a critical fixed point of two independent $`U(1)`$ Gaussian model and two independent SU(2) $`k=1`$ Wess-Zumino-Witten (WZW) theory. The terms in Eq.(14)($`\lambda _i,i=1,2,3,4,5`$) are the perturbations to the fixed point. We study how the perturbations influence the fixed point defined by Eq.(11) via the renormalization group method in the next section.
## III Renormalization Group Equations
It is well-known that the renormalization group (R.G.) equations for the fixed point Hamiltonian perturbed by a number of scaling operators, up to the second order in couplings, are determined by the scaling dimensions of the scaling operators and the operator product expansion (OPE) coefficents of the scaling operators . Explicitly,
$$\frac{d\lambda _i}{d\mathrm{ln}L}=(dx_i)\lambda _i\underset{jk}{}c_{ijk}\lambda _j\lambda _k+\mathrm{},$$
(19)
where $`\lambda _i`$ is the coupling constants of the scaling operator $`O_i`$, $`d=1+1=2`$ is the space-time dimension, and $`x_i`$ is the scaling dimension of the scaling operator. $`L`$ is a cut-off length or time scale. The OPE coefficents $`c_{ijk}`$ is determined by
$$O_i(r)O_j(0)\underset{k}{}\frac{c_{ijk}}{r^{x_i+x_jx_k}}O_k(0).$$
(20)
Among the perturbations in Eq.(14), only $`\lambda _4`$ has a scaling dimension different from 2, which is $`x_4=1+K_c`$.
The relevant operator product expansions for our case are
$`e^{i\alpha \varphi _{\rho l}(x,\tau )}e^{i\alpha \varphi _{\rho l}(0,0)}`$ $``$ $`{\displaystyle \frac{1}{|z_c|^{\alpha ^2K_c/2\pi }}}[1+i\alpha (z_c_{z_c}\varphi _{\rho l}(0,0)+\overline{z}_c_{\overline{z}_c}\varphi _{\rho l}(0,0))`$ (21)
$``$ $`\alpha ^2|z_c|^2_{z_c}\varphi _{\rho l}(0,0)\overline{z}_c_{\overline{z}_c}\varphi _{\rho l}(0,0){\displaystyle \frac{\alpha ^2}{2}}(z_c^2(_{z_c}\varphi _{\rho l})^2+\overline{z}_c^2(_{\overline{z}_c}\varphi _{\rho l})^2)+\mathrm{}],`$ (22)
with $`z_c=x+iv_c\tau `$, $`\overline{z}_c=xiv_c\tau `$.
The OPE of WZW model ($`k=1`$) are given by
$`J_{Ll}^a(z)J_{Lm}^b(w)`$ $``$ $`{\displaystyle \frac{1}{(2\pi )^2}}{\displaystyle \frac{k}{2}}{\displaystyle \frac{\delta _{lm}\delta _{ab}}{(z_sw_s)^2}}+{\displaystyle \frac{1}{2\pi }}{\displaystyle \frac{iϵ^{abc}}{z_sw_s}}J_{Ll}^c(w)+\mathrm{}`$ (23)
$`J_{Ll}^a(z)(x,\tau )g_m(0,0)`$ $``$ $`{\displaystyle \frac{\delta _{lm}}{2\pi z_s}}t^ag_m+\mathrm{}`$ (24)
$`\mathrm{tr}(g_l(x,\tau )\sigma ^a)\mathrm{tr}(g_m^{}(0,0)\sigma ^b)`$ $``$ $`{\displaystyle \frac{\delta _{lm}}{|z_s|}}[\delta _{ab}+iϵ^{abc}(z_sJ_{Ll}^c(0,0)+\overline{z}_sJ_{Rl}^c(0,0))`$ (25)
$`+`$ $`z_s^2{\displaystyle \underset{c}{}}J_{Ll}^c(0,0)J_{Ll}^c(0,0)+\overline{z}_s^2{\displaystyle \underset{c}{}}J_{Rl}^c(0,0)J_{Rl}^c(0,0)]+\mathrm{},`$ (26)
where $`z_s=x+iv_s\tau ,\overline{z}_s=xiv_s\tau `$, and $`t^a`$ are the generators of SU(2) algebra in the spin basis $`t^a=\frac{\sigma ^a}{2},\mathrm{tr}(t^at^b)=\frac{\delta ^{ab}}{2},[t^a,t^b]=iϵ^{abc}t^c`$. The terms of dimension higher than 2 are omitted since they are irrelevant in our computation. The OPE Eqs.(23,24) are the fundamental consequencs of the chiral gauge invariance of WZW model. The OPE Eq.(26) follows from the fusion rule of Kac-Moody primary fields . For level 1 $`(k=1)`$ SU(2) WZW theory, the product of spin 1/2 primary fields can only give rise to the identity field and its descedants . $`J`$ and $`J^2`$ terms in Eq.(26) are the descendant fields of the identity field.
Using Eqs.(19,21,23,24,26), the R.G. equations can be computed straighforwardly. The results are :
$$\frac{d\frac{1}{K_c}}{d\mathrm{ln}L}=\frac{\lambda _4^2}{\pi v_cv_s}f_1(\frac{v_c}{v_s}),f_1(u)=_0^\pi \frac{d\theta }{(\mathrm{cos}^2\theta +u^2\mathrm{sin}^2\theta )^{K_c1}}.$$
(27)
$$\frac{d\frac{1}{K_2}}{d\mathrm{ln}L}=\frac{\lambda _4^2}{\pi v_cv_s}f_1(\frac{v_c}{v_s}),$$
(28)
$$\frac{d\lambda _1}{d\mathrm{ln}L}=\frac{\lambda _1^2}{2\pi v_s}.$$
(29)
$$\frac{d\lambda _2}{d\mathrm{ln}L}=0.$$
(30)
$$\frac{d\lambda _3}{d\mathrm{ln}L}=\frac{1}{2\pi }\left(\frac{\lambda _4^2}{v_c}+\frac{\lambda _3^2}{v_s}\right).$$
(31)
$$\frac{d\lambda _4}{d\mathrm{ln}L}=(1K_c)\lambda _4+\frac{\lambda _1\lambda _4}{4\pi v_s}\frac{\lambda _3\lambda _4}{2\pi v_s}.$$
(32)
$$\frac{d\lambda _5}{d\mathrm{ln}L}=(1K_c)\lambda _5\frac{3}{16\pi v_s}\lambda _3\lambda _4.$$
(33)
We note the similarity of the above R.G. equations with those of 1D Kondo lattice problem . The Eq.(27,28) suggests a combination of fields $`\varphi _{\rho \pm }=\frac{1}{\sqrt{2}}(\varphi _{\rho 1}\pm \varphi _{\rho 2})`$. In terms of $`\varphi _{\rho \pm }`$, the sum of the charge parts of our action becomes
$`S_c`$ $`=`$ $`{\displaystyle \frac{v_c}{2}}({\displaystyle \frac{1}{K_c}}+{\displaystyle \frac{1}{K_2}}){\displaystyle 𝑑x𝑑\tau \left[\left(\frac{\varphi _{\rho +}}{v_c\tau }\right)^2+\left(\frac{\varphi _{\rho +}}{x}\right)^2\right]}`$ (34)
$`+`$ $`{\displaystyle \frac{v_c}{2}}({\displaystyle \frac{1}{K_c}}{\displaystyle \frac{1}{K_2}}){\displaystyle 𝑑x𝑑\tau \left[\left(\frac{\varphi _\rho }{v_c\tau }\right)^2+\left(\frac{\varphi _\rho }{x}\right)^2\right]}.`$ (35)
Then from Eq.(27,28),
$$\frac{(\frac{1}{K_c}+\frac{1}{K_2})}{\mathrm{ln}L}=0,\frac{(\frac{1}{K_c}\frac{1}{K_2})}{\mathrm{ln}L}=\frac{2\lambda _4^2}{\pi v_cv_s}f_1(\frac{v_c}{v_s}).$$
(36)
We investigate the properties of the above R.G. equations in the next section.
## IV R.G. flow and phase diagram
We consider only the ferromagnetic Hund’s coupling constants, so that the couplings $`\lambda _i,i=1,\mathrm{},5`$ should be taken to be positive in the physically relevant range. Given the initial values of $`K_c,\lambda _i,i=1,\mathrm{},4`$, the R.G. flows are uniquely determined. The initial values are determined by Fermi momentum $`k_F`$, $`U/t`$, and $`J_H/t`$. The derived R.G equations are valid until $`\mathrm{max}[\lambda _i/v_s]O(1)`$. If all $`\lambda _i`$ converge to finite values as $`t=\mathrm{ln}L\mathrm{}`$, the initial fixed point is stable and the R.G. equations are valid along the whole R.G. trajectory since the $`\lambda _i(\mathrm{})`$ can be made arbitrarily small by taking sufficiently small $`J_H`$ and $`U`$ . If any coupling constant diverges, the asymptotic behaviour of all coupling constants can be determined by R.G equations .
Some properties of the R.G. equations can be understood by simple inspection without numerical integration. We find that $`\lambda _2`$ is marginal from Eq.(30). The marginality of $`\lambda _2`$ is due to the chirality of interaction: it couples only the currents with the same chirality. The symmetric charge mode $`\varphi _{\rho +}`$ is not renormalized as can be seen in Eq.(36). This is physically obvious since only the relative charge fluctuations affect Hund’s coupling. On the contrary, we see from Eq.(36) that the Luttinger parameter of the asymmetric charge mode $`\varphi _\rho `$ renormalizes to zero when the R.G. equations are extended beyond its validity range. This implies that $`\varphi _\rho `$ is pinned, and it is consistent with the formation of a charge gap from $`\lambda _4`$ renormalization (see below). In Eq.(32), the first term in the right hand side is the most dominant for the range $`J_HU<t`$, in which case $`\lambda _4`$ is certainly relevant. It gives rise to both the charge gap of $`\varphi _\rho `$ and the spin gaps of $`\theta _\sigma `$ and $`\varphi _{\sigma +}`$, as can be seen by expressing $`\lambda _4`$ term in the effective action in Abelian bosonized form using Eq.(8). However, for the opposite range $`UJ_H<t`$, the first and the third terms in the right hand side compete, and only the numerical integrations can determine the limiting behaviour. In this range, $`\lambda _4`$ is expected to increase slower compared with that of the range $`J_HU<t`$. $`\lambda _5`$ starts to grow slowly initially since the initial value vanishes. But eventually the first term of Eq.(33) dominates and $`\lambda _5`$ becomes relevant.
The values of gaps can be estimated in the range $`J_HU<t`$ by integrating the Eq.(32) until the $`\lambda _4`$ grows up to the cut-off value $`ta`$.
$$\mathrm{\Delta }_{c,s}t\left[\frac{J_H}{t}\right]^{\frac{1}{1K_c}}.$$
(37)
Note that this estimated value of the gaps is larger for the stronger Hubbard repulsion. For the range $`UJ_H<t`$, no analytical expression is available. It is worth noting that the crucial $`\lambda _4`$ interaction is not of the current-current interaction type, thus it is not constrained by the current conservation law and allows the presence of the anomalous scaling dimension.
We have numerically integrated R.G. equations for a range of different initial values $`U/t,J_H/t`$. Some typical R.G. flows at $`U=0`$ and $`k_F=\pi /4`$ are shown in Fig.1 and Fig.2. The Fig.1 clearly shows the diverging $`\lambda _4`$ as $`t`$ increases for a range of initial values of $`J_H`$. After the initial transient period all flows merge, which implies that the first term in Eq.(32) dominates. As mentioned above, for non-zero $`U`$, the divergence of $`\lambda _4`$ is much faster. In Fig.2, the R.G flow for $`U=0`$ and a very small initial value of $`J_H=0.01`$ is shown. Initially, the $`\lambda _4`$ slightly decreases due to the last term in R.G Eq.(32), but soon it begins to increase due to the first term of Eq.(32), and eventually diverges. From the examination of R.G flows, we conclude that the charge and spin gaps open up for arbitrarily small initial values of $`J_H`$ irrespective of $`U`$. Our system Eq.(1) is in the massless phase only for a line $`J_H=0`$, which is simply two decoupled Hubbard chain, where each of Hubbard chain is in critical Luttinger liquid phase. In any other region with positive $`J_H>0`$, $`\lambda _4`$ is relevant and the charge and spin gaps are present.
## V Discussions
The spin sector of a single Hubbard chain at low energy is in the same universality class of antiferromagnetic(AF) XXX Heisenberg chain . Therefore, the spin sector of the system Eq.(1) can be approximately described as two spin 1/2 XXX chains coupled by ferromagnetic exchange coupling at low energy. In the strong coupling limit $`J_H\mathrm{}`$ we only need to consider the triplet combination of spins from each spin chain. When projected onto this triplet subspace, the coupled spin 1/2 chains become essentially single spin 1 chain. The spin 1 chain is well known to be gapped , and the gap is called Haldane gap . Also Strong and Millis studied the coupled spin 1/2 chains and showed that the sping gap opens even for infinitesimally small $`|J_H|`$, which is consistent with our result. The spin gap we have obtained can be also understood from the semiclassical point of view .
It is interesting to examine other parameter domains even though they are not directly related to our system. First, the antiferromagnetic exchange can be considered. It corresponds to a negative $`J_H`$, which implies the signs of $`\lambda _i,i=2,3,4,5`$ to be flipped. The R.G equation of charge Luttinger parameters, Eq.(36), is unchanged, but $`\lambda _3,\lambda _4,\lambda _5`$ now become relevant. $`\lambda _3`$ only generates a spin gap, while $`\lambda _4,\lambda _5`$ generate both the charge and spin gap. The relative initial values of the couplings determine the gap which is going to be opened the fastest, and R.G flows depend critically on the gap. The detailed investigation will not be presented here. Second, the attractive Hubbard interaction is another possibility (with positive $`J_H`$). In this case, the sign of $`\lambda _1`$ is flipped and the initial value of $`K_c`$ is greater than 1. In this case, $`\lambda _1`$ becomes marginally relevant, and $`\lambda _4`$ becomes irrelevant. The marginally relevant $`\lambda _1`$ would open up a spin gap for each band $`l=1,2`$. This is essentially the Luther-Emery spin gap , and the Hund’s coupling does not play any role here.
In this paper we have investigated two Hubbard chains coupled by Hund’s coupling using renormalization group method for arbitrary densities away from half filling. We find that in the weak coupling regime Hund’s coupling is always relevant, irrespective of the stength of $`U`$, and opens up gaps for both symmetric and antisymmetric spin modes and for antisymmetric charge mode.
ACKNOWLEDGEMENTS
H.C. Lee is grateful to Prof. A. Changrim for discussions. This work was supported by Korea science and engineering foundation through the Quantum-functional Semiconductor Research Center at Dongguk University.
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# Weak Phase 𝛾 and Strong Phase 𝛿 from CP Averaged 𝐵→𝜋𝜋 and 𝜋𝐾 Decays
## Abstract
Assuming $`SU(3)`$ symmetry for the strong phases in the four decay modes $`B\pi ^{}\pi ^+,\pi ^0\pi ^+,\pi ^{}K^+,\pi ^{}\overline{K}^0`$ and ignoring the relative small electroweak penguin effects in those decays, the weak phase $`\gamma `$ and the strong phase $`\delta `$ can be determined in a model independent way by the CP-averaged branching ratios of the four decay modes. It appears that the current experimental data for $`B\pi \pi `$ and $`\pi K`$ decays prefer a negative value of $`\mathrm{cos}\gamma \mathrm{cos}\delta `$. By combining with the other constraints from $`V_{ub}`$, $`B_{d,s}^0\overline{B}_{d,s}^0`$ mixings and indirect CP-violating parameter $`ϵ_K`$ within the standard model, two favorable solutions for the phases $`\gamma `$ and $`\delta `$ are found to lie in the region: $`35^{}\gamma 62^{}`$ and $`106^{}\delta 180^{}`$ or $`86^{}\gamma 151^{}`$ and $`0^{}\delta 75^{}`$ within 1$`\sigma `$ standard deviation. It is noted that if allowing the standard deviation of the data to be more than 1$`\sigma `$, the two solutions could approach to one solution with a much larger region for the phases $`\gamma `$ and $`\delta `$. Direct CP asymmetry $`a_{ϵ^{\prime \prime }}^{(\pi ^{}K^+)}`$ in $`B\pi ^{}K^+`$ decay can be as large as the present experimental upper bound. Direct CP asymmetry $`a_{ϵ^{\prime \prime }}^{(\pi ^+\pi ^{})}`$ in $`B\pi ^{}\pi ^+`$ decay can reach up to about $`40\%`$ at 1$`\sigma `$ level.
The study of CP-violation is one of the central topics in the present day elementary particle physics. In the standard model (SM), all the CP violating phenomena arise from a single complex phase of the Cabbibo-Kobayashi-Maskawa(CKM) matrix elements. If the CKM phase is the only source of CP violation, some unitarity relations such as $`V_{ud}V_{ub}^{}+V_{cd}V_{cb}^{}+V_{td}V_{tb}^{}=0`$ will hold . The unitarity relations can be represented geometrically by a set of triangles called unitarity triangles . The three angles in the triangle containing $`b`$ and $`d`$ quarks are often denoted by $`\alpha ,\beta `$ and $`\gamma `$ with $`\alpha +\beta +\gamma =\pi `$ in the SM. Thus one of the important issues is to precisely determine those angles and their sum. Any deviation of the sum from $`\pi `$ will be a signal of new physics.
Although angles $`\alpha `$ and $`\beta `$ may be well measured via the time dependent measurements of $`B\psi K_S`$ and $`B\pi \pi `$, the determination of angle $`\gamma `$ is a great challenge for both theorists and experimentists. In the recent years much work has been done on this issue. As it was first proposed in Ref., the angle $`\gamma `$ may be extracted through six $`BDK`$ decay rates. The difficulty of this method is that it needs tagging of the CP eigenstate $`D_{CP}^0`$ which is rather difficult in the experiment. It may also be extracted from the decay mode $`B_s^0(D_s^{}K^+,D_s^+K^{})\overline{B}_s^0`$ in a model-independent way since one only needs to extract the rephase invariant observables $`a_{ϵ+ϵ^{}}`$ and $`a_ϵ^{}`$ from a time-dependent measurement. Thus the weak phase is simply given by $`\mathrm{sin}\gamma =a_{ϵ+ϵ^{}}/\sqrt{1+a_ϵ^{}^2}`$. In the recent years an alternative way of using the CP averaged $`B\pi ^\pm K^0,\pi ^+\pi ^0`$ and the CP conjugate $`B^+\pi ^0K^+,B^{}\pi ^0K^{}`$ branching ratios has been aroused a great attention. However this method needs some theoretical input in evaluating the electroweak penguin (EWP) effects. At present, limited by the statistics the difference of CP conjugate rates can not be definitely established
Recently, the CLEO Collaboration has reported the first observation of rare decays $`B\pi ^{}\pi ^+`$ and $`\pi ^0\overline{K}^0`$ . The observation of $`\pi ^0\overline{K}^0`$ complete the set of measurements on $`B\pi K`$ decays. Other channels of $`\pi K`$ have also been largely improved. The most recent results reported by CLEO collab. are (in units of $`10^6`$),
$`Br(B\pi ^{}\pi ^+)`$ $`=`$ $`4.3_{1.4}^{+1.6}\pm 0.5`$ (1)
$`Br(B\pi ^0\pi ^+)`$ $`=`$ $`<12.7(5.6_{2.3}^{+2.6}\pm 1.7)`$ (2)
$`Br(B\pi ^{}K^+)`$ $`=`$ $`17.2_{2.4}^{+2.5}\pm 1.2`$ (3)
$`Br(B\pi ^{}\overline{K}^0)`$ $`=`$ $`18.2_{4.0}^{+4.6}\pm 1.6`$ (4)
$`Br(B\pi ^0K^+)`$ $`=`$ $`11.6_{2.71.3}^{+3.0+1.4}`$ (5)
$`Br(B\pi ^0\overline{K}^0)`$ $`=`$ $`14.6_{5.13.3}^{+5.9+2.4}`$ (6)
Although only the upper bound of $`\pi ^+\pi ^0`$ is given, the CLEO Collab. also quote a value of $`Br(B\pi ^0\pi ^+)=5.6_{2.3}^{+2.6}\pm 1.7`$. This will be improved by the future measurements. The relative small value of $`\pi ^{}\pi ^+`$, the almost equal $`K\pi `$ rates: $`\pi ^{}\overline{K}^0\pi ^{}K^+`$ and large $`\pi ^0\overline{K}^0`$ seem to be in conflict with the theoretical predictions. However, as it was pointed out in Ref., if one takes the weak phase $`\gamma `$ of the CKM matrix elements to be larger than $`90^{}`$ and include the EWP effects, the situation for $`\pi ^0K^+`$ may be improved greatly, but for $`\pi ^0\overline{K}^0`$ it may become worse as the EWP-SP (strong penguin) interference in $`\overline{K}^0\pi ^0`$ decay is likely to be destructive. Some alternative ways in solving this puzzle are also proposed, such as the small $`|V_{ub}|`$ in $`B\pi ^+\pi ^{}`$ and the use of different form factors and the possibility of large final state interaction phase. It may also be interesting to consider the new physics effects in those decay modes.
Note that the theoretical description on nonleptonic $`B`$ decays is model dependent. Although the short-distance effects are calculable from the Wilson coefficients, one has to assume factorization approach and adopt some models in evaluating the long-distance effects. It may then concern many phenomenological parameters, such as the decay constant of $`B`$ meson, the transition form factors as well as the effective color number $`N_C`$, which still suffer from large uncertainties. Thus the precision of theoretical calculations is unfortunately limited.
In this paper, we shall consider some less model dependent ways to extract both the weak phase $`\gamma `$ and the strong phase $`\delta `$ due to final state interactions. The basic point is to assume approximate relations among the strong phases and choose the decay modes with relative small EWP effects so that one could ignore their contributions as the first step approximation. For this purpose, we take the following four interesting decay modes: $`B\pi ^{}\pi ^+,\pi ^0\pi ^+,\pi ^{}K^+,\pi ^{}\overline{K}^0`$. It will be seen that, under the above assumptions and considerations, the four CP-averaged branching ratios could be used to extract the phases $`\gamma `$ and $`\delta `$ as well as the relative contributions between strong penguin (SP) graphs and tree graphs without additional theoretical inputs. Though such a treatment still suffers from some uncertainties, it could directly provide us useful constraints and insight on the phases $`\gamma `$ and $`\delta `$. We will show that at the 1$`\sigma `$ level of the current experimental data, there exist two correlated regions between $`\gamma `$ and $`\delta `$, which are corresponding to two solutions of negative $`\mathrm{cos}\gamma \mathrm{cos}\delta `$, i.e., one solution is with positive $`\mathrm{cos}\gamma `$ but negative $`\mathrm{cos}\delta `$, another with negative $`\mathrm{cos}\gamma `$ but positive $`\mathrm{cos}\delta `$. While at more than 1$`\sigma `$ level a much larger region for the phases $`\gamma `$ and $`\delta `$ is allowed.
Generally, the $`B`$ decay amplitude can be decomposed by several $`SU(3)`$ invariant Fenyman diagrams. In this decomposition one may see that the amplitudes of decay $`B\pi \pi `$ and $`B\pi K`$ are correlated. This can be used to study the penguin effectcs as well as the strong phases in those modes. In $`SU(3)`$ limits, the decay amplitudes of $`B\pi ^{}\pi ^+,\pi ^0\pi ^+,\pi ^{}K^+,\pi ^{}\overline{K}^0`$ have the following forms
$`𝒜(B\pi ^0\pi ^+)`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}V_{ud}V_{ub}^{}({\displaystyle \frac{T+C}{\sqrt{2}}})`$ (7)
$`𝒜(B\pi ^{}\pi ^+)`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}[V_{ud}V_{ub}^{}(TP)V_{cd}V_{cb}^{}P]`$ (8)
$`𝒜(B\pi ^{}\overline{K}^0)`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}V_{ts}V_{tb}^{}P^{}`$ (9)
$`𝒜(B\pi ^{}K^+)`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}[V_{us}V_{ub}^{}T^{}+V_{ts}V_{tb}^{}P^{}]`$ (10)
where the factor $`1/\sqrt{2}`$ in Eq.(7) comes from the $`\pi ^0`$ wave function. $`T,T^{}(C)`$ and $`P,P^{}`$ denote the Tree(Color suppressed) and QCD penguin amplitude with different strong phases:
$`T=|T|e^{i\delta _T}`$ $`,P=|P|^{i\delta _P}`$ (11)
$`T^{}=|T^{}|e^{i\delta _T^{}}`$ $`,P^{}=|P^{}|^{i\delta _P^{}}`$ (12)
In the expression for the amplitude $`𝒜(B\pi ^{}\pi ^+)`$ in Eq.(8), the unitarity relation of CKM matrix elements $`V_{ud}V_{ub}^{}+V_{cd}V_{cb}^{}+V_{td}V_{tb}^{}=0`$ has been used to remove the factor $`V_{td}V_{tb}^{}`$ which comes from the inner t-quark of the QCD penguins. This allows us to extract the weak phase $`\gamma `$ instead of $`\alpha `$, which is different from the usual treatments.
The charge conjugate decay amplitude can be obtained by simply inverting the sign of weak phase $`\gamma `$. We then get the CP- averaged branching ratios:
$`Br(B\pi ^+K^0)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(B^0\pi ^+K^0+\overline{B^0}\pi ^{}\overline{K}^0\right)|V_{ts}V_{tb}^{}|^2|P^{}|^2,`$ (13)
$`Br(B\pi ^+\pi ^0)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(B^0\pi ^+\pi ^0+\overline{B^0}\pi ^{}\pi ^0\right){\displaystyle \frac{1}{2}}|V_{ud}V_{ub}^{}|^2|T+C|^2,`$ (14)
$`Br(B\pi ^{}K^+)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(Br(B^0K^+\pi ^{})+Br(\overline{B}^0K^{}\pi ^+)\right)`$ (15)
$``$ $`|V_{us}V_{ub}^{}|^2\left|T^{}\right|^22|V_{us}V_{ub}^{}||V_{ts}V_{tb}^{}|\left|T^{}P^{}\right|\mathrm{cos}\delta \mathrm{cos}\gamma +|V_{ts}V_{tb}^{}|^2\left|P^{}\right|^2,`$ (16)
$`Br(B\pi ^{}\pi ^+)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(Br(B^0\pi ^+\pi ^{})+Br(\overline{B}^0\pi ^+\pi ^{})\right)`$ (17)
$``$ $`|V_{ud}V_{ub}^{}|^2|(\left|T\right|e^{i\delta }\left|P\right|)|^2`$ (20)
$`+2|V_{ud}V_{ub}^{}||V_{cd}V_{cb}^{}|\left|TP\right|(\mathrm{cos}\delta \mathrm{cos}\gamma \left|{\displaystyle \frac{P}{T}}\right|\mathrm{cos}\gamma )`$
$`+|V_{cd}V_{cb}^{}|^2\left|P\right|^2.`$
In writing down the above equations, we have neglected the EWP effects. The EWP effects are often thought to be very important in the $`B\pi K`$ decays, but it remains depending on different decay modes. In deed, the EWP effects are of crucial importance in the decay modes: $`B\pi ^0K^0`$ and $`\pi ^0\pi ^0`$ in which the contributions from the tree diagrams are color suppressed. While in the decay modes: $`B\pi ^{}\pi ^+,\pi ^0\pi ^+,\pi ^{}K^+,\pi ^{}\overline{K}^0`$, the EWP effects are relatively small as the contributions from tree diagrams are not color suppressed. As there remain large errors in the current experimental data, for simplicity, we may ignore the EWP effects in those four decay modes as a good approximation in comparison with the experimental uncertainties. To have a quantitative estimation of how good of the approximation, it may be seen from the model dependent calculations, where the contributions from the EWP graphs were found to be about $`1\%,5\%,5\%,8\%`$ in the decay modes $`B\pi ^{}\pi ^+,\pi ^0\pi ^+,\pi ^{}K^+,\pi ^{}\overline{K}^0`$, respectively. It is not difficult to recognize that the relative contributions of the EWP to SP graphs is about $`8\%`$, the relative contributions of the tree diagrams to the SP graphs is about $`40\%`$ in the $`B\pi ^{}K^+`$ decay and is dominated in the $`B\pi ^{}\pi ^+`$ decay.
It is useful to consider the ratios of the decay rates. Let us define
$`R_1{\displaystyle \frac{Br(B\pi ^{}\overline{K}^0)}{Br(B\pi ^0\pi ^+)}}`$ $`>`$ $`1.52(3.25\pm 1.94)`$ (21)
$`R_2{\displaystyle \frac{Br(B\pi ^{}\overline{K}^0)}{Br(B\pi ^{}K^+)}}`$ $`=`$ $`1.06\pm 0.32`$ (22)
$`R_3{\displaystyle \frac{Br(B\pi ^{}\pi ^+)}{Br(B\pi ^{}K^+)}}`$ $`=`$ $`0.25\pm 0.1`$ (23)
In a naive estimation, the ratio between color suppressed diagram and the tree diagram, i.e. $`|C/T|`$ is of the order $`𝒪(0.3)`$ from the color suppression. However, the model dependent calculation show a very small value:$`|C/T|a_2/a_10.05`$ when $`N_C`$ is near 3. By adopting the recent analysis from Ref. which is based on the heavy quark limit, we have $`|C/T|0.08`$.
In a good approximation, (i.e. neglecting the terms proportional to $`|V_{us}V_{ub}^{}T/(V_{ts}V_{tb}^{}P)|^2𝒪(10^2)`$ in $`\pi K`$ modes.) $`\mathrm{cos}\gamma \mathrm{cos}\delta `$ can be solved from the definitions of $`R_2`$ :
$$\mathrm{cos}\gamma \mathrm{cos}\delta \frac{1}{2}\left|\frac{V_{ts}V_{tb}^{}}{V_{us}V_{ub}^{}}\right|\left|\frac{P^{}}{T^{}}\right|\left(1\frac{1}{R_2}\right)$$
(24)
On with including the leading $`SU(3)`$ breaking factor $`f_\pi /f_K`$ in the sense of generalized factorization, one then has
$$\frac{|P|}{|T|}=\frac{|P^{}|}{|T^{}|},\frac{|T|}{|T^{}|}=\frac{f_\pi }{f_K}$$
(25)
and
$$\delta _T=\delta _T^{},\delta _P=\delta _P^{}$$
(26)
Under this approximation, it is then easily seen that the ratio $`|P/T|`$ can be estimated from $`R_1`$:
$$\left|\frac{P}{T}\right|1.09\times \frac{f_\pi }{f_K}\sqrt{\frac{R_1}{2}}\frac{|V_{ud}V_{ub}^{}|}{|V_{ts}V_{tb}^{}|}>0.055$$
(27)
when taking the central value for the mode $`Br(B\pi ^0\pi ^+)=5.6`$, we have $`|P/T|=0.08`$. The value of $`\left|P/T\right|`$ can also be evaluated from the effective Hamitonian and be simply given only by the short distance Wilson coefficients once adopting the factorization approach for the hadronic matrix elements
$$\frac{P}{T}=\frac{1}{a_1}\left[a_4+a_{10}+(a_6+a_8)\frac{2m_\pi ^2}{(m_bm_u)(m_u+m_d)}\right]$$
(28)
which is found to be 0.05 for $`N_C`$=3 and $`m_u+m_d=1.5`$ MeV. Since the validity of Eq.(28) only depends on the assumption of factorization, the ratio $`\left|P/T\right|`$ extracted in this way is helpful to examine how goodness of the factorization approach. It seems that the current experimental data prefer a larger $`\left|P/T\right|`$. This needs to be further confirmed by future experiments.
To naively see the changes of the sign of $`\mathrm{cos}\gamma \mathrm{cos}\delta `$ as $`R_3`$ and $`R_1`$, one may neglect the terms of the order $`𝒪(|P/T|^2)`$ in $`\pi \pi `$ decay modes and use the modified $`SU(3)`$ relations. Then $`\mathrm{cos}\gamma \mathrm{cos}\delta `$ can be simply given in terms of $`R_1`$ and $`R_3`$
$`\mathrm{cos}\gamma \mathrm{cos}\delta `$ $``$ $`1.09\times \left({\displaystyle \frac{\sqrt{2R_1}}{4}}\right){\displaystyle \frac{R_31.68/R_1}{\frac{|V_{cd}V_{cb}^{}||V_{ud}V_{ub}^{}|}{|V_{ts}V_{tb}|}+\frac{|V_{us}V_{ub}^{}|}{|V_{ud}V_{ub}^{}|}\frac{f_K}{f_\pi }R_3}}`$ (29)
$``$ $`1.09\times \left({\displaystyle \frac{\sqrt{2R_1}}{4\lambda }}\right){\displaystyle \frac{R_31.68/R_1}{\frac{f_K}{f_\pi }R_3+\lambda \left|\frac{V_{ub}}{V_{cb}}\right|}}`$ (30)
where $`\lambda =0.22`$ is the Wolfenstein parameter. This shows that $`\mathrm{cos}\gamma \mathrm{cos}\delta `$ will change sign when $`R_3`$ and $`R_1`$ satisfy the approximate relation $`R_31.68/R_1`$. The precise numerical values of $`R_3`$ and $`R_1`$ for changing the sign of $`\mathrm{cos}\gamma \mathrm{cos}\delta `$ may be seen from Fig.1. The values of $`R_3`$ is slightly higher than the ones by Eq.(30).
With the above considerations, the phases $`\gamma `$ and $`\delta `$ can be extracted from $`R_1`$, $`R_2`$ and $`R_3`$ . As the equations are quadratic in $`\mathrm{cos}\gamma `$ and $`\mathrm{cos}\delta `$, there exists a twofold ambiguity in determining these two phases. In Fig.1, we present a contour plot for $`R_2`$ and $`R_3`$ in the $`\mathrm{cos}\gamma \mathrm{cos}\delta `$ plane with $`R_1`$ being fixed at 3.0, 4.5, 6.0, 7.5. Where the solid and dashed contours correspond to different values of $`R_2`$ and $`R_3`$. The points at which the two kind of curves intersect are the solutions of $`\mathrm{cos}\gamma `$ and $`\mathrm{cos}\delta `$. It can be seen from Fig.1 that these contours change largely for different values of $`R_2`$ and $`R_3`$. When $`R_2<1`$ the contours of $`R_2`$ and $`R_3`$ are all in the II and IV quadrants. When $`R_2>1`$ the contours move into the I and III quadrants. This behavior can be understood from Eq.(24). Thus $`\mathrm{cos}\gamma \mathrm{cos}\delta `$ will change sign when $`R_2`$ moves across the point $`R_2=1`$. The changes of $`R_3`$ contours also have the similar reason. From the present data within 1$`\sigma `$ standard deviation, $`R_2`$ and $`R_3`$ are in the range $`0.74R_21.38`$ and $`0.15R_30.35`$, respectively . Since $`R_3`$ is smaller than 0.35 at $`1\sigma `$ level, $`\mathrm{cos}\gamma \mathrm{cos}\delta `$ will be negative for small $`R_1`$. Namely a negative $`\mathrm{cos}\gamma `$ corresponds to a strong phase $`\delta `$ in the first quadrant, for positive $`\mathrm{cos}\gamma `$, the angle $`\delta `$ becomes large and takes values in the second quadrant. From Fig.1, one may see that for large $`R_1`$, a positive solution of $`\mathrm{cos}\gamma \mathrm{cos}\delta `$ is also allowed. For $`R_10.75`$, the allowed range of $`\mathrm{cos}\gamma `$ and $`\mathrm{cos}\delta `$ becomes large and lies in the region: $`0.2\mathrm{cos}\gamma 1`$ and $`1\mathrm{cos}\delta 1`$ or $`1\mathrm{cos}\gamma 0.1`$ and $`1\mathrm{cos}\delta 1`$.
The constraints on the phase $`\gamma `$ may also come from other experiments, such as $`B_{d,s}^0\overline{B}_{d,s}^0`$ mixing, CP-violating parameter $`ϵ_K`$ in the kaon decay, and CKM matrix element $`V_{ub}`$ from semileptonic $`bu`$ decays. Combining all the constraints together and taking the branching ratio for the $`B\pi ^+\pi ^0`$ decay to be $`Br(B\pi ^+\pi ^0)=5.6_{2.3}^{+2.6}\pm 1.7`$, the allowed region for $`\gamma `$ is shown in Fig.2.
It is found that the allowed range for $`\gamma `$ is: $`35^{}\gamma 62^{}`$ or $`86^{}\gamma 151^{}`$, the corresponding values for the phase $`\delta `$ could range from $`106^{}`$ to $`180^{}`$ or from $`0^{}`$ to $`75^{}`$. The allowed regions for the phases $`\gamma `$ and $`\delta `$ are plotted in Fig. 3 and given by the two shadowed ones. One can see from the figure that large region for $`\mathrm{cos}\gamma `$ and $`\mathrm{cos}\delta `$ has been exluded from $`R_1`$, $`R_2`$ and $`R_3`$ when they are at 1$`\sigma `$ level.
Recently, CLEO collaboration has also reported the data on direct CP violation in $`B\pi ^{}K^+`$ decay. Let us now consider CP asymmetries in both $`B\pi \pi `$ and $`B\pi K`$ decays. They are defined as
$`a_{ϵ^{\prime \prime }}^{(\pi ^{}K^+)}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(\overline{B}^0\pi ^{}K^+)\mathrm{\Gamma }(B^0\pi ^+K^{})}{\mathrm{\Gamma }(\overline{B}^0\pi ^{}K^+)+\mathrm{\Gamma }(B^0\pi ^+K^{})}}`$ (31)
$`=`$ $`\left(2|V_{us}V_{ub}^{}V_{ts}V_{tb}^{}|\left|{\displaystyle \frac{P}{T}}\right|\mathrm{sin}\gamma \mathrm{sin}\delta \right)`$ (33)
$`\times \left(|V_{us}V_{ub}^{}|^22|V_{us}V_{ub}^{}V_{ts}V_{tb}^{}|\left|{\displaystyle \frac{P}{T}}\right|\mathrm{cos}\gamma \mathrm{cos}\delta +|V_{ts}V_{tb}^{}|^2\left|{\displaystyle \frac{P}{T}}\right|^2\right)^1,`$
$`a_{ϵ^{\prime \prime }}^{(\pi ^+\pi ^{})}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(\overline{B}^0\pi ^+\pi ^{})\mathrm{\Gamma }(B^0\pi ^+\pi ^{})}{\mathrm{\Gamma }(\overline{B}^0\pi ^+\pi ^{})+\mathrm{\Gamma }(B^0\pi ^+\pi )}}`$ (34)
$`=`$ $`\left(2|V_{ud}V_{ub}^{}V_{cd}V_{cb}^{}|\left|{\displaystyle \frac{P}{T}}\right|\mathrm{sin}\gamma \mathrm{sin}\delta \right)`$ (37)
$`\times (|V_{ud}V_{ub}^{}|^2(12|{\displaystyle \frac{P}{T}}|\mathrm{cos}\delta +|{\displaystyle \frac{P}{T}}|^2)+|V_{cd}V_{cb}^{}|^2|{\displaystyle \frac{P}{T}}|^2`$
$`+2|V_{ud}V_{ub}^{}V_{cd}V_{cb}^{}|\left|{\displaystyle \frac{P}{T}}\right|\mathrm{cos}\gamma (\mathrm{cos}\delta |{\displaystyle \frac{P}{T}}|))^1`$
Here we have used the notation for the general rephase-invariant CP-violating observables classified in . As $`|P/T|`$ is at order of $`𝒪(10^1)`$, for an approximate estimation, one may neglect the $`|P/T|`$ terms in the denominator, thus the above formulae are simplified
$`a_{ϵ^{\prime \prime }}^{(\pi ^{}K^+)}`$ $``$ $`2{\displaystyle \frac{|V_{us}V_{ub}^{}|}{|V_{ts}V_{tb}^{}|}}\left|{\displaystyle \frac{T}{P}}\right|\mathrm{sin}\gamma \mathrm{sin}\delta `$ (38)
$`a_{ϵ^{\prime \prime }}^{(\pi ^+\pi ^{})}`$ $``$ $`2{\displaystyle \frac{|V_{ud}V_{ub}^{}|}{|V_{cd}V_{cb}^{}|}}\left|{\displaystyle \frac{P}{T}}\right|\mathrm{sin}\gamma \mathrm{sin}\delta `$ (39)
$``$ $`0.59\times {\displaystyle \frac{f_\pi ^2}{f_K^2}}R_1a_{ϵ^{\prime \prime }}^{(\pi ^{}K^+)}`$ (40)
which implies that $`a_{ϵ^{\prime \prime }}^{(\pi ^+\pi ^{})}`$ may become large with $`R_1`$ increasing. From the data reported by the CLEO collaboration, no significant deviation from zero was observed: $`a_{ϵ^{\prime \prime }}^{(\pi ^{}K^+)}=0.04\pm 0.16`$. Even at $`90\%`$ CL, $`a_{ϵ^{\prime \prime }}^{(\pi ^{}K^+)}`$ is limited in the range \[-0.30,0.22\]. Incoporating this result, the allowed regions for the phases $`\gamma `$ and $`\delta `$ are further constrained , which is shown in Fig. 3. It is seen that some regions have further been excluded when $`a_{ϵ^{\prime \prime }}^{(\pi ^{}K^+)}`$ has the value within the 1$`\sigma `$ standard deviation. If the values of $`R_1`$, $`R_2`$ and $`R_3`$ are taken to be at 2$`\sigma `$ level, $`a_{ϵ^{\prime \prime }}^{(\pi ^{}K^+)}`$ could be as large as the experimental bound given at 90$`\%`$ CL. From Eq.(40) and $`(\text{27})`$, the maximum value of $`a_{ϵ^{\prime \prime }}^{(\pi ^+\pi ^{})}`$ is approximately given by
$`a_{ϵ^{\prime \prime }}^{(\pi ^+\pi ^{})}|_{max}`$ $``$ $`0.59{\displaystyle \frac{f_\pi ^2}{f_K^2}}R_1^{(max)}\times a_{ϵ^{\prime \prime }}^{(\pi ^{}K^+)}|_{max}`$ (41)
$``$ $`0.40,\text{ (at 1}\sigma \text{ level)}`$ (42)
Where $`R_1`$ is taken to be within the 1$`\sigma `$ standard deviation. The numerical results for $`a_{ϵ^{\prime \prime }}^{(\pi ^+\pi ^{})}`$ and $`a_{ϵ^{\prime \prime }}^{(\pi ^{}K^+)}`$ are plotted in Fig.4 and Fig.5 as functions of the ratios $`R_1`$ and $`R_3`$. It can be seen that for $`R_12.65`$ one has $`|a_{ϵ^{\prime \prime }}^{(\pi ^+\pi ^{})}|>|a_{ϵ^{\prime \prime }}^{(\pi ^{}K^+)}|`$.
In conclusion, assuming $`SU(3)`$ symmetry for the strong phases and ignoring the relative small EWP effects in the $`B\pi ^{}\pi ^+,\pi ^0\pi ^+,\pi ^{}K^+,\pi ^{}\overline{K}^0`$ decays, a model independent approach is proposed to extract the weak phase $`\gamma `$ and the strong final interacting phase $`\delta `$. From the present data a negative $`\mathrm{cos}\gamma \mathrm{cos}\delta `$ is favored. Two solutions for the phases $`\gamma `$ and $`\delta `$ have been obtained at 1$`\sigma `$ level of the current experimental data, though their allowed regions have been strongly restricted, there remain large uncertainties, two interesting allowed regions for the phases $`\gamma `$ and $`\delta `$ have been obtained at the 1$`\sigma `$ level. The numerical values of the phases $`\gamma `$ and $`\delta `$ have been found to lie in the regions: $`35^{}\gamma 62^{}`$ and $`106^{}\delta 180^{}`$ or $`86^{}\gamma 151^{}`$ and $`0^{}\delta 75^{}`$ We would like to point out that with large uncertainties of the current experimental data at more than 1$`\sigma `$ level, one cannot exclude solutions with a small strong phase $`\delta `$ and the values of $`\gamma `$ constrained from $`V_{ub}`$, $`B_{d,s}^0\overline{B}_{d,s}^0`$ mixings and indirect CP-violating parameter $`ϵ_K`$ within the standard model. The direct CP asymmetries $`a_{ϵ^{\prime \prime }}^{(\pi ^+\pi ^{})}`$ and $`a_{ϵ^{\prime \prime }}^{(\pi ^{}K^+)}`$ in $`B\pi ^+\pi ^{}`$ and $`\pi ^{}K^+`$ decays have also be estimated. within the errors of the measurement of $`a_{ϵ^{\prime \prime }}^{(\pi ^{}K^+)}`$, the maximum value of $`a_{ϵ^{\prime \prime }}^{(\pi ^+\pi ^{})}`$ could be as large as 40$`\%`$, a larger value may be possible if $`R_1`$ is large. The more precise experimental data in the $`B\pi \pi `$ and $`\pi K`$ decays will be very plausible for extracting the important weak phase $`\gamma `$ and strong phase $`\delta `$ as well as testing how good of the factorization approach. It may also provide us a possible window for new physics with new CP-violating sources which could change all the constraints arising from the $`B_{d,s}^0\overline{B}_{d,s}^0`$ mixings, radiative rare $`B`$ decays ($`bs\gamma `$) and observed CP-violating parameters in the kaon decays as well as from the decay amplitudes of hadrons. Finally, we would like to address that our current results have been obtained by assuming the SU(3) relations among the strong phases and ignoring the EWP effects in the considered four decay modes. To precisely extract the phases $`\gamma `$ and $`\delta `$, one needs to improve not only the experimental measurements but also theoretical approaches which is going to be investigated elsewhere.
Acknowledgments We would like to thank David E. Jaffe at CLEO for very helpful comments concerning the data for the direct CP violation. This work was supported in part by the NSF of China under grant No. 19625514.
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# References
IITK-HEP-00-01
hep-ph/0002079
HERA Constraint on Warped Quantum Gravity
Prasanta Das <sup>1</sup><sup>1</sup>1E-mail: pdas@iitk.ac.in, Sreerup Raychaudhuri <sup>2</sup><sup>2</sup>2E-mail: sreerup@iitk.ac.in and Saswati Sarkar Department of Physics, Indian Institute of Technology,
Kanpur 208 016, India.
Abstract
We study recent data on deep inelastic $`e^+p`$ scattering at HERA to constrain the parameters of a Randall-Sundrum-type scenario of quantum gravity with a small extra dimension and a non-factorable geometry.
February 2000
Theories with extra dimensions which predict observable consequences at the current high energy accelerators have lately attracted a great deal of interest. Following the original suggestion by Arkani-Hamed, Dimopoulos and Dvali (ADD) , there have been numerous studies in the literature which probe consequences of multiple Kaluza-Klein graviton exchange leading to interactions of electroweak strength. The fact that these theories predict quantum gravity effects at TeV scales has been suggested as a solution of the well-known hierarchy problem in the Standard Model (SM). Though novel and interesting, however, the model suggested by ADD, which is based on a factorable $`𝐑^4\times (𝐒^1)^d`$ geometry, $`d`$ being the number of extra compact dimensions, has the drawback of introducing large compactification radii (amounting to an energy scale as low as $`10^{13}`$ GeV), which effectively introduces a new hierarchy problem. Motivated by this, Randall and Sundrum (RS) have suggested a somewhat different mechanism to solve the hierarchy problem. Instead of writing a factorable metric
$$ds^2=\eta _{\mu \nu }dx^\mu dx^\nu +R_c^2d\varphi _id\varphi _i$$
(1)
where the $`\varphi _i(i=1,d)`$ are extra dimensions compactified with a common radius $`R1`$ mm, they write a non-factorable metric
$$ds^2=e^{𝒦R_c\varphi }\eta _{\mu \nu }dx^\mu dx^\nu +R_c^2d\varphi ^2$$
(2)
involving one extra dimension compactified with a radius $`R_c`$, which is assumed to be marginally greater than the Planck length $`10^{33}`$ cm, and an extra mass scale $`𝒦`$, which is related to the Planck scale $`M_P^{(5)}`$ in the five-dimensional bulk by $`𝒦\left[M_P^{(4)}\right]^2\left[M_P^{(5)}\right]^3`$. Such a ‘warped’ geometry is motivated by compactifying the extra dimension on a $`𝐒^1/𝐙_2`$ orbifold, with two $`D`$-branes at the orbifold fixed points, viz., one at $`\varphi =0`$ (‘Planck brane’ or ‘invisible brane’), and one at $`\varphi =\pi `$ (‘TeV brane’ or ‘visible brane’). It can then be shown that if we assume matter fields to be confined to these $`D`$-branes, one can solve the Einstein equations to obtain a metric of the above form. The interesting physical consequence of this geometry is that any mass scale $``$ on either brane gets scaled by the ‘warp factor $`e^{𝒦R_c\varphi }`$ on either brane. Thus, a mass scale on the Planck brane ($`\varphi =0`$) will remain unchanged, but any mass scale on the TeV brane will be scaled by a factor $`e^{\pi 𝒦R_c}`$. If we assume that the Planck scale is the only fundamental mass scale in the theory, all masses on the TeV brane will be scaled to
$$e^{\pi 𝒦R_c}M_P^{(4)}$$
(3)
It now requires $`𝒦R_c1112`$ to obtain $``$ of the order of the electroweak scale, which justifies the name ‘TeV brane’. Thus, in this model there is no hierarchy problem, since all the independent mass scales are close to the Planck scale. There still remains a minor problem: that of stabilizing the radius $`R_c`$ (which is marginally smaller than the Planck scale) against quantum fluctuations, but this is not so severe as in the model of ADD, where the compactification radius needs to be stabilized over as many as 30 orders of magnitude. A simple extension of the RS construction involving an extra bulk scalar field has been proposed to stabilize $`R_c`$ and this predicts light radion excitations with possible collider signatures . However, as these will not contribute to the processes of interest in this letter, this idea will not be discussed further. On the flip side, it is not as simple to embed the RS construction within the framework of string theories as it is for the ADD case. However, a first attempt has been made , and it may be hoped that future work will achieve this highly desirable goal.
Following the ingenious suggestion of a non-factorable geometry, the mass spectrum and couplings of the graviton in the RS model have been worked out, in Refs. . We do not describe the details of this calculation, but refer the reader to the original literature. It is worth noting that there are strong phenomenological constraints on bulk excitations of the SM fields . It suffices here to note that the effective Lagrangian density for graviton interactions on the TeV brane (which we identify with the observable world) has the form
$$_{eff}^{RS}=\frac{1}{\overline{M}_P}h_{\mu \nu }^0(x)T^{\mu \nu }(x)\frac{e^{\pi 𝒦R_c}}{\overline{M}_P}\underset{n=1}{\overset{\mathrm{}}{}}h_{\mu \nu }^n(x)T^{\mu \nu }(x)$$
(4)
where $`\overline{M}_PM_P^{(4)}/\sqrt{8\pi }`$ is the reduced Planck mass and the $`h_{\mu \nu }^n(x)`$ correspond to the Kaluza-Klein expansion of the massless graviton in five dimensions
$$h_{\mu \nu }(x,\varphi )=\underset{n=0}{\overset{\mathrm{}}{}}h_{\mu \nu }^n(x)\frac{\chi ^n(\varphi )}{\sqrt{R_c}}.$$
(5)
Equation (4) tells us that the massless Kaluza-Klein (KK) mode effectively decouples from ordinary matter since its interactions are suppressed by the Planck mass. On the other hand, the massive KK modes couple as the inverse of the Planck mass, scaled by $`e^{\pi 𝒦R_c}`$, which is an electroweak-strength interaction. Feynman rules to the lowest order for these modes, assuming a coupling $`\overline{M}_P^1`$ to all the modes have been worked out in Refs. ) and in the context of ADD-like scenarios. All we need to do to get the corresponding Feynman rules in the RS model is to multiply the couplings by the warp factor $`e^{\pi 𝒦R_c}`$ where necessary.
As shown in Ref. , the orbifold geometry forces the Fourier coefficients $`\chi ^n(\varphi )`$ to satisfy a Bessel equation, whence it may be shown that they are given by a linear combination of the Bessel and Neumann functions of order 2. The requirement that the first derivative of $`\chi ^n(\varphi )`$ be continuous at the orbifold fixed point $`\varphi =\pi `$ then requires $`J_1(x_n)=0`$. Using this, the masses $`M_n`$ of the graviton states can be written in terms of the zeros of the Bessel function of order unity as
$$M_n=x_n𝒦e^{\pi 𝒦R_c}x_nm_o$$
(6)
where $`m_0`$ sets the scale of graviton masses and is essentially a free parameter of the theory. It is also convenient to write
$$\frac{e^{\pi 𝒦R_c}}{\overline{M}_P}=\frac{c_0}{m_0}\sqrt{8\pi }$$
(7)
using (6) and introducing another undetermined parameter $`c_0𝒦/M_P^{(4)}`$. Ref. points out that ($`m_0,c_0`$) may conveniently be taken as the free parameters of the theory, and we follow their prescription in our work.
Though $`c_0`$ and $`m_0`$ are not precisely known, one can make estimates of their magnitude using theoretical ideas and phenomenological inputs. We note that the RS construction requires $`𝒦`$ to be at least an order of magnitude less than $`M_P^{(4)}`$, because $`𝒦^1`$ sets the scale for the curvature of the fifth dimension, and should therefore be large compared with the Planck length. The latter is necessitated by the requirement that fluctuations in the bulk gravitational field in the vicinity of the $`D`$-branes be small. The range of interest for $`c_0`$ is, therefore, about 0.01 to 0.1, the lower value being determined by naturalness considerations. Regarding $`m_0`$, Eq. (6) tells us that it is reduced from the scale $`𝒦`$ by the factor $`e^{\pi 𝒦R_c}`$. In the RS construction, one requires $`𝒦R_c11`$–12, which reduces $`m_0`$ to the electroweak scale. Hence, we may consider $`m_0`$ in the range of a few tens of GeV to a few TeV. Eq. (6) also tells us that the first massive graviton lies at $`M_1=x_1m_03.83m_0`$. Since no graviton resonances have been seen at LEP-2, running at energies upto 200 GeV, it is clear that we should expect $`m_0>52`$ GeV.
In this letter we report on a study of graviton effects, within the RS model, on $`e^+p`$ deep inelastic scattering (DIS) at HERA. At the leading order, there are two extra Feynman diagrams contributing to $`e^+pe^++X`$. One of these involves a $`t`$-channel exchange of a virtual (massive) graviton between the $`e^+`$ and a quark; the other involves a $`t`$-channel exchange of a virtual (massive) graviton between the $`e^+`$ and a gluon. The first one adds coherently with the corresponding SM diagrams with photon and $`Z`$-boson exchange; the second one has no SM analogue and hence adds incoherently. However, at HERA energies, we do not expect much contribution from the gluon-induced diagram because of the low gluon flux.
The cross-section for the above processes has been calculated for the case of the ADD model in Ref. and can be easily translated to the RS model using the replacement
$$\frac{\lambda }{M_S^4}\frac{8\pi c_0^2}{m_0^2}\underset{n}{}\frac{1}{|\widehat{t}|+M_n^2}$$
(8)
We have developed approximate analytic formulae for this sum, using the well-known properties of the zeros of the Bessel function $`J_1(x)`$. These will be presented elsewhere . Using this, we incorporate the calculated theoretical cross-section into a parton-level Monte Carlo event generator, with two free parameters, viz. the graviton mass scale $`m_0`$, and the coupling parameter $`c_0`$. Finally the simulation results are compared with data from the ZEUS Collaboration to constrain the ($`m_0`$$`c_0`$) plane.
Our numerical studies are founded on the latest results presented by the ZEUS Collaboration, which are based on 47.7 pb<sup>-1</sup> of data collected over the period 1994-1997. The ZEUS Collaboration uses the double-angle (DA) method to determine the DIS variables. In this method, one measures the polar angle $`\theta _e`$ of the scattered positron, and reconstructs the polar angle $`\gamma _h`$ of the struck quark in the naive parton model using all hadronic clusters which can be identified with a jet having the requisite $`p_T`$ balance with the positron. In terms of these observables and the energy $`E_e(E_p)`$ of the initial positron (proton) beam, one can reconstruct the standard DIS variables as
$`Q_{DA}^2`$ $`=`$ $`4E_e^2{\displaystyle \frac{\mathrm{sin}\gamma _h(1+\mathrm{cos}\theta _e)}{\mathrm{sin}\gamma _h+\mathrm{sin}\theta _e\mathrm{sin}(\gamma _h+\theta _e)}}`$ (9)
$`x_{DA}`$ $`=`$ $`{\displaystyle \frac{E_e}{E_p}}{\displaystyle \frac{\mathrm{sin}\gamma _h+\mathrm{sin}\theta _e+\mathrm{sin}(\gamma _h+\theta _e)}{\mathrm{sin}\gamma _h+\mathrm{sin}\theta _e\mathrm{sin}(\gamma _h+\theta _e)}}`$ (10)
$`y_{DA}`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}\theta _e(1\mathrm{cos}\gamma _h)}{\mathrm{sin}\gamma _h+\mathrm{sin}\theta _e\mathrm{sin}(\gamma _h+\theta _e)}}`$ (11)
Various triggers, acceptances and selection cuts have been used by the ZEUS Collaboration, of which we need to impose only the following in a parton-level analysis:
* If the final state positron has polar angle greater than $`17.2^0`$, it must have a total energy greater than 10 GeV.
* If the final state positron has polar angle less than $`17.2^0`$, it must have a transverse momentum greater than 30 GeV.
Of the DIS variables listed above, it is well-known that it is the first, namely $`Q_{DA}^2`$, which exhibits maximum sensitivity to most kinds of new physics. The ZEUS Collaboration has presented their data for 20 bins in $`Q_{DA}^2`$, ranging from $`Q_{DA}^2`$ = 400 GeV<sup>2</sup> to 51200 GeV<sup>2</sup>. The Born-level cross-section in each bin, obtained by suitably scaling out radiative effects, has been presented by the ZEUS Collaboration together with the SM expectation. We have checked that the latter, obtained using hadronization procedures incorporated in the standard HERACLES and ARIADNE program packages, are in excellent agreement (within a few per cent) with our parton-level analysis. Any small differences which persist can be removed by calibrating the cross-section binwise, so as to yield the actual ZEUS expectations. This procedure has the added merit of taking care of residual higher order effects such as initial state radiation, which should be roughly the same in the SM as in the case when the graviton exchanges are included.
In Fig. 1, we present a graph showing the variation in the $`Q_{DA}^2`$ distribution with the mass scale $`m_0`$ in the RS model. For this graph, we have plotted the ratio
$$R(m_0,c_0)\frac{d\sigma _{RS}/dQ_{DA}^2}{d\sigma _{SM}/dQ_{DA}^2}$$
(12)
of the cross-section predicted in the RS model with the prediction of the SM, for $`c_0`$ = 0.1 and $`m_0=80,90`$ and 100 GeV, together with the ZEUS data. It may be seen that the cross-section in the RS model (like the ADD model ) exhibits large deviations in the highest $`Q^2`$ bins. This, of course, rapidly approaches the SM (dotted line) if $`c_0`$ is chosen smaller. Interestingly, the RS model predictions seem to show a slight diminution for intermediate values of $`Q_{DA}^2`$, which are intriguingly like the trend shown by the data. However, the experimental errors are too large to enable us to attach any significance to this circumstance. Accordingly, we take the conservative viewpoint that the data fit the SM very well and can be used to constrain new physics.
Figure 1. Illustrating the ratio $`R(m_0,c_0)`$ of the $`Q_{DA}^2`$ distribution in the RS model to that in the SM (see Eq. 12), for $`c_0=0.1`$ and $`m_0=80,90`$ and 100 GeV. The dotted line corresponds to the SM. The ZEUS data are also shown.
Once the above simulation is set-up, we estimate the binwise cross-section $`\chi ^2(m_0,c_0)`$ for each value of $`m_0`$ and $`c_0`$ and use this distribution to calculate
$$\chi ^2(m_0,c_0)=\underset{i=1}{\overset{20}{}}\frac{\left[\sigma _i(m_0,c_0)\sigma _i^{(C.V.)}\right]^2}{ϵ_i^2}$$
(13)
where
$$ϵ_i=ϵ_1^i\theta [\sigma _i(m_0,c_0)\sigma _i^{(C.V.)}]+ϵ_2^i\theta [\sigma _i^{(C.V.)}\sigma _i(m_0,c_0)]$$
(14)
assuming that the experimental value in the $`i`$-th bin is given by $`[\sigma _i^{(C.V.)}]_{ϵ_2^i}^{+ϵ_1^i}`$. In this, it is assumed that $`ϵ_1^i`$ and $`ϵ_2^i`$ contain the statistical and systematic errors added in quadrature. The 95% C.L. bound is then obtained by requiring $`\chi ^2(m_0,c_0)<31.41`$, which is the expectation from random fluctuations.
Figure 2. Illustrating the constraint on the parameter space of the RS model arising from an analysis of ZEUS high-$`Q^2`$ data. The shaded region is ruled out at the 95% C.L. level.
In Figure 2, we show the 95% C.L. constraints on the $`m_0`$$`c_0`$ plane using the above technique. Since the effective graviton coupling is quadratic in $`c_0`$ we expect the cross-section to rise as $`c_0`$ increases — this is reflected in the fact that Figure 2 shows upper bounds on $`c_0`$. On the other hand, the $`m_0`$ dependence of the cross-section is very complicated, because of the summation over the KK states. However, as the figure makes clear, there is a sharp drop in the cross-section as $`m_0`$ increases, so that the ZEUS data become insensitive to the new physics beyond about $`m_0=120`$ GeV. This corresponds to $`M_1460`$ GeV, a value which is still not accessible to the generation of colliders running at present.
As the above figure and discussion makes clear, HERA data as presented by the ZEUS Collaboration provide somewhat modest, but nevertheless interesting constraints on the parameter space of the RS model of quantum gravity. Since gravitons couple to the energy-momentum tensor of the matter fields, one may expect considerable improvements in these results at machines running at higher energies, such as the LHC, the proposed NLC and possible muon colliders. In particular, it would be interesting to see if these machines could actually find graviton resonances, which one expects in the RS model , but not in the ADD theory. We have performed a preliminary study of the RS model in the light of HERA data, and we expect that the future will see many more detailed studies of this very interesting scenario.
The authors have benefited greatly from discussions with Gour Bhattacharya, Dilip K. Ghosh, Sudipta Mukherji, Gautam Sengupta, K. Sridhar and Ajit M. Srivastava. SS acknowledges financial support from the Council of Scientific and Industrial Research (Award No. 9/92(185)/95-EMR-I), Government of India.
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# 1 Introduction
## 1 Introduction
In this paper, we discuss two examples where non-trivial subsets of the scalar sectors of gauged supergravities are obtained by spherical reduction from a higher dimension. The first example is the embedding of the scalars in the $`U(1)^4`$ maximal abelian truncation of $`SO(8)`$ gauged $`N=8`$ maximal supergravity in $`D=4`$, arising from $`D=11`$ via compactification on $`S^7`$. The consistency of the full $`SO(8)`$ reduction on $`S^7`$ was proven in , although at a somewhat implicit level. The $`N=2`$ truncation includes a total of six scalar fields, comprising three dilaton/axion pairs. In terms of the original $`SO(8)`$ representations of the full theory, where there are 35 scalars in the $`35_v`$, and 35 pseudoscalars in the $`35_c`$ of $`SO(8)`$, the three dilatons come from the $`35_v`$, and the three axions come from the $`35_c`$. In , a further simplifying truncation was performed, in which the three axions were set to zero. The reduction Ansatz becomes considerably more complicated when axions are included, as was already seen in the case of the single dilaton/axion pair of the $`N=4`$ gauged $`SO(4)`$ truncation, discussed in . In the present example, the inclusion of the three axions as well as the three dilatons leads to a considerably more complicated structure in the reduction Ansatz.
The second example is a truncation of the $`SO(6)`$ gauged $`N=8`$ maximal supergravity in $`D=5`$, arising from type IIB via compactification on $`S^5`$. In this case there are 42 spin-0 fields in total, comprising 20 scalars in the $`20^{}`$ of $`SO(6)`$, 20 pseudoscalars in the $`10`$ and $`\overline{10}`$, and two singlets corresponding to the original dilaton and axion of the type IIB theory. The truncation we shall consider retains two spin-0 fields, comprising one scalar from the $`20^{}`$, and one pseudoscalar from the $`10`$ and $`\overline{10}`$. This particular truncation is of interest because it is large enough to include the fields that participate in two distinct supersymmetric vacua of the $`D=5`$ gauged theory , one with maximal $`N=8`$, $`SO(6)`$ symmetry, and the other with $`N=2`$, $`SU(1)\times U(1)`$ symmetry. Although an explicit interpolating solution is not known it is in principle describable within the truncation we are making.
In both of our examples, we shall concentrate on elucidating the geometrical structure of the embedding in $`D=11`$ or type IIB supergravity. Specifically, we shall concentrate on the Ansatz for the Kaluza-Klein reduction of the metric tensor. Strictly speaking, one can only be sure that the reduction is fully consistent with all the equations of motion of the higher-dimensional theory if one has the complete Ansatz for all the higher-dimensional fields, including the antisymmetric tensor field strengths. (Or, alternatively, if an “existence proof” for the consistency of the reduction Ansatz has independently been constructed.) Obtaining the Ansatz for the antisymmetric tensor fields is notoriously difficult, and we shall not complete this part of the analysis in this paper. In the case of our $`D=4`$ example, we can appeal to the results of , in which a complete proof of the consistency of the $`S^7`$ reduction is exhibited. In principle it allows one to read off the Ansatz for the 4-form field strength, although only an implicit procedure for its construction is presented. On the other hand, the general Anstaz for the metric tensor is rather explicit, and it is by making use of this expression that we are able to obtain the $`D=4`$ results in this paper. These results can be used in order to study the eleven-dimensional geometrical structure of general domain-wall solutions in $`D=4`$ supported both by the three dilatonic scalars and also the three accompanying axions. Such solutions can be constructed from the purely dilatonic ones by means of $`SL(2,R)`$ transformations.
In $`D=5`$ the situation is less clear, since no proof for the consistency of the full $`S^5`$ reduction to $`SO(6)`$ gauged $`N=8`$ maximal supergravity currently exists. A conjecture for the metric reduction Ansatz appears in , which is closely analogous to the known construction in $`D=4`$ given in , and it is this that we use in order to obtain an explicit expression for the metric embedding for our 2-scalar truncation. Again, the complexities of the antisymmetric tensor embedding have prevented us from obtaining a full non-linear result in that sector. Thus the status of our $`D=5`$ embedding is that, subject to the assumption of an ultimate consistency of the $`S^5`$ reduction scheme,<sup>1</sup><sup>1</sup>1Further evidence for the consistency of the $`S^5`$ reduction was obtained in , where certain scalar plus gravity truncations in Kaluza-Klein sphere reductions were proved to be consistent. Additionally, the complete consistent reductions of $`D=11`$ supergravity on $`S^4`$ and massive type IIA supergravity on $`S^4`$ have been constructed. Recently, more evidence for the consistency of the $`S^5`$ reduction was presented in . and subject to the assumption that the conjecture for the metric Ansatz in is correct, then our explicit results for the 2-scalar metric Ansatz is valid. In principle, our result can then be used to study the geometry of the RG flow describing the transition between the two supersymmetric extrema of the associated scalar potential.
## 2 $`N=2`$ $`U(1)^4`$ Gauged Supergravity in $`D=4`$ From $`D=11`$
### 2.1 The Three Dilaton/Axion Pairs in $`D=4`$
The $`35_v+35_c`$ of spin-0 fields in $`SO(8)`$ gauged supergravity in $`D=4`$ are described in terms of a 56-vielbein $`𝒱`$, with the block-diagonal form
$$𝒱=\left(\begin{array}{cc}u_{ij}^{IJ}& v_{ijKL}\\ v^{k\mathrm{}IJ}& u^k\mathrm{}_{KL}\end{array}\right),$$
(1)
which transforms under local $`SU(8)`$ and rigid $`E_7`$ . In terms of the quantities $`u_{ij}^{IJ}`$ and $`v_{ijKL}`$, it was shown in (having been previously proposed in ) that the Ansatz for the inverse of the internal $`S^7`$ compactifying metric is
$$\widehat{g}^{mn}(x,y)\widehat{\mathrm{\Delta }}^1g^{mn}(x,y)=\frac{1}{2}(K^{mIJ}K^{nKL}+K^{nIJ}K^{mKL})(u_{ij}{}_{}{}^{IJ}+v_{ijIJ})(u^{ij}{}_{KL}{}^{}+v^{ijKL}),$$
(2)
where $`K^{mIJ}`$ are the 28 Killing vectors on the round $`S^7`$, and
$$\widehat{\mathrm{\Delta }}^2=\frac{det(g_{mn}(x,y))}{det(g_{mn}(y))},$$
(3)
where $`g_{mn}(x,y)`$ is the inverse of $`g^{mn}(x,y)`$, and $`g_{mn}(y)`$ is $`g_{mn}(x,y)`$ with the scalar fields all set to zero, so that it becomes the round $`S^7`$ metric. The eleven-dimensional metric Ansatz will be given by
$$d\widehat{s}_{11}^2=\widehat{\mathrm{\Delta }}^1ds_4^2+g_{mn}(x,y)dy^mdy^n=\widehat{\mathrm{\Delta }}^1(ds_4^2+\widehat{g}_{mn}(x,y)dy^mdy^n),$$
(4)
where $`\widehat{g}_{mn}(x,y)=\widehat{\mathrm{\Delta }}g_{mn(x,y)}`$ is the inverse of $`\widehat{g}^{mn}(x,y)`$.<sup>2</sup><sup>2</sup>2For now, we shall leave out the Kaluza-Klein gauge fields from the construction of the metric. As discussed in , the truncation to three dilaton/axion pairs is naturally accompanied by the four $`U(1)`$ gauge fields of the maximal abelian $`U(1)^4`$ subgroup of $`SO(8)`$. These gauge fields are easily incorporated in the Kaluza-Klein Ansatz, and we shall add them in at the end of the derivation. We shall also set the gauge coupling constant $`g`$ equal to 1 for now, and restore it later.
We use the parameterisation of the $`u_{ij}^{IJ}`$ and $`v_{ijKL}`$ matrices described in . In particular, we introduce three scalars $`\lambda _i`$, and three associated pseudoscalars $`\sigma _i`$, whose kinetic Lagrangian is
$$=\frac{1}{2}\underset{i}{}\left((\lambda _i)^2+\mathrm{sinh}^2\lambda _i(\sigma _i)^2\right).$$
(5)
To shorten the subsequent formulae, we make the following definitions:
$$c_i\mathrm{cosh}\lambda _i,s_i\mathrm{sinh}\lambda _i.$$
(6)
Also, for future convenience, we introduce the “standard” dilaton/axion pairs $`(\phi _i,\chi _i)`$, related to $`(\lambda _i,\sigma _i)`$ by
$`\mathrm{cosh}\lambda _i`$ $`=`$ $`\mathrm{cosh}\phi _i+\frac{1}{2}\chi _i^2e^{\phi _i},`$
$`\mathrm{cos}\sigma _i\mathrm{sinh}\lambda _i`$ $`=`$ $`\mathrm{sinh}\phi _i\frac{1}{2}\chi _i^2e^{\phi _i},`$ (7)
$`\mathrm{sin}\sigma _i\mathrm{sinh}\lambda _i`$ $`=`$ $`\chi _ie^{\phi _i}.`$
In terms of these fields, the scalar kinetic terms are
$$=\frac{1}{2}\underset{i}{}\left((\phi _i)^2+e^{2\phi _i}(\chi _i)^2\right).$$
(8)
After some algebra, we find that $`u_{ij}^{IJ}`$ and $`v_{ijKL}`$ are given by
$`\frac{1}{4}u_{ij}{}_{}{}^{KL}P_{ij}^{}Q_{KL}=c_1(P_{a_1a_2}Q_{a_1a_2}+P_{a_3a_4}Q_{a_3a_4})`$ (9)
$`+c_2(P_{a_1a_3}Q_{a_1a_3}+P_{a_2a_4}Q_{a_2a_4})+c_3(P_{a_1a_4}Q_{a_1a_4}+P_{a_2a_3}Q_{a_2a_3})`$
$`+P_{12}(c_1c_2c_3Q_{12}+c_1s_2s_3e^{\mathrm{i}(\sigma _2+\sigma _3)}Q_{34}+c_2s_1s_3e^{\mathrm{i}(\sigma _1+\sigma _3)}Q_{56}+c_3s_1s_2e^{\mathrm{i}(\sigma _1+\sigma _2)}Q_{78})`$
$`+P_{34}(c_1c_2c_3Q_{34}+c_1s_2s_3e^{\mathrm{i}(\sigma _2+\sigma _3)}Q_{12}+c_2s_1s_3e^{\mathrm{i}(\sigma _1\sigma _3)}Q_{78}+c_3s_1s_2e^{\mathrm{i}(\sigma _1\sigma _2)}Q_{56})`$
$`+P_{56}(c_1c_2c_3Q_{56}+c_1s_2s_3e^{\mathrm{i}(\sigma _2\sigma _3)}Q_{78}+c_2s_1s_3e^{\mathrm{i}(\sigma _1+\sigma _3)}Q_{12}+c_3s_1s_2e^{\mathrm{i}(\sigma _1+\sigma _2)}Q_{34})`$
$`+P_{78}(c_1c_2c_3Q_{78}+c_1s_2s_3e^{\mathrm{i}(\sigma _2+\sigma _3)}Q_{56}+c_2s_1s_3e^{\mathrm{i}(\sigma _1+\sigma _3)}Q_{34}+c_3s_1s_2e^{\mathrm{i}(\sigma _1+\sigma _2)}Q_{12})`$
$`\frac{1}{4}v_{ijKL}P_{ij}Q_{KL}=`$ (10)
$`s_1(e^{\mathrm{i}\sigma _1}ϵ^{a_1b_1}ϵ^{a_2b_2}P_{a_1a_2}Q_{b_1b_2}+e^{\mathrm{i}\sigma _1}ϵ^{a_3b_3}ϵ^{a_4b_4}P_{a_3a_4}Q_{b_3b_4})`$
$`s_2(e^{\mathrm{i}\sigma _2}ϵ^{a_1b_1}ϵ^{a_3b_3}P_{a_1a_3}Q_{b_1b_3}+e^{\mathrm{i}\sigma _2}ϵ^{a_2b_2}ϵ^{a_4b_4}P_{a_2a_4}Q_{b_2b_4})`$
$`s_3(e^{\mathrm{i}\sigma _3}ϵ^{a_1b_1}ϵ^{a_4b_4}P_{a_1a_4}Q_{b_1b_4}+e^{\mathrm{i}\sigma _3}ϵ^{a_2b_2}ϵ^{a_3b_3}P_{a_2a_3}Q_{b_2b_3})`$
$`+P_{12}(s_1s_2s_3e^{\mathrm{i}(\sigma _1+\sigma _2+\sigma _3)}Q_{12}+s_1c_2c_3e^{\mathrm{i}\sigma _1}Q_{34}+s_2c_1c_3e^{\mathrm{i}\sigma _2}Q_{56}+s_3c_1c_2e^{\mathrm{i}\sigma _3}Q_{78})`$
$`+P_{34}(s_1s_2s_3e^{\mathrm{i}(\sigma _1\sigma _2\sigma _3)}Q_{34}+s_1c_2c_3e^{\mathrm{i}\sigma _1}Q_{12}+s_2c_1c_3e^{\mathrm{i}\sigma _2}Q_{78}+s_3c_1c_2e^{\mathrm{i}\sigma _3}Q_{56})`$
$`+P_{56}(s_1s_2s_3e^{\mathrm{i}(\sigma _1+\sigma _2\sigma _3)}Q_{56}+s_1c_2c_3e^{\mathrm{i}\sigma _1}Q_{78}+s_2c_1c_3e^{\mathrm{i}\sigma _2}Q_{12}+s_3c_1c_2e^{\mathrm{i}\sigma _3}Q_{34})`$
$`+P_{78}(s_1s_2s_3e^{\mathrm{i}(\sigma _1\sigma _2+\sigma _3)}Q_{78}+s_1c_2c_3e^{\mathrm{i}\sigma _1}Q_{56}+s_2c_1c_3e^{\mathrm{i}\sigma _2}Q_{34}+s_3c_1c_2e^{\mathrm{i}\sigma _3}Q_{12}).`$
Here, we have introduced $`P`$ and $`Q`$ simply as arbitrary antisymmetric tensors, in order to provide a compact way of summarising all the components of the $`u_{ij}^{IJ}`$ and $`v_{ijKL}`$ matrices. The index notation is as follows. Indices with a “1” subscript, such as $`a_1`$, range over the values $`(1,2)`$; similarly $`a_2`$ ranges over $`(3,4)`$, $`a_3`$ ranges over $`(5,6)`$ and $`a_4`$ ranges over $`(7,8)`$.
Next, we substitute these results into the Ansatz (2) for the inverse $`S^7`$ metric. It is advantageous to introduce a new parameterisation for the dilaton/axion pairs, as follows:
$$Y_ie^{\frac{1}{2}\phi _i},\stackrel{~}{Y}_i(1+\chi _i^2Y_i^4)^{\frac{1}{2}}Y_i^1,b_i\chi _iY_i^2,$$
(11)
and so
$`\mathrm{cosh}\lambda _i`$ $`=`$ $`\frac{1}{2}(Y_i^2+\stackrel{~}{Y}_i^2),`$
$`\mathrm{cos}\sigma _i\mathrm{sinh}\lambda _i`$ $`=`$ $`\frac{1}{2}(Y_i^2\stackrel{~}{Y}_i^2),`$ (12)
$`\mathrm{sin}\sigma _i\mathrm{sinh}\lambda _i`$ $`=`$ $`b_i.`$
It is also advantageous to redefine the $`SO(8)`$ basis relative to the one we have used so far. The action of transformation, which amounts to a triality rotation under which $`K_{ij}\frac{1}{2}(\mathrm{\Gamma }_{ij})^k\mathrm{}K_k\mathrm{}`$, is given explicitly in Appendix A. After doing this, we find that the inverse internal metric (2) takes the form<sup>3</sup><sup>3</sup>3The notation for writing the inverse metric is $`_s^2g^{mn}_m_n`$. The derivatives do not act on any other objects here; it is just a convenient way of writing all the components of $`g^{mn}`$ in one formula, exactly analogous to writing the downstairs metric as $`ds^2=g_{mn}dy^mdy^n`$. For example, the inverse of the 2-sphere metric $`ds^2=d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2`$ is written as
$$_s^2=_\theta ^2+\frac{1}{\mathrm{sin}^2\theta }_\varphi ^2.$$
$`\widehat{}_s^2\widehat{g}^{mn}_m_n`$ $`=`$ $`Y_1^2(K_{13}^2+K_{14}^2+K_{23}^2+K_{24}^2)+\stackrel{~}{Y}_1^2(K_{57}^2+K_{58}^2+K_{67}^2+K_{68}^2)`$ (13)
$`+Y_2^2(K_{15}^2+K_{16}^2+K_{25}^2+K_{26}^2)+\stackrel{~}{Y}_2^2(K_{37}^2+K_{38}^2+K_{47}^2+K_{48}^2)`$
$`+Y_3^2(K_{17}^2+K_{18}^2+K_{27}^2+K_{28}^2)+\stackrel{~}{Y}_3^2(K_{35}^2+K_{36}^2+K_{45}^2+K_{46}^2)`$
$`+Y_1^2Y_2^2Y_3^2K_{12}^2+Y_1^2\stackrel{~}{Y}_2^2\stackrel{~}{Y}_3^2K_{34}^2+\stackrel{~}{Y}_1^2Y_2^2\stackrel{~}{Y}_3^2K_{56}^2+\stackrel{~}{Y}_1^2\stackrel{~}{Y}_2^2Y_3^2K_{78}^2`$
$`2b_2b_3(Y_1^2K_{12}K_{34}\stackrel{~}{Y}_1^2K_{56}K_{78})`$
$`2b_1b_3(Y_2^2K_{12}K_{56}\stackrel{~}{Y}_2^2K_{34}K_{78})`$
$`2b_1b_2(Y_3^2K_{12}K_{78}\stackrel{~}{Y}_3^2K_{34}K_{56}).`$
In order to proceed further, it is useful to look at the geometry of the 7-sphere in some detail. Some useful results on this topic are collected in Appendix B.
### 2.2 The Metric Ansatz for the three dilaton/axion pairs
From the results in Appendix B, it follows that the inverse metric (13) for the system with 3 dilatons and 3 axions is a direct sum of a $`4\times 4`$ part involving the $`_{\varphi _i}`$ basis vectors, and a $`3\times 3`$ part involving the $`_{\mu _i}`$ basis vectors (which are constrained by the fact that $`\mu _i\mu _i=1`$):
$$\widehat{}_s^2=\widehat{}_{s_4}^2+\widehat{}_{s_3}^2.$$
(14)
For the $`4\times 4`$ inverse metric, we find
$`\widehat{}_{s_4}^2`$ $`=`$ $`{\displaystyle \underset{i}{}}\mu _i^2Q_i_{\varphi _i}^22b_2b_3(Y_1^2_{\varphi _1}_{\varphi _2}\stackrel{~}{Y}_1^2_{\varphi _3}_{\varphi _4})`$ (15)
$`2b_1b_3(Y_2^2_{\varphi _1}_{\varphi _3}\stackrel{~}{Y}_2^2_{\varphi _2}_{\varphi _4})2b_1b_2(Y_3^2_{\varphi _1}_{\varphi _4}\stackrel{~}{Y}_3^2_{\varphi _2}_{\varphi _3}),`$
where
$`Q_1`$ $`=`$ $`Y_1^2Y_2^2Y_3^2\mu _1^2+Y_1^2\mu _2^2+Y_2^2\mu _3^2+Y_3^2\mu _4^2,`$
$`Q_2`$ $`=`$ $`Y_1^2\stackrel{~}{Y}_2^2\stackrel{~}{Y}_3^2\mu _2^2+Y_1^2\mu _1^2+\stackrel{~}{Y}_3^2\mu _3^2+\stackrel{~}{Y}_2^2\mu _4^2,`$
$`Q_3`$ $`=`$ $`Y_2^2\stackrel{~}{Y}_1^2\stackrel{~}{Y}_3^2\mu _3^2+Y_2^2\mu _1^2+\stackrel{~}{Y}_3^2\mu _2^2+\stackrel{~}{Y}_1^2\mu _4^2,`$
$`Q_4`$ $`=`$ $`Y_3^2\stackrel{~}{Y}_1^2\stackrel{~}{Y}_2^2\mu _4^2+Y_3^2\mu _1^2+\stackrel{~}{Y}_2^2\mu _2^2+\stackrel{~}{Y}_1^2\mu _3^2.`$ (16)
For the $`3\times 3`$ part, we find
$`\widehat{}_{s_3}^2`$ $`=`$ $`Y_1^2(\mu _1_{\mu _2}\mu _2_{\mu _1})^2+Y_2^2(\mu _1_{\mu _3}\mu _3_{\mu _1})^2+Y_3^2(\mu _1_{\mu _4}\mu _4_{\mu _1})^2`$
$`+\stackrel{~}{Y}_1^2(\mu _3_{\mu _4}\mu _4_{\mu _3})^2+\stackrel{~}{Y}_2^2(\mu _2_{\mu _4}\mu _4_{\mu _2})^2+\stackrel{~}{Y}_3^2(\mu _2_{\mu _3}\mu _3_{\mu _2})^2,`$
where $`\mu _i\mu _i=1`$.
Because of the block-diagonal structure, we can invert the two parts separately. For the $`4\times 4`$ part, we straightforwardly invert the inverse metric to obtain
$`d\widehat{s}_4^2`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Xi }}}[{\displaystyle \underset{i}{}}\mu _i^2Z_id\varphi _i^2+2b_2b_3(\mu _1^2\mu _2^2d\varphi _1d\varphi _2\mu _3^2\mu _4^2d\varphi _3d\varphi _4)`$
$`+2b_1b_3(\mu _1^2\mu _3^2d\varphi _1d\varphi _3\mu _2^2\mu _4^2d\varphi _2d\varphi _4)`$
$`+2b_1b_2(\mu _1^2\mu _4^2d\varphi _1d\varphi _4\mu _2^2\mu _3^2d\varphi _2d\varphi _3)],`$
where
$`Z_1`$ $`=`$ $`\mu _1^2+\stackrel{~}{Y}_2^2\stackrel{~}{Y}_3^2\mu _2^2+\stackrel{~}{Y}_1^2\stackrel{~}{Y}_3^2\mu _3^2+\stackrel{~}{Y}_1^2\stackrel{~}{Y}_2^2\mu _4^2,`$
$`Z_2`$ $`=`$ $`\mu _2^2+Y_2^2Y_3^2\mu _1^2+\stackrel{~}{Y}_1^2Y_2^2\mu _3^2+\stackrel{~}{Y}_1^2Y_3^2\mu _4^2,`$
$`Z_3`$ $`=`$ $`\mu _3^2+Y_1^2Y_3^2\mu _1^2+Y_1^2\stackrel{~}{Y}_2^2\mu _2^2+Y_3^2\stackrel{~}{Y}_2^2\mu _4^2,`$
$`Z_4`$ $`=`$ $`\mu _4^2+Y_1^2Y_2^2\mu _1^2+Y_1^2\stackrel{~}{Y}_3^2\mu _2^2+Y_2^2\stackrel{~}{Y}_3^2\mu _3^2.`$ (19)
The function $`\mathrm{\Xi }`$ is given by
$`\mathrm{\Xi }`$ $`=`$ $`Y_1^2Y_2^2Y_3^2\mu _1^4+Y_1^2\stackrel{~}{Y}_2^2\stackrel{~}{Y}_3^2\mu _2^4+\stackrel{~}{Y}_1^2Y_2^2\stackrel{~}{Y}_3^2\mu _3^4+\stackrel{~}{Y}_1^2\stackrel{~}{Y}_2^2Y_3^2\mu _4^4`$ (20)
$`+(Y_2^2\stackrel{~}{Y}_2^2+Y_3^2\stackrel{~}{Y}_3^2)(Y_1^2\mu _1^2\mu _2^2+\stackrel{~}{Y}_1^2\mu _3^2\mu _4^2)`$
$`+(Y_1^2\stackrel{~}{Y}_1^2+Y_3^2\stackrel{~}{Y}_3^2)(Y_2^2\mu _1^2\mu _3^2+\stackrel{~}{Y}_2^2\mu _2^2\mu _4^2)`$
$`+(Y_1^2\stackrel{~}{Y}_1^2+Y_2^2\stackrel{~}{Y}_2^2)(Y_3^2\mu _1^2\mu _4^2+\stackrel{~}{Y}_3^2\mu _2^2\mu _3^2).`$
There remains the problem of inverting the $`3\times 3`$ part $`\widehat{}_{s_3}^2`$ of the inverse metric. Since we know the inverse metric in the form (2.2), expressed in terms of the four $`_{\mu _i}`$ basis vectors formed from the the constrained $`\mu _i`$, it is helpul first to solve the constraint $`\mu _i\mu _i=1`$ explicitly, by introducing three angular coordinates as follows:
$$\mu _1=c\mathrm{cos}\frac{1}{2}\theta ,\mu _2=c\mathrm{sin}\frac{1}{2}\theta ,\mu _3=s\mathrm{cos}\frac{1}{2}\stackrel{~}{\theta },\mu _4=s\mathrm{sin}\frac{1}{2}\stackrel{~}{\theta },$$
(21)
where $`c=\mathrm{cos}\xi `$, $`s=\mathrm{sin}\xi `$. It then follows that
$`_\theta `$ $`=`$ $`\frac{1}{2}(\mu _1_{\mu _2}\mu _2_{\mu _1}),`$
$`_{\stackrel{~}{\theta }}`$ $`=`$ $`\frac{1}{2}(\mu _3_{\mu _4}\mu _4_{\mu _3}),`$
$`_\xi `$ $`=`$ $`sc^1(\mu _1_{\mu _1}+\mu _2_{\mu _2})+cs^1(\mu _3_{\mu _3}+\mu _4_{\mu _4}).`$ (22)
Substituting into (2.2), the inverse metric is then expressed in terms of the three unconstrained basis vectors $`(_\xi ,_\theta ,_{\stackrel{~}{\theta }})`$, and hence it can be straightforwardly inverted. Having done so, the downstairs metric can then be re-expressed elegantly in terms of the redundant set of four $`d\mu _i`$ differentials, in the form
$`d\widehat{s}_3^2`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Xi }}}[{\displaystyle \underset{i}{}}Z_id\mu _i^2+\frac{1}{2}b_1^2((\mu _1d\mu _1+\mu _2d\mu _2)^2+(\mu _3d\mu _3+\mu _4d\mu _4)^2)`$ (23)
$`+\frac{1}{2}b_2^2\left((\mu _1d\mu _1+\mu _3d\mu _3)^2+(\mu _2d\mu _2+\mu _4d\mu _4)^2\right)`$
$`+\frac{1}{2}b_3^2((\mu _1d\mu _1+\mu _4d\mu _4)^2+(\mu _2d\mu _2+\mu _3d\mu _3)^2)].`$
Finally, adding this to the $`4\times 4`$ metric $`d\widehat{s}_4^2`$ given in (2.2), we obtain the result for the downstairs 7-metric, $`d\widehat{s}_7^2=d\widehat{s}_4^2+d\widehat{s}_3^2`$:
$`d\widehat{s}_7^2`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Xi }}}[{\displaystyle \underset{i}{}}Z_i(d\mu _i^2+\mu _i^2d\varphi _i^2)+2b_2b_3(\mu _1^2\mu _2^2d\varphi _1d\varphi _2\mu _3^2\mu _4^2d\varphi _3d\varphi _4)`$ (24)
$`+2b_1b_3(\mu _1^2\mu _3^2d\varphi _1d\varphi _3\mu _2^2\mu _4^2d\varphi _2d\varphi _4)+2b_1b_2(\mu _1^2\mu _4^2d\varphi _1d\varphi _4\mu _2^2\mu _3^2d\varphi _2d\varphi _3)`$
$`+\frac{1}{2}b_1^2\left((\mu _1d\mu _1+\mu _2d\mu _2)^2+(\mu _3d\mu _3+\mu _4d\mu _4)^2\right)`$
$`+\frac{1}{2}b_2^2\left((\mu _1d\mu _1+\mu _3d\mu _3)^2+(\mu _2d\mu _2+\mu _4d\mu _4)^2\right)`$
$`+\frac{1}{2}b_3^2((\mu _1d\mu _1+\mu _4d\mu _4)^2+(\mu _2d\mu _2+\mu _3d\mu _3)^2)].`$
We can now work out the eleven-dimensional metric Ansatz, given by (4). To do this, we first note that the determinant of (24), where it is understood that the $`\mu _i`$ coordinates are expressed in terms of $`(\xi ,\theta ,\stackrel{~}{\theta })`$ using (21), is
$$det(\widehat{g}_{mn})=\left(\frac{\mu _1^2\mu _2^2\mu _3^2\mu _4^2}{\mathrm{\Xi }^2}\right)\left(\frac{s^2c^2}{16\mathrm{\Xi }}\right)=\frac{\mu _1^2\mu _2^2\mu _3^2\mu _4^2s^2c^2}{16\mathrm{\Xi }^3},$$
(25)
where in the first expression, the first factor is the determinant of $`4\times 4`$ block involving the $`\varphi _i`$ coordinates, and the second factor is from the $`3\times 3`$ block involving the $`(\xi ,\theta ,\stackrel{~}{\theta })`$ coordinates. From (3), it follows that
$$\widehat{\mathrm{\Delta }}=\mathrm{\Xi }^{\frac{1}{3}},$$
(26)
and hence from (4) that the Ansatz for the eleven-dimensional metric takes the following rather explicit form:
$`d\widehat{s}_{11}^2`$ $`=`$ $`\mathrm{\Xi }^{\frac{1}{3}}ds_4^2+\mathrm{\Xi }^{\frac{1}{3}}d\widehat{s}_7^2`$ (27)
$`=`$ $`\mathrm{\Xi }^{\frac{1}{3}}ds_4^2+g^2\mathrm{\Xi }^{\frac{2}{3}}[{\displaystyle \underset{i}{}}Z_i(d\mu _i^2+\mu _i^2d\varphi _i^2)+2b_2b_3(\mu _1^2\mu _2^2d\varphi _1d\varphi _2\mu _3^2\mu _4^2d\varphi _3d\varphi _4)`$
$`+2b_1b_3(\mu _1^2\mu _3^2d\varphi _1d\varphi _3\mu _2^2\mu _4^2d\varphi _2d\varphi _4)+2b_1b_2(\mu _1^2\mu _4^2d\varphi _1d\varphi _4\mu _2^2\mu _3^2d\varphi _2d\varphi _3)`$
$`+\frac{1}{2}b_1^2\left((\mu _1d\mu _1+\mu _2d\mu _2)^2+(\mu _3d\mu _3+\mu _4d\mu _4)^2\right)`$
$`+\frac{1}{2}b_2^2\left((\mu _1d\mu _1+\mu _3d\mu _3)^2+(\mu _2d\mu _2+\mu _4d\mu _4)^2\right)`$
$`+\frac{1}{2}b_3^2((\mu _1d\mu _1+\mu _4d\mu _4)^2+(\mu _2d\mu _2+\mu _3d\mu _3)^2)].`$
Note that we have reinstated the gauge-coupling constant $`g`$ in this expression.
Having obtained the Kaluza-Klein metric Ansatz for the three dilaton/axion pairs, it is a simple matter to incorporate also the associated $`U(1)^4`$ gauge fields that naturally accompany this truncation of the maximal supergravity. Denoting their potentials by $`A_{\left(1\right)}^i`$, for $`i=1,2,3,4`$, we simply replace each occurrence of $`d\varphi _i`$ in (27) by
$$d\varphi _id\varphi _igA_{\left(1\right)}^i.$$
(28)
Finally in this section, we may note that our result (27) is consistent with previously-obtained special cases. In particular, if we set the three axions $`\chi _i`$ to zero, then the function $`\mathrm{\Xi }`$ reduces to
$$\mathrm{\Xi }=\mathrm{\Delta }^2,$$
(29)
where
$$\mathrm{\Delta }=Y_1Y_2Y_3\mu _1^2+\frac{Y_1}{Y_2Y_3}\mu _2^2+\frac{Y_2}{Y_1Y_3}\mu _3^2+\frac{Y_3}{Y_1Y_2}\mu _4^2.$$
(30)
In the absence of axions, it is natural to define
$$X_1=Y_1Y_2Y_3,X_2=\frac{Y_1}{Y_2Y_3},X_3=\frac{Y_2}{Y_1Y_3},X_4=\frac{Y_3}{Y_1Y_2},$$
(31)
implying that we shall have
$$\mathrm{\Delta }=\underset{i}{}X_i\mu _i^2,Z_i=\mathrm{\Delta }X_i^1.$$
(32)
It can be seen that the metric Ansatz (27) therefore indeed reduces to the one given in if the axions are set to zero.
### 2.3 The Ansatz for the 4-form Field Strength
In principle, we should like to obtain also the Ansatz for the 4-form field strength $`\widehat{F}_{\left(4\right)}`$ of eleven-dimensional supergravity. In spherical Kaluza-Klein reductions it is always much more difficult to obtain the Ansatz for antisymmetric tensors than for the metric, and the present case is no exception. Unfortunately, one can only obtain limited guidance from those results that are presented in . In other truncations, simpler than the case in hand, it has been possible to determine the field-strength Ansatz by brute-force methods, and up to a point, this technique is still useful here. (This method was used successfully in , where the complete and explicit Ansätze for the $`S^7`$ reduction to the bosonic sector of $`N=4`$, $`SO(4)`$ gauged supergravity in $`D=4`$ were obtained.) The contributions to the 4-form Ansatz can be organised into different sectors, and in all except one of these we have obtained complete results. Since these are instructive and useful in their own right, it seems to be worthwhile to present those results that we have obtained here.
We begin with a summary of the four-dimensional theory comprising gravity, the three dilaton/axion pairs, and the associated $`U(1)^4`$ gauge fields.
#### 2.3.1 $`D=4`$ Lagrangian
The complete Lagrangian for four-dimensional $`N=8`$ $`SO(8)`$-gauged supergravity was obtained in . In , the truncation to the $`N=2`$ $`U(1)^4`$-gauged subsector was discussed. Adapting these results to the notation of this paper, we find that the four-dimensional bosonic Lagrangian for this $`N=2`$ truncation is given by
$$_4=R\text{1}\mathrm{l}\frac{1}{2}\underset{i=1}{\overset{3}{}}(d\phi _id\phi _i+e^{2\phi _i}d\chi _id\chi _i)V\text{1}\mathrm{l}+_{Kin}+_{CS},$$
(33)
where $`V`$ is the potential for the scalar fields, and $`_{Kin}`$ and $`_{CS}`$ are the kinetic terms and the Chern-Simons terms for the four $`U(1)`$ gauge fields $`F_{\left(2\right)}^i=dA_{\left(1\right)}^i`$. The scalar potential is given by
$$V=4g^2\underset{i=1}{\overset{3}{}}(Y_i^2+\stackrel{~}{Y}_i^2).$$
(34)
The kinetic terms for the gauge fields are
$`_{Kin}`$ $`=`$ $`\frac{1}{2}|W|^2[P_0(\stackrel{~}{Y}_1^2\stackrel{~}{Y}_2^2\stackrel{~}{Y}_3^2F_{\left(2\right)}^1F_{\left(2\right)}^1+\stackrel{~}{Y}_1^2Y_2^2Y_3^2F_{\left(2\right)}^2F_{\left(2\right)}^2`$ (35)
$`+Y_1^2\stackrel{~}{Y}_2^2Y_3^2F_{\left(2\right)}^3F_{\left(2\right)}^3+Y_1^2Y_2^2\stackrel{~}{Y}_3^2F_{\left(2\right)}^4F_{\left(2\right)}^4)`$
$`+2P_1b_2b_3(\stackrel{~}{Y}_1^2F_{\left(2\right)}^1F_{\left(2\right)}^2Y_1^2F_{\left(2\right)}^3F_{\left(2\right)}^4)`$
$`+2P_2b_1b_3(\stackrel{~}{Y}_2^2F_{\left(2\right)}^1F_{\left(2\right)}^3Y_2^2F_{\left(2\right)}^2F_{\left(2\right)}^4)`$
$`+2P_3b_1b_2(\stackrel{~}{Y}_3^2F_{\left(2\right)}^1F_{\left(2\right)}^4Y_3^2F_{\left(2\right)}^2F_{\left(2\right)}^3)],`$
where
$`P_01+b_1^2+b_2^2+b_3^2,WP_02\mathrm{i}b_1b_2b_3,`$
$`P_11b_1^2+b_2^2+b_3^2,P_21+b_1^2b_2^2+b_3^2,P_31+b_1^2+b_2^2b_3^2.`$ (36)
Finally, the Chern-Simons terms for the gauge fields are
$`_{CS}`$ $`=`$ $`|W|^2[b_1b_2b_3(\stackrel{~}{Y}_1^2\stackrel{~}{Y}_2^2\stackrel{~}{Y}_3^2F_{\left(2\right)}^1F_{\left(2\right)}^1+\stackrel{~}{Y}_1^2Y_2^2Y_3^2F_{\left(2\right)}^2F_{\left(2\right)}^2`$ (37)
$`+Y_1^2\stackrel{~}{Y}_2^2Y_3^2F_{\left(2\right)}^3F_{\left(2\right)}^3+Y_1^2Y_2^2\stackrel{~}{Y}_3^2F_{\left(2\right)}^4F_{\left(2\right)}^4)`$
$`+b_1(P_0+2b_2^2b_3^2)(\stackrel{~}{Y}_1^2F_{\left(2\right)}^1F_{\left(2\right)}^2Y_1^2F_{\left(2\right)}^3F_{\left(2\right)}^4)`$
$`+b_2(P_0+2b_1^2b_3^2)(\stackrel{~}{Y}_2^2F_{\left(2\right)}^1F_{\left(2\right)}^3Y_2^2F_{\left(2\right)}^2F_{\left(2\right)}^4)`$
$`+b_3(P_0+2b_1^2b_2^2)(\stackrel{~}{Y}_3^2F_{\left(2\right)}^1F_{\left(2\right)}^4Y_3^2F_{\left(2\right)}^2F_{\left(2\right)}^3)].`$
From (33), we find that the equations of motion for the gauge fields are
$$d(|W|^2R_i)=0,$$
(38)
for $`i=1,2,3,4`$, where
$`R_1`$ $`=`$ $`\stackrel{~}{Y}_1^2\stackrel{~}{Y}_2^2\stackrel{~}{Y}_3^2[P_0F_{\left(2\right)}^1+2b_1b_2b_3F_{\left(2\right)}^1]+\stackrel{~}{Y}_1^2[P_1b_2b_3F_{\left(2\right)}^2+b_1(P_0+2b_2^2b_3^2)F_{\left(2\right)}^2]`$
$`+\stackrel{~}{Y}_2^2[P_2b_1b_3F_{\left(2\right)}^3+b_2(P_0+2b_1^2b_3^2)F_{\left(2\right)}^3]+\stackrel{~}{Y}_3^2[P_3b_1b_2F_{\left(2\right)}^4+b_3(P_0+2b_1^2b_2^2)F_{\left(2\right)}^4],`$
$`R_2`$ $`=`$ $`\stackrel{~}{Y}_1^2Y_2^2Y_3^2[P_0F_{\left(2\right)}^2+2b_1b_2b_3F_{\left(2\right)}^2]+\stackrel{~}{Y}_1^2[P_1b_2b_3F_{\left(2\right)}^1+b_1(P_0+2b_2^2b_3^2)F_{\left(2\right)}^1]`$
$`Y_2^2[P_2b_1b_3F_{\left(2\right)}^4+b_2(P_0+2b_1^2b_3^2)F_{\left(2\right)}^4]Y_3^2[P_3b_1b_2F_{\left(2\right)}^3+b_3(P_0+2b_1^2b_2^2)F_{\left(2\right)}^3],`$
$`R_3`$ $`=`$ $`Y_1^2\stackrel{~}{Y}_2^2Y_3^2[P_0F_{\left(2\right)}^3+2b_1b_2b_3F_{\left(2\right)}^3]Y_1^2[P_1b_2b_3F_{\left(2\right)}^4+b_1(P_0+2b_2^2b_3^2)F_{\left(2\right)}^4]`$
$`+\stackrel{~}{Y}_2^2[P_2b_1b_3F_{\left(2\right)}^1+b_2(P_0+2b_1^2b_3^2)F_{\left(2\right)}^1]Y_3^2[P_3b_1b_2F_{\left(2\right)}^2+b_3(P_0+2b_1^2b_2^2)F_{\left(2\right)}^2],`$
$`R_4`$ $`=`$ $`Y_1^2Y_2^2\stackrel{~}{Y}_3^2[P_0F_{\left(2\right)}^4+2b_1b_2b_3F_{\left(2\right)}^4]Y_1^2[P_1b_2b_3F_{\left(2\right)}^3+b_1(P_0+2b_2^2b_3^2)F_{\left(2\right)}^3]`$
$`Y_2^2[P_2b_1b_3F_{\left(2\right)}^2+b_2(P_0+2b_1^2b_3^2)F_{\left(2\right)}^2]+\stackrel{~}{Y}_3^2[P_3b_1b_2F_{\left(2\right)}^1+b_3(P_0+2b_1^2b_2^2)F_{\left(2\right)}^1].`$
#### 2.3.2 The Ansatz for $`\widehat{F}_{\left(4\right)}`$
In previous papers the Ansatz for the 4-form field strength $`\widehat{F}_{\left(4\right)}`$ was obtained for the $`U(1)^4`$ truncation in absence of the three axions , and for the $`N=4`$ gauged $`SO(4)`$ truncation, in which there is one scalar and one axion . Based on those results, it can be seen to be natural to write the Ansatz for $`\widehat{F}_{\left(4\right)}`$ as the sum of three terms, each with its own characteristic contribution to the whole.
Thus we are led to the following construction for the 4-form field strength:
$`\widehat{F}_{\left(4\right)}`$ $`=`$ $`2gUϵ_{\left(4\right)}+\widehat{F}_{\left(4\right)}^{}+\widehat{F}_{\left(4\right)}^{\prime \prime }`$ (40)
$`+\frac{1}{2g}(2Y_1^1dY_1\chi _1Y_1^4d\chi _1)d(\mu _1^2+\mu _2^2)`$
$`+\frac{1}{2g}(2Y_2^1dY_2\chi _2Y_2^4d\chi _2)d(\mu _1^2+\mu _3^2)`$
$`+\frac{1}{2g}(2Y_3^1dY_3\chi _3Y_3^4d\chi _3)d(\mu _1^2+\mu _4^2),`$
where
$$U=Y_1^2(\mu _1^2+\mu _2^2)+\stackrel{~}{Y}_1^2(\mu _3^2+\mu _4^2)+Y_2^2(\mu _1^2+\mu _3^2)+\stackrel{~}{Y}_2^2(\mu _2^2+\mu _4^2)+Y_3^2(\mu _1^2+\mu _4^2)+\stackrel{~}{Y}_3^2(\mu _2^2+\mu _3^2),$$
(41)
and $`ϵ_{\left(4\right)}`$ denotes the volume form on the four-dimensional spacetime. The term $`\widehat{F}_{\left(4\right)}^{\prime \prime }`$ will be given by
$$\widehat{F}_{\left(4\right)}^{\prime \prime }=\frac{1}{2g^2}|W|^2\underset{i}{}d\mu _i^2(d\varphi _igA_{\left(1\right)}^i)R_i.$$
(42)
(We shall justify these expressions below.) The remaining term is $`\widehat{F}_{\left(4\right)}^{}`$. This will be written in terms of a potential $`\widehat{A}_{\left(3\right)}^{}`$, as $`\widehat{F}_{\left(4\right)}^{}=d\widehat{A}_{\left(3\right)}^{}`$. It will be the determination of $`\widehat{A}_{\left(3\right)}^{}`$ that presents the greatest difficulty.
It will be noted that $`\widehat{F}_{\left(4\right)}`$ does not identically satisfy $`d\widehat{F}_{\left(4\right)}=0`$. This feature was already seen in the truncations in and . It is not possible, at least within the usual second-order formulation of eleven-dimensional supergravity, to write an Ansatz for $`\widehat{F}_{\left(4\right)}`$ in the $`S^7`$ reduction that identically satisfies $`d\widehat{F}_{\left(4\right)}=0`$. An implication from this is that one cannot write the Ansatz directly on the potential $`\widehat{A}_{\left(3\right)}`$, which in turn means that one cannot write an Ansatz that can be substituted directly into the eleven-dimensional action. One must work at the level of the equations of motion.
In fact the requirement that $`\widehat{F}_{\left(4\right)}`$ must satisfy the Bianchi identity $`d\widehat{F}_{\left(4\right)}=0`$ provides us with very important clues as to the correct form of the reduction Ansatz, and we used this in writing down our results in (40) and (42). The point is that the Bianchi identity will be satisfied by virtue of the $`D=4`$ equations of motion for the scalar fields and the $`U(1)`$ gauge fields being satisfied. (To be precise, the scalar equations of motion in question here are those of the three dilatons $`\phi _i`$, in combination with certain non-linear admixtures of the three axion equations of motion.) Of course the contribution to $`\widehat{F}_{\left(4\right)}`$ from $`\widehat{A}_{\left(3\right)}^{}`$, whose precise form we have not been able to determine, does not enter into the discussion of the Bianchi identity, since it gives a contribution $`\widehat{F}_{\left(4\right)}^{}`$ that identically satisfies $`d\widehat{F}_{\left(4\right)}^{}=0`$.
To see how the Bianchi identity $`d\widehat{F}_{\left(4\right)}=0`$ implies the four-dimensional equations of motion for the scalars and the gauge fields, we note from the structure of (40) and (42) that after acting with $`d`$ we shall have two distinct classes of terms. First, there will be terms of the form $`d\mu _i^2\omega _{\left(4\right)}`$, where $`\omega _{\left(4\right)}`$ is a 4-form living entirely in the four-dimensional spacetime. ($`\omega _{\left(4\right)}`$ will comprise terms of the form $`ϵ_{\left(4\right)}`$, and of the form $`ddY_i`$, etc. Of course they are all proportional to $`ϵ_{\left(4\right)}`$.) The requirement of the vanishing of these terms will imply the scalar equations of motion. Secondly, there will be terms of the form $`d\mu _i^2(d\varphi _i\frac{1}{2}gA_{\left(1\right)}^i)\omega _{\left(3\right)}`$ coming from the action of $`d`$ on $`\widehat{F}_{\left(4\right)}^{\prime \prime }`$, where $`\omega _{\left(3\right)}`$ is a 3-form living in the four-dimensional spacetime. The vanishing of these terms will imply the four-dimensional equations of motion for the gauge fields.
Let us consider the second type of contribution first, since it is the simpler one. The terms of this type come only from $`d\widehat{F}_{\left(4\right)}^{\prime \prime }`$, and give
$$\underset{i}{}d\mu _i^2(d\varphi _igA_{\left(1\right)}^i)d(|W|^2R_i)=0.$$
(43)
This can immediately be seen to imply precisely the equations of motion for the four $`U(1)`$ gauge fields, given in (38).
It remains to check that the terms of the form $`d\mu _i^2\omega _{\left(4\right)}`$ coming from the Bianchi identity vanish by virtue of the four-dimensional scalar equations of motion. The kinetic terms of these scalar equations come from the action of $`d`$ on the final three lines in (40). Clearly, we get the combinations of the form
$$d(2Y_1^1dY_1\chi _1Y_1^4d\chi _1),$$
(44)
arising (with similar independent expressions involving the $`(Y_2,\chi _2)`$ and $`(Y_3,\chi _3)`$ pairs). This is a combination of the $`\phi _1`$ and the $`\chi _1`$ equations of motion. In fact it is
$$[dd\phi _1+e^{2\phi _1}d\chi _1d\chi _1]\chi _1[d(e^{2\phi _1}d\chi _1)],$$
(45)
where the first quantity in square brackets is the dilaton equation of motion, and the second quantity in square brackets is the axion equation of motion.
This particular combination, of the dilaton equation plus an admixture of the axion equation, is an especially simple one to compare with the scalar equations of motion coming from the four-dimensional Lagrangian (33). It means that we are looking at the combination that comes from the following variation of the $`D=4`$ Lagrangian:
$$\widehat{\delta }_4\frac{\delta _4}{\delta \phi _1}\chi _1\frac{_4}{\delta \chi _1}.$$
(46)
If we define a symbol $`\widehat{\delta }`$ to denote this specific combination of field variations, i.e.
$$\widehat{\delta }\frac{\delta }{\delta \phi _1}\chi _1\frac{\delta }{\delta \chi _1},$$
(47)
then we find the great simplification that
$$\widehat{\delta }Y_1^2=Y_1^2,\widehat{\delta }\stackrel{~}{Y}_1^2=\stackrel{~}{Y}_1^2,\widehat{\delta }b_1=0.$$
(48)
(Of course since we are focusing on the scalars with the index $`i=1`$ at the moment, all of the scalar quantities with $`i=2`$ or $`i=3`$ labels are invariant under this transformation.) The last equation in (48), $`\widehat{\delta }b_1=0`$, leads to an enormous simplification when we vary $`_{Kin}`$ and $`_{CS}`$ given by (35) and (37). It means that $`|W|`$, the $`P_a`$, and all the $`b_i`$ are invariant. We need only consider $`Y_1`$ and $`\stackrel{~}{Y}_1`$, and these just vary by the very simple rules given in (48).
With these observations, it becomes a relatively straightforward matter to verify that the terms of the form $`d\mu _i^2\omega _{\left(4\right)}`$ that arise in the Bianchi identity for $`\widehat{F}_{\left(4\right)}`$ vanish precisely as a consequence of the scalar equations of motion following from (33), to all orders in scalar fields and gauge field strengths. Note that the contributions to the scalar equations of motion from the potential $`V`$ given in (34) arise from the action of the exterior derivative on the term $`2gUϵ_{\left(4\right)}`$ in (40). This part of the calculation can be seen quite easily, and can be examined in isolation from the more complicated contributions from the four-dimensional gauge fields.
The contribution $`\widehat{F}_{\left(4\right)}^{}=\widehat{A}_{\left(3\right)}^{}`$ in (40) remains undetermined. We know some aspects of it structure, for example that it is of the general from
$$\widehat{A}_{\left(3\right)}^{}=\underset{ij}{}h_{ij}(\mu _i^2d\mu _j^2\mu _j^2d\mu _i^2)(d\varphi _igA_{\left(1\right)}^i)(d\varphi _jgA_{\left(1\right)}^j),$$
(49)
where the functions $`h_{ij}`$ depend on the scalars $`\phi _i`$ and $`\chi _i`$, and the direction cosines $`\mu _i`$. At leading order, these terms will give rise to the linearised Ansatz for the axions $`\chi _i`$. If explicit expressions for the complete Ansatz for the $`N=8`$ $`SO(8)`$ gauged supergravity embedding were available, $`A_{\left(3\right)}^{}`$ could in principle be determined by substituting the expressions for $`u_{ij}^{KL}`$ and $`v_{ijKL}`$ appearing in (9) and (10) into them. To the extent that such expressions are implicit in the work of , a procedure in principle exists for reading off $`A_{\left(3\right)}^{}`$. It is not clear that attempting such a substitution would be simpler than a brute-force direct attack on the problem, of the type that has proved successful in previous (simpler) cases .
### 2.4 Domain wall solutions and their oxidation
The four-dimensional $`U(1)^4`$ Lagrangian (33) supports a four-charge AdS black hole solution . In the extremal limit, the four $`U(1)`$ gauge fields decouple and the solution becomes AdS domain wall, supported by the scalar fields only. It is given by
$`ds_4`$ $`=`$ $`(gr)^4(H_1H_2H_3H_4)^{1/2}dx^\mu dx_\mu +(H_1H_2H_3H_4)^{1/2}{\displaystyle \frac{dr^2}{g^2r^2}},`$
$`e^{\phi _i}`$ $`=`$ $`Y_1^2=f_i,\chi _i=0,`$ (50)
where
$`f_1`$ $`=`$ $`{\displaystyle \frac{(H_3H_4)^{1/2}}{(H_1H_2)^{1/2}}},f_2={\displaystyle \frac{(H_2H_4)^{1/2}}{(H_1H_3)^{1/2}}},`$
$`f_3`$ $`=`$ $`{\displaystyle \frac{(H_2H_3)^{1/2}}{(H_1H_4)^{1/2}}},H_i=1+{\displaystyle \frac{\mathrm{}_i^2}{r^2}}.`$ (51)
This solution can be oxidised back to $`D=11`$ , where it acquires the interpretation of being a continuous ellipsoidal distribution of M2-branes.
The scalar kinetic terms in the Lagrangian (33) are invariant under global $`SL(2,R)^3`$ transformations, corresponding to the usual fractional-linear group action on each of the axion/dilaton pairs. The scalar potential in (33), on the other hand, is invariant only under the $`SO(2)^3`$ subgroup transformations
$$\tau _i\tau _i^{}=\frac{\mathrm{cos}\lambda _i\tau +sin\lambda _i}{\mathrm{sin}\lambda _i\tau +\mathrm{cos}\lambda _i}.$$
(52)
where $`\tau _i=\chi _i+\mathrm{i}e^{\phi _i}`$. Applying these global transformations to the original domain walls we obtain new solutions, with
$$Y_i^2=e^{\phi _i}=\frac{1}{f_i}(f_i^2\mathrm{cos}^2\lambda _i+\mathrm{sin}^2\lambda _i),\chi _i=\frac{\frac{1}{2}(f_i1)}{f_i^2\mathrm{cos}^2\lambda _i+\mathrm{sin}^2\lambda _i}.$$
(53)
The $`\stackrel{~}{Y}_i`$ are hence given by
$$\stackrel{~}{Y}_i^2=\frac{f^2+\frac{1}{4}(f1)^2\mathrm{sin}^2(2\lambda _i)}{f_i(f_i^2\mathrm{cos}^2\lambda _i+\mathrm{sin}^2\lambda _i)}.$$
(54)
Having obtained the $`SO(2)^3`$ rotated domain-wall solutions, they can be oxidised back to $`D=11`$. The eleven-dimensional metric is given by substituting the solution into (27). These solutions with non-vanishing $`\chi _i`$ no longer simply describe distributed M2-branes. To see this we note from (49) that with non-vanishing axions the field strength $`F_{\left(4\right)}`$ will involve components lying purely in the internal $`S^7`$. By contrast, in a distributed M2-brane solution one has $`F_{\left(4\right)}=d^3xdH^1`$, where $`H`$ is the harmonic function in the transverse space. Thus for a distributed M2-brane the field strength $`F_{\left(4\right)}`$ always carries three world-volume indices.
## 3 The 2-scalar $`D=5`$ embedding in type IIB
In this section, we consider the embedding of the 2-scalar truncation of $`D=5`$ gauged supergravity discussed in the introduction, and its embedding in the type IIB theory via an $`S^5`$ reduction. In the early stages of the derivation, we retain all four of the scalar fields of the truncation discussed in .
### 3.1 The metric reduction Ansatz
The set of 42 spin-0 fields in the complete $`SO(6)`$ gauged $`N=8`$ supergravity in $`D=5`$ are described by a 27-bein $`𝒱`$, which transforms under local $`USp(8)`$ and global $`E_6`$. The truncation to four spin-0 fields is described in , in terms of an $`SL(6,R)\times SL(2,R)`$ basis, for which the components of the vielbein are decomposed as $`(𝒱^{IJab},𝒱_{I\alpha }{}_{}{}^{ab})`$. In terms of this decomposition, the following conjecture for the inverse $`S^5`$ metric has been proposed :
$$\widehat{g}^{mn}(x,y)\widehat{\mathrm{\Delta }}^{\frac{2}{3}}g^{mn}(x,y)=2K_{IJ}^mK_{KL}^n\stackrel{~}{𝒱}_{IJab}\stackrel{~}{𝒱}_{KLcd}\mathrm{\Omega }^{ac}\mathrm{\Omega }^{bd},$$
(55)
where $`\stackrel{~}{𝒱}`$ is the inverse of the vielbein $`𝒱`$, $`\widehat{\mathrm{\Delta }}^2=det(g_{mn}(x,y))/det(g_{mn}(y))`$, and $`g_{mn}(y)`$ is the undeformed round $`S^5`$ metric where the scalar fields are set to zero. The ten-dimensional metric Ansatz will then be
$$d\widehat{s}_{10}^2=\widehat{\mathrm{\Delta }}^{\frac{2}{3}}ds_5^2+g_{mn}(x,y)dy^mdy^n=\widehat{\mathrm{\Delta }}^{\frac{2}{3}}(ds_5^2+\widehat{g}_{mn}(x,y)dy^mdy^n).$$
(56)
The process of making the 4-scalar truncation in the vielbein $`𝒱`$ has been described in detail in . Substituting this into the metric Ansatz (55) is a mechanical exercise that is most conveniently implemented by computer. Since the final result is considerably simpler than the intermediate stages we shall, without further ado, present the final answer. We find that the inverse 5-sphere metric $`\widehat{}_{s_5}^2\widehat{g}^{mn}_m_n`$ is given by
$`\widehat{}_{s_5}^2`$ $`=`$ $`X^1(\mathrm{cosh}2y_2(\mathrm{cosh}2r\mathrm{sin}\theta \mathrm{sinh}2r)(K_{15}^2+K_{25}^2+K_{35}^2+K_{45}^2)`$ (57)
$`+\mathrm{cosh}2y_2(\mathrm{cosh}2r+\mathrm{sin}\theta \mathrm{sinh}2r)(K_{16}^2+K_{26}^2+K_{36}^2+K_{46}^2)`$
$`+2\mathrm{cos}\theta \mathrm{sinh}2r\mathrm{sinh}2y_2(K_{26}K_{35}K_{25}K_{36}+K_{16}K_{45}K_{15}K_{46}))`$
$`+X^2(\frac{1}{4}(3\mathrm{cos}\theta +2\mathrm{cos}^2\theta \mathrm{cosh}4r)(K_{12}^2+K_{34}^2)+(K_{14}^2+K_{23}^2)`$
$`+\mathrm{cosh}^22y_2(K_{13}^2+K_{24}^2)+2\mathrm{cos}^2\theta \mathrm{sinh}^22rK_{12}K_{34}2\mathrm{sinh}^22y_2K_{13}K_{24})`$
$`+X^4K_{56}^2.`$
The scalars $`(X,r,y_2,\theta )`$ are related to the quantities $`(\rho ,\phi _1,\phi _2,\varphi )`$ appearing in by
$$\rho =X^{\frac{1}{2}},r=\frac{1}{2}(\phi _2\phi _1),y_2=\frac{1}{2}(\phi _1+\phi _2),\theta =2\varphi .$$
(58)
Note that the $`D=5`$ scalar Lagrangian for this truncation is
$$=2\underset{i=1}{\overset{3}{}}(\phi _i)^2\mathrm{sinh}^2(\phi _1\phi _2)(\theta )^2V,$$
(59)
where $`X=e^{\sqrt{6}\phi _3/2}`$, and the scalar potential $`V`$ takes the form
$`V`$ $`=`$ $`g^2(X^2[1\mathrm{cos}^2\theta (\mathrm{sinh}^2\phi _1\mathrm{sinh}^2\phi _2)]+X^1[\mathrm{cosh}2\phi _1+\mathrm{cosh}2\phi _2]`$
$`+\frac{1}{16}X^4[2+2\mathrm{sin}^2\theta 2\mathrm{sin}^2\theta \mathrm{cosh}(2(\phi _1\phi _2))\mathrm{cosh}4\phi _1\mathrm{cosh}4\phi _2]).`$
At this stage, we impose the further truncation to the 2-scalar subsector that we really want to consider. This corresponds to setting $`\theta =0`$ and $`\phi _2=0`$ . It is easily verified from (59) and (3.1) that this is a consistent truncation. Thus we shall have $`r=\frac{1}{2}\phi `$, and $`y_2=\frac{1}{2}\phi `$, where we now drop the “1” subscript on $`\phi _1`$. The potential (3.1) reduces to
$$V=\mathrm{cosh}^2\phi \left[X^2(2\mathrm{cosh}^2\phi )+2X^1\frac{1}{2}X^4\mathrm{sinh}^2\phi \right].$$
(61)
It is convenient also at this stage to perform a labelling of indices on the Killing vectors $`K_{ij}`$ in (57), under which the index values $`(2,3,4)`$ are cycled: $`23`$, $`34`$ and $`42`$.
We now adopt a description of the round 5-sphere that is precisely analogous to the one that we introduced in Appendix B for $`S^7`$. This time, we shall end up with three “direction cosines” $`\mu _i`$, subject to the condition $`\mu _i\mu _i=1`$, and three azimuthal angles $`\varphi _i`$. After manipulations similar to those in section 2, we arrive at the following expression for the inverse 5-sphere metric $`\widehat{}_5^2`$:
$$\widehat{}_{s_5}^2=\widehat{}_{s_2}^2+\widehat{}_{s_3}^2,$$
(62)
where the $`2\times 2`$ and $`3\times 3`$ blocks are given by
$`\widehat{}_{s_2}^2`$ $`=`$ $`\mathrm{cosh}^2\phi \left(X^1[(\mu _1_{\mu _3}\mu _3_{\mu _1})^2+(\mu _2_{\mu _3}\mu _3_{\mu _2})^2]+X^2(\mu _1_{\mu _2}\mu _2_{\mu _1})^2\right),`$
$`\widehat{}_{s_3}^2`$ $`=`$ $`\mathrm{\Delta }\mathrm{cosh}^2\phi \left[X(\mu _1^2_{\varphi _1}^2+\mu _2^2_{\varphi _2}^2)+X^2\mu _3^2_{\varphi _3}^2\right]`$ (63)
$`\mathrm{sinh}^2\phi \left(X(_{\varphi _1}+_{\varphi _2})X^2_{\varphi _3}\right)^2,`$
and
$$\mathrm{\Delta }(\mu _1^2+\mu _2^2)X+\mu _3^2X^2.$$
(64)
Note that the $`2\times 2`$ inverse metric $`\widehat{}_{s_2}^2`$ is just equal to the metric for the single-scalar truncation when $`\phi =0`$, multiplied by a factor of $`\mathrm{cosh}^2\phi `$. The $`3\times 3`$ inverse metric is equal to $`\mathrm{cosh}^2\phi `$ times the $`\phi =0`$ metric, with the correction term appearing in its second line.
The inverse of the $`3\times 3`$ block $`\widehat{}_{s_3}^2`$ is straightforward to calculate, and we find
$$d\widehat{s}_3^2=\frac{\mathrm{sech}^2\phi }{\mathrm{\Delta }}\left(X^1(\mu _1^2d\varphi _1^2+\mu _2^2d\varphi _2^2)+X^2\mu _3^2d\varphi _3^2\right)+\frac{\mathrm{tanh}^2\phi }{\mathrm{\Delta }^2}(\mu _1^2d\varphi _1+\mu _2^2d\varphi _2\mu _3^2d\varphi _3)^2.$$
(65)
Note that the determinant of $`d\widehat{s}_3^2`$ is given by $`(\mu _1\mu _2\mu _3)^2/(\mathrm{\Delta }^3\mathrm{cosh}^4\phi )`$.
For the $`2\times 2`$ block, the inversion gives the metric
$$d\widehat{s}_2^2=\frac{\mathrm{sech}^2\phi }{\mathrm{\Delta }}\left(X^1(d\mu _1^2+d\mu _2^2)+X^2d\mu _3^2\right).$$
(66)
It is helpfull at this stage to reparameterise the direction cosines $`\mu _i`$, and make redefinitions of the azimuthal angles $`(\varphi _1,\varphi _2)`$ as follows:
$`\mu _1=\mathrm{cos}\xi \mathrm{cos}\frac{1}{2}\vartheta ,\mu _2=\mathrm{cos}\xi \mathrm{sin}\frac{1}{2}\vartheta ,\mu _3=\mathrm{sin}\xi ,`$
$`\varphi _1=\frac{1}{2}(\psi +\varphi ),\varphi _2=\frac{1}{2}(\psi \varphi ).`$ (67)
In fact $`(\vartheta ,\varphi ,\psi )`$ are just the Euler angles on $`S^3`$. One can define left-invariant 1-forms $`\sigma _i`$, as
$$\sigma _1+\mathrm{i}\sigma _2=e^{\mathrm{i}\psi }(d\vartheta +\mathrm{i}\mathrm{sin}\vartheta d\varphi ),\sigma _3=d\psi +\mathrm{cos}\vartheta d\varphi .$$
(68)
These satisfy $`d\sigma _1=\sigma _2\sigma _3`$, and cyclically. Defining also
$$c\mathrm{cos}\xi ,s\mathrm{sin}\xi ,$$
(69)
we find that the 5-dimensional internal metric $`d\widehat{s}_5^2\widehat{g}_{mn}(x,y)dy^mdy^n=d\widehat{s}_2^2+d\widehat{s}_3^2`$ becomes
$`d\widehat{s}_5^2`$ $`=`$ $`{\displaystyle \frac{\mathrm{sech}^2\phi }{\mathrm{\Delta }}}\left[X\mathrm{\Delta }d\xi ^2+\frac{1}{4}X^1c^2(\sigma _1^2+\sigma _2^2+\sigma _3^2)+X^2s^2d\varphi _3^2\right]`$ (70)
$`+{\displaystyle \frac{\mathrm{tanh}^2\phi }{4\mathrm{\Delta }^2}}(c^2\sigma _32s^2d\varphi _3)^2,`$
where
$$\mathrm{\Delta }=Xc^2+X^2s^2.$$
(71)
In the absence of the pseudoscalar field $`\phi `$, this reduces to the metric Ansatz encountered in the $`N=4`$ gauged $`SU(2)\times U(1)`$ supergravity embedding obtained in . In that case, the scalar field $`X`$ parameterises inhomogeneous deformations of $`S^5`$ viewed as a foliation of $`S^3\times S^1`$ surfaces.
With the pseudoscalar $`\phi `$ non-vanishing, it is advantageous to rewrite the metric (70) as the sum of squares of just five quantities, by completing the square. After doing this, we obtain the result
$$d\widehat{s}_5^2=\frac{X}{\mathrm{cosh}^2\phi }d\xi ^2+\frac{c^2X^1}{4\mathrm{\Delta }\mathrm{cosh}^2\phi }(\sigma _1^2+\sigma _2^2)+\frac{c^2X}{4\mathrm{\Omega }}\sigma _3^2+\frac{s^2\mathrm{\Omega }}{\mathrm{\Delta }^2\mathrm{cosh}^2\phi }\left(d\varphi _3\frac{c^2\mathrm{sinh}^2\phi }{2\mathrm{\Omega }}\sigma _3\right)^2,$$
(72)
where
$$\mathrm{\Omega }X^3c^2+s^2\mathrm{cosh}^2\phi .$$
(73)
This expression reduces to the one found in if $`\phi =0`$. In that case, the scalar $`X`$ parameterises deformations of $`S^5`$ corresponding to inhomogeneities of codimension 1 of the foliation by $`S^3\times S^1`$. When the pseudoscalar $`\phi `$ is included too, the inhomogeneities remain of codimension 1, but with a slightly more complicated structure. In addition, there is a sort of “twist” in the $`S^3\times S^1`$ product structure of the homogeneous foliating surfaces, as indicated by the cross-term between the interval $`d\varphi _3`$ on $`S^1`$, and the 1-form $`\sigma _3`$ on $`S^3`$.
Finally, substituting our result for the internal hatted metric $`d\widehat{s}_5^2`$ into (56), we arrive at the conjectured ten-dimensional metric Ansatz for this two-scalar truncation:
$`d\widehat{s}_{10}^2`$ $`=`$ $`\mathrm{\Delta }^{\frac{1}{2}}\mathrm{cosh}\phi ds_5^2+{\displaystyle \frac{X\mathrm{\Delta }^{\frac{1}{2}}}{\mathrm{cosh}\phi }}d\xi ^2+{\displaystyle \frac{c^2X^1}{4\mathrm{\Delta }^{\frac{1}{2}}\mathrm{cosh}\phi }}(\sigma _1^2+\sigma _2^2)`$ (74)
$`+{\displaystyle \frac{c^2X\mathrm{\Delta }^{\frac{1}{2}}\mathrm{cosh}\phi }{4\mathrm{\Omega }}}\sigma _3^2+{\displaystyle \frac{s^2\mathrm{\Omega }}{\mathrm{\Delta }^{\frac{3}{2}}\mathrm{cosh}\phi }}\left(d\varphi _3{\displaystyle \frac{c^2\mathrm{sinh}^2\phi }{2\mathrm{\Omega }}}\sigma _3\right)^2.`$
### 3.2 The field-strength Ansätze
There does not seem to be any straightforward way to determine the Ansatz for the Kaluza-Klein reduction other fields of the ten-dimensional type IIB theory, in this two-scalar reduction. We know that when $`\phi `$ is taken to be zero, the Ansatz must reduce to one encompassed by the results in . In particular, the remaining scalar field $`X`$ enters in the Ansatz for the self-dual 5-form, whilst the dilaton, axion and 3-form field strengths of the type IIB theory vanish when $`\phi =0`$. Since it is a pseudoscalar, the field $`\phi `$ enters at the linearised level in the Ansatz for the NS-NS and R-R 2-form potentials $`\widehat{A}_{\left(2\right)}\widehat{A}_{\left(2\right)}^{\mathrm{NS}}`$ and $`\widehat{A}_{\left(2\right)}^{\mathrm{RR}}`$ .
The relevant bosonic equations of motion of the type IIB theory are
$`\widehat{R}_{MN}`$ $`=`$ $`\frac{1}{96}\widehat{H}_{MN}^2+\frac{1}{4}\left((\widehat{F}_{\left(3\right)}^1)_{MN}^2\frac{1}{12}(\widehat{F}_{\left(3\right)}^1)^2\widehat{g}_{MN}\right)+\frac{1}{4}\left((\widehat{F}_{\left(3\right)}^2)_{MN}^2\frac{1}{12}(\widehat{F}_{\left(3\right)}^2)^2\widehat{g}_{MN}\right),`$
$`d\widehat{}\widehat{F}_{\left(3\right)}`$ $`=`$ $`\mathrm{i}\widehat{H}_{\left(5\right)}\widehat{F}_{\left(3\right)},`$ (75)
$`d\widehat{H}_{\left(5\right)}`$ $`=`$ $`\frac{\mathrm{i}}{2}\widehat{F}_{\left(3\right)}\widehat{\overline{F}_{\left(3\right)}},\widehat{H}_{\left(5\right)}=\widehat{}\widehat{H}_{\left(5\right)},`$
where we have introduced the notation that
$$\widehat{A}_{\left(2\right)}\widehat{A}_{\left(2\right)}^{\mathrm{NS}}+\mathrm{i}\widehat{A}_{\left(2\right)}^{\mathrm{RR}}.$$
(76)
We are assuming here that the dilaton and axion of the type IIB theory vanish in the reduction. For this to be consistent with the type IIB equations of motion, it is necessary that
$$\widehat{}\widehat{F}_{\left(3\right)}\widehat{\overline{F}_{\left(3\right)}}=0,\widehat{}\widehat{F}_{\left(3\right)}\widehat{F}_{\left(3\right)}=\widehat{}\widehat{\overline{F}_{\left(3\right)}}\widehat{\overline{F}_{\left(3\right)}}.$$
(77)
We shall restrict our discussion from now on to the linearised level.
In the notation that we are using here, the linearised Ansatz for pseudoscalars $`\phi `$ will be of the form
$$\widehat{A}_{\left(2\right)}=\phi Y_{\left(2\right)},$$
(78)
where $`Y_{\left(2\right)}`$ is a complex 2-form spherical harmonic satisfying
$$dY_{\left(2\right)}=\mathrm{i}\lambda Y_{\left(2\right)}$$
(79)
on the unit round 5-sphere. The Ansatz for the self-dual 5-form $`\widehat{H}_{\left(5\right)}\widehat{G}_{\left(5\right)}+\widehat{}\widehat{G}_{\left(5\right)}`$ includes a Freunnd-Rubin term $`\widehat{G}_{\left(5\right)}=4ϵ_{\left(5\right)}`$ (we have set the gauge coupling $`g=1`$ here). Substituting into the type IIB equations of motion, one finds that the pseudoscalar $`\phi `$ satisfies the linearised equation of motion
$$[dd\phi +\lambda (\lambda 4)\phi ϵ_{\left(5\right)}]Y_{\left(2\right)}=0.$$
(80)
A 2-form harmonic with eigenvalue $`\lambda `$ gives a pseudoscalar $`\phi `$ with $`m^2=\lambda (\lambda 4)`$. We want the mass for the $`10`$ and $`\overline{10}`$ members of the massless multiplet, namely $`m^2=3`$, which therefore requires $`\lambda =1`$ or $`\lambda =3`$. In fact, the required harmonics are those with $`\lambda =3`$ (there are none with $`\lambda =1`$).
There are ten such harmonics on $`S^5`$, which can be written in terms of the Killing spinors. There are Killing spinors $`\eta _\pm `$ satisfying $`D_a\eta _\pm =\pm \frac{\mathrm{i}}{2}\mathrm{\Gamma }_a\eta _\pm `$. It turns out that the required 2-form harmonics are given by the construction
$$Y_{ab}=\overline{\eta }_{}\mathrm{\Gamma }_{ab}\eta _+,$$
(81)
where $`\eta _{}`$ and $`\eta _+`$ are any two Killing spinors of the minus and plus kinds respectively. Solving for the Killing spinors, and substituting into (81), we find that one of the ten harmonics has a structure that is particularly naturally adapted to our parameterisation of the sphere, namely
$$Y_{\left(2\right)}=e^{\mathrm{i}\varphi _3}\left(cd\xi \sigma _3+\frac{1}{2}sc^2\sigma _1\sigma _2\mathrm{i}sc^2\sigma _3d\varphi _3\right).$$
(82)
One may expect that this harmonic, or a closely related construction, will play a significant rôle in the construction of the reduction Ansatz at the full non-linear order, but we have not yet completed this investigation.
### 3.3 Oxidation of five-dimensional solutions
Given the conjectured metric reduction Ansatz, we can oxidise the metric in any solution of the two-scalar truncation of five-dimensional maximal gauged supergravity back to a solution of type IIB supergravity in $`D=10`$. In principle, one can solve the equations of motion in this two-scalar sector to obtain a supersymmetric domain wall solution, which has an interpretation as the RG-flow equations on the strongly coupled field theory side, as discussed in . Unfortunately the equations seem not to allow an explicit solution in terms of elementary functions.
One simple oxidation that we can perform is to take the $`D=5`$ solution corresponding to the second (non-trivial) supersymmetric stationary point of the potential. This corresponds to the stationary point of (61) with
$$X=2^{\frac{1}{3}},\mathrm{sinh}\phi =\frac{1}{\sqrt{3}}.$$
(83)
(The fully-supersymmetric stationary point is at $`X=1`$, $`\phi =0`$.) Substituting into (72), we find that the internal 5-sphere metric $`d\widehat{s}_5^2`$ at this stationary point is given by
$$d\widehat{s}_5^2=\frac{3}{2^{7/3}}\left[d\xi ^2+\frac{c^2}{2(1+s^2)}(\sigma _1^2+\sigma _2^2)+\frac{2c^2}{3+5s^2}\sigma _3^2+\frac{s^2(3+5s^2)}{3(1+s^2)^2}\left(d\varphi _3\frac{c^2}{3+5s^2}\sigma _3\right)^2\right].$$
(84)
## 4 Conclusion
In this paper, we have obtained the metric Ansätze for two examples of Kaluza-Klein sphere reductions, both of which involve pseudoscalar as well as scalar fields. The first example is the $`S^7`$ reduction of eleven-dimensional supergravity, with a truncation from $`N=8`$ to the $`N=2`$ theory with $`U(1)^4`$ gauge fields, three dilatons and three axions. Among other uses this reduction allows one to study the eleven-dimensional geometries corresponding to the lifting of the four-dimensional BPS AdS black hole and domain-wall solutions of gauged supergravity. Our results generalise those obtained previously in , where the problem was studied in the absence of the three axionic scalars.
Our second example is a truncation of five-dimensional maximal gauged supergravity, to a subsector in which two spin-0 fields are retained, one of which is a scalar, and the other a pseudoscalar. This truncation retains the fields necessary for describing a second supersymmetric vacuum in $`D=5`$, with $`N=2`$ supersymmetry and $`SU(2)\times U(1)`$ invariance, in addition to the maximally-supersymmetric one with $`SO(6)`$ invariance . The metric reduction Ansatz that we obtain here allows one to study the ten-dimensional geometries corresponding to the lifting of solutions of the five-dimensional theory. In principle, this can include the renormalisation-group flow associated with the second supersymmetric extremum, although the explicit form of this five-dimensional solution is not known.
## Acknowledgement
We are grateful to Jim Liu, Krzysztof Pilch, Tuan Tran and Nick Warner for discussions.
## Appendix A
In this appendix, we present the explicit form of the $`SO(8)`$ triality rotation that we used in section 2.1 in order to simplify the Kaluza-Klein metric reduction Ansatz:
$`K_{12}\frac{1}{2}(K_{12}+K_{34}+K_{56}+K_{78}),K_{13}\frac{1}{2}(K_{13}K_{24}+K_{57}K_{68}),`$
$`K_{14}\frac{1}{2}(K_{14}+K_{23}+K_{58}+K_{67}),K_{15}\frac{1}{2}(K_{15}K_{26}+K_{37}K_{48}),`$
$`K_{16}\frac{1}{2}(K_{16}+K_{25}+K_{38}+K_{47}),K_{17}\frac{1}{2}(K_{17}K_{28}+K_{35}K_{46}),`$
$`K_{18}\frac{1}{2}(K_{18}+K_{27}+K_{36}+K_{45}),K_{23}\frac{1}{2}(K_{23}+K_{14}K_{58}K_{67}),`$
$`K_{24}\frac{1}{2}(K_{24}K_{13}+K_{57}K_{68}),K_{25}\frac{1}{2}(K_{25}+K_{16}K_{38}K_{47}),`$
$`K_{26}\frac{1}{2}(K_{26}K_{15}+K_{37}K_{48}),K_{27}\frac{1}{2}(K_{27}+K_{18}K_{36}K_{45}),`$
$`K_{28}\frac{1}{2}(K_{28}K_{17}+K_{35}K_{46}),K_{34}\frac{1}{2}(K_{34}+K_{12}K_{56}K_{78}),`$
$`K_{35}\frac{1}{2}(K_{35}+K_{17}+K_{28}+K_{46}),K_{36}\frac{1}{2}(K_{36}+K_{18}K_{27}K_{45}),`$
$`K_{37}\frac{1}{2}(K_{37}+K_{15}+K_{26}+K_{48}),K_{38}\frac{1}{2}(K_{38}+K_{16}K_{25}K_{47}),`$
$`K_{45}\frac{1}{2}(K_{45}+K_{18}K_{27}K_{36}),K_{46}\frac{1}{2}(K_{46}+K_{35}K_{17}K_{28}),`$
$`K_{47}\frac{1}{2}(K_{47}+K_{16}K_{25}K_{38}),K_{48}\frac{1}{2}(K_{48}+K_{37}K_{15}K_{26}),`$
$`K_{56}\frac{1}{2}(K_{56}+K_{12}K_{34}K_{78}),K_{57}\frac{1}{2}(K_{57}+K_{13}+K_{24}+K_{68}),`$
$`K_{58}\frac{1}{2}(K_{58}+K_{14}K_{23}K_{67}),K_{67}\frac{1}{2}(K_{67}+K_{14}K_{23}K_{58}),`$
$`K_{68}\frac{1}{2}(K_{68}+K_{57}K_{13}K_{24}),K_{78}\frac{1}{2}(K_{78}+K_{12}K_{34}K_{56}).`$ (85)
## Appendix B
In this Appendix, we collect some results on the geometry of the 7-sphere. We can describe $`S^7`$ as the unit sphere in $`R^8`$, with 8 real coordinates $`x_I`$;
$$x_Ix_I=1.$$
(86)
As such, it has a manifest $`SO(8)`$ symmetry, with 28 Killing vectors $`K_{IJ}`$ given by
$$K_{IJ}=x^I\frac{}{x_J}x^J\frac{}{x_I}.$$
(87)
We can also describe $`S^7`$ as the unit sphere in $`C^4`$, with 4 complex coordinates $`z_i`$:
$$\overline{z}_iz_i=1.$$
(88)
We can relate these complex coordinates to the previous real ones as follows:
$$z_1=x_1+\mathrm{i}x_2,z_2=x_3+\mathrm{i}x_4,z_3=x_5+\mathrm{i}x_6,z_4=x_7+\mathrm{i}x_8.$$
(89)
We can parameterise these complex coordinates as
$$z_1=\mu _1e^{\mathrm{i}\varphi _1},z_2=\mu _2e^{\mathrm{i}\varphi _2},z_3=\mu _3e^{\mathrm{i}\varphi _3},z_4=\mu _4e^{\mathrm{i}\varphi _4},$$
(90)
where (88) implies that
$$\underset{i=1}{\overset{4}{}}\mu _i^2=1.$$
(91)
These $`(\mu _i,\varphi _i)`$ coordinates are precisely the ones used for describing higher-dimensional rotating black holes in , and in the $`S^7`$ reduction Ansatz obtained in .
From the coordinate transformations above, it is straightforward to establish that the real derivatives $`/x_I`$ that appear in the Killing vectors (87) are given by
$$\frac{}{x_1}=\mathrm{cos}\varphi _1\frac{}{\mu _1}\frac{\mathrm{sin}\varphi _1}{\mu _1}\frac{}{\varphi _1},\frac{}{x_2}=\mathrm{sin}\varphi _1\frac{}{\mu _1}+\frac{\mathrm{cos}\varphi _1}{\mu _1}\frac{}{\varphi _1},$$
(92)
with analogous expressions involving $`(\mu _2,\varphi _2)`$, $`(\mu _3,\varphi _3)`$ and $`(\mu _4,\varphi _4)`$ for the pairs $`(x_3,x_4)`$, $`(x_5,x_6)`$ and $`(x_7,x_8)`$ respectively. It is easy to see from this that the four Killing vectors $`K_{12}`$, $`K_{34}`$, $`K_{56}`$ and $`K_{78}`$ are simply of the form:
$$K_{12}=\frac{}{\varphi _1},K_{34}=\frac{}{\varphi _2},K_{56}=\frac{}{\varphi _3},K_{78}=\frac{}{\varphi _4}.$$
(93)
These are the four commuting $`U(1)`$ generators. It is convenient to write them as $`_{\varphi _1}`$, etc.
We also note that the Killing-vector bilinears in the top 3 lines in (13) are also relatively simple, when expressed in terms of the $`\mu _i`$ and $`\varphi _i`$ coordinates. After some algebra we find, for example, that
$$K_{13}^2+K_{14}^2+K_{23}^2+K_{24}^2=(\mu _1_{\mu _2}\mu _2_{\mu _1})^2+\frac{\mu _2^2}{\mu _1^2}_{\varphi _1}^2+\frac{\mu _1^2}{\mu _2^2}_{\varphi _2}^2$$
(94)
with analogous results for the other five combinations.
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# Contents
## Chapter 1 Introduction
Neutrino physics has become very popular during the past decade. The main impetus in this field has come from the strong improvements in detecting neutrinos; there is now strong evidence from two natural neutrino sources, the solar core and the Earth’s atmosphere , that neutrinos of a certain flavor disappear on their way to the detector. Furthermore, a neutrino beam produced in the laboratory has been found to change flavor in the LSND experiment .
The most natural explanation for these anomalies is achieved by extending the Standard Model with neutrino masses. Together with the Cabibbo-Kobayashi-Maskawa (CKM) matrix for the lepton sector, the masses imply flavor mixing; neutrino oscillations have been born . Such oscillations can be fully parameterized by the differences between the squared masses $`\delta m_{ij}^2m_j^2m_i^2`$ and the mixing angles $`\theta _{ij}`$. Also, a possible phase in the CKM matrix could induce CP violation. With these neutrino oscillations, the experiments can be explained by the solutions presented in Table 1.1.
However, with this straightforward extension of the Standard Model it is not possible to explain all three experiments simultaneously, simply because the trivial condition
$$\delta m_{ij}^2=(m_3^2m_2^2)+(m_2^2m_1^2)+(m_1^2m_3^2)=0$$
(1.1)
cannot be fulfilled by the data. Hence, many interesting new models have been suggested to explain the three different experiments. One of the most appealing possibilities is the existence of a fourth neutrino. Since the $`Z^0`$ decay width excludes more than three standard neutrinos, this new flavor must be inert with respect to the standard weak interaction.
Although the hypothesis of such a sterile neutrino may seem highly speculative, its possible existence is the most far-reaching implication of the current experimental situation. It is certainly worthwhile to investigate its consequences! Sterile neutrinos would have a strong impact on certain astrophysical environments due to their mixing with active neutrinos. These effects would become especially important under extreme conditions, like in core-collapse supernovae and in the Early Universe.
It is the impact of sterile neutrinos in the Early Universe, and in particular the outcome of the primordial nucleosynthesis, that has been the subject of an intense recent debate. By incorporating sterile neutrinos into the well understood mechanism of standard Big Bang nucleosynthesis (BBN), two frontiers of physics would be connected in one phenomenon, making this a powerful tool for predictions.
The production of primordial nuclei depends sensitively on the expansion rate of the Universe, which in turn depends on the number density of sterile neutrinos. This simple correlation is well known ; the actual challenge is to deduce how much the sterile neutrino sector is populated. Since the only connection between sterile neutrinos and ordinary matter, apart from gravity, is given by oscillations, this phenomenon must be the key for creating a relation between sterile neutrino parameters and Big Bang nucleosynthesis predictions.
In the beginning of the 1990’s, the number density for sterile neutrinos at the time of primordial nucleosynthesis was calculated depending on the parameters for oscillations between an active and a sterile neutrino, $`\delta m^2`$ and $`\theta `$ . The idea was to constrain the allowed parameter space, since strong observational results (which have weakened since) set an upper bound on the number density of the $`\nu _\mathrm{s}`$. In these calculations, the influence of the thermal plasma on the mixing parameters was included, but the small CP asymmetric contributions were neglected. The constraints were so strong that they clearly excluded the $`\nu _\mu \nu _\mathrm{s}`$ solution for the atmospheric neutrino anomaly.
In 1995, Foot, Thomson and Volkas found an interesting effect by including the CP asymmetric contributions; the asymmetry could induce different population rates for the sterile neutrinos and antineutrinos, or equivalently, different depopulation rates for the active neutrinos and antineutrinos. This would change the neutrino asymmetry and thus would have an influence on the population rates. They found that this back-reaction could amplify an initially small CP asymmetry by several orders of magnitude for a large range of parameter space. (Actually, this back-reaction had been discussed much earlier by Barbieri and Dolgov , but they erroneously found the effect to be small.)
The implications from this new effect were manifold. The most interesting conclusion was given by a model in which $`\nu _\tau \nu _\mathrm{s}`$ oscillations created a large $`\nu _\tau `$ asymmetry, which in turn suppressed $`\nu _\mu \nu _\mathrm{s}`$ oscillations \[12–17\]. In such a model, the atmospheric neutrino anomaly could be explained by $`\nu _\mu \nu _\mathrm{s}`$ oscillations without the sterile neutrinos coming into thermal equilibrium. Another thoroughly discussed implication treated the impact of a large $`\nu _\mathrm{e}`$ asymmetry on primordial nucleosynthesis ; the large chemical potential $`\mu _\mathrm{e}`$ would change the neutron-to-proton ratio and thus the nuclear abundances.
Clearly, sterile neutrinos can only be included correctly into BBN if the mechanism producing large neutrino asymmetries is revealed. Unfortunately, this has proven to be a very difficult task. A series of papers have been published on this subject by several groups \[12–36\], but the situation remains unclear. The main problem is that the system of differential equations describing the mixing is very complex: one has to treat neutrinos of different momentum separately, and one especially has to take care of the non-linear terms in the equations. Therefore, most works tried to solve the problem with numerical calculations. Thereby, they often used either the adiabatic approximation, which simplifies the equations for a given momentum, or they neglected the momentum distribution. Only during the past year some authors claim to have solved the full system of differential equations numerically .
Due to these difficulties in solving the system, the different works present contradictory solutions: when neglecting the momentum distribution, the neutrino asymmetry shows an oscillating behavior . The system then ends up with a calculable value of the $`\nu `$ asymmetry, but with unpredictable sign. Such a scenario would introduce domains in the Early Universe with different signs of neutrino asymmetry . On the other hand, applying the adiabatic approximation is questionable since the neutrino oscillations are not adiabatic close to the resonance for the parameter space of interest. Here the neutrino asymmetry does not oscillate. Since the numerical calculation using the full system of differential equations is very CPU-time consuming, one also has to question whether the calculations are done with sufficient accuracy and whether the numerics are stable at all. After all, non-linear systems tend to show chaotic behavior. Besides, different results achieved by these calculations are contradictory on the point of oscillating neutrino asymmetry .
All these numerical works have in common that they predict a similar final absolute neutrino asymmetry (if effects of neutrino domains are ignored) of orders $`10^2\text{}1`$. But even this outcome has been questioned very recently in an analytical work by Dolgov et al. , who end up with a final asymmetry several orders of magnitude lower than the former results. Thus, the only point on which all works agree is that oscillations between sterile and active neutrinos in the Early Universe have the potential of creating a neutrino asymmetry at least of order $`10^5`$, a number which is still large compared to the baryon asymmetry, $`\eta =𝒪(10^9)`$.
The main problem is to get a grip on the exact behavior of the non-linear system. Our work will try to bring some clarity into this subject. Most of the former papers were strongly based on numerical calculations. In this paper, the discussion will be based on a thorough analytical treatment, which will be supported by simple numerical calculations. We will describe the evolution of the system, neglecting the momentum distribution, and compare our results with a numerical solution for the parameters $`(\delta m^2,\mathrm{sin}2\theta _0)=(1\mathrm{eV}^2,10^4)`$. It should be the job of future investigations to include the effects from the momentum distribution.
In Chapter 2, those aspects of the Early Universe and Big Bang nucleosynthesis will be summarized which are important for our system. Chapter 3 contains all relevant information about neutrino oscillations in matter. In Chapter 4, we analyze the system in the simple two flavor case. After introducing initial conditions, we prove that our system creates oscillations of the lepton asymmetry. In Chapter 5, we summarize our findings.
## Chapter 2 Big Bang Nucleosynthesis
We study the influence of neutrinos, in particular their effective number of flavors and the $`\nu _\mathrm{e}`$ chemical potential, on Big Bang nucleosynthesis (BBN). To this end we first describe the expansion of the Early Universe before and during BBN. We define number density and lepton asymmetry. Next, we give a short summary of BBN, concentrating on the main product, <sup>4</sup>He, and we estimate the influence the neutrino sector can have on its abundance. Finally, we summarize the observational results.
### 2.1 Dynamics in the Radiation Epoch
In order to discuss the influence of neutrinos on BBN we need to introduce some general concepts pertaining to the relevant cosmological epoch. Naturally, we only need to look at the Universe at times before and during BBN, i.e. at temperatures above 0.1 MeV. On the other hand, for the mixing parameters we will consider, neutrino oscillations are strongly damped at temperatures much higher than the QCD phase transition scale $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ of around $`200\mathrm{MeV}`$. The temperatures of interest are therefore between $`200\mathrm{MeV}`$ and $`0.1\mathrm{MeV}`$, corresponding to cosmological time scales from $`10^5`$ s to $`3`$ min.
The Universe during this epoch is flat, homogeneous and isotropic, and can therefore be described by the Robertson-Walker metric. The Friedmann equation for these conditions is simply
$$H^2(t)=\frac{8\pi G}{3}\rho (t),$$
(2.1)
where $`H(t)=\dot{R}(t)/R(t)`$ is the Hubble parameter or expansion rate, $`R`$ the cosmic scale factor, $`G=m_{\mathrm{pl}}^2`$ Newton’s constant, and $`m_{\mathrm{pl}}=1.22\times 10^{19}\mathrm{GeV}`$ the Planck mass. Because the energy density $`\rho `$ is dominated by radiation, it scales as $`R^4`$, which upon integrating (2.1) gives $`H=\frac{1}{2}t^1`$.
The radiation density is usually expressed in the form
$$\rho =\frac{\pi ^2}{30}g_{}T^4,$$
(2.2)
where $`T`$ is the photon temperature and $`g_{}`$ the total effective number of relativistic degrees of freedom,
$$g_{}=\underset{i=\mathrm{bosons}}{}g_i\left(\frac{T_i}{T}\right)^4+\frac{7}{8}\underset{i=\mathrm{fermions}}{}g_i\left(\frac{T_i}{T}\right)^4.$$
(2.3)
Here, $`g_i`$ is the number of internal degrees of freedom of particle species $`i`$ and $`T_i`$ is its temperature. If we insert (2.2) into (2.1) and use $`t=1/2H`$, we finally derive a relationship between time and temperature:
$$t=\frac{1}{2}\sqrt{\frac{90}{8\pi ^3}}g_{}^{1/2}\frac{m_{\mathrm{pl}}}{T^2}=0.301g_{}^{1/2}\frac{m_{\mathrm{pl}}}{T^2}.$$
(2.4)
As another consequence, we note that $`TR^1`$ as long as $`g_{}=\text{const}`$.
In the range $`1\mathrm{MeV}\mathrm{\Gamma }<T\mathrm{\Gamma }<\mathrm{\hspace{0.17em}100}\mathrm{MeV}`$, the only particles which contribute significantly to the radiation density are photons ($`g_\gamma =2`$), electrons and positrons ($`g_e=4`$), and three left-handed neutrino families with $`g_\nu =2`$ each. For all of them $`T_i=T`$ applies as long as they remain in thermal equilibrium so that $`g_{}=10.75`$. Any additional neutrino species would add another $`7/8`$ per internal degree of freedom. For $`T\mathrm{\Gamma }>\mathrm{\hspace{0.17em}100}\mathrm{MeV}`$, $`g_{}`$ is higher due to the presence of muons and pions, and for $`T\mathrm{\Gamma }>\mathrm{\Lambda }_{\mathrm{QCD}}200\mathrm{MeV}`$ many gluon and quark degrees of freedom are excited. For $`T\mathrm{\Gamma }<\mathrm{\hspace{0.17em}1}\mathrm{MeV}`$, the electrons and positrons become non-relativistic and annihilate, reducing $`g_{}`$ by $`(7/8)g_e=7/2`$. This effect heats the photons (the neutrinos are already decoupled) so that $`T_\nu /T`$ is reduced to $`(4/11)^{1/3}`$, leading to $`g_{}=3.36`$.
In our subsequent discussion we will ignore these effects and always use the value $`g_{}=10.75`$, implying
$$H=\frac{1}{2}t^1=5.5\frac{T^2}{m_{\mathrm{pl}}}.$$
(2.5)
This is a good approximation as long as the oscillation parameters are such that the crucial events happen in the range $`1\mathrm{MeV}\mathrm{\Gamma }<T\mathrm{\Gamma }<\mathrm{\hspace{0.17em}100}\mathrm{MeV}`$.
### 2.2 Number Densities and <br>Fermion Asymmetries
Because we want to study the influence of neutrino number densities and matter-antimatter asymmetries, we introduce appropriate measures for them. It is convenient to normalize these quantities to the photon number density
$$n_\gamma =2\frac{\zeta _3}{\pi ^2}T^3,$$
(2.6)
where $`\zeta `$ is the Riemann zeta function with $`\zeta _31.202`$.
Fermions may have non-vanishing chemical potentials $`\mu _i`$ which affect their number densities. In addition, a non-vanishing $`\mu _i`$ implies a CP asymmetry for species $`i`$ which we parameterize as
$$A_i\frac{n_i\overline{n}_i}{n_\gamma },$$
(2.7)
where $`\overline{n}_i`$ refers to the antiparticle density.
At temperatures below the QCD phase transition virtually no hadronic antimatter exists. Since there are no baryon-number violating interactions at this late epoch, the total baryon number is conserved. Therefore, the number density of baryons, $`n_\mathrm{B}`$, scales as $`R^3`$ or $`T^3`$. The baryon asymmetry
$$A_\mathrm{B}=\frac{n_\mathrm{B}\overline{n}_\mathrm{B}}{n_\gamma }\frac{n_\mathrm{B}}{n_\gamma }$$
(2.8)
is thus constant under the cosmic expansion.
The commonly used present-time baryon asymmetry $`\eta `$ is related to our $`A_\mathrm{B}`$ through
$$\eta =\frac{4}{11}A_\mathrm{B}=1\text{ to }6\times 10^{10}.$$
(2.9)
The difference is due to photon heating when electrons and positrons annihilate, an effect which modifies all asymmetries. We will always take $`A_i`$ to refer to the fermion asymmetry of species $`i`$ at the epoch before $`\mathrm{e}^+\mathrm{e}^{}`$ annihilation.
For $`T<\mathrm{\Lambda }_{\mathrm{QCD}}`$, almost all baryons are either neutrons n or protons p, which implies that $`n_\mathrm{B}=n_\mathrm{n}+n_\mathrm{p}`$. The non-trivial evolution of $`n_\mathrm{n}`$ and $`n_\mathrm{p}`$ will be discussed in the next section.
We will assume that the Universe is charge neutral. Neglecting muons and pions, we therefore use
$$A_\mathrm{e}=A_\mathrm{p}$$
(2.10)
for the electron asymmetry.
Turning to neutrinos, we can safely neglect their mass. Therefore, their number densities are given by
$$n_\nu =\frac{d^3𝐩}{(2\pi )^3}f(p,T,\mu ),$$
(2.11)
where $`E=p=|𝐩|`$, $`f(p,T,\mu )=[e^{(p\mu )/T}+1]^1`$ is the Fermi-Dirac phase-space distribution function, and $`\mu `$ is the chemical potential. For very small asymmetries ($`\mu T`$) we find
$$n_\nu +\overline{n}_\nu =\frac{3}{4}n_\gamma +𝒪\left((\mu /T)^2\right)$$
(2.12)
and
$$A_\nu =\frac{\pi ^2}{12\zeta _3}\frac{\mu }{T}+𝒪\left((\mu /T)^3\right)0.68\frac{\mu }{T}+𝒪\left((\mu /T)^3\right),$$
(2.13)
where we have used that the chemical potential for antineutrinos is $`\overline{\mu }=\mu `$. Of course, these equations are only valid as long as the neutrinos are in thermal equilibrium.
### 2.3 Helium Abundance
Big Bang nucleosynthesis is a rather complex system depending on a number of parameters, including the baryon asymmetry $`\eta `$ and the effective number of neutrino families $`N_\nu ^{\mathrm{eff}}`$, or more generally, $`g_{}`$. This system has been analyzed very thoroughly, including exhaustive numerical calculations. Here we will give a short summary, concentrating on the production of <sup>4</sup>He.
The outcome of BBN will be abundances of the different species $`(A,Z)`$ with $`Z`$ protons and $`AZ`$ neutrons. We will express these abundances as the mass fraction contributed by the species, i.e.
$$X_{A,Z}\frac{An_{A,Z}}{n_\mathrm{B}},$$
(2.14)
where $`n_{A,Z}`$ and $`n_\mathrm{B}`$ are the number densities of the nuclear species $`(A,Z)`$ and all baryons, respectively.
In strict thermodynamic equilibrium, the abundances are given by nuclear statistical equilibrium (NSE) so that
$$(X_{A,Z})_{\mathrm{NSE}}\eta ^{A1}\mathrm{exp}\left(\frac{B_{A,Z}}{T}\right),$$
(2.15)
where $`B_{A,Z}`$ is the binding energy of the nuclear species $`(A,Z)`$. We see that the NSE abundances are strongly suppressed by the small baryon asymmetry $`\eta 10^9`$. This suppression is compensated by the exponential factor in (2.15) at low temperatures $`TB_{A,Z}`$. For <sup>4</sup>He, $`(X_{\mathrm{He}})_{\mathrm{NSE}}𝒪(1)`$ at $`T0.3\mathrm{MeV}`$. Heavier nuclei have significant NSE abundances only at even lower temperatures.
The NSE abundance for an element X will freeze out when the rates $`\mathrm{\Gamma }(\mathrm{ab}\mathrm{X})`$ of producing it from the lighter elements a and b become smaller than the expansion rate $`H`$. We have
$$\mathrm{\Gamma }(\text{ab}\text{X})n_\mathrm{a}n_\mathrm{b}\mathrm{exp}\left[2\left(\frac{k}{T_{\mathrm{MeV}}}\right)^{1/3}\right],$$
(2.16)
where $`kZ_\mathrm{a}^2Z_\mathrm{b}^2A_\mathrm{a}A_\mathrm{b}/(A_\mathrm{a}+A_\mathrm{b})`$ and $`T_{\mathrm{MeV}}=T/\mathrm{MeV}`$. The first two factors are the number densities of the nuclear species a and b, respectively, the last factor represents the Coulomb-barrier suppression, which increases with $`A_i`$ and $`Z_i`$. The freeze-out temperatures increase with the Coulomb-barrier suppression.
As the temperature decreases, the NSE abundances increase, but at the same time the nuclear reactions begin to freeze out. Therefore, heavy nuclei are not produced during BBN because they freeze out long before their NSE abundances have become significant.
Apart from traces of other nuclei, BBN produces primarily <sup>4</sup>He so that it is a good approximation to assume that all neutrons end up in <sup>4</sup>He. Then
$$YX_{\mathrm{He}}\frac{4n_{\mathrm{He}}}{n_\mathrm{B}}=\frac{4(n_\mathrm{n}/2)}{n_\mathrm{n}+n_\mathrm{p}}=\frac{2(n_\mathrm{n}/n_\mathrm{p})_{\mathrm{BBN}}}{1+(n_\mathrm{n}/n_\mathrm{p})_{\mathrm{BBN}}}.$$
(2.17)
Therefore, the all-important helium mass fraction $`Y`$ depends primarily on the n/p ratio at the time when <sup>4</sup>He freezes out of NSE, which happens at $`t_{\mathrm{BBN}}=`$ 13 minutes.
To find $`(n_\mathrm{n}/n_\mathrm{p})_{\mathrm{BBN}}`$, we have to follow the evolution of the n/p ratio from the beginning. At temperatures above $`1\mathrm{MeV}`$, reactions of the type
$$\mathrm{n}+\nu _\mathrm{e}\mathrm{p}+\mathrm{e},$$
(2.18)
maintain chemical equilibrium with the rate $`\mathrm{\Gamma }G_\mathrm{F}T^5`$. Therefore
$$\left(\frac{n_\mathrm{n}}{n_\mathrm{p}}\right)=\mathrm{exp}\left[\frac{Q}{T}+\frac{\mu _\mathrm{e}}{T}\frac{\mu _{\nu _\mathrm{e}}}{T}\right],$$
(2.19)
where $`Q=1.293\mathrm{MeV}`$ is the mass difference between neutrons and protons, and $`\mu _\mathrm{e}`$ and $`\mu _{\nu _\mathrm{e}}`$ are the chemical potentials of the electrons and electron neutrinos, respectively. Surely, $`\mu _\mathrm{e}`$ can be neglected since $`\mu _\mathrm{e}/T\eta 10^9`$. For the moment we will set the neutrino asymmetry to zero which is equivalent to $`\mu _{\nu _\mathrm{e}}=0`$. Then the n/p ratio depends only on the temperature.
At $`T_{\mathrm{fr}}0.8\mathrm{MeV}`$, the n/p ratio freezes out and thus
$$\left(\frac{n_\mathrm{n}}{n_\mathrm{p}}\right)_{\mathrm{fr}}\mathrm{exp}(Q/T_{\mathrm{fr}})0.2.$$
(2.20)
More exact calculations give $`(n_\mathrm{n}/n_\mathrm{p})_{\mathrm{fr}}1/6`$.
After freeze-out, the n/p ratio continues to decrease slowly due to neutron decay,
$$\mathrm{n}\mathrm{p}+\mathrm{e}^{}+\overline{\nu }_\mathrm{e}.$$
(2.21)
Therefore,
$$\left(\frac{n_\mathrm{n}}{n_\mathrm{p}}\right)_{\mathrm{BBN}}\left(\frac{n_\mathrm{n}}{n_\mathrm{p}}\right)_{\mathrm{fr}}\mathrm{exp}\left[\mathrm{ln}(2)\frac{t_{\mathrm{BBN}}}{\tau _\mathrm{n}}\right]\frac{1}{7},$$
(2.22)
where $`\tau _\mathrm{n}=886.7\pm 1.9\mathrm{s}`$ is the neutron half-life. If we insert this value into (2.17), we get $`Y0.25`$.
### 2.4 Non-Standard Neutrinos
We now want to derive the influence of neutrinos on the helium abundance. According to (2.19), a non-zero electron neutrino asymmetry changes the n/p ratio by a factor
$$\mathrm{exp}(\mu _{\nu _\mathrm{e}}/T_{\mathrm{fr}})\mathrm{exp}(1.5A_{\nu _\mathrm{e}})(11.5A_{\nu _\mathrm{e}})$$
(2.23)
provided that $`A_{\nu _\mathrm{e}}1`$. As an example, we take $`A_{\nu _\mathrm{e}}=\pm 0.01`$, a realistic value according to . Then $`(n_\mathrm{n}/n_\mathrm{p})`$ is altered by a factor $`(10.015)`$. Of course, this is a very rough estimate, but more detailed works find the same order of magnitude.
A higher effective neutrino number also alters $`(n_\mathrm{n}/n_\mathrm{p})`$, since it raises $`g_{}`$, which results in a higher expansion rate $`H`$ and thus in a higher freeze-out temperature. Taking $`N_\nu ^{\mathrm{eff}}=4`$ instead of 3 as an example, $`T_{\mathrm{fr}}`$ is increased by a factor of
$$\left(\frac{g_{}(N_\nu ^{\mathrm{eff}}=4)}{g_{}(N_\nu ^{\mathrm{eff}}=3)}\right)^{1/6}=\left(\frac{10.75+2\frac{7}{8}}{10.75}\right)^{1/6}=1.025,$$
(2.24)
where the power $`1/6`$ comes from $`\mathrm{\Gamma }/HT^3/\sqrt{g_{}}`$. This changes the n/p ratio by a factor of 1.04.
We insert these changes into (2.17) and expand to get
$$Y0.25+0.01(N_\nu ^{\mathrm{eff}}3)0.33A_{\nu _\mathrm{e}}$$
(2.25)
if $`A_{\nu _\mathrm{e}}1`$ and $`|N_\nu ^{\mathrm{eff}}3|\mathrm{\Gamma }<\mathrm{\hspace{0.17em}1}`$. We see that $`N_\nu ^{\mathrm{eff}}=4`$ or $`|A_{\nu _\mathrm{e}}|=10^2`$ change the helium abundance by several percent, which is within the present experimental precision. More detailed numerical calculations give
$$Y0.225+0.025\mathrm{log}(\eta /10^{10})+0.012(N_\nu ^{\mathrm{eff}}3),$$
(2.26)
where we have also included the influence of the baryon asymmetry $`\eta `$, since the main uncertainty in Standard BBN comes from this parameter.
We should mention that in the literature, the variable $`N_\nu ^{\mathrm{eff}}`$ has often been used not only to account for the effective number of neutrinos, but also for the influence of $`A_{\nu _\mathrm{e}}`$. This was done by adding a new term $`\delta N_\nu ^{\mathrm{eff}}(A_{\nu _\mathrm{e}})`$ to $`N_\nu ^{\mathrm{eff}}`$.
### 2.5 Observational Results
The best measurements of the primordial <sup>4</sup>He abundance, Y, come from observing extra-galactic regions of ionized H. In these systems, the abundance of heavier elements, which are not created in BBN, is very low, so we can assume that the abundances in these regions are close to their primordial values. The present estimate is
$$Y=0.238\pm 0.002\pm 0.005,$$
(2.27)
where the two errors are the statistical and systematic errors, respectively. The $`2\sigma `$ range is then estimated to be 0.228–0.248.
Direct present-time measurements of the cosmic baryon abundance are very uncertain due to the dark matter problem. Much better estimates arise from measurements of the primordial Deuterium abundance which depends sensitively on the baryon asymmetry $`\eta `$. The best measurements come from the absorption of quasar light by high-redshift, low-metallicity hydrogen clouds. Two main results have been published ,
$$(\mathrm{D}/\mathrm{H})_{\mathrm{low}}=(3.4\pm 0.3)\times 10^5\text{and}(\mathrm{D}/\mathrm{H})_{\mathrm{high}}=(2\pm 0.5)\times 10^4,$$
(2.28)
which are mutually inconsistent. These two measurements give ranges for the baryon asymmetry of
$$\eta _{\mathrm{low}}=4.2\text{}6.3\times 10^{10}\text{and}\eta _{\mathrm{high}}=1.2\text{}2.8\times 10^{10},$$
(2.29)
at a nominal $`2\sigma `$ level.
From the measurements of D/H and $`Y`$, one obtains bounds on the effective neutrino number. Olive et al. derived
$$(N_\nu ^{\mathrm{eff}}3)_{\mathrm{low}}<0.3\text{and}(N_\nu ^{\mathrm{eff}}3)_{\mathrm{high}}<1.8,$$
(2.30)
provided that the $`\nu _\mathrm{e}`$ chemical potential can be neglected.
The low-D result today appears to be strongly favored. If it should be confirmed, sterile neutrinos would be forbidden to come into equilibrium for negligible $`\mu _{\nu _\mathrm{e}}`$. However, a large positive $`A_{\nu _\mathrm{e}}`$ can compensate the effect of an increased $`N_\nu ^{\mathrm{eff}}`$, and can thus circumvent this constraint.
### 2.6 Summary
We have shown that a deviation of the effective number of neutrinos of order 1, as well as a $`\nu _\mathrm{e}`$ asymmetry exceeding about $`10^2`$, will have a measurable effect on the outcome of BBN. Sterile neutrinos have the potential to change both $`N_\nu ^{\mathrm{eff}}`$ and $`A_{\nu _\mathrm{e}}`$. If we knew how these two variables depend on the mixing parameters of the sterile neutrinos, we could use the measured primordial element abundances to derive constraints on the mixing parameters. On the other hand, if the existence of sterile neutrinos was proven by future experiments such as MiniBooNE , a detailed understanding of their impact on primordial nucleosynthesis would be necessary to constrain the free parameters of BBN, in particular the baryon asymmetry $`\eta `$.
## Chapter 3 Neutrino Oscillations
Neutrino oscillations play a crucial role in our considerations of the Early Universe involving sterile neutrinos. In this chapter, we derive the density matrix formalism for neutrino oscillations in media between any two neutrino flavors, active or sterile. This includes the direct influence of matter on the vacuum neutrino oscillations as well as scattering processes which tend to destroy the coherence of the oscillations. Our analysis will only be applicable for temperatures between a few$`\mathrm{MeV}`$ and $`100\mathrm{MeV}`$ and for neutrino asymmetries $`A_\nu 1`$.
### 3.1 Equation of Motion
Neutrino oscillations occur because the basis of the neutrino weak eigenstates $`\nu _\alpha `$, $`\alpha =\mathrm{e},\mu ,\tau ,\mathrm{s},\mathrm{}`$, is different from the basis of the neutrino mass eigenstates $`\nu _i`$, $`i=1,2,3,\mathrm{}`$. In other words, a $`\nu _\alpha `$ that is produced in a weak-interaction process does not propagate like a free particle, but as a superposition of neutrinos $`\nu _i`$ with different masses $`m_i`$, respectively. Thus, when measured, the propagated neutrino contains contributions of weak eigenstates other than the original $`\nu _\alpha `$. From the start we include the possibility that neutrinos exist beyond the active states $`\nu _\mathrm{e}`$, $`\nu _\mu `$ and $`\nu _\tau `$. These additional flavors would have to be sterile with regard to the weak interaction.
The two bases are connected by a unitary transformation $`\mathrm{\Psi }_\mathrm{W}=U\mathrm{\Psi }_\mathrm{M}`$, where
$$\mathrm{\Psi }_\mathrm{W}\left(\begin{array}{c}\mathrm{\Psi }_{\nu _\mathrm{e}}\\ \mathrm{\Psi }_{\nu _\mu }\\ \mathrm{\Psi }_{\nu _\tau }\\ \mathrm{}\end{array}\right)\text{and}\mathrm{\Psi }_\mathrm{M}\left(\begin{array}{c}\mathrm{\Psi }_{\nu _1}\\ \mathrm{\Psi }_{\nu _2}\\ \mathrm{\Psi }_{\nu _3}\\ \mathrm{}\end{array}\right)$$
(3.1)
are the field vectors represented in the weak and mass basis, respectively, and $`U`$ is the unitary transformation matrix.
We are mainly interested in the evolution of the weak eigenstates, as these are the particles we can produce and measure by weak-interaction processes. Therefore, we write the equation of motion in the basis of the weak eigenstates. The Klein-Gordon equation is
$$(_t^2^2+M_W^2)\mathrm{\Psi }_W=0,$$
(3.2)
where $`M_\mathrm{W}=UM_\mathrm{M}U^{}`$ and $`M_\mathrm{M}=\text{diag}(m_1,m_2,\mathrm{})`$ is the mass matrix in the weak and mass basis, respectively. Of course, $`M_\mathrm{W}^2=UM_\mathrm{M}U^{}UM_\mathrm{M}U^{}=UM_\mathrm{M}^2U^{}`$ and $`M_\mathrm{W}`$ are not diagonal in the case of mixing.
Since the Early Universe is homogeneous, we are only interested in the time evolution of the neutrino fields. It is therefore convenient to expand them in plane waves $`\mathrm{\Psi }_\mathrm{W}=\mathrm{\Psi }_p(t)e^{i\mathrm{𝐩𝐱}}`$. Note that usually neutrino oscillations are considered in environments which are spatially inhomogeneous but stationary, e.g. in experiments involving solar or atmospheric neutrinos. In such cases, $`\mathrm{\Psi }_\mathrm{W}`$ would be expanded in components of fixed energy, $`\mathrm{\Psi }_\mathrm{W}=\mathrm{\Psi }_E(𝐱)e^{iEt}`$, instead of components of fixed momentum.
Since the neutrinos are highly relativistic, $`m_\nu Ep`$, we can linearize the Klein-Gordon equation. Then we get the usual Schrödinger-type equation
$$i_t\mathrm{\Psi }_p=\mathrm{\Omega }_p\mathrm{\Psi }_p,\mathrm{\Omega }_p=p+\frac{M_\mathrm{W}^2}{2p}.$$
(3.3)
We see that the off-diagonal elements of $`M_\mathrm{W}^2`$ couple the fields of the different neutrino flavors, leading to oscillations.
An equivalent equation is given by
$$i_t\rho =[\rho ,\mathrm{\Omega }],$$
(3.4)
where $`\rho _{\alpha \beta }N\mathrm{\Psi }_\alpha \mathrm{\Psi }_\beta ^{}`$ is the flavor density matrix, $`N`$ is a normalization factor, and we have dropped the index $`p`$. The advantage of the density matrix formalism is that we later can include effects that destroy the coherence of the neutrino oscillations.
For later convenience, we normalize the density matrix such that for a CP symmetric neutrino species $`\nu _\alpha `$ in equilibrium, $`\rho _{\alpha \alpha }=1`$ for all momentum modes, i.e.
$$n_{\nu _\alpha }=\frac{d^3𝐩}{(2\pi )^3}\mathrm{\Psi }_\alpha \mathrm{\Psi }_\alpha ^{}=\frac{d^3𝐩}{(2\pi )^3}f_0(p,T)\rho _{\alpha \alpha }(p),$$
(3.5)
where $`f_0(p,T)=f(p,T,\mu =0)=[1+\mathrm{exp}(p/T)]^1`$. Thus, the normalization factor $`N=1/f_0(p,T)`$.
We have neglected the expansion of the Universe in the derivation of the equation of motion. It could be included by adding a term $`\mathrm{\Omega }_H=iHp_p`$ to the Hamiltonian. The term disappears if we expand the field vector in comoving plane waves, $`\mathrm{\Psi }_W=\mathrm{\Psi }_q(t)e^{iT\mathrm{𝐪𝐱}}`$, where $`q=p/T`$, instead of plane waves with fixed momentum. Therefore, later in Chapter 4 we will not describe the evolution of flavor density matrices of fixed momentum $`p`$, but of fixed comoving momentum $`q`$.
### 3.2 Medium Effects
In the Early Universe, we are confronted with a thermal medium which interacts with the neutrinos. Therefore, we need to include the effects of these interactions on the neutrino oscillations. In this section, we restrict ourselves to the discussion of the refractive effects .
The medium contributions to the neutrino oscillations enter the Schrödinger equation (3.3) through the weak-potential term $`V\text{diag}(V_{\nu _\alpha },V_{\nu _\beta },\mathrm{})`$ in
$$\mathrm{\Omega }=p+\frac{M_\mathrm{W}^2}{2p}+V.$$
(3.6)
For each neutrino weak eigenstate the contributions can be split into two terms,
$$V_{\nu _\alpha }=\pm V_{\nu _\alpha }^AV_{\nu _\alpha }^T.$$
(3.7)
Here, the plus sign is valid for neutrinos, while the minus sign applies to antineutrinos.
The first term accounts for the fermion asymmetries. For neutrinos of species $`\alpha `$ it is given by
$$V_{\nu _\alpha }^A=\sqrt{2}G_\mathrm{F}n_\gamma \stackrel{~}{A}_{\nu _\alpha },$$
(3.8)
where $`G_\mathrm{F}=1.166\times 10^5\mathrm{GeV}^2`$ is Fermi’s constant and $`\stackrel{~}{A}_{\nu _\alpha }`$ a weighted sum over all fermion asymmetries $`A_i`$. Generally,
$`\stackrel{~}{A}_{\nu _\alpha }`$ $`=`$ $`A_{\nu _\alpha }+A_{\nu _\mathrm{e}}+A_{\nu _\mu }+A_{\nu _\tau }`$ (3.9)
$`+A_\alpha {\displaystyle \frac{1}{2}}(14\mathrm{sin}^2\theta _\mathrm{W})(A_\mathrm{e}+A_\mu +A_\tau )`$
$`+{\displaystyle \frac{1}{2}}(14\mathrm{sin}^2\theta _\mathrm{W})A_\mathrm{p}{\displaystyle \frac{1}{2}}A_\mathrm{n},`$
where $`\theta _\mathrm{W}`$ is the Weinberg angle. For our assumption of a charge-neutral Universe we have $`A_\mathrm{e}+A_\mu +A_\tau =A_\mathrm{p}`$ and therefore
$$\stackrel{~}{A}_{\nu _\alpha }=A_{\nu _\alpha }+A_{\nu _\mathrm{e}}+A_{\nu _\mu }+A_{\nu _\tau }\frac{1}{2}A_\mathrm{n}+A_\alpha .$$
(3.10)
Since the muons and tauons are non-relativistic, their asymmetries are negligible, so that we will use $`A_\mathrm{e}=A_\mathrm{p}`$, $`A_\mu =0`$, and $`A_\tau =0`$. In our analysis, all of the asymmetries but for one type of neutrino will remain constant. Therefore, we write
$$\stackrel{~}{A}_{\nu _\alpha }=2A_{\nu _\alpha }+A_\mathrm{c},$$
(3.11)
where all asymmetries other than $`A_{\nu _\alpha }`$ have been absorbed in a constant $`A_\mathrm{c}`$.
The second term in (3.7) represents the low-energy tail of the $`W^\pm `$ and $`Z^0`$ resonances. It has the remarkable feature that it is independent of the CP asymmetry of the background medium and has the same effect on neutrinos and antineutrinos. Since it is of order $`G_\mathrm{F}^2`$, $`V^T`$ can only compete with $`V^A`$ in the Early Universe, where $`A1`$. For neutrino temperatures far below the $`W^\pm `$ and $`Z^0`$-masses it can be written in the form
$$V_{\nu _\alpha }^T=z_{\nu _\alpha }G_\mathrm{F}^2pTn_\gamma .$$
(3.12)
Here, the coefficients are
$`z_{\nu _\mathrm{e}}`$ $`=`$ $`{\displaystyle \frac{7\pi ^3}{45\zeta _3\alpha }}\mathrm{sin}^2\theta _\mathrm{W}\left(\mathrm{cos}^2\theta _\mathrm{W}+2\right)350`$
$`z_{\nu _\mu ,\nu _\tau }`$ $`=`$ $`{\displaystyle \frac{7\pi ^3}{45\zeta _3\alpha }}\mathrm{sin}^2\theta _\mathrm{W}\mathrm{cos}^2\theta _\mathrm{W}97,`$ (3.13)
where $`\alpha =1/137`$ is the fine-structure constant.
Since we will be discussing sterile neutrinos, we stress their special role in a medium. They do not interact, so $`V_{\nu _\mathrm{s}}=0`$. Likewise, a possible asymmetry $`A_{\nu _\mathrm{s}}`$ does not contribute to $`V_{\nu _\alpha }^A`$ for the active flavor.
### 3.3 Density Matrix Formalism for <br>Two-Flavor Oscillations
We will now treat the simplest case of neutrino oscillations between two types of neutrinos, so $`\mathrm{\Psi }_\mathrm{W}=(\mathrm{\Psi }_{\nu _\alpha },\mathrm{\Psi }_{\nu _\beta })`$. The unitary transformation can be written in the simple form
$$U=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right),$$
(3.14)
where $`\theta `$ is the vacuum mixing angle. Note that the absolute phase of the off-diagonal terms is arbitrary; our convention differs from some of the literature. In the weak basis we get
$$\mathrm{\Omega }=\frac{1}{2}bI\frac{1}{2}\left(\begin{array}{cc}c\mathrm{\Delta }+\delta V& s\mathrm{\Delta }\\ s\mathrm{\Delta }& (c\mathrm{\Delta }+\delta V)\end{array}\right),$$
(3.15)
where $`b=2p+\frac{m_1^2+m_2^2}{2p}(V_{\nu _\alpha }+V_{\nu _\beta })`$, $`I`$ is the $`2\times 2`$ unit matrix, and $`\delta V=V_{\nu _\alpha }V_{\nu _\beta }`$. Furthermore, we have defined
$$s\mathrm{sin}2\theta ,c\mathrm{cos}2\theta ,\mathrm{\Delta }\frac{\delta m^2}{2p}$$
(3.16)
with the mass squared difference $`\delta m^2m_2^2m_1^2`$.
To find a convenient form for the equation of motion, we represent $`\rho `$ and $`\mathrm{\Omega }`$ by the unit matrix $`I`$ and the Pauli matrices $`\sigma _i`$, $`i=1,2,3`$ . Then
$$\rho =\frac{1}{2}P_0I+\frac{1}{2}𝐏𝝈,\mathrm{\Omega }=\frac{1}{2}bI\frac{1}{2}𝐁𝝈,$$
(3.17)
where $`𝐏`$ is a so-called flavor polarization vector, and
$$𝐁=\left(\begin{array}{c}s\mathrm{\Delta }\\ 0\\ c\mathrm{\Delta }+\delta V\end{array}\right).$$
(3.18)
Inserting $`\rho `$ and $`\mathrm{\Omega }`$ into (3.4) and using the commutation relation for the Pauli matrices, $`[\frac{1}{2}\sigma _i,\frac{1}{2}\sigma _j]=iϵ_{ijk}\frac{1}{2}\sigma _k`$, we derive
$$_t𝐏=𝐁\times 𝐏.$$
(3.19)
This equation is equivalent to the precession of a spin vector in a magnetic field. The physical observables are now contained in $`P_z=\rho _{\alpha \alpha }\rho _{\beta \beta }`$ and $`P_0=\rho _{\alpha \alpha }+\rho _{\beta \beta }`$. $`f_0(p,T)P_0`$ is the total occupation number of both $`\nu _\alpha `$ and $`\nu _\beta `$ of momentum mode p, while $`P_z`$ represents the difference between the two neutrino types; if e.g. $`P_z=1`$ for all modes, then $`n_{\nu _\beta }=0`$. However, if $`P_z=0`$ for all modes, $`n_{\nu _\alpha }=n_{\nu _\beta }`$. The other two components $`P_{x,y}`$ represent the phase of the oscillations.
### 3.4 Scattering Processes
We must still take into account that the weak neutrino eigenstates scatter on the background plasma with the rate
$$\mathrm{\Gamma }_{\nu _\alpha }=y_{\nu _\alpha }G_\mathrm{F}^2pT^4,$$
(3.20)
where $`y_{\nu _\mathrm{e}}1.13`$, $`y_{\nu _\mu ,\nu _\tau }0.79`$ and $`y_{\nu _\mathrm{s}}0`$. To be precise, $`y_{\nu _\alpha }`$ also depends on the momentum due to the Pauli blocking factors. We will use these momentum averaged values .
Scattering keeps the active neutrinos in thermal and chemical equilibrium at temperatures above a few MeV. Therefore, the integrand in (3.5), $`f_0(p,T)\rho _{\alpha \alpha }(p)`$, should be equal to $`f(p,T,\mu )`$. However, $`\rho _{\alpha \alpha }(p)`$ changes due to the neutrino oscillations. The scattering processes compensate this change by re- or depopulating the number densities through $`\dot{\rho }_{\alpha \alpha }(p)=R_{\nu _\alpha }(p)`$, where the rate is approximately given by
$$R_{\nu _\alpha }(p)\mathrm{\Gamma }_{\nu _\alpha }\left[\frac{f(p,T,\mu )}{f_0(p,T)}\rho _{\alpha \alpha }(p)\right].$$
(3.21)
The approximation holds as long as the deviation from equilibrium is small. Of course, $`R_{\nu _\mathrm{s}}=\mathrm{\Gamma }_{\nu _\mathrm{s}}=0`$.
A much more interesting effect appears if the scattering rates are different for the two neutrino types considered. Then scattering distinguishes between the two $`\nu `$ types and thus the coherence of the neutrino oscillations is destroyed. Therefore, we need to include a term $`D𝐏_{}(p)`$, where $`𝐏_{}=(P_x,P_y,0)`$ and $`D`$ is called the damping rate. In the case of an active neutrino oscillating with a sterile neutrino ($`\beta =s`$) we have $`D=\frac{1}{2}\mathrm{\Gamma }_{\nu _\alpha }`$. The lengthy derivation of this term can be read in .
In summary, for each neutrino momentum mode $`p`$ we have derived a system of differential equations
$`_t𝐏`$ $`=`$ $`𝐁\times 𝐏D𝐏_{}+(R_{\nu _\alpha }R_{\nu _\beta })\widehat{𝐳},`$ (3.22)
$`_tP_0`$ $`=`$ $`(R_{\nu _\alpha }+R_{\nu _\beta }),`$ (3.23)
$`_t\overline{𝐏}`$ $`=`$ $`\overline{𝐁}\times \overline{𝐏}D\overline{𝐏}_{}+(\overline{R}_{\nu _\alpha }\overline{R}_{\nu _\beta })\widehat{𝐳},`$ (3.24)
$`_t\overline{P}_0`$ $`=`$ $`(\overline{R}_{\nu _\alpha }+\overline{R}_{\nu _\beta }),`$ (3.25)
where $`\widehat{𝐳}=(0,0,1)`$. In the next chapter we will study the highly non-trivial behavior of this equation of motion.
## Chapter 4 Oscillating Lepton Asymmetry
We analyze the system of flavor oscillations between active and sterile neutrinos before BBN in a simplified model. First we discuss our approximations and present a typical numerical solution of the simplified model. Next, we introduce a more convenient coordinate system for the analytical treatment. With its help, we describe the evolution of our system from very early times up to the resonance. At resonance, we show that we can describe the evolution of the neutrino asymmetry with a simple differential equation, and that the solution will indeed oscillate for some of the parameter space.
### 4.1 Simplified Model
We describe the oscillations between two flavors of neutrinos, one active and one sterile, in an expanding medium. This means that we neglect all other neutrino mixings, which is a good approximation if all other effective mixing angles are small. Thus, we can describe our oscillations by the system of differential equations (3.22)–(3.25) derived in the previous chapter. We restrict ourselves to a simplified model which we present in this section. For definiteness, we will analyze the case where the active neutrino is $`\nu _\tau `$; the analysis is analogous for $`\nu _\mathrm{e}`$ and $`\nu _\mu `$.
Our most important simplification is that we neglect the momentum distribution. This is surely not a good approximation, since then all neutrinos encounter the oscillation resonance at the same time. In reality only a small fraction of the neutrinos will be simultaneously at resonance, especially when the resonance width is small, i.e. the vacuum mixing angle is small. However, the complete system is very complicated, as can be seen from the controversial literature on its solution, and we therefore have decided to analyze this simplified model. Of course, it is desirable to include the effects of the momentum distribution in future investigations.
So from now on, all neutrinos are taken to have the same momentum. We choose this momentum to be the average of a fermion species with vanishing chemical potential,
$$p=\frac{7\zeta _4}{2\zeta _3}T3.15T,$$
(4.1)
where $`\zeta `$ is the Riemann zeta function with $`\zeta _31.202`$ and $`\zeta _4=\pi ^4/90`$. As a consequence, we only have two density matrices, $`\rho _{\alpha \alpha }`$ for the neutrinos and $`\overline{\rho }_{\alpha \alpha }`$ for the antineutrinos. We normalize these in analogy to our previous normalization in Section 3.1, i.e. such that the number densities are $`n_{\nu _\tau }=n_\nu ^{\mathrm{eq}}\rho _{\tau \tau }`$, etc. where $`n_\nu ^{\mathrm{eq}}=\frac{3}{8}n_\gamma `$ is the equilibrium neutrino number density for vanishing chemical potential.
We will also neglect the repopulation terms $`R_{\nu _\tau }`$ and $`\overline{R}_{\nu _\tau }`$ defined in (3.21). They will be small if the $`\nu _\tau `$ and $`\overline{\nu }_\tau `$ are not depopulated significantly by neutrino oscillations, and if $`A_{\nu _\tau }1`$ so that the chemical potential has no significant influence on the equilibrium number density. This simplification is valid for some part of the parameter space, which we will determine later. As a result of this approximation, $`P_0`$ and $`\overline{P}_0`$ are constant in time.
So now, the system of differential equations has been simplified to
$`_t𝐏`$ $`=`$ $`𝐁\times 𝐏D𝐏_{},`$
$`_t\overline{𝐏}`$ $`=`$ $`\overline{𝐁}\times \overline{𝐏}D\overline{𝐏}_{}.`$ (4.2)
The tau neutrino asymmetry
$$A_{\nu _\tau }=\frac{n_{\nu _\tau }^{\mathrm{eq}}}{n_\gamma }(\rho _{\tau \tau }\overline{\rho }_{\tau \tau })=\frac{3}{8}\times \frac{1}{2}(P_z\overline{P}_z+P_0\overline{P}_0)$$
(4.3)
is an important variable since it depends on $`𝐏`$, enters into $`𝐁`$, and thus causes the system to be nonlinear. It will be more convenient to use the equivalent variable
$$\delta P_z\frac{4}{3}\stackrel{~}{A}_{\nu _\tau }=\frac{1}{2}(P_z\overline{P}_z)+P_\mathrm{c},$$
(4.4)
where $`\stackrel{~}{A}_{\nu _\tau }`$ was given in (3.10) and $`P_\mathrm{c}\frac{4}{3}(A_{\nu _\mathrm{e}}+A_{\nu _\mu }\frac{1}{2}A_\mathrm{n})+\frac{1}{2}(P_0\overline{P}_0)`$ is constant.
In terms of this variable, the coefficients in the differential equations are given by
$`𝐁(T,\delta P_z)`$ $`=`$ $`\left(\begin{array}{c}b_s\\ 0\\ b_c+b_Tb_A\delta P_z\end{array}\right),`$ (4.8)
$`\overline{𝐁}(T,\delta P_z)`$ $`=`$ $`\left(\begin{array}{c}b_s\\ 0\\ b_c+b_T+b_A\delta P_z\end{array}\right),`$ (4.12)
where we use
$`b_s`$ $``$ $`{\displaystyle \frac{s\delta m^2}{2p}}{\displaystyle \frac{s\delta m^2}{6.3}}T^1,`$
$`b_c`$ $``$ $`{\displaystyle \frac{c\delta m^2}{2p}}{\displaystyle \frac{c\delta m^2}{6.3}}T^1,`$
$`b_A`$ $``$ $`V_{\nu _\tau }^A/\delta P_z=k_AT^3,`$
$`b_T`$ $``$ $`V_{\nu _\tau }^T=k_TT^5.`$ (4.13)
Recall that $`s=\mathrm{sin}2\theta `$, $`c=\mathrm{cos}2\theta `$, and $`\delta m^2=m_{\nu _\mathrm{s}}^2m_{\nu _\tau }^2`$. The constants are
$`k_A`$ $`=`$ $`{\displaystyle \frac{3\zeta _3}{4\pi ^2}}\sqrt{2}G_\mathrm{F}1.51\times 10^{24}\mathrm{eV}^2,`$
$`k_T`$ $`=`$ $`{\displaystyle \frac{14\pi }{45\alpha }}\mathrm{sin}^2\theta _\mathrm{W}\mathrm{cos}^2\theta _\mathrm{W}{\displaystyle \frac{p}{T}}G_\mathrm{F}^21.02\times 10^{44}\mathrm{eV}^4,`$ (4.14)
where $`\alpha =1/137`$ is the fine-structure constant and $`\theta _\mathrm{W}`$ the Weinberg angle. The coefficient $`D`$ is proportional to $`b_T`$, so we will often use the relation
$$D=k_Db_T,$$
(4.15)
where
$$k_D=\frac{\pi ^2y_{\nu _\tau }}{2\zeta _3z_{\nu _\tau }}1/60,$$
(4.16)
and we have used $`y_{\nu _\tau }0.79`$ and $`z_{\nu _\tau }97`$ from Section 3.2. For $`\nu _\mu `$, the constants are the same. For $`\nu _\mathrm{e}`$, we have $`k_T3.66\times 10^{44}\mathrm{eV}^4`$ and $`k_D1/151`$.
We will be considering neutrino oscillations with small vacuum mixing angles, i.e. $`sc1`$. As a consequence, the first component of $`𝐁`$ will be much smaller than the third component for most of the time. The system becomes interesting when the third component disappears. Then the neutrino oscillations are at resonance, i.e. the neutrinos mix maximally. If we assume the initial asymmetry to be negligible, the resonance condition is given by $`b_c=b_T`$. This condition has a solution only if $`\delta m^2<0`$. Then resonance occurs when
$$T_{\mathrm{res}}=\left(\frac{|\delta m^2|c}{6.3k_T}\right)^{1/6}15.8\mathrm{MeV}|\delta m_{\mathrm{eV}}^2|^{1/6},$$
(4.17)
where $`\delta m_{\mathrm{eV}}^2=\delta m^2/\mathrm{eV}^2`$. The resonance will be the crucial feature of the system. We will see that shortly after the resonance, the neutrino oscillations will create a large asymmetry, an effect which is driven by the non-linear term $`b_A\delta P_z`$ in the system of differential equations. We will therefore only consider $`\delta m^2<0`$, i.e. the sterile neutrino is lighter than the tau neutrino.
### 4.2 Numerical Solution
We have solved the system (4.2) of differential equations numerically and find results similar to those in and . We have plotted the evolution of the effective asymmetry $`\delta P_z`$ for $`(\delta m^2,s)=(1\mathrm{eV}^2,10^4)`$ in Fig. 4.2. Its behavior is representative for a large region of parameter space. The evolution of the system falls into five distinct phases.
* At high temperatures, the asymmetry $`\delta P_z`$ is constant, since damping totally destroys the coherence of the oscillations. Thus, no flavor transition occurs at all.
* As the temperature decreases, the mixing angle increases, while the damping rate decreases. As a consequence, there will be a small amount of flavor transition induced by the damping. Due to the asymmetric term $`b_A\delta P_z`$, this transition rate will be different for neutrinos and antineutrinos, washing out the asymmetry $`\delta P_z`$. However, there remains a small relic asymmetry.
* Until the system reaches the resonance, the relic asymmetry changes slowly.
* When the system has passed the resonance, the difference in the transition rates for neutrinos and antineutrinos has the opposite effect than before. In other words, a small asymmetry is no longer washed out, but amplified, which leads to an exponential increase of $`\delta P_z`$.
* When $`\delta P_z=𝒪(10^6)`$, it starts oscillating and thereby even changes its sign. For our numerical example, we have plotted the evolution of $`\delta P_z`$ during this step in Fig. 4.2.
* Eventually, the neutrino asymmetry stops oscillating. After that, it slowly increases.
In the past, the underlying mechanism for the oscillatory behavior in step 5 was not fully understood. Merely the authors of discussed this phenomenon qualitatively. Our main achievement will be to analyze the system analytically and to prove that it indeed oscillates. Thus, we invalidate the frequently used argument that the oscillations in this phase are an artifact caused by a numerical instability.
### 4.3 Spherical Coordinates
We will now change to a more convenient coordinate system. The two phase-related coordinates $`P_x`$ and $`P_y`$ are not very practical, so we prefer to use spherical coordinates. As a measure for the length $`P=|𝐏|`$ of the polarization vector we use $`P_z`$ because it enters directly into (4.4). This choice becomes problematic when $`P_z=0,P0`$, but as mentioned above, $`P_z`$ will not change significantly. The new coordinate system consists of $`P_z`$ and the angles $`\vartheta `$ and $`\phi `$. These new variables are related to the old ones as
$`P_x`$ $`=`$ $`P_z\mathrm{tan}\vartheta \mathrm{sin}\phi ,`$
$`P_y`$ $`=`$ $`P_z\mathrm{tan}\vartheta \mathrm{cos}\phi ,`$
$`P_z`$ $`=`$ $`P_z.`$ (4.18)
$`P_z`$ represents the relation between the number densities of the active and sterile neutrinos, $`\vartheta `$ is a measure of how much the polarization vector deviates from the $`z`$-axis, and $`\phi `$ is the rotation angle around the $`z`$-axis. Then the differential equations (4.2) change to
$`\dot{P}_z`$ $`=`$ $`b_sP_z\mathrm{tan}\vartheta \mathrm{cos}\phi ,`$
$`\dot{\vartheta }`$ $`=`$ $`D\mathrm{sin}\vartheta \mathrm{cos}\vartheta +b_s\mathrm{cos}\phi ,`$
$`\dot{\phi }`$ $`=`$ $`b_cb_T+b_A\delta P_zb_s\mathrm{cot}\vartheta \mathrm{sin}\phi ,`$ (4.19)
and
$`\dot{\overline{P}_z}`$ $`=`$ $`b_s\overline{P}_z\mathrm{tan}\overline{\vartheta }\mathrm{cos}\overline{\phi },`$
$`\dot{\overline{\vartheta }}`$ $`=`$ $`D\mathrm{sin}\overline{\vartheta }\mathrm{cos}\overline{\vartheta }+b_s\mathrm{cos}\overline{\phi },`$
$`\dot{\overline{\phi }}`$ $`=`$ $`b_cb_Tb_A\delta P_zb_s\mathrm{cot}\overline{\vartheta }\mathrm{sin}\overline{\phi },`$ (4.20)
where over-barred quantities, as usual, refer to antineutrinos. With (4.4) we have a closed set of differential equations.
The only difference between the equations for neutrinos and antineutrinos is the sign of the term $`b_A\delta P_z`$. Therefore, if this term is negligible, the variables for neutrinos and antineutrinos develop equally. To be more sensitive to the small differences induced by the asymmetry, we perform another substitution and use the variables
$`P_z^\pm `$ $``$ $`\frac{1}{2}(P_z\pm \overline{P}_z),`$
$`\vartheta ^\pm `$ $``$ $`\frac{1}{2}(\vartheta \pm \overline{\vartheta }),`$
$`\phi ^\pm `$ $``$ $`\frac{1}{2}(\phi \pm \overline{\phi })`$ (4.21)
instead. Thus $`P_z^+P_z^{}`$, $`\vartheta ^+\vartheta ^{}`$ and $`\phi ^+\phi ^{}`$ when the fermion asymmetry is small. Note that (4.4) becomes
$$\delta P_z(t)=P_z^{}(t)+P_\mathrm{c}$$
(4.22)
so that $`\frac{d}{dt}\delta P_z=\dot{P}_z^{}`$. We will make use of this simple relation by considering $`\delta P_z`$ instead of $`P_z^{}`$ when it is convenient.
In terms of the new variables, our system of differential equations is
$`\dot{P}_z^+`$ $`=`$ $`b_sP_z^+{\displaystyle \frac{\mathrm{sin}2\vartheta ^+\mathrm{cos}\phi ^+\mathrm{cos}\phi ^{}\mathrm{sin}2\vartheta ^{}\mathrm{sin}\phi ^+\mathrm{sin}\phi ^{}}{\mathrm{cos}2\vartheta ^++\mathrm{cos}2\vartheta ^{}}}`$ (4.23)
$`b_sP_z^{}{\displaystyle \frac{\mathrm{sin}2\vartheta ^+\mathrm{sin}\phi ^+\mathrm{sin}\phi ^{}+\mathrm{sin}2\vartheta ^{}\mathrm{cos}\phi ^+\mathrm{cos}\phi ^{}}{\mathrm{cos}2\vartheta ^++\mathrm{cos}2\vartheta ^{}}},`$
$`\dot{P}_z^{}`$ $`=`$ $`b_sP_z^+{\displaystyle \frac{\mathrm{sin}2\vartheta ^+\mathrm{sin}\phi ^+\mathrm{sin}\phi ^{}+\mathrm{sin}2\vartheta ^{}\mathrm{cos}\phi ^+\mathrm{cos}\phi ^{}}{\mathrm{cos}2\vartheta ^++\mathrm{cos}2\vartheta ^{}}}`$ (4.24)
$`b_sP_z^{}{\displaystyle \frac{\mathrm{sin}2\vartheta ^+\mathrm{cos}\phi ^+\mathrm{cos}\phi ^{}\mathrm{sin}2\vartheta ^{}\mathrm{sin}\phi ^+\mathrm{sin}\phi ^{}}{\mathrm{cos}2\vartheta ^++\mathrm{cos}2\vartheta ^{}}},`$
$`\dot{\vartheta }^+`$ $`=`$ $`{\displaystyle \frac{1}{2}}D\mathrm{sin}2\vartheta ^+\mathrm{cos}2\vartheta ^{}+b_s\mathrm{cos}\phi ^+\mathrm{cos}\phi ^{},`$ (4.25)
$`\dot{\vartheta }^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}D\mathrm{cos}2\vartheta ^+\mathrm{sin}2\vartheta ^{}b_s\mathrm{sin}\phi ^+\mathrm{sin}\phi ^{},`$ (4.26)
$`\dot{\phi }^+`$ $`=`$ $`b_cb_T`$ (4.27)
$`+b_s{\displaystyle \frac{\mathrm{sin}2\vartheta ^+\mathrm{sin}\phi ^+\mathrm{cos}\phi ^{}+\mathrm{sin}2\vartheta ^{}\mathrm{cos}\phi ^+\mathrm{sin}\phi ^{}}{\mathrm{cos}2\vartheta ^{}\mathrm{cos}2\vartheta ^+}},`$
$`\dot{\phi }^{}`$ $`=`$ $`b_A(P_z^{}+P_\mathrm{c})`$
$`+b_s{\displaystyle \frac{\mathrm{sin}2\vartheta ^+\mathrm{cos}\phi ^+\mathrm{sin}\phi ^{}+\mathrm{sin}2\vartheta ^{}\mathrm{sin}\phi ^+\mathrm{cos}\phi ^{}}{\mathrm{cos}2\vartheta ^{}\mathrm{cos}2\vartheta ^+}}.`$
As a last approximation, we assume $`\vartheta ^\pm 1`$, which effectively means that the polarization vectors will stay close to the $`z`$-axis. We expand $`\mathrm{sin}2\vartheta ^\pm `$ and $`\mathrm{cos}2\vartheta ^\pm `$ and only take into account the leading order terms. Then we get
$`\dot{P}_z^+`$ $`=`$ $`b_sP_z^+\left(\vartheta ^+\mathrm{cos}\phi ^+\mathrm{cos}\phi ^{}\vartheta ^{}\mathrm{sin}\phi ^+\mathrm{sin}\phi ^{}\right)`$ (4.29)
$`b_sP_z^{}\left(\vartheta ^+\mathrm{sin}\phi ^+\mathrm{sin}\phi ^{}+\vartheta ^{}\mathrm{cos}\phi ^+\mathrm{cos}\phi ^{}\right),`$
$`\dot{P}_z^{}`$ $`=`$ $`b_sP_z^+\left(\vartheta ^+\mathrm{sin}\phi ^+\mathrm{sin}\phi ^{}+\vartheta ^{}\mathrm{cos}\phi ^+\mathrm{cos}\phi ^{}\right)`$ (4.30)
$`b_sP_z^{}\left(\vartheta ^+\mathrm{cos}\phi ^+\mathrm{cos}\phi ^{}\vartheta ^{}\mathrm{sin}\phi ^+\mathrm{sin}\phi ^{}\right),`$
$`\dot{\vartheta }^+`$ $`=`$ $`D\vartheta ^++b_s\mathrm{cos}\phi ^+\mathrm{cos}\phi ^{},`$ (4.31)
$`\dot{\vartheta }^{}`$ $`=`$ $`D\vartheta ^{}b_s\mathrm{sin}\phi ^+\mathrm{sin}\phi ^{},`$ (4.32)
$`\dot{\phi }^+`$ $`=`$ $`b_cb_T`$ (4.33)
$`+{\displaystyle \frac{b_s}{\vartheta ^+}}\left(\mathrm{sin}\phi ^+\mathrm{cos}\phi ^{}+{\displaystyle \frac{\vartheta ^{}}{\vartheta ^+}}\mathrm{cos}\phi ^+\mathrm{sin}\phi ^{}\right),`$
$`\dot{\phi }^{}`$ $`=`$ $`b_A(P_z^{}+P_\mathrm{c})`$ (4.34)
$`+{\displaystyle \frac{b_s}{\vartheta ^+}}\left(\mathrm{cos}\phi ^+\mathrm{sin}\phi ^{}+{\displaystyle \frac{\vartheta ^{}}{\vartheta ^+}}\mathrm{sin}\phi ^+\mathrm{cos}\phi ^{}\right).`$
From (4.31) we see that $`\dot{\vartheta }^+D\vartheta ^++b_s`$. Thus $`\vartheta ^+b_s/D`$ would automatically mean that $`\dot{\vartheta }^+0`$. Therefore, we can safely say that $`\vartheta ^+<b_s/D`$ at all times. We will see that $`\vartheta ^+`$ becomes maximal at resonance, so $`\vartheta _{\mathrm{max}}^+(b_s/D)_{\mathrm{res}}=s/(ck_D)`$, where we have used $`b_s/s=b_c/c`$, $`(b_c)_{\mathrm{res}}=(b_T)_{\mathrm{res}}`$ and $`D=k_Db_T`$. Thus our approximation is only valid if $`sk_D=1/60`$, $`c1`$, i.e. for small mixing angles.
### 4.4 Initial Conditions
We need to derive the initial conditions before we can consider the evolution of the system of differential equations. We begin with temperatures far above the resonance. For $`T\mathrm{}`$, the coefficients $`D,b_T,R_{\nu _\tau }T^5\mathrm{}`$, $`b_AT^3\mathrm{}`$, while $`b_s,b_cT^10`$. If we compare these coefficients with $`HT^2`$, we conclude that $`b_s`$ and $`b_c`$ can be neglected at very high temperatures. Thus, from (4.29)–(4.32) we see that regardless of the initial conditions, $`\vartheta ^+`$ and $`\vartheta ^{}`$ are both exponentially damped to zero, while $`P_z^\pm b_s`$ do not change. However, the scattering processes equilibrate $`\nu _\tau `$ and $`\overline{\nu }_\tau `$, implying $`\rho _{\tau \tau }`$ and $`\overline{\rho }_{\tau \tau }1`$. Since the sterile neutrinos do not interact, $`\rho _{\mathrm{ss}}`$ and $`\overline{\rho }_{\mathrm{ss}}`$ remain constant. They are only diluted whenever massive particles become non-relativistic and annihilate into the still relativistic particles, which heats up the plasma. For example, when the temperature decreases from $`T\text{TeV}`$ (just before the electro-weak symmetry breaking) to $`T\text{MeV}`$ (shortly after the quark-hadron phase transition), $`\rho _{\mathrm{ss}}`$ and $`\overline{\rho }_{\mathrm{ss}}`$ are diluted by a factor of $`[g_{}(T300\mathrm{MeV})/g_{}(T\text{TeV})]^{4/3}0.07`$. So if the $`\nu _\mathrm{s}`$ were in equilibrium with the thermal plasma at very early times by some unknown mechanism, they would have been strongly diluted by the time they become interesting for us. We will simply assume that their initial density is zero. Then the initial conditions are $`P_0=\overline{P}_0=1`$ and $`P_z=\overline{P}_z=1`$. Including the small initial neutrino asymmetries, $`A_{\nu _\tau }`$ and $`A_{\nu _\mathrm{s}}`$, we find
$`P_z^+`$ $`=`$ $`1,`$
$`P_z^{}`$ $`=`$ $`{\displaystyle \frac{4}{3}}(A_{\nu _\tau }A_{\nu _\mathrm{s}}),`$
$`{\displaystyle \frac{1}{2}}(P_0\overline{P}_0)`$ $`=`$ $`{\displaystyle \frac{4}{3}}(A_{\nu _\tau }+A_{\nu _\mathrm{s}})=\mathrm{const},`$
$`P_\mathrm{c}`$ $`=`$ $`{\displaystyle \frac{4}{3}}(A_\mathrm{c}+A_{\nu _\tau }+A_{\nu _\mathrm{s}})=\mathrm{const},`$ (4.35)
where in $`P_z^+`$ we have neglected $`A_{\nu _\tau }`$ and $`A_{\nu _\mathrm{s}}1`$.
### 4.5 Quasi-Static Solutions
We will now follow the time evolution of the system. We first stress the important role that damping plays in our system. We consider a variable $`x`$, which follows the differential equation
$$\frac{}{t}x(t)=d(t)x(t)+f(t),$$
(4.36)
where $`d(t)`$ is some damping coefficient and $`f(t)`$ is a function. If $`d`$ and $`f`$ are constant, $`x`$ will relax to a static value $`x_{\mathrm{st}}=\frac{f}{d}`$. In our case, $`d`$ and $`f`$ will slowly vary in time, and so will the static value. But if its rate of change $`r_x\dot{x}_{\mathrm{st}}/x_{\mathrm{st}}`$ is smaller than the damping coefficient $`d`$, the damping will force the variable to follow its static value. We will make use of this approximation; since it is not static in the strict sense of the meaning, we will call these solutions quasi-static.
We have calculated the rates of change of the quasi-static solutions for our variables in Appendix A. For the discussion in this section it is sufficient to know that $`r=𝒪(H)`$.
We can now determine the variables for large temperatures. We can easily assume that $`b_s\vartheta ^\pm D\vartheta ^\pm ,b_s,(b_cb_T),b_s/\vartheta ^\pm `$, so before $`P_z^\pm `$ change significantly at all, the other four variables will have relaxed to their quasi-static values. If we assume that $`\dot{P}_z^{}=0`$, the differential equations for neutrinos and antineutrinos (4.2) decouple. For small $`\vartheta `$, the angular equations in (4.19) become
$`\dot{\vartheta }`$ $`=`$ $`D\vartheta +b_s\mathrm{cos}\phi `$ (4.37)
$`\dot{\phi }`$ $`=`$ $`(b_cb_T+b_A\delta P_z){\displaystyle \frac{b_s}{\vartheta }}\mathrm{sin}\phi .`$ (4.38)
Since $`\vartheta >0`$ by definition, the first term in the first equation and the second term in the second equation are both damping terms. The two variables relax to $`\dot{\vartheta }\dot{\phi }0`$ if the quasi-static conditions are fulfilled, i.e. $`D>H`$ and $`b_s/\vartheta >H`$, respectively. The first condition is fulfilled for $`T>\mathrm{few}\mathrm{MeV}`$. When $`\vartheta `$ has relaxed, we get from (4.37) that $`\vartheta b_s\mathrm{cos}\phi /D<b_s/D`$, which fulfills the second condition, $`b_s/\vartheta >D>H`$. Setting $`\dot{\vartheta }=\dot{\phi }=0`$, we get the quasi-static solutions
$`\vartheta _{\mathrm{qs}}`$ $`=`$ $`{\displaystyle \frac{b_s}{\sqrt{D^2+(b_cb_T+b_A\delta P_z)^2}}},`$ (4.39)
$`\phi _{\mathrm{qs}}`$ $`=`$ $`\mathrm{arctan}\left({\displaystyle \frac{b_cb_T+b_A\delta P_z}{D}}\right).`$ (4.40)
The same is valid for the antineutrinos, but with $`\delta P_z\delta P_z`$. Since $`b_A\delta P_zb_cb_T`$, we can already say that the variables $`\vartheta ^{}`$ and $`\phi ^{}`$ will be much smaller than $`\vartheta ^+`$ and $`\phi ^+`$, respectively. Thus,
$`\vartheta _{\mathrm{qs}}^+`$ $`=`$ $`\frac{1}{2}(\vartheta _{\mathrm{qs}}+\overline{\vartheta }_{\mathrm{qs}}){\displaystyle \frac{b_s}{\sqrt{D^2+(b_cb_T)^2}}},`$ (4.41)
$`\phi _{\mathrm{qs}}^+`$ $`=`$ $`\frac{1}{2}(\phi _{\mathrm{qs}}+\overline{\phi }_{\mathrm{qs}})\mathrm{arctan}\left({\displaystyle \frac{b_cb_T}{D}}\right).`$ (4.42)
The same quasi-static solutions are also derived from (4.31) and (4.33) if we neglect the last term in latter.
Now it is also easy to calculate $`\vartheta ^{}`$ and $`\phi ^{}`$ for a given $`\delta P_z`$ by setting $`\dot{\vartheta }^{}=0`$ in (4.32) and $`\dot{\phi }^{}=0`$ in (4.34), inserting $`\vartheta _{\mathrm{qs}}^+`$ and $`\phi _{\mathrm{qs}}^+`$ and setting $`\mathrm{cos}\phi ^{}=1`$. Then
$`\phi _{\mathrm{qs}}^{}`$ $`=`$ $`\mathrm{arcsin}\left({\displaystyle \frac{Db_A}{D^2+(b_cb_T)^2}}\delta P_z\right)`$ (4.43)
$`\vartheta _{\mathrm{qs}}^{}`$ $`=`$ $`{\displaystyle \frac{b_s(b_cb_T)b_A}{[D^2+(b_cb_T)^2]^{3/2}}}\delta P_z.`$ (4.44)
Again, we need to show that the quasi-static conditions are fulfilled, which we again find to be $`D>H`$.
If we insert all these variables into $`\dot{P}_z^{}`$, and neglect the fourth term in (4.30), we obtain
$$\frac{d}{dt}\delta P_z=\dot{P}_z^{}=\kappa (\delta P_z+\epsilon P_z^{}),$$
(4.45)
where the damping coefficient is
$$\kappa \left|\frac{2P_z^+b_s^2(b_cb_T)Db_A}{[D^2+(b_cb_T)^2]^2}\right|T^9,$$
(4.46)
and
$$\epsilon \left|\frac{D^2+(b_cb_T)^2}{2P_z^+b_A(b_cb_T)}\right|,$$
(4.47)
where $`\epsilon 1`$ for $`T<\mathrm{GeV}`$ and $`TT_{\mathrm{res}}`$. Analogously, we get
$$\dot{P}_z^+=\kappa \epsilon P_z^+T^7,$$
(4.48)
where we have only taken into account the first term in (4.29), since the other terms are all $`𝒪(P_z^{})`$ smaller. This confirms that $`P_z^{}`$ and $`P_z^+`$ will not relax to their quasi-static values at high temperatures.
### 4.6 Evolution Towards the Resonance
When $`\kappa \mathrm{\Gamma }>H`$, $`\delta P_z`$ starts changing. The damping coefficients for $`\phi ^{}`$ and $`\vartheta ^{}`$ (which are $`\mathrm{\Gamma }>D`$) are much larger than H, so even though $`\delta P_z`$ changes, (4.43) and (4.44) remain valid.
From (4.45) we see that $`\delta P_z`$ is damped towards 0, i.e. $`P_z^{}P_c`$. When $`|\delta P_z||P_z^{}|`$, the second term in (4.45) becomes important. We set $`\dot{P}_z^{}=0`$ to get
$$(P_z^{})_{\mathrm{qs}}=P_\mathrm{c}\frac{1}{1+\epsilon },(\delta P_z)_{\mathrm{qs}}=P_\mathrm{c}\frac{\epsilon }{1+\epsilon }.$$
(4.49)
We have to prove that $`\delta P_z`$ takes on its quasi-static value before the system passes the resonance. The solution of (4.45) is given by
$$\delta P_z(t)=\delta P_z(0)\mathrm{exp}\left(\underset{0}{\overset{t}{}}\kappa (t)𝑑t\right).$$
(4.50)
If we substitute $`x=T(t)/T_{\mathrm{res}}`$, we can write the integral as
$$F_{}(t)\underset{0}{\overset{t}{}}\kappa 𝑑t=\frac{k_Dm_{\mathrm{pl}}2P_z^+k_A}{6.3^{1/6}5.5k_T^{1/6}}\frac{s^2|\delta m^2|^{1/6}}{c^{11/6}}I_{}(t),$$
(4.51)
where
$$I_{}(t_{\mathrm{res}})=\underset{1}{\overset{\mathrm{}}{}}\frac{x^6(x^61)}{[k_D^2x^{12}+(x^61)^2]^2}𝑑x293.$$
(4.52)
Assuming that initially $`\delta P_z`$ and $`P_z^{}`$ are of the same order of magnitude, we know that $`\delta P_z`$ will decrease by a factor of order $`\epsilon b_T/b_A𝒪(10^6)`$, which corresponds to $`F_{}(t_{\mathrm{res}})6\mathrm{ln}(10)`$. Thus, we obtain the condition
$$s^2|\delta m^2|^{1/6}\mathrm{\Gamma }>\mathrm{\hspace{0.17em}10}^{12}\mathrm{eV}^{1/3}$$
(4.53)
using $`P_z^+1`$.
For $`P_z^+`$ we can perform a similar calculation. However, here we demand that $`P_z^+`$ does not change significantly to prevent that the repopulation terms become important. Therefore we demand that $`F_+(t_{\mathrm{res}})1`$, where
$$P_z^+(t)=P_z^+(0)\mathrm{exp}\left[F_+(t)\right].$$
(4.54)
Thus,
$$F_+(t)=\underset{0}{\overset{t}{}}\kappa (t)\epsilon (t)𝑑t=\frac{k_Dm_{\mathrm{pl}}k_T^{1/2}}{6.3^{1/2}5.5}\frac{s^2|\delta m^2|^{1/2}}{c^{3/2}}I_+(t),$$
(4.55)
where
$$I_+(t_{\mathrm{res}})=\underset{1}{\overset{\mathrm{}}{}}\frac{x^2}{k_D^2x^{12}+(x^61)^2}𝑑x15.2.$$
(4.56)
We then get the condition
$$s^4|\delta m_{\mathrm{eV}}^2|2\times 10^9,$$
(4.57)
where $`\delta m_{\mathrm{eV}}^2=\delta m^2/\mathrm{eV}^2`$. $`F_+(t_{\mathrm{res}})1`$ is also the condition that the $`\nu _\mathrm{s}`$ do not come into equilibrium. Our result here is in good agreement with previous constraints .
The quasi-static approximation is of course not exact, actually the variables will always be a bit behind their quasi-static values. This delay becomes important when the quasi-static value of a variable passes zero, or when the rate of change of the quasi-static value becomes of order of the damping rate of the variable. At resonance, i.e. when $`(b_cb_T)=0`$, we have $`\vartheta _{\mathrm{qs}}^{}=0`$, $`\phi _{\mathrm{qs}}^+=0`$, $`(P_z^{})_{\mathrm{qs}}=0`$. Furthermore, all the quasi-static values change relatively fast near the resonance.
Here is a qualitative discussion of the behavior of the variables close to resonance. The term that is responsible for the resonance, $`\mathrm{\Delta }bb_cb_T`$, only enters into the differential equation of $`\phi ^+`$. When $`\mathrm{\Delta }b`$ changes sign, $`\phi _{\mathrm{qs}}^+`$ does too. Therefore, after a short delay, $`\phi ^+`$ will also change its sign. The other variables do not depend directly on $`\mathrm{\Delta }b`$; they only feel its change via $`\phi ^+`$. Thus, they are not sensitive to the time delay between $`\phi ^+`$ and $`\phi _{\mathrm{qs}}^+`$. Therefore, setting $`\phi ^+=\phi _{\mathrm{qs}}^+`$ will merely shift the events by a negligible amount of time.
Close to resonance, we have $`\mathrm{sin}\phi ^+0`$ and $`\mathrm{cos}\phi ^+1`$. Therefore, $`\vartheta ^+`$ becomes maximal close to the resonance, i.e. $`\vartheta ^+b_s/D`$, while $`\vartheta ^{}`$ becomes minimal so that we can neglect it in the equations. $`\phi ^{}`$ will stay close to its quasi-static value since $`\delta P_z`$ does not change significantly: $`\delta P_z`$ freezes out when the rate of change in $`(\delta P_z)_{\mathrm{qs}}`$ becomes larger than the damping rate $`\kappa `$, i.e. $`r_{\delta P_z}>\kappa `$, where $`r_{\delta P_z}`$ is given in (A.7). Since $`\mathrm{\Delta }bb_cb_TD`$ and $`\epsilon 1`$ at freeze-out, we get
$$r_{\delta P_z}6H\frac{b_c}{\mathrm{\Delta }b},\kappa \frac{2P_z^+b_s^2b_A\mathrm{\Delta }b}{D^3}.$$
(4.58)
Using $`T=T_{\mathrm{res}}`$, we find that $`\mathrm{\Delta }b_{\mathrm{fr}}=8.1\times 10^9s^1|\delta m_{\mathrm{eV}}^2|^{1/12}c^{11/12}b_T`$, and thus $`\epsilon _{\mathrm{fr}}=7.2\times 10^3s|\delta m_{\mathrm{eV}}^2|^{5/12}c^{7/12}1`$, so that $`(\delta P_z)_{\mathrm{fr}}\epsilon _{\mathrm{fr}}P_\mathrm{c}`$ is indeed very small.
In summary, at resonance we can still approximate the variables $`\vartheta ^+`$, $`\phi ^+`$ and $`\phi ^{}`$ by their static values, while we can neglect $`\vartheta ^{}`$. $`\delta P_z`$ is given by its freeze-out value.
### 4.7 Evolution at Resonance
We now come to the most interesting feature of the system: $`\mathrm{sin}\phi ^+`$ changes its sign and becomes positive. Therefore, the back-reaction of $`\phi ^{}`$ on $`\delta P_z`$, which is given by the coefficients of the first terms in (4.30) and (4.34), $`b_Ab_sP_z^+\vartheta ^+\mathrm{sin}\phi ^+`$, also changes its sign to positive. So while $`\delta P_z`$ until now was washed out through this back-reaction, the same back-reaction now amplifies $`\delta P_z`$. We will now derive the consequences of this effect. To this end, we will assume that the variables $`\vartheta ^\pm `$ and $`\phi ^\pm `$ will be given approximately by their static values. As time passes, these approximations will break down and we will need to find other approximations.
#### 4.7.1 Static Approximation after Resonance
The change of $`\delta P_z`$ becomes interesting again when $`r_{\delta P_z}\kappa `$ after the resonance, i.e. when $`\mathrm{\Delta }b_0=\mathrm{\Delta }b_{\mathrm{fr}}`$. We will denote the time when this happens with $`t_0`$. So until then, $`\delta P_z(t_0)(\delta P_z)_{\mathrm{fr}}=\epsilon _{\mathrm{fr}}P_\mathrm{c}`$. Now (4.45) has changed to
$$\frac{d}{dt}\delta P_z=+\kappa (\delta P_z\epsilon P_z^{})+\kappa (\delta P_z+\epsilon P_\mathrm{c}).$$
(4.59)
If we neglect the second term, we get the solution
$$\delta P_z(t)=\delta P_z(t_0)\mathrm{exp}\left(\underset{t_0}{\overset{t}{}}\kappa 𝑑t\right).$$
(4.60)
To solve the integral, we substitute $`t`$ with $`\xi =1x1`$, where $`x=T(t)/T_{\mathrm{res}}`$. Then
$$\mathrm{\Delta }b=b_cb_T=b_c^{\mathrm{res}}x^1b_T^{\mathrm{res}}x^5=b_T^{\mathrm{res}}x^1(1x^6)6b_T^{\mathrm{res}}\xi ,$$
(4.61)
and
$$d\xi =\frac{1}{T_{\mathrm{res}}}dT=\frac{5.5T_{\mathrm{res}}^2x^3}{m_{\mathrm{pl}}}dt\frac{5.5T_{\mathrm{res}}^2}{m_{\mathrm{pl}}}dt,$$
(4.62)
where we have expanded $`\xi `$ and have only taken into account the leading order. The integral becomes
$`{\displaystyle \underset{t_0}{\overset{t}{}}}\kappa 𝑑t`$ $`=`$ $`{\displaystyle \frac{m_{\mathrm{pl}}}{5.5T_{\mathrm{res}}^2}}{\displaystyle \underset{\xi _0}{\overset{\xi }{}}}\kappa 𝑑\xi {\displaystyle \frac{m_{\mathrm{pl}}}{5.5T_{\mathrm{res}}^2}}{\displaystyle \underset{\xi _0}{\overset{\xi }{}}}{\displaystyle \frac{2P_z^+b_s^2b_A6b_T^{\mathrm{res}}\xi }{D^3}}𝑑\xi `$ (4.63)
$`=`$ $`{\displaystyle \frac{2P_z^+b_s^2b_A6b_Tm_{\mathrm{pl}}}{D^35.5T^2}}|_{T=T_{\mathrm{res}}}{\displaystyle \underset{\xi _0}{\overset{\xi }{}}}\xi 𝑑\xi `$
$`=`$ $`5.5\times 10^{17}s^2|\delta m_{\mathrm{eV}}^2|^{1/6}c^{11/6}\frac{1}{2}(\xi ^2\xi _0^2),`$
where we have assumed that $`\mathrm{\Delta }bD`$ and have again expanded $`\xi `$ to leading order. Furthermore, $`\xi _0=\mathrm{\Delta }b_0/6b_T^{\mathrm{res}}=1.35\times 10^9s^1|\delta m_{\mathrm{eV}}^2|^{1/12}c^{11/12}`$.
The solution (4.60) is only valid as long as $`\phi ^{}`$ and $`\vartheta ^{}`$ can follow their static approximations, which are given by (4.43) and (4.44). We see that as long as $`\mathrm{\Delta }b=b_cb_TD`$, the changes in $`\phi _{\mathrm{qs}}^{}`$ and $`\vartheta _{\mathrm{qs}}^{}`$ will mainly be due to $`\delta P_z`$, i.e. $`\dot{\phi }_{\mathrm{qs}}^{}/\phi _{\mathrm{qs}}^{}\dot{\vartheta }_{\mathrm{qs}}^+/\vartheta _{\mathrm{qs}}^+\frac{d}{dt}\delta P_z/\delta P_z`$. Then the validity of (4.60) breaks down when $`D=\frac{d}{dt}\delta P_z/\delta P_z=\kappa `$, which happens when $`\xi =\xi _1=1.1\times 10^{14}s^2|\delta m_{\mathrm{eV}}^2|^{1/3}c^{7/3}`$. We see that the above solution already breaks down very early, for a large parameter range we even have $`\xi _1<\xi _0`$. To give an estimate of how big the integral (4.63) will be, we set $`\xi _0=0`$ and get $`_{t_{\mathrm{res}}}^{t_1}\kappa 𝑑t=0.8\times 10^{11}s^2|\delta m_{\mathrm{eV}}^2|^{5/6}c^{17/6}`$, which means that the change in $`\delta P_z`$ between $`t_{\mathrm{res}}`$ and $`t_1`$ is negligible. So we can immediately give up the static approximations for $`\phi ^{}`$ and $`\vartheta ^{}`$.
#### 4.7.2 Mathematical Pendulum
After the quasi-static approximation for $`\phi ^{}`$ breaks down, the first term in (4.34) will be the dominant one due to the rapid growth of $`\delta P_z`$, so we neglect the other terms. In (4.30), we can neglect the terms with the slowly varying $`\vartheta ^{}`$. Since $`b_Ab_s\vartheta ^+\mathrm{sin}\phi ^+`$, $`P_z^{}`$ will have a much larger effect on $`\dot{\phi }^{}`$ than on $`d(\delta P_z)/dt`$, so that soon $`\phi ^{}P_z^{}`$, and we can thus neglect the third term in (4.30). Therefore we get the simplified equations
$`{\displaystyle \frac{d}{dt}}\delta P_z`$ $``$ $`b_s(\vartheta ^+\mathrm{sin}\phi ^+)\mathrm{sin}\phi ^{}`$ (4.64)
$`{\displaystyle \frac{d}{dt}}\phi ^{}`$ $``$ $`b_A\delta P_z.`$ (4.65)
Together, they give the second order equation
$$\ddot{\phi }^{}g\mathrm{sin}\phi ^{},$$
(4.66)
where $`g=b_Ab_s(\vartheta ^+\mathrm{sin}\phi ^+)`$. We have now found a very simple description of the system. In the next two subsections, we will discuss the features of this equation. Afterwards, we discuss the consequences for our system.
We take $`t=t_0`$ to be the initial time. Then the initial conditions are given by
$$\phi _0^{}=\phi _{\mathrm{qs}}^{}\frac{b_A}{D}(\delta P_z)_{\mathrm{fr}}1,\dot{\phi }^{}0.$$
(4.67)
For definiteness, we take $`\phi _0^{}`$ to be positive, which implies $`\delta P_z>0`$.
Let us first assume that $`g`$ is constant. Then the differential equation corresponds to a mathematical pendulum, where $`g`$ is the acceleration of gravity, see Fig. 4.3. Usually, we could apply the small angle expansion around the stable minimum of the potential energy. However, in our case $`\phi ^{}=0`$ corresponds to the meta-stable maximum of the potential energy. So our system will perform large-amplitude oscillations. We denote the amplitude of the oscillations with $`\phi _{\mathrm{max}}^{}`$, so small $`\phi _{\mathrm{max}}^{}`$ corresponds to large amplitude.
For illustration, we can define the analogy of a potential energy $`E_{\mathrm{pot}}=g(1\mathrm{cos}\phi ^{})`$ and a kinetic energy $`E_{\mathrm{kin}}=\frac{1}{2}(\dot{\phi }^{})^2`$. Conservation of energy implies that the total energy $`E_{\mathrm{tot}}=E_{\mathrm{pot}}(t_0)+E_{\mathrm{kin}}(t_0)=\mathrm{const}`$. Since $`E_{\mathrm{tot}}(t_0)<0`$, we know that the system will oscillate around the stable point $`\phi ^{}=\pi `$ with constant amplitude $`\phi _{\mathrm{max}}^{}=|\phi _0^{}|`$, and that $`\phi ^{}`$ will never pass the meta-stable point $`\phi ^{}=0`$. We can also state that due to the non-linearity of $`\mathrm{sin}\phi ^{}`$, the oscillation frequency of the system, $`\nu `$, will be smaller than the oscillation frequency $`\sqrt{g}/2\pi `$ for the linear small-angle approximation. In fact, for $`E_{\mathrm{tot}}0`$, the oscillation frequency goes to zero. In Fig. 4.4, we have plotted the ratio $`2\pi \nu /\sqrt{g}`$ as a function of the amplitude $`\phi _{\mathrm{max}}^{}`$.
Another important fact is that the time average of $`\mathrm{cos}\phi ^{}`$, $`\mathrm{cos}\phi ^{}`$, is greater than 0 for small $`\phi _{\mathrm{max}}^{}`$, i.e. for large amplitudes. This results from the fact that the system develops relatively slowly close to the turning points $`\pm \phi _{\mathrm{max}}^{}`$, where the kinetic energy is small. We have plotted $`\mathrm{cos}\phi ^{}`$ as a function of the amplitude $`\phi _{\mathrm{max}}^{}`$ in Fig. 4.5.
Next, we discuss the differential equation with time dependent $`g`$. It is sufficient if we deduce the effect on the amplitude $`\phi _{\mathrm{max}}^{}`$; we can then use $`\phi _{\mathrm{max}}^{}`$ and Figs. 4.4 and 4.5 to derive the oscillation frequency and $`\mathrm{cos}\phi ^{}`$.
We will assume that $`\dot{g}`$ is approximately constant over the time of half the oscillation period, $`\tau /2`$, where $`\tau =1/\nu 2\pi /\sqrt{g}`$. Then it is easy to calculate the change in the amplitude $`d\phi ^{}`$ during the time $`dt=\tau /2`$. Taking the differential of the potential energy, we get
$$dE_{\mathrm{pot}}^{\mathrm{max}}=dg(1\mathrm{cos}\phi _{\mathrm{max}}^{})+gd(\mathrm{cos}\phi ^{})$$
(4.68)
Since the acceleration of gravity, $`g`$, is no longer constant, the total energy is not conserved, so that
$`\dot{E}_{\mathrm{tot}}`$ $`=`$ $`\dot{E}_{\mathrm{kin}}+\dot{E}_{\mathrm{pot}}=\dot{\phi }^{}\ddot{\phi }^{}g\mathrm{sin}\phi ^{}\dot{\phi }^{}\dot{g}(1\mathrm{cos}\phi ^{})`$ (4.69)
$`=`$ $`\dot{g}(1\mathrm{cos}\phi ^{}),`$
where the first two terms cancel due to the differential equation (4.66). We get the differential
$$dE_{\mathrm{tot}}dg\left(1\mathrm{cos}\phi ^{}\right),$$
(4.70)
where we have time averaged $`\mathrm{cos}\phi ^{}`$. Equating (4.68) and (4.70), we get
$$d(\mathrm{cos}\phi _{\mathrm{max}}^{})=\frac{dg}{g}\left(\mathrm{cos}\phi _{\mathrm{max}}^{}\mathrm{cos}\phi ^{}\right),$$
(4.71)
If we use $`dg=\dot{g}dt=\dot{g}\tau /2`$ and $`d(\mathrm{cos}|\phi ^{}|)=d(\mathrm{cos}|\phi ^{}|)=\mathrm{sin}|\phi ^{}|d|\phi ^{}|`$, we finally get
$$d|\phi _{\mathrm{max}}^{}|=\frac{\dot{g}\tau /2}{g\mathrm{sin}|\phi _{\mathrm{max}}^{}|}\left(\mathrm{cos}\phi _{\mathrm{max}}^{}\mathrm{cos}\phi ^{}\right)$$
(4.72)
as the change in the amplitude during one half oscillation, where of course $`\mathrm{cos}\phi _{\mathrm{max}}^{}\mathrm{cos}\phi ^{}`$. We see that $`\dot{g}>0`$ will result in a decreasing amplitude, i.e. increasing $`|\phi _{\mathrm{max}}^{}|`$. Analogously, $`\dot{g}<0`$ will result in an increasing amplitude. If thereby $`|\phi _{\mathrm{max}}^{}|+d|\phi _{\mathrm{max}}^{}|<0`$, $`E_{\mathrm{tot}}>0`$, so the system passes the meta-stable point and accelerates on the other side instead of turning. Then it is better to use (4.70) instead of the differential for the amplitude to describe how the system evolves.
Our approximation breaks down if $`|\phi _{\mathrm{max}}^{}|+d|\phi _{\mathrm{max}}^{}|1`$, since then $`(\mathrm{cos}\phi _{\mathrm{max}}^{}\mathrm{cos}\phi ^{})0`$, $`\mathrm{sin}^1\phi _{\mathrm{max}}^{}\mathrm{}`$ and $`\tau \mathrm{}`$. Also, $`\dot{g}=\mathrm{const}`$ might not apply any longer.
Now we can already make a statement about the neutrino asymmetry, represented by $`\delta P_z=\dot{\phi }^{}/b_A`$. If $`E_{\mathrm{tot}}<0`$, the asymmetry will change sign at the turning points $`\pm \phi _{\mathrm{max}}^{}`$, i.e. the asymmetry oscillates. If $`E_{\mathrm{tot}}>0`$, the asymmetry will oscillate between $`\sqrt{2E_{\mathrm{tot}}}/b_A`$ and $`\sqrt{2(E_{\mathrm{tot}}+g)}/b_A`$, but will not change sign.
#### 4.7.3 The Factor $`𝒈`$
The next step is to derive $`\dot{g}`$. The factors $`b_A`$ and $`b_s`$ in $`g`$ only vary slowly with time, and we will take them to be constant. Thus the main contribution to $`\dot{g}`$ comes from $`\vartheta ^+\mathrm{sin}\phi ^+`$. For convenience, we define two new variables, $`\alpha =\vartheta ^+\mathrm{sin}\phi ^+`$ and $`\beta =\vartheta ^+\mathrm{cos}\phi ^+`$. Then $`g=b_sb_A\alpha `$, so that $`\dot{g}\dot{\alpha }`$. We get
$`\dot{\alpha }`$ $`=`$ $`D\alpha +(b_cb_T)\beta `$ (4.73)
$`\dot{\beta }`$ $`=`$ $`D\beta (b_cb_T)\alpha +b_s\mathrm{cos}\phi ^{},`$ (4.74)
where we have neglected all terms dependent on $`\vartheta ^{}`$.
Initially, $`\phi ^{}1`$, so we will first consider the simpler case where $`\mathrm{cos}\phi ^{}=1`$, $`\mathrm{sin}\phi ^{}=0`$. Then we can apply the quasi-static approximation for $`\vartheta ^+`$ and $`\phi ^+`$. Thus,
$`\alpha _{\mathrm{qs}}`$ $`=`$ $`{\displaystyle \frac{b_s(b_cb_T)}{D^2+(b_cb_T)^2}}`$ (4.75)
$`\beta _{\mathrm{qs}}`$ $`=`$ $`{\displaystyle \frac{b_sD}{D^2+(b_cb_T)^2}}.`$ (4.76)
Now we can determine $`\dot{g}=b_sb_A\dot{\alpha }_{\mathrm{qs}}`$ with the help of
$$\dot{\alpha }_{\mathrm{qs}}=H\frac{b_s(b_T\mathrm{\Delta }b^2+D^2\mathrm{\Delta }b+D^2b_c)}{[D^2+\mathrm{\Delta }b^2]^2},$$
(4.77)
where $`\mathrm{\Delta }b=b_cb_T`$.
For small $`\mathrm{\Delta }b`$, $`\dot{\alpha }_{\mathrm{qs}}`$ is positive. However, $`\dot{\alpha }`$ changes sign when the numerator becomes zero, i.e.
$$\mathrm{\Delta }b_\alpha =\frac{D^2}{2b_T}\left(1\sqrt{1+4\frac{b_cb_T}{D^2}}\right).$$
(4.78)
Using $`b_cb_T=D/k_D`$, we find that the term $`b_cb_T/D^2k_D^21`$. Therefore,
$$\mathrm{\Delta }b_\alpha D=k_Db_T.$$
(4.79)
This result can be expressed in terms of the dimensionless variable $`\xi =1T(t)/T`$. Then we get $`\xi _\alpha k_D/6`$, see (4.61).
In summary, $`g`$ will increase at first and then decrease for $`\mathrm{\Delta }b>\mathrm{\Delta }b_\alpha `$. We have plotted the parameter-independent dimensionless variables $`\alpha _{\mathrm{qs}}c/s`$ and $`\frac{d}{d\xi }\alpha _{\mathrm{qs}}c/s=\dot{\alpha }_{\mathrm{qs}}T_{\mathrm{res}}c/HTs`$ as a function of $`\xi `$ in Figs. 4.7 and 4.7, respectively.
We can see that $`\dot{\alpha }_{\mathrm{qs}}`$ is constant to lowest order for very small $`\xi `$. Therefore, we expand $`g`$ in terms of $`\xi `$,
$$g=6\frac{b_Ab_s^2b_T}{D^2}|_{T_{\mathrm{res}}}\xi +𝒪(\xi ^2).$$
(4.80)
This approximation is valid for $`\xi <\xi _\alpha `$.
Now we consider the situation where $`\phi ^{}`$ changes in time. We assume that the oscillation period $`\tau =𝒪(1/\sqrt{g})`$ is much smaller than the damping time scales $`1/D`$ and $`1/\mathrm{\Delta }b`$. Then we can approximate $`\mathrm{sin}\phi ^{}`$ and $`\mathrm{cos}\phi ^{}`$ by their time average values, $`0`$ and $`\mathrm{cos}\phi ^{}`$, respectively, when describing the evolution of $`\alpha `$ and $`\beta `$. Thus
$`\dot{\alpha }`$ $``$ $`D\alpha +\mathrm{\Delta }b\beta `$ (4.81)
$`\dot{\beta }`$ $``$ $`D\beta \mathrm{\Delta }b\alpha +b_s\mathrm{cos}\phi ^{}.`$ (4.82)
Note that here we can easily neglect $`\vartheta ^{}`$, since $`\vartheta _{\mathrm{qs}}^{}\mathrm{sin}\phi ^{}=0`$. The coefficient $`\mathrm{\Delta }b`$ in (4.81) and (4.82) induces an oscillation of the two variables, which will be damped by the coefficient $`D`$ toward their quasi-static approximations, given by
$`\alpha _{\mathrm{qs}}`$ $`=`$ $`{\displaystyle \frac{b_s\mathrm{\Delta }b}{D^2+\mathrm{\Delta }b^2}}\mathrm{cos}\phi ^{},`$ (4.83)
$`\beta _{\mathrm{qs}}`$ $`=`$ $`{\displaystyle \frac{b_sD}{D^2+\mathrm{\Delta }b^2}}\mathrm{cos}\phi ^{}.`$ (4.84)
#### 4.7.4 Proof for Oscillations
It is now easy to describe the evolution of our system. We will restrict ourselves to proving that the pendulum, and therefore $`\delta P_z`$, will oscillate for a certain range of mixing parameters. $`\delta P_z`$ will change sign at least once if $`\dot{g}`$ is positive during the first oscillation of $`\phi ^{}`$, since then the amplitude $`\phi _{\mathrm{max}}^{}`$ decreases.
We start with the initial values given in (4.67). For small $`\phi ^{}`$, we can linearize the differential equation (4.66), so that
$$\ddot{\phi }^{}g\phi ^{}.$$
(4.85)
Furthermore, we linearize $`g`$ as we have done in (4.80) and substitute $`dt`$ by $`d\xi `$, as shown in (4.62). Then we get
$$\frac{d^2\phi ^{}}{d\xi ^2}=h\xi \phi ^{},$$
(4.86)
where $`h=1.02\times 10^{20}s^2c^{4/3}|\delta m_{\mathrm{eV}}^2|^{2/3}`$. The solutions for this equation are the Airy-functions.
When $`\phi ^{}=𝒪(1)`$, the linear approximation breaks down, and $`\phi ^{}`$ will start oscillating with a frequency $`\nu \mathrm{\Gamma }<\sqrt{g}/2\pi `$. We want $`\dot{g}\dot{\alpha }`$ to be positive, so this first oscillation has to start before $`\xi _\alpha `$. Thus we demand $`\xi _\phi ^{}<\xi _\alpha `$, where $`\xi _\phi ^{}`$ is defined by $`\phi ^{}(\xi _\phi ^{})=1`$. We find numerically that we can fit the solution of this condition with
$$s^{3.1}|\delta m_{\mathrm{eV}}^2|10^{15.1}$$
(4.87)
for $`|\delta m^2|=10^9\text{}10^5\mathrm{eV}^2`$ with an error smaller than 3%.
Of course, we need to take into account that the quasi-static solution for $`\alpha `$ given in 4.75 breaks down, since $`\phi ^{}<1`$ during oscillations. However, as we can see in (4.81) and (4.82), $`\dot{\alpha }`$ does not depend on $`\mathrm{cos}\phi ^{}`$ directly, but indirectly through $`\mathrm{\Delta }b\beta `$. Thus, the change in $`\mathrm{cos}\phi ^{}`$ will first affect $`\dot{\alpha }`$ after a time scale $`1/\mathrm{\Delta }b`$. So if the condition $`\mathrm{\Delta }b\sqrt{g}`$ holds at $`\xi =\xi _\phi ^{}`$, then $`\dot{\alpha }`$ stays positive during the first few oscillation periods, and thus the amplitude will decrease.
Since $`\mathrm{\Delta }b/\sqrt{g}\sqrt{\xi }`$, we can use the stronger condition $`\mathrm{\Delta }b/\sqrt{g}1`$ for $`\xi =\xi _\alpha `$. Then we get
$$s|\delta m_{\mathrm{eV}}^2|^{1/6}3\times 10^6.$$
(4.88)
We have now found the conditions for which $`\delta P_z`$ oscillates.
We can also estimate the amplitude with which the neutrino asymmetry $`\delta P_z`$ oscillates. An upper bound is given if we use the maximal value for $`g`$, i.e. at $`\xi =\xi _\alpha `$, and assume $`E_{\mathrm{tot}}=0`$. Then $`\delta P_z=\dot{\phi }^{}/b_A`$ is maximal at $`\phi ^{}=\pi `$, i.e. $`E_{\mathrm{kin}}=E_{\mathrm{pot}}=2g`$. We derive
$`(\delta P_z)_{\mathrm{max}}`$ $`=`$ $`{\displaystyle \frac{1}{b_A}}(\dot{\phi }^{})_{\mathrm{max}}={\displaystyle \frac{1}{b_A}}\sqrt{2(E_{\mathrm{kin}})_{\mathrm{max}}}`$ (4.89)
$`=`$ $`{\displaystyle \frac{1}{b_A}}\sqrt{4g(\xi _\alpha )}{\displaystyle \frac{2}{b_A}}\sqrt{{\displaystyle \frac{b_s^2b_A}{2D}}}`$
$``$ $`1.4\times 10^2s|\delta m_{\mathrm{eV}}^2|^{1/6}c^{5/6}.`$
We can compare this result with the numerical solution given in Section 4.2. Numerically we find during the period where $`\delta P_z`$ oscillates that $`\delta P_z7.53\times 10^7`$ for the parameters $`\delta m^2=1\mathrm{eV}^2`$, $`s=10^4`$. This is in very good agreement with our analytical value $`(\delta P_z)_{\mathrm{max}}=1.4\times 10^6`$. The difference is due to the fact that when $`g`$ is maximal, the amplitude of $`\phi ^{}`$ is not maximal, and thus a factor of $`\sqrt{2}`$ has to be replaced by $`\sqrt{1+\mathrm{cos}\phi _{\mathrm{max}}^{}}`$. Furthermore, $`\alpha `$ is slightly decreased due to $`\mathrm{cos}\phi ^{}<1`$.
## Chapter 5 Conclusions
We have analytically examined neutrino oscillations between a sterile and an active neutrino in the Early Universe, neglecting the neutrino momentum distribution. Our main achievement has been to prove that this system, which is known to create a large neutrino asymmetry, exhibits oscillations of this neutrino asymmetry for a large range of mixing parameters, as shown in Fig. 5.1. We conclude that these asymmetry oscillations, which have been encountered in only some of the numerical calculations, can not arise from numerical instabilities.
With the methods presented in our work, the complete analytical treatment of this system seems to have become feasible. Naturally, the next step will be to include the effects from the neutrino momentum distribution on the oscillatory behavior. Then the analytical approach will clearly be superior to the numerical one, which continues to encounter many difficulties . Finally, the duration of the asymmetry oscillations and the power law of the asymmetry growth after the oscillations have ceased should be derived.
Understanding the mechanism of $`\nu _\alpha `$$`\nu _\mathrm{s}`$ oscillations in the Early Universe is very important. On the one hand, if the primordial abundances are determined with sufficient precision, it becomes possible to constrain the mixing parameters of the $`\nu _\alpha `$$`\nu _\mathrm{s}`$ system. As a result, some of the models which have been proposed to explain the current experimental situation could be excluded.
On the other hand, if future neutrino experiments, such as MiniBooNE , prove the existence of a sterile neutrino, it is necessary to understand the impact of the $`\nu _\mathrm{s}`$ on the primordial abundances to test whether BBN is a consistent theory.
## Appendix A Quasi-Static Approximation and its Range of Validity
The quasi-static approximation of a variable $`x`$ which follows the differential equation
$$\frac{}{t}x(t)=D(t)x(t)+f(t)$$
(A.1)
is given by
$$x_{\mathrm{qs}}=\frac{f}{D}.$$
(A.2)
If $`f`$ and $`D`$ depend on time, $`x_{\mathrm{qs}}`$ also changes its value. Therefore, $`x_{\mathrm{qs}}`$ will only be a good approximation if the quasi-static condition is fulfilled, i.e.
$$|\dot{x}_{\mathrm{qs}}|<Dx_{\mathrm{qs}}.$$
(A.3)
We have derived all values $`r_x\dot{x}_{\mathrm{qs}}/x_{\mathrm{qs}}`$ for our system of differential equations (4.29)–(4.34) assuming that all variables except for $`P_z^+`$ are at their static values, where $`H=\dot{T}/T`$.
$`r_{\vartheta ^+}`$ $``$ $`{\displaystyle \frac{\dot{\vartheta }_{\mathrm{qs}}^+}{\vartheta _{\mathrm{qs}}^+}}=6H\left(1+{\displaystyle \frac{b_c(b_cb_T)}{D^2+(b_cb_T)^2}}\right),`$ (A.4)
$`r_{\phi ^+}`$ $``$ $`{\displaystyle \frac{\dot{\phi }_{\mathrm{qs}}^+}{\mathrm{sin}\phi _{\mathrm{qs}}^+}}=6H{\displaystyle \frac{b_c}{b_cb_T}}{\displaystyle \frac{D}{\sqrt{D^2+(b_cb_T)^2}}}.`$ (A.5)
Here we have divided by $`\mathrm{sin}\phi _{\mathrm{qs}}^+`$ instead of $`\phi _{\mathrm{qs}}^+`$, since the differential equation for $`\phi ^+`$ has the form $`\frac{}{t}x(t)=D(t)\mathrm{sin}[x(t)]+f(t)`$. In the case of $`\phi ^{}`$, we can neglect the $`\mathrm{sin}`$ as long as $`\phi ^{}1`$. Furthermore,
$`r_{P_z^{}}`$ $``$ $`{\displaystyle \frac{(\dot{P}_z^{})_{\mathrm{qs}}}{(P_z^{})_{\mathrm{qs}}}}={\displaystyle \frac{\dot{\epsilon }}{\epsilon }}{\displaystyle \frac{\epsilon }{1+\epsilon }},`$ (A.6)
$`r_{\delta P_z}`$ $``$ $`{\displaystyle \frac{\frac{d}{dt}(\delta P_z)_{\mathrm{qs}}}{(\delta P_z)_{\mathrm{qs}}}}={\displaystyle \frac{\dot{\epsilon }}{\epsilon }}{\displaystyle \frac{1}{1+\epsilon }},`$ (A.7)
$`r_\vartheta ^{}`$ $``$ $`{\displaystyle \frac{\dot{\vartheta }_{\mathrm{qs}}^{}}{\vartheta _{\mathrm{qs}}^{}}}=r_{\vartheta ^+}+r_{P_z^{}}`$ (A.8)
$`=`$ $`H\left(86{\displaystyle \frac{b_c}{b_cb_T}}\left(13{\displaystyle \frac{(b_cb_T)^2}{D^2+(b_cb_T)^2}}\right)\right)+r_{\delta P_z},`$
$`r_\phi ^{}`$ $``$ $`{\displaystyle \frac{\dot{\phi }_{\mathrm{qs}}^{}}{\phi _{\mathrm{qs}}^{}}}=6H{\displaystyle \frac{b_c}{b_cb_T}}+r_{P_z^{}}`$ (A.9)
$`=`$ $`H\left(2+12{\displaystyle \frac{b_c(b_cb_T)}{D^2+(b_cb_T)^2}}\right)+r_{\delta P_z},`$
where
$`\epsilon `$ $`=`$ $`{\displaystyle \frac{D^2+(b_cb_T)^2}{2P_z^+b_A(b_cb_T)}},`$ (A.10)
$`{\displaystyle \frac{\dot{\epsilon }}{\epsilon }}`$ $`=`$ $`H\left(2+6{\displaystyle \frac{b_c}{b_cb_T}}\left(12{\displaystyle \frac{(b_cb_T)^2}{D^2+(b_cb_T)^2}}\right)\right).`$ (A.11)
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# 𝑑- and 𝑝-wave superconductivity mediated by spin fluctuations in two- and three-dimensional single-band repulsive Hubbard model
## 1 Introduction
There has been an increasing fascination with electronic mechanisms for pairing since the high-$`T_\mathrm{c}`$ superconductivity was discovered in the cuprates.$`^{\text{?}\text{)}}`$ Among various mechanisms, the pairing mediated by spin fluctuations has been proposed$`^{\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}\text{)}}`$ and have been examined intensively.
The pairing interaction, denoted by $`V^{(2)}`$, resulting from exchanging spin fluctuations is dominated by the spin susceptibility which has a strong wavenumber dependence, so that $`V^{(2)}`$ in turn has a strong wavenumber dependence. This determines the symmetry of the pairing, since this dictates the dominant pair-hopping processes on the Fermi surface. There, the dominant mode of spin fluctuations, hence the symmetry of the pairing, must be sensitively correlated with the lattice structure and filling in general. For example, the respulsive electron correlation causes the face centered cubic (FCC) lattice with low density of electrons to have strong ferromagnetic fluctuations, or the half-filled body centerd cubic (BCC) lattice to have strong antiferromagnetic fluctuations as recently confirmed.$`^{\text{?}\text{)}}`$ Thus, it is an interesting question to ask how the superconductivity would be correlated with the lattice structures.
As for the singlet, anisotropic (e.g., d-wave) pairing in 2D systems, calculations on a microscopic level have been performed: For the repulsive Hubbard model, which is a simplest possible model of the electron correlation, the fluctuation exchange approximation (FLEX) developed by Bickers et al.,$`^{\text{?}\text{}\text{?}\text{)}}`$ has been employed to show the occurrence of the $`d`$-wave superconductivity in a square lattice.$`^{\text{?}\text{}\text{?}\text{)}}`$ As for other 2D lattices, several authors have suggested the occurrence of $`d`$-wave superconductivity in the Hubbard model on 2D anisotropic triangular lattice, which represents an organic superconductor (BEDT-TTF)<sub>2</sub>X .$`^{\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}\text{)}}`$ These results indicate that $`T_\mathrm{c}`$ for the superconductivity near the antiferromagnetic instability in 2D is $`O(0.01t)`$ ($`t`$: transfer integral), i.e., two orders of magnitude smaller than the original electronic energy, but still ‘high $`T_\mathrm{c}`$$`O(100`$ K) if we take $`tO(1`$ eV) for the cuprates (while $`T_\mathrm{c}O(10`$ K) if we take $`tO(0.1`$ eV) for the organic conductors).
As for 3D systems, Scalapino et al$`^{\text{?}\text{)}}`$ showed that paramagnon exchange near a spin-density wave instability gives rise to a strong $`d`$-wave pairing interaction for the 3D Hubbard model on the simple cubic lattice, but $`T_\mathrm{c}`$ was not discussed there. Nakamura et al$`^{\text{?}\text{)}}`$ extended Moriya’s spin fluctuation theory of superconductivity$`^{\text{?)}}`$ to 3D systems, and concluded that $`T_\mathrm{c}`$ is similar between the 2D and 3D cases provided that common parameter values (scaled by the band width) are taken. However, the parameters there are phenomelogical ones, so it is not clear whether the result remains valid for microscopic models.
The possibility of the triplet pairing, on the other hand, has been investigated since the 1960’s for superfluid $`{}_{}{}^{3}\mathrm{He}`$,$`^{\text{?}\text{)}}`$ a heavy fermion system $`\mathrm{UPt}_3`$,$`^{\text{?}\text{)}}`$ or most recently, an oxide $`\mathrm{Sr}_2\mathrm{RuO}_4`$.$`^{\text{?}\text{}\text{?}\text{)}}`$ The experimental results suggesting $`p`$-wave pairing in these materials have stimulated theoretical studies for electron-repulsion originated superconductivity. For the electron gas model, Chubukov extended the Kohn-Luttinger theorem,$`^{\text{?}\text{)}}`$ which asserts that the repulsively interacting electron gas should be instable against pairing formations at low enough temperatures, to $`p`$-wave pairing for 2D and 3D electron gas in the dilute limit by analyzing the singularity of the scattering amplitude.$`^{\text{?}\text{}\text{?}\text{)}}`$ Takada$`^{\text{?}\text{)}}`$ discussed the possibility of $`p`$-wave superconductivity in the dilute electron gas with the Kukkonen-Overhauser model.$`^{\text{?}\text{)}}`$ This model considers the effective electron-electron interaction composed of the bare interaction and the interactions mediated by charge/spin fluctuations that contains the so-called local-field correction.
As for lattice systems, which is the subject of the present paper, 2D Hubbard model with large enough next-nearest-neighbor hopping $`(t^{})`$ has been shown to exhibit $`p`$-pairing for dilute band fillings.$`^{\text{?}\text{)}}`$ Hlubina$`^{\text{?}\text{)}}`$ reached a similar conclusion by evaluating the superconducting vertex in a perturbative way.$`^{\text{?}\text{)}}`$ However, the energy scale of the $`p`$-pairing in the Hubbard model, i.e., $`T_\mathrm{c}`$, has not been questioned so far. Using a phenomenological approach, Monthoux and Lonzarich$`^{\text{?}\text{)}}`$ have recently concluded for 2D systems that the $`d`$-wave pairing is much stronger than $`p`$-wave pairing.
Given these backgrounds, from microscopic view of point, we investigate in this paper various 2D and 3D lattice structures (square, triangular, simple cubic, BCC, FCC) with systematically varied next-nearest neighbor hopping to tune the dispersion and varied band filling. Specifically, we address the following fundamental questions:
(i) Can the pairing instability in 3D be stronger than that in 2D?
(ii) Can the pairing instability with other symmetry, i.e., spin-triplet $`p`$-pairing in the presence of ferromagnetic spin fluctuations, become strong?
A part of the present study has been briefly reported,$`^{\text{?}\text{)}}`$ while here we extend the calculation to various cases in order to extensively confirm our previous conclusions. Namely, we study possibility of $`d`$-wave superconductivity for nearly half-filled square, simple cubic (SC), and BCC lattices where strong antiferromagnetic fluctuations are present, along with possibility of $`p`$-wave superconductivity for low density square lattice with significant next nearest-neighbor hopping, quarter-filled triangular lattice, and low-density FCC lattice where ferromagnetic spin fluctuations should be dominant.
We employ the FLEX approximation, which enables us to handle strong spin fluctuations. As seen from the result, to determine the best situation for superconductivity in the Hubbard model we have to consider various factors such as the form of the pairing interaction and the energy/momentum-dependence of Green’s functions (lifetime of quasi-particles, etc), so that the way in which the pairing instability is correlated with the lattice structure is a highly nontrivial problem which is by no means predictable from the outset.
Here we shall show that
(i) $`d`$-wave instability mediated by antiferromagnetic spin fluctuations is stronger than $`p`$-wave instability mediated by ferromagnetic spin fluctuations both in 2D and 3D, and
(ii) pairing instability in 2D is much stronger than that in 3D.
Thus the ‘best’ situation for the spin fluctuation mediated pairing is suggested to be the 2D case with dominant antiferromagnetic fluctuations as far as the single-band Hubbard model on ordinary lattices are concerned.
## 2 Formulation
### 2.1 Model Hamiltonian
We consider the single-band Hubbard model,
$$=\underset{i,j\sigma }{}t_{ij}c_{i\sigma }^{}c_{j\sigma }+U\underset{i}{}n_in_i,$$
(1)
where $`c_{i\sigma }^{}`$ creates an electron at the $`i`$-th site with spin $`\sigma `$, $`n_{i\sigma }c_{i\sigma }^{}c_{i\sigma }`$ is the number operator. We consider the transfer between second-nearest neighbors, $`t_{ij}=t^{}`$, along with $`t(=1`$ hereafter) for nearest neighbors.
The energy dispersions for square and triangular lattices are,
$`\epsilon ^{\mathrm{}}(𝐤)`$ $`=`$ $`2t{\displaystyle \underset{i=1}{\overset{2}{}}}\mathrm{cos}(k_i)+4t^{}\mathrm{cos}(k_1)\mathrm{cos}(k_2),`$ (2)
$`\epsilon ^{\mathrm{}}(𝐤)`$ $`=`$ $`2t{\displaystyle \underset{i=1}{\overset{2}{}}}\mathrm{cos}(k_i)+2t\mathrm{cos}(k_1+k_2),`$ (3)
respectively, where we take a square Brillouin zone for the latter case as well by inserting diagonal transfers in the square lattice.
The dispersions for SC, BCC and FCC lattices are given as
$`\epsilon ^{\mathrm{SC}}(𝐤)`$ $`=`$ $`2t{\displaystyle \underset{i=1}{\overset{3}{}}}\mathrm{cos}k_i+4t^{}{\displaystyle \underset{i<j}{}}\mathrm{cos}k_i\mathrm{cos}k_j,`$ (4)
$`\epsilon ^{\mathrm{BCC}}(𝐤)`$ $`=`$ $`8t\mathrm{cos}k_1\mathrm{cos}k_2\mathrm{cos}k_3+2t^{}{\displaystyle \underset{i=1}{\overset{3}{}}}\mathrm{cos}(2k_i),`$ (5)
$`\epsilon ^{\mathrm{FCC}}(𝐤)`$ $`=`$ $`4t{\displaystyle \underset{i<j}{}}\mathrm{cos}k_i\mathrm{cos}k_j+2t^{}{\displaystyle \underset{i=1}{\overset{3}{}}}\mathrm{cos}(2k_i),`$ (6)
respectively, where $`(k_1,k_2,k_3)(k_x,k_y,k_z)`$. To facilitate the calculation here we take a cubic Brillouin zone ($`\pi <k_i\pi `$) by considering a four(two) equivalent, interpenetrating BCC(FCC) lattices.
The density of states for the non-interacting case is displayed in Fig. 1 for the square lattice with $`t^{}=0.0,0.5`$, triangular lattice, SC lattice with $`t^{}=0.0,0.1,0.3`$, FCC lattice with $`t^{}=0.0,0.5`$ and BCC lattice with $`t^{}=0.0,0.1`$.
### 2.2 Method
The FLEX, introduced by Bickers et al.,$`^{\text{?, ?)}}`$ starts from a set of skeleton diagrams for the thermodynamic potential, $`\mathrm{\Phi }`$, introduced by Luttinger and Ward. $`\mathrm{\Phi }`$ is a functional of Green’s function, $`G`$, and a ($`k`$-dependent) self energy can be computed by a functional derivative of $`\mathrm{\Phi }`$ with $`G`$. For the calculation of $`\mathrm{\Phi }`$ the idea of Baym and Kadanoff$`^{\text{?}\text{}\text{?}\text{)}}`$ of taking an important series of diagrams is employed. Hence the FLEX approximation is a self-consistent perturbation approximation with respect to on-site interaction $`U`$. The FLEX approximation is suitable for the analysis of Fermi liquid with strong spin fluctuations.
The self energy is obtained as
$`\mathrm{\Sigma }(k)={\displaystyle \frac{1}{N}}{\displaystyle \underset{q}{}}G(kq)V^{(1)}(q),`$ (7)
in which RPA-type bubble and ladder diagrams are collected for the interaction,
$`V^{(1)}(q)`$ $`=`$ $`{\displaystyle \frac{1}{2}}U^2\chi _{\mathrm{irr}}(q)\left[{\displaystyle \frac{1}{1+U\chi _{\mathrm{irr}}(q)}}\right]`$
$`+{\displaystyle \frac{3}{2}}U^2\chi _{\mathrm{irr}}(q)\left[{\displaystyle \frac{1}{1U\chi _{\mathrm{irr}}(q)}}\right]U^2\chi _{\mathrm{irr}}(q)`$
with
$$\chi _{\mathrm{irr}}(q)=\frac{1}{N}\underset{k}{}G(k+q)G(k).$$
Here we have denoted $`q(𝐪,iϵ_\nu )`$ and $`k(𝐤,i\omega _n)`$, $`ϵ_\nu =2\pi \nu T`$ is the Matsubara frequency for bosons while $`\omega _n=(2n1)\pi T`$ for fermions, and $`N`$ is the total number of sites. For simplicity, we neglect the diagrams in the particle-particle channel.
The Dyson equation is written as
$$G(𝐤,\omega _n)^1=G^0(𝐤,\omega _n)^1\mathrm{\Sigma }(𝐤,\omega _n),$$
(8)
where $`G^0`$ is the bare Green’s function,
$$G^0(𝐤,\omega _n)=\frac{1}{\mathrm{i}\omega _n+\mu \epsilon _\mathrm{k}^0},$$
(9)
with $`\epsilon _\mathrm{k}^0`$ being the energy of a free electron. We have solved the equations (7) $``$ (8) by setting the chemical potential $`\mu `$ so as to fix the density of electrons.
To obtain $`T_\mathrm{c}`$, we solve, with the power method, the eigenvalue (Éliashberg) equation,
$`\lambda \mathrm{\Sigma }^{(2)}(k)`$ $`=`$ $`{\displaystyle \frac{T}{N}}{\displaystyle \underset{k^{}}{}}\mathrm{\Sigma }^{(2)}(k^{})|G(k^{})|^2V^{(2)}(k,k^{}),`$ (10)
where $`G(k)`$ the dressed Green’s function, and $`\mathrm{\Sigma }^{(2)}(k)`$ the anomalous self energy, and $`T=T_\mathrm{c}`$ corresponds to the point at which the maximum eigenvalue $`\lambda _{\mathrm{Max}}`$ reaches unity.
The interaction $`V^{(2)}`$ originates from the transverse spin fluctuations, longitudinal spin fluctuations and charge fluctuations, namely,
$`V^{(2)}(k,k^{})`$ $`=`$ $`U^2[{\displaystyle \frac{1}{2}}\chi _{\mathrm{ch}}(kk^{})`$
$`{\displaystyle \frac{1}{2}}\chi ^{zz}(kk^{})+\chi ^\pm (k+k^{})]`$
$`=`$ $`\left[{\displaystyle \frac{U^3\chi _{\mathrm{irr}}^2(kk^{})}{1U^2\chi _{\mathrm{irr}}^2(kk^{})}}\right]\left[{\displaystyle \frac{U^2\chi _{\mathrm{irr}}(k+k^{})}{1U\chi _{\mathrm{irr}}(k+k^{})}}\right]`$
where $`\chi _{\mathrm{ch}}`$ is the charge susceptibility, while $`\chi ^{zz}(\chi ^\pm )`$ is the longitudinal (transverse) spin susceptibility in the RPA form where the dressed Green’s function is used.
Since we have $`\mathrm{\Sigma }^{(2)}(k)=\mathrm{\Sigma }^{(2)}(k)`$ for the spin-singlet pairing, whereas $`\mathrm{\Sigma }^{(2)}(k)=\mathrm{\Sigma }^{(2)}(k)`$ for the spin-triplet pairing, $`V^{(2)}(k,k^{})`$ becomes a function of $`kk^{}=q`$ with
$`V^{(2)}(q)={\displaystyle \frac{3}{2}}\left[{\displaystyle \frac{U^2\chi _{\mathrm{irr}}(q)}{1U\chi _{\mathrm{irr}}(q)}}\right]+{\displaystyle \frac{1}{2}}\left[{\displaystyle \frac{U^2\chi _{\mathrm{irr}}(q)}{1+U\chi _{\mathrm{irr}}(q)}}\right]`$ (11)
for the singlet pairing and
$`V^{(2)}(q)={\displaystyle \frac{1}{2}}\left[{\displaystyle \frac{U^2\chi _{\mathrm{irr}}(q)}{1U\chi _{\mathrm{irr}}(q)}}\right]+{\displaystyle \frac{1}{2}}\left[{\displaystyle \frac{U^2\chi _{\mathrm{irr}}(q)}{1+U\chi _{\mathrm{irr}}(q)}}\right]`$ (12)
for the triplet pairing. We take $`N=64^2`$ sites with $`n_c=2048`$ Matsubara frequencies for 2D square lattice and with $`n_c=1024`$ Matsubara frequencies for 2D triangular lattice, and $`N=32^3`$ with $`n_c=1024`$ for 3D.
## 3 Results
### 3.1 Square lattice with antiferromagnetic spin fluctuations
Let us start with the case of nearly half-filled square lattice, which has strong antiferromagnetic fluctuations. In this case, $`T_\mathrm{c}0.02`$ for $`d`$-wave pairing should be obtained within the FLEX approximation as mentioned in the Introduction. We first present the result for $`t^{}=0.0`$, $`n=0.85`$ (0.15 holes doped) and $`U=4`$, which will serve as a reference for other lattices.
In Fig. 2, we show $`|G(𝐤,\mathrm{i}\pi k_\mathrm{B}T)|^2`$, a quantity which appears in the right-hand side of the Éliashberg equation (10). We can see that $`|G|^2`$ takes large values ($`8.0`$) around the Fermi surface. Note that the smaller the self energy correction, the peak of $`|G(𝐤,\mathrm{i}\pi k_\mathrm{B}T)|^2`$ becomes larger and the Fermi-liquid picture becomes more valid. Also, a large $`|G(𝐤,\mathrm{i}\pi k_\mathrm{B}T)|^2`$ favors superconductivity as we can see from the Éliashberg equation (10).
In Fig. 3, we plot the susceptibility in the RPA form, $`\chi (𝐤,0)=\chi _{\mathrm{irr}}/(1U\chi _{\mathrm{irr}})`$, as a function of the wave number for $`T=0.03`$. Dominant antiferromagnetic spin fluctuations are seen as $`\chi `$ peaked around $`𝐤=(\pi ,\pi )`$. A large peak in $`\chi (𝐤,0)`$ should imply a large pairing interaction as seen in the Éliashberg equation (10). The spread of $`\chi (𝐤,0)`$ around the peak in the momentum space is also important as we shall come back later, because it measures the fraction of effective channels that contribute to $`V^{(2)}`$ in the Éliashberg equation (10).
We also plot in Fig. 3 Im$`\chi (𝐤_{\mathrm{Max}},\omega )`$ against $`\omega `$, where $`𝐤_{\mathrm{Max}}`$ is the momentum for which $`\chi (𝐤,0)`$ becomes maximum and we have normalized Im$`\chi (𝐤_{\mathrm{Max}},\omega )`$ with the maximum value. We have obtained the dependence on real $`\omega `$ with an analytic continuation in the Padé approximation.$`^{\text{?}\text{)}}`$ The spread of Im$`\chi (𝐤_{\mathrm{Max}},\omega )`$ around the peak in the $`\omega `$ sector is an electronic couterpart to the Debye frequency $`\omega _D`$ in the BCS theory for the electron-phonon system, so this quantity is another important factor in the pairing instability.
In Fig. 4, we plot $`\lambda _{\mathrm{Max}}`$ as a function of temperature $`T`$ (normalized by $`t`$) along with the reciprocal of the peak value of $`\chi (𝐤,0)`$. The pairing instability is measured by how $`\lambda _{\mathrm{Max}}`$ is close to unity, while $`1/\chi 0`$ signifies the magnetic ordering. $`\lambda _{\mathrm{Max}}`$ is seen to be extrapolated to unity at $`T0.02`$, in accord with previous results.$`^{\text{?)}}`$
Here let us comment on a finite $`T_\mathrm{c}`$ in 2D systems. As usually done, it is taken to be a measure of $`T_\mathrm{c}`$ when the layers are stacked with the Josephson coupling. However, if we take account of the superconducting fluctuations rigorously, $`T_\mathrm{c}`$ for a purely 2D system must be zero according to Mermin’s theorem.$`^{\text{?}\text{}\text{?}\text{}\text{?}\text{)}}`$ To judge whether a mean field treatment of pairing instabilities is adequate will require an evaluation of the coherence length of the pairing. Although this has been done for the conventional phonon mechanism of superconductivity with Gor’kov’s argument,$`^{\text{?}\text{)}}`$ such an evaluation is rather difficult in the present case of an electron mechanism, since the pairing potential strongly depends on frequency and wave-number.
We can on the other hand study the effect of a weak three dimensionality on the $`\lambda _{\mathrm{Max}}`$ within the present formalism. Specifically, we introduce an inter-layer hopping, $`t_z`$. Figure 5 shows the results for $`\lambda _{\mathrm{Max}}`$ for a 3D anisotropic cubic lattice with the inter-layer hopping varied over $`0.0<t_z<1.0`$ on top of $`t_x=t_y=1.0`$. We can see that weak but finite $`t_z<0.3`$ does not appreciably change the behavior of $`\lambda _{\mathrm{Max}}`$. Thus, provided the coherence length in $`z`$-direction is large enough for $`t_z0.3`$ to validate the mean-field treatment, $`T_\mathrm{c}`$ obtained here for 2D can indeed be used as a measure of $`T_\mathrm{c}`$ when a weak three dimensionality is present.
### 3.2 $`t`$-$`t^{}`$ square lattice with ferromagnetic spin fluctuations
We move on to a 2D case with dominant ferromagnetic spin fluctuations, where triplet pairing is expected. The situation for which the ferromagnetic fluctuations become dominant has extensively been investigated for the Hubbard model with various approaches, and one guiding principle is that a large density of states at the Fermi level located near the bottom of the band should favor ferromagnetism for sufficiently strong electron-electron repulsion.
For the square lattice, such a situation may be realized for relatively large $`t^{}(0.5)`$ for dilute electron densities. In this case, the divergence (van Hove singularity) in the density of states, which is at the center and has a functional form $`|\mathrm{ln}E|`$ for $`t^{}=0`$, shifts toward the band bottom with $`D(E)1/\sqrt{(2+E)}|\mathrm{ln}(2+E)|`$ for $`t^{}=0.5`$, so that the density of states at the Fermi level becomes large for the dilute case, favoring ferromagnetism.$`^{\text{?}\text{)}}`$ It has in fact been shown from quantum Monte Carlo studies that the ground state is fully spin-polarized for $`t^{}=0.47`$, $`n0.4`$.$`^{\text{?}\text{, ?)}}`$ So we explore the possibility of $`p`$-wave instability associated with this ferromagnetism, and compare the result with that of $`d`$-wave instability associated with the antiferromagnetism discussed in the previous subsection.
First, in Fig. 6, we show $`\lambda _{\mathrm{Max}}`$ for the density varied over $`0.3n0.6`$ and $`t^{}`$ varied over $`0.3t^{}0.6`$ for $`U=4,6`$ with $`T=0.03`$. We can see that $`\lambda _{\mathrm{Max}}`$ becomes largest for $`n=0.3`$ and $`t^{}=0.5`$ for both $`U=4`$ and 6, so we take these parameter sets.
In Fig. 7, we plot $`\chi (𝐤,0)`$ for $`U=4`$ as a function of momentum along with Im$`\chi (𝐤_{\mathrm{Max}},\omega )`$ as a function of $`\omega `$. The peak of $`\chi (𝐤,0)`$ is indeed located at the $`\mathrm{\Gamma }`$ point $`(𝐤=(0,0))`$ indicating ferromagnetic fluctuations. The frequency spread is similar to the case of the nearly half-filled $`t^{}=0`$ square lattice. In Fig.8, we plot $`\lambda _{\mathrm{Max}}`$ as a function of $`T`$. We can see that $`\lambda _{\mathrm{Max}}`$ is much smaller than that in the antiferromagnetic case, Fig.4.
### 3.3 Why is p-pairing weaker than d-pairing?
The present result that the p-pairing has a lower $`T_\mathrm{c}`$ contrasts with a naive expectation from the BCS picture, in which the $`T_\mathrm{c}`$ should be high for a large density of states at the Fermi level. We may trace back the reason why this does not apply as follows.
First, if we look at the dominant ($`1/[1U\chi _{\mathrm{irr}}(q)]`$) term in the pairing potential $`V^{(2)}`$ itself in eqs. (11) and (12), the triplet pairing interaction is only one-third of the singlet pairing interaction. So this should be one reason.
On the other hand, the large density of states, namely the flatness of the band around the Fermi level, is reflected to the fact that $`|G|^2`$ (Fig.9) forms an almost flat plateau in a large portion of the Brillouin zone. If the maximum value of $`|G|^2`$ is large enough to compensate the disadvantage of the one-third $`V^{(2)}`$, a large $`\lambda `$ may emerge. However, we can see that $`|G|^2`$ is much smaller than that in the antiferromagnetic case (Fig.2), which implies that the self energy correction is large. Even when we take a larger repulsion $`U`$ to increase the triplet pairing attraction (i.e., to increase the susceptibility), this makes the self-energy correction even greater, resulting in only a slight change in $`\lambda `$.
### 3.4 2D triangular lattice
Next, we discuss the case of 2D triangular lattice. As mentioned in the Introduction, the half-filled case with dominant antiferromagnetic fluctuations has already been discussed by a number of authors. $`^{\text{?, ?, ?, ?, ?, ?)}}`$
Here, we focus on the quarter-filled isotropic triangular lattice, where we can expect ferromagnetic fluctuations, and hence $`p`$-wave pairing. As we can see in Fig. 1, the density of states for the triangular lattice has a sharp peak (with $`D(E)|\mathrm{ln}(E+2)|`$) near the bottom of the band, so a dominant ferromagnetic fluctuation is expected if the Fermi level is located at the peak. Hanish et al.$`^{\text{?}\text{)}}`$ studied the instability of the fully-ferromagnetic state of Hubbard model on the triangular lattice for large $`U`$ and concluded that the ferromagnetic state is stable for $`n0.5`$ (quarter-filled).
In this situation, we have calculated $`\lambda _{\mathrm{Max}}`$ for $`T=0.03`$, $`U=4,8,12`$. Since $`\lambda _{\mathrm{Max}}`$ takes similar values for $`U=4,8,12`$, we take $`U=8`$.
In Fig. 10, we plot the wave-number dependence of $`\chi (𝐤)`$ and frequency dependence of Im$`\chi (𝐤_{\mathrm{Max}},\omega )`$. We can see that the peak of $`\chi (𝐤,0)`$ is located around the $`\mathrm{\Gamma }`$ point. In Fig. 11, we plot $`\lambda _{\mathrm{Max}}`$ along with the reciprocal of the peak value of $`\chi (𝐤,0)`$ as a function of $`T`$. We can see that the $`p`$-wave instability here is again much weaker than the $`d`$-wave instability for the nearly half-filled square lattice (Fig. 4). The one-third factor discussed in the previous subsection should again be one factor. The second factor is also involved, namely, the maximum value of $`U^2|G(𝐤,ik_\mathrm{B}T)|^2`$ is small ($`1.0\times 8^2`$, see Fig. 12) compared to that in Fig. 2 ($`8.0\times 4^2`$), so the quasi-particles are short-lived.
We note that the ferromagnetic fluctuation in the triangular case is much weaker than in the square lattice with large $`t^{}`$ (compare Figs. 10(b) and 7(b)), although we have taken a larger $`U`$ here. This may be because the peak in the density of states in the triangular lattice is situated slightly above the bottom (compare Figs.1 (a) and (b)). Onoda and present authors have suggested in ref.$`^{\text{?)}}`$ that closer the peak of the density of states is to the band bottom, the stronger the ferromagnetic fluctuations tend to be.
Despite this difference in ferromagnetic fluctuation, the $`\lambda _{\mathrm{Max}}`$ in the present case is comparable to that in the square lattice with $`t^{}`$. This may be because the frequency spread in the present case is quite large as seen in Fig.10(a).
### 3.5 3D simple cubic lattice
Let us now move on to the case of 3D systems. we first discuss the possibility of $`d`$-wave pairing in the nearly half-filled Hubbard model on the 3D SC lattice. In this case, we find that the $`\mathrm{\Gamma }_3^+`$ representation ($`x^2y^2`$ etc, with the gap function $`\mathrm{\Delta }(𝐤)\mathrm{cos}k_x\mathrm{cos}k_y`$ etc) of O<sub>h</sub> group$`^{\text{?}\text{)}}`$ has the largest $`\lambda _{\mathrm{Max}}`$, so we concentrate on this pairing symmetry.
In Fig. 13, we show $`\lambda _{\mathrm{Max}}`$ for the density varied over $`0.75n0.9`$ and $`t^{}`$ varied over $`0.5t^{}+0.4`$ for $`U=4,6,8,10,12`$ with $`T=0.03`$. We can see that among these parameter sets, $`\lambda _{\mathrm{Max}}`$ becomes largest for $`n=0.8`$, $`t^{}=0.20.3`$ and $`U=810`$, so we take these parameter sets.
In Fig. 14(a), we show the momentum dependence of $`\chi (𝐤,0)`$ for $`U=8`$, $`n=0.8`$, $`t^{}=0.3`$, $`T=0.03`$. We can see that the peak of $`\chi (𝐤,0)`$ is located around the K point $`(\pi ,\pi ,\pi )`$ as expected for the antiferromagnetism.
In Fig. 15, we plot $`\lambda _{\mathrm{Max}}`$ along with the reciprocal of the peak value of $`\chi (𝐤,0)`$ as a function of $`T`$ for $`t^{}=0.2,0.3`$ ,$`U=8`$ and $`n=0.8`$. We see that $`\lambda _{\mathrm{Max}}`$ does not become very close to unity in the range calculated here. For $`T<0.02`$, the result obtained for $`N=32^3`$ and $`n_c=1024`$ is not convergent enough with respect to the system size and the number of Matsubara frequencies, so we show the result for a larger $`n_c=2048`$ in the inset, which suggests that $`\lambda _{\mathrm{Max}}`$ tends to increase with $`n_c`$. We have also performed a calculation for $`N=16^3`$, and found that $`\lambda _{\mathrm{Max}}`$ also increases with $`N`$, so a finite $`T_\mathrm{c}`$ ($`<0.01`$) may be obtained at least for $`t^{}=0.3`$, $`U=8`$, $`n=0.8`$ in the limit of large $`N`$ and $`n_c`$. At any event, $`T_\mathrm{c}`$ is significantly smaller than in the square lattice with strong antiferromagnetic fluctuations (Fig. 4).
### 3.6 Why is $`d`$-pairing stronger in 2D than in 3D?
So the $`d`$-wave instability is decidedly stronger in 2D than in 3D as far as the square and simple cubic lattices are concernded, and the question is: what are physical reasons for that. We can pinpoint the origin by looking at the factors involved in the Éliashberg equation (10), i.e., (a) the factor $`U^2|G|^2`$, (b) the summation over the frequency, and (c) the summation over the momentum. In addition, the factor $`V^{(2)}`$ is of course important in the equation, but in the following we compare 2D and 3D in the situation where the maximum value of $`\chi (𝐤,0)`$ (that determines $`V^{(2)}`$) is similar between the two cases to concentrate on the factors (a)(b)(c).
In Fig. 16 we plot $`|G|^2`$ in 3D for $`k_z=0,\pi /2,\pi `$ as a function of $`k_x`$ and $`k_y`$ for $`U=8`$, $`n=0.8`$, $`t^{}=0.3`$, $`T=0.03`$. We can see that the maximum value of $`U^2|G|^2`$ is greater in 3D than in 2D.
Fig. 14(b) displays $`\mathrm{Im}\chi (𝐤_{\mathrm{Max}},\omega )`$ as a function of $`\omega `$. The figure compares the SC lattice ($`t^{}=0.2,n=0.8,U=8`$) with a typical square lattice with $`t^{}=0`$, $`n=0.85`$ and $`U=4`$ having a similar magnitude of $`\chi `$ at $`T=0.03`$. We can see that the frequency spread of $`\mathrm{Im}\chi (𝐤_{\mathrm{Max}},\omega )`$ is similar between 3D and 2D, so the factors (a)(b) can be excluded from the reason for the 2D-3D difference. Note that if the frequency spread of the susceptibility is scaled not by $`t`$ but by the band width, as Nakamura et al$`^{\text{?)}}`$ have assumed, $`\lambda _{\mathrm{Max}}`$ would have become larger.
If we turn to the momentum sector in the susceptibility, $`\chi (𝐤,0)`$, Fig. 14(a) shows that the width, $`a`$, of the $`\chi (𝐤,0)`$ peak in each momentum component is similar to those in 2D displayed in Fig.3, where the main contribution of $`V^{(2)}`$ to $`\lambda `$ is confined around $`(\pi ,\pi )`$ in 2D or $`(\pi ,\pi ,\pi )`$ in 3D. Since the right-hand side of the Éliashberg equation (10) is normalized by $`NL^D`$ with $`L`$ being the linear dimension of the system, $`\lambda `$ is proportional to $`(a/L)^D`$, so that if $`a`$ has similar values between 2D and 3D we end up with a smaller $`\lambda `$ in 3D than that in 2D. So we can conclude that this is the main reason why the 2D case is more favorable than the 3D case.
### 3.7 BCC lattice
Let us turn to the BCC lattice. For BCC lattice, the density of states diverges around the center with $`D(E)[\mathrm{ln}(E)]^2`$, and the antiferromagnetic fluctuation is dominant near half-filling,$`^{\text{?)}}`$ so we focus on the possibility of $`d`$-wave superconductivity. For the d-wave, we found that $`\mathrm{\Gamma }_5^+`$ representation ($`xy`$ etc, with the gap function $`\mathrm{\Delta }(𝐤)\mathrm{cos}(k_x+k_y+k_z)\mathrm{cos}(k_x+k_yk_z)`$, etc) of $`O_h`$ group has the largest $`\lambda _{\mathrm{Max}}`$.
In Fig. 17, we show $`\lambda _{\mathrm{Max}}`$ for the density varied over $`0.75n0.9`$ and $`t^{}`$ varied over $`0.4t^{}0.4`$ for $`U=4,6,8`$ with $`T=0.03`$. Antiferromagnetic spin fluctuations are much stronger in BCC lattice than in the SC lattice as pointed out in ref. .$`^{\text{?)}}`$ In fact, in this figure, the truncated curves for $`\lambda _{\mathrm{Max}}`$ means that the system becomes antiferromagnetic (i.e., $`U\chi _0`$ becomes unity) there. Hereafter, we focus on the case of $`U=6`$, $`t^{}=0.1`$, $`n=0.75`$.
In Fig. 18(a), we plot the momentum dependence of $`\chi (𝐤,0)`$. We can see that the peak of $`\chi (𝐤,0)`$ is located around $`K`$-point ($`𝐤=(\pi ,\pi ,\pi ))`$. In Fig. 18(b), we plot the frequency dependence of Im$`\chi (𝐤_{\mathrm{Max}},\omega )`$ which shows that the frequency spread for the BCC lattice is similar to the case of the square lattice or the SC lattice. In Fig. 19, we plot $`\lambda _{\mathrm{Max}}`$ for $`d`$-wave pairing as a function of $`T`$. $`\lambda _{\mathrm{Max}}`$ is again much smaller than that for the nearly half-filled square lattice, and is even smaller than that for the SC lattice. The reason for the former is mainly due to the $`(a/L)^D`$ factor discussed above, while the reason for the latter is because the maximum in $`U^2|G|^2`$ in BCC lattice is smaller (see Fig. 20) than that in SC lattice (Fig. 16).
### 3.8 FCC lattice
We finally come to the FCC lattice. The density of states of FCC lattice diverges at the bottom of the band with $`D(E)|\mathrm{ln}(E+4)|`$ for $`t^{}=0`$ and $`1/\sqrt{E+3}`$ for $`t^{}=0.50`$, so we can expect large ferromagnetic spin fluctuations for low densities of electrons. In fact, the possibility of ferromagnetic ground state has been discussed using various approaches.$`^{\text{?, }\text{?}\text{, ?)}}`$ According to our previous study,$`^{\text{?)}}`$ the ferromagnetic spin fluctuation is most dominant for $`n0.2`$ and $`t^{}0.5`$ in the weak coupling regime. Here, we focus on the possibility of $`p`$-wave pairing at low densities. In this case, we found that the case of the gap function $`\mathrm{\Delta }(𝐤)\mathrm{sin}(k_x+k_y)`$ has the largest $`\lambda _{\mathrm{Max}}`$. (Since the Hubbard model has SU(2) symmetry, the $`T_\mathrm{c}`$ does not depend on the direction of the d-vector, so that we may concentrate on the wave-number dependence.)
In Fig. 21, we show the maximum eigenvalue of the Éliashberg equation for the density varied over $`0.2n0.4`$ and $`t^{}`$ varied over $`0.0t^{}0.6`$ for $`U=2,4,6,8`$ with $`T=0.03`$. We can see that $`\lambda _{\mathrm{Max}}`$ is much smaller than unity even around $`t^{}0.5`$.
To probe the origin of this behavior, we take the case of $`U=2`$, $`t^{}=0.5`$, and $`n=0.3`$, which is also convenient becase the maximum value of $`\chi (𝐤,0)`$ (see below) in this case takes a similar value as in the square lattice with $`t^{}=0.5`$ (Fig.7(b)), which facilitates the comparison between the two cases. In Fig. 22(b), we show the momentum dependence of $`\chi (𝐤,0)`$. We can see that the peak is indeed located around the ferromagnetic point ($`\mathrm{\Gamma }`$).
In Fig. 23, we plot $`\lambda _{\mathrm{Max}}`$ along with the reciprocal of the peak value of $`\chi (𝐤,0)`$ as a function of $`T`$. $`\lambda _{\mathrm{Max}}`$ in this case is much smaller than that for the $`d`$-wave pairing in the SC lattice. This is again mainly due to the one-third $`V^{(2)}`$. $`\lambda _{\mathrm{Max}}`$ is smaller even when compared with that for the $`p`$-wave pairing in the square lattice with $`t^{}=0.5`$, although the maximum value of $`\chi (𝐤,0)`$, the frequency spread of $`\mathrm{Im}\chi (𝐤_{\mathrm{Max}},\omega )`$ (Fig. 22(a)), and $`U^2|G|^2`$ (see Fig. 24) all take similar values between the two. In fact, the reason for this 2D-3D discrepancy can again be traced back to the $`(a/L)^D`$ factor discussed in the previous subsection.
## 4 Discussions and Summary
To summarize, we have studied the possibility of spin-fluctuation mediated superconductivity in the single-band, repulsive Hubbard model for the $`d`$-wave channel on the (a) square lattice, (b) SC lattice, (c) BCC lattice, and for the $`p`$-wave channel on the (d) square lattice with large second-nearest neighbor hopping, (e) triangular lattice, and (vi) FCC lattice. We have shown that (i) $`d`$-wave instability mediated by antiferromagnetic spin fluctuations is stronger than $`p`$-wave instability mediated by ferromagnetic spin fluctuations both in 2D and 3D, and (ii)$`d`$-wave instability in 2D is much stronger than that in 3D.
We have given the physical reasons why the triplet p-pairing is unfavored as (i) the pairing interaction $`V^{(2)}`$ for triplet pairing is only $`1/3`$ of those for the singlet pairing, (ii) the self energy correction is large (i.e., quansi-particles are short-lived).
We have also traced back physical reasons why the superconducting instability in 3D is weaker than in 2D: if the momentum spread of $`\chi `$ (that determines $`V^{(2)}`$) in each momentum direction take similar values in 2D and 3D, it makes the higher-dimensional 3D disadvantageous because of a structure of the Éliashberg equation. Thus, our conclusion is that, so far as the single-band Hubbard model on ordinary lattices are concerned, $`d`$-wave pairing in 2D square lattice is the “best” situation for the spin fluctuation mediated superconductivity, where $`T_\mathrm{c}`$ can reach $`O(0.001W)`$ if we measure $`T_\mathrm{c}`$ in units of the band width, $`W`$. In this sense, the layer-type cuprates do seem to hit upon the right situation.
However, our conclusion has been obtained for the single-band Hubbard model. If we look more extensively at 3D superconductors, the heavy fermion system, in which the pairing is also thought to be meditated by spin fluctuations, $`T_\mathrm{c}`$ has a similar order of magnitude $`O(0.01W0.001W)`$, i.e., $`T_\mathrm{c}1`$ K with $`W=`$ a few hundred K.$`^{\text{?)}}`$ Since the present result indicates that $`T_\mathrm{c}`$ in the 3D Hubbard model should be much smaller than that for the 2D Hubbard model, we can envisage that the heavy fermion system exploits a more favorable situation than in the single-band Hubbard model where the frequency and/or momentum spreads in $`\chi (𝐤,\omega )`$ are larger than those in the 3D Hubbard model. In fact, the standard models, such as the multiband periodic Anderson model, employed to describe the heavy fermion system are more complicated than the single-band Hubbard model, and the frequency/momentum spreads in $`\chi `$ in such a model should be large enough to explain the superconductivity in the heavy fermion systems in the present context. It is an appealing future problem to explore how this can be so.
## 5 Acknowledgment
We would like to thank Professor K. Ueda and Dr H. Kontani for illuminating discussions. R.A. would like to thank Dr S. Koikegami for discussions on the FLEX. R.A. is supported by JSPS, while K.K. acknowledges a Grant-in-Aid for Scientific Research from the Ministry of Education of Japan. Numerical calculations were performed at the Supercomputer Center, ISSP, University of Tokyo.
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# Conductance fluctuations in diffusive rings: Berry phase effects and criteria for adiabaticity
## I Introduction and Overview
Since its discovery, the Berry phase has been a subject of continued interest. As this geometrical phase emerges from the very basic laws of quantum mechanics, it has implications for a broad range of physical systems. Even though the Berry phase has been observed in single-particle experiments, its manifestation in condensed matter systems is still under investigation. Some settings were proposed, where the Berry phase, resulting from the motion of a spin-carrying particle through an inhomogeneous magnetic field $`𝐁(𝐱)`$, can be observed in mesoscopic structures. The expected effects are measurable as persistent currents as well as in the magnetoconductance and the universal conductance fluctuations (UCFs). The first experiments reporting such effects were realized with semiconductor structures: the conductance was investigated in an InAs sample, where the Berry phase can emerge through the Rashba effect, in a very similar way as produced by an inhomogeneous field. Magnetoconductance measurements were performed where a ferromagnetic dot, placed slightly above a GaAs sample, produced an inhomogeneous field. Measurements on metallic systems also showed effects, which have been explained in terms of the Berry phase. Further experiments on metallic systems are in progress. An additional scenario was proposed, where domain walls of mesoscopic ferromagnets lead to a Berry phase.
During orbital motion in a magnetic field, a spin acquires a Berry phase in a similar way as a charge collects an Aharonov-Bohm phase. Thus, these two phases lead to similar implications for interference phenomena in mesoscopic samples. However, in the first case the phase originates from the change in local field direction, whereas in the second case it results from an enclosed magnetic flux. As these field properties can be varied individually, the interplay of the two phases yields a rich variety of behaviour. These quantum phases are distinguished by another important difference: while Aharonov-Bohm effects appear for arbitrarily small magnitudes $`B`$ of the magnetic field, Berry phase effects appear to their full extent only in the adiabatic limit, i.e. for large enough fields (specified below). The physical situation required for this limit to be satisfied can be pictured as a spin which must complete many precessions $`\omega _\mathrm{B}t_o/2\pi `$ around the local magnetic field, while it moves during a time $`t_o`$ through a region of size $`\mathrm{}_B`$ over which the direction of the field changes significantly. Here we have introduced the Bohr frequency $`\omega _\mathrm{B}=g\mu _\mathrm{B}B/2\mathrm{}`$, where $`g`$ is the Landé g-factor and $`\mu _\mathrm{B}`$ is the Bohr magneton. For ballistic motion as it occurs in clean semiconductors, one has $`v_\mathrm{F}t_o\mathrm{}_B`$ and there is general consensus about the criterion for adiabaticity, i.e. $`\omega _\mathrm{B}\mathrm{}_B/v_\mathrm{F}2\pi `$, with $`v_\mathrm{F}`$ being the Fermi velocity. However, for diffusive systems there were recently some discussions whether $`t_o`$ can be correctly set as the diffusion time $`t_d=\mathrm{}_B^2/D`$ or if one should replace it by the elastic scattering time $`\tau `$. The first criterion is more optimistic, in the sense that much lower field magnitudes are required to reach adiabaticity, as due to diffusive motion the electrons effectively move more slowly (compared to the ballistic motion) through the changing magnetic field and thus have more time to adjust their spins to the local field orientation. For magnetoconductance quantitative values for the required field magnitudes have been obtained. Solving the special case of a cylindrically symmetrical texture exactly, it was confirmed that the more favorable criterion is indeed sufficient. We remark that, if the ballistic criterion was appropriate for diffusive systems, the large fields required for adiabaticity would imply a strong curvature of the semiclassical trajectories (apart from the case of very large $`g`$ factors). This curvature in turn is in conflict with the approximation of the orbital motion by its zero-field value and therefore an approach beyond weak localization theory would be required for a self-consistent theory. At this point it should also be noted that Berry phase effects occur even if the adiabatic limit is not fully reached; there is no sharp cutoff where the Berry phase disappears completely. Thus, calculations without assuming adiabaticity are very desirable, as they can be used to study how the Berry phase effects gradually emerge while the magnetic field is increased from low to adiabatic strengths. The adiabatic limit can still be taken at the end of the calculation, so the formal appearance of the Berry phase and the associated dephasing can be identified.
Besides having a spin following the direction of an inhomogeneous external field, there is another scenario which produces a Berry phase: spin-orbit coupling. If an electron moves through an electrical field perpendicular to the ring plane, an effective magnetic field, which is produced in the rest frame of the electron, couples to the electron spin. As this effective field is in radial direction of the ring and perpendicular to the direction of motion, the field rotates while the electron moves around the ring and can therefore produce a Berry phase. By switching on, in addition, an external magnetic field, an arbitrary tilt angle of the total effective field can be realized and so this Berry phase can be tuned. For ballistic motion, the Berry phase manifests itself in precisely the same way as in the case with an inhomogeneous external magnetic field. However, for diffusive motion the situation becomes more complicated, as the change of the direction of motion of the electron due to a elastic scattering event abruptly changes the effective field direction. Now the picture of a spin, moving adiabatically through a slowly varying field, is no longer valid and needs to be modified. This leads to a new physical situation which has to be considered separately from the situation with inhomogeneous fields.
The outline of this paper is as follows. In Sec. II we study the conductance fluctuations $`\delta g^{(2)}`$ of quasi-1D diffusive rings in inhomogeneous magnetic fields. While $`\delta g^{(2)}`$ has already been calculated within the adiabatic approximation, i.e. for strong magnetic fields, the behavior outside the adiabatic limit and the influence of inhomogeneous fields on dephasing were not dicussed so far. We address these issues in the present work, starting in Sec. II A with a calculation of an exact expression for $`\delta g^{(2)}`$ (i.e. allowing arbitrarily small field magnitudes) for a special texture \[see Eq. (1)\] of the magnetic field. In this process we derive a new form of the diffuson differential equation, which includes inhomogeneous magnetic fields. We evaluate the adiabatic limit of the UCFs, $`\delta g_{\mathrm{ad}}^{(2)}`$, in Sec. II B and compare our results with those derived in previous work. Further, we investigate in Sec. II C the finite temperature behavior of the conductance fluctuations. In Sec. III the effects of the Berry phase on the UCFs and their dependence on magnetic field strengths are discussed in detail. We identify in Sec. III A a new effect of the Berry phase by showing that the amplitudes of the $`h/2e`$ Aharonov-Bohm oscillations depend directly on the value of the Berry phase. In particular, we find some magic tilt angles of the magnetic field, where these Aharonov-Bohm oscillations are completely suppressed. This effect provides a tool for experimental searches of the Berry phase. We use this observation to illustrate the gradually appearing effects of the Berry phase for increasing field strengths and thus give a direct demonstration of the onset of adiabaticity. Then, in Sec. III B, we give quantitative values of the fields strengths needed for reaching adiabaticity. We show that the criterion for adiabaticity is less stringent for diffusive than for ballistic motion. An exact evaluation of magnetoconductance $`\delta g_{SO}`$ and conductance fluctuations $`\delta g_{SO}^{(2)}`$ in the presence of spin-orbit coupling and homogeneous magnetic fields is given in Sec. IV. These results show how the amplitudes of the Aharonov-Bohm oscillations in $`\delta g_{SO}^{(2)}`$ depend non-monotoneously on the direction of an effective field, similarly as it is the case for inhomogeneous magnetic fields. In Sec. V A we show how frequency shift of the Aharonov-Bohm oscillations appear in $`\delta g`$ and $`\delta g^{(2)}`$ caused by the Berry phase. We then point out in Sec. V B that the Zeeman term can also produce frequency shifts even in the case of homogeneous fields. In Sec. V C we plot and discuss the exact expressions for $`\delta g`$ and $`\delta g^{(2)}`$ for inhomogeneous fields and for spin-orbit coupling as well as the corresponding power spectra. In three appendices we provide details of our calculations.
## II Conductance Fluctuations
As foundation for further discussions of Berry phase effects and adiabaticity, we will first calculate the conductance fluctuations $`\delta g^{(2)}`$ in the weak-localization regime. To motivate the analysis of the conductance fluctuations, we would like to emphasize the advantage of studying the UCFs instead of the magnetoconductance. The latter quantity has only contributions from the cooperon, which are suppressed by moderately large magnetic fields penetrating the ring arms. This suppression is in direct competition with the requirement of having large fields to satisfy adiabaticity. In contrast, the conductance fluctuations also have contributions from the diffuson, which is only sensitive to the difference of the two magnetic fields, for which the conductance correlator is considered. Therefore, if both fields are taken of similar magnitude, Aharonov-Bohm oscillations and Berry phase effects in the UCFs will still be visible at high magnetic fields where the adiabatic criterion is certainly satisfied.
### A Exact solution
We shall concentrate on rings with circumference $`L`$ and study the conductance-conductance correlator $`\delta g^{(2)}(𝐁,\stackrel{~}{}B)=g_𝐁g_{\stackrel{~}{}B}g_𝐁g_{\stackrel{~}{}B}`$, where we have two different magnetic fields $`𝐁`$ and $`\stackrel{~}{}B`$. We consider a special texture for which we obtain exact results (i.e. without making the adiabatic assumption of strong magnetic fields). We assume the magnetic fields to be applied in such a way that they wind $`f`$ times around the $`z`$-axis in one turn around the ring, with tilt angles $`\eta `$, $`\stackrel{~}{\eta }`$, see Fig. 1. The position along the direction of the ring is described by the coordinate $`x`$, varying from $`0`$ to $`L`$, so the special texture of the magnetic field is expressed as
$`𝐁`$ $`=`$ $`B𝐧`$ (1)
$`=`$ $`B(\mathrm{sin}\eta \mathrm{cos}(\frac{2\pi fx}{L}+\theta ),\mathrm{sin}\eta \mathrm{sin}(\frac{2\pi fx}{L}+\theta ),\mathrm{cos}\eta ),`$ (2)
and similarly for $`\stackrel{~}{}B`$. We have introduced $`\theta `$, so we can describe the textures with a field component radial to the ring, i.e. $`\theta =0`$, as well as textures with a field component tangential to the ring, i.e. $`\theta =\pi /2`$.
The starting point of our calculation is the conductance correlator derived in Ref. and given by
$`\delta g^{(2)}`$ $`=`$ $`\left({\displaystyle \frac{2e^2D}{hL^2}}\right)^2{\displaystyle }dϵdϵ^{}n^{}(ϵ)n^{}(ϵ^{})\{{\displaystyle \frac{1}{d}}\mathrm{Tr}\widehat{\chi }_\omega ^C\widehat{\chi }_\omega ^C`$ (4)
$`+2\mathrm{Re}\mathrm{Tr}\widehat{\chi }_\omega ^C\widehat{\chi }_\omega ^C+[\widehat{\chi }^C\widehat{\chi }^D]\},`$
where $`n^{}(ϵ)`$ is the derivative of the Fermi function and $`\mathrm{}\omega =ϵϵ^{}`$. The dimensionality of the system with respect to the diffusive motion is denoted by $`d`$, which describes the relation of the mean free path $`\mathrm{}`$ to the diffusion coefficient $`D`$, i.e. $`D=v_\mathrm{F}\mathrm{}/d`$. The propagators $`\widehat{\chi }^{C/D}`$ can be evaluated explicitly by using the operator equation Eq. (A10):
$$\widehat{\chi }^{C/D}=\frac{L^2}{(2\pi )^2D}\frac{1}{iw+\gamma ^{C/D}h^{C/D}}.$$
(5)
We have defined $`w=(L/2\pi L_T)^2(ϵϵ^{})/kT`$, with the thermal diffusion length $`L_T=\sqrt{D\mathrm{}\beta }.`$ The (non-hermitian) Hamiltonian is given by
$$h^{C/D}=\frac{L^2}{(2\pi )^2}\frac{^2}{x^2}+i\kappa 𝐧𝝈_1i\stackrel{~}{\kappa }\stackrel{~}{}n𝝈_2^{()},$$
(6)
where the star means complex conjugation in $`h^D`$ and where we have introduced an adiabaticity parameter
$$\kappa =\frac{\omega _\mathrm{B}}{D}\frac{L^2}{(2\pi )^2},$$
(7)
and equivalently for $`\stackrel{~}{\kappa }`$ and $`\omega _{\stackrel{~}{\mathrm{B}}}`$. We have inserted a phenomenological damping constant $`\gamma ^{C/D}=(L/2\pi L_{C/D})^2`$ expressed in terms of the magnetic dephasing length $`L_{C/D}`$:
$$\gamma ^{C/D}=\frac{L^2}{(2\pi )^2L_\phi ^2}+\frac{1}{3(4\pi )^2}\left(\frac{A|B_z\pm \stackrel{~}{B}_z|}{2\pi \varphi _0}\right)^2.$$
(8)
The first term of this damping constant incorporates the loss of phase due to inelastic scattering events. The second term takes into account magnetic flux penetration into the arms of the ring with a finite width $`a`$ and a surface area $`A=aL`$, while the height $`b`$ is assumed to be small compared to $`a`$. This field penetration leads to averaging over closed paths of different lengths, each of which collects a different Aharonov-Bohm phase, resulting finally in dephasing.
Next we define the basis in which we evaluate the Hamiltonian $`h^{C/D}`$. As done in Ref. for the cooperon propagator, we now introduce the operators
$$J^{C/D}=\frac{L}{2\pi i}\frac{}{x}+\frac{1}{2}f(\sigma _{1z}\pm \sigma _{2z}),$$
(9)
which commute with $`h^{C/D}`$.
We will now go to the basis of eigenvectors $`|j,\alpha \beta _{C/D}`$ of $`J^{C/D}`$. This basis is orthonormal with the following wave functions:
$$x,\alpha ^{}\beta ^{}|j,\alpha \beta _{C/D}=\frac{\delta _{\alpha ^{}\alpha }\delta _{\beta ^{}\beta }}{\sqrt{L}}\mathrm{exp}\left\{\frac{2\pi ix}{L}(j\frac{f}{2}\alpha \frac{f}{2}\beta )\right\}.$$
(10)
Because of the periodic boundary conditions in $`x`$, the eigenvalues $`j`$ of $`J^{C/D}`$ have to be integers. The matrix elements of $`h^{C/D}`$ in the basis $`\{|j,_{C/D},|j,_{C/D},|j,_{C/D},|j,_{C/D}\}`$ become:
$$_{C/D}j,\alpha \beta \left|h^{C/D}\right|j^{},\alpha ^{}\beta ^{}{}_{C/D}{}^{}=\delta _{jj^{}}(h_j^{C/D}+h_\sigma ^{C/D}),$$
(11)
where $`h_j^C`$ and $`h_j^D`$ are diagonal 4$`\times `$4 matrices with the entries $`\{(jf)^2,j^2,j^2,(j+f)^2\}`$, and $`\{j^2,(jf)^2,(j+f)^2,j^2\}`$, resp., and the $`\eta `$, $`\stackrel{~}{\eta }`$ dependent matrices are
$$h_\sigma ^{C/D}=\left(\begin{array}{cccc}i\kappa \mathrm{cos}\eta i\stackrel{~}{\kappa }\mathrm{cos}\stackrel{~}{\eta }& i\stackrel{~}{\kappa }e^{i\theta }\mathrm{sin}\stackrel{~}{\eta }& i\kappa e^{i\theta }\mathrm{sin}\eta & 0\\ i\stackrel{~}{\kappa }e^{\pm i\theta }\mathrm{sin}\stackrel{~}{\eta }& i\kappa \mathrm{cos}\eta +i\stackrel{~}{\kappa }\mathrm{cos}\stackrel{~}{\eta }& 0& i\kappa e^{i\theta }\mathrm{sin}\eta \\ i\kappa e^{i\theta }\mathrm{sin}\eta & 0& i\kappa \mathrm{cos}\eta i\stackrel{~}{\kappa }\mathrm{cos}\stackrel{~}{\eta }& i\stackrel{~}{\kappa }e^{i\theta }\mathrm{sin}\stackrel{~}{\eta }\\ 0& i\kappa e^{i\theta }\mathrm{sin}\eta & i\stackrel{~}{\kappa }e^{\pm i\theta }\mathrm{sin}\stackrel{~}{\eta }& i\kappa \mathrm{cos}\eta +i\stackrel{~}{\kappa }\mathrm{cos}\stackrel{~}{\eta }\end{array}\right).$$
(12)
To take the Aharonov-Bohm flux into account, we replace $`jm=j(\varphi /\varphi _0\pm \stackrel{~}{\varphi }/\varphi _0)`$, where $`\varphi ,\stackrel{~}{\varphi }`$ are the fluxes of the fields $`𝐁,\stackrel{~}{}B`$ through the ring and $`\varphi _0=h/e`$ is the magnetic flux quantum. Now it is straightforward to evaluate the exact conductance fluctuations $`\delta g^{(2)}`$ by calculating the propagators by matrix inversion and inserting the result into Eq. (4). This can be done with the help of the computer program Mathematica, which however leads to lengthy expressions which we will not reproduce here. We merely point out that the phase factors in $`\theta `$ cancel each other in $`\delta g^{(2)}`$ and $`\delta g`$.
### B Adiabatic Approximation
To evaluate the adiabatic limit, we shall consider the regime of large magnetic fields with $`B`$ and $`\stackrel{~}{B}`$ of similar magnitude. If we define $`\mathrm{\Delta }\kappa =\stackrel{~}{\kappa }\kappa `$, this adiabatic regime is described by
$$\kappa 1\text{and}\kappa |\mathrm{\Delta }\kappa |.$$
(13)
The exact propagators $`\chi ^{C/D}`$ turn out to be rational functions which are of order two in $`\kappa `$ in both numerator and denominator. Now we will keep only the terms of highest order in $`\kappa `$; terms with large $`j`$ can be neglected as the sum over $`j`$ converges rapidly. This leads us to the UCFs in the adiabatic regime:
$`\delta g_{\mathrm{ad}}^{(2)}`$ $`=`$ $`\left({\displaystyle \frac{e^2}{h}}\right)^2{\displaystyle \frac{1}{4\pi ^4}}{\displaystyle 𝑑ϵ𝑑ϵ^{}n^{}(ϵ)n^{}(ϵ^{})\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}\underset{\alpha =\pm 1}{}\left(G_{\alpha ,C}^{\mathrm{ad}}+G_{\alpha ,D}^{\mathrm{ad}}\right)}`$ (14)
$`G_{\alpha ,C/D}^{\mathrm{ad}}`$ $`=`$ $`{\displaystyle \frac{1}{d}}\{(w\alpha \mathrm{\Delta }\kappa )^2+\delta _\alpha ^{C/D}(j)^2+P\}\times \{[(w\alpha \mathrm{\Delta }\kappa )^2+\delta _\alpha ^{C/D}(j)^2P]`$ (16)
$`[(w+\alpha \mathrm{\Delta }\kappa )^2+\delta _\alpha ^{C/D}(j)^2P]+4P[w^2+f^2m^2(\mathrm{cos}\eta \pm \mathrm{cos}\stackrel{~}{\eta })^2]\}^1`$
$`+`$ $`2\mathrm{R}\mathrm{e}\mathbf{[}\{[iwi\alpha \mathrm{\Delta }\kappa +\delta _\alpha ^{C/D}(j)]^2+P\}`$ (18)
$`\times \{[iwi\alpha \mathrm{\Delta }\kappa +\delta _\alpha ^{C/D}(j)][iw+i\alpha \mathrm{\Delta }\kappa +\delta _\alpha ^{C/D}(j)]P\}^2\mathbf{]},`$
where
$`P`$ $`=`$ $`{\displaystyle \frac{f^4}{4}}\mathrm{sin}^2\eta \mathrm{sin}^2\stackrel{~}{\eta },`$ (19)
$`\delta _\alpha ^{C/D}(j)`$ $`=`$ $`\stackrel{~}{\gamma }_{\eta ,\stackrel{~}{\eta }}^{C/D}+\left(m{\displaystyle \frac{f}{2}}\alpha \mathrm{cos}\eta {\displaystyle \frac{f}{2}}\alpha \mathrm{cos}\stackrel{~}{\eta }\right)^2,`$ (20)
with
$`\stackrel{~}{\gamma }_{\eta ,\stackrel{~}{\eta }}^{C/D}`$ $`=`$ $`\gamma ^{C/D}+{\displaystyle \frac{f^2}{4}}\mathrm{sin}^2\eta +{\displaystyle \frac{f^2}{4}}\mathrm{sin}^2\stackrel{~}{\eta }.`$ (21)
The sum over $`\alpha `$ has been introduced here artificially to facilitate the following interpretation. As it is also seen in Ref. for the case of the magnetoconductance $`\delta g`$, the terms $`f^2(\mathrm{sin}^2\eta +\mathrm{sin}^2\stackrel{~}{\eta })/4`$ in Eq. (21) act as additional dephasing sources and are here absorbed in the phenomenological dephasing parameter $`\stackrel{~}{\gamma }_{\eta ,\stackrel{~}{\eta }}^{C/D}`$. However, in Eq. (18) there are further $`\eta `$, $`\stackrel{~}{\eta }`$-dependent terms $`P`$, which cannot be formally absorbed in $`\stackrel{~}{\gamma }_{\eta ,\stackrel{~}{\eta }}^{C/D}`$. $`P`$ reduces the effect of the additional dephasing terms in Eq. (21), as we can see by the following numerical evaluation. We consider equal fields $`𝐁=\stackrel{~}{}B`$ and low temperatures, thus $`\mathrm{\Delta }\kappa ,\omega =0`$, and assume $`\eta `$, $`\stackrel{~}{\eta }`$ to be close to $`\pi /2`$. Then we estimate the amplitude of the Aharonov-Bohm oscillations by taking the difference between the values of $`G_{\alpha ,C}^{\mathrm{ad}}`$ \[Eq. (18)\] for the two phases $`m=0`$ and $`m=\pm 1/2`$ (i.e. we are considering only the main contributions in the sum over $`j`$ \[Eq. (14)\]). We then see by numerical evaluation that the oscillations are suppressed if we set $`P=0`$ instead of using Eq. (19), thus $`P`$ indeed reduces dephasing.
We can compare now with previous calculations where the UCFs $`\delta g_{\mathrm{LSG}}^{(2)}`$ have been derived for arbitrary textures and adiabatic evolution of the spin. These results can be recovered from Eq. (14) by the replacement $`\stackrel{~}{\gamma }_{\eta ,\stackrel{~}{\eta }}^{C/D}\gamma _{\mathrm{LSG}}^{C/D}`$ and $`P0`$. The dephasing terms due to inhomogeneous fields coupling to the spin \[see Eqs. (19), (21)\] were not explicitly given in Ref. ; to account for such dephasing these terms must be included in the phenomenological parameter $`\gamma _{\mathrm{LSG}}^{C/D}`$, and thus $`\gamma _{\mathrm{LSG}}^{C/D}\gamma ^{C/D}`$ and $`\gamma _{\mathrm{LSG}}^{C/D}\stackrel{~}{\gamma }_{\eta ,\stackrel{~}{\eta }}^{C/D}`$ in general.
We also recognize a strong simplification in the special case where one field is homogeneous, $`\eta =0`$, i.e., $`P`$ vanishes. Thus the comparison of $`\delta g_{\mathrm{ad}}^{(2)}`$ with the solution for arbitrary textures $`\delta g_{\mathrm{LSG}}^{(2)}`$ yields the simple relation $`\gamma _{\mathrm{LSG}}^{C/D}=\stackrel{~}{\gamma }_{0,\stackrel{~}{\eta }}^{C/D}`$. Finally we note that in this case the dephasing due to the orientational inhomogenity of $`\stackrel{~}{}B`$ measured by the winding $`f`$ grows like $`f^2\mathrm{sin}^2\stackrel{~}{\eta }`$ \[cf. Eq. (21)\].
### C Finite Temperatures
Now we consider the effects of finite temperatures $`T>0`$ on the UCFs $`\delta g_{\mathrm{ad}}^{(2)}`$ in the adiabatic regime. In the case of $`\eta =0`$, i.e. $`P=0`$, the factors containing $`\delta _\alpha ^{C/D}(j)`$ in Eq. (18) cancel, so we obtain
$`G_{\alpha ,C/D}^{\mathrm{ad}}|_{\eta =0}`$ $`=`$ $`{\displaystyle \frac{1}{d}}\left\{\left(w+\alpha \mathrm{\Delta }\kappa \right)^2+\delta _\alpha ^{C/D}(j)^2\right\}^1`$ (23)
$`+2\mathrm{R}\mathrm{e}\left\{iw+i\alpha \mathrm{\Delta }\kappa +\delta _\alpha ^{C/D}(j)\right\}^2.`$
This strong simplification allows us to evaluate the integrals over $`ϵ`$ and $`ϵ^{}`$ in Eq. (14) explicitly by using standard Matsubara techniques, as described in App. B, and we obtain for the UCFs $`\delta g_{\mathrm{ad}}^{(2)}=\delta g_{\mathrm{ad},C}^{(2)}+\delta g_{\mathrm{ad},D}^{(2)}`$,
$`\delta g_{\mathrm{ad},C/D}^{(2)}|_{\eta =0}`$ $`=\left({\displaystyle \frac{e^2}{h}}\right)^2{\displaystyle \frac{1}{8\pi ^6}}\left({\displaystyle \frac{L^2}{L_T^2}}\right)^2\mathrm{Re}{\displaystyle \underset{\alpha =\pm 1}{}}{\displaystyle \underset{j,n,m}{}^{}}`$ (26)
$`\{{\displaystyle \frac{1}{d\delta _\alpha ^{C/D}(j)\left[\frac{L^2}{4\pi L_T^2}(m+n)+\delta _\alpha ^{C/D}(j)i\alpha \mathrm{\Delta }\kappa \right]^3}}`$
$`+{\displaystyle \frac{6}{\left[\frac{L^2}{4\pi L_T^2}(m+n)+\delta _\alpha ^{C/D}(j)i\alpha \mathrm{\Delta }\kappa \right]^4}}\}.`$
Here $`n`$ and $`m`$ are odd, positive integers. For plotting, it is advantageous to calculate the sum in Eq. (26) analytically, which gives an expression containing Psi-functions.
We can now obtain a qualitative criterion when the thermal dephasing effects can be ignored. If we ignore thermal effects, i.e. assume low temperatures, we can simplify our calculation leading to Eq. (26) by replacing $`n^{}(ϵ)`$ by a delta function $`\delta (ϵ)`$ in Eq. (14). This yields for $`\eta =0`$ the same result as applying Poisson’s summation formula to Eq. (26) in order to replace the summations over $`n`$ and $`m`$ by integrations. We are only allowed to perform this step if the summand varies slowly in $`n`$$`m`$, which is the case for $`L_T^22\pi L_{C/D}^2`$. From a physical point of view, this is an evident requirement: the smearing of the conductance fluctuations due to nonzero temperatures, described by the thermal diffusion length $`L_T`$, can only be neglected if the dephasing lengths related to inelastic scattering or penetrating magnetic fields are much shorter than $`L_T`$.
In App. C we evaluate the dephasing behavior of the UCFs $`\delta g_{\mathrm{hom}}^{(2)}`$ for homogeneous fields and finite temperatures. Then we confirm the result of Ref. \[Eq. (4)\] and show that our calculation in the homogeneous limit indeed reproduces known results.
## III Berry phase and Adiabaticity
### A Magic Angles—Qualitative Criterion for Adiabaticity
We now consider the qualitative effects of the Berry phase on the conductance fluctuations $`\delta g^{(2)}`$. They emerge from the Berry phase in $`\delta _\alpha ^{C/D}(j)`$ in the adiabatic solution \[Eq. (14)\] and lead to vanishing Aharonov-Bohm oscillations at special “magic” tilt angles of the magnetic fields. This effect has some similarities with the phenomenon of beating, where the superposition of two oscillations with different but fixed frequencies leads to a periodic vanishing of the envelope. However, in our case we have two frequencies which will change when the perpendicular field $`B_z`$ is increased, since then the Berry phase is altered, too. Thus a suppression of the Aharonov-Bohm oscillations can only be observed at two special tilt angles of the magnetic field, i.e. the Berry phase has a highly non-periodic effect on the envelope of these oscillations as a function of $`B_z`$.
From now on we shall only study the experimentally realizable field texture with one winding, $`f=1`$. The other configurations with $`f>1`$ are solely of academic interest. To illustrate expected experimental results, we will use some material parameters recently determined. The sample Au-1 given in Table I of Ref. has the values $`D=9\times 10^3\mathrm{m}^2\mathrm{s}^1`$ and $`L_\phi =\sqrt{D\tau _\phi }=5.54\mu \mathrm{m}`$. We assume a ring with diameter of $`4\mu \mathrm{m}`$, so $`L=12.6\mu \mathrm{m}`$, and an arm width $`a=60\mathrm{nm}`$, which lies well within present-day experimental reach. Finally we assume low temperatures, i.e. $`L_TL,L_\phi `$, so we can ignore the dephasing due to thermal fluctuations.
Now we shall consider two equal fields, so no phase terms appear in the diffuson contribution $`\delta g_D^{(2)}`$. The cooperon contribution $`\delta g_C^{(2)}`$ is $`h/2e`$ periodic in the magnetic flux, as a shift of $`m=\varphi /\varphi _0+\stackrel{~}{\varphi }/\varphi _0+j=2\varphi /\varphi _0+j`$ by 1 is absorbed in the sum over $`j`$ in Eq. (14). For the next argument we take the dephasing due to inhomogeneous fields only phenomenologically into account, i.e. we use the result $`\delta g_{\mathrm{LSG}}^{(2)}`$ from Ref. or equivalently set $`P=0`$ \[Eqs. (14, 23)\], so the factors containing $`\delta _\alpha ^{C/D}(j)`$ cancel in Eq. (18). If the tilt angle $`\eta `$ is such that $`\mathrm{cos}\eta =1/4`$, the phase dependent term in $`\delta _\alpha ^{C/D}(j)`$ \[Eq. (20)\] becomes $`m\alpha /4`$. One sees that in this special case shifting $`m`$ by $`1/2`$ does not affect the value of $`\delta g_C^{(2)}`$, as it leads solely to an exchange $`\alpha \alpha `$. The very same argument applies to $`\mathrm{cos}\eta =3/4`$. Thus, for these magic angles $`\eta `$, where $`\mathrm{cos}\eta =1/4,3/4`$, the UCFs $`\delta g^{(2)}`$ are $`h/4e`$ periodic and therefore their power spectrum shows a vanishing $`h/2e`$ amplitude. If we take the exact solution in the adiabatic regime $`\delta g_{\mathrm{ad}}^{(2)}`$ instead of $`\delta g_{\mathrm{LSG}}^{(2)}`$, the magic angles are still present, but at shifted values. The angle at $`\mathrm{cos}\eta =3/4`$ is nearly unaffected, as $`P0.05`$ is very small at this angle. The suppression of the Aharonov-Bohm oscillations is illustrated in Fig. 2 (see also Sec. V C and Fig. 9) by plotting the $`h/2e`$ amplitude of the exact solution $`\delta g^{(2)}`$ with varying tilt angle $`\eta `$ and for different radial field components. As one can readily see from Fig. 2, the effect described here is fully developed for $`B200\mathrm{G}`$. For smaller fields, the $`h/2e`$ amplitude does not completely vanish at the magic angles, as adiabaticity is not yet reached. It should be noted that even if the adiabatic regime is not fully reached, an effect of the Berry phase is still visible as a distinct non-monotonic behavior of the UCFs $`\delta g^{(2)}`$ as a function of the tilt angle $`\eta `$, unlike the UCFs for a configuration with a homogeneous field texture (also shown in Fig. 2).
Another interesting situation arises for $`B\stackrel{~}{B}`$. Now, phase effects from the diffuson contribution to $`\delta g^{(2)}`$ emerge and remain present even for large fields, since the dephasing due to flux penetrating the arms of the ring depends only on the difference of the fields and not on the sum as for the cooperon contribution, see Eq. (8). For illustration, we consider the configuration where $`𝐁`$ is homogeneous with $`\eta =0`$. The other field $`\stackrel{~}{}B`$ is assumed to have a radial component so that for a tilt angle $`\stackrel{~}{\eta }=\pi /3`$ the magnitudes of both fields are equal, i.e. $`\stackrel{~}{B}_{||}=(\sqrt{3}/2)B_z`$. In the adiabatic approximation $`\delta g_{\mathrm{ad}}^{(2)}`$ \[Eq. (14)\] $`P`$ vanishes, yielding the simple relation Eq. (21) between the dephasing due to the inhomogeneous field textures and $`\gamma ^{C/D}`$: the effective dephasing will be increased by $`3/16`$ at the most interesting angle, $`\stackrel{~}{\eta }=\pi /3`$, in the situation considered here. The contribution of the penetrating fields to $`\gamma ^{C/D}`$ will be three times larger for the cooperon than for the diffuson, as can be seen from Eq. (8). Varying $`\stackrel{~}{B}_z`$ changes the Aharonov-Bohm phase $`\stackrel{~}{\varphi }/\varphi _0`$, while $`\varphi /\varphi _0=\mathrm{const}.`$, leading to $`h/e`$ oscillations. At $`\stackrel{~}{B}_z=B_z/2`$ two features are worth mentioning. First, the magnitudes of both fields become equal, therefore $`\mathrm{\Delta }\kappa `$ vanishes and so the second part of the criterion in Eq. (13) is fulfilled and we can use the adiabatic approximation $`\delta g_{\mathrm{ad}}^{(2)}`$ \[Eq. (14)\]. Second, we have $`\mathrm{cos}\stackrel{~}{\eta }=1/2`$, so the phase dependent terms $`m\alpha /4`$ arise in $`\delta _\alpha ^{C/D}(j)`$, as can be seen from Eq. (20). With the same argument as above, the UCFs $`\delta g^{(2)}`$ become $`h/2e`$ periodic at this magic angle $`\pi /3`$, so the $`h/e`$ amplitude vanishes in the power spectrum. We note that, in the adiabatic regime, this magic angle is exact, since for the configuration $`\eta =0`$ we have $`\delta g_{\mathrm{ad}}^{(2)}=\delta g_{\mathrm{LSG}}^{(2)}`$. This is shown in Fig. 3, again as a function of the tilt angle $`\stackrel{~}{\eta }=\mathrm{cot}(\stackrel{~}{B}_z/B_z)`$, see also Sec. V C and Fig. 10.
### B Quantitative Criterion for Adiabaticity
In order to obtain a quantitative criterion for adiabaticity, we numerically compare the exact solution of the conductance fluctuations $`\delta g^{(2)}`$ with the adiabatic approximation $`\delta g_{\mathrm{ad}}^{(2)}`$ \[Eq. (14)\]. We take equal magnitudes for both fields, i.e. $`B=\stackrel{~}{B}`$. We search for a minimal $`\kappa _{\mathrm{min}}`$ so that the relative difference $`\left|\delta g^{(2)}\delta g_{\mathrm{ad}}^{(2)}\right|/\delta g^{(2)}`$ is below a certain value. This is done with a bisection algorithm (in $`\kappa `$) and by sampling over the parameter subspace $`[0,\pi /2]^2\times [0,1]^2\times [\frac{1}{100},10]^2\{(\eta ,\stackrel{~}{\eta },\varphi /\varphi _0,\stackrel{~}{\varphi }/\varphi _0,\gamma ^C,\gamma ^D)\}`$ with a grid resolution of 10 intersections in the first four dimensions. A finer resolution has been chosen for $`\gamma ^{C/D}`$. As can be seen from Fig. 4, for $`0.01\gamma ^D1,\gamma ^D\gamma ^C`$ and a field strength such that $`\kappa 3`$, the numerical values for $`\delta g^{(2)}`$ and $`\delta g_{\mathrm{ad}}^{(2)}`$ are already within five percent of each other.
However, as we are interested in the Aharonov-Bohm oscillations rather than in the absolute value of the UCFs $`\delta g^{(2)}`$, we now use a different method of comparison: We consider the oscillations in the conductance fluctuations resulting from different Aharonov-Bohm fluxes through the ring. As a measure for accuracy we take the relative error of these amplitudes, i.e.
$`\mathrm{\Delta }(\kappa ,\gamma ^C,\gamma ^D,\eta ,\stackrel{~}{\eta })`$ (27)
$`=`$ $`{\displaystyle \frac{\underset{\varphi ,\stackrel{~}{\varphi }}{\mathrm{max}}|(\delta g^{(2)}\delta g^{(2)}|_{\varphi =\stackrel{~}{\varphi }=0})(\delta g_{\mathrm{ad}}^{(2)}\delta g_{\mathrm{ad}}^{(2)}|_{\varphi =\stackrel{~}{\varphi }=0})|}{\underset{\varphi ,\stackrel{~}{\varphi }}{\mathrm{max}}|\delta g^{(2)}\delta g^{(2)}|_{\varphi =\stackrel{~}{\varphi }=0}|_{\eta ,\stackrel{~}{\eta }=0}}}.`$ (28)
Again we search for a minimal $`\kappa _{\mathrm{min}}`$ so that $`\mathrm{\Delta }`$ is bounded from above by a certain percentage over the whole parameter subspace. We notice from the results shown in Fig. 5 that in the regime with only moderate damping $`\gamma ^C=\gamma ^D=0.1`$, adiabaticity is already reached at $`\kappa 2`$. If we put this in the context of the experimental parameters given in the beginning of Sec. III A, we expect adiabaticity to be fully reached at magnetic fields of magnitude larger than $`500\mathrm{G}`$. By comparing this value with Fig. 2, we note that the qualitative effect of the Berry phase can already be seen for fields which are an order of a magnitude smaller, i.e. for $`B,\stackrel{~}{B}50\mathrm{G}`$.
We now discuss the effects of different parameters on $`\kappa `$ and on the minimal magnetic fields required to reach adiabaticity, thus indicating favorable experimental setups. If we consider rings of increasing circumference $`L`$, we can see from Eq. (7) that the minimal magnetic field strength needed decreases as $`B_{\mathrm{ad}}L^2`$. However, to observe the Berry phase, dephasing must not be too strong, so the condition $`L2L_{C/D}`$ should still be met. We note that for two equal fields, the first term of $`\gamma ^CL_C^2`$ in Eq. (8) depends on $`L^2`$, which restrains us from taking $`L>2L_\phi `$, whereas the second one depends for $`B=B_{\mathrm{ad}}`$ on $`L^2`$. So not only the high magnetic fields needed for adiabaticity, but also the small arm widths $`a`$ required to minimize strong dephasing due to the penetrating flux, disfavors experimental setups with very small $`L`$.
Introducing more impurities and thus decreasing the diffusion coefficient $`D`$ leads to slower motion of the electrons around the ring, giving their spins more time to adjust to the local magnetic texture. Thus, the field strengths required for adiabaticity to occur decrease as $`B_{\mathrm{ad}}D`$, which can be seen from Eq. (7). However, such slow diffusion also leads to shorter dephasing lengths $`L_T,L_\phi D^{1/2}`$; assuming that $`\tau _\phi `$ remains constant. To avoid such an additional dephasing, i.e. leaving $`\gamma ^{C/D}`$ unaffected, the sample size must also be decreased as $`LD^{1/2}`$. Thus, because of $`\kappa D^1L^2`$, no net decrease of the required fields for adiabaticity can be gained by decreasing the diffusion coefficient.
## IV Exact calculations with Spin-Orbit Interaction in Diffusive Limit
We turn now to the discussion of Berry phases induced by spin-orbit interaction. Instead of considering an inhomogeneous field, we use here an effective (non-hermitian) Hamiltonian
$`h_{SO}^{C/D}`$ $`=`$ $`{\displaystyle \frac{L}{(2\pi )^2}}{\displaystyle \frac{^2}{x^2}}+i\kappa \sigma _{1z}i\stackrel{~}{\kappa }\sigma _{2z}`$ (30)
$`+i{\displaystyle \frac{\alpha }{\mathrm{}^2}}{\displaystyle \frac{L^2}{D(2\pi )^2}}(𝐞_z\times 𝝈^{()})𝐩,`$
with spin-orbit interaction, using a coupling constant $`\alpha `$ as defined in Ref. , and with a Zeeman term from an external magnetic field, which is perpendicular to the ring plane. One arrives at this Hamiltonian by starting from the Feynman path integral representation of the transition amplitude with spin-orbit coupling, as it is given in Ref. . One can then formally decouple orbital and spin motion, and following the steps given in App. A of Ref. , one arrives at the effective Schrödinger equation for the cooperon propagator with the Hamiltonian $`h_{SO}^C`$. The equation with $`h_{SO}^D`$ for the diffuson, which will be required in Sec. IV B, can be obtained by applying the techniques explained in App. A.
Note that in Eq. (30) the momentum operator is still in the Cartesian coordinate system. Now we adopt a polar coordinate system, with $`(x^{},y^{})=(r\mathrm{cos}\frac{2\pi x}{L},r\mathrm{sin}\frac{2\pi x}{L})`$ and $`(_x^{},_y^{})=(\frac{1}{2}\{\mathrm{sin}\frac{2\pi x}{L},_x\},\frac{1}{2}\{\mathrm{cos}\frac{2\pi x}{L},_x\})`$, where $`x`$ denotes the position along the ring and runs from $`0`$ to $`L`$. The curly braces denote the anticommutator, which ensures the hermiticity of the momentum operator. We now have
$`h_{SO}^{C/D}`$ $`=`$ $`{\displaystyle \frac{L^2}{(2\pi )^2}}{\displaystyle \frac{^2}{x^2}}+i\kappa \sigma _{1z}i\stackrel{~}{\kappa }\sigma _{2z}`$ (33)
$`+{\displaystyle \frac{\alpha }{\mathrm{}}}{\displaystyle \frac{L^2}{D(2\pi )^2}}{\displaystyle \frac{1}{2}}\{\sigma _{1x}\mathrm{cos}\frac{2\pi x}{L}+\sigma _{1y}\mathrm{sin}\frac{2\pi x}{L}`$
$`\sigma _{2x}\mathrm{cos}\frac{2\pi x}{L}\sigma _{2y}\mathrm{sin}\frac{2\pi x}{L},{\displaystyle \frac{}{x}}\}.`$
To diagonalize the Hamiltonian, we follow the ideas used above and use the operators defined in Eq. (9), but now with $`f=\stackrel{~}{f}=1`$:
$$J^{C/D}:=\frac{L}{2\pi i}\frac{}{x}+\frac{1}{2}\sigma _{1z}\pm \frac{1}{2}\sigma _{2z},$$
(34)
which commute with the Hamiltonians $`h_{SO}^{C/D}`$, as can be seen using $`[\{n(x),_x\},_x]=\{n^{}(x),_x\}`$. We can now calculate the matrix elements of $`h_{SO}^{C/D}`$ in the basis defined in Eq. (10), with $`f=\stackrel{~}{f}=1`$, as
$$j,\alpha \beta \left|h_{SO}^C\right|j^{},\alpha ^{}\beta ^{}=\delta _{jj^{}}\left(\begin{array}{cccc}(j1)^2+i\kappa i\stackrel{~}{\kappa }& iS\left(j\frac{1}{2}\right)& iS\left(j+\frac{1}{2}\right)& \mathrm{\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}0}\\ iS\left(j\frac{1}{2}\right)& j^2+i\kappa +i\stackrel{~}{\kappa }& \mathrm{\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}0}& iS\left(j\frac{1}{2}\right)\\ iS\left(j+\frac{1}{2}\right)& \mathrm{\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}0}& j^2i\kappa i\stackrel{~}{\kappa }& iS\left(j+\frac{1}{2}\right)\\ \mathrm{\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}0}& iS\left(j\frac{1}{2}\right)& iS\left(j+\frac{1}{2}\right)& (j+1)^2i\kappa +i\stackrel{~}{\kappa }\end{array}\right),$$
(35)
and
$$j,\alpha \beta \left|h_{SO}^D\right|j^{},\alpha ^{}\beta ^{}=\delta _{jj^{}}\left(\begin{array}{cccc}j^2+i\kappa i\stackrel{~}{\kappa }& iS\left(j\frac{1}{2}\right)& iS\left(j\frac{1}{2}\right)& \mathrm{\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}0}\\ iS\left(j\frac{1}{2}\right)& (j1)^2+i\kappa +i\stackrel{~}{\kappa }& \mathrm{\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}0}& iS\left(j+\frac{1}{2}\right)\\ iS\left(j\frac{1}{2}\right)& \mathrm{\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}0}& (j+1)^2i\kappa i\stackrel{~}{\kappa }& iS\left(j+\frac{1}{2}\right)\\ \mathrm{\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}0}& iS\left(j+\frac{1}{2}\right)& iS\left(j+\frac{1}{2}\right)& j^2i\kappa +i\stackrel{~}{\kappa }\end{array}\right).$$
(36)
In Eqs. (35) and (36), we have introduced a dimensionless spin-orbit coupling parameter
$$S=\frac{\alpha }{\mathrm{}D}\frac{L}{2\pi }.$$
(37)
By comparing Eqs. (7) and (37), we note that while $`\kappa `$ is quadratic in $`L`$, the parameter $`S`$ is only linearly dependent on $`L`$. If we define an effective field angle for diffusive motion with spin-orbit coupling
$$\mathrm{tan}\eta _{SO}=S/\kappa ,$$
(38)
and anticipate the Berry phase to be of the form $`\mathrm{\Phi }^g=\mathrm{cos}\eta `$, we obtain for $`S\kappa `$ the dependency $`\mathrm{\Phi }^g\kappa /SL`$. Thus the phase can now be enhanced by increasing the size of the ring. However, the phase cannot be increased arbitrarily; for large $`L`$, the assumption $`S\kappa `$ becomes invalid.
### A Magnetoconductance
We shall now calculate the magnetoconductance with the formula from Ref.
$$\delta g_{SO}=\frac{e^2}{\pi \mathrm{}}\frac{L}{(2\pi )^2}\underset{\alpha ,\beta =\pm 1}{}x,\alpha ,\beta \left|\frac{1}{\gamma h_{SO}^C}\right|x,\beta ,\alpha .$$
(39)
With Eq. (35), we obtain the magnetoconductance
$`\delta g_{SO}`$ $`=`$ $`{\displaystyle \frac{e^2}{\pi \mathrm{}}}{\displaystyle \underset{j=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{2\left[4\kappa ^2+\left(m^2+\gamma \right)^2\right]\left(m^2+\gamma +1\right)+S^2\left[8m^4+2m^2\left(4\gamma 1\right)+2\gamma +1\right]}{\left[4\kappa ^2+\left(m^2+\gamma \right)^2\right]\left[m^4+2m^2\left(\gamma 1\right)+\left(\gamma +1\right)^2\right]+S^2\left(m^2+\gamma \right)\left[4m^4+m^2\left(4\gamma 3\right)+\gamma +1\right]}},`$ (40)
where $`m=j2\varphi /\varphi _0`$ contains the Aharonov-Bohm flux. In Sec. V we will see that in the “adiabatic” limit $`\kappa ,S1`$ the magnetoconductance $`\delta g_{SO}`$ will show some similar properties as for inhomogeneous fields, in particular a peak-splitting in the power spectrum, see Fig. 11.
### B Conductance Fluctuations
We turn now to a discussion of the recent experiment by Morpurgo et al. by specifying the parameters of the effective Hamiltonian $`h_{SO}^{C/D}`$, as given in Eqs. (30), (35), and (36). In Ref. , conductance measurements were performed on an InAs ring, with nearly ballistic transport. For the parameters given, $`\alpha =5.5\times 10^{10}\mathrm{eV}\mathrm{cm}`$, $`L=6.6\mu \mathrm{m}`$, $`v_\mathrm{F}=9.8\times 10^7\mathrm{cm}/\mathrm{s}`$, $`\mathrm{}=1.0\mu \mathrm{m}`$, and $`D=v_\mathrm{F}\mathrm{}/2=4.9\times 10^3\mathrm{cm}^2/\mathrm{s}`$, we calculate with Eq. (37) a numerical value of $`S1/50`$. Compared to this, the strength of the Zeeman term $`\kappa 1/2`$ (with $`|g|=15`$) is much larger. Within the diffusive approximation, this spin-orbit coupling $`S\kappa `$ gives only a negligible contribution to the effective Hamiltonian $`h^{C/D}`$ \[Eq. (30)\] and thus does not produce any Berry phase effects. This very same finding has also been obtained in Ref. , based on a slightly different reasoning. Still, we show in Sec. V that a spin-splitting produced by spin-orbit interaction can be obtained in the “adiabatic regime” $`\kappa ,S1`$, which, however, is in the opposite limit to the one reported in Ref. . So although we cannot give a quantitative explanation of the experiment here, we can offer a qualitative interpretation, see Fig. 13. Further, there is an uncertainty in the spin-orbit coupling parameter $`\alpha `$ in InAs, as it was recently pointed out, and more experiments might be needed to clarify this issue.
To this end we calculate the exact, i.e. without assuming any form of adiabaticity, expression for the conductance fluctuations $`\delta g_{SO}^{(2)}`$ in the presence of spin-orbit interaction. With the block-diagonalization of the Hamiltonian $`h_{SO}^{C/D}`$ \[Eqs. (35), (36)\] we obtain the propagators required in the formula for the conductance correlator \[Eq. (4)\]. We use Mathematica to obtain an explicit algebraic expression for $`\delta g_{SO}^{(2)}`$ (which is lengthy and thus not reproduced here) and plot it in Fig. 6 (see also Figs. 12 and 13). From this plot we deduce that in a configuration with spin-orbit coupling, the Aharonov-Bohm oscillations vanish for certain values of $`S`$ and $`\kappa `$. It is remarkable that this happens, for $`S2`$, at the fixed ratios $`\kappa /S=0.2`$ and $`0.5`$, which can be ascribed again to some effective magic angles. Thus we see that Berry phase-like effects occur in $`\delta g_{SO}^{(2)}`$ as the amplitudes of the Aharonov-Bohm oscillations become dependent on $`\kappa /S`$. This resembles the case for inhomogeneous fields, where the amplitudes of the Aharonov-Bohm oscillations became dependent on the tilt angle $`\eta `$ of the magnetic field due to the Berry phase, as it was shown in Sec. III A.
## V Peak Splittings in Power Spectra
### A Frequency Shifts in $`\delta g`$ and $`\delta g^{(2)}`$
We discuss now the emergence of the Berry phase in terms of a splitting of the frequencies of the Aharonov-Bohm oscillations in the magnetoconductance $`\delta g`$ and in the UCFs $`\delta g^{(2)}`$, which can be made visible in the power spectrum. Both quantities depend on the spin-dependent total phase $`\mathrm{\Phi }_\alpha `$, given here for the special case of the texture defined in Eq. (1) and for two equal fields $`𝐁=\stackrel{~}{}B`$,
$`\mathrm{\Phi }_{\pm 1}`$ $`=`$ $`2\varphi /\varphi _0\pm \mathrm{cos}\eta =2\varphi /\varphi _0\pm {\displaystyle \frac{1}{\sqrt{1+(B_{}/B_z)^2}}}`$ (41)
$``$ $`2\varphi /\varphi _0\pm B_z/B_{}=B_z\left(2B_{\varphi _0}^1\pm B_{}^1\right).`$ (42)
The approximation used here is valid for small perpendicular fields $`B_zB_{}`$. We have introduced $`B_{\varphi _0}=\varphi _0/A`$ as the perpendicular field which produces a flux of one flux quantum $`\varphi _0`$ through the ring, i.e. the period of an Aharonov-Bohm oscillation in $`\varphi `$. The Berry phase is not sensitive to the area enclosed by the ring; thus we prefer here to describe oscillations in $`B_z`$ rather than in $`\varphi `$. As both $`\delta g`$ and $`\delta g^{(2)}`$ contain periodic terms in $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_1`$, they exhibit oscillations in $`B_z`$ with the Aharonov-Bohm frequency for homogeneous fields, $`2B_{\varphi _0}^1`$, shifted (at $`B_z=0`$) by the frequency
$$\frac{1}{\mathrm{\Delta }B_1}=\pm \frac{1}{B_{}},$$
(43)
which results in a peak splitting in the power spectrum.
These splittings are, however, generally on the order of the resolution of the spectrum, which makes it difficult to make them visible. If the perpendicular field is varied from $`B_{\mathrm{max}}`$ to $`B_{\mathrm{max}}`$, the discrete Fourier transform (DFT) of such an interval has a resolution of $`1/2B_{\mathrm{max}}`$, i.e. the sampling frequencies are separated by this value. Thus, the peak-splitting term can only be made visible if this resolution is high enough, i.e. $`1/2B_{\mathrm{max}}1/B_{}`$, or
$$B_{\mathrm{max}}\frac{1}{2}B_{}.$$
(44)
We note that this restriction is still consistent with the approximation made in Eq. (41), since for $`B_z=B_{}/2`$ the approximated value of the Berry phase is larger than the exact value by only a factor of $`\sqrt{5}/21.1`$.
Now we consider the case beyond the above approximation. Here, an estimate for the frequency shifts can be obtained by counting the additional oscillations upon increasing $`B_z`$. In this estimation we again neglect the change in frequency of the Aharonov-Bohm oscillations while $`B_z`$ is increased. However, now we take the mean value of the frequency instead of the frequency at $`B_z=0`$ as in Eq. (43). Varying $`B_z`$ from $`0`$ to $`B_{\mathrm{max}}`$ changes the Berry phase contribution to $`\mathrm{\Phi }_{\pm 1}`$ \[Eq. (41)\] from $`0`$ to $`\pm \mathrm{cos}\eta |_{B_z=B_{\mathrm{max}}}`$, and so we obtain the mean frequency shift
$`{\displaystyle \frac{1}{\mathrm{\Delta }B_2}}`$ $`=`$ $`\pm {\displaystyle \frac{1}{\sqrt{B_{\mathrm{max}}^2+B_{}^2}}}\pm {\displaystyle \frac{1}{B_{}}}\left(1{\displaystyle \frac{B_{\mathrm{max}}^2}{2B_{}^2}}\right).`$ (45)
When we have calculated the DFT of $`\delta g`$ and $`\delta g^{(2)}`$, we have confirmed the predictions given above, i.e. we do not observe a peak splitting in the $`2B_{\varphi _0}^1`$ frequency for low $`B_{\mathrm{max}}`$, due to an insufficient resolution of the DFT. However, we do see a peak splitting in the DFT for higher fields (see Figs. 89), which vanishes again for $`B_{\mathrm{max}}B_{}`$. Since studies of the DFT suffer from a restricted resolution, it might be more promising to search for the Berry phase via the effects discussed in Sec. III A.
Finally, we point out that an anisotropic $`g`$ factor affects the size of the frequency splitting. If the $`g`$ factor perpendicular to the ring, $`g_z`$, is larger than the one in the plane of the ring, $`g_{}`$, the Berry phase dependence on $`B_z`$ increases while the Aharonov-Bohm phase remains unaffected. As the total phase is $`\mathrm{\Phi }_{\pm 1}2\varphi /\varphi _0\pm g_zB_z/g_{}B_{}`$, the frequency splitting is increased by a factor of $`g_z/g_{}`$.
### B Frequency shifts in $`\delta g_{\mathrm{hom}}^{(2)}`$ for homogeneous fields
At this point it is important to realize that frequency shifts can also appear in the conductance fluctuations $`\delta g^{(2)}`$ for homogeneous fields, i.e. even when there is no Berry phase present. For homogeneous fields the evaluation of Eq. (4) is straightforward, as $`h^{C/D}`$ \[Eq. (12)\] becomes diagonal, see also App. C. We evaluate the DOS terms, i.e. the terms containing $`\mathrm{Re}\mathrm{Tr}\widehat{\chi }_\omega \widehat{\chi }_\omega `$ in Eq. (4), in the low temperature limit for $`\eta =\stackrel{~}{\eta }=0`$:
$`\delta g_{\genfrac{}{}{0pt}{}{DOS,}{C/D}}^{(2)}\mathrm{Re}{\displaystyle \underset{\genfrac{}{}{0pt}{}{j}{\alpha ,\stackrel{~}{\alpha }=\pm 1}}{}}{\displaystyle \frac{1}{\left[\gamma +(j\mathrm{\Phi }^{C/D})^2+i(\alpha \kappa +\stackrel{~}{\alpha }\stackrel{~}{\kappa })\right]^2}}`$ (46)
$`{\displaystyle \frac{2\pi }{\gamma ^{3/2}}}+{\displaystyle \underset{\alpha ,\stackrel{~}{\alpha }}{}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{2\pi ^2n}{\gamma }}e^{2\pi n\sqrt{\gamma }}\mathrm{cos}[2\pi n(\mathrm{\Phi }^{C/D}+{\displaystyle \frac{\alpha \kappa +\stackrel{~}{\alpha }\stackrel{~}{\kappa }}{2\sqrt{\gamma }}})],`$ (47)
where we have defined $`\mathrm{\Phi }^{C/D}=\varphi /\varphi _0\pm \stackrel{~}{\varphi }/\varphi _0`$. The approximation on the second line of Eq. (46) is valid for $`\gamma 1/4\pi ^2,\alpha \kappa +\stackrel{~}{\alpha }\stackrel{~}{\kappa }`$. From Eq. (46), we see that the Zeeman term itself already leads to a frequency splitting. So, for instance, if we take the Fourier transform of $`\delta g^{(2)}(B_z,B_z)`$ with respect to $`B_z`$, we can observe a frequency splitting of the $`h/e`$ oscillations of the diffuson contribution in the DOS term $`\delta g_{DOS,D}^{(2)}`$, given by
$$\frac{1}{\mathrm{\Delta }B_{Zeeman}}=\pm \frac{g\mu _\mathrm{B}}{4\mathrm{}D}\frac{L_DL}{2\pi }.$$
(48)
We checked numerically that the estimated frequency splitting \[Eq. (48)\] is correct within 20 percent even for parameters beyond the assumptions made for the second line of Eq. (46). It is important to keep this property of the conductance fluctuations $`\delta g^{(2)}`$ in mind, when searching for Berry phase effects. If vanishing Aharonov-Bohm oscillations or peak splittings in the power spectrum are used to identify the presence of a Berry phase, one has to rule out effects coming from the Zeeman term in the UCFs, e.g. by comparison with the results for homogeneous fields.
### C Numerical Evaluations
We shall now numerically evaluate the magnetoconductance $`\delta g`$ for a ring in an inhomogeneous field. We base our analysis on the calculations from Ref. . In Fig. 7 we show the Aharonov-Bohm oscillations for different tilt angles $`\eta `$ of the external field $`𝐁`$, which is set so strong that we are well within the adiabatic regime. We can readily see that for $`\eta \pi /3`$ a phase shift of $`\pi `$ occurs, which comes directly from the Berry phase, compared to the oscillations at $`\eta =0`$ and $`\eta =\pi /2`$. For the intermediate tilt angles the effect of the Berry phase is only visible in the amplitude of the Aharonov-Bohm oscillations, as the phase shifts for the two spin directions occur with opposite signs and thus—if both spin directions contribute equally—no phase-shift effect is visible.
As such a phase shift at $`\pi /3`$ might not be easy to observe, studying signs in the power spectrum provides an interesting alternative, even though it requires a sufficiently high resolution, as discussed in Sec. V A. Indeed, we can observe a peak splitting in the spectrum of the magnetoconductance, as shown in the inset of Fig. 8. We notice an even more distinct feature: the Aharonov-Bohm oscillations vanish at two magic tilt angles, $`\mathrm{cos}\eta =0.4,\mathrm{\hspace{0.17em}0.75}`$, of the field. The mechanism for this effect is exhibited in Fig. 7, where it is shown how the two contributions of the different spins suppress the oscillations.
At this point, we would like to stress that the peak splitting depends strongly on the different dephasing terms. In particular, one cannot rely on calculations where the dephasing due to the inhomogeneous fields is not properly taken into account. So if the dephasing $`\gamma `$ due to homogeneous effects is very small, e.g. on the order of $`1/100`$, the amplitude of the oscillations gets reduced drastically as soon as the tilt angle $`\eta `$ changes from $`\pi /2`$ to a smaller, nonzero value, since the field inhomogenity causes additional dephasing. Thus the Fourier transform of such oscillations has a dominant contribution only from the first few oscillations close to $`\pi /2`$. This suppression of the remaining oscillations acts as a narrowing of the data window and leads to a widening of the peaks in the power spectrum, masking the peak splitting. The oscillations are further suppressed by the additional dephasing arising from an increasing perpendicular field, which penetrates the ring arms. Of course, it is possible to remove this unwanted over-emphasizing of certain oscillations from experimental data in a post-processing step; using a standard windowing function (we used the Hann window for the inset of Fig. 8) for DFTs greatly reduces this problem, in addition to the usual reduction of components leakage of neighboring frequencies in the power spectrum.
For the conductance fluctuations $`\delta g^{(2)}`$, we will further illustrate the effects of the two configurations discussed in Section III A. In Fig. 9 we show the Aharonov-Bohm oscillations occurring in $`\delta g^{(2)}`$ when the fields are equal, i.e. $`𝐁=\stackrel{~}{}B`$. Taking the discrete Fourier transform of $`\delta g^{(2)}`$ over the range $`B_z=0,\mathrm{},1\mathrm{T}`$, yields a clear peak splitting of the contribution of the $`h/2e`$ oscillations to the power spectrum, see left inset in Fig. 9. We notice a splitting into four peaks of the contribution of the $`h/4e`$ oscillations (right inset of Fig. 9). They only occur in the exact solution $`\delta g^{(2)}`$, whereas $`\delta g_{\mathrm{ad}}^{(2)}`$ exhibits only two peaks if we ignore the $`\eta `$, $`\stackrel{~}{\eta }`$-dependent dephasing, i.e. set $`\stackrel{~}{\gamma }_{\eta ,\stackrel{~}{\eta }}^{C/D}\gamma ^{C/D}`$ and $`P0`$ in Eq. (14). We point out that the frequency shifts for the $`n`$th harmonics of the Aharonov-Bohm oscillations increase with $`n`$ and are thus are better resolved in the power spectrum with increasing $`n`$.
We plot $`\delta g^{(2)}(\stackrel{~}{}B)`$ in Fig. 10 for the special case $`\stackrel{~}{}B=(0,\mathrm{\hspace{0.17em}0},\stackrel{~}{B}_z)`$ homogeneous (see also Sec. III A).
Finally, we consider the power spectrum of the magnetoconductance $`\delta g_{SO}`$ in the presence of spin-orbit coupling. We use Eq. (40) and ignore for simplicity dephasing due to the external magnetic fields penetrating the arms of the ring. Indeed, taking the Fourier transform of the magnetoconductance, a spin splitting can be observed. However, the splitting is not as pronounced as in the case for inhomogeneous fields. Especially important, the splitting is only visible for sufficiently large dephasing parameters $`\gamma `$ (produced by inelastic scattering), which can be seen in Fig. 11. In contrast to the effects discussed before, using a windowing function was not sufficient to identify a peak splitting for moderately small dephasing parameters $`\gamma 0.3`$. Qualitatively, however, the power spectra of the magnetoconductance for inhomogeneous magnetic fields and for spin-orbit coupling agree, with both showing a peak splitting.
The UCFs with spin-obit interaction $`\delta g_{SO}^{(2)}`$ are plotted in Fig. 12 as a function of the perpendicular fields $`B_z=\stackrel{~}{B}_z`$. We observe a Berry phase-like frequency splitting in the power spectrum. However, as this splitting is rather small, it is only visible in the $`h/4e`$ oscillations, where the splitting is twice as large as in the $`h/2e`$ oscillations. Again, the suppression of the Aharonov-Bohm oscillations at $`\kappa /S0.25`$ is a distinct feature of a Berry phase-like effect.
A quantity, which was subject of recent studies, is the disorder-averaged squared power spectrum of the conductance
$$\left|g(\nu )\right|^2=\left|g(\nu )\right|^2+\left|g(\nu )g(\nu )\right|^2.$$
(49)
On the one hand, we recognize that the first term contains the Fourier transform of the (averaged) magnetoconductance $`\delta g`$, which has frequency contributions from its $`h/2e`$ oscillations. On the other hand, the second term of Eq. (49) is given through the conductance fluctuations $`\delta g^{(2)}`$ as $`𝑑B_z𝑑\stackrel{~}{B}_z\mathrm{exp}\{2\pi i\nu (B_z\stackrel{~}{B}_z)\}\delta g^{(2)}(𝐁,\stackrel{~}{}B)`$. This term contributes frequencies corresponding to $`h/e`$ oscillations, coming from the diffuson term $`\delta g_D^{(2)}`$ in the conductance fluctuations. Thus, if we now investigate $`h/e`$ oscillations, we can restrict our studies to the second term of Eq. (49). We have evaluated $`|g(\nu )|^2`$ for inhomogeneous fields, with the parameters given in the caption of Fig. 9. A splitting of the frequency corresponding to the $`h/e`$ oscillations was observed and was identified not to result from the Berry phase but from the Zeeman term already present in the case of homogeneous fields \[Eq. (48)\]. Then we examined $`|g_{SO}(\nu )|^2`$ with spin-orbit coupling for various parameters. An additional peak splitting to the one produced by the Zeeman term \[Eq. (48)\] appears for some specific parameters, i.e. for $`S`$ large enough to reach “adiabaticity” and for large enough sampling intervals of $`B_z`$, $`\stackrel{~}{B}_z`$ to obtain a sufficiently high resolution in the power spectrum. In Fig. 13 we see such a splitting of the $`h/e`$ contribution into four peaks. However, using the parameters given in Ref. , we have $`S1/50`$ and $`\kappa 1/2`$ (see Sec. IV B) and in this regime we do not observe any peak splitting, in accordance with Ref. .
## VI Conclusion
We have calculated the exact conductance fluctuations $`\delta g^{(2)}`$ for a special texture \[Eq. (1)\] and given its adiabatic approximation $`\delta g_{\mathrm{ad}}^{(2)}`$. In addition to the already known differential equations for the cooperon we have derived the ones for the diffuson in inhomogeneous magnetic fields (App. A). With the result $`\delta g_{\mathrm{ad}}^{(2)}`$ the dephasing due to inhomogeneous fields became explicit and could be compared with previous calculations where adiabatic eigenstates were used and this dephasing was only implemented with a phenomenological parameter. Then we have described some magic tilt angles of the magnetic field at which the Berry phase suppresses the Aharonov-Bohm oscillations. We have used this effect to illustrate how the adiabatic criterion becomes gradually satisfied. We have calculated numerically the required magnetic field strength for which the adiabatic approximation becomes valid and have shown that the adiabatic criterion is less stringent for diffusive than for ballistic motion, thus confirming previous findings.
Furthermore, we have calculated the magnetoconductance and the conductance fluctuations for a diffusive conductor in the presence of spin-orbit coupling. A numerical analysis revealed a non-monotonic behavior of the amplitudes of the Aharonov-Bohm oscillations and peak-splittings in the power spectrum—observations that are similar to the Berry phase effects we have found for inhomogeneous magnetic fields.
Finally, we have described the mechanisms which lead to peak splittings in the power spectrum of magnetoconductance and UCFs and have discussed numerical requirements to make such peaks splittings visible.
###### Acknowledgements.
We would like to thank G. Burkard, R. Häussler, and E.V. Sukhorukov for many fruitful discussions. This work has been supported in part by the Swiss National Science Foundation.
## A Differential equations for Cooperon and Diffuson
Here we will transform the exact conductance correlator for diffusive systems and arbitrary magnetic textures to make a Schrödinger equation approach possible. Further we will derive the explicit differential equation for the diffuson propagator (the one for the cooperon has been derived previously ).
The conductance correlator has been derived in Ref. , using diagrammatic techniques, and is given by
$`\delta g^{(2)}`$ $`=`$ $`\left({\displaystyle \frac{2e^2D}{hL^2}}\right)^2{\displaystyle }dϵdϵ^{}n^{}(ϵ)n^{}(ϵ^{}){\displaystyle }d𝐱d𝐱^{}{\displaystyle \underset{\alpha _1,\alpha _2,\alpha _3,\alpha _4}{}}\left\{{\displaystyle \frac{1}{d}}\right|\chi _{\alpha _1\alpha _2,\alpha _3\alpha _4}^C(𝐱,𝐱^{},\omega )|^2`$ (A1)
$`+`$ $`2\mathrm{Re}\left[\chi _{\alpha _1\alpha _2,\alpha _3\alpha _4}^C(𝐱,𝐱^{},\omega )\chi _{\alpha _2\alpha _1,\alpha _4\alpha _3}^C(𝐱^{},𝐱,\omega )\right]+[\chi ^C\chi ^D]\},`$ (A2)
where $`n^{}(ϵ)`$ is the derivative of the Fermi function, $`\mathrm{}\omega =ϵϵ^{}`$, and $`d`$ describes the dimension of the system with respect to the mean free path $`l`$. The inverse Fourier transform of the cooperon/diffuson propagators $`\chi ^{C/D}(𝐱^{},𝐱,w)`$ were obtained as
$`\chi _{\alpha _1\alpha _2,\alpha _3\alpha _4}^{C/D}(𝐱^{},𝐱;t^{},t)`$ $`=`$ $`\theta (t^{}t){\displaystyle _{𝐑(t)=𝐱}^{𝐑(t^{})=𝐱^{}}}𝒟𝐑\mathrm{exp}\left\{{\displaystyle \frac{1}{4D}}{\displaystyle _t^t^{}}𝑑\tau |\dot{}R|^2\right\}`$ (A3)
$`\times `$ $`\mathrm{exp}\left\{i{\displaystyle \frac{e}{\mathrm{}}}{\displaystyle _t^t^{}}𝑑\tau \left[\dot{}R𝐀^{\mathrm{em}}(𝐑(\tau ))+\dot{}R^\pm \stackrel{~}{}A^{\mathrm{em}}(𝐑^\pm (\tau ))\right]\right\}`$ (A4)
$`\times `$ $`\alpha _4\alpha _2\left|𝒯\mathrm{exp}\left\{i{\displaystyle \frac{g\mu _\mathrm{B}}{2\mathrm{}}}{\displaystyle _t^t^{}}𝑑\tau \left[𝐁(𝐑(\tau ))𝝈_1\stackrel{~}{}B(𝐑^\pm (\tau ))𝝈_2\right]\right\}\right|\alpha _3\alpha _1,`$ (A5)
where $`𝐑^{}(\tau )=𝐑(t^{}+t\tau )`$ is the time-reversed path of $`𝐑^+𝐑`$.
For explicit evaluation it is convenient to transform this path-integral representation into a differential equation. In the case of the diffuson we first have to eliminate the time-reversed paths. As a result of reverting the time integration, an additional sign appears in the second term of the electromagnetic vector potential. For the Zeeman interaction we can use the relation
$`\alpha _2\left|𝒯\mathrm{exp}\left\{i{\displaystyle \frac{g\mu _\mathrm{B}}{2\mathrm{}}}{\displaystyle _t^t^{}}𝑑\tau \stackrel{~}{}B(𝐑^{}(\tau ))𝝈\right\}\right|\alpha _1=\alpha _1\left|𝒯\mathrm{exp}\left\{i{\displaystyle \frac{g\mu _\mathrm{B}}{2\mathrm{}}}{\displaystyle _t^t^{}}𝑑\tau \stackrel{~}{}B(𝐑(\tau ))𝝈\right\}\right|\alpha _2^{}`$ (A6)
$`=\alpha _1\left|𝒯\mathrm{exp}\left\{i{\displaystyle \frac{g\mu _\mathrm{B}}{2\mathrm{}}}{\displaystyle _t^t^{}}𝑑\tau \stackrel{~}{}B(𝐑(\tau ))𝝈^{}\right\}\right|\alpha _2.`$ (A7)
The latter equation can be proven by writing the time-ordered product as a Dyson series and by inserting a resolution of unity in spin space between all products $`(𝐁(x_j)𝝈)(𝐁(x_{j+1})𝝈)`$, thereby arriving at an expression with terms of the form $`\alpha \left|B_i(x_j)\sigma _i\right|\beta ^{}`$. Such terms are the complex conjugate of Pauli matrix elements multiplied by the real number $`B_i(x_j)`$. So we can rewrite them as $`\alpha \left|B_i(x_j)\sigma _i^{}\right|\beta `$, remove the previously inserted unities, and go back to the time-ordered product.
Now we can give the differential equations for the propagators
$`\mathbf{(}{\displaystyle \frac{}{t^{}}}+D[i{\displaystyle \frac{}{𝐱^{}}}{\displaystyle \frac{e}{\mathrm{}}}[𝐀^{\mathrm{em}}(𝐱^{})\pm \stackrel{~}{}A^{\mathrm{em}}(𝐱^{})]]^2`$ (A8)
$`i{\displaystyle \frac{g\mu _\mathrm{B}}{2\mathrm{}}}[𝐁(𝐱^{})𝝈_1\stackrel{~}{}B(𝐱^{})𝝈_2^{()}]\mathbf{)}\widehat{\chi }^{C/D}(𝐱^{},𝐱;t^{},t)=\delta (𝐱^{}𝐱)\delta (t^{}t)\widehat{11},`$ (A9)
where $`\widehat{\chi }^{C/D}(𝐱^{},𝐱;t^{},t)`$ is a matrix in four-dimensional spin space. The upper sign is for the cooperon , the lower sign and the complex conjugate of $`𝝈_2`$ for the diffuson. Passing to Fourier space and operator notation, the above equation becomes
$$\left(i\omega D\frac{(2\pi )^2}{L^2}h^{C/D}\right)\widehat{\chi }_\omega ^{C/D}=\widehat{11},$$
(A10)
where the effective Hamiltonian $`h^{C/D}`$ is defined in Eq. (6).
Finally we express the conductance correlation in terms of the operators $`\chi _\omega ^{C/D}`$. We note that with $`\chi _{\alpha _1\alpha _2,\alpha _3\alpha _4}^C(𝐱,𝐱^{},\omega )^{}=𝐱^{},\alpha _4\alpha _2|\widehat{\chi }_\omega ^C|𝐱,\alpha _3\alpha _1^{}=𝐱,\alpha _3\alpha _1|\widehat{\chi }_\omega ^C|𝐱^{},\alpha _4\alpha _2`$ , and $`\chi _{\alpha _1\alpha _2,\alpha _3\alpha _4}^D(𝐱,𝐱^{},\omega )^{}=𝐱^{},\alpha _4\alpha _1|\widehat{\chi }_\omega ^D|𝐱,\alpha _3\alpha _2^{}=𝐱,\alpha _3\alpha _2|\widehat{\chi }_\omega ^D|𝐱^{},\alpha _4\alpha _1`$ we can simplify the terms in Eq. (A1):
$$𝑑𝐱𝑑𝐱^{}\underset{\alpha _1,\mathrm{},\alpha _4}{}\left|\chi _{\alpha _1\alpha _2,\alpha _3\alpha _4}^{C/D}(𝐱,𝐱^{},\omega )\right|^2=\mathrm{Tr}\widehat{\chi }_\omega ^{C/D}\widehat{\chi }_\omega ^{C/D},$$
(A11)
and
$$𝑑𝐱𝑑𝐱^{}\underset{\alpha _1,\mathrm{},\alpha _4}{}\chi _{\alpha _1\alpha _2,\alpha _3\alpha _4}^{C/D}(𝐱,𝐱^{},\omega )\chi _{\alpha _2\alpha _1,\alpha _4\alpha _3}^{C/D}(𝐱^{},𝐱,\omega )=\mathrm{Tr}\widehat{\chi }_\omega ^{C/D}\widehat{\chi }_\omega ^{C/D}.$$
(A12)
## B Finite Temperature Integrals
We shall explain here the integrations performed to obtain Eq. (26). We are interested in
$$I=𝑑ϵ^{}n^{}(ϵ^{})J=𝑑ϵ^{}n^{}(ϵ^{})𝑑ϵn^{}(ϵ)\left(\frac{1}{d}\frac{1}{(ϵϵ^{}+a)^2+c^2}+2\mathrm{R}\mathrm{e}\frac{1}{(iϵiϵ^{}+iac)^2}\right)$$
(B1)
with $`a,c`$ real and $`c>0`$. We consider a rectangular integration contour $`\mathrm{\Gamma }`$ with one side lying on the real axis, extending $`M=2\pi l/\beta `$ towards the positive imaginary axis and the same amount on each side of the real axis. For any positive integer $`l`$, the absolute value of the Fermi function is bounded above on such a contour: $`|n(z)||_\mathrm{\Gamma }<2`$. The integrands considered further below are a product of the Fermi function and a rational function decaying with at least $`|z|^2`$. The integral of these products over the section of $`\mathrm{\Gamma }`$ in the upper half plane, will thus vanish for $`M\mathrm{}`$, as we have $`|z|M`$ on this contour. We further note, that the complex expansion of the Fermi function $`n(z)`$ has its poles at $`z=i\omega _n`$, where $`\omega _n=\pi n/\beta `$ are the Matsubara frequencies and $`n`$ is an odd integer.
We expand the first rational function in Eq. (B1) into partial fractions and then integrate $`J`$ by parts:
$`J`$ $`=`$ $`{\displaystyle 𝑑ϵn(ϵ)\left\{\frac{1}{d}\frac{1}{2ic}\left(\frac{1}{(ϵϵ^{}+aic)^2}\frac{1}{(ϵϵ^{}+a+ic)^2}\right)+2\mathrm{R}\mathrm{e}\frac{2}{(ϵϵ^{}+a+ic)^3}\right\}}`$ (B2)
$`=`$ $`\mathrm{Re}{\displaystyle 𝑑ϵn(ϵ)\left\{\frac{1}{d}\frac{i}{c(ϵϵ^{}+a+ic)^2}\frac{4}{(ϵϵ^{}+a+ic)^3}\right\}}.`$ (B3)
We now evaluate the integral along the contour described above. As the poles of the rational functions in Eq. (B3) are in the lower half plane at $`ϵ^{}aic`$, they are not within the integration contour. Applying Cauchy’s residue theorem and accounting for the residues of the Fermi function $`\mathrm{res}n(z)|_{z=i\omega _n}=(1/\beta )`$ yields
$$J=\frac{2\pi }{\beta }\mathrm{Re}\underset{n\mathrm{odd}>0}{}\left\{\frac{1}{d}\frac{1}{c(i\omega _nϵ^{}+a+ic)^2}+\frac{4i}{(i\omega _nϵ^{}+a+ic)^3}\right\}.$$
(B4)
For the second integration in Eq. (B1), we replace the expression in braces in the above equation by its complex conjugate. As before, we first integrate by parts over $`ϵ^{}`$ and apply the residue theorem afterwards. This results in
$`I`$ $`=`$ $`{\displaystyle \frac{4\pi }{\beta }}\mathrm{Re}{\displaystyle \underset{n\mathrm{odd}>0}{}}{\displaystyle 𝑑ϵ^{}n(ϵ^{})\left\{\frac{1}{d}\frac{1}{c(i\omega _n+ϵ^{}a+ic)^3}+\frac{6i}{(i\omega _n+ϵ^{}a+ic)^4}\right\}}`$ (B5)
$`=`$ $`{\displaystyle \frac{8\pi ^2}{\beta ^2}}\mathrm{Re}{\displaystyle \underset{n,m\mathrm{odd}>0}{}}\left\{{\displaystyle \frac{1}{d}}{\displaystyle \frac{1}{c(\omega _n+\omega _m+ia+c)^3}}+{\displaystyle \frac{6}{(\omega _n+\omega _m+ia+c)^4}}\right\}.`$ (B6)
## C UCFs $`\delta g_{\mathrm{hom}}^{(2)}`$ for homogeneous fields
For homogeneous fields we have $`\eta =\stackrel{~}{\eta }=0`$ and $`f=0`$, thus the hamiltonians $`h^{C/D}`$ \[Eq. (11)\] become diagonal with the matrix elements $`j^2+i\alpha \kappa i\stackrel{~}{\alpha }\stackrel{~}{\kappa }`$. Now we evaluate the propagators $`\widehat{\chi }^{C/D}`$ \[Eq. (5)\] and by evaluating the integrals over the Fermi functions in Eq. (4) explicitly by using standard Matsubara techniques, as explained in App. B. We obtain $`\delta g_{\mathrm{hom}}^{(2)}=\delta g_{\mathrm{hom},C}^{(2)}+\delta g_{\mathrm{hom},D}^{(2)}`$, where
$`\delta g_{\genfrac{}{}{0pt}{}{\mathrm{hom},}{C/D}}^{(2)}`$ $`=\left({\displaystyle \frac{e^2}{h}}\right)^2{\displaystyle \frac{1}{8\pi ^6}}\left({\displaystyle \frac{L^2}{L_T^2}}\right)^2\mathrm{Re}{\displaystyle \underset{\alpha ,\stackrel{~}{\alpha }=\pm 1}{}}{\displaystyle \underset{j=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n,m}{}^{}}`$ (C3)
$`\{{\displaystyle \frac{1}{d(\gamma ^{C/D}+j^2)\left[b_{nm}+\gamma ^{C/D}+j^2+i(\alpha \kappa \stackrel{~}{\alpha }\stackrel{~}{\kappa })\right]^3}}`$
$`+{\displaystyle \frac{6}{\left[b_{nm}+\gamma ^{C/D}+j^2+i(\alpha \kappa \stackrel{~}{\alpha }\stackrel{~}{\kappa })\right]^4}}\},`$
and we have introduced $`b_{nm}=(n+m)(L/L_T)^2/4\pi `$. Here $`n`$ and $`m`$ are positive, odd integers. The Aharonov-Bohm flux is implemented by replacing $`jj(\varphi /\varphi _0\pm \stackrel{~}{\varphi }/\varphi _0)`$. For further evaluation we now set $`\kappa =\stackrel{~}{\kappa }`$. We describe the summation of cooperon and diffuson terms with a prefactor $`\beta `$, which is $`1`$ if both terms contribute and $`2`$ if time-reversal symmetry is broken, so the cooperon contribution vanishes. Thus we have $`\delta g_{\mathrm{hom}}^{(2)}(2/\beta )\delta g_{\mathrm{hom},D}^{(2)}`$ and from now on we only consider the dephasing parameter $`\gamma =\gamma ^D=L^2/(2\pi L_\phi )^2`$, according to Eq. (8).
If the spin-channel mixing is suppressed (i.e. $`\kappa \gamma `$) in Eq. (C3), we can replace the sum over the spins $`_{\alpha \stackrel{~}{\alpha }}`$ by the number of spin states $`g_s`$. For weaker magnetic fields ($`\kappa \gamma `$) we have full spin degeneracy and obtain the factor $`g_s^2`$. Accounting for valley degeneracy yields a factor $`g_v^2`$.
Since we will check our results against the ones given in Ref. , where one-dimensional systems were considered, we take $`d=1`$. Since we will evaluate some limiting cases below, where $`L2\pi L_\phi `$, we have $`\gamma 1`$ and can therefore replace the $`j`$-sum in Eq. (C3) by an integral. The Aharonov-Bohm phase can then be removed by shifting the integration variable $`j`$ and we obtain
$`\delta g_{\mathrm{hom}}^{(2)}`$ $`=`$ $`\left({\displaystyle \frac{e^2}{h}}\right)^2{\displaystyle \frac{1}{8\pi ^6}}{\displaystyle \frac{2g_s^2g_v^2}{\beta }}\left({\displaystyle \frac{L^2}{L_T^2}}\right)^2{\displaystyle \underset{n,m}{}^{}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑j`$ (C5)
$`\left\{{\displaystyle \frac{1}{d(\gamma +j^2)\left[b_{nm}+\gamma +j^2\right]^3}}+{\displaystyle \frac{6}{\left[b_{nm}+\gamma +j^2\right]^4}}\right\}.`$
In the limit $`(2\pi )^2L_\phi ^2L^2,2\pi L_T^2`$, we have
$$\frac{2\pi L_T^2}{L^2}(\gamma +j^2)\frac{L_T^2}{2\pi L_\phi ^2}1.$$
(C6)
Thus, we can use Poisson’s summation formula to replace the summation over $`n`$ and $`m`$ in Eq. (C5) by integration to arrive at
$`\delta g_{\mathrm{hom}}^{(2)}`$ $`=`$ $`{\displaystyle \frac{3}{4\pi ^4}}{\displaystyle \frac{g_s^2g_v^2}{\beta }}\left({\displaystyle \frac{e^2}{h}}\right)^2{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{dj}{\left(\gamma +j^2\right)^2}}`$ (C8)
$`=3{\displaystyle \frac{g_s^2g_v^2}{\beta }}\left({\displaystyle \frac{e^2}{h}}\right)^2\left({\displaystyle \frac{L_\phi }{L}}\right)^3.`$
We now consider another limit, $`2\pi L_T^2L^2,(2\pi )^2L_\phi ^2`$. Again, we start from Eq. (C5), but now we first calculate the integral over $`j`$, which has the dominant contribution $`\pi \gamma ^{1/2}b^3`$ since $`1\gamma b_{nm}`$. Thus we obtain
$`\delta g_{\mathrm{hom}}^{(2)}`$ $`=`$ $`{\displaystyle \frac{4}{\pi }}{\displaystyle \frac{g_s^2g_v^2}{\beta }}\left({\displaystyle \frac{e^2}{h}}\right)^2{\displaystyle \frac{L_T^2}{L^2}}{\displaystyle \frac{L_\phi }{L}}{\displaystyle \underset{n,m}{}^{}}{\displaystyle \frac{1}{\left(\frac{1}{2}(n+m)\right)^3}}`$ (C9)
$`=`$ $`{\displaystyle \frac{2\pi }{3}}{\displaystyle \frac{g_s^2g_v^2}{\beta }}\left({\displaystyle \frac{e^2}{h}}\right)^2{\displaystyle \frac{L_T^2}{L^2}}{\displaystyle \frac{L_\phi }{L}}.`$ (C10)
Indeed, our results $`\delta g_{\mathrm{hom}}^{(2)}`$ given in Eqs. (C8) and (C9) agrees with these of Ref. . Thus, on the one hand, we have confirmed that the result from Ref. \[used in Eq. (4)\] is consistent with earlier calculations. On the other hand, in Eq. (C3) we have given an explicit formula (not known before as far as we are aware of) describing how the spin-channel mixing becomes suppressed for increasing magnetic fields, such that $`\delta g_{\mathrm{hom}}^{(2)}`$ contains a prefactor $`g_s^2`$ for low and $`g_s`$ for high magnetic fields.
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# Dimension estimates for Hilbert schemes and effective base point freeness on moduli spaces of vector bundles on curves
## 1. Background
The underlying idea for studying linear series on the moduli space $`SU_X(r,e)`$ has its roots in the paper of Faltings , where a construction of the moduli space based on theta divisors is given. A very nice introduction to the subject is provided in .
Fix $`r`$ and $`e`$ and denote $`h=\mathrm{gcd}(r,e)`$, $`r_1=\frac{r}{h}`$ and $`e_1=\frac{e}{h}`$. Consider a vector bundle $`F`$ of rank $`pr_1`$ and degree $`p(r_1(g1)e_1)`$. Generically such a choice determines (cf. 0.2) a *theta divisor* $`\mathrm{\Theta }_F`$ on $`SU_X(r,e)`$, supported on the set
$$\mathrm{\Theta }_F=\{E|h^0(EF)0\}.$$
All the divisors $`\mathrm{\Theta }_F`$ for $`FU_X(pr_1,p(r_1(g1)e_1))`$ belong to the linear system $`|^p|`$, where $``$ is the *determinant* line bundle $``$. We have the following well-known:
###### Lemma 1.1.
$`ESU_X(r,e)`$ is not a base point for $`|^p|`$ if there exists a vector bundle $`F`$ of rank $`pr_1`$ and degree $`p(r_1(g1)e_1)`$ such that $`h^0(EF)=0`$.
It is easy to see that such an $`F`$ must necessarily be semistable (cf. (2.5)). It is also a simple consequence of the existence of Jordan-Hölder filtrations that one has to check the condition in the above lemma only for $`E`$ stable. We sketch the proof for convenience:
###### Lemma 1.2.
If for any stable bundle $`V`$ of rank $`r^{}r`$ and slope $`e/r`$ there exists $`FU_X(pr_1,p(r_1(g1)e_1))`$ such that $`h^0(VF)=0`$, then the same is true for every $`ESU_X(r,e)`$.
###### Proof.
Assume that $`E`$ is strictly semistable. Then it has a Jordan-Hölder filtration:
$$0=E_0E_1\mathrm{}E_n=E$$
such that $`E_i/E_{i1}`$ are stable for $`i\{1,\mathrm{},n\}`$ and $`\mu (E_i/E_{i1})=\frac{e}{r}`$. By assumption there exist $`F_iU_X(pr_1,p(r_1(g1)e_1))`$ such that $`h^0(E_i/E_{i1}F_i)=0`$ and so if we denote
$$\mathrm{\Theta }_{E_i/E_{i1}}:=\{F|h^0(E_i/E_{i1}F)0\}U_X(pr_1,p(r_1(g1)e_1)),$$
this is a proper subset for every $`i`$. It is clear that any
$$FU_X(pr_1,p(r_1(g1)e_1))\underset{i=1}{\overset{𝑛}{}}\mathrm{\Theta }_{E_i/E_{i1}}$$
satisfies $`h^0(EF)=0`$. ∎
We also record a simple lemma which will be useful in §4. It is most certainly well known, but we sketch the proof for convenience (cf. also 1.1).
###### Lemma 1.3.
Consider a sheaf extension:
$$0FEG0.$$
If $`E`$ is stable, then $`h^0(G^{}F)=0`$.
###### Proof.
Assuming the contrary, there is a nonzero morphism $`GF`$. Composing this with the maps $`EG`$ to the left and $`FE`$ to the right, we obtain a nontrivial endomorphism of $`E`$, which contradicts the stability assumption. ∎
As a final remark, note that we are always slightly abusing the notation by using vector bundles instead of $`S`$-equivalence classes. This is harmless, since it is easily seen that it is enough to check the assertions for any representative of the equivalence class.
## 2. An upper bound on the dimension of Hilbert schemes
The goal of this paragraph is to prove a result (see Theorem 2.2 below) giving an upper bound on the dimension of the Hilbert schemes of coherent quotients of fixed rank and degree of a given vector bundle, optimal at least in the case corresponding to line subbundles. For the general theory of Hilbert schemes the reader can consult for example §4.
Concretely, fix a vector bundle $`E`$ of rank $`r`$ and degree $`e`$ on $`X`$ and denote by $`\mathrm{Quot}_{rk,ed}(E)`$ the Hilbert scheme of coherent quotients of $`E`$ of rank $`rk`$ and degree $`ed`$. We can (and will) identify $`\mathrm{Quot}_{rk,ed}(E)`$ to the set of subsheaves of $`E`$ of rank $`k`$ and degree $`d`$. Consider also:
$$d_k:=\mathrm{max}\{\mathrm{deg}(F)|FE,\mathrm{rk}(F)=k\}$$
and
$$M_k(E)=\{F|FE,\mathrm{rk}(F)=k,\mathrm{deg}(F)=d_k\}$$
the set of maximal subbundles of rank $`k`$. Clearly any $`FM_k(E)`$ has to be a vector subbundle of $`E`$. Note that the number $`ed_k`$ is exactly the minimal degree of a quotient bundle of $`E`$ of rank $`rk`$. By §2 we have the following basic result:
###### Proposition 2.1.
The following hold and are equivalent:
(i) dim $`M_k(E)k(rk)`$
(ii) for any $`xX`$ and any $`WE(x)`$ $`k`$-dimensional subspace of the fiber of $`E`$ at $`x`$, there are at most finitely many $`FM_k(E)`$ such that $`F(x)=W`$.
Part $`(i)`$ above thus gives an upper bound for the dimension of the Hilbert scheme in the case $`d=d_k`$. The next result is a generalization in the case of arbitrary degree $`d`$, which turns out to give an optimal result (see Example 2.13 below).
###### Theorem 2.2.
With the notation above, we have:
$$\mathrm{dim}\mathrm{Quot}_{rk,ed}(E)k(rk)+(d_kd)k(rk+1).$$
###### Remark 2.3.
To avoid any confusion, we emphasize here that the notation is slightly different from that used in the introduction, in the sense that we are replacing $`k`$ by $`rk`$, $`d`$ by $`ed`$ and $`f_k`$ by $`ed_k`$. This is done for consistency in rewriting everything in terms of subbundles, but note that the statement is exactly the same.
The proof will proceed by induction on the difference $`d_kd`$. In order to perform this induction we have to use a special case of the notion of *elementary transformation* along a zero-dimensional subscheme of arbitrary length. We call this construction simply elementary transformation since there is no danger of confusion.
###### Definition 2.4.
Let $`\tau `$ be a zero-dimensional subscheme of $`X`$ supported on the points $`P_1,\mathrm{},P_s`$. An *elementary transformation* of $`E`$ along $`\tau `$ is a vector bundle $`E^{}`$ defined by a sequence of the form:
$$0E^{}E\stackrel{\mathit{\varphi }}{}\tau 0.$$
where the morphism $`\varphi `$ is determined by giving surjective maps $`E_{P_i}\stackrel{\varphi _i}{}𝐂_{P_i}^{a_i}`$ induced by specifying $`a_i`$ distinct hyperplanes in $`E(P_i)`$ (whose intersection is the kernel of $`\varphi _i`$) , $`i\{1,\mathrm{},s\}`$. We call $`m=a_1+\mathrm{}+a_s`$ the *length* of $`\tau `$ and $`a_i`$ the *weight* of $`P_i`$.
Let us briefly remark that this is not the most general definition, since we are imposing a condition on the choice of hyperplanes. We prefer to work with this notion because it is sufficient for our purposes and allows us to avoid some technicalities. Note though that the space parametrizing these transformations is not compact. One could equally well work with the general definition, when the hyperplanes could come together, and obtain a compact parameter space, which can be shown to be irreducible (it is basically a Hilbert scheme of rank zero quotients of fixed length).
In fact it is an immediate observation that the elementary transformations of $`E`$ of length $`m`$, in the sense of the definition above, are parametrized by $`Y:=(𝐏E)_m\mathrm{\Delta }`$, where $`(𝐏E)_m`$ is the $`m`$-th symmetric product of the projective bundle $`𝐏E`$ and $`\mathrm{\Delta }`$ is the union of all its diagonals. There is an obvious forgetful map
$$\pi :YX_m,$$
where $`X_m`$ is the $`m`$-th symmetric product of the curve $`X`$. We will denote by $`Y_m(𝐏E)_m\mathrm{\Delta }`$ the open subset $`(𝐏E)_m\pi ^1(\delta )`$, where $`\delta `$ is the union of the diagonals in $`X_m`$. Its points correspond to the elementary transformations of length $`m`$ supported at $`m`$ distinct points of $`X`$.
###### Definition 2.5.
Let $`V`$ be a subbundle of $`E`$. An elementary transformation
$$0E^{}E\stackrel{\mathit{\varphi }}{}\tau 0$$
is said to *preserve* $`V`$ if the inclusion $`VE`$ factors through the inclusion $`E^{}E`$.
###### Lemma 2.6.
If $`E^{}`$ is determined by the hyperplanes $`H_i^1,\mathrm{},H_i^{a_i}E(P_i)`$ for $`i\{1,\mathrm{},s\}`$ and $`V_i:=\underset{j=1}{\overset{a_i}{}}H_i^j`$, then $`V`$ is preserved by $`E^{}`$ if and only if $`V(P)V_i,i\{1,\mathrm{},s\}`$.
###### Proof.
We have an induced diagram
where $`\alpha `$ is the composition of $`\varphi `$ with the inclusion of $`V`$ in $`E`$. It is clear that $`E^{}`$ preserves $`V`$ if and only if $`\alpha `$ is identically zero. The lemma follows then easily from the definitions. ∎
In general it is important to know the dimension of the set of elementary transformations of a certain type preserving a given subbundle. This is given by the following simple lemma:
###### Lemma 2.7.
Let $`VE`$ be a subbundle of rank $`k`$. Consider the set of elementary transformations of $`E`$ along a zero dimensional subscheme of length $`m`$ belonging to an irreducible subvariety $`W`$ of $`X_m\delta `$ that preserve $`V`$:
$$𝒵_V:=\{E^{}|VE^{}\}\pi ^1(W),$$
where $`\pi :Y_mX_m`$ is the natural projection. Then $`𝒵_V`$ is irreducible of dimension $`m(rk1)+\mathrm{dim}W`$.
###### Proof.
An elementary transformation at $`m`$ points $`x_1,\mathrm{},x_m`$ is given by a choice of hyperplanes $`H_iE(x_i)`$ for each $`i`$. By the previous lemma, such a transformation preserves $`V`$ if and only if $`V(x_i)H_i`$ for all $`i`$. We have a natural diagram
where $`\pi `$ is the restriction to $`\pi ^1(W)`$ and $`p`$ is the composition with the inclusion of $`𝒵_V`$ in $`\pi ^1(W)`$. For $`D=x_1+\mathrm{}+x_mW`$, we have:
$$p^1(D)\{(H_1,\mathrm{},H_m)|V(x_i)H_iE(x_i),i=1,\mathrm{},m\}$$
$$𝐏^{rk1}\times \mathrm{}\times 𝐏^{rk1}$$
where the product is taken $`m`$ times. So $`p^1(D)`$ is irreducible of dimension $`m(rk1)`$ and this gives that $`𝒵_V`$ is irreducible of dimension $`m(rk1)+\mathrm{dim}W`$. ∎
The following proposition will be the key step in running the inductive argument. It computes “how fast” we can eliminate all the maximal subbundles of $`E`$ while preserving a fixed nonmaximal subbundle.
###### Proposition 2.8.
Let $`VE`$ be a subbundle of rank $`k`$ and degree $`d`$, not maximal. Then if $`mrk+1`$, there exists an elementary transformation of length $`m`$
$$0E^{}E\tau 0$$
such that $`VE^{}`$, but $`FE^{}`$ for any $`FM_k(E)`$. In other words $`E^{}`$ preserves $`V`$, but does not preserve any maximal $`F`$. If we fix a point $`PX`$, then $`\tau `$ can be chosen to have weight $`m1`$ at $`P`$ and weight $`1`$ at a generic point $`QX`$.
###### Proof.
Fix a point $`PX`$. We can consider an elementary transformation of $`E`$ of length $`rk`$, supported only at $`P`$:
$$0E^{}E\tau 0,$$
such that $`\mathrm{Im}(E^{}(P)E(P))=V(P)`$. Then as in 2.6, the only maximal subbundles $`F`$ that are preserved by this transformation are exactly those such that $`F(P)=V(P)`$. By Proposition 2.1 this implies that only at most a finite number of $`F`$’s can be preserved.
If none of the maximal subbundles actually survive in $`E^{}`$, then any further transformation at one point would do. Otherwise clearly for a generic $`QX`$ we have $`F(Q)V(Q)`$ for all the $`F`$’s that are preserved and so we can choose a hyperplane $`V(Q)HE(Q)`$ such that $`F(Q)H`$ for any such $`F`$. The elementary transformation of $`E^{}`$ at $`Q`$ corresponding to this hyperplane satisfies then the required property. ∎
###### Remark 2.9.
(1) It can definitely happen that all the maximal subbundles are killed by elementary transformations of length less than $`rk+1`$ which preserve $`V`$. In any case, as it was already suggested in the proof above, by further elementary transforming we obviously don’t change the property that we are interested in, so $`rk+1`$ is a bound that works in all situations.
(2) By Lemma 2.7, the set $`𝒵_V`$ of all elementary transformations of length $`m`$ preserving $`V`$ is irreducible of dimension $`m(rk)`$. On the other hand the condition of preserving at least one maximal subbundle is closed, so once the lemma above is true for one elementary transformation, it applies for an open subset of $`𝒵_V`$.
Finally we have all the ingredients necessary to prove the theorem. To simplify the formulations, it is convenient to introduce the following ad-hoc definition:
###### Definition 2.10.
An irreducible component $`𝒬\mathrm{Quot}_{rk,ed}(E)`$ is called *non-special* if its generic point corresponds to a locally free quotient of $`E`$ and *special* if all its points correspond to non-locally free quotients. For any $`𝒬`$, denote by $`𝒬_0`$ the open subset parametrizing locally free quotients and consider $`\mathrm{Quot}_{rk,ed}^0(E):=\underset{𝒬}{}𝒬_0`$.
###### Proof.
(of 2.2) Denote by $`𝒬`$ an irreducible component of $`\mathrm{Quot}_{rk,ed}(E)`$ (recall that we are thinking now of this Hilbert scheme as parametrizing subsheaves of rank $`k`$ and degree $`d`$). The first step is to observe that it is enough to prove the statement when $`𝒬`$ is non-special. To see this, note that every nonsaturated subsheaf $`FE`$ determines a diagram:
where $`F^{}`$ is the saturation of $`F`$, $`G`$ is a quotient vector bundle and $`\tau `$, the torsion subsheaf of $`G^{}`$, is a nontrivial zero-dimensional subscheme, say of length $`a`$. We can stratify the set of all such $`F`$’s according to the value of the parameter $`a`$, which obviously runs over a finite set. If we denote by $`\{F\}_a`$ the subset corresponding to a fixed $`a`$, this gives then:
$$\mathrm{dim}\{F\}_a\mathrm{dim}\mathrm{Quot}_{rk,eda}^0(E)+ka.$$
The right hand side is clearly less than $`k(rk)+(d_kd)k(rk+1)`$ if we assume that the statement of the theorem holds for $`\mathrm{Quot}_{rk,eda}^0(E)`$.
Let us then restrict to the case when $`𝒬`$ is a non-special component. The proof goes by induction on $`d_kd`$. If $`d_k=d`$, the statement is exactly the content of 2.1. Assume now that $`d_k>d`$ and that the statement holds for all the pairs where this difference is smaller. Recall that $`𝒬_0𝒬`$ denotes the open subset corresponding to vector subbundles and fix $`V𝒬_0`$. Then by Proposition 2.8, there exists an elementary transformation
$$0E^{}E\tau 0$$
of length $`rk+1`$, such that $`VE^{}`$, but $`FE^{}`$ for any $`FM_k(E)`$. Then necessarily $`d_k(E^{})<d_k(E)=d_k`$ (consider the saturation in $`E`$ of a maximal subbundle of $`E^{}`$) and so $`d_k(E^{})d<d_kd`$. This means that we can apply the inductive hypothesis for any non-special component of the set of subsheaves of rank $`k`$ and degree $`d`$ of $`E^{}`$. To this end, consider the correspondence:
By Lemma 2.7 and Remark 2.9(b), for any $`V𝒬_0`$, the fiber $`p_1^1(V)`$ is a (quasi-projective irreducible) variety of dimension $`(rk+1)(rk)`$ and so:
(1)
$$\mathrm{dim}𝒲=\mathrm{dim}𝒬_0+(rk+1)(rk).$$
On the other hand, for $`E^{}\mathrm{Im}(p_2)`$, the inductive hypothesis implies that
$$\mathrm{dim}p_2^1(E^{})k(rk)+(d_k(E^{})d)k(rk+1)$$
$$k(rk)+(d_kd1)k(rk+1)$$
and since $`\mathrm{dim}Y_{rk+1}=r(rk+1)`$ we have:
(2)
$$\mathrm{dim}𝒲r(rk+1)+k(rk)+(d_kd1)k(rk+1).$$
Combining (1) and (2) we get:
$$\mathrm{dim}𝒬_0k(rk)+(d_kd)k(rk+1)$$
and of course the same holds for $`𝒬=\overline{𝒬_0}`$. This completes the proof. ∎
The formulation and the proof of the theorem give rise to a few natural questions and we address them in the following examples.
###### Example 2.11.
It is easy to construct special components of Hilbert schemes. For example consider for any $`X`$ the Hilbert scheme of quotients of $`𝒪_X^2`$ of rank $`1`$ and degree $`1`$. There certainly exist such quotients which have torsion, like
$$𝒪_X^2𝒪_X𝒪_P0,$$
where $`P`$ is any point of $`X`$, but for obvious cohomological reasons there can be no sequence of the form
$$0L^1𝒪_X^2L0$$
with $`\mathrm{deg}(L)=1`$. So in this case there are actually no non-special components.
###### Example 2.12.
Going one step further, there may exist special components whose dimension is greater than that of any of the non-special ones. Note though that the proof shows that in this case that the bound cannot be optimal. To see an example, consider quotients of $`𝒪_X^2`$ of rank $`1`$ and degree $`1dg2`$ on a nonhyperelliptic curve $`X`$. Any such locally free quotient $`L`$ gives a sequence:
$$0L^1𝒪_X^2L0$$
and so the dimension of any component of the Hilbert scheme containing it is bounded above by $`h^0(L^2)`$. Now Clifford’s theorem says that $`h^0(L^2)d+1`$, but our choices make the equality case impossible, so in fact $`h^0(L^2)d`$.
On the other hand consider an effective divisor $`D`$ of degree $`d`$. Then a point in the same Hilbert scheme is determined by a natural sequence:
$$0𝒪_X(D)𝒪_X^2𝒪_X𝒪_D0$$
and it is not hard to see that the dimension of the Hilbert scheme at this point is equal to $`d+1`$ (essentially $`d`$ parameters come from $`D`$ and one from the sections of $`𝒪_X^2`$). This gives then a special component whose dimension is greater than that of any non-special one.
###### Example 2.13.
More significantly, the bound given in the theorem is optimal. Consider for this a line bundle $`L`$ of degree $`4`$ on a curve $`X`$ of genus $`2`$ and a generic extension:
$$0𝒪_XEL0.$$
By standard arguments one can see that such an extension must be stable. Since $`\mu (E)=2`$, by the classical theorem of Nagata we get that $`d_1(E):=\underset{ME}{\mathrm{max}}\mathrm{deg}(M)=1`$ and so for the sequence above $`d_1d=1`$. The theorem then tells us that the dimension of any component of the Hilbert scheme containing the given quotient is bounded above by $`3`$. But on the other hand $`h^1(L)=0`$, so this gives a smooth point and the dimension of the component is $`h^0(L)`$, which by Riemann-Roch is exactly $`3`$.
This example turns out to be a special case of a general pattern, as suggested by M. Teixidor. In fact in it is shown that whenever $`E`$ is a generic stable bundle, the invariant $`d_k`$ is the largest integer $`d`$ that makes the expression $`kerdk(rk)(g1)`$ nonnegative (cf. also ). Also the dimension of the Hilbert scheme can be computed exactly in this case (see 0.2), and for example under the numerical assumptions above it is precisely equal to $`3`$. Thus in fact for every generic stable bundle of rank $`2`$ and degree $`4`$ on a curve of genus $`2`$, we have equality in the theorem. Much more generally, it can be seen analogously that for any $`r`$ and $`g`$ equality is satisfied for a generic stable bundle as long as $`d_1`$ satisfies a certain numerical condition.
The proof of the theorem given above can be slightly modified towards a more natural and compact form. We chose to follow the longer approach because it emphasizes very clearly what is the phenomenon involved, but below we would also like to briefly sketch this parallel argument, which grew out of conversations with I. Coandă.
We will use the same notations as above. There exists a natural specialization map:
$$X\times \mathrm{Quot}_{rk,ed}^0(E)𝐆_{rk}(E),$$
where $`𝐆_{rk}(E)`$ is the Grassmann bundle of $`rk`$ dimensional quotients of the fibers of $`E`$. Of course in the case $`d=d_k`$, $`\mathrm{Quot}_{rk,ed}^0(E)`$ is compact and the morphism above is finite. Fix now $`PX`$ and $`w𝐆_{rk}(E(P))`$ a point corresponding to a quotient $`E(P)W0`$. The choice of $`P`$ determines a map:
$$\varphi :\mathrm{Quot}_{rk,ed}^0(E)𝐆_{rk}(E(P))$$
and we would like to bound the dimension of $`\varphi ^1(w)`$. There is a natural induced sequence:
$$0FEW𝐂_P0$$
and it is not hard to see that $`\varphi ^1(w)`$ embeds as an open subset in $`\mathrm{Quot}_{rk,edk}^0(F)`$. Every locally free quotient of $`F`$ has degree $`ed_kk`$, and there are at most a finite number of quotients having precisely this degree (they come exactly from the minimal degree quotients of $`E`$ having fixed fiber $`W`$ at $`P`$). Let $`G_1,\mathrm{},G_m`$ be these quotients, sitting in exact sequences:
$$0F_iFG_i0.$$
The variety $`Y:=𝐏F\underset{i=1}{\overset{𝑚}{}}𝐏G_i`$ parametrizes then the one-point elementary transformations of $`F`$ that do not preserve any of the $`F_i`$’s. Consider the natural incidence
$$𝒵\mathrm{Quot}_{rk,edk}^0(F)\times Y$$
parametrizing the pairs $`(FQ0,F^{})`$, where $`F^{}`$ is an elementary transformation in $`Y`$ and $`Q`$ is not preserved as a quotient of $`F^{}`$ (in other words the corresponding kernel is preserved). The fiber of $`𝒵`$ over $`FQ0`$ is isomorphic to $`𝐏QY`$ and so
$$\mathrm{dim}𝒵=\mathrm{dim}\mathrm{Quot}_{rk,edk}^0(F)+rk.$$
On the other hand the fiber of $`𝒵`$ over $`F^{}Y`$ is $`\mathrm{Quot}_{rk,edk1}^0(F^{})`$ Now for $`F^{}`$ the minimal degree of a quotient of rank $`rk`$ is smaller, hence inductively as before:
$$\mathrm{dim}\mathrm{Quot}_{rk,edk1}^0(F^{})k(rk)+(d_kd1)k(rk+1).$$
This immediately implies that
$$\mathrm{dim}\mathrm{Quot}_{rk,edk}^0(F)k(rk)+(d_kd1)k(rk+1)+k.$$
As this consequently holds for every fiber of the map $`\varphi `$, we conclude that
$$\mathrm{dim}\mathrm{Quot}_{rk,ed}^0(E)k(rk)+(d_kd)k(rk+1),$$
which finishes the proof.
## 3. Warm up for effective base point freeness: the case of $`SU_X(2)`$
In this section we give a very simple proof of a theorem which first appeared in (see also ). It completely takes care of the case of $`SU_X(2)`$. Although the specific technique (based on the Clifford theorem for line bundles) is different from the methods that will be used in Section 4 to prove the main result, the general computational idea already appears here, in a particularly transparent form. This is the reason for including the proof.
###### Theorem 3.1.
The linear system $`||`$ on $`SU_X(2)`$ has no base points.
###### Proof.
Recall from 1.1 and 1.2 that the statement of the theorem is equivalent to the following fact: for any stable bundle $`ESU_X(2)`$, there exists a line bundle $`L\mathrm{Pic}^{g1}(X)`$ such that $`H^0(EL)=0`$. This is certainly an open condition and it is sufficient to prove that the algebraic set
$$\{L\mathrm{Pic}^{g1}(X)|H^0(EL)0\}\mathrm{Pic}^{g1}(X)$$
has dimension strictly less than $`g`$.
A nonzero map $`E^{}L`$ comes together with a diagram of the form:
where $`M`$ is just the image in $`L`$. Then we have $`M=L(D)`$ for some effective divisor $`D`$. Since $`E`$ is stable, the degree of $`M`$ can vary from $`1`$ to $`g1`$ and we want to count all these cases separately. So for $`m=1,\mathrm{},g1`$, consider the following algebraic subsets of $`\mathrm{Pic}^{g1}(X)`$:
$$A_m:=\{L\mathrm{Pic}^{g1}(X)|0\varphi :E^{}L\mathrm{with}M=\mathrm{Im}(\varphi ),\mathrm{deg}(M)=m\}.$$
The claim is that $`\mathrm{dim}A_mg1`$ for all such $`m`$. Then of course
$$A_1\mathrm{}A_{g1}\mathrm{Pic}^{g1}(X)$$
and any $`L`$ outside this union satisfies our requirement. To prove the claim, denote by $`\mathrm{Quot}_{1,m}(E)`$ the Hilbert scheme of coherent quotients of $`E`$ of rank $`1`$ and degree $`m`$. The set of line bundle quotients $`E^{}M0`$ of degree $`m`$ is a subset of $`\mathrm{Quot}_{1,m}(E)`$. On the other hand every $`LA_m`$ can be written as $`L=M(D)`$, with $`M`$ as above and $`D`$ effective of degree $`g1m`$. This gives the obvious bound:
$$\mathrm{dim}A_m\mathrm{dim}\mathrm{Quot}_{1,m}(E)+g1m.$$
To bound the dimension of the Hilbert scheme in question, fix an $`M`$ as before and consider the exact sequence that it determines:
$$0M^{}EM0.$$
Note that the kernel is isomorphic to $`M^{}`$ since $`E`$ has trivial determinant. Now we use the well known fact from deformation theory that $`\mathrm{dim}\mathrm{Quot}_{1,m}(E)h^0(M^2)`$. To estimate $`h^0(M^2)`$, one uses all the information provided by Clifford’s theorem. The initial bound that it gives is $`h^0(M^2)m+1`$ (note that $`\mathrm{deg}(M)g1`$). If actually $`h^0(M^2)m`$, then we immediately get $`\mathrm{dim}A_mg1`$ as required. On the other hand if $`h^0(M^2)=m+1`$, by the equality case in Clifford’s theorem (see e.g. III, §1) one of the following must hold: $`M^2𝒪_X`$ or $`M^2\omega _X`$ or $`X`$ is hyperelliptic and $`M^2mg_2^1`$. The first case is impossible since $`\mathrm{deg}(M)>0`$. In the second case $`M`$ is a theta characteristic and we are done either by the fact that these are a finite number or by other overlapping cases. The third case can also happen only for a finite number of $`M`$’s and if we’re not in any of the other cases then of course $`\mathrm{dim}A_mg1m<g1`$. This concludes the proof of the theorem. ∎
###### Remark 3.2.
Note that the key point in the proof above is the ability to give a convenient upper bound on the dimension of certain Hilbert schemes. This will essentially be the main ingredient in the general result proved in Section 4, and the needed estimate was provided in the previous section.
## 4. Base point freeness for pluritheta linear series on $`SU_X(r,e)`$
Using the dimension bound given in Section 2, we are now able to prove the main result of this paper, namely an effective base point freeness bound for pluritheta linear series on $`SU_X(r,e)`$. The proof is computational in nature and the roots of the main technique involved have already appeared in 3.1. Let $`r2`$ and $`e`$ be arbitrary integers and let $`h=\mathrm{gcd}(r,e),r=r_1h`$ and $`e=e_1h`$. For the statement it is convenient to introduce another invariant of the moduli space. If $`ESU_X(r,e)`$ and $`1kr1`$, define $`s_k(E):=kerd_k`$, where $`d_k`$ is the maximum degree of a subbundle of $`E`$ of rank $`k`$ (cf.). Note that if $`E`$ is stable one has $`s_k(E)h`$ and we can further define $`s_k=s_k(r,e):=\underset{E\mathrm{stable}}{\mathrm{min}}s_k(E)`$ and $`s=s(r,e):=\underset{1kr1}{\mathrm{min}}s_k`$. Clearly $`sh`$ and it is also an immediate observation that $`s(r,e)=s(r,e)`$.
###### Theorem 4.1.
The linear series $`|^p|`$ on $`SU_X(r,e)`$ is base point free for
$$p\mathrm{max}\{\frac{(r+1)^2}{4r}h,\frac{r^2}{4s}h\}.$$
###### Remark 4.2.
Note that the bound given in the theorem is always either a quadratic or a linear function in the rank $`r`$. It should also be said right away that although this bound works uniformly, in almost any particular situation one can do a little better. Unfortunately there doesn’t seem to be a better uniform way to express it, but we will comment more on this at the end of the section (cf. Remark 4.5).
###### Proof.
(of 4.1) Let us denote for simplicity $`M:=\mathrm{max}\{\frac{(r+1)^2}{4r}h,\frac{r^2}{4s}h\}`$. Since the problem depends only on the residues of $`e`$ modulo $`r`$, there is no loss of generality in looking only at $`SU_X(r,e)`$ with $`0er1`$. The statement of the theorem is implied by the following assertion, as described in 1.1 and 1.2:
$$ForanystableESU_X(r,e)andanypM,thereisan$$
$$FU_X(pr_1,p(r_1(g1)+e_1))suchthath^0(EF)=0.$$
Fix a stable bundle $`ESU_X(r,e)`$. Note that only in this proof, as opposed to the rest of the paper, $`e`$ in fact denotes the degree of $`E^{}`$, and not that of $`E`$. If for some $`FU_X(pr_1,p(r_1(g1)+e_1))`$ there is a nonzero map $`E^{}\stackrel{\mathit{\varphi }}{}F`$, then this comes together with a diagram of the form:
where the vector bundle $`V`$ is the image of $`\varphi `$. The idea is essentially to count all such diagrams assuming that the rank and degree of $`V`$ are fixed and see that the $`F`$’s involved in at least one of them cover only a proper subset of the whole moduli space. Denote as before by $`\mathrm{Quot}_{k,d}(E^{})`$ the Hilbert scheme of quotients of $`E^{}`$ of rank $`k`$ and degree $`d`$ and for any $`1kr`$ and any $`d`$ in the suitable range (given by the stability of $`E`$ and $`F`$) consider its subset:
$$A_{k,d}:=\{V\mathrm{Quot}_{k,d}(E^{})|FU_X(pr_1,p(r_1(g1)+e_1)),0\varphi :E^{}F\mathrm{with}V=\mathrm{Im}(\varphi )\}.$$
The theorem on Hilbert schemes stated in the introduction then gives us the dimension estimate:
(3)
$$\mathrm{dim}A_{k,d}k(rk)+(df_k)(k+1)(rk),$$
where $`f_k=f_k(E^{})`$ is the minimum possible degree of a quotient bundle of $`E^{}`$ of rank $`k`$ (which is the same as $`d_k`$). Define now the following subsets of $`U_X(pr_1,p(r_1(g1)+e_1))`$:
$$U_{k,d}:=\{F|VA_{k,d}\mathrm{with}VF\}U_X(pr_1,p(r_1(g1)+e_1)).$$
The elements of $`U_{k,d}`$ are all the $`F`$’s that appear in diagrams as above for fixed $`k`$ and $`d`$. The claim is that
$$\mathrm{dim}U_{k,d}<(pr_1)^2(g1)+1,$$
which would imply that $`U_{k,d}U_X(pr_1,p(r_1(g1)+e_1))`$. Assuming that this is true, and since $`k`$ and $`d`$ run over a finite set, any $`FU_X(pr_1,p(r_1(g1)+e_1))\underset{k,d}{}U_{k,d}`$ satisfies the desired property that $`h^0(EF)=0`$, which gives the statement of the theorem. It is easy to see, and in fact a particular case of the computation below, that in the case $`k=r`$ (i.e. $`V=E^{}`$) $`U_{k,d}`$ has dimension exactly $`(pr_1)^2(g1)`$.
Let us concentrate then on proving the claim above for $`1kr1`$. Note that the inclusions $`VF`$ appearing in the definition of $`U_{k,d}`$ are valid in general only at the sheaf level. Any such inclusion determines an exact sequence:
(4)
$$0VFG^{}0,$$
where $`G^{}=G\tau _a`$, with $`G`$ locally free and $`\tau _a`$ a zero dimensional subscheme of length $`a`$. We stratify $`U_{k,d}`$ by the subsets
$$U_{k,d}^a:=\{F|F\mathrm{given}\mathrm{by}\mathrm{an}\mathrm{extension}\mathrm{of}\mathrm{type}(4)\}U_{k,d},$$
where $`a`$ runs over the obvious allowable finite set of integers. A simple computation shows that $`G`$ has rank $`pr_1k`$ and degree $`p(r_1(g1)+e_1)da`$. Denote by $`T_{k,d}^a`$ the set of all vector bundles $`G`$ that are quotients of some $`FU_X(pr_1,p(r_1(g1)+e_1))`$. These can be parametrized by a relative Hilbert scheme (see e.g. §8.6) over (an étale cover of) $`U_X(pr_1,p(r_1(g1)+e_1))`$ and so they form a bounded family. We invoke a general result, proved in 4.1 and 4.2, saying that the dimension of such a family is always at most what we get if we assume that the generic member is stable. Thus we get the bound:
$$\mathrm{dim}T_{k,d}^a(pr_1k)^2(g1)+1.$$
Now we only have to compute the dimension of the family of all possible extensions of the form (4) when $`V`$ and $`G`$ are allowed to vary over $`A_{k,d}`$ and $`T_{k,d}^a`$ respectively and $`\tau _a`$ varies over the symmetric product $`X_a`$. Any such extension induces a diagram
If we denote by $`A_{k,d}^a`$ the set of isomorphism classes of vector bundles $`V^{}`$ that are (inverse) elementary transformations of length $`a`$ of vector bundles in $`A_{k,d}`$, then we have the obvious:
$$\mathrm{dim}A_{k,d}^a\mathrm{dim}A_{k,d}+ka.$$
On the other hand any $`F`$ is obtained as an extension of a bundle in $`T_{k,d}^a`$ by a bundle in $`A_{k,d}^a`$. Denote by $`𝒰A_{k,d}^a\times T_{k,d}^a`$ the open subset consisting of pairs $`(V^{},G)`$ such that there exists an extension
$$0V^{}FG0$$
with $`F`$ stable. Note that by Lemma 1.3 for any such pair we have $`h^0(G^{}V^{})=0`$ and so by Riemann-Roch $`h^1(G^{}V^{})`$ is constant, given by:
(5)
$$h^1(G^{}V^{})=2kpr_1(g1)k^2(g1)+kpe_1pr_1dpr_1a.$$
In this situation it is a well known result (see e.g. (2.4) or §4) that there exists a universal space of extension classes $`𝐏(𝒰)𝒰`$ whose dimension is computed by the formula:
$$\mathrm{dim}𝐏(𝒰)=\mathrm{dim}A_{k,d}^a+\mathrm{dim}T_{k,d}^a+h^1(G^{}V^{})1.$$
There is an obvious forgetful map:
$$𝐏(𝒰)U_X(pr_1,p(r_1(g1)+e_1))$$
whose image is exactly $`U_{k,d}^a`$. Thus by putting together all the inequalities above we obtain:
$$\mathrm{dim}U_{k,d}^a\mathrm{dim}A_{k,d}^a+\mathrm{dim}T_{k,d}^a+h^1(G^{}V^{})1$$
$$k(rk)+(df_k)(k+1)(rk)+ka+(pr_1k)^2(g1)+1$$
$$+2kpr_1(g1)k^2(g1)+kpe_1pr_1dpr_1a1$$
$$k(rk)+(df_k)(k+1)(rk)+(pr_1)^2(g1)+kpe_1pr_1d+kapr_1a$$
$$k(rk)+(df_k)(k+1)(rk)+(pr_1)^2(g1)+kpe_1pr_1d,$$
where the last inequality is due to the obvious fact that $`kr1<pr_1`$ if $`pM`$. Since $`a`$ runs over a finite set, to conclude the proof of the claim it is enough to see that $`\mathrm{dim}U_{k,d}^a(pr_1)^2(g1)`$. By the inequality above this is true if
$$p(r_1dke_1)k(rk)+(df_k)(k+1)(rk),$$
or equivalently if
$$p(rdke)k(rk)h+(df_k)(k+1)(rk)h$$
for any $`k`$ and $`d`$. This can be rewritten in the following more manageable form:
$$p(r(df_k)+rf_kke)k(rk)h+(df_k)(k+1)(rk)h.$$
The first case to look at is $`d=f_k`$, when we should have $`p(rf_kke)k(rk)h`$ and this should hold for every $`k`$. But clearly $`rf_kke=s_{rk}(E^{})=s_k(E)`$, defined above in terms of maximal subbundles, and $`h|s_k(E)`$. Since $`E`$ is stable we then have $`s_k(E)h`$ for all $`k`$ and so $`sh`$ as mentioned before. Note though that in general one cannot do better (cf. Remark 4.6). In any case, this says that the inequality $`p\frac{r^2}{4s}h`$ must be satisfied (which would certainly hold if $`p\frac{r^2}{4}`$).
When $`d>f_k`$, it is convenient to collect together all the terms containing $`df_k`$. The last inequality above then reads:
$$(df_k)(pr(k+1)(rk)h)+ps_k(E)k(rk)h.$$
For $`p`$ as before it is then sufficient to have $`pr(k+1)(rk)h`$, which again by simple optimization is satisfied for $`p\frac{(r+1)^2}{4r}h`$. Concluding, the desired inequality holds as long as $`pM`$. ∎
The most important instances of this theorem are the cases of vector bundles of degree $`0`$ (more generally $`d0(\mathrm{mod}r)`$) and degree $`1`$ or $`1`$ (more generally $`d\pm 1(\mathrm{mod}r)`$). In the second situation the moduli space in question is smooth. It is somewhat surprising that the results obtained in these cases have different orders of magnitude.
###### Corollary 4.3.
$`|^p|`$ is base point free on $`SU_X(r)`$ for $`p\frac{(r+1)^2}{4}`$.
###### Proof.
This is clear since $`h=r`$. ∎
###### Corollary 4.4.
$`|^p|`$ is base point free on $`SU_X(r,1)`$ and $`SU_X(r,1)`$ for $`pr1`$.
###### Proof.
Note that by duality it suffices to prove the claim for one of the moduli spaces, say $`SU_X(r,1)`$. In this case $`h=1`$ and for any $`ESU_X(r,1)`$, $`s_{rk}(E^{})=rf_kkrk`$. Following the proof of the theorem we see thus that it suffices to have
$$p\mathrm{max}\{r1,\frac{(r+1)^2}{4r}\},$$
which is equal to $`r1`$ for $`r3`$. For $`r=2`$ one can slightly improve the last inequality in the proof of the theorem (actually this is true whenever $`r`$ is even) to see that $`p=1`$ already works. ∎
###### Remark 4.5.
Corollary 4.4 can be seen as a generalization of the well-known fact that $`||`$ is base point free on $`SU_X(2,1)`$. Also, as already noted in its proof, the general bound obtained in the theorem can be slightly improved in each particular case, due to the fact that the two optimization problems do not simultaneously have integral solutions. Thus for example if $`r`$ is even, the proof of the theorem actually gives that $`|^p|`$ is base point free on $`SU_X(r)`$ for $`pr(r+2)/4`$.
###### Remark 4.6.
As already noted, in any given numerical situation the bound given by the theorem is either linear or quadratic in the rank $`r`$. One may thus hope that at least in the case $`h=1`$ (i.e. $`r`$ and $`e`$ coprime), a closer study of the number $`s(r,e)`$ might always produce by this method a linear bound. Examples show though that this is not the case: one can take $`r=4l`$, $`e=2l1`$, $`k=2l+1`$ and $`f_k=l`$ (this works for special vector bundles) for a positive integer $`l`$, which implies $`s=s_k=1`$.
## 5. Effective base point freeness on $`U_X(r,e)`$ and some conjectures
The deformation theoretic methods used in and allow one to prove results similar to Theorem 4.1 for pluritheta linear series on $`U_X(r,e)`$ (with some extra effort due to the fact that in this case one has to control the determinant of the complementary vector bundle). Since the method used in this paper is of a different nature, a generalization along those lines is not immediately apparent. Instead we propose the formalism of Verlinde bundles, which we developed in . This comes with the advantage that it applies automatically as soon as one has results on $`SU_X(r,e)`$ and also suggests what the optimal bounds should be. Moreover, the method equally applies to other linear series on $`U_X(r,e)`$, as we will see shortly.
Fix a generic vector bundle $`FU_X(r_1,r_1(g1)e_1)`$, where as usual $`r_1=r/h`$ and $`e_1=e/h`$. To it we can associate the theta divisor $`\mathrm{\Theta }_F`$ on $`U_X(r,e)`$ supported on the set
$$\mathrm{\Theta }_F=\{EU_X(r,e)|h^0(EF)0\}.$$
Fix also $`L\mathrm{Pic}^e(X)`$. The $`(r,e,k)`$-*Verlinde bundle* $`E_{r,e,k}`$ associated to these choices is defined as (cf. §6):
$$E_{r,e,k}(=E_{r,e,k}^{F,L}):=\pi _{L}^{}{}_{}{}^{}𝒪_U(k\mathrm{\Theta }_F),$$
where $`\pi _L`$ is the composition:
$$\pi _L:U_X(r,e)\stackrel{det}{}\mathrm{Pic}^e(X)\stackrel{L^1}{}J(X).$$
This is a vector bundle on $`J(X)`$ of rank equal to the Verlinde number $`h^0(SU_X(r,e),^k)`$. The following results are proved in :
###### Theorem 5.1.
( 6.4 and 5.3) $`𝒪_U(k\mathrm{\Theta }_F)`$ is globally generated on $`U_X(r,e)`$ as long as $`^k`$ is globally generated on $`SU_X(r,e)`$ and $`E_{r,e,k}`$ is globally generated. Moreover, $`𝒪_U(k\mathrm{\Theta }_F)`$ is not globally generated for $`kh`$.
###### Proposition 5.2.
( 5.2) $`E_{r,e,k}`$ is globally generated for $`kh+1`$ and this bound is optimal.
We immediately obtain by using 4.1 the following bound, where $`s`$ is the invariant defined in the previous section:
###### Theorem 5.3.
$`𝒪_U(k\mathrm{\Theta }_F)`$ is globally generated on $`U_X(r,e)`$ for $`k\mathrm{max}\{\frac{(r+1)^2}{4r}h,\frac{r^2}{4s}h\}`$.
In fact the theorem is a special case of the more general 5.9 that we will treat at the end of the section. Right now it is interesting to see how these bounds relate to possible optimal bounds and discuss some conjectures and questions in this direction. Given the shape of the result, we will carry out this discussion in the case of $`SU_X(r)`$ and $`SU_X(r,\pm 1)`$, based on the results 4.3 and 4.4. A similar analysis can be applied to any other case, but we will not give any details here.
We begin by looking at degree $`0`$ vector bundles, where global generation is attained for $`k\frac{(r+1)^2}{4}`$, with the improvement 3.1 in the case of rank $`2`$, when $`k1`$ suffices. In view of 5.1, the bound in 5.3 is optimal in the case of rank $`2`$ and rank $`3`$ vector bundles.
###### Corollary 5.4.
Let $`N\mathrm{Pic}^{g1}(X)`$. Then:
(i) $`𝒪_U(3\mathrm{\Theta }_N)`$ is globally generated on $`U_X(2,0)`$.
(ii) $`𝒪_U(4\mathrm{\Theta }_N)`$ is globally generated on $`U_X(3,0)`$.
This could be seen as a natural extension of the classical fact that $`𝒪_J(2\mathrm{\Theta }_N)`$ is globally generated on $`J(X)U_X(1,0)`$. In presence of this evidence it is natural to conjecture that this is indeed the case for any rank:
###### Conjecture 5.5.
For any $`r1`$, $`𝒪_U(k\mathrm{\Theta }_N)`$ is globally generated on $`U_X(r,0)`$ for $`kr+1`$.
This is the best that one can hope for and there is a possibility that it might be a little too optimistic, or in other words that Corollary 5.4 might be an accident of low values of a quadratic function. On the other hand if that is the case, the theorem should be very close to being optimal. Turning to $`SU_X(r)`$, in §3 we showed that, granting the strange duality conjecture, the optimal bound for the global generation of $`^k`$ should also go up as we increase the rank $`r`$. The underlying reason (without specifying the actual numbers) is the following: assume that we are given a vector bundle $`E`$ such that $`h^0(E\xi )0`$ for all $`\xi \mathrm{Pic}^0(X)`$ (for examples see , or ). If we choose some complementary bundle $`F`$ (i.e. $`\chi (EF)=0`$), of rank $`t`$, then a theorem of Lange and Mukai-Sakai (see and ) asserts that $`F`$ admits a line subbundle of degree $`\mu (F)g+g/t=g/t1\mu (E)`$. For small $`t`$ (with respect to $`r`$), in most examples mentioned above it happens that this number is positive, which automatically implies that $`h^0(EF)0`$ for all such $`F`$. This prevents the global generation of a certain multiple of $``$ depending on the rank of $`F`$. The case of rank $`2`$ vector bundles 3.1 suggests though that we could ask for a slightly better result than for $`U_X(r,0)`$, but unfortunately further evidence is still missing:
###### Conjecture/Question 5.6.
Is $`^k`$ globally generated on $`SU_X(r)`$ for $`kr1`$?
Note also that in view of 5.1 and 5.2 any positive answer in the range $`\{r1,r,r+1\}`$ would imply the optimal conjecture 5.5.
In the case of $`SU_X(r,\pm 1)`$ 4.4 and 5.1 give that $`𝒪_U(k\mathrm{\Theta }_F)`$ is globally generated for $`k\mathrm{max}\{2,r1\}`$, while $`𝒪_U(\mathrm{\Theta }_F)`$ cannot be. We obtain thus again optimal bounds for rank $`2`$ and rank $`3`$ vector bundles.
###### Corollary 5.7.
$`𝒪_U(2\mathrm{\Theta }_F)`$ is globally generated on $`U_X(2,1)`$ and $`U_X(3,\pm 1)`$.
Note also that for all the examples of special vector bundles constructed in , and we have $`h1`$, therefore theoretically an optimal bound that does not depend on the rank $`r`$ is still possible. It is natural to ask if the best possible result always holds:
###### Question 5.8.
Is $`𝒪_U(2\mathrm{\Theta }_F)`$ on $`U_X(r,\pm 1)`$, and so also $`^2`$ on $`SU_X(r,\pm 1)`$, globally generated? More generally, is this true whenever $`r`$ and $`d`$ are coprime?
We conclude the section with a generalization of Theorem 5.3. For simplicity we present it only in the degree $`0`$ case, but the extension to other degrees is immediate. Recall from Theorem C that for $`N\mathrm{Pic}^{g1}(X)`$, $`\mathrm{Pic}(U_X(r,0))𝐙𝒪(\mathrm{\Theta }_N)\text{det}^{}\mathrm{Pic}(J(X))`$. The method provided by the Verlinde bundles allows one to study effective global generation for “mixed” line bundles of the form $`𝒪(k\mathrm{\Theta }_N)\text{det}^{}L`$ with $`L\mathrm{Pic}(J(X))`$. Concretely we have the following cohomological criterion (assume $`r2`$):
###### Theorem 5.9.
$`𝒪(k\mathrm{\Theta }_N)\mathrm{det}^{}L`$ is globally generated if $`k\frac{(r+1)^2}{4}`$ and
$$h^i(𝒪_J((krr^2)\mathrm{\Theta }_N)L^{r^2}\alpha )=0,i>0,\alpha \mathrm{Pic}^0(J(X)).$$
###### Proof.
By the projection formula, for every $`i>0`$ we have:
$$R^i\text{det}_{}(𝒪_U(k\mathrm{\Theta }_N)\text{det}^{}L)R^i\text{det}_{}𝒪_U(k\mathrm{\Theta }_N)L=0.$$
Also the restriction of $`𝒪_U(k\mathrm{\Theta }_N)\text{det}^{}L`$ to any fiber $`SU_X(r,\xi )`$ of the determinant map is isomorphic to $`^k`$ and so globally generated for $`k\frac{(r+1)^2}{4}`$. It is a simple consequence of general machinery, described in (7.1), that in these conditions the statement holds as soon as
$$\text{det}_{}(𝒪_U(k\mathrm{\Theta }_N)\text{det}^{}L)E_{r,k}L$$
is globally generated on $`J(X)`$, where $`E_{r,k}`$ is a simplified notation for $`E_{r,0,k}`$. To study this we make use, as in , of a cohomological criterion for global generation of vector bundles on abelian varieties due to Pareschi (2.1). In our particular setting it says that $`E_{r,k}L`$ is globally generated if there exists some ample line bundle $`A`$ on $`J(X)`$ such that
$$h^i(E_{r,k}LA^1\alpha )=0,i>0,\alpha \mathrm{Pic}^0(J(X)).$$
We chose $`A`$ to be $`𝒪_J(\mathrm{\Theta }_N)`$, where $`\mathrm{\Theta }_N`$ is the theta divisor on $`J(X)`$ associated to $`N`$. The cohomology vanishing that we need is true if it holds for the pullback of $`E_{r,k}L𝒪_J(\mathrm{\Theta }_N)\alpha `$ by any finite cover of $`J(X)`$. But recall from (2.3) that $`r_J^{}E_{r,k}𝒪_J(kr\mathrm{\Theta }_N)`$, where $`r_J`$ is the multiplication by $`r`$. Since $`r_J^{}LL^{r^2}`$, via pulling back by $`r_J`$ the required vanishing certainly holds if
$$h^i(𝒪_J((krr^2)\mathrm{\Theta }_N)L^{r^2}\alpha )=0,i>0,\alpha \mathrm{Pic}^0(J(X)).$$
###### Corollary 5.10.
If $`l𝐙`$, $`𝒪(k\mathrm{\Theta }_N)\mathrm{det}^{}𝒪_J(l\mathrm{\Theta }_N)`$ is globally generated for
$$k\mathrm{max}\{r+1lr,\frac{(r+1)^2}{4}\}.$$
## 6. An application to surfaces à la Le Potier
Another, in some sense algorithmic, application of the effective bound 4.1 can be given following the paper of Le Potier . By a simple use of a restriction theorem due to Flenner (cf. also §11), Le Potier shows that effective results for the determinant bundle $``$ induce effective results for the Donaldson determinant line bundles on moduli spaces of semistable sheaves on surfaces. For the appropriate definitions and basic results, the reader can consult §8. Using the uniform bound $`k(r+1)^2/4`$ that works on every moduli space $`SU_X(r,d)`$, the result can be formulated as follows:
###### Theorem 6.1.
Let $`(X,𝒪_X(1))`$ be a polarized smooth projective surface and $`L`$ a line bundle on $`X`$. Let $`M=M_X(r,L,c_2)`$ be the moduli space of semistable sheaves of rank $`r`$, fixed determinant $`L`$ and second Chern class $`c_2`$ on $`X`$ and denote $`n=\mathrm{deg}(X)=𝒪_X(1)^2`$ and $`d=n[\frac{r^2}{2}]`$. If $`𝒟`$ is the Donaldson determinant line bundle on $`M`$, then $`𝒟^p`$ is globally generated for $`pd\frac{(r+1)^2}{4}`$ divisible by $`d`$.
Note that it is not true that $`𝒟`$ is ample, which accounts for the formulation of the theorem. The significance of the map to projective space given by some multiple of $`𝒟`$ is well known. Its image is the moduli space of $`\mu `$-semistable sheaves and in the rank $`2`$ and degree $`0`$ case this is homeomorphic to the Donaldson-Uhlenbeck compactification of the moduli space of $`ASD`$-connections in gauge theory, the map realizing the transition between the Gieseker and Uhlenbeck points of view (see e.g. §8.2 for a survey). Better bounds for the global generation of the Donaldson line bundle thus limit the dimension of an ambient projective space for this moduli space. The main improvement brought by the results in the present paper comes from the fact that our result is not influenced by the genus of the curve given by Flenner’s theorem. Effectively that reduces the bound given in §3.2 by an order of four, namely from a polynomial of degree $`8`$ in the rank $`r`$ to a polynomial of degree $`4`$.
*Sketch of proof of 6.1.*(cf. 3.6) In analogy with the curve situation, given $`EM`$, the problem is to find a complementary $`1`$-dimensional sheaf $`F`$ on $`X`$ such that $`h^1(EF)=0`$. Flenner’s theorem says that there exists a smooth curve $`C`$ belonging to the linear series $`|𝒪_X(d)|`$ such that $`E_{|C}`$ is semistable. By theorem 4.1, on $`C`$ one can find for any $`k(r+1)^2/4`$ a vector bundle $`V`$ of rank $`kr_1`$ such that $`h^1(EV)=0`$. The $`F`$ that we are looking for is obtained by considering $`V`$ as a $`1`$-dimensional sheaf on $`X`$ and a simple computation shows that if $`p=dk`$, this gives the global generation of $`𝒟^p`$.
###### Remark 6.2.
Depending on the values of the invariants involved, this bound may sometimes be improved to a polynomial of degree $`3`$ in $`r`$, according to the precise statement of Theorem 4.1.
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# 1 Introduction
## 1 Introduction
A most challenging and potentially important problem is to construct the representation theory of local Lie algebras in $`N`$-dimensional spacetime. The name indicates that the generators are localized in spacetime, i.e. that the structure constants are proportional to finitely many derivatives of delta functions. This class includes the current algebra $`map(N,𝔤)`$, the diffeomorphism algebra $`diff(N)`$, as well as algebras of divergence free, Hamiltonian, or contact vector fields. When $`N=1`$, the projective representations are described by affine Kac-Moody and Virasoro algebras, respectively, but much less is known in higher dimensions. The reason for this is the fundamental observation that functions and normal ordering are incompatible except in one dimension. Typically, the classical representations of local Lie algebras are functions over spacetime, with values in some finite-dimensional vector space such as $`𝔤`$ or $`gl(N)`$ modules. A naïve strategy to construct projective representations would be to start from such functions, add canonically conjugate momenta, normal order, and hope to obtain a realization on Fock space. However, this approach only works in one dimension; in higher dimensions, new infinities arise.
In view of this fundamental incompatibility between functions and normal ordering, there are two philosophically distinct strategies. One is to keep functions and do something about normal ordering; papers following this route typically contain the keywords “further regularization”. The most ambitious program in this direction has been carried out by Mickelsson and collaborators, targeting the Mickelsson-Faddeev (MF) algebra . Although representations in an abstract sense have been reported, concrete representations (on a separable Hilbert space) seem to be missing .
The logical alternative is to keep normal ordering and do something about functions; more precisely, functions can be replaced by trajectories in the space of finite jets, which can be viewed as the coefficients of truncated Taylor expansions. This route was first entered by Moody, Eswara-Rao and Yokonoma , whereas the geometrical understanding was provided by myself . This approach immediately leads to concrete Fock representations of abelian extensions of current algebras. However, these cocycles are not of MF type, but rather of the higher-dimensional Kac-Moody type described by Kassel and rediscovered in . In particular, they involve one-chains rather than three-chains. Applied to the diffeomorphism algebra, the same method leads to the higher-dimensional Virasoro algebras of .
It is thus natural to ask if the MF algebra could also be represented using such methods. The answer appears to be negative. Cederwall et al. found a natural realization of a “classical MF algebra”, where the inhomogeneous term in the connection’s transformation law has been dropped. However, this term is not recovered by normal ordering; even worse, it spoils the Jacobi identities for the realization mentioned above. It is shown in the present paper that the classical MF algebra also admits other cocycles, but these are of Kac-Moody type. These new MF-like algebras possess lowest-energy representations, described in section 3. Moreover, they can be intertwined with the diffeomorphism algebra, but the extensions are then no longer central, since they do not commute with diffeomorphisms.
## 2 Embeddings
Let $`𝔤`$ be a semisimple finite-dimensional Lie algebra with basis $`J^a`$, totally skew structure constants $`f^{abc}`$, Killing metric $`\delta ^{ab}`$, and brackets $`[J^a,J^b]=f^{abc}J^c`$. As usual, set $`d^{abc}=\mathrm{tr}J^a\{J^b,J^c\}\mathrm{tr}J^{(a}J^bJ^{c)}`$, where the trace is evaluated in some representation and paranthesized indices are symmetrized. The following identities hold: $`f^{aed}f^{bcd}+f^{acd}f^{bed}+f^{abd}f^{ced}=0`$, $`f^{aed}d^{bcd}+f^{acd}d^{bed}+f^{abd}d^{ced}=0`$, $`f^{bac}=f^{abc}`$, $`f^{bca}=f^{abc}`$, $`d^{bac}=d^{abc}`$, and $`d^{bca}=d^{abc}`$.
Denote by $`𝒥^a(m)=\mathrm{exp}(imx)J^a`$ the generators of $`map(N,𝔤)`$, the algebra of maps from $`^N`$ to $`𝔤`$, where $`x=(x^\mu )`$ and $`m=(m_\mu )`$. Moreover, let $`A=A_\mu ^a(x)J^adx^\mu `$ be the connection one-form, with Fourier coefficients $`A_\mu ^a(m)`$. The Mickelsson-Faddeev (MF) algebra is the following Lie algebra extension of $`map(3,𝔤)`$:
$`[𝒥^a(m),𝒥^b(n)]`$ $`=`$ $`f^{abc}𝒥^c(m+n)+d^{abc}m_\mu n_\nu ϵ^{\mu \nu \rho }A_\rho ^c(m+n),`$
$`[𝒥^a(m),A_\nu ^b(n)]`$ $`=`$ $`f^{abc}A_\nu ^c(m+n)+\delta ^{ab}m_\nu \delta (m+n),`$ (1)
$`[A^\mu (m),A^b\nu (n)]`$ $`=`$ $`0,`$
where
$$\begin{array}{ccccc}\hfill ϵ^{\mu \nu \rho }& =& +1,\hfill & & \mu \nu \rho \text{ positive permutation of }123,\hfill \\ & =& 1,\hfill & & \mu \nu \rho \text{ negative permutation of }123,\hfill \\ & =& 0,\hfill & & \text{otherwise},\hfill \end{array}$$
(2)
is the totally anti-symmetric symbol in three dimensions.
Set $`^{a\mu \nu }(m)=ϵ^{\mu \nu \rho }A_\rho ^a(m)`$. Then the MF algebra takes the form
$`[𝒥^a(m),𝒥^b(n)]`$ $`=`$ $`f^{abc}𝒥^c(m+n)+d^{abc}m_\mu n_\nu ^{c\mu \nu }(m+n),`$
$`[𝒥^a(m),^{b\mu \nu }(n)]`$ $`=`$ $`f^{abc}^{c\mu \nu }(m+n)+\delta ^{ab}m_\rho S_3^{\mu \nu \rho }(m+n),`$
and all other brackets vanish. Moreover, $`S_3^{\mu \nu \rho }`$ is totally antisymmetric and subject to the additional condition
$$m_\rho S_3^{\mu \nu \rho }(m)0,$$
(4)
which geometrically means that it is a closed three-chain. In three dimensions, the unique solution to (4) is $`S_3^{\mu \nu \rho }(m)=ϵ^{\mu \nu \rho }\delta (m)`$, but the present formulation holds in any number of dimensions $`3`$, and in two dimensions with $`S_3^{\mu \nu \rho }(m)0`$.
Eq. (LABEL:MF2) can be embedded in the following algebra:
$`[𝒥^a(m),𝒥^b(n)]`$ $`=`$ $`f^{abc}𝒥^c(m+n),`$
$`[𝒥^a(m),𝒢^{b\mu }(n)]`$ $`=`$ $`f^{abc}𝒢^{c\mu }(m+n),`$
$`[𝒥^a(m),^{b\mu \nu }(n)]`$ $`=`$ $`f^{abc}^{c\mu \nu }(m+n)+\delta ^{ab}m_\rho S_3^{\mu \nu \rho }(m+n),`$
$`[𝒢^{a\mu }(m),𝒢^{b\nu }(n)]`$ $`=`$ $`d^{abc}^{c\mu \nu }(m+n),`$
and all other brackets vanish. Explicitly, this is accomplished by means of the redefinition
$$𝒥^a(m)𝒥^a(m)+m_\mu 𝒢^{a\mu }(m).$$
(6)
In the absence of the closed three-chain $`S_3^{\mu \nu \rho }`$, this embedding was first described by .
However, there is one big problem with (LABEL:emb1): in the presence of the three-chain, it is not a Lie algebra, because the following Jacobi identity fails:
$$[𝒥^a(m),[𝒢^{b\mu }(n),𝒢^{c\nu }(r)]]+\text{cycl.}=d^{abc}m_\rho S_3^{\mu \nu \rho }(m+n+r)0.$$
(7)
I consider this as a strong indication that the MF algebra has no Fock representations with $`S_3^{\mu \nu \rho }(m)`$ non-zero. At least, it can not be possible to isolate operators $`𝒢^{a\mu }(m)`$ as in (6), since that would violate the Jacobi identities.
If we skip the three-chain, we obtain the “classical MF algebra”, which has another abelian extension:
$`[𝒥^a(m),𝒥^b(n)]`$ $`=`$ $`f^{abc}𝒥^c(m+n)k\delta ^{ab}m_\rho S_1^\rho (m+n)+`$
$`+d^{abc}m_\mu n_\nu ^{c\mu \nu }(m+n),`$
$`[𝒥^a(m),^{b\mu \nu }(n)]`$ $`=`$ $`f^{abc}^{c\mu \nu }(m+n),`$
where $`S_1^\rho (m)`$ is a closed one-form, satisfying
$`[𝒥^a(m),S_1^\rho (n)]`$ $`=`$ $`[^{a\mu \nu }(m),S_1^\rho (n)]=0,`$
$`m_\rho S_1^\rho (m)`$ $``$ $`0.`$
This algebra can be embedded into the following algebra by means of the same redefinition (6).
$`[𝒥^a(m),𝒥^b(n)]`$ $`=`$ $`f^{abc}𝒥^c(m+n)k\delta ^{ab}m_\rho S_1^\rho (m+n),`$
$`[𝒥^a(m),𝒢^{b\mu }(n)]`$ $`=`$ $`f^{abc}𝒢^{c\mu }(m+n),`$
$`[𝒥^a(m),^{b\mu \nu }(n)]`$ $`=`$ $`f^{abc}^{c\mu \nu }(m+n),`$
$`[𝒢^{a\mu }(m),𝒢^{b\nu }(n)]`$ $`=`$ $`d^{abc}^{c\mu \nu }(m+n),`$
where $`𝒢^{a\mu }(m)`$ also commutes with the one-chain $`S_1^\rho `$.
## 3 Fock representations
Consider the following Lie algebra:
$`[J^a(s),J^b(t)]`$ $`=`$ $`f^{abc}J^c(s)\delta (st)+{\displaystyle \frac{k}{2\pi i}}\delta ^{ab}\dot{\delta }(st),`$
$`[J^a(s),G^{b\nu }(t)]`$ $`=`$ $`f^{abc}G^{c\nu }(s)\delta (st),`$
$`[G^{a\mu }(s),G^{b\nu }(t)]`$ $`=`$ $`d^{abc}H^{c\mu \nu }(s)\delta (st),`$ (11)
$`[J^a(s),H^{b\mu \nu }(t)]`$ $`=`$ $`f^{abc}H^{c\mu \nu }(s)\delta (st),`$
$`[G^{a\mu }(s),H^{b\nu \rho }(t)]`$ $`=`$ $`[H^{a\mu \nu }(s),H^{b\sigma \tau }(t)]=0,`$
where $`s,tS^1`$. Note that the first relation is the affine Kac-Moody algebra with central charge $`k`$. Moreover, introduce $`N`$ bosonic oscillators $`q^\mu (t)`$. Then the following expressions yield a realization of (LABEL:emb2) and thus of the modified MF algebra (LABEL:MF3).
$`𝒥^a(m)`$ $`=`$ $`{\displaystyle 𝑑t\mathrm{e}^{imq(t)}J^a(t)},`$
$`𝒢^{a\mu }(m)`$ $`=`$ $`{\displaystyle 𝑑t\mathrm{e}^{imq(t)}G^{a\mu }(t)},`$
$`^{a\mu \nu }(m)`$ $`=`$ $`{\displaystyle 𝑑t\mathrm{e}^{imq(t)}H^{a\mu \nu }(t)},`$
$`S_1^\mu (m)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle 𝑑t\dot{q^\mu }(t)\mathrm{e}^{imq(t)}}.`$
Moreover, the value of the central charge $`k`$ is the same in both formulas.
The problem of finding Fock representations of the modified MF algebra has thus been reduced to representing (11). This may be done e.g. by introducing oscillators $`\varphi ^a(t)`$, $`\psi ^{a\mu }(t)`$, and $`\zeta ^{a\mu \nu }(t)`$, together with their canonical conjugate momenta $`\overline{\varphi }^a(t)`$, $`\overline{\psi }_\mu ^a(t)`$, and $`\overline{\zeta }_{\mu \nu }^a(t)`$. Moreover, $`\zeta ^{a\mu \nu }(t)`$ and $`\overline{\zeta }_{\mu \nu }^a(t)`$ are assumed to be symmetric in $`\mu \nu `$. The canonical commutation relation read
$`[\overline{\varphi }^a(s),\varphi ^b(t)]`$ $`=`$ $`\delta ^{ab}\delta (st),`$
$`[\overline{\psi }_\mu ^a(s),\psi ^{b\mu }(t)]`$ $`=`$ $`\delta ^{ab}\delta _\mu ^\nu \delta (st),`$ (13)
$`[\overline{\zeta }_{\mu \nu }^a(s),\zeta ^{b\sigma \tau }(t)]`$ $`=`$ $`\delta ^{ab}\delta _\mu ^{(\sigma }\delta _\nu ^{\tau )}\delta (st),`$
and all other brackets vanish. Then the following operators
$`J^a(t)`$ $`=`$ $`f^{abc}(:\varphi ^c(t)\overline{\varphi }^b(t):+:\psi ^{c\mu }(t)\overline{\psi }_\mu ^b(t):+:\zeta ^{c\mu \nu }(t)\overline{\zeta }_{\mu \nu }^b(t):),`$
$`G^{a\mu }(t)`$ $`=`$ $`f^{abc}\psi ^{c\mu }(t)\overline{\varphi }^b(t)+d^{abc}\zeta ^{c\mu \nu }(t)\overline{\psi }_\nu ^b(t),`$ (14)
$`H^{a\mu \nu }(t)`$ $`=`$ $`f^{abc}\zeta ^{c\mu \nu }(t)\overline{\varphi }^b(t),`$
satisfy (11). The double dots in the first expression indicate standard one-dimensional normal ordering with respect to frequency.
## 4 Diffeomorphisms
The algebra (LABEL:emb2), and thus also the modified MF algebra (LABEL:MF3), admits an intertwining action of an extension of the diffeomorphism algebra $`diff(N)`$. The additional brackets read
$`[_\mu (m),_\nu (n)]`$ $`=`$ $`n_\mu _\nu (m+n)m_\nu _\mu (m+n)`$
$`+(c_1m_\nu n_\mu +c_2m_\mu n_\nu )m_\rho S_1^\rho (m+n),`$
$`[_\mu (m),𝒥^a(n)]`$ $`=`$ $`n_\mu 𝒥^a(m+n),`$
$`[_\mu (m),𝒢^{a\nu }(n)]`$ $`=`$ $`n_\mu 𝒢^{a\nu }(m+n)+\delta _\mu ^\nu m_\rho 𝒢^{a\rho }(m+n),`$ (15)
$`[_\mu (m),^{a\nu \rho }(n)]`$ $`=`$ $`n_\mu ^{a\nu \rho }(m+n)`$
$`+\delta _\mu ^\nu m_\sigma ^{a\sigma \rho }(m+n)+\delta _\mu ^\rho m_\sigma ^{a\nu \sigma }(m+n),`$
$`[_\mu (m),S_1^\nu (n)]`$ $`=`$ $`n_\mu S_1^\nu (m+n)+\delta _\mu ^\nu m_\rho S_1^\rho (m+n),`$
where $`_\mu (m)`$ are the $`diff(N)`$ generators and the cocycles multiplied by $`c_1`$ and $`c_2`$ were first found by Eswara-Rao and Moody and myself , respectively. For the classification of $`diff(N)`$ cocycles, see and . To construct a representation of (15), introduce $`N`$ oscillators $`p_\mu (t)`$ which are the canonical momenta of $`q^\mu (t)`$, i.e.
$$[p_\mu (s),q^\nu (t)]=\delta _\mu ^\nu \delta (st),[p_\mu (s),p_\nu (t)]=[q^\mu (s),q^\nu (t)]=0.$$
(16)
The $`diff(N)`$ generators have the realization
$$_\mu (m)=dt(i:\mathrm{e}^{imq(t)}p_\mu (t):+m_\nu \mathrm{e}^{imq(t)}T_\mu ^\nu (t)),$$
(17)
where $`T_\nu ^\mu (t)`$ are the generators of the Kac-Moody algebra $`\widehat{gl(N)}`$. The relevant relations read
$`[T_\nu ^\mu (s),T_\tau ^\sigma (t)]`$ $`=`$ $`(\delta _\nu ^\sigma T_\tau ^\mu (s)\delta _\tau ^\mu T_\nu ^\sigma (s))\delta (st)`$
$`{\displaystyle \frac{1}{2\pi i}}(k_1\delta _\tau ^\mu \delta _\nu ^\sigma +k_2\delta _\nu ^\mu \delta _\tau ^\sigma )\dot{\delta }(st),`$
$`[T_\nu ^\mu (s),q^\rho (t)]`$ $`=`$ $`[T_\nu ^\mu (s),p_\rho (t)]=0,`$ (18)
$`[T_\nu ^\mu (s),J^a(t)]`$ $`=`$ $`0,`$
$`[T_\nu ^\mu (s),G^{a\sigma }(t)]`$ $`=`$ $`\delta _\nu ^\sigma G^{a\mu }(s)\delta (st),`$
$`[T_\nu ^\mu (s),H^{a\sigma \tau }(t)]`$ $`=`$ $`(\delta _\nu ^\sigma H^{a\mu \tau }(s)+\delta _\nu ^\tau H^{a\sigma \mu }(s))\delta (st),`$
and the abelian charges in (15) take the values $`c_1=1+k_1`$, $`c_2=k_2`$. $`\widehat{gl(N)}`$ acts on the same Fock space as the operators in (14), by means of the following expression:
$`T_\nu ^\mu (t)`$ $`=`$ $`\delta _\nu ^\mu (:\varphi ^a(t)\overline{\varphi }^a(t):+:\psi ^{a\sigma }(t)\overline{\psi }_\sigma ^a(t):+:\zeta ^{a\sigma \tau }(t)\overline{\zeta }_{\sigma \tau }^a(t):)+`$
$`+:\psi ^{a\mu }(t)\overline{\psi }_\nu ^a(t):+:\zeta ^{a\mu \rho }(t)\overline{\zeta }_{\nu \rho }^a(t):.`$
As was noted in , we can actually represent a larger algebra on the same Fock space. Namely, there is a natural action of an additional $`diff(1)`$ algebra, which classically commutes with both $`diff(N)`$ and the MF algebra. Geometrically, this algebra describes reparametrizations of the observer’s trajectory.
## 5 Conclusion
Originally, the present study had two goals: to construct projective Fock representations of the classical MF algebra, and to find representations where the cocycle had precisely the MF form (LABEL:MF2).
The first goal was easily reached using the formalism developped in . It is clear that much more complicated representations can be written down along the same lines. Geometrically, the oscillators can be viewed as zero-jets, i.e. the value of fields like $`\varphi ^a(x)`$ on the observer’s trajectory $`x^\mu =q^\mu (t)`$. One can generalize to $`p`$-jets, with basis $`_{\nu _1}\mathrm{}_{\nu _r}\varphi ^a(q(t)))`$ for all $`rp`$. This is a genuine $`N`$-dimensional object which probes not only the value of the fields along the trajectory, but also finitely many transverse derivatives. On the other hand, the presence of the MF term $`^{a\mu \nu }(m)`$ is quite uninteresting, since it can be disentangled using (6). The interesting quantum (normal ordering) effect is the Kac-Moody extension for the $`𝒥𝒥`$ bracket in (LABEL:emb2).
However, my second goal failed. Indeed, the fact that the true MF algebra (LABEL:MF2) can be naturally embedded into the non-Lie algebra (LABEL:emb1) is a serious obstruction against the existence of Fock modules. I am convinced that the technique of combining normal ordering with jet space trajectories can only produce one-chain (Kac-Moody) cocycles. This is a serious problem because this technique has so far been the only viable method to produce concrete Fock modules in more than one dimension. Two conclusions are possible: either the MF algebra lacks Fock modules altogether, or it points to a new type of representation theory which is not yet understood.
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# Quantum Monte Carlo calculations of 𝐴=8 nuclei
## I INTRODUCTION
The $`A=8`$ nuclei have many interesting properties that we would like to understand on the basis of the bare interactions between individual nucleons. The three strong-stable nuclei, <sup>8</sup>He, <sup>8</sup>Li, and <sup>8</sup>B, decay by weak processes with half-lives of $`1`$ second down to <sup>8</sup>Be, which then immediately fissions to two <sup>4</sup>He nuclei; thus there are no long-lived $`A=8`$ nuclei. These facts have tremendous consequences for the nature of the universe we live in: they make the production of elements beyond $`A=7`$ in the early universe very difficult and help give stars like the sun a long stable lifetime. The ground and excited states of the $`A=8`$ nuclei also provide a very sensitive testing ground for models of nuclear forces. In particular, <sup>8</sup>He is the most neutron-rich strong-stable nucleus, and thus an ideal place to study the isospin dependence of the three-nucleon force. The $`T=1,2`$ isomultiplets are also a good place to look at charge-independence breaking, while some excited states in <sup>8</sup>Be display significant isospin mixing.
Previously we have reported variational Monte Carlo (VMC) and Green’s function Monte Carlo (GFMC) calculations for the $`A=6,7`$ nuclei . In this paper we extend these calculations to the ground and excited states of $`A=8`$ nuclei. (Some preliminary $`A=8`$ results were given in a number of conference proceedings .) We first construct a trial function, $`\mathrm{\Psi }_T`$, with the proper $`(J^\pi ;T)`$ quantum numbers and antisymmetry, and optimize the energy expectation value. We then use $`\mathrm{\Psi }_T`$ as the starting point for a GFMC calculation, which projects out the exact lowest-energy state by the Euclidean propagation $`\mathrm{\Psi }_0=lim_\tau \mathrm{}\mathrm{exp}[(HE_0)\tau ]\mathrm{\Psi }_T`$. We believe the resulting energy estimates are accurate to 1 to 2% of the binding energy in most cases.
In our previous work, the Argonne $`v_{18}`$ plus Urbana IX Hamiltonian was fairly successful in generating the excitation spectra of the $`A=6,7`$ nuclei, but did not give quite enough binding in the lithium isotopes, while <sup>6</sup>He was unstable against breakup. One of our main goals in studying the $`A=8`$ nuclei is to continue to test the Hamiltonian, and lay the ground work for studies of more sophisticated and accurate three-nucleon forces. The Hamiltonian is reviewed in Sec. II. The variational wave functions and calculations are described in Sec. III. Our GFMC calculations make use of a new, more efficient algorithm called constrained-path propagation, which is described in Sec. IV. Results for the ground-state energies, excitation spectra, isomultiplet differences, and isospin-mixing matrix elements are given in Sec. V. In Sec. VI we present various density distributions, and in Sec. VII we discuss the intrinsic shapes of these nuclei. Our conclusions are given in Sec. VIII.
## II HAMILTONIAN
These calculations all use the same realistic Hamiltonian, which includes a nonrelativistic one-body kinetic energy, the Argonne $`v_{18}`$ two-nucleon potential and the Urbana IX three-nucleon potential :
$$H=\underset{i}{}K_i+\underset{i<j}{}v_{ij}+\underset{i<j<k}{}V_{ijk}.$$
(1)
The kinetic energy operator has charge-independent (CI) and charge-symmetry-breaking (CSB) components, the latter due to the difference in proton and neutron masses,
$$K_i=K_i^{CI}+K_i^{CSB}\frac{\mathrm{}^2}{4}(\frac{1}{m_p}+\frac{1}{m_n})_i^2\frac{\mathrm{}^2}{4}(\frac{1}{m_p}\frac{1}{m_n})\tau _{zi}_i^2.$$
(2)
The Argonne $`v_{18}`$ model is one of a class of new, highly accurate $`NN`$ potentials that fit both $`pp`$ and $`np`$ scattering data up to 350 MeV with a $`\chi ^2/`$datum near 1. The potential can be written as a sum of electromagnetic and one-pion-exchange terms and a shorter-range phenomenological part:
$$v_{ij}=v_{ij}^\gamma +v_{ij}^\pi +v_{ij}^R.$$
(3)
The electromagnetic terms include one- and two-photon-exchange Coulomb interactions, vacuum polarization, Darwin-Foldy, and magnetic moment terms, with appropriate proton and neutron form factors. The one-pion-exchange part of the potential includes the charge-dependent (CD) terms due to the difference in neutral and charged pion masses. The shorter-range part has about 40 parameters which are adjusted to fit the $`pp`$ and $`np`$ scattering data, the deuteron binding energy, and also the $`nn`$ scattering length. The one-pion-exchange and the remaining phenomenological part of the potential can be written as a sum of eighteen operators,
$$v_{ij}^\pi +v_{ij}^R=\underset{p=1,18}{}v_p(r_{ij})O_{ij}^p.$$
(4)
The first fourteen are charge-independent,
$`O_{ij}^{p=1,14}=[1,\sigma _i\sigma _j,S_{ij},𝐋𝐒,𝐋^2,𝐋^2\sigma _i\sigma _j,(𝐋𝐒)^2][1,\tau _i\tau _j],`$ (5)
and the last four,
$$O_{ij}^{p=15,18}=[1,\sigma _i\sigma _j,S_{ij}]T_{ij},\tau _{zi}+\tau _{zj},$$
(6)
break charge independence. We will refer to the potential from the $`p=1517`$ terms as $`v^{CD}`$ and from the $`p=18`$ term as $`v^{CSB}`$. We note that in the context of isospin symmetry the CI, CSB and CD terms are respectively isoscalar, isovector and isotensor.
The two-nucleon potential is supplemented by a three-nucleon interaction from the Urbana series of models , including both long-range two-pion exchange and a short-range phenomenological component:
$$V_{ijk}=V_{ijk}^{2\pi }+V_{ijk}^R.$$
(7)
The two-pion-exchange term can be expressed simply as a sum of anticommutator and commutator terms,
$$V_{ijk}^{2\pi }=\underset{cyclic}{}V_{ij;k}^{2\pi ,A}+V_{ij;k}^{2\pi ,C}.$$
(8)
Here
$$V_{ij;k}^{2\pi ,A}=A_{2\pi }\{X_{ik}^\pi ,X_{jk}^\pi \}\{\tau _i\tau _k,\tau _j\tau _k\},$$
(9)
and
$$V_{ij;k}^{2\pi ,C}=\frac{1}{4}A_{2\pi }[X_{ik}^\pi ,X_{jk}^\pi ][\tau _i\tau _k,\tau _j\tau _k],$$
(10)
with $`X_{ij}^\pi =Y(r_{ij})\sigma _i\sigma _j+T(r_{ij})S_{ij}`$ as the basic one-pion exchange operator. The $`V_{ijk}^R`$ has no spin-isospin dependence, and in the Urbana model IX its strength and $`A_{2\pi }`$ are adjusted to reproduce the binding energy of <sup>3</sup>H and give a reasonable saturation density in nuclear matter when used with Argonne $`v_{18}`$.
The CD and CSB terms in $`H`$ are fairly weak, so we can treat them conveniently as a first-order perturbation and use a wave function of good isospin, which is significantly more compact. Also, direct GFMC calculations with the spin-dependent terms that involve the square of the momentum operator can have large statistical fluctuations, as discussed in Ref. . Thus we construct the GFMC propagator with a simpler isoscalar Hamiltonian,
$$H^{}=\underset{i}{}K_i^{CI}+\underset{i<j}{}v_{ij}^{}+\underset{i<j<k}{}V_{ijk}^{},$$
(11)
where $`v_{ij}^{}`$ is defined as
$$v_{ij}^{}=\underset{p=1,8}{}v_p^{}(r_{ij})O_{ij}^p+v_C^{}(r_{ij}).$$
(12)
The interaction $`v_{ij}^{}`$ has only eight operator terms, with operators $`[1,\sigma _𝐢\sigma _𝐣,S_{ij},𝐋𝐒][1,\tau _𝐢\tau _𝐣]`$, chosen such that it equals the isoscalar part of the full interaction in all $`S`$ and $`P`$ waves as well as in the $`{}_{}{}^{3}D_{1}^{}`$ wave and its coupling to the $`{}_{}{}^{3}S_{1}^{}`$. The isoscalar part of the $`pp`$ Coulomb interaction, $`v_C^{}`$, is also included in $`H^{}`$. Detailed expressions are given in Ref. . The $`v_{ij}^{}`$ is a little more attractive than $`v_{ij}`$, so we compensate by using a $`V_{ijk}^{}`$ that is adjusted to keep $`H^{}H`$; this should help prevent the GFMC propagation from producing excessively large densities due to overbinding. The small contribution of $`(HH^{})`$ is calculated perturbatively.
## III VARIATIONAL MONTE CARLO
The variational method can be used to obtain approximate solutions to the many-body Schrödinger equation, $`H\mathrm{\Psi }=E\mathrm{\Psi }`$, for a wide range of nuclear systems, including few-body nuclei, light closed-shell nuclei, nuclear matter, and neutron stars . A suitably parameterized wave function, $`\mathrm{\Psi }_V`$, is used to calculate an upper bound to the exact ground-state energy,
$$E_V=\frac{\mathrm{\Psi }_V|H|\mathrm{\Psi }_V}{\mathrm{\Psi }_V|\mathrm{\Psi }_V}E_0.$$
(13)
The parameters in $`\mathrm{\Psi }_V`$ are varied to minimize $`E_V`$, and the lowest value is taken as the approximate ground-state energy.
Upper bounds to excited states are also obtainable, either from standard VMC calculations if they have different quantum numbers from the ground state, or from small-basis diagonalizations if they have the same quantum numbers. The corresponding $`\mathrm{\Psi }_V`$ can then be used to calculate other properties, such as electromagnetic form factors and spectroscopic factors , or it can be used as the starting point for a Green’s function Monte Carlo calculation. In this section we first describe our ansatz for $`\mathrm{\Psi }_V`$ for the $`A=8`$ nuclei and then briefly review how the expectation value, Eq. (13), is evaluated and the parameters of $`\mathrm{\Psi }_V`$ are fixed.
### A Wave Function
Our best variational wave function for the nuclei studied here has the form
$$|\mathrm{\Psi }_V=\left[1+\underset{i<j<k}{}(U_{ijk}+U_{ijk}^{TNI})+\underset{i<j}{}U_{ij}^{LS}\right]|\mathrm{\Psi }_P,$$
(14)
where the pair wave function, $`\mathrm{\Psi }_P`$, is given by
$$|\mathrm{\Psi }_P=𝒮\underset{i<j}{}(1+U_{ij})|\mathrm{\Psi }_J.$$
(15)
The $`U_{ij}`$, $`U_{ij}^{LS}`$, $`U_{ijk}`$, and $`U_{ijk}^{TNI}`$ are noncommuting two- and three-nucleon correlation operators, and the $`𝒮`$ is a symmetrization operator. The $`U_{ij}`$ includes spin, isospin, and tensor operators, while the $`U_{ij}^{LS}`$ has spin-orbit operators, reflecting the operator structure of the two-nucleon interaction, Eq. (4). The $`U_{ijk}`$ is a nontrivial operator in spin-isospin space also induced by $`v_{ij}`$, while the $`U_{ijk}^{TNI}`$ reflects the structure of the three-nucleon interaction. All these correlations are discussed fully in Refs. . The two-body correlations are generated by the solution of coupled differential equations with embedded variational parameters. We have found that the parameters optimized for the $`\alpha `$-particle are near optimal for use in the light p-shell nuclei. Likewise, the best parameters for the three-body correlations are remarkably constant for different s- and p-shell nuclei, so they have not been changed from the previous work.
The form of the totally antisymmetric Jastrow wave function, $`\mathrm{\Psi }_J`$, depends on the nuclear state under investigation. For s-shell nuclei we use the simple form
$$|\mathrm{\Psi }_J=\underset{i<j<k}{}f_{ijk}^c\underset{i<j}{}f_c(r_{ij})|\mathrm{\Phi }_A(JMTT_3).$$
(16)
Here $`f_c(r_{ij})`$ and $`f_{ijk}^c`$ are central two- and three-body correlation functions and for the $`\alpha `$-particle,
$$|\mathrm{\Phi }_4(0000)=𝒜|ppnn.$$
(17)
The Jastrow wave function for the light p-shell nuclei is significantly more complicated due to the requirements of antisymmetry. Expressions for $`A=6,7`$ nuclei are given in Ref. ; the present $`\mathrm{\Psi }_J`$ is a straightforward extension. We use $`LS`$ coupling to obtain the desired $`JM`$ value of a given state, as suggested by standard shell-model studies . We also need to specify the spatial symmetry $`[n]`$ of the angular momentum coupling of four p-shell nucleons . Different possible $`LS[n]`$ combinations lead to multiple components in the Jastrow wave function. We allow for the possibility that the central correlations $`f_c(r_{ij})`$ could depend upon the shells ($`s`$ or $`p`$) occupied by the particles and on the $`LS[n]`$ coupling. The Jastrow wave function is taken as
$`|\mathrm{\Psi }_J`$ $`=`$ $`𝒜\{{\displaystyle \underset{i<j<k4}{}}f_{ijk}^c{\displaystyle \underset{i<j4}{}}f_{ss}(r_{ij}){\displaystyle \underset{k4}{}}f_{sp}(r_{k5})f_{sp}(r_{k6})f_{sp}(r_{k7})f_{sp}(r_{k8})`$ (19)
$`{\displaystyle \underset{LS[n]}{}}(\beta _{LS[n]}{\displaystyle \underset{5l<m8}{}}f_{pp}^{[n]}(r_{lm})|\mathrm{\Phi }_8(LS[n]JMTT_3)_{1234:5678})\}.`$
The operator $`𝒜`$ indicates an antisymmetric sum over all possible partitions of the eight particles into 4 s-shell and 4 p-shell ones. For the two-body correlations we use $`f_{ss}(r)=f_c(r)`$ from the <sup>4</sup>He wave function, while
$`f_{sp}(r)`$ $`=`$ $`f_c(r)+c_{sp}(1\mathrm{exp}[(r/d_{sp})^2]),`$ (20)
and
$`f_{pp}^{[n]}(r)`$ $`=`$ $`f_c(r)+c_{pp}^{[n]}(1\mathrm{exp}[(r/d_{pp})^2]),`$ (22)
where we have supplemented the $`f_c(r)`$ with a long-range tail. The $`c_{sp}`$, $`d_{sp}`$, etc., are variational parameters, whose values are given in Table I. For the three-body correlations, our best present trial function has the $`f_{ijk}^c`$ acting only within the s-shell.
The $`LS[n]`$ components of the single-particle wave function are given by:
$`|\mathrm{\Phi }_8(LS[n]JMTT_3)_{1234:5678}`$ $`=`$ $`|\mathrm{\Phi }_4(0000)_{1234}\varphi _p^{LS}(R_{\alpha 5})\varphi _p^{LS}(R_{\alpha 6})\varphi _p^{LS}(R_{\alpha 7})\varphi _p^{LS}(R_{\alpha 8})`$ (26)
$`\{[Y_{1m_l}(\mathrm{\Omega }_{\alpha 5})Y_{1m_l^{}}(\mathrm{\Omega }_{\alpha 6})Y_{1m_l^{\prime \prime }}(\mathrm{\Omega }_{\alpha 7})Y_{1m_l^{\prime \prime \prime }}(\mathrm{\Omega }_{\alpha 8})]_{LM_L}`$
$`\times [\chi _5({\displaystyle \frac{1}{2}}m_s)\chi _6({\displaystyle \frac{1}{2}}m_s^{})\chi _7({\displaystyle \frac{1}{2}}m_s^{\prime \prime })\chi _8({\displaystyle \frac{1}{2}}m_s^{\prime \prime \prime })]_{SM_S}\}_{JM}`$
$`\times [\nu _5({\displaystyle \frac{1}{2}}t_3)\nu _6({\displaystyle \frac{1}{2}}t_3^{})\nu _7({\displaystyle \frac{1}{2}}t_3^{\prime \prime })\nu _8({\displaystyle \frac{1}{2}}t_3^{\prime \prime \prime })]_{TT_3}.`$
The $`\varphi _p^{LS}(R_{\alpha k})`$ are $`p`$-wave solutions of a particle of reduced mass $`\frac{4}{5}m_N`$ in an effective $`\alpha `$-$`N`$ potential:
$$V_{\alpha N}(r)=V_{\alpha N}^{WS}(r)+V_{\alpha N}^C(r).$$
(27)
They are functions of the distance between the center of mass of the $`\alpha `$ core (which contains particles 1-4 in this partition) and nucleon $`k`$, and again may be different for different $`LS`$ components. For each state considered in the present work, we have used bound-state asymptotic conditions for the $`\varphi _p^{LS}`$, even if the state is particle unstable. The Woods-Saxon potential
$$V_{\alpha N}^{WS}(r)=V_p^{LS}[1+exp(\frac{rR_p}{a_p})]^1,$$
(28)
has variational parameters $`V_p^{LS}`$, $`R_p`$, and $`a_p`$ whose values are given in Table I. The Coulomb potential is obtained by folding over nuclear form factors:
$`V_{\alpha N}^C(r)=2Q{\displaystyle \frac{e^2}{r}}\{1`$ $``$ $`{\displaystyle \frac{1}{2}}exp(x_\alpha )[2+x_\alpha +{\displaystyle \frac{4}{1y^2}}][1y^2]^2`$ (29)
$``$ $`{\displaystyle \frac{1}{2}}exp(x_p)[2+x_p+{\displaystyle \frac{4}{1y^2}}][1y^2]^2\}.`$ (30)
Here $`x_\alpha =\sqrt{12}r/r_\alpha `$, $`x_p=\sqrt{12}r/r_p`$, and $`y=r_\alpha /r_p`$, with the charge radii $`r_\alpha =1.65`$ fm and $`r_p=0.81`$ fm. This additional potential term can be used with strength $`Q=0`$, $`\frac{1}{4}`$, or $`\frac{1}{2}`$ for <sup>8</sup>He, <sup>8</sup>Li, or <sup>8</sup>Be, respectively, corresponding to the average Coulomb interaction between the $`\alpha `$ core and a p-shell nucleon. The wave function is translationally invariant, hence there is no spurious center of mass motion.
The experimental spectra for $`A=8`$ nuclei are shown in Fig. 1. The ground state of <sup>8</sup>He is strong stable, but decays by $`\beta ^{}`$ emission with a half life of 119 ms. One excited state is identified at 3.59 $`\pm `$ 0.05 MeV , above the threshold for decay to <sup>6</sup>He+2$`n`$. In the shell model, the $`(J^\pi ;T)=(0^+;2)`$ ground state is predominantly a $`{}_{}{}^{2S+1}L[n]=^1`$S state, where we use spectroscopic notation to denote the total $`L`$ and $`S`$ and the Young pattern $`[n]`$ to indicate the spatial symmetry. The $`(2^+;2)`$ first excited state is predominantly a <sup>1</sup>D state. We also allow for possible <sup>3</sup>P admixtures, which are the only other available configurations in the p-shell, in our $`\mathrm{\Psi }_J`$. After other parameters in the trial function have been optimized, we make a series of calculations in which the $`\beta _{LS[n]}`$ may be different in the left- and right-hand-side wave functions to obtain the diagonal and off-diagonal matrix elements of the Hamiltonian and the corresponding normalizations and overlaps. We diagonalize the resulting matrices to find the $`\beta _{LS[n]}`$ eigenvectors. The shell-model wave functions are orthonormal, but the correlated $`\mathrm{\Psi }_V`$ are not. Hence the diagonalizations use generalized eigenvalue routines including overlap matrices. We also calculate the position of the three predominantly <sup>3</sup>P states, with $`(J^\pi ;T)=(2^+;2)`$, $`(1^+;2)`$, and $`(0^+;2)`$; none of these have been identified experimentally. The normalized $`\beta _{LS[n]}`$ for these different states are given in Table II.
The ground state of <sup>8</sup>Li is a strong stable $`(2^+;1)`$ state that decays by $`\beta ^{}`$ emission with a half life of 838 ms, and is predominantly <sup>3</sup>P in character. The $`(1^+;1)`$ first excited state at 0.98 MeV excitation is also strong stable with a 12 fs $`\gamma `$-decay, and is primarily a mix of <sup>1,3</sup>P configurations. The $`(3^+;1)`$ second excited state at 2.26 MeV is just above the threshold for breakup into <sup>7</sup>Li+$`n`$ and is fairly narrow with a width of 33 keV, decaying by both $`\gamma `$ and $`n`$ emission. A number of higher states have been identified, most having fairly large widths. Two special cases are the $`(4^+;1)`$ stretched state at 6.53 MeV excitation with a width of 35 keV, which can only come from the <sup>3</sup>F p-shell configuration, and the $`(0^+;2)`$ isobaric analog of <sup>8</sup>He that occurs at 10.82 MeV with a width less than 12 keV. The possible p-shell components in the $`T=1`$ states include the <sup>1,3</sup>P, <sup>1,3</sup>D, <sup>1,3</sup>F, <sup>3</sup>S, <sup>3</sup>D, and <sup>1,3,5</sup>P configurations. We include all but the lowest-symmetry components in our $`\mathrm{\Psi }_J`$ and calculate all possible first, second and some third states of given $`(J^\pi ;T)`$ through the diagonalization procedure discussed above. Table III gives a summary of the $`\beta _{LS[n]}`$ amplitudes.
The $`(0^+;0)`$ ground state of <sup>8</sup>Be is 92 keV above the threshold for breakup into two $`\alpha `$-particles, with a width of only 7 eV; it is an almost pure <sup>1</sup>S configuration. The first and second excited states are very broad: a $`(2^+;0)`$ <sup>1</sup>D state at 3.04 MeV and a $`(4^+;0)`$ <sup>1</sup>G state at 11.4 MeV — the spacing of an almost rigid rotor. Indeed, as discussed below, <sup>8</sup>Be appears to have the intrinsic deformation of a $`2\alpha `$ molecule. Higher in the spectrum is the famous pair of isospin-mixed $`(2^+;0+1)`$ states at 16.63 and 16.92 MeV which have widths of order 100 keV; the $`T=1`$ component is the isobaric analog of the <sup>8</sup>Li ground state. There are similar $`(1^+;0+1)`$ and $`(3^+;0+1)`$ pairs near 18 and 19 MeV which have widths less than 300 keV. Many additional states have been identified above 18 MeV, most of them with large widths, up to the $`(0^+;2)`$ isobaric analog of <sup>8</sup>He at 27.49 MeV. There are also some negative-parity states in this region, which we have not attempted to calculate. In constructing $`\mathrm{\Psi }_J`$, we use all p-shell states of symmetry and , but neglect those of symmetry and , as shown in Table IV.
The known experimental spectrum for <sup>8</sup>B is similar to <sup>8</sup>Li, except that with the extra Coulomb repulsion, the $`(2^+;1)`$ ground state is just barely strong stable, decaying by $`\beta +`$ emission with a half life of 770 ms. The $`(1^+;1)`$ first excited state is seen as a narrow resonance in <sup>7</sup>Be+$`p`$ scattering, and the $`(3^+;1)`$ second excited state is much broader than its <sup>8</sup>Li analog. The only other observed state is the $`(0^+;2)`$ isobaric analog to <sup>8</sup>He at 10.62 MeV. The ground state for <sup>8</sup>C is unstable against several possible breakup channels, having a width of $``$230 keV, but its mass excess is known within 20 keV. We have calculated the energies of these states in <sup>8</sup>B and <sup>8</sup>C so we can study the energy differences in the $`T=1`$ and 2 isomultiplets.
The full $`A=8`$ wave function is constructed by acting on the $`\mathrm{\Psi }_J`$, Eq. (19), with the $`U_{ij}`$, $`U_{ij}^{LS}`$, $`U_{ijk}`$, and $`U_{ijk}^{TNI}`$ correlations. Because of the tensor and other correlations in $`U_{ij}`$, many additional symmetry components, beyond the explicit p-shell states discussed above, are built up in the wave function. In principle, the $`U_{ij}`$ could be generalized to be different according to whether particles $`i`$ and $`j`$ are in the s- or p-shell, but this would require a larger sum over the different partitions and would increase the computational cost by an order of magnitude.
For input to the GFMC algorithm, it is more efficient to use the somewhat simplified trial function
$`|\mathrm{\Psi }_T=𝒮{\displaystyle \underset{i<j}{}}\left(1+U_{ij}+{\displaystyle \underset{ki,j}{}}\stackrel{~}{U}_{ij;k}^{TNI}\right)|\mathrm{\Psi }_J,`$ (31)
where $`\stackrel{~}{U}_{ij;k}^{TNI}`$ is a truncated three-nucleon interaction correlation based on the short-range $`V_{ijk}^R`$ term and on the anticommutator part of the two-pion exchange, $`V_{ij;k}^{2\pi ,A}`$, which can be reduced to operators that depend only on the spins and isospins of nucleons $`i`$ and $`j`$. Thus the sum over $`k`$ can be made, leaving a two-body spin-isospin operator that can be combined with $`U_{ij}`$; the result is calculable with only a little more effort than just $`U_{ij}`$ alone. This trial function gets the bulk of the energy, as shown below, but for about half the computational effort of the full $`\mathrm{\Psi }_V`$.
### B Energy Evaluation
The energy expectation value of Eq. (13) is evaluated using Monte Carlo integration. A detailed technical description of the methods used here can be found in Refs. . Monte Carlo sampling is done both in configuration space and in the discrete order of operators in the symmetrized product of the pair wave function by following a Metropolis random walk. The expectation value for an operator $`O`$ with the full wave function $`\mathrm{\Psi }_V`$is given by
$$O=\frac{_{p,q}𝑑𝐑\left[\mathrm{\Psi }_{V,p}^{}(𝐑)O\mathrm{\Psi }_{V,q}(𝐑)/W_{pq}(𝐑)\right]W_{pq}(𝐑)}{_{p,q}𝑑𝐑\left[\mathrm{\Psi }_{V,p}^{}(𝐑)\mathrm{\Psi }_{V,q}(𝐑)/W_{pq}(𝐑)\right]W_{pq}(𝐑)},$$
(32)
where we have introduced a probability distribution, $`W_{pq}(𝐑)`$, based on the approximate wave function $`\mathrm{\Psi }_P`$ of Eq. (15),
$$W_{pq}(𝐑)=|\mathrm{Re}(\mathrm{\Psi }_{P,p}^{}(𝐑)\mathrm{\Psi }_{P,q}(𝐑))|.$$
(33)
The subscripts $`p`$ and $`q`$ specify the order of operators on the left and right hand side of the pair wave functions, while the integration runs over the particle coordinates $`𝐑=(𝐫_1,𝐫_2,\mathrm{},𝐫_A)`$. This probability distribution is much less expensive to compute than one using the full wave function of Eq. (14) with its spin-orbit and operator-dependent three-body correlations, but the denominator of Eq. (32) is typically within 1-2% of unity. Expectation values have a statistical error which can be estimated by the standard deviation $`\sigma `$ in either gaussian approximation or by using block averaging schemes.
Our wave functions are vectors of $`2^A\times I(A,T)`$ complex numbers,
$$\mathrm{\Psi }(𝐑)=\underset{\alpha }{}\psi _\alpha (𝐑)|\alpha ,$$
(34)
where the $`\psi _\alpha (𝐑)`$ are the coefficients of each possible spin-isospin state $`|\alpha `$ with specific third components of the spins of each nucleon and the desired total isospin. The CD and CSB force components are sufficiently small that we do not worry about isospin mixing in the wave function, and the expression for the number of isospin states, $`I(A,T)`$, is given in Ref. . This gives arrays with 3584, 5120, and 7168 elements for <sup>8</sup>Be, <sup>8</sup>He, and <sup>8</sup>Li, respectively. The spin, isospin, and tensor operators $`O_{ij}^{p=2,6}`$ contained in the two-body correlation operator $`U_{ij}`$, and in the Hamiltonian are sparse matrices in this basis.
Expectation values of the kinetic energy and spin-orbit potential require the computation of first derivatives and diagonal second derivatives of the wave function. These are obtained by evaluating the wave function at $`6A`$ slightly shifted positions of the coordinates $`𝐑`$ and taking finite differences, as discussed in Ref. . Potential terms quadratic in L require mixed second derivatives, which can be obtained by additional wave function evaluations and finite differences. A rotation trick can be used to reduce the number of additional locations at which the wave function must be evaluated .
In addition to calculating energies, we evaluate $`J^2`$ and $`J_z`$ expectation values to verify that our wave functions truly have the specified quantum numbers. Another check is made on the antisymmetry of the Jastrow wave function by evaluating, at an initial randomized position,
$$\frac{\mathrm{\Psi }_J^{}[1+P_{ij}^xP_{ij}^\sigma P_{ij}^\tau ]\mathrm{\Psi }_J}{\mathrm{\Psi }_J^{}\mathrm{\Psi }_J},$$
(35)
where $`P_{ij}^{x,\sigma ,\tau }`$ are the space, spin, and isospin exchange operators. This value should be exactly zero for an antisymmetric wave function, and it is in fact less than $`10^9`$ for each pair of particles in each nuclear state that we study.
A major problem arises in minimizing the variational energy for p-shell nuclei using the above wave functions: there is no variational minimum that gives reasonable rms radii. For example, the variational energy for <sup>6</sup>Li is slightly more bound than for <sup>4</sup>He, but is not more bound than for separated <sup>4</sup>He and <sup>2</sup>H nuclei, so the wave function is not stable against breakup into $`\alpha +d`$ subclusters. Consequently, the energy can be lowered toward the sum of <sup>4</sup>He and <sup>2</sup>H energies by making the wave function more and more diffuse. Such a diffuse wave function would not be useful for computing other nuclear properties, or as a starting point for the GFMC calculation.
On the basis of our work in $`A`$ = 6–7 nuclei, we believe that part of the problem is a fault of the Hamiltonian, but most of it is due to the presence of small admixtures of highly excited states in the trial function . In that work, we constrained our search for optimal variational parameters by requiring the resulting point proton rms radius, $`r_p`$, to be close to the experimental values for <sup>6</sup>Li and <sup>7</sup>Li ground states. Then we allowed only small variations in the construction of the <sup>6</sup>He and <sup>7</sup>Be ground states and all the excited or resonant states, for which there are no experimental measurements of the charge radii.
In the present work, we have the additional complication that all the $`A=8`$ nuclei are sufficiently short-lived that precise experimental determinations of their radii exist do not exist. Consequently, the variational parameters in Table I are chosen very close to those of our previous work, with only a systematic reduction in the depth of the Woods-Saxon potential well as $`L`$ increases, and an increase in the tail of the $`f_{pp}^{[n]}(r)`$ correlation as the spatial symmetry declines, as shown in Table I. These gradual changes help to insure that the radii of excited states increase as the excitation energy increases.
The last step is always the diagonalization to determine the $`\beta _{LS[n]}`$ mixing coefficients of Tables II-IV. Shell model lore tells us that the lowest state of any given $`(J^\pi ;T)`$ will be the state with maximal spatial symmetry and smallest $`L`$ that can be formed from the allowed couplings, e.g., the <sup>1</sup>S ground state in <sup>8</sup>Be or the <sup>3</sup>P ground state in <sup>8</sup>Li. For the purposes of obtaining a variational upper bound and a GFMC starting point, we could settle for a $`\mathrm{\Psi }_V`$ constructed using only that $`LS[n]`$ component. However, by using more components, we can gain a significant amount of energy in some cases and this gain persists in our GFMC propagations. For the <sup>8</sup>He ground state, when the dominant <sup>1</sup>S piece is supplemented by the <sup>3</sup>P term, the energy gain is 1.2 MeV. In the case of <sup>8</sup>Li ground state, where the dominant term is <sup>3</sup>P, addition of the three other contributing symmetry states lowers the energy 0.8 MeV; addition of the one state (the next highest spatial symmetry) gives only an additional 0.1 MeV. However, for the <sup>8</sup>Be ground state, with the dominant <sup>1</sup>S term, addition of the one symmetry piece gives only 0.1 MeV. For this reason, we feel reasonably confident in truncating the <sup>8</sup>Li and <sup>8</sup>Be p-shell bases after the top two symmetry states.
## IV GREEN’S FUNCTION MONTE CARLO
A detailed description of the nuclear GFMC method, and many tests of its accuracy, are given in Ref. . In this section we present a brief review of the method and then describe two improvements that have been made since that publication. In most of this section, we will not make the distinction between $`H^{}`$ and $`H`$; the reader should remember that in fact we use the simpler $`H^{}`$ in our GFMC propagator and evaluate $`HH^{}`$ perturbatively.
### A Review
The GFMC method starts with the trial wave function, $`\mathrm{\Psi }_T`$ of Eq. (31), and projects out of it the exact lowest energy state with the same quantum numbers, $`\mathrm{\Psi }_0`$:
$$\mathrm{\Psi }_0=\underset{\tau \mathrm{}}{lim}\mathrm{\Psi }(\tau ),$$
(36)
$$\mathrm{\Psi }(\tau )=e^{(HE_0)\tau }\mathrm{\Psi }_T=\left[e^{(HE_0)\mathrm{}\tau }\right]^n\mathrm{\Psi }_T.$$
(37)
Here we have sliced the imaginary propagation time, $`\tau `$, into a number of small time steps, $`\mathrm{}\tau =\tau /n`$. The small-time-step Green’s function, $`G_{\alpha \beta }(𝐑,𝐑^{})`$, is a matrix function of $`𝐑`$ and $`𝐑^{}`$ in spin-isospin space, with matrix elements defined as
$`G_{\alpha \beta }(𝐑,𝐑^{})=𝐑,\alpha |e^{(HE_0)\mathrm{}\tau }|𝐑^{},\beta .`$ (38)
Then $`\mathrm{\Psi }(𝐑_n,\tau )`$ is given by
$`\mathrm{\Psi }(𝐑_n,\tau )={\displaystyle G(𝐑_n,𝐑_{n1})\mathrm{}G(𝐑_1,𝐑_0)\mathrm{\Psi }_T(𝐑_0)𝑑𝐑_{n1}\mathrm{}𝑑𝐑_1𝑑𝐑_0}.`$ (39)
The small-time-step propagator used in Ref. is
$`G_{\alpha \beta }(𝐑,𝐑^{})`$ $`=`$ $`e^{E_0\mathrm{}\tau }G_0(𝐑,𝐑^{})\mathrm{exp}\left[{\displaystyle \frac{\mathrm{}\tau }{2}}{\displaystyle (V_{ijk}^R(𝐑)+V_{ijk}^R(𝐑^{}))}\right]`$ (41)
$`\times \alpha |I_3(𝐑)|\gamma \gamma |\left[𝒮{\displaystyle \underset{i<j}{}}{\displaystyle \frac{g_{ij}(𝐫_{ij},𝐫_{ij}^{})}{g_{0,ij}(𝐫_{ij},𝐫_{ij}^{})}}\right]|\delta \delta |I_3(𝐑^{})|\beta ,`$
where $`g_{ij}(𝐫_{ij},𝐫_{ij}^{})`$ is the exact two-body propagator, $`g_{0,ij}(𝐫_{ij},𝐫_{ij}^{})`$ is the free two-body propagator, $`G_0(𝐑,𝐑^{})`$ is the free many-body propagator, $`\alpha ,\beta ,\gamma ,\delta `$ are spin-isospin state indices, and summation over $`\gamma ,\delta `$ is implied. There is also an implicit sampling of the order of pairs in the symmetrized product of Eq. (41). The construction of the exact two-body propagator is described in Ref. . The influence of the three-nucleon potential on the many-body propagator is broken into two pieces: the scalar $`V_{ijk}^R`$ which is easily exponentiated, and the $`V_{ijk}^{2\pi }`$ which is a more complicated operator in spin-isospin space. The simplest treatment of this term in the TNI is to expand to first order in $`\mathrm{}\tau `$,
$$I_3(𝐑)=1\frac{\mathrm{}\tau }{2}V_{ijk}^{2\pi }(𝐑);$$
(42)
in fact we use a more efficient procedure described below in subsection E.
The integrals in Eq. (39) are evaluated stochastically by averaging over a set of $`n`$-step paths, $`𝐏_n=𝐑_0,𝐑_1,\mathrm{},𝐑_n`$. The paths are chosen by first sampling a set of positions, $`𝐑_0`$, using a probability function based on $`\mathrm{\Psi }_T(𝐑_0)`$, and then sequentially sampling the free Green’s functions, $`G_0(𝐑_{i+1},𝐑_i)`$, to generate $`𝐑_{i+1}`$ from $`𝐑_i`$. We thus obtain a spin-isospin vector $`\mathrm{\Psi }(𝐏_n)`$ which is one sample of the integrand of Eq. (39). A GFMC “configuration” consists of the position, $`𝐑_n`$, and the vector $`\mathrm{\Psi }(𝐏_n)`$. Branching and importance sampling, described in detail in Ref. , are used to obtain samples with probability proportional to the scalar importance function $`I`$:
$$I[\mathrm{\Psi }(𝐏_n),\mathrm{\Psi }_{T,p}(𝐑_n)]=|\mathrm{Re}[\underset{\alpha }{}\psi _\alpha (𝐏_n)^{}\psi _{T,p,\alpha }(𝐑_n)]|+ϵ\underset{\alpha }{}|[\psi _\alpha (𝐏_n)^{}\psi _{T,p,\alpha }(𝐑_n)]|,$$
(43)
where the sums run over the spin-isospin states $`\alpha `$, and $`\mathrm{\Psi }_{T,p}`$ is the trial wave function evaluated with a specific choice of pair operator orders $`p`$. Here $`ϵ`$ is a small constant ($`0.01`$) that ensures a positive-definite importance function so that diffusion can take place across nodal surfaces.
The GFMC method allows one to compute “mixed” expectation values:
$`O_{Mixed}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Psi }_T|O|\mathrm{\Psi }(\tau )}{\mathrm{\Psi }_T|\mathrm{\Psi }(\tau )}},`$ (44)
$`=`$ $`{\displaystyle \frac{𝑑𝐑_n\mathrm{\Psi }_T^{}(𝐑_n)O\mathrm{\Psi }(𝐑_n,\tau )}{𝑑𝐑_n\mathrm{\Psi }_T^{}(𝐑_n)\mathrm{\Psi }(𝐑_n,\tau )}}.`$ (45)
Because $`H`$ commutes with the propagator, $`H(\tau )_{Mixed}`$ is an upper bound to $`E_0`$ and approaches $`E_0`$ from above. However expectation values of operators that do not commute with $`H`$ are extrapolated using
$`O(\tau )`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Psi }(\tau )|O|\mathrm{\Psi }(\tau )}{\mathrm{\Psi }(\tau )|\mathrm{\Psi }(\tau )}},`$ (46)
$``$ $`O(\tau )_{Mixed}+[O(\tau )_{Mixed}O_T],`$ (47)
where
$$O_T=\frac{\mathrm{\Psi }_T|O|\mathrm{\Psi }_T}{\mathrm{\Psi }_T|\mathrm{\Psi }_T}.$$
(48)
In the following we address refinements which have been made to the GFMC algorithm since Ref. . These refinements are important in order to make calculations of larger nuclei feasible.
### B Constrained Path Algorithm
Diffusion or Green’s function Monte Carlo simulations of many-fermion systems generally suffer from the so-called “fermion-sign problem”. In essence this results from stochastically evaluating matrix elements of the form encountered in Eq. (45). The Monte Carlo techniques used to calculate the path integrals leading to $`\mathrm{\Psi }(𝐑_n,\tau )`$ involve only local properties, while antisymmetry is a global property. This leads to integrands in Eq. (45) that have oscillating signs at large $`\tau `$, which cause the statistical error to grow exponentially with imaginary time. The problem is not insurmountable for light systems, because the states are fairly well separated and one can propagate for a substantial imaginary time without having unacceptable statistical errors. However, the sign problem also grows exponentially with particle number, as the interchange of any pair of nucleons causes a change in sign for the matrix element.
The sign problem significantly limits the maximum $`\tau `$ that we can use in the simulations; hence we invoke an approximate technique to deal with it. The approximation involves keeping only a subset of the paths in evaluating the integrals, and using the knowledge gained in the VMC calculations to choose the subset of paths. It is closely related to methods used previously in condensed matter and elsewhere; they generally go by the name of constrained-path techniques . Some details of the algorithm, however, are special to the nuclear physics case.
The basic idea of the constrained-path method is to discard those configurations that, in future generations, will contribute only noise to expectation values. If we knew the exact ground state $`|\mathrm{\Psi }_0`$, we could discard any configuration for which:
$$\mathrm{\Psi }(𝐏_n)^{}\mathrm{\Psi }_0(𝐑_n)=0,$$
(49)
where a sum over spin-isospin states is implied. The sum of these discarded configurations can be written as a state $`|\mathrm{\Psi }_d`$, which obviously has zero overlap with the ground state
$$\mathrm{\Psi }_d|\mathrm{\Psi }_0=0.$$
(50)
The $`\mathrm{\Psi }_d`$ contains only excited states and should decay away as $`\tau \mathrm{}`$, thus discarding it is justified. However, in general the ground state $`\mathrm{\Psi }_0`$ is not known exactly and hence the constraint is imposed approximately using $`\mathrm{\Psi }_T`$ in the place of $`\mathrm{\Psi }_0`$.
In Green’s function or Auxiliary Field Monte Carlo (AFMC) , the overlap of the configuration with the trial state evolves smoothly with time. The change of the configurations $`\mathrm{\Psi }(𝐏_n)`$ per time step scales with $`\sqrt{\mathrm{\Delta }\tau }`$, which can be made arbitrarily small. If the wave function of the system is a purely real scalar quantity, any configuration which yields a negative overlap must first pass through a point at which $`\mathrm{\Psi }_T`$ and hence the overlap is zero. Discarding configurations at this point is sufficient to stabilize the simulation and produce an approximate solution $`\mathrm{\Psi }_C`$ to the many-fermion problem. It solves the many-body Schrödinger equation with the boundary conditions imposed by the nodes of $`\mathrm{\Psi }_T`$, and is known as the fixed-node approximation . The approximate $`\mathrm{\Psi }_C`$ is the best wave function (in the sense of lowest energy) with the same nodes as $`\mathrm{\Psi }_T`$. The discarded configurations are orthogonal not only to the trial state but also to the solution of the constrained problem $`\mathrm{\Psi }_C`$, and it has been shown that this method produces an upper bound to the ground state energy .
More generally, and particularly in nuclei, the trial wave function $`\mathrm{\Psi }_T`$ is a vector in spin-isospin space, and there are no coordinates for which all the spin-isospin amplitudes are zero. This is also true in AFMC, where the propagated configurations describe a fully antisymmetric product of single-particle wave functions. In both these cases, the evolution of the wave function remains smooth. In AFMC it is still possible to use a constraint of the form given in Eq. (50). However, the discarded configurations, while orthogonal to the trial state $`\mathrm{\Psi }_T`$, are not necessarily orthogonal to the the solution in the constrained space $`\mathrm{\Psi }_C`$. For this reason, the method does not produce an upper bound to the true ground-state energy .
A new difficulty specific to the nuclear problem is that the overlap $`\mathrm{\Psi }_{T,p}(𝐑_n)^{}\mathrm{\Psi }(𝐏_n)`$ is complex, and is estimated stochastically with a randomly selected order, denoted by the subscript $`p`$, of the operators in $`\mathrm{\Psi }_T`$, Eq. (31). These overlaps do not evolve smoothly and pass through zero. Therefore we can satisfy the constraint Eq. (50) only on the sum of discarded configurations ($`\mathrm{\Psi }_d`$), but not for individual configurations as in Eq. (49).
The fluctuations in $`\mathrm{\Psi }_{T,p}`$, due to the sampling of the pair operator orders, are small compared to the wave function itself, so we define an algorithm for discarding configurations which resembles as much as possible the fixed-node or constrained-path algorithms described above. Configurations in the GFMC are obtained with probability proportional to the importance function $`I_{T,p}=I[\mathrm{\Psi }(𝐏_n),\mathrm{\Psi }_{T,p}]`$, Eq. (43), which depends upon $`p`$. For each of them the overlap $`O_{T,p}`$ is defined as:
$$O_{T,p}=\mathrm{Re}[\mathrm{\Psi }(𝐏_n)^{}\mathrm{\Psi }_{T,p}(𝐑_n)],$$
(51)
so that the required constraint is the sum over discarded configurations of $`O_{T,p}/I_{T,p}=0`$. We define a probability $`P[\mathrm{\Psi }(𝐏_n),\mathrm{\Psi }_{T,p}(𝐑_n)]`$ for discarding a configuration in terms of the ratio $`O_{T,p}/I_{T,p}`$:
$$\begin{array}{ccccccccc}\hfill P[\mathrm{\Psi }(𝐏_n),\mathrm{\Psi }_{T,p}(𝐑_n)]& =& 0\hfill & & \hfill O/I& >& \alpha _c& & \\ & =& \frac{\alpha _cO/I}{\alpha _c\beta _c}\alpha _c\hfill & \hfill >& \hfill O/I& >& \beta _c& & \\ & =& 1\hfill & & \hfill O/I& <& \beta _c.& & \end{array}$$
(52)
According to this algorithm configurations with $`O/I`$ less than $`\beta _c`$ are always discarded, configurations with $`O/I`$ greater than $`\alpha _c`$ are never discarded, and there is a linear interpolation in between. The values of $`\alpha _c`$ and $`\beta _c`$ are held constant during the simulation. At each step, a configuration is discarded with a probability that depends upon the order of pair operators $`p`$. Hence an unbiased estimate of the overlap $`\mathrm{\Psi }_d|\mathrm{\Psi }_T`$ is required. It is obtained by choosing independent samples $`p^{}`$ of the pair operator order to compute the overlap:
$$\mathrm{\Psi }_d|\mathrm{\Psi }_T=\underset{\stackrel{\mathrm{discarded}}{\mathrm{configurations}}}{}O_{T,p^{}}/I_{T,p^{}}.$$
(53)
One or two preliminary runs are sufficient to adjust the constants $`\alpha _c`$ and $`\beta _c`$ such that this overlap is zero within statistical errors. When this condition is satisfied, it is trivial to show that the growth estimate of the energy, obtained from the growth or decay of the population with imaginary time, is identical to the mixed estimate. These two energy estimators are described in detail in Ref. .
Note that in absence of the term with $`ϵ`$ in the importance function $`I_{T,p}`$, Eq. (43), the ratio $`O_{T,p}/I_{T,p}`$ is $`\pm 1`$. In constrained path calculations, a significantly larger value of $`ϵ`$, typically 0.15, is used to reduce the fluctuations in $`O_{T,p}/I_{T,p}`$. For light nuclei, we have tried a variety of $`\alpha _c`$ and $`\beta _c`$ adjusted to the constraint without finding any statistically significant differences in the results. Indeed, even preliminary estimates of $`\alpha _c`$ and $`\beta _c`$ which do not precisely satisfy the constraint condition Eq. (50) yield results indistinguishable from the results with finely tuned constraint parameters. Typical values of $`\alpha _c`$ and $`\beta _c`$ for $`A`$ = 8 are 0.1 and 0.2 respectively.
In principle, two parameters are not needed to adjust the constraint to yield zero overlap. We were motivated to try this approach by the fact that, in standard fixed-node calculations, one defines a propagator which goes exactly to zero at the node. This results in some configurations being discarded which have small positive overlaps, as well as all configurations which have negative overlaps. In the fixed node scheme, this produces a result of higher order in the time step than simply discarding configurations with negative overlaps. For this reason both $`\alpha _c`$ and $`\beta _c`$ are chosen to be greater than zero. The algorithm described above is related to one proposed by Sorella in the context of condensed matter simulations , but simpler in that it uses only the overlap and the constraint parameters are held fixed during the simulation.
As discussed in the next subsection this method yields results with stable statistical errors independent of $`\tau `$. For $`A7`$, the calculated energy is very close to that obtained without constraints in most cases; however in some cases, particularly those with poor $`\mathrm{\Psi }_T`$’s, the calculated energy can be significantly below or above the true unconstrained result. The amplitudes of the excited states $`\mathrm{\Psi }_i`$, with energies $`E_i`$, in $`\mathrm{\Psi }_T`$ decay exponentially as $`\mathrm{exp}((E_iE_0)\tau )`$ in the unconstrained $`\mathrm{\Psi }(\tau )`$ without any change in their phase. Therefore these give a positive contribution to the mixed energy, making it $`E_0`$. The equality is obtained in the limit $`\tau 0`$. The constrained energy can have additional errors if, due to a poor choice of $`\mathrm{\Psi }_T`$, high-energy excitations get reintroduced in $`\mathrm{\Psi }(\tau )\mathrm{\Psi }_d`$. Such excitations do not necessarily have the same phase as in $`\mathrm{\Psi }_T`$ and can give contributions to the constrained $`E(\tau )`$ of either sign. Such errors can be easily detected by propagating $`\mathrm{\Psi }(\tau )\mathrm{\Psi }_d`$ without constraint for a limited number of steps $`n_u`$. Most of the $`\mathrm{\Psi }_T`$ dependence of the calculated energy is eliminated by fairly small $`n_u`$ without a significant increase in the statistical error. Thus the wave function $`\mathrm{\Psi }_C(\tau ,n_u)`$ used to estimate energies and other observables in the present calculations, is obtained from:
$$\mathrm{\Psi }_C(\tau =n\mathrm{}\tau ,n_u)=\{\mathrm{exp}[Hn_u\mathrm{}\tau ]\}_u\{\mathrm{exp}[H(nn_u)\mathrm{}\tau ]\}_c\mathrm{\Psi }_T,$$
(54)
where $`\{\mathrm{}\}_c`$ signifies propagation with the constraint and $`\{\mathrm{}\}_u`$ indicates normal propagation without constraints. As is shown in the next subsection, the simulations are stable for arbitrary $`n`$, and a fairly small $`n_u`$ ($`n_u\mathrm{}\tau 0.01`$ MeV<sup>-1</sup>) is sufficient to eliminate the $`\mathrm{\Psi }_T`$ dependence. Evaluation of matrix elements of this wave function can be implemented very easily on the computer by, after each propagation step, labeling those configurations which are to be discarded and then retaining them in the simulation for $`n_u`$ more steps. In principle this could be used to evaluate the overlap of the discarded configurations with $`\mathrm{\Psi }_T`$ after the unconstrained propagation. This could result in better values of the constraint parameters, however we find no change in the overlap for the values of $`n_u`$ we have considered.
### C Tests of Constrained Path
We have tested the constrained path algorithm in a variety of light nuclear systems, studying the dependence upon constraining wave function and also the convergence of the results obtained by relaxing the constraint. In some cases bad trial wave functions were used to test the algorithm under extreme conditions. In this subsection we present results for the <sup>6</sup>Li and <sup>8</sup>He ground states and for 8 neutrons bound in an external well. The tests for <sup>6</sup>Li and 8 neutrons were made using $`H=H^{}`$, Eq. (11), with no three-body potential. This eliminates uncertainties from the extrapolation of $`HH^{}`$, and allows us to have just the exact two-body propagator, $`g_{ij}`$. Eliminating the three-body potential also results in faster calculations allowing smaller statistical errors to be achieved. In the following all energies are in MeV and imaginary times in MeV<sup>-1</sup>; for simplicity the units are omitted in most cases.
Figure 2 shows the ground state energy of <sup>6</sup>Li calculated with two choices of $`\mathrm{\Psi }_T`$ \[Eq. (31)\]: (1) $`\mathrm{\Psi }_t`$, which is the full wave function with no $`\stackrel{~}{U}_{ij;k}^{TNI}`$ because there is no $`V_{ijk}`$, and (2) $`\mathrm{\Psi }_\sigma `$, a simplified wave function obtained by removing the tensor pair correlations leaving only spin-spin and isospin correlations. The subscript $`t`$ on $`\mathrm{\Psi }_t`$ emphasizes that it has the essential tensor correlations. In this and the following two figures, constrained-propagation results with zero and finite $`n_u`$ are shown as open symbols while solid symbols show completely unconstrained results or unconstrained continuations of constrained propagations. Averages of the last few energies are shown by solid lines with the corresponding statistical errors indicated by the surrounding dashed lines; the length of the line indicates the range of $`\tau `$ included in the average. Consider first the results using the full ($`\mathrm{\Psi }_t`$) wave function shown in the upper part of Fig. 2; these are typical of our production calculations. The solid circles show an unconstrained propagation from $`\tau `$ = 0 to $`\tau `$ = 0.1 of 200,000 configurations; the statistical errors on the last points are growing rapidly, but the energy seems to be fairly independent of $`\tau `$ beyond 0.05. The average from $`\tau `$ = 0.06 to $`\tau `$ = 0.1 is –28.16(12), as shown by horizontal solid and dashed lines. It should be close to the correct answer for this Hamiltonian. The open triangles show constrained results with $`n_u`$ = 0 for the same case. This simulation is stable with much smaller statistical errors than the unconstrained calculation and could be extended to an arbitrarily large $`\tau `$; here we find no significant change after $`\tau `$ = 0.05. The average over 0.14 $`\tau `$ 0.20 is –27.94(2). We released the constraint at $`\tau `$ = 0.20 and continued the propagation to $`\tau `$ = 0.27, as shown by solid triangles. The errors grow, and the average over this entire region is –28.15(5). It is likely that this result, which is consistent with the unconstrained result of –28.16(12) but with smaller statistical error, is the most accurate available. A constrained propagation with $`n_u`$ = 20 is shown by the open squares; only 50,000 configurations were used. This simulation is very similar to the $`n_u`$ = 0 one; the average is –28.01(5). These results suggest that constrained propagation with our typical $`\mathrm{\Psi }_t`$, using $`n_u`$ = 20, may lead to an error of +0.14(7) (or $``$0.5%) in the binding energy of <sup>6</sup>Li.
The results in the lower part of Fig. 2 were obtained using the $`\mathrm{\Psi }_\sigma `$ trial wave function. This is an extremely bad $`\mathrm{\Psi }_T`$ because, with no tensor correlations, this wave function results in an identically zero expectation value for the tensor potentials. Thus while the $`\tau `$ = 0.0 expectation value for $`\mathrm{\Psi }_t`$ is –23.95(4), that for $`\mathrm{\Psi }_\sigma `$ is +31.1(1). The solid circles show an unconstrained propagation of 400,000 configurations for this $`\mathrm{\Psi }_\sigma `$ case; the GFMC has managed to improve the $`\tau `$ = 0.0 energy by 59 MeV to -28.0(2), but the statistical errors are large even after using twice the number of configurations as for the unconstrained propagation with $`\mathrm{\Psi }_t`$. The open triangles show $`n_u`$ = 0 constrained propagation of 200,000 configurations. This propagation becomes stable after $`\tau `$ = 0.12, but it has converged to the significantly overbound result of –30.2(1). The solid triangles show the results of releasing the constraint at $`\tau `$ = 0.20; the energy immediately goes up to the correct value with an average of –28.0(2). Finally the open squares show $`n_u`$ = 20 constrained propagation of 400,000 configurations. The repeated use of 20 unconstrained steps gives a stable propagation that is quite accurate; the average is –28.27(7) with an error of –0.12(9). The unconstrained continuation of this calculation, shown by the solid squares, makes no significant change. These results show the need to have unconstrained steps before evaluation of the energy or any other observable.
The second test case is a system of 8 neutrons bound in an external one-body potential. We have previously reported results for such systems as a basis for comparing Skyrme models with microscopic calculations based on realistic interactions . As in the previous example, the neutrons interact via the $`v_{ij}^{}`$, with no three-nucleon interaction. Because systems of neutrons are not self binding, the Hamiltonian also includes an external one-body potential of Woods-Saxon form,
$$V_1(r)=\underset{i}{}\frac{V_0}{1+\mathrm{exp}[(r_ir_0)/a_0]};$$
(55)
the parameters are $`V_0=20`$ MeV, $`r_0=3.0`$ fm, and $`a_0=0.65`$ fm. Neither the external well nor the internal $`v_{ij}^{}`$ potential are individually attractive enough to produce a bound state of eight neutrons, however the combination does produce binding.
The Jastrow $`\mathrm{\Psi }_J`$ for neutron drops is just
$$|\mathrm{\Psi }_J=\underset{i<j}{}f_c(r_{ij})|\mathrm{\Phi }_A(JM).$$
(56)
where $`\mathrm{\Phi }_A`$ is a Slater determinant of single-particle orbitals. We compare results obtained with the full $`\mathrm{\Psi }_T`$, again refered to as $`\mathrm{\Psi }_t`$ because of the tensor operators, shown in the upper part of Fig. 3, with those obtained with the $`\mathrm{\Psi }_J`$, shown in the lower part.
The solid circles show results of unconstrained propagation using 400,000 configurations of the $`\mathrm{\Psi }_t`$ and 560,000 configurations of the $`\mathrm{\Psi }_J`$ trial wave functions, respectively. The errors are growing rapidly by $`\tau `$ = 0.06, and the calculations do not appear to have converged. Averages of the last four energies are –39.3(1) and –38.9(1) for $`\mathrm{\Psi }_t`$ and $`\mathrm{\Psi }_J`$ respectively, but these are clearly just upper bounds. Constrained propagations using the $`\mathrm{\Psi }_J`$ trial wave function and $`n_u`$ = 0 and $`n_u`$ = 20 are shown by the open triangles and squares, respectively in the lower part of the figure. These have stabilized beyond $`\tau `$ = 0.1, with end averages of respectively –39.58(4) and –39.79(6). Both propagations were continued without constraint; their results are shown by solid triangles and squares. The energy remains stable but the errors grow rapidly; the averages are –39.86(11) and –39.85(34). In this case constrained propagation with $`n_u`$ = 20 gives results that are below the best unconstrained (starting from $`\tau =0`$) upper bound with $`\mathrm{\Psi }_t`$ by 0.5(1). The constrained result is consistent with unconstrained continuations beyond $`\tau =0.25`$, indicating its reliability.
A different situation obtains with the variationally better $`\mathrm{\Psi }_t`$ trial wave function; the $`\tau `$ = 0, VMC energy for $`\mathrm{\Psi }_t`$, –35.30(4), is lower than the –30.62(9) given by $`\mathrm{\Psi }_J`$, but the difference is not as large as in the case of <sup>6</sup>Li. Constrained propagation using $`\mathrm{\Psi }_t`$ and $`n_u`$ = 0 is shown by the open triangles in the upper part; it stabilizes beyond $`\tau `$ = 0.1, but at a value that is 1 MeV too high. The results of releasing the constraint are shown by the solid triangles; the energy immediately drops 1 MeV and the final average is –39.71(10). Finally constrained propagation for $`\mathrm{\Psi }_t`$ and $`n_u`$ = 20 is shown by the open squares. The results are stable with an average of –39.74(5), in agreement with the best results obtained in the previous paragraph. Releasing the constraint does not result in any significant change, as is shown by the solid squares. This case also confirms the need to release the constraint before measurements, with the surprising result that a variationally better wave function may not necessarily provide better constraints to guide the GFMC. Nevertheless, constrained propagation with $`n_u`$ = 20 gives results with presumably less than 1% error, for both trial wave functions.
Finally, in Fig. 4 we show results for the <sup>8</sup>He ground state with the full Argonne $`v_{18}`$ plus Urbana IX three-nucleon interaction. We are presenting this case because we find that <sup>8</sup>He is the most difficult nucleus we have studied, in terms of obtaining reliable error estimates, and the need for $`n_u>0`$. All of the calculations are made with the full $`\mathrm{\Psi }_T`$ of Eq. (31). In this figure the results of a standard calculation without constraint (solid circles) are compared to those with constrained-path $`n_u`$ = 0 (open triangles) and $`n_u`$ = 20 (open squares), and their unconstrained continuation beyond $`\tau `$ = 0.2 shown by solid triangles and squares. The average unconstrained energy in the $`0.03\tau 0.06`$ range, –26.1(3), is clearly an unconverged upper bound. The $`n_u`$ = 0 and 20 end averages, –26.89(9) and –27.16(15), are below it by $``$1 MeV. The averages of the results obtained after the constraint is released at $`\tau `$ = 0.2, –27.5(4) and –26.9(2), are not significantly different from the constrained averages.
From these, and many other, tests we conclude that constrained propagation, including unconstrained steps prior to measurement, yields results that are reliable. In the $`A=8`$ nuclei and neutron drops they are up to 4% below the unconverged upper bounds that can be obtained by unconstrained propagation up to $`\tau =0.06`$, and have smaller statistical errors. The constrained path results with $`n_u10`$ to 20 are stable within 1% with respect to reasonable, and in case of <sup>6</sup>Li even unreasonable, changes in $`\mathrm{\Psi }_T`$ used to constrain the paths. Our present practice is to use the full $`\mathrm{\Psi }_T`$, Eq. (31) in constrained path calculations with $`n_u`$ = 20 for the neutron-rich He isotopes and neutron drops and $`n_u`$ = 10 for all other nuclei where even $`n_u=0`$ calculations seem to be fairly accurate.
### D Resonance states in constrained-path algorithms
Since the constrained-path algorithm is stable to large imaginary time, it is useful to consider the large-$`\tau `$ behavior of the energy of resonant states. The bound-state simulations described above yield asymptotically stable energies out to very large imaginary times. Resonant states, however, are more delicate. In principle these states should decay to separated clusters, and in fact this does occur with the constrained-path algorithm. The rate of this decay presumably depends not only on the resonance energy but also on the width of the state, and the trial state wave function used to impose the constraint.
We have studied several cases including the two unbound p-wave states in <sup>5</sup>He. The GFMC energy consistently decreases with $`\tau `$ for both the $`J^\pi =1/2^{}`$ and the $`3/2^{}`$ states. The 1/2<sup>-</sup> state, lying higher in energy and with a larger width, decays more quickly than the $`3/2^{}`$ state. We have verified that the system breaks apart into a separated $`\alpha `$ plus $`n`$ by studying the $`pp`$ pair distribution function, which depends only upon the internal structure of the $`\alpha `$-particle. This distribution remains constant out to the largest imaginary times (3 MeV<sup>-1</sup>) studied, while the $`pn`$ and $`nn`$ distributions steadily become broader.
We have also studied three low-lying states in <sup>6</sup>Li out to large $`\tau `$. The $`J^\pi =1^+`$ ground state is stable, the $`3^+`$ excitation is a narrow resonance, while the $`2^+`$ state lies higher in energy and is broader. The GFMC energies of these states are displayed in Fig. 5 as a function of $`\tau `$. The ground state is clearly stable out to the largest imaginary times. The narrow $`3^+`$ shows a plateau in $`E(\tau )`$, decreasing only modestly in energy, though of course it will eventually also decay to the $`\alpha `$+d threshold energy. The $`2^+`$ state is much broader and decays much more quickly in imaginary time. Therefore, comparison of the energies obtained for such broad states with experiment could be misleading. The $`\alpha `$+d threshold energy for this Hamiltonian is shown in the figure; it is clear that both excited states are far from convergence to this energy.
It is possible to compute energies that may be directly compared with experiment for resonant states that have only a single two-body channel available for breakup . Scattering observables, including scattering length, effective range, and phase shifts, can be calculated directly by imposing scattering boundary conditions on the asymptotic wave function. This is quite important for studying a variety of interesting low energy reactions, including parity-violation and important astrophysical reactions, but has not been considered in this work.
### E Three-body Propagator
The small-time-step propagator of Eq. (41) involves two complete evaluations of the three-body potential (two sums over all triples) and one product of all pair propagators for each time step. Thus as the number of nucleons is increased, the time spent in the three-body part of the propagation becomes a larger and larger fraction of the total time for the calculation. For this reason it is desirable to find a less costly treatment of the three-body propagator.
It was noted in Eq. (31) that the full $`U_{ijk}^{TNI}`$ can be replaced with a $`\stackrel{~}{U}_{ij;k}^{TNI}`$ that omits the commutator part of $`V_{ijk}^{2\pi }`$ with very little degradation of the variational energy, and that the resulting correlation is an operator in the spins and isospins of only two nucleons. This led us to consider the combined two- and three-body propagator (we omit the spin-isospin indices):
$`\stackrel{~}{G}(𝐑,𝐑^{})`$ $`=`$ $`e^{E_o\mathrm{}\tau }𝒮{\displaystyle \underset{i<j}{}}{\displaystyle \frac{\stackrel{~}{g}_{ij}(𝐫_{ij},𝐫_{ij}^{})}{g_{0,ij}(𝐫_{ij},𝐫_{ij}^{})}}`$ (57)
$`\stackrel{~}{g}_{ij}(𝐫_{ij},𝐫_{ij}^{})`$ $`=`$ $`\left[1{\displaystyle \frac{\mathrm{}\tau }{2}}{\displaystyle \underset{ki,j}{}}\alpha V_{ij;k}^{2\pi ,A}(𝐑)\right]g_{ij}(𝐫_{ij},𝐫_{ij}^{})\left[1{\displaystyle \frac{\mathrm{}\tau }{2}}{\displaystyle \underset{ki,j}{}}\alpha V_{ij;k}^{2\pi ,A}(𝐑^{})\right]`$ (59)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\mathrm{}\tau }{2}}{\displaystyle \underset{ki,j}{}}(V_{ijk}^R(𝐑)+V_{ijk}^R(𝐑^{}))\right].`$
The $`V_{ij;k}^{2\pi ,A}`$ is defined in Eq. (9) as the anticommutator part of $`V_{ijk}^{2\pi }`$, and $`\alpha `$ is chosen so that
$`\alpha V^{2\pi ,A}=V^{2\pi };`$ (60)
typically $`\alpha =1.6`$. This $`\stackrel{~}{g}_{ij}`$ can be reduced to a single $`4\times 4`$ operator in $`i,j`$ spin space for each of the two isospin states, and thus takes little more time to evaluate than just $`g_{ij}`$. Propagation with just $`\stackrel{~}{G}`$ gives much the same results as with the much more costly propagator of Eq. (41), but we attempt to make a more reliable propagation by using the following sequence as the basic propagation step:
$`\stackrel{~}{G}(𝐑_{i+n},𝐑_{i+n1})\stackrel{~}{G}(𝐑_{i+n1},𝐑_{i+n2})\mathrm{}\stackrel{~}{G}(𝐑_{i+n/2+1},𝐑_{i+n/2})`$ (61)
$`\times \{1n\mathrm{}\tau {\displaystyle \underset{i<j<k}{}}[V_{ijk}^{2\pi }(𝐑_{i+n/2})\alpha V_{ij;k}^{2\pi ,A}(𝐑_{i+n/2})]\}`$ (62)
$`\times \stackrel{~}{G}(𝐑_{i+n/2},𝐑_{i+n/21})\mathrm{}\stackrel{~}{G}(𝐑_{i+1},𝐑_i),`$ (63)
where $`n`$ is a small (typically 4) number of steps. Here we go from $`𝐑_i`$ to $`𝐑_{i+n}`$ by making $`n/2`$ steps with $`\stackrel{~}{G}`$; then applying $`n`$ times the correction due to the difference of the complete $`V^{2\pi }`$ and the approximate $`\alpha V^{2\pi ,A}`$, both computed at the position $`𝐑_{i+n/2}`$; and finally making another $`n/2`$ steps to $`𝐑_{i+n}`$.
Table V shows the reliability of this method for <sup>6</sup>Li(gs). The first line gives results with the old method, Eq. (41), and the next three lines use the new method with increasing values of $`n`$ \[$`n=\mathrm{}`$ means that no correction for the approximation in $`\stackrel{~}{G}`$ is made with Eq. (63)\]. The last line shows results using just a two-body propagator; that is, the three-body potential is included only perturbatively. Clearly a three-body propagator is essential for obtaining results with $`<1\%`$ error, however the $`\stackrel{~}{G}`$ is reliable to better than 1%.
### F Numerical Evaluations
The p-shell GFMC calculations reported here were made with constrained-path propagation to $`\tau =0.2`$ MeV<sup>-1</sup> using 400 steps of $`\mathrm{}\tau =0.0005`$ MeV<sup>-1</sup>. Energies and other observables were evaluated every 20 steps and the last 7 ($`\tau 0.14`$ MeV<sup>-1</sup>) values were averaged. At least 10 unconstrained steps were taken before the observables were computed (20 steps were used for <sup>8</sup>He). The $`A=8`$ results used 10,000 to 20,000 initial configurations. Once the propagation has stabilized (typically by $`\tau =0.1`$ MeV<sup>-1</sup>), the constraint removes from 0.7% \[for <sup>8</sup>Be(gs)\] to 1.7% \[for <sup>8</sup>Be(3<sup>+</sup>)\] of the configurations at each propagation step (the percentages for <sup>8</sup>He and <sup>8</sup>Li are in this range). The removed configurations are replaced by new ones generated by branching with the average branching probability chosen to maintain an approximately constant population.
Most of the GFMC calculations reported here were made on the 128-node SGI Origin 2000 in the Mathematics and Computer Science division of Argonne National Laboratory. The individual nodes in this machine are 250 MHz R10000 processors which (when the processors are not being shared with other users and memory remains local to each processor) deliver sustained speeds of $`200`$ MFLOPS for the 8-nucleon calculations. This results in a 10,000 configuration calculation for <sup>8</sup>Be taking about 460 node-hours. The approximate times required for other nuclei can be determined from Table VI. The columns of this table show the number of nucleons, the number of pairs, and the size of the spin-isospin vector. We find that the total calculational time is proportional to the product of these three numbers; this product (scaled to 1 for <sup>8</sup>Be) is given in the last column.
Typically we have computed the $`M=J`$ state of a given nucleus; this allows us to directly evaluate the spectroscopic quadrupole moment and magnetic moment. Recently we have realized that if we compute the $`M=0`$ state of even-$`A`$ nuclei, the size of the spin-isospin vector can be reduced by a factor of two. This is done by observing that for even $`A`$,
$$\mathrm{\Psi }_{s_1,s_2,\mathrm{},s_A}(J,M=0)=(1)^{\frac{1}{2}A+JM_S}\mathrm{\Psi }_{s_1,s_2,\mathrm{},s_A}^{}(J,M=0),$$
(64)
where $`s_i`$ is the spin-projection of nucleon $`i`$ and $`M_S=s_i`$. Thus only half of the spin-isospin vector needs to be computed and stored, resulting in a saving of half the computer time for even $`A8`$ nuclei. This saving is not included in the times discussed in the previous paragraph. Unfortunately, the corresponding relation in odd-$`A`$ nuclei relates $`\mathrm{\Psi }_{s_i}(J,M)`$ to $`\mathrm{\Psi }_{s_i}(J,M)`$ and thus does not reduce the computational effort.
## V ENERGY RESULTS
### A Ground-State and Excitation Energies
In this Section we present VMC and GFMC results for the Hamiltonian consisting of Argonne $`v_{18}`$ plus Urbana IX. Figure 6 shows the $`E(\tau )`$ for the lowest $`T=0`$, $`J^\pi =0^+,2^+,4^+,1^+,3^+`$ states of <sup>8</sup>Be. The solid and dashed lines show the average energy and its statistical error; these numbers are reported in Tables VII and VIII. The $`E(\tau )`$ are
$$E(\tau )=H^{}(\tau )_{Mixed}+H(\tau )H^{}(\tau ),$$
(65)
where the $`HH^{}`$ is perturbatively extrapolated by Eq. (47). Figures 7 and 8 show the corresponding <sup>8</sup>Li and <sup>8</sup>He results. In all cases for which the states are experimentally narrow or stable, the $`E(\tau )`$ rapidly decreases for small $`\tau `$ and stabilizes before $`\tau =0.1`$ MeV<sup>-1</sup>. It is clear that $`E(\tau )`$ has not converged by $`\tau =0.2`$ MeV<sup>-1</sup> for the experimentally very broad <sup>8</sup>Be(4<sup>+</sup>) state, making it not possible to determine an accurate excitation energy for this state. The other broad states \[<sup>8</sup>Be(2<sup>+</sup>) and <sup>8</sup>He(2<sup>+</sup>)\], and the experimentally unknown <sup>8</sup>He(1<sup>+</sup>) state seem reasonably converged.
Table VII shows the computed and experimental ground-state energies. The errors shown in parentheses are only the Monte Carlo statistical errors; the systematic errors discussed in the previous section and in Ref. could add an additional 1% to the GFMC error. This paper is our first formal publication of $`A`$ = 8 results. Results for $`A=6`$ and 7 were published in Ref. ; most of the corresponding values in Table VII have been recomputed using improved VMC wave functions and GFMC propagation. Most of the VMC values have not changed significantly from those in Ref. ; the exception is <sup>7</sup>He which is now constrained to what we think is a more reasonable rms radius and actually has a higher energy. The $`A`$ = 6 and 7 GFMC energies have now all been computed with constrained propagation to $`\tau =0.2`$ MeV<sup>-1</sup> instead of the unconstrained propagation to only $`\tau =0.06`$ MeV<sup>-1</sup> used in Ref. . This has resulted in a lowering of <sup>6,7</sup>He energies by 0.47(17) and 0.63(23) MeV respectively; other energies changed by less than 1%. For $`A`$ = 8 nuclei, the GFMC improves the $`\mathrm{\Psi }_T`$ energy by $``$10 MeV; our best $`\mathrm{\Psi }_V`$ provides only $``$1.5 MeV of this improvement.
Table VIII shows excitation energies for the $`A`$ = 8 nuclei. The errors are the combined statistical errors of the energies computed for the given state and the ground state. We see that, while the absolute VMC energies have substantial errors as compared to the GFMC values, the VMC excitation energies are typically within an MeV of the GFMC values.
Figure 9 shows the energies of nuclear states for $`4A8`$ and Fig. 10 shows the corresponding excitation spectra. As is the case for the ground-state values, most of the $`A`$ = 6 and 7 excited-state energies have been recomputed with improved VMC wave functions and all of the GFMC propagations have been made to $`\tau =0.2`$ MeV<sup>-1</sup>. The VMC excitation energies of the lowest <sup>7</sup>Li($`\frac{1}{2}^{}`$) and <sup>7</sup>Li($`\frac{5}{2}^{}`$) states have been significantly reduced from the values reported in Ref. . These are the result of improved mixings of the different symmetry states in the variational wave function.
The figures confirm the principal $`A`$ = 6 and 7 results of Ref. for $`A`$ = 8: the Hamiltonian consisting of Argonne $`v_{18}`$ plus Urbana IX underbinds nuclei in the p-shell with the underbinding becoming worse as one increases $`A`$ or $`NZ`$. However the predictions of the excitation spectra are generally reasonable; the rms error of all of the 16 excitation energies computed by GFMC is only 650 keV; for the 9 states with experimental width less than 200 keV it is 540 keV. Considering that no parameters were adjusted to fit these p-shell energies, this is quite respectable. We do note that, where there is an experimental value to compare to, the VMC excitation energies of second excited states of a given $`J^\pi `$, which are very difficult to compute by GFMC, are usually too high.
Although the Argonne $`v_{18}`$ plus Urbana IX Hamiltonian consistently underbinds the p-shell nuclei, the errors are small compared to the magnitudes of the potential energies. Table IX shows the perturbatively extrapolated GFMC expectation values for the kinetic energy and potential energy terms. As is discussed in Ref. , the perturbatively extrapolated GFMC terms do not add up to the total energy, which is the most reliable number computed by this method. The total two-body potential, $`v_{ij}`$ is dominated by the one-pion exchange term, $`v_{ij}^\pi `$. There is a large cancellation between the kinetic and two-body potential energies so that their sum is only 15% to 20% of the two-body potential. The total three-body potential is typically 4% to 5% of the two-body potential. The repulsive $`V_{ijk}^R`$ typically cancels 45% of the $`V_{ijk}^{2\pi }`$. As an indication of the smallness of the errors produced by this Hamiltonian, consider the ground state of <sup>8</sup>Be: the difference of the computed and experimental energies is 2.1 MeV; the two-body potential is $``$301 MeV, $`V_{ijk}^{2\pi }`$ = $``$27 MeV, and $`V_{ijk}^R`$ = +12.3 MeV. Thus the corrections that have to be made to $`V_{ijk}`$ are significantly smaller than the terms already included.
### B Isobaric Analog States
Energy differences of isobaric analog states are sensitive probes of the charge-independence-breaking parts of the Hamiltonian. To study these it is useful to express the energies in an isobaric multiplet, characterized by $`A`$ and $`T`$, in terms of the isospin multipole operators of order $`n`$:
$$E_{A,T}(T_z)=\underset{n2T}{}a_{A,T}^{(n)}Q_n(T,T_z).$$
(66)
The $`Q_n(T,T_z)`$ are orthogonal functions for projecting out isovector, isotensor, and higher-order terms ; the first terms are $`Q_0=1`$, $`Q_1=T_z`$, and $`Q_2=\frac{1}{2}[3T_z^2T(T+1)]`$. The coefficients $`a^{(n)}`$ are then obtained from the calculated energies:
$$a_{A,T}^{(n)}=\underset{T_z}{}Q_n(T,T_z)E_{A,T}(T_z)/\underset{T_z}{}Q_n^2(T,T_z),$$
(67)
or perturbatively from expectation values of the isomultipole operators present in the Hamiltonian.
In first-order perturbation theory, the electromagnetic interaction contributes to the $`a^{(n)}`$ for $`n=0`$, 1, and 2, the kinetic energy to $`n=0`$ and 1, the nuclear CSB potential to $`n=1`$, and the nuclear CD potential to $`n=2`$. Because a significant portion of the $`v_{ij}^{CD}`$ comes from one-pion exchange, there should also be a CD component to $`V_{ijk}^{2\pi }`$; however, a plausible extension of the Urbana IX model gives negligible contributions of 3 keV or less to the $`n=2`$ terms. The $`a^{(n)}`$ for higher $`n`$ are zero in first order with our Hamiltonian, and there is little experimental evidence for $`n3`$ terms in nuclei . We have made VMC calculations of the $`a^{(1)}`$ and $`a^{(2)}`$ in first order by using a CI wave function of good isospin, $`T`$, and simply varying $`T_z`$ to compute the $`E_{A,T}(T_z)`$; in the GFMC calculations we evaluate expectations of the isomultiplet operators for the $`T_z=T`$ nuclei. The CSB and CD parts of the Hamiltonian can induce corresponding changes in the nuclear wave functions, leading to higher-order perturbative corrections to the splitting of isospin mass multiplets. However, it is difficult for us to estimate these higher order effects reliably in either the VMC or GFMC calculations.
Table X shows results for the $`T=1`$ and 2 isovector and isotensor coefficients in $`A=8`$ nuclei compared to experiment. In obtaining the experimental $`a_{8,1}^{(2)}`$ coefficient, we have used the average of the two isospin-mixed (2<sup>+</sup>;0+1) states; our calculation of this mixing is discussed in the next subsection. The contributions of the complete $`v^\gamma `$, $`v^{CD}`$, $`v^{CSB}`$, and $`K^{CSB}`$ terms to the $`a^{(n)}`$ are also given. The results show that the present Hamiltonian and VMC wave function match the experimental CSB and CD of the $`T=1`$ multiplet fairly well, but the calculated value for $`a_{8,2}^{(1)}`$ is about 10% too small. The GFMC consistently reduces the coefficients, mostly by reducing the Coulomb term. For the isovector coefficients, this worsens the comparison with experiment, however the case of $`a_{8,2}^{(2)}`$ is significantly improved. In all these calculations, the Coulomb interaction between protons is the dominant contribution to the $`a^{(n)}`$, but the strong interaction CSB terms serve to bring the final results closer to experiment. However, in view of the fact that the $`A=8`$ nuclei are underbound with the present Hamiltonian, it is premature to use these calculations as a precise test of the charge-independence-breaking components of the interaction.
### C Isospin Mixed States
The (2<sup>+</sup>;1) isobaric analog of the <sup>8</sup>Li ground state in <sup>8</sup>Be is very close in energy to the second (2<sup>+</sup>;0) excitation, as shown in Fig. 1, and is in fact experimentally observed to be isospin-mixed with it. There are also fairly close (1<sup>+</sup>;0,1) and (3<sup>+</sup>;0,1) pairs at slightly higher energies in the <sup>8</sup>Be spectrum. Our VMC calculations do not get the (2<sup>+</sup>;0,1) states so close together, although the other pairs do come out quite near each other. However, we can calculate the isospin-mixing matrix elements that connect these pairs of states:
$$E_{01}(J)=\mathrm{\Psi }(J^+,0)|H|\mathrm{\Psi }(J^+,1).$$
(68)
This is done by increasing the model space to include both $`T=0`$ and 1 components, which corresponds to a wave function vector with 10,752 terms, and employing the same off-diagonal evaluations used to determine the $`\beta _{LS[n]}`$ wave function components.
Results for the $`E_{01}(J)`$ are given in Table XI. The experimental values are determined from the observed decay widths and energies . The dominant contribution, from the Coulomb potential, typically accounts for less than half of the matrix element. We find the magnetic moment part of the electromagnetic interaction and the strong CSB interaction can provide a significant boost, although we still underpredict the experimental mixing by $``$20%. The spatial symmetry components of the different wave functions, as given in the Tables III and IV, for the (2<sup>+</sup>;0,1) and (3<sup>+</sup>;0,1) pairs are indeed fairly close to each other, since they have similar sizes and signs for the largest components. However, the (1<sup>+</sup>;0,1) states are not so similar, particularly due to the large $`T=1`$ <sup>1</sup>P component that is not available to the $`T=0`$ state. This may be why the $`E_{01}(1)`$ is noticeably smaller, which shows up through the change in sign of the magnetic moment contribution.
## VI MOMENTS AND DENSITY DISTRIBUTIONS
The proton rms radii, magnetic moments, and quadrupole moments for the $`A=8`$ nuclear ground states are given in Table XII. The calculations have been made in impulse approximation for both the initial VMC wave function and perturbatively after GFMC constrained path propagation. Comparing the VMC and GFMC radii, it appears that the variational wave functions are a little too compact for this Hamiltonian. However, given that the present interaction underbinds these nuclei, the actual radii should be smaller than the GFMC values. Attempts have been made to determine the matter radius of <sup>8</sup>He by proton-scattering experiments in inverse kinematics. Interpretation of the data is not model-independent, however, and has resulted in estimates ranging from 2.45 fm to 2.6 fm . The corresponding VMC and GFMC matter radii from Table XII are 2.69 fm and 2.92 fm, respectively. In the not too distant future, it may be possible to trap <sup>8</sup>He, <sup>8</sup>Li, and <sup>8</sup>B long enough to determine their charge radii to high precision by atomic means.
The magnetic moments of <sup>8</sup>Li and <sup>8</sup>B have been measured by $`\beta `$-radiation detection of implanted polarized ions . Our IA calculations should be supplemented by meson-exchange currents (MEC); their contributions are dominantly isovector, and have been shown to change the magnetic moments of <sup>3</sup>H and <sup>3</sup>He by $`\pm `$ 0.4 $`\mu _N`$ . The isoscalar average of the calculated magnetic moments is close the experimental value, and it is plausible that MEC contributions will bring both magnetic moments into good agreement with experiment. The quadrupole moments of <sup>8</sup>Li and <sup>8</sup>B have been measured by $`\beta `$-NMR and nuclear quadrupole resonance techniques . The VMC and GFMC results are fairly close for <sup>8</sup>Li, and just a little above the experimental value. The calculations for <sup>8</sup>B are significantly further apart, despite the proton rms radii differing by only a small amount; nevertheless they bracket the experimental value. We have also calculated the hexadecapole moment for these two ground states, and find it to be consistent with 0.
One- and two-nucleon density distributions have been calculated in both VMC and GFMC. The GFMC densities are a little more spread out and less peaked than the corresponding VMC densities, as reflected in the charge radii differences; here we show only the GFMC densities. Single-nucleon density distributions for the $`A=8`$ nuclei are shown in Fig. 11; they are normalized such that the integrated value equals the appropriate total value of $`N`$ or $`Z`$. The two protons in <sup>8</sup>He are the most peaked distribution, which should be expected on the grounds that they are mostly confined to the $`\alpha `$ core, while the six neutrons in <sup>8</sup>He have the broadest distribution. The proton distribution in <sup>8</sup>Li is also fairly peaked near the origin, but is broader than in <sup>8</sup>He since there is one additional proton in the p-shell. The neutron distribution in <sup>8</sup>Li is comparable to the <sup>8</sup>He neutrons near the origin, is slightly larger in the 1-2 fm range, and then falls below at larger distances; the intermediate-range excess may be due to the significantly greater binding of <sup>8</sup>Li compared to <sup>8</sup>He. In contrast, the <sup>8</sup>Be proton and neutron distributions are much less peaked at the origin and are rather flat out past 1 fm. This could be because <sup>8</sup>Be has a significant 2$`\alpha `$ component in its intrinsic structure, with the two $`\alpha `$’s sitting side-by-side; as discussed in the next section.
A logarithmic plot of the single-nucleon densities for <sup>4,6,8</sup>He isotopes is shown in Fig. 12. The neutron halos in <sup>6</sup>He and <sup>8</sup>He are clearly evident, while the proton cores of these nuclei are nearly identical. The peak neutron and proton distributions in <sup>6</sup>He and <sup>8</sup>He are much reduced compared to <sup>4</sup>He, because of the motion of the $`\alpha `$ core against the center-of-mass. Two-proton density distributions for <sup>4,6,8</sup>He are shown in Fig. 13. The relative proton-proton density in <sup>6,8</sup>He is not effected by the motion of the $`\alpha `$ relative to the center of mass; thus $`\rho _{pp}`$ is an indicator of changes in the internal structure of the $`\alpha `$ core in <sup>6,8</sup>He. Here we note that there is only a small change in the $`\alpha `$ core of <sup>6,8</sup>He, i.e., about a 10% reduction in the peak compared to <sup>4</sup>He, which must be compensated in the long-range tail, since the total integral under the curves is unity in each case. This reduction could be due to swelling of the $`\alpha `$ core, or due to charge-exchange interactions between the protons and p-shell neutrons.
We have also calculated the single-nucleon momentum distributions shown in Fig. 14, using the VMC wave functions and the method of Ref. , with some algorithmic improvements. The $`N(k)`$ are normalized to $`N/(2\pi )^3`$ or $`Z/(2\pi )^3`$. All the distributions show a remarkably similar structure: a low-momentum core attributable to s- and p-shell orbitals, followed by a high-momentum tail beyond 2 fm<sup>-1</sup> that is the sum of many small-amplitude higher-orbital contributions. This high-momentum tail is already evident in the <sup>4</sup>He momentum distribution and is intimately related to the D-state of <sup>4</sup>He, present due to the strong tensor forces and corresponding correlations in the wave function. The <sup>8</sup>He proton distribution is very similar to that of <sup>4</sup>He, showing that the $`\alpha `$ core is not much altered, while the neutrons exhibit a p-shell peak near 0.5 fm<sup>-1</sup>. The <sup>8</sup>Li proton distribution is a little broader, as might be expected by the addition of a p-shell proton, while the neutrons show a smaller p-shell peak. The <sup>8</sup>Be momentum distribution is almost exactly double that of <sup>4</sup>He except at very low momenta, indicating a significant 2$`\alpha `$ structure.
## VII INTRINSIC SHAPES
The ground state and first two excited states of <sup>8</sup>Be have an approximate rotational energy spectrum, and are assumed to be well approximated as two $`\alpha `$’s rotating around their common center of mass. This structure is not manifest in the shell model part of the VMC wave functions, in which the <sup>8</sup>Be states are constructed from an $`\alpha `$-like core surrounded by four p-shell nucleons coupled to the appropriate total angular momentum. In this section we describe an attempt to study the intrinsic structure of the <sup>8</sup>Be states described by the VMC wave functions including correlations.
The standard Monte Carlo method for computing one-body densities, $`\rho (𝐫)`$, is to make a random walk that samples $`|\mathrm{\Psi }(𝐫_1,𝐫_2,\mathrm{},𝐫_A)|^2`$ and to bin $`𝐫_1,𝐫_2,\mathrm{},𝐫_A`$ for each configuration in the walk. The density is then proportional to the number of samples in each bin. In the case of a $`J`$ = 0 nucleus, this “laboratory” density will necessarily be spherically symmetric.
We can attempt to find the intrinsic density in body-fixed coordinates by computing the moment of inertia matrix:
$$=\underset{i=1}{\overset{A}{}}\left(\begin{array}{ccc}x_i^2& x_iy_i& x_iz_i\\ y_ix_i& y_i^2& y_iz_i\\ z_ix_i& z_iy_i& z_i^2\end{array}\right),$$
(69)
for each configuration. We then find the eigenvalues and eigenvectors of $``$, rotate to those principal axes, and bin the resulting $`𝐫_1^{},𝐫_2^{},\mathrm{},𝐫_A^{}`$. The eigenvector with the largest eigenvalue is chosen as the $`𝐳^{}`$ axis; further choices of the $`\pm `$ direction along $`𝐳^{}`$ and which eigenvectors to assign to $`𝐲^{}`$ and $`𝐱^{}`$ may also be made or averaged over. This procedure will not produce a spherically symmetric distribution, even if there is no underlying deformed structure, because almost every random configuration will have principal axes of different lengths and the rotation will always orient the longest principal axis in the $`𝐳^{}`$ direction. We have made a number of tests using simple wave functions with no internal correlations and also nuclei like <sup>4</sup>He which, in our models, should have no intrinsic deformations. We find that the projected “intrinsic” density for such cases is always prolate but that no other artificial structure is introduced.
When the above procedure is applied to the <sup>8</sup>Be rotational states, a dramatic intrinsic structure is revealed. Figure 15 shows this calculation for the ground state. The figure shows contours of constant density plotted in cylindrical coordinates, with the z-axis being the axis of quantization. The calculation was made for the VMC wave function. The left side of the figure shows the standard density calculation; we can think of this as the density in the “laboratory” frame. For the $`J`$ = 0 ground state, this is spherically symmetric as shown. The right side of the figure shows the intrinsic density computed as described above. In this case the orientation (up or down) along $`𝐳^{}`$ and around $`𝐳^{}`$ ($`r=\sqrt{x^2+y^2}`$) was averaged over. It is clear that the intrinsic density has two peaks, with the neck between them having only one-third the peak density; we regard these as two $`\alpha `$’s.
Figure 16 shows the corresponding calculation for the $`J`$ = 4<sup>+</sup>, $`M_J`$ = 4, state of <sup>8</sup>Be. In this case the the laboratory density does not have to be spherically symmetric and there is evidence of two $`\alpha `$’s rotating around the $`𝐳`$ axis in the equatorial plane. The projection to the intrinsic frame rotates these to the $`𝐳^{}`$ axis and results in an intrinsic density that is insignificantly different from that obtained for the ground state. A calculation for the $`J`$ = 2<sup>+</sup> state also produces the same intrinsic density.
These results, obtained for the VMC wave functions, suggest that the 0<sup>+</sup>, 2<sup>+</sup>, and 4<sup>+</sup> wave functions for <sup>8</sup>Be have the structure of a deformed rotor consisting of two $`\alpha `$’s. The structure is not manifest in the shell model part of the VMC wave functions; it is induced by the correlations. The optimum spatial correlations between pairs of s-shell and p-shell nucleons are similar, but those between one s-shell and one p-shell nucleon are different. Calculations using GFMC configurations give very similar results except that, especially for the $`J`$ = 4<sup>+</sup> state, the rms radius is still growing at the end of the GFMC propagation.
If the $`0^+`$, $`2^+`$, and $`4^+`$ states are generated by rotations of a common deformed structure, then their electromagnetic moments and transition strengths should all be related to the intrinsic moments. Table XIII shows the VMC computed values of the quadrupole and hexadecapole moments, and of the B(E2) and B(E4) strengths for these states. In the rotational model these are related by simple Clebsch Gordon coefficients to the intrinsic quadrupole, $`Q_0`$, or $`M_4`$ moment depending on the multipolarity, $`\lambda `$. These moments are defined as
$$Q_0=\sqrt{\frac{16\pi }{5}}𝑑𝐫^{}\rho (𝐫^{})𝐫_{}^{}{}_{}{}^{2}Y_{20}(\widehat{𝐫}^{}),$$
(70)
and
$$M_4=𝑑𝐫^{}\rho (𝐫^{})𝐫_{}^{}{}_{}{}^{4}Y_{40}(\widehat{𝐫}^{}),$$
(71)
where $`\rho `$ is the point proton density and $`𝐫^{}`$ refers to the intrinsic (body-fixed) frame. The last column gives the extracted values of these intrinsic moments; we see that, with the exception $`M_4`$ for the $`J=2^+`$ state, these are remarkably constant.
We can also compute values of $`Q_0`$ and $`M_4`$ by integrating over the projected body-fixed densities. These values might be too large because, as described above, the projection method can introduce an excessive deformation in the intrinsic shape. The values for the $`0^+`$, $`2^+`$, and $`4^+`$ states are $`Q_0`$ = 26.2, 27.9, and 26.7, and $`M_4`$ = 55, 62, 64; the $`Q_0`$ are in good agreement with the values given in the Table, which were obtained from the spectroscopic quadrupole moments. We note that the ratios of these $`Q_0`$ and $`M_4`$ values are in reasonable agreement with the ratio obtained for a diatomic model of <sup>8</sup>Be assuming two point $`\alpha `$’s separated by 4 fm.
We can attempt to project out an intrinsic structure of the $`T`$ = 0, $`J`$ = 1<sup>+</sup> and $`J`$ = 3<sup>+</sup>, states of <sup>8</sup>Be. In this case we also find two peaks but of somewhat smaller density than for the 0<sup>+</sup>, 2<sup>+</sup>, and 4<sup>+</sup> states. Clearly these peaks cannot be due to two $`\alpha `$’s rotating around each other. They presumably reflect occurrence of $`\alpha +t+p`$ and $`\alpha +^3`$He$`+n`$ structures in these states.
## VIII CONCLUSIONS
We have made quantum Monte Carlo calculations for the ground states and low-lying excitations of $`A=8`$ nuclei interacting by realistic two- and three-nucleon potentials. These calculations have been made practical by the development of a constrained-path algorithm for the complex spatial and spin-isospin wave functions needed to describe these nuclei. This algorithm greatly reduces the “fermion-sign problem” that quantum Monte Carlo methods are subject to, allowing us to obtain binding energies for a given Hamiltonian that are accurate to 1 to 2%.
The Hamiltonian we have used, consisting of the Argonne $`v_{18}`$ two-nucleon and Urbana IX three-nucleon potentials, gets the general features of the light p-shell nuclei fairly well, including the bulk of the experimentally observed binding and the correct ordering and approximate spacing of the excitation spectrum. However, with accurate calculations, we can identify specific failings of this Hamiltonian, including a few percent underbinding in $`N=Z`$ nuclei, which gets progressively worse as the neutron-proton asymmetry increases, and spin-orbit splittings in the excitation spectrum that are too small. Nevertheless, we can see that the energy differences with experiment are much smaller than the magnitude of the short-range (and least-well-known) part of our three-nucleon interaction, so it is plausible that new three-nucleon force models may be able to reproduce the ground and excited states near the 1% level.
We have also been able to study energy differences between isobaric analog states and the isospin-mixing matrix elements in <sup>8</sup>Be. While final conclusions should be reserved until the bulk energies are corrected by improved three-nucleon potentials, we can see that the charge-dependent and charge-symmetry-breaking components of the Argonne $`v_{18}`$ potential are making significant contributions that improve the agreement with experiment.
Finally, we have also studied the moments, densities, momentum distributions, and intrinsic shapes of these nuclei. It appears that the observed static moments may be understood with our present microscopic model after we include meson-exchange current contributions. We also see evidence for strong clustering in the light p-shell nuclei, particularly the 2$`\alpha `$ character in <sup>8</sup>Be. This clustering is built into the variational wave functions by the strong pair correlations, which depend upon the nucleon orbits and seem to be preserved in the GFMC propagation.
With the present method, and the continuing rapid increases in computational power that massively parallel machines are bringing, we are confident that we will be able to extend our calculations to the $`A=9,10`$ nuclei in the near future, and to <sup>12</sup>C in a few year’s time.
###### Acknowledgements.
The authors thank Dr. Dieter Kurath for many useful suggestions. The many-body calculations were made possible by generous grants of time on the IBM SP and SGI Origin 2000 of the Mathematics and Computer Science Division, Argonne National Laboratory. The GFMC calculations of <sup>8</sup>Li excited states were made with early-user time on the IBM SP at the National Energy Research Scientific Center. The work of SCP and RBW is supported by the U. S. Department of Energy, Nuclear Physics Division, under contract No. W-31-109-ENG-38, that of JC by the U.S. Department of Energy under contract W-7405-ENG-36, and that of VRP by the U.S. National Science Foundation via grant PHY98-00978.
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# 1 Introduction
## 1 Introduction
There are at least two excellent reasons to study the in-medium modifications of the pion-pion interaction in the scalar-isoscalar channel both being related to fundamental questions in present day nuclear physics. The first refers the binding energy of nuclear matter since a modification of the correlated two-pion exchange may have some deep consequences on the saturation mechanism. The second one is the direct connection with chiral symmetry restoration. From very general arguments based on QCD and model calculations, partial chiral symmetry restoration is expected to occur in nuclear matter. Hence, there must be a softening of a collective scalar-isoscalar mode, usually called the sigma meson, which becomes degenerate with its chiral partner i.e. the pion at full restoration density. This also implies that at some density the sigma-meson spectral function should exhibit a significant enhancement near the two-pion threshold. This effect can be seen as a precursor of chiral symmetry restoration associated with large fluctuations of the quark condensate near the chiral phase transition .
However, the first proposed medium effect was the modification of the two-pion propagator and the unitarized $`\pi \pi `$ interaction from the softening of the pion dispersion relation by p-wave coupling to $`ph`$ and $`\mathrm{\Delta }h`$ states. The existence of collective pionic modes produces a strong accumulation of strength near the two-pion threshold in the scalar-isoscalar channel . According to recent calculations , this reshaping of the strength may provide a partial explanation of the $`\pi 2\pi `$ data obtained on various nuclei by the CHAOS collaboration at TRIUMF . These results have been questioned in a recent paper where it is found that pion absorption forces the reaction to occur in the nuclear surface, i.e. at very low density . It is clear that the effect of chiral symmetry restoration has to be included on top of p-wave pionic effects to get a better explanation of the data. One attempt to combine both effects is based on the linear sigma model (implemented with a form factor fitted to phase shifts) in which the sigma mass is dropped through a Brown-Rho scaling relation . It is found that chiral symmetry restoration increases the strength of the threshold enhancement by about a factor four as reported in these proceedings . Here we will discuss another attempt based on the Nambu-Jona-Lasinio model . From this NJL model we derive an in-medium pion-pion potential, formally equivalent to the linear sigma model, but with parameters ($`m_\sigma ,m_\pi ,f_\pi `$) replaced by their in-medium values. Hence, at variance with a pure dropping sigma mass scenario, the basic $`\sigma \pi \pi `$ and $`4\pi `$ couplings are also modified in a chirally consistent framework. P-wave coupling of the pion is incorporated within a standard nuclear-matter approach since the NJL model completely misses (at least in its present treatment) the phenomenologically well-established strong screening effect from short-range correlations ($`g^{}`$ parameter). The underlying philosophy can be summarized in stating that the medium-modified soft physics linked to chiral symmetry ($`m_\pi ,f_\pi `$, low energy $`\pi \pi `$ potential ) is calculated within the NJL model while p-wave physics yielding pionic nuclear collective modes is described through standard nuclear phenomenology.
## 2 NJL and Density Dependent Linear Sigma Model
We start with the $`SU(2)`$ version of the NJL model : $`=\overline{\psi }(im)\psi +g[(\overline{\psi }\psi )^2+(\overline{\psi }i\gamma _5\tau _j\psi )^2]`$. At finite density, the gap equation which determines the constituent quark $`M`$ mass reads :
$`M`$ $`=`$ $`m+4N_cN_fg{\displaystyle _{k_F}^\mathrm{\Lambda }}{\displaystyle \frac{k^2dk}{2\pi ^2}}{\displaystyle \frac{M}{E}}`$ (1)
where $`\mathrm{\Lambda }`$ is the (three-momentum) cutoff. The pion and the sigma meson can be constructed within the standard RPA approximation. Limiting ourselves to zero momentum, the meson propagators are obtained according to :
$`g_{\pi qq}^2D_\pi (\omega )`$ $`=`$ $`\left(\omega ^2I(\omega )m_\pi ^2I(m_\pi )\right)^1`$
$`g_{\sigma qq}^2D_\sigma (\omega )`$ $`=`$ $`\left((\omega ^24M^2)I(\omega )m_\pi ^2I(m_\pi )\right)^1`$
$`\text{with}I(\omega )`$ $`=`$ $`2N_cN_f{\displaystyle _{k_F}^\mathrm{\Lambda }}{\displaystyle \frac{p^2dp}{2\pi ^2}}{\displaystyle \frac{1}{E_p(4E_p^2\omega ^2)}}`$ (2)
The in-medium masses of the pion and the sigma meson are found as the pole of the above propagators. It is important to notice that they are pure collective $`q\overline{q}`$ states since, at zero momentum, there is no contribution from particle-hole excitations (Fermi sea). According to what is said above, the particle-hole sector is better treated within a standard nuclear physics approach. In addition, for reason of simplicity, we use in the following a simplified scheme where $`I(\omega )`$ is “frozen” at $`\omega =0`$. We have verified numerically that this approximation gives almost the same result than the exact calculation for the pion mass, the sigma-meson mass and also for the pion decay constant (see ). The parameters of the model ($`m,\mathrm{\Lambda },g`$) have been adjusted to obtain the vacuum values $`f_\pi =93MeV`$, $`m_\pi =139MeV`$ and $`m_\sigma =1GeV`$ which is precisely the value of the bare sigma-meson mass, systematically used in our previous works . The s-wave optical potential is somewhat too small ($`4MeV`$ at $`\rho _0`$) but the incorporation of a phenonenologically more appropriate s-wave optical potential ($`10MeV`$) has only minor influences on the medium effects presented in this paper.
We are now in position to construct vacuum and in-medium $`\pi \pi `$ potentials following the method developed in but adapted to finite density. The basic diagrams are box diagrams with four internal quark lines for the direct $`4\pi `$ interaction and three internal quark lines for the $`\pi \pi \sigma `$ coupling from which one obtains the sigma exchange diagram (Fig. 1). Keeping the full momentum dependence makes the problem of subsequent unitarization hopelessly complex. As in , we first limit the calculation to the case where the four pions have momenta $`p_i`$ such that $`p_i^2=m_\pi ^2`$ (what we call the exact scheme in ). However to make the unitarization tractable we further simplify the calculation by freezing the momentum to $`p=0`$ in the quark-loop integral (simplified scheme). We have checked numerically that the resulting low-energy $`\pi \pi `$ amplitudes are almost identical in the simplified and exact schemes . We obtain for the coupling constants $`\lambda _{4\pi }\lambda =(m_\sigma ^2m_\pi ^2)/2f_\pi ^2`$ and $`\lambda _{\sigma \pi \pi }=\lambda f_\pi `$. This is the well-known result of the linear sigma model but with medium-modified parameters ($`m_\sigma ,m_\pi ,f_\pi `$). The linear sigma model, seen as a $`O(N+1)`$ model with $`N=3`$, is treated to leading order in the $`1/N`$ expansion which fulfills all the constraints (Ward identities) of chiral symmetry . We supplement the model by adding a one-parameter form factor $`v(k)`$ to fit the phase shifts in vacuum once the scattering amplitude is unitarized. The (in-medium) unitarized scalar-isoscalar $`\pi \pi `$ T matrix (in the CM frame and for total energy $`E`$ of the pion pair) is :
$$𝐤,𝐤|T(E)|𝐤^{},𝐤^{}=v(k)v(k^{})\frac{6\lambda (E^2m_\pi ^2)}{13\lambda \mathrm{\Sigma }(E)}D_\sigma (E)$$
(3)
where the unitarized sigma propagator (i.e with two-pion loop) $`D_\sigma (E)`$ is :
$$D_\sigma (E)=\left(E^2m_\sigma ^2\frac{6\lambda ^2f_\pi ^2\mathrm{\Sigma }(E)}{13\lambda \mathrm{\Sigma }(E)}\right)^1$$
(4)
Chiral symmetry restoration is accounted for by the in-medium pion- as well as sigma-meson masses and $`f_\pi `$. The p-wave collective effects are embedded in the two-pion loop :
$$\mathrm{\Sigma }(E)=\frac{d𝐪}{(2\pi )^3}v(q)\frac{idq_0}{2\pi }D_\pi (𝐪,q_0)D_\pi (𝐪,Eq_0).$$
(5)
The pion propagator $`D_\pi (𝐪,q_0)`$ is calculated in a standard nuclear matter approach and incorporates the p-wave coupling of the pion to delta-hole states with short-range screening described by the usual $`g_{\mathrm{\Delta }\mathrm{\Delta }}^{}=0.5`$ parameter. It is particularly interesting to inspect the sigma-meson spectral function. The result of the calculation is shown on Fig. 2. At twice normal nuclear-matter, which is close to the density where the quasi-pole in $`D_\sigma `$ is $`E=2m_\pi `$, we find a very sharp peak which can be understood as a precursor effect of chiral symmetry restoration. At normal nuclear-matter density the threshold peak is enhanced by a factor three as compared to a pure p-wave calculation. This is in qualitative agreement with the dropping sigma-meson mass calculation where a factor four is obtained . At density $`0.5\rho _0`$, more relevant for the CHAOS experiment, we still have a sizable low-energy reshaping which could help to explain the data.
## 3 Discussion and Interpretation
Four-quark condensates are fundamental quantities of non-pertubative QCD. They are known to play an important role in QCD sum rule analyses of hadron spectral functions in the vacuum and in matter. As an example the density evolution of the four-quark condensate appearing in the rho-meson case is an important input for the medium modification of this vector meson in connection with dilepton production in relativistic heavy ion collisions. Here we will concentrate on to the scalar four-quark condensate $`(\overline{q}q)^2`$ and study how much it deviates from $`\overline{q}q^2`$ to asses the evolution of quantum fluctuations with increasing density. Inserting a complete set of states, the scalar four-quark condensate is given as :
$$0|(\overline{q}q)^2|0=0|\overline{q}q|00|\overline{q}q|0+\underset{n}{}0|\overline{q}q|nn|\overline{q}q|0$$
(6)
The first contribution to the sum are scalar-isoscalar two-pion states, able to build a collective sigma meson. To evaluate (at least qualitatively) this quantity one can use low-energy effective theories.
Linear sigma model. In the linear sigma model with chiral-symmetry breaking piece $`_{\chi SB}=f_\pi m_\pi ^2\sigma `$, $`\sigma `$ plays the role of the condensate and $`\sigma ^2`$ relates to the four-quark condensate. Here we do not aim to estimate the absolute value of this four-quark condensate but restrict our study to its evolution with density. We thus define the quantity $`Q_1=0|(\overline{q}q)^2|0(\rho )/\overline{q}q_{vac}^2`$. Introducing the fluctuating part of the sigma field $`s=\sigma f_\pi `$, one obtains :
$$Q_1(\rho )=\frac{(\overline{q}q)^2(\rho )}{\overline{q}q_{vac}^2}=\frac{\sigma ^2(\rho )}{f_\pi ^2}=1+\mathrm{\hspace{0.17em}2}\frac{s(\rho )}{f_\pi }+\frac{s^2(\rho )}{f_\pi ^2}$$
(7)
A better measure of the fluctuations of the condensate is provided by the ’kappa factor’ defined in QCD sum rule analyses :
$$\kappa _1(\rho )=\frac{(\overline{q}q)^2(\rho )}{\overline{q}q^2(\rho )}=\frac{\sigma ^2(\rho )}{\sigma ^2(\rho )}=\frac{1+\mathrm{\hspace{0.17em}2}s(\rho )/f_\pi +s^2(\rho )/f_\pi ^2}{1+\mathrm{\hspace{0.17em}2}s(\rho )/f_\pi +s^2(\rho )/f_\pi ^2}$$
(8)
To first order in the density $`s(\rho )`$ is governed by the pion-nucleon sigma term. Using a value compatible with the model (ignoring the pionic contribution), we take $`s(\rho )=0.18\rho /\rho _0`$. From a standard dispersive analysis, $`s^2`$ can be expressed in terms of a phase-space integral of the sigma-meson spectral function as :
$$s^2(\rho )=_0^{\mathrm{\Lambda }_P}\frac{d𝐏}{(2\pi )^3}_0^{\mathrm{}}𝑑E\left(\frac{1}{\pi }\right)ImD_\sigma (E,𝐏)$$
(9)
where we have introduced a momentum cutoff $`\mathrm{\Lambda }_P`$ defining the range of validity of the effective approach. Taking $`\mathrm{\Lambda }_P1GeV`$, and making a simple estimate with a sharp sigma meson of mass 1 GeV, we obtain in the vacuum a kappa factor of the order or larger than 2 which is very close (most probably by accident) to $`\kappa =2.36`$, generally used in the rho-meson channel. Finally, since we do not know the full momentum dependence of the sigma spectral function, we assume covariance for $`D_\sigma `$. We have checked that, using another extreme assumption (static approximation), the results are qualitatively similar. To reduce the uncertainty on the cutoff we prefer to present the results in Tabl. 1 for the quantities $`\mathrm{\Delta }Q_1(\rho )=Q_1(\rho )(Q_1)_{vac}`$ and $`\mathrm{\Delta }\kappa _1(\rho )=\kappa _1(\rho )(\kappa _1)_{vac}`$. For the sigma-meson propagator we use the form (4) but with vacuum values for $`m_\sigma `$, $`m_\pi `$ and $`f_\pi `$, hence keeping only medium effects from p-wave pionic collective modes. In the actual calculation we have adopted the cutoff parameter $`\mathrm{\Lambda }_P=1.2GeV`$ to (arbitrarily) fix the kappa factor in the vacuum to $`\kappa _{vac}=2.36`$ keeping in mind that the density evolution should not be very sensitive to this particular choice.
NJL model. In the NJL model the four-quark condensate can be directly calculated. It can be expressed in terms of a phase-space integral of the full scalar-isoscalar response function and subsequently in terms of the $`q\overline{q}`$ sigma meson spectral function :
$$\kappa _2(\rho )=\frac{(\overline{q}q)^2(\rho )}{\overline{q}q^2(\rho )}=1+\frac{1}{f_\pi ^2(\rho )}_0^{\mathrm{\Lambda }_P}\frac{d𝐏}{(2\pi )^3}_0^{\mathrm{}}𝑑E\left(\frac{1}{\pi }\right)ImD_\sigma (E,𝐏)$$
(10)
$$Q_2(\rho )=\frac{(\overline{q}q)^2(\rho )}{\overline{q}q)_{vac}^2}=\frac{\left(f_\pi ^2m_\pi ^2\right)(\rho )}{\left(f_\pi ^2m_\pi ^2\right)_{vac}}\kappa _2(\rho )$$
(11)
The sigma propagator is then unitarized by incorporating a dressed pion loop (eq. 4). Notice that, in the vacuum, we exactly recover the linear sigma model results. We show in Tab. 1 the quantities $`\mathrm{\Delta }Q_2(\rho )`$ and $`\mathrm{\Delta }\kappa _2(\rho )`$ which now contain the effect of chiral symmetry restoration on top of the p-wave pionic collective modes. One sees that the four-quark condensate $`Q_2`$ increases with density when chiral symmetry restoration is incorporated, contrary to the pure p-wave case ($`Q_1`$). We also see, by looking at the kappa factor, that chiral symmetry restoration considerably increases the fluctuations of the condensate. This demonstrates that the sharp structure near $`2m_\pi `$, which is mainly associated with the dropping sigma-meson mass, is intimately related to precritical effects with a strong enhancement of chiral fluctuations. In a second-order phase transition these would actually diverge in the chiral limit, $`m_\pi 0`$.
Acknowledgements:
We wish to thank P. Schuck for many fruitful discussions and constant interest in this work.
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# On 𝛼-Critical Edges in König-Egerváry Graphs
## 1 Introduction
Throughout this paper $`G=(V,E)`$ is a simple (i.e., a finite, undirected, loopless and without multiple edges) graph with vertex set $`V=V(G)`$, edge set $`E=E(G)`$, and order $`n(G)=|V(G)|`$. If $`XV`$, then $`G[X]`$ is the subgraph of $`G`$ spanned by $`X`$. By $`GW`$ we mean the subgraph $`G[VW]`$ , if $`WV(G)`$. For $`FE(G)`$, by $`GF`$ we denote the partial subgraph of $`G`$ obtained by deleting the edges of $`F`$, and we use $`Ge`$, if $`W`$ $`=\{e\}`$. If $`A,B`$ $`V`$ and $`AB=\mathrm{}`$, then $`(A,B)`$ stands for the set $`\{e=ab:aA,bB,eE\}`$. The neighborhood of a vertex $`vV`$ is the set $`N(v)=\{w:wV`$ and $`vwE\}`$, and $`N(A)=\{N(v):vA\}`$, $`N[A]=AN(A)`$ for $`AV`$.
A set $`S`$ of vertices is stable if no two vertices from $`S`$ are adjacent. A stable set of maximum size will be referred to as a maximum stable set of $`G`$. The stability number of $`G`$, denoted by $`\alpha (G)`$, is the cardinality of a maximum stable set of $`G`$. Let $`\mathrm{\Omega }(G)`$ denotes the set $`\{S:S`$ is a maximum stable set of $`G\}`$, $`\sigma (G)=\left|\{VS:S\mathrm{\Omega }(G)\}\right|`$ and $`\xi (G)=\left|core(G)\right|`$, where $`core(G)=\{S:S\mathrm{\Omega }(G)\}`$, . In other words, $`\xi (G)`$ equals the number of $`\alpha `$-critical vertices of $`G`$, (a vertex $`vV(G)`$ is $`\alpha `$-critical provided $`\alpha (Gv)<\alpha (G)`$).
By $`P_n,C_n,K_n`$ we mean the chordless path on $`n3`$, the chordless cycle on $`n`$ $`4`$ vertices, and respectively the complete graph on $`n1`$ vertices.
A matching (i.e., a set of non-incident edges of $`G`$) of maximum cardinality $`\mu (G)`$ is a maximum matching, and a perfect matching is one covering all vertices of $`G`$. An edge $`eE(G)`$ is $`\mu `$-critical provided $`\mu (Ge)<\mu (G)`$. By their definition, $`\mu `$-critical edges of $`G`$ belong to all maximum matchings of $`G`$.
If $`\alpha (G)+\mu (G)=n(G)`$, then $`G`$ is called a König-Egerváry graph, , . Properties of these graphs were presented in several papers, like of Sterboul , Deming , Lovász and Plummer , Korach , Bourjolly and Pulleyblank , Paschos and Demange , Levit and Mandrescu , . It is worth observing that a disconnected graph is of König-Egerváry type if and only if all its connected components are König-Egerváry graphs. In this paper, by ”graph” we mean a connected graph having at least one edge.
An edge $`eE(G)`$ is $`\alpha `$-critical whenever $`\alpha (Ge)>\alpha (G)`$. Let denote by $`\eta (G)`$ the number of $`\alpha `$-critical edges of $`G`$. Notice that there are graphs in which: ($`a`$) any edge is $`\alpha `$-critical (so-called $`\alpha `$-critical graphs); e.g., all $`C_{2n+1}`$ for $`n3`$; ($`b`$) no edge is $`\alpha `$-critical; e.g., all $`C_{2n}`$ for $`n2`$. More generally, Haynes et al., , have proved that a graph $`G`$ has no $`\alpha `$-critical edge if and only if $`\left|N(x)S\right|2`$ holds for any $`S\mathrm{\Omega }(G)`$ and every $`xV(G)S`$.
Beineke, Harary and Plummer, , have shown that any two incident $`\alpha `$-critical edges of a graph lie on an odd cycle, and hence, they deduce that no two $`\alpha `$-critical edges of a bipartite graph can have a common endpoint. Independently, Zito, , has proved the same result for trees using a different technique. Some variations and strengthenings of these results are discussed in , , and .
In this paper we generalize the above assertion to König-Egerváry graphs. We also show that $`\alpha `$-critical edges are $`\mu `$-critical in a König-Egerváry graph, and that they coincide in bipartite graphs. As a corollary, we obtain one result of Zito, , stating that a vertex $`v`$ is in some but not in all maximum stable sets of a tree $`T`$ if and only if $`v`$ is an endpoint of an $`\alpha `$-critical edge of $`T`$. In the sequel, we analyze other relationships between $`\alpha `$-critical edges and $`\mu `$-critical edges in a König-Egerváry graph, and its corresponding implications to equalities and inequalities linking $`\alpha (G)`$, $`\xi (G)`$, $`\eta (G)`$, $`\sigma (G)`$, and $`\mu (G)`$. Eventually, we infer that $`\alpha (T)=\xi (T)+\eta (T),\sigma (T)+\eta (T)=\mu (T)`$ and $`\xi (T)+2\eta (T)+\sigma (T)=n(T)`$ holds for any tree $`T`$, and characterize the König-Egerváry graphs having these properties.
## 2 $`\alpha `$-Critical and $`\mu `$-Critical Edges
According to a well-known result of König, , and Egerváry, , any bipartite graph is a König-Egerváry graph. It is easy to see that this class includes also some non-bipartite graphs (see, for instance, the graph $`K_3+e`$ in Figure 1).
If $`G_i=(V_i,E_i),i=1,2`$, are two disjoint graphs, then $`G=G_1G_2`$ is defined as the graph with $`V(G)=V(G_1)V(G_2)`$, and
$$E(G)=E(G_1)E(G_2)\{xy:forsomexV(G_1)andyV(G_2)\}.$$
Clearly, if $`H_1,H_2`$ are subgraphs of a graph $`G`$ such that $`V(G)=V(H_1)V(H_2)`$ and $`V(H_1)V(H_2)=`$ $`\mathrm{}`$, then $`G=H_1H_2`$, i.e., any graph of order at least two admits such decompositions. However, some particular cases are of special interest. For instance, if: $`E(H_i)=\mathrm{},i=1,2,`$ then $`G=H_1H_2`$ is bipartite; $`E(H_1)=\mathrm{}`$ and $`H_2`$ is complete, then $`G=H_1H_2`$ is a split graph .
The following result shows that the König-Egerváry graphs are, in this sense, between these two ”extreme” situations. The equivalence of the first and the third parts of this proposition was proposed by Klee and included in without proof (private communication).
###### Proposition 2.1
The following assertions are equivalent:
($`i`$) $`G`$ is a König-Egerváry graph;
($`\mathrm{𝑖𝑖}`$) $`G=H_1H_2`$, where $`V(H_1)=S\mathrm{\Omega }(G)`$ and $`n(H_1)\mu (G)=n(H_2)`$;
($`\mathrm{𝑖𝑖𝑖}`$) $`G=H_1H_2`$, where $`V(H_1)=S`$ is a stable set in $`G,`$ $`\left|S\right|n(H_2)`$ and $`(S,V(H_2))`$ contains a matching $`M`$ with $`\left|M\right|=n(H_2)`$.
In the sequel, we shall often represent a König-Egerváry graph $`G`$ as $`G=SH`$, where $`S\mathrm{\Omega }(G)`$ and $`H=G[VS]`$ has $`n(H)=\mu (G)`$.
###### Lemma 2.2
If $`G=(V,E)`$ is a König-Egerváry graph, then any maximum matching of $`G`$ is contained in $`(S,VS)`$, where $`S\mathrm{\Omega }(G)`$.
Clearly, Lemma 2.2 is not valid for any graph. For instance, $`K_4`$ is a counterexample. Moreover, $`K_4`$ has $`\alpha `$-critical edges that are incident. Nevertheless, there are graphs having only non-incident $`\alpha `$-critical edges.
###### Theorem 2.3
If $`G`$ is a König-Egerváry graph, then the following assertions hold:
($`i`$) for any $`\alpha `$-critical edge $`e`$ of $`G`$, the graph $`Ge`$ is still a König-Egerváry graph;
($`\mathrm{𝑖𝑖}`$) any $`\alpha `$-critical edge of $`G`$ is also $`\mu `$-critical;
($`\mathrm{𝑖𝑖𝑖}`$) the $`\alpha `$-critical edges of $`G`$ form a matching.
Proof. ($`i`$) If $`e=xy`$ is an $`\alpha `$-critical edge $`G`$, then there is some $`S\mathrm{\Omega }(G)`$ such that either $`N(x)S=\{y\}`$ or $`N(y)S=\{x\}`$. Suppose that $`yS`$. Since $`S\mathrm{\Omega }(G)`$, we get, by Proposition 2.1, that $`G=SH`$, where $`H=G[VS]`$ has $`\mu (G)=n(H)=\left|M\right|`$ and $`M`$ is a maximum matching of $`G`$, included, by Lemma 2.2, in $`(S,V(G)S)`$. Hence, it follows that $`Ge=S^{}V(H^{})`$, where $`S^{}=S\{x\}\mathrm{\Omega }(Ge)`$ and $`n(H^{})=\left|M\{e\}\right|`$. According to Proposition 2.1($`\mathrm{𝑖𝑖𝑖}`$), we infer that $`Ge`$ is also a König-Egerváry graph.
($`\mathrm{𝑖𝑖}`$) If $`eE(G)`$ is an $`\alpha `$-critical edge of $`G`$, then according to ($`i`$) we obtain:
$$n(G)=\alpha (G)+\mu (G)\alpha (Ge)+\mu (Ge)=\alpha (G)+1+\mu (Ge)=n(Ge),$$
and this implies $`\mu (G)=1+\mu (Ge)`$, i.e., $`e`$ is also $`\mu `$-critical.
($`\mathrm{𝑖𝑖𝑖}`$) Let $`e_1,e_2`$ be two $`\alpha `$-critical edges of $`G`$. We have to show that they are not incident. According to second part ($`\mathrm{𝑖𝑖}`$), both edges are also $`\mu `$-critical. Hence, it follows that $`e_1,e_2\{M:MisamaximummatchingofG\}`$ and this ensures that $`e_1,e_2`$ have no common endpoint. Consequently, the set of all $`\alpha `$-critical edges of $`G`$ yields a matching.
Notice that:
($`a`$) Theorem 2.3($`i`$) is not true for any $`\mu `$-critical edge of a König-Egerváry graph; e.g., the edge $`e`$ of $`G=K_3+e`$ is $`\mu `$-critical, but $`Ge`$ is not a König-Egerváry graph;
($`b`$) Theorem 2.3($`\mathrm{𝑖𝑖}`$) is not true for any graph; e.g., all the edges of $`K_3`$ are $`\alpha `$-critical, but none is also $`\mu `$-critical;
($`c`$) the converse of Theorem 2.3($`\mathrm{𝑖𝑖}`$) is not valid for any König-Egerváry graph; e.g., the edge $`e`$ of graph $`K_3+e`$ is $`\mu `$-critical, but is not also $`\alpha `$-critical. However, as we shall see later, (namely Proposition 2.6), the $`\mu `$-critical edges are also $`\alpha `$-critical in the case of bipartite graphs.
###### Corollary 2.4
A König-Egerváry graph is $`\alpha `$-critical if and only if it is isomorphic to $`K_2`$.
Since any bipartite graph is also a König-Egerváry graph, we obtain the following statement, due to Beineke, Harary and Plummer.
###### Theorem 2.5
No two $`\alpha `$-critical edges of a bipartite graph are incident.
###### Proposition 2.6
If $`G`$ is a bipartite graph, then its $`\alpha `$-critical edges coincide with its $`\mu `$-critical edges.
Proof. By Theorem 2.3($`\mathrm{𝑖𝑖}`$), it suffices to show that any $`\mu `$-critical edge $`e`$ of $`G`$ is also $`\alpha `$-critical. Since $`Ge`$ is still bipartite, and hence, also a König-Egerváry graph, it follows that $`\alpha (Ge)+\mu (Ge)=n(G)=\alpha (G)+\mu (G)=\alpha (G)+1+\mu (Ge)`$, and this implies $`\alpha (Ge)>\alpha (G)`$, i.e., $`e`$ is an $`\alpha `$-critical edge of $`G`$.
In Theorem 4.2 we will meet another type of König-Egerváry graphs with this property. Notice that there are also non-bipartite König-Egerváry graphs in which their $`\mu `$-critical edges are $`\alpha `$-critical (see the graph in Figure 2).
It is well-known that if a tree has a perfect matching, then it is unique. Consequently, we obtain:
###### Corollary 2.7
A tree has a perfect matching if and only if the set of its $`\alpha `$-critical edges forms a maximal matching of the tree.
Using the definition of König-Egerváry graphs and the fact that $`\mu (G)n(G)/2`$ is true for any graph $`G`$, we get:
###### Lemma 2.8
If $`G`$ admits a perfect matching, then $`G`$ is a König-Egerváry graph if and only if $`\alpha (G)=\mu (G)`$. If $`G`$ is a König-Egerváry graph, then $`\mu (G)\alpha (G)`$.
Combining Corollary 2.7 and Lemma 2.8, we get the following result from .
###### Corollary 2.9
If a tree $`T`$ has a perfect matching $`M`$, then all the edges of $`M`$ are $`\alpha `$-critical and $`2\alpha (T)=n(T)`$.
###### Proposition 2.10
If $`G=(V,E)`$ is a König-Egerváry graph, then the following assertions are true:
($`i`$) any $`S\mathrm{\Omega }(G)`$ meets each $`\mu `$-critical edge in exactly one vertex;
($`\mathrm{𝑖𝑖}`$) any $`S\mathrm{\Omega }(G)`$ meets each $`\alpha `$-critical edge in exactly one vertex;
($`\mathrm{𝑖𝑖𝑖}`$) if $`G`$ has a maximal matching consisting of only $`\alpha `$-critical edges, then it is the unique perfect matching of $`G`$.
Proof. ($`i`$) $`\mathrm{𝑎𝑛𝑑}`$ ($`\mathrm{𝑖𝑖}`$) By Theorem 2.3($`\mathrm{𝑖𝑖}`$), any $`\alpha `$-critical edge of $`G`$ is also $`\mu `$-critical. Consequently, we infer that
$$\{eE:eis\alpha critical\}\{M:MisamaximummatchingofG\}(S,VS)$$
holds for any $`S\mathrm{\Omega }(G)`$, according to Lemma 2.2. It follows that if $`e=xy`$ is an $`\alpha `$-critical or a $`\mu `$-critical edge of $`G`$, then any $`S\mathrm{\Omega }(G)`$ contains one of $`x`$ and $`y`$, (since clearly, no stable set may contain both $`x`$ and $`y`$).
($`\mathrm{𝑖𝑖𝑖}`$) Let $`M`$ be a maximal matching of $`G`$ consisting of only $`\alpha `$-critical edges. By Theorem 2.3, all the edges of $`M`$ are also $`\mu `$-critical. Therefore, we infer that $`M`$ is included in any maximum matching of $`G`$, and because $`M`$ is a maximal matching, it results that $`M`$ is the unique maximum matching of $`G`$. Suppose, on the contrary, that $`M`$ is not perfect, and let $`S\mathrm{\Omega }(G)`$. According to Proposition 2.1, $`G`$ can be written as $`G=SH`$, with $`n(H)=\left|M\right|=\mu (G)`$, and by Lemma 2.2 we have that $`M(S,VS)`$. Since $`G`$ is a König-Egerváry graph without perfect matchings, Lemma 2.8 implies $`\left|S\right|=\alpha (G)>\mu (G)=\left|M\right|`$. Hence, it follows that there are at least two vertices $`v_1,v_2S`$ having a common neighbor $`wV(H)`$ and such that one of them, say $`v_1`$, is unmatched by $`M`$ and $`v_2wM`$. Thus, $`M\{v_1w\}\{v_2w\}`$ is another maximum matching of $`G`$, in contradiction with the uniqueness of $`M`$. Consequently, $`M`$ must be also perfect.
For trees, Proposition 2.10($`\mathrm{𝑖𝑖}`$) was proved by Zito in .
Notice that the matching in Proposition 2.10($`\mathrm{𝑖𝑖𝑖}`$) is not necessarily formed by pendant edges; e.g., $`P_6`$ has such a matching. Concerning the uniqueness of this matching, it is worth mentioning that: ($`a`$) if $`G`$ is not a König-Egerváry graph, then it may have several different maximum matchings consisting of only $`\alpha `$-critical edges (e.g., $`C_5`$); ($`b`$) if a König-Egerváry graph has a unique perfect matching, then it may contain non-$`\alpha `$-critical edges (e.g., the edge $`e`$ of $`K_3+e`$ is not $`\alpha `$-critical, but it belongs to the unique perfect matching of $`K_3+e`$).
## 3 Equalities and Inequalities between Parameters
If $`vN(core(G))`$, then clearly follows that $`vV(G)S`$, for any $`S\mathrm{\Omega }(G)`$, that is $`N(core(G))\{VS:S\mathrm{\Omega }(G)\}`$ holds for any graph $`G`$.
###### Lemma 3.1
If $`G=(V,E)`$ is a König-Egerváry graph, then
$`N(core(G))=\{VS:S\mathrm{\Omega }(G)\}.`$
Notice that there are graphs that do not enjoy the above equality, for example, the graph $`G`$ in Figure 3($`a`$) has $`N(core(G))=\mathrm{}`$ and $`\{VS:S\mathrm{\Omega }(G)\}=\{v\}`$. There exist non-König-Egerváry graphs for which $`N(core(G))=\{VS:S\mathrm{\Omega }(G)\}`$, (see, for instance, the graph $`G`$ from Figure 3($`b`$)).
###### Proposition 3.2
If $`G=(V,E)`$ is a König-Egerváry graph, $`G_0=GN[core(G)]`$ and $`S\mathrm{\Omega }(G)`$, then the following assertions are true:
($`i`$) $`\left|core(G)\right|\left|N(core(G))\right|`$;
($`\mathrm{𝑖𝑖}`$) $`\left|Score(G)\right|=\left|VSN(core(G))\right|`$;
($`\mathrm{𝑖𝑖𝑖}`$) $`G_0`$ has a perfect matching and it is also a König-Egerváry graph.
Proof. According to Proposition 2.1, $`G`$ can be written as $`G=SH`$, where $`H=G[VS]`$ has $`n(H)=\mu (G)`$. Let denote $`A=Score(G)`$ and $`B=V(H)N(core(G))`$. In it has been proved that $`\left|A\right|\left|B\right|`$ holds for any graph $`G`$. Since $`\{VS:S\mathrm{\Omega }(G)\}V(H)`$, and $`N(core(G))=\{VS:S\mathrm{\Omega }(G)\}`$ (see Lemma 3.1), we obtain $`B=V(H)\{VS:S\mathrm{\Omega }(G)\}`$.
($`i`$) Since $`\left|A\right|+\left|core(G)\right|=\alpha (G)\mu (G)=n(H)=\left|B\right|+|N(core(G)|`$ and, on the other hand $`\left|A\right|\left|B\right|`$, it follows that $`\left|core(G)\right|\left|N(core(G))\right|`$.
($`\mathrm{𝑖𝑖}`$) Let $`M`$ be a maximum matching in $`G`$. Since $`G`$ is a König-Egerváry graph, Lemma 2.2 ensures that $`M`$ is included in $`(S,V(H))`$, and $`\left|M\right|=\mu (G)=n(H)`$. The matching $`M`$ matches $`B`$ into $`A`$, because there are no edges connecting $`B`$ and $`core(G)`$. Hence, $`\left|B\right|\left|A\right|`$. Together with $`\left|A\right|\left|B\right|`$ $`\left|A\right|\left|B\right|`$, it implies $`\left|A\right|=\left|B\right|`$, i.e., $`\left|Score(G)\right|=\left|VSN(core(G))\right|`$, and that $`M(A,B)`$ is a perfect matching of $`G[AB]`$.
($`\mathrm{𝑖𝑖𝑖}`$) Since, in fact, $`G_0=G[AB]`$, it follows necessarily that $`G_0`$ has a perfect matching. In addition, because $`A`$ is stable, we get $`\alpha (G_0)\mu (G_0)=\left|A\right|\alpha (G_0)`$, i.e., $`\alpha (G_0)=\mu (G_0)`$, and according to Lemma 2.8, $`G_0`$ must be also a König-Egerváry graph.
###### Corollary 3.3
If $`G`$ is a König-Egerváry graph, then $`\alpha (G)+\sigma (G)=\mu (G)+\xi (G)`$.
Proof. By Lemma 3.1, $`N(core(G))=\{VS:S\mathrm{\Omega }(G)\}`$ and according to Proposition 3.2($`\mathrm{𝑖𝑖}`$), $`\left|Score(G)\right|=\left|VSN(core(G))\right|`$. Hence, we obtain that $`\alpha (G)\xi (G)=\left|Score(G)\right|=\left|VSN(core(G))\right|=\mu (G)\sigma (G)`$.
Let us observe that there exist non-König-Egerváry graphs satisfying the equality $`\alpha (G)+\sigma (G)=\mu (G)+\xi (G)`$ (see graph $`W_1`$ in Figure 6). It is also interesting to notice that there exists a non-König-Egerváry graph enjoying the property that its subgraph $`G_0=GN[core(G)]`$ has a perfect matching (see Figure 8). Figure 4 shows a non-König-Egerváry graph $`G`$ whose $`G_0`$ has no perfect matching.
###### Lemma 3.4
Let $`G=(V,E)`$ and $`G_0=GN[core(G)]`$. Then the following assertions are valid:
($`i`$) no $`\alpha `$-critical edge in $`G`$ has an endpoint in $`N[core(G)]`$;
($`\mathrm{𝑖𝑖}`$) $`\alpha (G)=\alpha (G_0)+\xi (G),\mathrm{\Omega }(G_0)=\{SV(G_0):S\mathrm{\Omega }(G)\},core(G_0)=\mathrm{}`$;
($`\mathrm{𝑖𝑖𝑖}`$) $`e=xy`$ is an $`\alpha `$-critical edge of $`G`$ if and only if $`e`$ is an $`\alpha `$-critical edge of $`G_0`$.
Proof. ($`i`$) Let $`e=xy`$ be an $`\alpha `$-critical edge in $`G`$, and let $`\overline{S}\mathrm{\Omega }(Ge)`$. Since $`\left|\overline{S}\right|=\alpha (Ge)>\alpha (G)`$, it follows that $`x,y\overline{S}`$ and $`\overline{S}\{x\},\overline{S}\{y\}\mathrm{\Omega }(G)`$. Now, the inclusion $`N(core(G))\{VS:S\mathrm{\Omega }(G)\}`$ completes the proof that no $`\alpha `$-critical edge in $`G`$ has an endpoint in $`N(core(G))`$, and respectively, in $`core(G)`$.
($`\mathrm{𝑖𝑖}`$) By definition of $`G_0`$, if $`S\mathrm{\Omega }(G)`$, then $`Score(G)=SV(G_0)`$, and therefore
$$\alpha (G)\xi (G)=\left|Score(G)\right|\alpha (G_0).$$
For any $`S_{G_0}\mathrm{\Omega }(G_0)`$ we have that $`S_{G_0}core(G)`$ is stable, and hence
$$\left|S_{G_0}core(G)\right|=\alpha (G_0)+\xi (G)\alpha (G).$$
Consequently, we get $`\alpha (G)=\alpha (G_0)+\xi (G)`$. Now it is easy to check that $`\mathrm{\Omega }(G_0)=\{SV(G_0):S\mathrm{\Omega }(G)\}`$ and $`core(G_0)=\mathrm{}`$.
($`\mathrm{𝑖𝑖𝑖}`$) Let $`e=xy`$ be an $`\alpha `$-critical edge of $`G`$. By ($`i`$), we infer that $`eE(G_0)`$, and as we saw above, there is some stable set $`S_{xy}`$ such that $`S_{xy}\{x\},S_{xy}\{y\}\mathrm{\Omega }(G)`$ and $`S_{xy}\{x,y\}\mathrm{\Omega }(Ge)`$. Hence, ($`\mathrm{𝑖𝑖}`$) implies that
$$V(G_0)(S_{xy}\{x\}),V(G_0)(S_{xy}\{y\})\mathrm{\Omega }(G_0)andV(G_0)(S_{xy}\{x,y\})\mathrm{\Omega }(G_0e),$$
because $`V(G_0)(S_{xy}\{x,y\})`$ is stable in $`G_0e`$ and larger than $`V(G_0)(S_{xy}\{x\})`$. Therefore, $`e`$ is $`\alpha `$-critical in $`G_0`$, as well. Similarly, we can show that any $`\alpha `$-critical edge of $`G_0`$ is $`\alpha `$-critical in $`G`$ too.
###### Proposition 3.5
If $`G`$ is a König-Egerváry graph, then
($`i`$) $`\xi (G)+\eta (G)\alpha (G)`$;
($`\mathrm{𝑖𝑖}`$) $`\sigma (G)+\eta (G)\mu (G)`$;
($`\mathrm{𝑖𝑖𝑖}`$) $`\xi (G)+2\eta (G)+\sigma (G)n(G)`$.
Proof. For any $`S\mathrm{\Omega }(G)`$, we have that $`core(G)S`$, and by Lemma 3.4($`i`$), no $`\alpha `$-critical edge has an endpoint in $`core(G)`$. In addition, according to Proposition 2.10($`\mathrm{𝑖𝑖}`$), $`S`$ meets each $`\alpha `$-critical edge in exactly one vertex. Hence, it follows that $`\xi (G)+\eta (G)\alpha (G)`$, and using Corollary 3.3 we obtain ($`\mathrm{𝑖𝑖}`$). Clearly, ($`\mathrm{𝑖𝑖𝑖}`$) follows from ($`i`$) and ($`\mathrm{𝑖𝑖}`$).
Notice that $`\xi (K_3+e)+\eta (K_3+e)=\alpha (K_3+e)`$ and also $`\eta (K_3+e)+\sigma (K_3+e)=\mu (K_3+e)`$, but there are König-Egerváry graphs satisfying $`\xi (G)+\eta (G)<\alpha (G)`$ and $`\eta (G)+\sigma (G)<\mu (G)`$. For instance, $`G=C_6`$, and also the graph $`W`$ in Figure 5 is a König-Egerváry non-bipartite graph that has $`\eta (W)=\left|\{e\}\right|=1,\xi (W)=\left|\{a\}\right|=1=\sigma (W),\alpha (W)=\mu (W)=4`$.
Observe that Proposition 3.5 is not true for general graphs; e.g., the graph $`W_1`$ in Figure 6 has $`\alpha (W_1)=3,\mu (W_1)=2,\eta (W_1)=3,\xi (W_1)=2,\sigma (W_1)=1`$. However, there are non-König-Egerváry graphs satisfying $`\xi (G)+\eta (G)<\alpha (G)`$ and $`\eta (G)+\sigma (G)<\mu (G)`$, for example, the graph $`W_2`$ in Figure 6 has $`\alpha (W_2)=3,\eta (W_2)=\left|\{ab,cd\}\right|,\xi (W_2)=\sigma (W_2)=0`$. There also exist non-König-Egerváry graphs satisfying $`\xi (G)+\eta (G)=\alpha (G)`$ and $`\eta (G)+\sigma (G)=\mu (G)`$, e.g., the graph $`W_3`$ in Figure 6. Nevertheless, $`\xi (K_5e)+\eta (K_5e)=\alpha (K_5e)`$, but $`\eta (K_5e)+\sigma (K_5e)>\mu (K_5e)`$.
###### Proposition 3.6
If $`G`$ is a König-Egerváry graph, then the following assertions are equivalent:
($`i`$) $`\xi (G)+\eta (G)=\alpha (G)`$;
($`\mathrm{𝑖𝑖}`$) $`\sigma (G)+\eta (G)=\mu (G)`$;
($`\mathrm{𝑖𝑖𝑖}`$) $`\xi (G)+2\eta (G)+\sigma (G)=n(G)`$.
Proof. Suppose that $`\xi (G)+\eta (G)=\alpha (G)`$. According to Corollary 3.3, we get that $`\mu (G)=\alpha (G)+\sigma (G)\xi (G)=\xi (G)+\eta (G)+\sigma (G)\xi (G)=\eta (G)+\sigma (G).`$ The converse is proven in the same way.
Suppose $`\xi (G)+2\eta (G)+\sigma (G)=n(G)`$. Proposition 3.5 claims that $`\xi (G)+\eta (G)\alpha (G)`$ and $`\sigma (G)+\eta (G)\mu (G)`$. Together with $`\alpha (G)+\mu (G)=n(G)`$, which is true for König-Egerváry graphs, it gives us the two equalities needed. Conversely, if, for instance, $`\xi (G)+\eta (G)=\alpha (G)`$ then, as we already proved, $`\sigma (G)+\eta (G)=\mu (G)`$. Summing these two equalities we obtain $`\xi (G)+2\eta (G)+\sigma (G)=n(G)`$.
## 4 König-Egerváry Graphs for which $`\xi +\eta =\alpha `$
###### Lemma 4.1
Let $`G`$ be a König-Egerváry graph and $`G_0=GN[core(G)]`$. If $`G_0`$ has a unique perfect matching then its $`\alpha `$-critical edges coincide with its $`\mu `$-critical edges.
Proof. By Theorem 2.3, it is enough to show that all the edges of $`M`$ (the unique perfect matching of $`G_0`$) are also $`\alpha `$-critical.
According to Proposition 2.1, we may write $`G`$ as $`G=SH`$, where $`S\mathrm{\Omega }(G)`$ and $`H=G[VS]`$ has $`n(H)=\mu (G)`$. By virtue of Lemma 3.4($`\mathrm{𝑖𝑖}`$), $`G_0`$ has $`\alpha (G_0)=\left|Score(G)\right|=q`$ and $`core(G_0)=\mathrm{}`$. Let $`M=\{a_ib_i:1iq\}`$ and suppose that $`\{a_i:1iq\}=AS`$. We shall show that any $`a_ib_iM`$ is $`\alpha `$-critical, by exhibiting a maximum stable set $`S_0`$ in $`G_0`$ that satisfies: $`b_iS_0`$ and $`S_0N(a_i)=\{b_i\}`$. For the sake of simplicity, let us take $`i=1`$. In the sequel, if $`DV\left(G_0\right)`$, then by $`M(D)`$ we mean the set of vertices, which $`D`$ is matched onto.
*Claim 1.* There exists some $`S_0\mathrm{\Omega }(G_0)`$ with $`b_1S_0`$.
Otherwise, any $`W\mathrm{\Omega }(G_0)`$ contains $`a_1`$, because $`\left|M\right|=\left|W\right|`$ and $`\left|W\{a_j,b_j\}\right|=1`$ holds for every $`j\{1,2,\mathrm{},q\}`$. Hence, it follows that $`a_1core(G_0)`$, in contradiction with $`core(G_0)=\mathrm{}`$.
*Claim 2.* The following procedure gives rise to some $`S_0\mathrm{\Omega }(G_0)`$ that contains $`b_1`$.
Input: $`G_0,A=\{a_1,a_2,\mathrm{},a_q\},b_1B=\{b_1,b_2,\mathrm{},b_q\}=M(A)`$;
Output: $`b_1S_0\mathrm{\Omega }(G_0)`$;
$`S_0:=\{b_1\}`$;
$`D:=\{b_1\}`$;
while $`(N(D)A)M(S_0)\mathrm{}`$ do
begin
Step 1. $`S_1:=S_0`$;
Step 2. $`S_0:=S_0M((N(D)A)M(S_0))`$;
Step 3. $`D:=S_0S_1`$;
end
Step 4. $`S_0:=S_0M(BS_0)`$.
Clearly, $`\left|S_0\right|=q`$ and no edge of $`G_0`$ joins some $`a_lS_0`$ to any $`b_jS_0`$, according to building procedure of $`S_0`$. Any maximum stable set $`W\mathrm{\Omega }(G_0)`$ that contains $`b_1`$ must contain also all $`b_jS_0`$, because $`\left|W\right|=\left|M\right|`$ and $`\left|W\{a_j,b_j\}\right|=1`$ holds for every $`j\{1,2,\mathrm{},q\}`$. Hence, the set $`\{b_j:b_jS_0\}`$ is stable, and consequently, we obtain that $`S_0\mathrm{\Omega }(G_0)`$. An example of $`S_0\mathrm{\Omega }(G_0)`$ obtained by this procedure is illustrated in Figure 7.
*Claim 3.* $`S_0\{a_1\}\mathrm{\Omega }(G_0a_1b_1)`$, and hence, the edge $`a_1b_1`$ is $`\alpha `$-critical in $`G_0`$.
Firstly, no $`a_iS_0`$ is adjacent to $`a_1`$, because $`a_i,a_1A`$. Secondly, no $`b_jS_0\{b_1\}`$ is adjacent to $`a_1`$, otherwise there exists an even cycle $`C`$, with half of its edges belonging to $`M`$, which means that $`\left(ME\left(C\right)\right)\left(E\left(C\right)M\right)`$ is another perfect matching in $`G_0`$, in contradiction with the premises on $`G_0`$. Therefore, $`S_0\{a_1\}\mathrm{\Omega }(G_0a_1b_1)`$ and this implies that the edge $`a_1b_1`$ is $`\alpha `$-critical in $`G_0`$. Since $`a_1b_1`$ is an arbitrary edge of $`M`$, we may conclude that all the edges of $`M`$ are $`\alpha `$-critical in $`G_0`$.
It is interesting to notice that if $`G_0`$ were bipartite for every König-Egerváry graph $`G`$, then it would be possible to prove Lemma 4.1 using only Proposition 2.6. Figures 2, 7 show that Proposition 2.6 is not enough for our purposes, because there exist non-bipartite König-Egerváry graphs $`G`$ with nonempty cores and whose $`G_0=GN[core(G)]`$ have a unique perfect matching.
###### Theorem 4.2
Let $`G`$ be a König-Egerváry graph and $`G_0=GN[core(G)]`$. Then the following assertions are equivalent:
($`i`$) $`G_0`$ has a unique perfect matching;
($`\mathrm{𝑖𝑖}`$) $`\alpha `$-critical edges of $`G_0`$ form a maximal matching in $`G_0`$;
($`\mathrm{𝑖𝑖𝑖}`$) $`\xi (G)+\eta (G)=\alpha (G)`$;
($`\mathrm{𝑖𝑣}`$) $`\sigma (G)+\eta (G)=\mu (G)`$;
($`v`$) $`\xi (G)+2\eta (G)+\sigma (G)=n(G)`$.
Proof. According to Proposition 3.2, $`G_0`$ is also a König-Egerváry graph and has a perfect matching, say $`M_0`$.
($`i`$) $``$ ($`\mathrm{𝑖𝑖}`$) If $`M_0`$ is the unique perfect matching of $`G_0`$, all its edges are $`\mu `$-critical and, by Lemma 4.1, $`\alpha `$-critical, as well. In other words, the $`\alpha `$-critical edges of $`G_0`$ form a maximal matching. The converse is true according to Proposition 2.10($`\mathrm{𝑖𝑖𝑖}`$).
($`i`$) $``$ ($`\mathrm{𝑖𝑖𝑖}`$) Assume that $`M_0`$ is the unique perfect matching of $`G_0`$. By Lemma 3.4($`\mathrm{𝑖𝑖}`$), it follows that $`\alpha (G_0)=\alpha (G)\xi (G)`$. Lemma 3.4($`\mathrm{𝑖𝑖𝑖}`$) and the uniqueness of $`M`$ imply that $`\alpha (G_0)=\eta (G_0)=\eta (G)`$. Hence, it results in $`\xi (G)+\eta (G)=\alpha (G)`$.
($`\mathrm{𝑖𝑖𝑖}`$) $``$ ($`\mathrm{𝑖𝑣}`$) $``$($`v`$) It is the claim of Proposition 3.6.
($`v`$) $``$ ($`\mathrm{𝑖𝑖}`$) By Proposition 3.1,
$$|N(core(G))|=|\{VS:S\mathrm{\Omega }(G)\}|=\sigma (G).$$
Hence, $`n(G_0)=n(G)\xi (G)\sigma (G)`$. Now, our premise claims that $`2\eta (G)=n(G_0)`$. By Lemma 3.4($`\mathrm{𝑖𝑖𝑖}`$) we obtain $`2\eta (G_0)=n(G_0)`$. According to Theorem 2.3($`\mathrm{𝑖𝑖𝑖}`$) the set of $`\alpha `$-critical edges of $`G`$ form a matching, say $`M`$. Applying again Lemma 3.4($`\mathrm{𝑖𝑖𝑖}`$), we see that $`M_0=M`$ and it consists of $`\alpha `$-critical edges of $`G_0`$.
Notice that Theorem 4.2 fails for non-König-Egerváry graphs. In Figure 8 is presented a non-König-Egerváry graph $`G`$ having $`\xi (G)=\left|\{v\}\right|=1,\eta (G)=10`$, (all the edges of the two $`C_5`$ are $`\alpha `$-critical), $`\alpha (G)=5<\mu (G)=6`$, but $`G_0=GN[core(G)]`$ owns a unique perfect matching.
Now using Theorem 4.2 we are giving a new characterization of the bipartite graphs that have a unique perfect matching (see some previous discussions of this topic in and ). This result generalizes Corollary 2.9.
###### Corollary 4.3
Let $`G`$ be a bipartite graph. Then the following assertions are equivalent:
($`i`$) $`G`$ has a unique perfect matching;
($`\mathrm{𝑖𝑖}`$) $`\alpha `$-critical edges of $`G`$ form a maximal matching;
($`\mathrm{𝑖𝑖𝑖}`$) $`\eta (G)=\alpha (G)`$;
($`\mathrm{𝑖𝑣}`$) $`\eta (G)=\mu (G)`$;
($`v`$) $`2\eta (G)=n(G)`$.
Proof. ($`i`$) $``$ ($`\mathrm{𝑖𝑖}`$) If $`M`$ is the unique perfect matching of $`G`$, all its edges are $`\mu `$-critical and, by Proposition 2.6, $`\alpha `$-critical, as well. In other words, the $`\alpha `$-critical edges of $`G`$ form a maximal matching. The converse is true according to Proposition 2.10($`\mathrm{𝑖𝑖𝑖}`$).
The other equivalences follow from Theorem 4.2, and the observation that if a bipartite graph has a perfect matching, then the two stable sets of its standard partition are maximum, and, consequently, $`\xi (G)=0`$.
It is interesting to notice that the equality $`2\alpha (G)=n(G)`$ mentioned in Corollary 2.9 follows from Corollary 4.3, but it can not join the above series of equivalences (see, for example, $`C_4`$).
Let us also observe that for the bipartite graph $`G`$ in Figure 9, the subgraph $`G_0=GN[core(G)]`$ has more than one perfect matching.
###### Proposition 4.4
If $`G`$ is a König-Egerváry graph and there is some $`S\mathrm{\Omega }(G)`$ such that the set $`W=(S,V(G)S)`$ generates a forest, then
$$\xi (G)+\eta (G)=\alpha (G),\sigma (G)+\eta (G)=\mu (G),and\xi (G)+2\eta (G)+\sigma (G)=n(G).$$
Proof. If $`G_0=GN[core(G)],A=Score(G),B=V(G)SN(core(G))`$, then Proposition 3.2($`\mathrm{𝑖𝑖𝑖}`$) implies that $`G_0`$ is also a König-Egerváry graph and has a perfect matching, say $`M`$. Let $`G_1`$ be the partial graph of $`G_0`$ having $`WE(G_0)`$ as edge set. Then, $`M`$ is a perfect matching in $`G_1`$, as well. Since $`G_1`$ is a forest, $`M`$ is unique. By Lemma 2.2, any maximum matching of $`G_0`$ is contained in $`(A,B)`$, and since the edges from $`(A,B)`$ yield a unique perfect matching, namely $`M`$, it follows that $`M`$ is the unique perfect matching of $`G_0`$ itself. Hence, according to Theorem 4.2, we obtain that $`\xi (G)+\eta (G)=\alpha (G)`$. By Proposition 3.5($`\mathrm{𝑖𝑖𝑖}`$), it implies $`\sigma (G)+\eta (G)=\mu (G)`$, and immediately $`\xi (G)+2\eta (G)+\sigma (G)=n(G)`$.
It is worth observing that if $`(S,V(G)S)`$ generates a forest for some $`S\mathrm{\Omega }(G)`$, this is not necessarily true for all maximum stable sets of $`G`$. For example, the graph $`G`$ presented in Figure 10($`i`$) and Figure 10($`\mathrm{𝑖𝑖}`$) has the partition $`\{S_1=\{a,b,c,d\}\mathrm{\Omega }(G),V(G)S_1\}`$ such that $`(S_1,V(G)S_1)`$ does not generate a forest, (see Figure 10($`i`$)), while for the partition $`\{S_2=\{a,b,y,z\}\mathrm{\Omega }(G),V(G)S_2\}`$ the set $`(S_2,V(G)S_2)`$ generates a forest (see Figure 10($`\mathrm{𝑖𝑖}`$)). Let us also remark that the converse of Proposition 4.4 is not generally true. For instance, the graph in Figure 10($`\mathrm{𝑖𝑖𝑖}`$) is a counterexample.
###### Corollary 4.5
If $`T`$ is a tree, then
$$\xi (T)+\eta (T)=\alpha (T),\sigma (T)+\eta (T)=\mu (T),and\xi (T)+2\eta (T)+\sigma (T)=n(T).$$
As a consequence of Corollary 4.5, we obtain:
###### Corollary 4.6
If $`T`$ is a tree, then a vertex $`vV(T)`$ is in some but not in all maximum stable sets of $`T`$ if and only if $`v`$ is an endpoint of an $`\alpha `$-critical edge.
Proof. If $`vV(T)`$ is in some but not in all maximum stable sets of $`T`$, then there exists $`S\mathrm{\Omega }(T)`$ such that $`vScore(T)`$. By Theorem 2.3, $`\alpha `$-critical edges of $`T`$ form a matching. Proposition 2.6 ensures that they are also $`\mu `$-critical, because $`T`$ is bipartite. Consequently, these edges belong to any maximum matching, which, according to Lemma 2.2, is included in $`(S,V(T)S)`$. Since, by Lemma 3.4($`i`$), no $`\alpha `$-critical edge has an endpoint in $`N[core(T)]`$, and Corollary 4.5 ensures that $`\eta (T)=\alpha (T)\xi (T)=|Score(T)|`$, we infer that $`v`$ must be an endpoint of an $`\alpha `$-critical edge.
Conversely, let $`e=vw`$ be an $`\alpha `$-critical edge in $`T`$ and $`\overline{S}\mathrm{\Omega }(Te)`$. Since $`\left|\overline{S}\right|=\alpha (Te)>\alpha (T)`$, it follows that $`v,w\overline{S}`$ and therefore, $`\overline{S}\{v\},\overline{S}\{w\}\mathrm{\Omega }(T)`$. Hence, $`v`$ is in some, namely, in $`\overline{S}\{w\}`$, but not in all maximum stable sets of $`T`$, namely, not in $`\overline{S}\{v\}`$.
Notice that Corollary 4.5 and Corollary 4.6 are not valid for general bipartite graphs (see, for instance, the graph in Figure 9).
## 5 Conclusions
In this paper we state several properties of $`\alpha `$-critical and $`\mu `$-critical edges belonging to König-Egerváry graphs. These findings generalize some previously known results for trees and bipartite graphs. We have proved that for bipartite graphs and for some special König-Egerváry graphs, their sets of $`\alpha `$-critical edges and $`\mu `$-critical edges coincide. It seems to be interesting to characterize all the graphs having this property. From the other point of view, since the $`\alpha `$-critical edges of a König-Egerváry graph span disjoint cliques of order two, one may be interested in describing the type of graphs where their $`\alpha `$-critical edges span disjoint cliques of order larger than two. Another challenging problem is to describe classes of non-König-Egerváry graphs $`G`$ satisfying $`\xi (G)+\eta (G)=\alpha (G)`$, $`\xi (G)+\eta (G)\alpha (G)`$, and/or $`\alpha (G)+\sigma (G)=\mu (G)+\xi (G)`$.
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# Physics in Ultra-strong Magnetic Fields
## Electrons at $`𝑩\mathbf{>}𝑩_𝑸`$
The significance of the quantum electrodynamic field strength, $`B_Q`$, can be understood via a simple, semi-classical argument. A classical electron gyrating in a magnetic field satisfies $`\dot{p}=ev\times B/c`$, where $`p=\gamma m_ev`$ is the momentum. Substituting $`\dot{p}=\omega p`$ and $`v=\omega r`$ in this equation and cancelling factors of $`\omega `$ (along with orbital phase factors), one finds a radius of gyration $`r=cp/eB`$, where $`p`$ is the transverse momentum ($`𝐁`$). Quantum mechanics implies $`rp\mathrm{}`$ in the ground state, thus the semi-classical gyration radius is $`r_{\mathrm{gyr}}\lambda _e(B/B_Q)^{1/2}`$, where $`\lambda _e\mathrm{}/m_ec`$ is the electron Compton wavelength. The associated momentum is $`p(\mathrm{}/r_{\mathrm{gyr}})m_ec(B/B_Q)^{1/2}`$.
This shows that electrons gyrate relativistically in fields $`B>B_Q`$. One thus expects excitation energies in excess of $`m_ec^2`$. This is borne out by the solution to the Dirac equation for an electron in a homogeneous magnetic field. The Dirac spinors are proportional to Hermite polynomials, and the energy levels or “Landau levels” are $`E_n=\left[m_e^2+p_z^2+m_e^2n(2B/B_Q)\right]^{1/2},`$ (1)in units with $`\mathrm{}=c=1`$ (adopted also in many equations that follow). The first term in the square brackets is the rest energy. The $`p_z`$-term gives the energy of motion parallel to the field, which can take a continuum of values. The discrete energy levels are given by $`n=0,1,2\mathrm{}`$ These states are also eigenstates of spin, with the $`n=0`$ ground state always having $`s=\frac{1}{2}`$. For $`p_z=0`$, the ground state energy is $`E_o=m_e`$, independent of $`B`$. In a semi-classical picture, one could say that the negative spin-alignment energy in the ground state cancels with the zero-point gyration energy. Excited Landau levels are two-fold degenerate in $`s`$. The first Landau-level excitation energy is $`\omega _B(1)=E_1E_o(2B/B_Q)^{1/2}m_e`$ for $`BB_Q`$. Because this energy is so large, electrons almost always remain in the ground state for processes thought to occur near the surfaces of ultra-magnetized neutron stars.
Electron self-interactions resolve the degeneracies of the Landau levels, and shift the ground state energy. This was first demonstrated by Schwinger, who estimated the “anomalous” magnetic moment of the electron schw48 . The relevant Feynman diagram is shown in Figure 1: a free electron (traveling upward on the page) emits a virtual photon, interacts with the magnetic field, then reabsorbs the photon. The electron’s effective spin magnetic moment is enhanced by $`(1+\alpha /2\pi )`$ to first order in $`\alpha =e^2/\mathrm{}c=1/137`$, the fine-structure constant. This results in a ground-state energy shift $`E_o=m_e[\mathrm{\hspace{0.17em}1}(\alpha /2\pi )(B/B_Q)]^{1/2}.`$ (2)
If extrapolated to $`B>B_Q`$, this formula would imply that the ground-state energy of an electron goes to zero at $`B=(2\pi /\alpha )B_Q4\times 10^{16}`$ Gauss. For stronger $`B`$, the vacuum would become unstable to pair production, with dramatic astrophysical consequencesocon68 . But eq. (2) is actually only valid in the sub-$`B_Q`$ regime. More generally, the electron’s self-energy is determined by the sum of Feynman diagrams shown in Figure 2 (ref. deme53 ). The triple line on the left-hand side represents the physical electron propagator (i.e., the probability amplitude for an electron to move from point A to point B). The double lines on the right are bare propagators for an electron in the presence of a magnetic field, corresponding to basis states with energies given by eq. (1).<sup>1</sup><sup>1</sup>1Self-interactions also occur when $`B=0`$. The double-line propagators of Fig. 2 are then replaced with single-line, free-electron propagators (plane-wave states), and the resultant energy shifts—formally divergent—are absorbed into the electron’s known rest mass by the renormalization of quantum electrodynamics. A strong magnetic field changes the self-energy when the same renormalization prescriptions are used. Note that the Schwinger diagram of Fig. 1 is included in the second diagram on the right of Fig. 2: when $`B<B_Q`$, the double-line propagators can be approximated as single-line, free electron propagators undergoing discrete, perturbative interactions with the magnetic field. Positron intermediate states are included; they correspond to a subset of the vertex time-orderings which are summed over. The lowest-order tadpole diagram gives no contribution in a homogeneous magnetic field.
When the calculation is done, it is found that the electron ground-state energy diminishes according to eq. (2) as $`B`$ increases within the Schwinger domain $`BB_Q`$, but it reaches a minimum value of $`(14.6\times 10^5)m_e`$ at $`B=0.25B_Q`$ and then rises cons72 ; gepr94 . At $`B>0.65B_Q`$ the electron grows heavier than $`m_e`$, but only slowly. The asymptotic fractional enhancement, valid at very large $`B`$, isjanc69 $`(E_om_e)/m_e=(\alpha /4\pi )\left(\left[\mathrm{ln}(2B/B_Q)\xi \frac{3}{2}\right]^2+\beta \right)`$ $`BB_Q`$ (3)to first order in $`\alpha `$, where $`\xi =0.577`$ is Euler’s constant, and $`\beta 3.9`$ is a numerical constant (estimated here from the numerical integrations of ref. cons72 ).
Thus, an electron’s ground-state energy is doubled, $`E_o2m_e`$, only at $`B10^{32}`$ Gauss. (Higher-order corrections might change this result somewhat.) Of course, the maximum fields attained in neutron stars fall far short of this. The dynamical saturation field for convective motions in nascent neutron stars is $`10^{16}`$ G; and $`B3\times 10^{17}`$ G is possible if the free energy of differential rotation in a rapidly-rotating, newborn neutron star is efficiently converted by a post-collapse dynamodt92 ; td93 . But if $`B>(8\pi PY_e)^{1/2}10^{17}`$ G, where $`P`$ is the pressure and $`Y_e`$ the electron fraction in the liquid interior of a neutron star, then buoyancy overcomes stable stratification and an inhomogeneous field is dynamically lostgold92 ; td93 .
For $`B10^{17}`$ G, eq. (3) implies $`E_om_e0.03m_e`$. Thus, magnetic self-energy corrections for electrons and positrons are probably not important over the range of magnetic fields and at the level of accuracy typically attained in neutron star astrophysics.
## Atoms and Molecules at $`𝑩\mathbf{>}𝑩_𝑸`$
At sufficiently low temperatures, a magnetar’s surface will be covered with atoms and molecules. This surface structure can have consequences for the star’s quiescent X-ray emissions, because it determines the work function for removing charged particles from the surface, as necessary for maintaining currents in the magnetospherethom00 . Such currents may result from magnetically-driven crustal deformations such as twists of circular patches of the crust.<sup>2</sup><sup>2</sup>2Crustal twists, with spiral patterns of shear strain, may be a common type of magnetically-driven deformation. The pressure in the crust is due mostly to degenerate particles (relativistic electrons, and free neutrons at densities above neutron drip), but the shear modulus is due only to relatively weak Coulomb forces of the lattice. Hence the crust is relatively incompressible, and pure shear deformations allow the largest range of motion, with the greatest energy transfer between the crust and the magnetic field. If a bundle of field lines, describing an arch in the magnetosphere, has one footpoint twisted (with the motion driven from below by the evolving field), then a current must flow along the arch to maintain the twisted exterior field, since $`𝐁𝐝\mathrm{}=4\pi I/c`$. Surface impacts of the flowing charges create hot spots at the arch’s footpoints and ultimately dissipate the exterior magnetic energy of the twist, with implications for SGR and AXP X-ray light curves and their time-variationsthom00 . Here we focus on the atomic and molecular physics that comes into play, following a paper by Ruderman rude74 and extending the arguments to $`B>B_Q`$.
The Bohr radius of a hydrogen atom is $`r_o=\lambda _e/\alpha `$. The quantum gyration radius, $`r_{\mathrm{gyr}}=\lambda _e(B/B_Q)^{1/2}`$, is smaller than $`r_o`$ for $`B>\alpha ^2B_Q=2.4\times 10^9`$ G. This is the characteristic field strength at which magnetism radically alters the atomic structure of matter.<sup>3</sup><sup>3</sup>3The largest field you are ever likely to encounter personally is $`10^4`$ G if you have an medical MRI scan. Fields $`10^9`$ G would be instantly lethal. At $`B>\alpha ^2B_Q`$, an atomic electron is constrained to gyrate along a cylinder which lies entirely within the spherical volume that the unmagnetized atom would occupy. Electrostatic attraction binds the electron strongly to the central nucleus. At $`B\alpha ^2B_Q`$ the cylinder becomes very long and narrow, and atomic binding energies are adequately given by eigenvalues of the one-dimensional Schrödinger equation. A simple, intuitive estimate—which gives a good estimate of the ground state energy despite its lack of rigor—involves idealizing the atom as a line-charge of length $`2\mathrm{}`$. For linear charge density $`e/2\mathrm{}`$, the electrostatic energy is $`\epsilon =(e^2/\mathrm{})\mathrm{ln}[\mathrm{}/r_{\mathrm{gyr}}]`$. A lower cutoff $`r_{\mathrm{gyr}}`$ is necessary because the charge distribution does not resemble a line when you get within $`r_{\mathrm{gyr}}`$ of the nucleus. It is more like a sphere, contributing an energy $`qe/r_{\mathrm{gyr}}`$ where $`q=er_{\mathrm{gyr}}/\mathrm{}`$; but this contribution can be neglected in the limit $`\mathrm{}r_{\mathrm{gyr}}`$ or $`B\alpha ^2B_Q`$. Thus, the ground state energy, including the energy of non-relativistic motion parallel to $`𝐁`$, is $`_o(\mathrm{})=(\mathrm{}^2/2m_e\mathrm{}^2)(e^2/\mathrm{})\mathrm{ln}[\mathrm{}/r_{\mathrm{gyr}}]`$. Minimizing this according to $`d_o/d\mathrm{}=0`$, we find $`\mathrm{}r_o[\mathrm{ln}(r_o/r_{\mathrm{gyr}})]^1`$. This shows that the length of the thin cylindrical atom is less than the Bohr diameter, but only by a modest, logarithmic factor. The ground state hydrogen binding energy is then $`_o(ϵ_o/4)[\mathrm{ln}(B/\alpha ^2B_Q)]^2\text{for}B\alpha ^2B_Q,`$ (4)where $`ϵ_o=\alpha ^2m_e/2=13.6`$ eV is one Rydberg. Note that $`E[\mathrm{ln}B]^2`$ energy scalings are ubiquitous in ultra-magnetized systems (cf. eqs. 3,4,5).
As $`B`$ increases beyond $`BB_Q`$, the radius of the atomic cylinder shrinks to less than the Compton wavelength but eq. (4) remains a reasonably good approximation. This is because the electron’s inertia for longitudinal motion ($`𝐁`$) stays close to $`m_e`$ in the ground-state Landau level even at $`B>B_Q`$. Equation (4) would become invalid if the longitudinal motion became relativistic. But this would require $`\mathrm{}<\lambda _e`$, which occurs only at $`B>\alpha ^2\text{exp}(2/\alpha )B_Q10^{115}`$ G. Magnetic fields can never get this strong. We will show that the vacuum breaks down at smaller $`B`$.
Equation (4) implies that the binding energy of hydrogen near the surface of a magnetar with $`B10B_Q`$, is $`_o0.5`$ keV. This is comparable to the surface temperatures of some young magnetar candidates td96 ; heyl97 .
There are two classes of hydrogenic excitations. Longitudinal excited states are well-approximated as multi-nodal eigenfunctions of the 1-D Schrödinger equation; e.g., the first excited state has a node at the position of the nucleus. Transverse excited states involve transverse displacements of the center of electron gyration away from the nucleus. Semi-classically, the electron then experiences $`𝐄\times 𝐁`$ drift, and its center of gyration moves in a circular orbit around the nucleus. (See ref. rude74 for details.) Of course, Landau-level excitations are also possible, but the excitation energy is enormous for $`B>B_Q`$. Atoms generally become unbound when such free energies are present.
Longitudinal excitations tend to require more energy than transverse, so in ultra-magnetized multi-electron atoms, orbitals corresponding to transverse hydrogenic states fill up before longitudinal. In fact, for $`B>Z^3\alpha ^2B_Q(Z/26)^3B_Q`$, where $`Z`$ is the electron number, no orbitals with longitudinal nodes are occupied, and the atomic structure is very simple rude74 ; lieb92 . Note that Fe<sup>56</sup>, which is likely to be the dominant nuclear species on a clean neutron star surface, has $`Z=26`$. Thus, this condition is satisfied on magnetars, but not on radio pulsars with fields $`10^{12}`$ G. The atomic binding energy is then $`_o(Z)(7/24)Z^3ϵ_o[\mathrm{ln}(B/Z^3\alpha ^2B_Q)]^2.`$ (5) When $`BZ^3\alpha ^2B_Q(Z/26)^3B_Q`$, atoms on a neutron star’s surface form long polymer-like molecular chains parallel to $`𝐁`$, bound by the electrostatic attraction of shared electrons. The molecular binding energy per nucleus is rude74 ; neuh87 ; lai92 $`\mathrm{\Delta }(3/2)Z^3ϵ_o(B/Z^3\alpha ^2B_Q)^{0.37}.`$ (6)Together, these results determine the approximate work function for ionic emission from a magnetar’s surfacethom00 .
## Vacuum Polarization and Radiative Processes
Photon modes in the magnetized vacuum include the extraordinary mode or E-mode, with oscillating electric vector $`𝐄_𝐄𝐁`$, and the ordinary mode or O-mode, with $`𝐄_𝐎𝐄_𝐄`$. Both electric vectors are also orthogonal to $`𝐤`$, the direction of propagation.<sup>4</sup><sup>4</sup>4Photon eigenmodes are linearly polarized, as described here, except in narrow zones of $`𝐤`$-space where the angle between $`𝐤`$ and $`\pm 𝐁`$ satisfies $`\theta _{kB}(\omega /m_e)^{1/2}(B/B_Q)^{1/2}`$ and $`\mathrm{}\omega `$ is the photon energy. For propagation along $`\pm 𝐁`$ within these zones, the E and O modes are elliptically polarized; and circularly polarized for $`𝐤\pm 𝐁`$. Due to the process shown in Fig. 3, where the double-lines are propagators for a magnetized, virtual $`e^+e^{}`$ pair, the indices of refraction of the two modes are very different at $`B>B_Q`$: <sup>5</sup><sup>5</sup>5The affect of a magnetized plasma on the eigenmodes and indices of refraction is small in comparison to the magnetic vacuum polarization so long as $`\omega \omega _{c2}(3\pi /\alpha )^{1/2}(B/B_Q)^{1/2}\omega _p`$ for $`BB_Q`$, where $`\omega _p=(4\pi N_ee^2/m_e)^{1/2}`$ is the plasma frequency and $`N_e`$ is the electron density (see buli97 and references cited therein). This is satisfied in many or most applications to observable phenomena in magnetar magnetospheres since $`\omega _{c2}=0.13(N_e/10^{23}\text{cm}^3)^{1/2}(B/10B_Q)^{1/2}`$ keV.
$`n_O=1+(\alpha /6\pi )\mathrm{sin}^2\theta _{kB}(B/B_Q)`$
$`n_E=1+(\alpha /6\pi )\mathrm{sin}^2\theta _{kB}`$.
If $`n_On_E(k\mathrm{}_B)^1`$, where $`k`$ is the wavenumber and $`\mathrm{}_B`$ is the scale-length of variation of the magnetic field, then the modes adiabatically track: E stays E and O stays O as photons move through the changing field geometry. This condition is generally satisfied for X-rays in the magnetospheres of magnetars. Shaviv, Heyl & Lithwick shav99 used geometrical optics to model magnetic lensing in the vicinity of a magnetar with a pure dipole field and a uniformly bright surface. They found O-mode image distortion and amplification, varying with viewing angle. This remarkable effect may be hard to observe in practice because fields $`B(6\pi /\alpha )B_Q10^{17}`$ G are required to produce strong lensing effects. Observations may also be complicated by non-uniform surface brightness, gravitational lensingpage95 , higher-order magnetic multipoles, photon splitting (see below) and X-ray emission from a magnetar’s diffuse, Alfvén-heated halo.
When the excitation energy of the first Landau-level is much greater than the photon energy, $`\omega _B(1)m_e[(1+2B/B_Q)^{1/2}1]\omega `$, then photon scattering off electrons is strongly suppressed in the E-mode. Semi-classically, this is easy to understand: the radiation electric field ($`𝐄_𝐄𝐁`$) is unable to significantly drive electron recoil. Paczyński first notedpacz92 that this greatly accelerates X-ray diffusion in the vicinity of magnetars, facilitating hyper-Eddington burst and flare emissions. The E-mode scattering cross section, relative to Thomson, is $`\sigma (E)/\sigma _T(\omega /m_e)^2(B/B_Q)^2`$ in the regime of possible relevance for soft gamma repeater (SGR) bursts; see §3.1 of ref. td95 for more details.
Photon splitting and merging, another important radiative effect, is depicted in Figure 4, with time advancing from left to right for splitting, and right to left for merging. These processes are kinematically forbidden in free space, but they operate at $`B>B_Q`$ because the field acts as an efficient sink of momentum. (Note the double-line, magnetized $`e^{}`$ and $`e^+`$ propagators in Fig. 4.) The dominant splitting channel is $`EOO`$. The rate for $`B>B_Q`$ and $`\omega <m_e`$ isadle71 ; td93b $`\mathrm{\Gamma }_{\mathrm{sp}}=(\alpha ^3/2160\pi ^2)\mathrm{sin}^6\theta _{kB}(\omega /m_e)^5m_e`$. (7)(Splitting $`EOE`$ also occurs, but at a lower rate.) Note that $`\mathrm{\Gamma }_{\mathrm{sp}}`$ increases steeply with increasing photon energy; but it is independent of $`B`$ for $`B>B_Q`$. At $`B<B_Q`$ the process shuts down abruptly: $`\mathrm{\Gamma }_{\mathrm{sp}}(B/B_Q)^6`$.
How does this process affect SGR burst spectra? Simple splitting cascade modelsbari95 are illustrative but not realistic since O-mode photons do not split. Realistically, one must consider the subtle interplay of splitting/merging and Compton scatteringtd95 . In particular, E-mode splitting outside the E-mode scattering photosphere produces O-mode photons which are isotropized by rapid Compton scattering. Subsequent mergers $`OOE`$ yield a quasi-isotropic E-mode source function. Only at $`B<B_Q`$ and outside the O-mode photosphere do the modes truly decouple and all photons stream outward; see §6 of ref. td95 for many more details.
## The Ultra-magnetized Pair Gas
Ultra-strong magnetic fields also have profound thermodynamic effects. The magnetized photon-pair gas gives an example. Such a gas may be created during an SGR burst or flare, when a crust fracture or other magnetically-driven instability suddenly injects a large quantity of energy into the magnetospheretd95 . The result is an optically-thick trapped fireball, confined by closed field lines, anchored to the star’s surface. The gas inside this fireball has remarkable properties, as illustrated in Fig. 5. (This figure is included here courtesy of A. Kudarikuda96 .) The figure shows the ratio of pair energy density to the photon energy density, $`\mathrm{\Lambda }U_{e^+e^{}}/U_\gamma `$, as a function of $`T`$ and $`B`$. For $`Tm_e`$ and $`T\omega _B(1)`$, the magnetic field has little effect on the ultra-relativistic pairs: $`U_{e^+e^{}}=2(7/8)aT^4`$, so $`\mathrm{\Lambda }=(7/4)`$. This should hold true across the whole right-hand side of Fig. 5, but only 1000 Landau levels were used in making this graph, so the ratio falls artificially below (7/4) at high $`T`$ and low $`B`$.
The striking peak in Fig. 5 is real, however. It occurs for pairs with non-relativistic longitudinal motion, $`Tm_e`$, and $`B>B_Q`$. In this regime, only the first Landau level is occupied: $`T\omega _B(1)`$. The peak occurs because electrons and positrons are strongly localized in directions transverse to the field: $`r_{\mathrm{gyr}}=\lambda _e(B/B_Q)^{1/2}`$. This allows more of them to be packed into a given volume of ultra-magnetized gas. Formulae for thermodynamic parameters of a pair-photon gas in various limits are given in ref. kuda96 and in §3.3 of ref. td95 .
Note that the trapped fireball of a common SGR burst, with energy $`\mathrm{\Delta }E10^{41}`$ ergs confined within a volume of order $`(\mathrm{\Delta }R)^3(10\text{km})^3`$ at $`B10B_Q`$ has a temperature $`T160`$ keV td95 ; kuda96 . This puts it right on the peak in Fig. 5 !
## Magnetic Vacuum Breakdown
We argued above that a uniformly magnetized vacuum is stable against spontaneous electron-positron pair production. Nevertheless, at sufficiently high $`B`$ the vacuum must break down. Magnetic monopoles with mass $`m_\eta `$ and magnetic charge $`\eta `$ are spontaneously created when the energy they acquire in falling across a monopole Compton wavelength, $`\epsilon \eta B(\mathrm{}/m_\eta c)`$, exceeds their rest energy $`m_\eta c^2`$. Dirac showed that a monopole charge is an integral multiple of $`\eta =(\mathrm{}c/2e)`$ from the condition that an electron wavefunction must be single-valued in the field of a monopole dira31 . Thus, magnetic fields can never get stronger than $`B_{\mathrm{max}}\alpha (m_\eta /m_e)^2B_Q`$ . (8)A firm upper bound is $`B_{max}10^{55}`$ G for Planck-mass monopoles, $`m_\eta =10^{19}`$ GeV. GUT theories predict $`m_\eta =10^{16}`$ GeV or $`B_{\mathrm{max}}10^{49}`$ G. Superstring/M-theory predicts intermediate values: $`m_\eta =\alpha _s^{1/2}=10^{17}`$$`10^{18}`$ GeV, where $`\alpha _s`$ is the string tension, thus $`B_{\mathrm{max}}10^{51}10^{53}\text{G}.`$ (9) New work shows that the energy scale for quantum gravity could be as low as $`M_o1`$ TeV if there exist “large” extra dimensions to space arka98 . The extra dimensions are wrapped in closed geometries (e.g., circles) of size $`L1(M_o/1\text{TeV})^2`$ millimeter for two extra dimensions, or $`L\mathrm{}_p(M_o/M_p)^{(n+2)/n}\mathrm{}_p`$ for $`n`$ extra dimensions; where $`M_p`$ the Planck mass and $`\mathrm{}_p`$ is the Planck length. This would imply a small limiting field strength: $`B_{\mathrm{max}}10^{23}(m_\eta /1\text{TeV})^2`$ G. However, there is no experimental evidence for large extra dimensions at the present time. The most plausible upper limit is given by eq. (9).
Thus, a vast range of tremendous field strengths are possible in Nature. We don’t yet know any objects that generate such fields, but some possibilities have been suggested. For example, superconducting cosmic strings—if they exist—could generate fields $`10^{30}`$ G in their vicinitiesostr87 . Perhaps future astrophysicists will regard neutron star magnetic fields as mild !
## Magnetar Spindown
In this final section, I consider a topic of great current interest, namely recent observations of soft gamma repeater spindown histories kouv98 ; hur99a ; kouv99 ; mura99 ; mars99 ; wood99 , and their interpretation in the context of the magnetar model. At present, the most promising scenario involves episodic, wind-aided spindown (§4 in ref. thom00 ). This is based upon several background developments. In 1995 Thompson and I proposed that frequent, small-scale fractures in the crust of a young magnetar produce quasi-steady seismic and magnetic vibrations, energizing the magnetosphere and driving a diffuse, relativistic outflow of particles and Alfvén waves (§7.1.2 in ref. td95 ). A year later we made a first estimate of this outflow’s power td96 . Thompson & Blaes subsequently noted that a magnetar’s rate of spindown is greatly accelerated by such a wind (§VII B in ref. tb98 ). All of this work pre-dated the discovery of X-ray pulsations from SGRs kouv98 .
How strong is the wind? The wind luminosity, $`L_W`$ scales roughly with the magnetic energy density in the deep crust,<sup>6</sup><sup>6</sup>6Most of the free energy in magnetars is stored in the internal magnetic field, probably in toroidal and high-multipole components, so $`B_{\mathrm{crust}}`$ is usually much greater than $`B_{\mathrm{dipole}}`$. $`B_{\mathrm{crust}}^2`$ td96 . But if $`B_{\mathrm{crust}}>(4\pi \mu )^{1/2}6\times 10^{15}`$ G, where $`\mu `$ is the shear modulus in the deep crust, then evolving magnetic stresses overwhelm lattice stresses and the crust deforms plastically instead of fracturing, choking off the Alfvén-powered wind. This suggests an upper limit $`L_W5\times 10^{36}`$ erg s<sup>-1</sup> for a $`10^4`$-year-old magnetar td96 . In 1996, we proposed that a wind operating near this upper limit could account for radio synchrotron nebula that seemed to surround SGR 1806$``$20 kulk94 . However, we now know that the SGR is not coincident with this nebula hurl99 . There is no direct observational evidence for a quasi-steady wind from any SGR. Magnetar winds must be mild enough to produce no detected radio emission, with $`L_W`$ probably much less than the theoretical upper limit of ref. td96 , because this limiting value assumed optimal conditions, including the dubious application of a formula at the edge of the regime where it breaks down (i.e., $`B_{\mathrm{crust}}B_\mu `$). It is likely that $`L_W`$ is comparable to the steady X-ray luminosity emitted by the hot stellar surface and Alfvén-energized halo: $`L_W10^{35}`$$`10^{36}`$ erg s<sup>-1</sup>.
Rothschild, Marsden and Lingenfelter have plotted two graphs, included in this volume, which nicely elucidate constraints on SGR spindown for constant $`L_W`$ and $`B_{dipole}`$, based upon formulae derived independently in refs. hard99 ; thom00 .<sup>7</sup><sup>7</sup>7These references correct an inaccuracy in the original wind-aided spindown rate given by ref. tb98 . These plots show that wind luminosities $`L_W10^{36}`$ erg s<sup>-1</sup> imply $`B_{\mathrm{dipole}}10^{14}`$ G in order to match the observed values of $`P`$ and $`\dot{P}`$; but the implied stellar ages are then moderately shorter than the estimated ages of the putative associated supernova remnants (SNRs). Of course, the SNR associations or ages may be unreliable, since the SGRs lie far from the SNR centers, and the ages are only rough order-of-magnitude estimates. However, we favor a different interpretation: the wind is probably episodic, so the effective spindown age of the star is less than the SNR age. In particular, we proposed (in §4.1–4.2 of ref. thom00 ) that strong winds and rapid spindown prevail only during limited episodes of a young magnetar’s life, when it is magnetically active and observable as an SGR. This fits in nicely with observations of anomalous X-ray pulsars (AXPs). These objects have spindown ages $`P/2\dot{P}`$ that are comparable or longer than the ages of their associated SNRsgott99 ; kasp99 , suggestive of young magnetars observed during their non-windy, inactive episodes. A fully consistent scenario is possible thom00 .
Note, incidentally, that $`B_{\mathrm{dipole}}B_Q`$ is possible in a magnetar if the lowest-order magnetic moment decays quickly, e.g., via the Flowers-Ruderman instabilityfr77 (see §14.2 and 15.2 in ref. td93 ; §7.1.2 in ref. td95 ). This is because a magnetar is a magnetically-powered neutron star: its emissions depend upon the total free energy and configuration of its magnetic field, not simply upon its exterior dipole moment. The light curve of the 1998 August 27 giant flare gives evidence for strong higher-order multipole moments in SGR 1900+14fero00 . SGR bursts give evidence for magnetically-powered activity (e.g., refs. td95 ; fero00 ; palm99 ; gogu99 ).
Marsden et al. mars99 also suggested that the spindown rate of SGR1900+14 was enhanced by a factor $`2`$ during the summer of 1998. In the context of the magnetar model, this could mean that $`L_W`$ increased by a factor $`4`$. Possible evidence for this comes from $`\dot{P}`$ measurements during RXTE runs immediately preceding and following the interval in questionkouv99 ; but it should be noted that RXTE was observing the SGR at those times as a “target of opportunity” because the star was emitting hundreds of bursts kouv99 ; wood99 . Transient accelerated spindown during episodes of vigorous bursting can occur in the magnetar model, because the relativistic outflow may be enhanced. But only a handful of bursts were detected by BATSE during mid-summer of 1998 wood99 , and the RXTE/ASCA determination of $`\dot{P}`$ between 1998 Aug. 28 and Sept. 17 was $`6.2\times 10^{11}`$ s/swood99 (a number that was rounded up to $`1.\times 10^{10}`$ in ref. mura99 ). Furthermore, spindown rates measured over short time intervals, such as during single RXTE runs, can be affected by other transient or periodic effects (e.g., free precession: see §4.3 in ref. thom00 ; also ref. mela99 ).
Although an increase in $`\dot{P}`$ by $`2`$ during the summer of 1998 cannot be ruled out, we suspect that the average spindown rate was similar to that which prevailed at other times during the past few years, and the observed shift in the spindown history was due to an abrupt spindown glitch during the extraordinary giant flare of 1998 August 27th. Such a glitch could be caused by the unpinning of crustal superfluid vortices in a magnetar with a crust that has been deformed by evolving magnetic stressesthom00 .
In conclusion, we have come a long way from the days of ref. dt92 when simple magnetic dipole radiation seemed to be an adequate idealization for SGR spindown! More observations are needed to determine whether glitches really occur in SGRs and AXPs td96 ; heyl99 ; thom00 and what sign they may have; whether these stars exhibit free precession (which could give us the first direct measure of a magnetar’s internal fieldmela99 ; thom00 ); and to further test and constrain models of these fascinating stars.
Acknowledgments: This work was supported by NASA grant NAG5-8381 and Texas Advanced Research Project grant ARP-028.
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# Jet directions in Seyfert galaxies: B and I imaging data
## 1 Introduction
We have recently shown (Schmitt et al. 1997; Clarke, Kinney & Pringle 1998; see also Nagar & Wilson 1999) that there is no correlation between the position angle of radio jets and disk major axes in Seyfert galaxies, confirming previous results based on smaller samples (Ulvestad & Wilson 1984; Brindle et al. 1990; Baum et al. 1993). Clarke et al. (1998) and Nagar & Wilson (1999) showed, using a statistical inversion technique, that the observed values of $`\delta `$ (the difference between the position angle of the jet and the host galaxy disk major axis) and $`i`$ (inclination of the galaxy disk relative to the line of sight) can be reproduced by a homogeneous distribution of angles $`\beta `$ between the jet and the galaxy disk axis.
These results contradict the expectation that the jets should be aligned perpendicular to the galaxy disk. The simplest assumption about the feeding of the accretion disk and the black hole suggests that the gas comes from the host galaxy disk, so it is natural to expect both disks to be aligned and have the same angular momentum vector. Since jets are emitted perpendicular to the accretion disk, we would expect them to be aligned with the host galaxy minor axis, which is not observed. These studies give us information about the inner workings of Seyferts and may shed some light on the processes involved in the feeding of the AGN.
Although the results from Clarke et al. (1998) and Nagar & Wilson (1999) were statistically significant, they had two major limitations, their samples and most of their measurements were obtained from the literature. This indicates that their results could be biased by selection effects, like the preferential selection of galaxies which have jets shining into the plane of the galaxy, resulting in brighter radio emission and narrow line regions, which would be easier to detect. From the point of view of the data, using measurements collected from the literature can also influence the results, since different authors are likely to measure the position angle of radio jets, the disk inclination and the position angle of the host galaxy major axis using different techniques and data of different quality.
In order to improve the data relative to previous studies, we obtained radio continuum maps at 3.6cm, optical broad band images and spectroscopy for a sample of Seyfert galaxies selected from a mostly isotropic property, the flux at 60$`\mu `$m. In this way we avoid selection effects and create a homogeneous database, with measurements done using a consistent technique.
Another possible improvement which will be used in the analysis paper (Kinney et al. 2000), is the distinction between which side of the galaxy minor axis is closer to Earth. According to Clarke et al. (1998), this information can improve the statistical determination of the $`\beta `$-distribution by a factor of 2. One way to obtain this information is from the inspection of dust lanes in the galaxies’ images. Dust lanes can be seen in the near side of the galaxy, because they are highlighted against the background bulge light. Due to this fact, we decided to obtain images in the B and I bands, with a large wavelength separation, which will allow us to search for dust lanes. Another way to obtain this information is from the rotation curve of the galaxy. Knowing which side of the galaxy is approaching us and assuming that the spiral arms are trailing, we can determine which side of the minor axis is closer to Earth. In order to do this, we obtained long-slit spectra, with the slit aligned close to the host galaxy major axis, for several objects in our sample.
In this paper we present the broad-band B and I imaging data. The radio continuum observations and optical spectroscopic data will be presented elsewhere. In Section 2 we present the samples used in our study. The description of the observations and reductions is given in Section 3, and the measurements are presented in Section 4. A summary is given in Section 5.
## 2 Sample
### 2.1 60$`\mu `$m sample
In order to avoid selection effects as much as possible, we have chosen a sample from a mostly isotropic property, the flux at 60$`\mu `$m. According to the torus models of Pier & Krolik (1992), which are the most anisotropic and hence the most conservative models, the circumnuclear torus radiates nearly isotropically at 60$`\mu `$m.
Our sample includes 88 Seyfert galaxies (29 Seyfert 1’s and 59 Seyfert 2’s), which correspond to all galaxies from the de Grijp et al. (1987, 1992) sample of warm IRAS galaxies with redshift z$`0.031`$. The galaxies in this sample were selected based on the quality of the 60$`\mu `$m flux, Galactic latitude $`|b|>20^{}`$, and 25$`\mu `$m$`60\mu `$m color in the range $`1.5<\alpha (25/60)<0`$, chosen to exclude starburst galaxies as much as possible. The candidate AGN galaxies were all observed spectroscopically (de Grijp et al. 1992) to confirm their activity class as being Seyfert 1 or Seyfert 2 and not a lower level of activity such as starburst or LINER. The distance limit of z$`0.031`$ is large enough to encompass a statistically significant number of objects yet close enough to ensure that radio features can be resolved.
Table 1 presents the galaxies in the de Grijp et al. (1987) catalog, selected for our study. We list their catalog numbers, names, coordinates, the total exposure times in the B and I bands, and the observing runs in which the galaxies were observed.
### 2.2 Additional sample
Parts of the study presented by Kinney et al. (2000) will also use an additional sample of 53 Seyfert galaxies selected from the literature. This sample comprises Seyferts known to have extended radio emission, used in previous studies (such as Schmitt et al. 1997; or Nagar et al. 1999) but which are not in the 60$`\mu `$m sample. For 20 of these galaxies (7 Seyfert 1’s and 13 Seyfert 2’s) we were able to obtain B and/or I images during our observing runs. Table 2 gives the names of the galaxies, their coordinates, total exposure times in B and I bands and the observing run in which they were observed.
Some of the galaxies in the additional sample were used in previous papers, but we now consider that they should not be included in this analysis. The reasons to exclude them are the fact that they are in interacting systems, mergers, or the radio emission is not extended enough to allow a reliable measurement of the position angle of the jet. For these galaxies, Column 8 (Comments) of Table 2 gives the reasons why they are excluded.
## 3 Observations and reductions
The data presented in this paper were obtained in 5 different observing runs, using 3 different observatories. The dates of these observing runs, corresponding telescopes and instruments are shown in Table 3.
The CTIO observations were done in the 0.9m telescope with focal ratio f/13.5, using the detector T2K6 for run $`a`$ and detector T2K3 for run $`d`$. Both CCD’s have the same plate scale, which gives a pixel size of 0.384$`{}_{}{}^{\prime \prime }pixel^1`$. The images in run $`a`$ were obtained using the whole CCD area of 2048$`\times `$2048 pixels, reading it out using 4 different amplifiers, which gives a field of view of $`13^{}\times 13^{}`$. For run $`d`$ we used only a 1024$`\times `$1024 section of the CCD, reading it out using one amplifier, which gives a field of view of $`7.5^{}\times 7.5^{}`$.
The observations at Lick Observatory were done in the 1.0m Nickel telescope with focal ratio f/17, using Dewar #5 in both runs ($`b`$ and $`c`$), which gives a pixel size of 0.248$`{}_{}{}^{\prime \prime }pixel^1`$. We used the whole CCD area (1024$`\times `$1024 pixels) for these observations, which gives a field of view of $`4.8^{}\times 4.8^{}`$. The KPNO observations were done in the 0.9m telescope with focal ratio f/7.5, using the detector T2KA, which gives a pixel size of 0.688$`{}_{}{}^{\prime \prime }pixel^1`$. We used only a 1024$`\times `$1024 section of the CCD, which gives a field of view of $`11.7^{}\times 11.7^{}`$.
We followed the same observing procedure for each one of the runs. For each night we obtained a series of bias images (between 20 and 50 exposures), dome flats (between 15 and 30 exposures per filter) and sky flats (between 5 and 10 exposures per filter). We did not obtain dark images, because our exposure times were short enough that the contribution of dark current was negligible.
The reductions were done following standard IRAF procedures. The individual images were overscanned, bias subtracted and divided by the normalized flat field. Tests showed that, for each observing run, there was no significant differences between calibration frames from individual nights. Therefore, all frames were combined and we used the resulting images, which had a higher S/N, for the data reduction. The images were flat-fielded using only the sky flats, since tests showed that the dome flats had inhomogeneous illumination.
To calibrate the images in the Cousins system, each night we observed several standard star fields from Graham (1982) and Landolt (1992). We estimate that the photometric accuracy of our observations is of the order of 0.05 mag. To avoid the saturation of the nuclear region and to eliminate cosmetic defects and cosmic rays, the images were dithered using 3 or more exposures of 400s or less. For three of the galaxies in the 60$`\mu `$m sample it was possible to obtain images in only one of the bands (I for MRK1040 and B for IRAS16382-0613 and UGC10683B). Furthermore, runs $`a`$ and $`e`$ took place close to full moon, resulting in shallower B images.
Since our images will be used to compare the position angles in the radio and optical, it is important to determine the orientation of the CCD’s relative to the equatorial plane. This was done using stars in the images, which showed that the fields are not rotated. The final orientation of the images is N up and E to the left, with an uncertainty of $`1^{}`$.
## 4 Measurements
In Figure 1 we present the I band images of the galaxies, organized following the same order of Tables 1 and 2. In the case of UGC10683B and IRAS16382-0613, for which we were not able to obtain I band images, the B band images are shown instead.
Measurements of disk ellipticities and major axis position angles (defined as the angle measured from N to E) were obtained fitting ellipses to the isophotes of the galaxies, using the routine “ellipse” in the STSDAS package of IRAF. We have chosen to fit the ellipses over the I band images because they were deeper, and also because this band is more sensitive to old stars, so the outer isophotes are not disturbed by HII regions like they can be in the B images.
The ellipses were fitted from the inner $`0.7^{\prime \prime }`$ of the I images, out to the level where the surface brightness reached the 3$`\sigma `$ level above the background (this limiting value is listed in Tables 4 and 5). The background level and its standard deviation ($`\sigma `$) were determined from several blank regions around the galaxy. For some galaxies, with bright stars close to the low surface brightness isophotes (e.g. NGC3783), the ellipse fitting procedure was truncated before it reached the 3$`\sigma `$ level, to avoid the disturbance of the fit by these stars. Ellipses centers were hold fixed at the nuclear position, and were fitted using a constant increment of the semi-major axis, 2 pixels for runs a, b, c and $`d`$ and 1 pixel for run $`e`$, which corresponds to $`0.77^{\prime \prime }`$ for runs $`a`$ and $`d`$, $`0.5^{\prime \prime }`$ for runs $`b`$ and $`c`$ and $`0.69^{\prime \prime }`$ for run $`e`$. The surface brightness of the 3$`\sigma `$ level in I was typically 21-22.5 mag arcsec<sup>-2</sup>, which corresponds to 23-24.5 mag arcsec<sup>-2</sup> in the B band, assuming that the mean color of spiral disks is (B-I)$`2`$ (Héraudeau, Simien & Mamon 1996; see also our own measurements in Figure 2).
The ellipse parameters obtained from the fit of the I band images were used to measure the surface brightness profile of the B images, thus allowing a direct comparison between the two measurements. In Figure 2 we show the surface brightness profiles, major axis position angle (PA<sub>MA</sub>) and disk ellipticity (e), defined as e=1-b/a, where b/a is the ratio between the minor and major axis. Notice that the ellipse parameters in the inner 1$`{}_{}{}^{\prime \prime }2^{\prime \prime }`$ are unreliable, because they were made on scales smaller or comparable to the seeing.
In Tables 4 and 5, for the 60$`\mu `$m and additional samples, respectively, we present the size of the ellipse major axis at the 3$`\sigma `$ level above the background, the surface brightness of this level, the integrated magnitude inside this region and the seeing during the observations, for both B and I bands. The integrated magnitudes were calculated by integrating the flux inside the ellipses corresponding to the 3$`\sigma `$ isophote, using the major axis lengths given in Tables 4 and 5, the ellipticities and PA<sub>MA</sub>’s given in Tables 6 and 7 for the 60$`\mu `$m and additional samples, respectively. Notice that, since there is a difference between the 3$`\sigma `$ level of the B and I band images, the B images are usually shallower and not as extended as the I images, the ellipse parameters used to measure the integrated magnitudes of these two bands are slightly different. This is the reason why Tables 6 and 7 give different values of ellipticity and PA<sub>MA</sub> for B and I bands.
The I band ellipticities and PA<sub>MA</sub>’s were obtained by averaging the results from the ellipses fitted between the isophotes 3$`\sigma `$ and 4$`\sigma `$ above the background (4 to 10 points, depending on the galaxy). We adopted this procedure to eliminate large spurious variations, since these values will be combined with radio measurements to study the orientation of radio jets relative to the host galaxy disk axis in these galaxies (Kinney et al. 2000). An inspection of the radial profiles in Figure 2 shows that there is not too much variation in the ellipticities and PA<sub>MA</sub>’s of the outer isophotes, with the exception usually being the galaxies close to face-on and interacting systems.
## 5 Summary
We presented B and I band images for a sample of 88 Seyfert galaxies selected from a mostly isotropic property, the flux at 60$`\mu `$m, as well as for an additional 20 Seyfert galaxies with extended radio emission. The isophotes of the I band images were fitted with ellipses to determine the surface brightness profiles, the ellipticities and position angles of the host galaxy major axis. The parameters obtained with these fits were used to measure the surface brightness profiles in the B band. These images were also used to measure the integrated B and I magnitudes of the galaxies. These measurements will be combined with information from radio observations to study the orientation of radio jets relative to the host galaxy disk (Kinney et al. 2000).
We would like to acknowledge the hospitality and help from the staff at CTIO, KPNO and Lick Observatories during the observations. We also would like to thank Blaise Canzian for useful comments on the measurement of ellipticities and position angles in spiral galaxies, as well as the anonymous referee for helpful comments. This work was supported by NASA grants NAGW-3757, AR-5810.01-94A, AR-6389.01-94A and the HST Director Discretionary fund D0001.82223. This research made use of the NASA/IPAC Extragalactic Database (NED), which is operated by the Jep Propulsion Laboratory, Caltech, under contract with NASA. We also used the Digitized Sky Survey, which was produced at the Space telescope Science Institute under U.S. Government grant NAGW-2166.
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# I Introduction
## I Introduction
The phase transition of the $`U(1)`$-Higgs theory in 3+1 dimensions at finite temperature provides an important model for cosmological phase transitions. In the high temperature limit, it reduces to the purely $`3d`$ Abelian Higgs model describing the superconducting phase transition , or certain nematic to smectic-A phase transitions in liquid crystals . The phase transition in this model is governed by the infra-red region of its spectrum of fluctuations. The nature of the phase transition depends primarily on the ratio $`m_\mathrm{H}/m_A`$ between the scalar and the gauge field mass. For superconductors, these mass scales correspond to the inverse correlation length and the inverse London penetration depth, respectively. For small values of the Higgs field mass the phase transition is strongly enough first order to cut-off long range fluctuations. This corresponds to the good type-I region for standard superconductors, $`m_\mathrm{H}/m_A<1`$. On the other hand, the type-II region corresponds to $`m_\mathrm{H}/m_A>1`$. Here, it is expected that the phase transition changes from first to second order.
A proper treatment of the long range fluctuations is decisive for an understanding of the $`U(1)`$-Higgs phase transition as they change the effective interactions between the fields. The ‘microscopic’ physics in the ultra-violet is characterised by the couplings at the short distance length scale $`1/\mathrm{\Lambda }`$ (or $`1/T`$, with $`T`$ the temperature). In turn, the physics close to a first-order phase transition depends typically on the (small) photon mass $`m_A`$, and thus requires the knowledge of the couplings at scales $`T`$.<sup>1</sup><sup>1</sup>1When speaking of ‘scales’ we have always ‘mass scales’ or ‘momentum scales’ $`k`$ in mind. The corresponding ‘length scales’ are given as $`k^1`$.
A field-theoretical approach which in principle is able to deal with the effects of long range fluctuations and which describes the related scaling of the couplings is given by the Wilsonian renormalisation group . This procedure is based on integrating-out infinitesimal momentum shells about some ‘coarse-graining’ scale $`k`$ within a (Euclidean) path-integral formulation. The infra-red effective theory obtains upon integrating the resulting flow w.r.t. $`k0`$. This way, the characteristic scaling behaviour (or ‘running’) of the couplings as functions of $`k`$, and in particular the running of the Abelian charge $`e(k)`$, is taken into account. A Wilsonian approach thus improves on perturbative resummations in that the perturbative expansion parameter $`e^2T/m_A`$ becomes now scale dependent, $`e^2(k)T/m_A(k)`$. While the former diverges close to a second-order phase transition where the photon mass vanishes, the latter remains finite in the infra-red even for $`m_A(k)0`$ due to the non-trivial scaling of the Abelian charge. The crucial role of running couplings in finite temperature phase transitions has been discussed in pure scalar theories.
In the present article we employ the Wilsonian renormalisation group to the type-I regime of the $`U(1)`$-Higgs phase transition. Our main contributions are two-fold. First, we take into account the non-trivial scaling of the Abelian charge $`e^2(k)`$, characterized by an effective Abelian fixed point which is kept as a free parameter. The infra-red effects lead the Abelian gauge coupling to cross-over<sup>2</sup><sup>2</sup>2This cross-over is not to be confused with the qualitatively different ‘cross-over’ observed in the type-II regime of 3+1 dimensional $`SU(2)`$+Higgs theory. from its slow logarithmic running in the ultra-violet (effectively $`4d`$) to a strong linear running in the infra-red (effectively $`3d`$). The characteristic scale for this cross-over depends on the precise infra-red (IR) behaviour of the Abelian charge, and is decisive for both the strength of the transition and the properties of the phase diagram. This is currently the least well understood part of the problem. Equally important is to retain the full field dependence of the effective potential (no polynomial approximation), for which an analytical expression is given in the sharp cut-off case. We obtain all thermodynamical quantities related to the first-order phase transition and study their dependence on the cross-over behaviour. Second, we present a detailed quantitative analysis of the ‘coarse-graining’ dependence of our results. This is an important consistency check for the method and the approximations involved. We give quantitative evidence for an intimate link between a truncation of the effective action, and the dependence on the coarse graining scheme, which can simply be displayed as additional ‘error-bars’ due to the scheme dependence.
The $`3d`$ $`U(1)`$-Higgs phase transition has been studied previously using flow equations , and within perturbation theory . Recent results from lattice simulations for both type-I and type-II regions have been reported as well . In , the RG flow has been studied for the type-II regime within a local polynomial approximation for the effective potential about the asymmetric vacuum up to order $`\varphi ^8`$ in order to establish the phase diagram, the relevant fixed points and the related critical indices. The polynomial approximation is expected to give reliable results for the scaling solution close to a second-order fixed point. The type-I regime has been discussed for the full potential, using a matching argument for the running of the Abelian charge. In , the large -$`N`$ limit and its extrapolation down to $`N=1`$ has been considered as well. It was pointed out that the local polynomial approximation becomes questionable close to a first-order phase transition or a tri-critical fixed point at about $`N4`$ or smaller. This was later confirmed by Tetradis , who in addition abandoned the local polynomial approximation. The present study, aiming particularly at the type-I region of the phase diagram, improves on in that the full field dependence of the effective potential will be taken into account. A quantitative description of thermodynamical observables at the phase transition requires a good accuracy for the effective potential in the first place. Our study also goes beyond the work of in three important aspects. We study the dependence of physical observables on the value of the effective Abelian fixed point. In addition, explicit analytical solutions to approximate flow equations are given, as well as a discussion of the scheme dependence.
This article is organised as follows. We introduce the Wilsonian flow equations and the particular Ansatz used for the $`U(1)`$-Higgs theory. The flows for the Abelian charge and the free energy are explained, as well as the further approximations involved (Sect. II). We then proceed with the thermal initial conditions as obtained from perturbative dimensional reduction (Sect. III) and a discussion of the phase diagram and the critical line (Sect. IV). This is followed by a computation of all relevant thermodynamical quantities at the first-order phase transition as functions of the effective Abelian fixed point, a computation of the corresponding characteristic scales, and a discussion of the approximations made (Sect. V). A quantitative study of the scheme dependence on the main characteristics of the phase transition is given (Sect. VI), followed by a summary and an outlook (Sect. VII). Three Appendices contain some more technical aspects of our analysis.
## II Flow equations
### A Wilsonian flows
Wilsonian flow equations are based on the idea of a successive integrating-out of momentum modes of quantum fields within a path-integral formulation of quantum field theory . This procedure, in turn, can also be interpreted as the step-by-step averaging of the corresponding fields over larger and larger volumes, hence the notion of ‘coarse-graining’. The modern way of implementing a coarse-graining within a path-integral formalism goes by adding a suitable regulator term $`\varphi R_k(q)\varphi `$, quadratic in the fields, to the action . This additional term introduces a new scale parameter $`k`$, the ‘coarse graining’ scale. A Wilsonian flow equation describes how the coarse grained effective action $`\mathrm{\Gamma }_k`$ changes with the scale parameter $`k`$, relating this scale dependence to the second functional derivative of $`\mathrm{\Gamma }_k`$ and the scale dependence of the IR regulator function $`R_k`$. The boundary conditions for the flow equation are such that the flow relates the microscopic action $`S=lim_k\mathrm{}\mathrm{\Gamma }_k`$ with the corresponding macroscopic effective action $`\mathrm{\Gamma }=lim_{k0}\mathrm{\Gamma }_k`$, the generating functional of 1PI Green functions.
To be more explicit, we follow the ‘effective average action’ approach as advocated in and consider the flow equation
$$\frac{}{t}\mathrm{\Gamma }_k[\mathrm{\Phi }]=\frac{1}{2}\mathrm{Tr}\left\{\left(\mathrm{\Gamma }_k^{(2)}[\mathrm{\Phi }]+R_k\right)^1\frac{R_k}{t}\right\}.$$
(1)
Here, $`\mathrm{\Phi }`$ denotes bosonic fields and $`t=\mathrm{ln}k`$ the logarithmic scale parameter. The length scale $`k^1`$ can be interpreted as a coarse-graining scale . The right hand side of eq. (1) contains the regulator function $`R_k`$ and the second functional derivative of the effective action with respect to the fields. The trace denotes a summation over all indices and integration over all momenta. The above flow interpolates between the classical and quantum effective action due to some properties of the regulator functions $`R_k`$ (see Sect. VI B). It is important to realise that the integrand of the flow equation (1), as a function of momenta $`q`$, is peaked about $`q^2k^2`$, and suppressed elsewhere. Consequently, at each infinitesimal integration step $`kk\mathrm{\Delta }k`$ only a narrow window of momentum modes contribute to the change of $`\mathrm{\Gamma }_k\mathrm{\Gamma }_{k\mathrm{\Delta }k}`$. In particular, modes with momenta $`qk`$ do no longer contribute to the running at the scale $`k`$. It is this property which justifies the interpretation of $`\mathrm{\Gamma }_k`$ as a coarse-grained effective action with modes $`qk`$ already integrated out.
For gauge theories, the flow equation (1) has to be accompanied by a modified Ward identity which has to be satisfied at each scale $`k`$. Such a requirement is necessary to guarantee that the physical Green functions obtained for $`k0`$ obey the usual Ward identity . Here, we use the background field formalism, as employed in in the context of the effective average action with a covariant gauge fixing (Landau gauge).
The flow equation couples the infinite number of operators describing an effective action with its second functional derivate. In order to solve (1), one has to truncate $`\mathrm{\Gamma }_k`$ to some finite number of operators relevant for the problem under investigation. Some systematic expansions for the flow equations are known. Apart from a weak coupling expansion, which is known to reproduce the standard perturbative loop expansion, one can use expansions in powers of the fields, derivative expansions, or combinations thereof. These latter expansions have the advantage of not being necessarily restricted to a small coupling regime. A discussion on the use of a derivative expansion in Wilsonian RG is presented in .
We now turn to our Ansatz for the Abelian Higgs model. The most important information regarding the phase structure of the model is encoded in the effective potential (or coarse grained free energy) $`U_k`$, from which all further thermodynamical quantities are derived. Equally important is the wave function renormalisation factor of the gauge fields $`Z_F`$, which encodes the non-trivial running of the Abelian charge. In turn, the wave function renormalisation factor $`Z_\phi `$ for the scalar fields is less important because the scalar field anomalous dimension remains small in the type-I region of phase transitions. Hence, we approximate the effective action $`\mathrm{\Gamma }_k`$ to leading order(s) in a derivative expansion through the following operators
$$\mathrm{\Gamma }_k[\varphi ,A]=d^dx\left\{U_k(\overline{\rho })+\frac{1}{4}Z_{F,k}F_{\mu \nu }F_{\mu \nu }+Z_{\phi ,k}(D_\mu [A]\phi )^{}D_\mu [A]\phi \right\}$$
(2)
where $`\overline{\rho }=\phi ^{}\phi `$, $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu `$ is the field strength of the electromagnetic field, and $`D_\mu `$ denotes the covariant derivative $`_\mu i\overline{e}A_\mu `$.
In principle, the flow equation can be used directly (starting with initial parameters of the 4$`d`$ theory at $`T=0`$) to compute the corresponding critical potential at finite temperature within the imaginary time formalism, or, like in , using a real-time formulation of the Wilsonian RG . Our strategy in the present case is slightly simpler. We are interested in the region of parameter space where the $`4d`$ couplings are small enough to allow a perturbative integrating-out of the super-heavy and heavy modes, i.e. the non-zero Matsubara modes for all the fields and the Debye mode. In this case, we can rely on the dimensional reduction scenario and employ the results of , who computed the initial conditions perturbatively. The result is then a purely three-dimensional theory for the remaining light degrees of freedom, whose infra-red behaviour is then studied applying the above Wilsonian renormalisation group. In the sequel, we will therefore need the flow equations for $`U_k`$ and $`e^2(k)`$ in $`3d`$. At the scale of dimensional reduction, that is the starting ultra-violet (UV) scale $`\mathrm{\Lambda }`$ of the $`3d`$-flow, we normalize the wave function factors to one, and the initial effective potential $`U_\mathrm{\Lambda }`$ is obtained from dimensional reduction.
### B Cross-over of the gauge coupling
We now consider the case $`d=3`$, and discuss the flow for the Abelian coupling. A main feature of the Abelian Higgs theory in $`3d`$ is that the Abelian charge scales in a non-trivial manner with the coarse graining scale $`k`$. The dimensionless Abelian charge in $`3d`$ is defined as
$$e_3^2(k)=\frac{\overline{e}_3^2(\mathrm{\Lambda })}{Z_F(k)k}\frac{\overline{e}_3^2(k)}{k},$$
(3)
and its scale dependence is related to the gauge field anomalous dimension $`\eta _F=_t\mathrm{ln}Z_F(k)`$ (here a function of $`k`$ and the fields) through
$$\frac{de_3^2}{dt}=e_3^2(1\eta _F).$$
(4)
The first term in (4) comes from the intrinsic dimension of the charge squared (proportional to $`k`$), while the second term proportional to the gauge field anomalous dimension accounts for the non-trivial running of the coupling. The flow (4) has always the (trivial) Gaussian fixed point given by $`e_3^2=0`$. In addition, one might encounter further non-trivial fixed points which are given implicitly through the solutions of $`\eta _F=1`$.
Both the scalar and the gauge field anomalous dimensions $`\eta _\phi `$ and $`\eta _F`$ are perturbatively small near the Gaussian fixed point, i.e. $`|\eta _\phi |`$ and $`|\eta _F|1`$. This holds true at the initial scale for $`k=\mathrm{\Lambda }`$ in the effective 3d running to be specified later. It follows that the running of the dimensionful Abelian charge is negligible near the Gaussian fixed point, $`\overline{e}^2(k)\overline{e}^2(\mathrm{\Lambda })`$. Here, the dimensionless coupling scales as $`e_3^2(k)\mathrm{\Lambda }/k`$. In this regime it is expected that standard perturbation theory gives a reliable estimate of the effective potential in this region of the parameter space .
However, for $`\eta _F<1`$ the Gaussian fixed point is IR unstable, as follows directly from (4). Therefore, when approaching the infra-red, the dimensionless Abelian charge will unavoidably grow large, scaling away from the Gaussian fixed point. In particular, it can enter into a region where $`\eta _F(e^2)`$ is no longer $`1`$. When a non-trivial fixed point is approached, i.e. $`\eta _F1`$, the suppression factor $`(1\eta _F)`$ in (4) becomes important. A strong linear running of $`\overline{e}^2k`$ (the IR region is effectively $`3d`$) will ultimately set in as soon as the deviation from the Gaussian fixed point becomes sizeable . In this regime, we expect some quantitative modifications of the predictions by perturbation theory due to the non-trivial running of the Abelian charge.
### C Abelian fixed point
The anomalous dimension $`\eta _F`$ has been calculated in . It is, in general, a complicated function of the gauge coupling, the fields, and the further parameter describing the effective action in a given approximation, like the coarse grained potential (cf. (113) of ). However, $`\eta _F`$ is proportional to $`e_3^2`$ itself, and we write it as
$$\eta _F(\overline{\rho })=\frac{e_3^2}{e_{}^2(\overline{\rho })}.$$
(5)
Given the anomalous dimension, (5) provides a definition of $`e_{}^2(\overline{\rho })`$. Our current understanding of the IR behaviour of the gauge sector hinges on the precise properties of $`\eta _F(\overline{\rho })`$, and hence of $`e_{}^2(\overline{\rho })`$.
Let us recall a few cases where $`e_{}^2`$ is approximately known. First, within standard perturbation theory, the dimensionful gauge coupling $`\overline{e}_3^2=e_3^2k=const`$ throughout. Within our formalism, the ‘no-running’ corresponds to the limit $`e_{}^2\mathrm{}`$. In this limit, the effective fixed point is independent of the fields and we can expect to be close to the results from perturbation theory, as long as additional effects due to the scalar anomalous dimension can be neglected.<sup>3</sup><sup>3</sup>3In the region where $`\lambda e^2`$ (e.g. strongly type-II superconductors), the critical behaviour of the limit $`e_{}^2\mathrm{}`$ corresponds to an effective scalar theory which belongs to a different universality class than the $`O(2N)`$ scalar theory obtained for $`e_{}^20`$ .
Second, consider the large -$`N`$ limit of the $`U(1)`$-Higgs model, where $`N`$ denotes the number of complex scalar fields. In this limit, the flow (4) is dominated by the contributions of the Goldstone modes. They overwhelm those due to the radial mode. Therefore, $`e_{}^2`$ becomes
$$e_{}^2\frac{6\pi ^2}{N}$$
(6)
close to the minimum of the effective potential.<sup>4</sup><sup>4</sup>4In (6), ‘$``$’ means equality up to a regulator scheme dependent coefficient of $`𝒪(1)`$. In particular, (6) does no longer depend on the quartic scalar coupling or the location of the v.e.v. because the massive (radial) mode is suppressed. Extrapolating (6) down to the physically relevant case $`N=1`$ corresponds to replacing the radial mode by a massless one. This yields $`e_{}^26\pi ^2`$ in accordance with the leading order result from the $`ϵ`$-expansion. This value serves as a reference value for our subsequent considerations.
Third, we recall the findings of and , where the function $`\eta _F`$ has been studied numerically for different $`N`$ within a local polynomial approximation of the flow about the non-trivial minimum at $`\overline{\rho }=\overline{\rho }_0`$ (up to $`\varphi ^8`$). It was found that the implicit solutions to $`\eta _F(e_{}^2)=1`$ for small $`N`$ (in particular $`N=1`$) can deviate considerably from the large -$`N`$ extrapolation $`6\pi ^2`$. This deviation is due to the decoupling effects of the massive mode. Still, the qualitative form of (5), where the function $`e_{}^2`$ is replaced by an effective field-independent fixed point, remains a good approximation to (5). This simplified picture persists if the field derivatives $`\mathrm{ln}[e_{}^2(\overline{\rho })]/\overline{\rho }`$ remain small within the non-convex region of the effective potential (see also the discussion in Sect. V E). This implies that the threshold effects of the radial mode for $`N=1`$ act on (4) as varying the number of scalar fields in (6).
Hence, the qualitative structure of the flow (4), to leading order, is determined by (5) with $`e_{}^2`$ given by some number, e.g. the appropriate effective fixed point. For the present purpose it is sufficient to study the flow (4) with $`e_{}^2`$ as a free parameter. The properties of the first-order phase transition depend on the size of $`e_{}^2`$. However, as we shall see in detail below, the dependence turns out to be very small for large $`e_{}^2`$: this part of the phase diagram can be studied without having a complete understanding of the underlying fixed point structure. In turn, we find a strong dependence within regions where the effective fixed point is small. For this case, a more refined analysis is required in order to provide more reliable predictions.
### D Cross-over scale
Within the remaining part of the article we approximate the anomalous dimension as described above. Hence, the eqs. (4) and (5) are easily solved by
$$e_3^2(k)=\frac{e_{}^2}{1+k/k_{\mathrm{cr}}}.$$
(7)
We note the appearance of a characteristic cross-over scale
$$k_{\mathrm{cr}}=\frac{\mathrm{\Lambda }e_3^2(\mathrm{\Lambda })}{e_{}^2e_3^2(\mathrm{\Lambda })}.$$
(8)
It describes the cross-over between the Gaussian and the Abelian fixed point, and depends on the initial conditions. For $`k>k_{\mathrm{cr}}`$ the running is very slow and dominated by the Gaussian fixed point, $`\overline{e}_3^2(k)const`$ . This corresponds also to the limit $`e_{}^2\mathrm{}`$. On the other hand, for $`k<k_{\mathrm{cr}}`$ the running becomes strongly linear and the Abelian fixed point governs the scale dependence, $`\overline{e}_3^2(k)k`$. The question as to how strong the first-order phase transition is affected by this cross-over depends on whether the cross-over scale is much larger (strong effect) or much smaller (weak effect) than the typical scales of the transition (see Sect. V D).
The cross-over scale turns negative if the initial value $`e_3^2(\mathrm{\Lambda })`$ is too big. This simply means that the flow would never be dominated by the Gaussian fixed point (see Fig. 1) in the first place (no cross-over). Although this case is interesting in its own right, this region will not be discussed any further.
### E The running potential
We now turn to the flow equation for the effective potential, which can be obtained from the flow equation (1) using the Ansatz given by (2). The resulting flow equation is a second-order non-linear partial differential equation. It has been derived originally in and reads in 3$`d`$
$$\frac{4\pi ^2}{k^2}\frac{d}{dk}U_k(\overline{\rho })=(2N1)\mathrm{}_0^3\left(\frac{U_k^{}(\overline{\rho })}{k^2}\right)+\mathrm{}_0^3\left(\frac{U_k^{}(\overline{\rho })+2\overline{\rho }U_k^{\prime \prime }(\overline{\rho })}{k^2}\right)+2\mathrm{}_0^3\left(\frac{2\overline{e}_3^2(k)\overline{\rho }}{k^2}\right)$$
(9)
for the case of $`N`$ complex scalar fields. Similar flow equations are obtained for the wave function factors $`Z_\phi `$ and $`Z_F`$, and thus for the anomalous dimensions $`\eta _\phi =_t\mathrm{ln}Z_\phi `$ and $`\eta _F=_t\mathrm{ln}Z_F`$. Here, $`\mathrm{}_0^3(\omega )`$ denotes a scheme dependent threshold function defined as
$$\mathrm{}_n^d(\omega )=\left(\delta _{n,0}+n\right)_0^{\mathrm{}}𝑑y\frac{r^{}(y)y^{1+{\scriptscriptstyle \frac{d}{2}}}}{[y(1+r)+\omega ]^{n+1}}.$$
(10)
These functions have a pole at some $`\omega <0`$ and vanish for large arguments. The function $`r(q^2/k^2)`$ is related to the regulator function $`R_k`$ introduced in (1) through
$$R_k(q^2)=Zq^2r(q^2/k^2),$$
(11)
where $`Z`$ denotes either the scalar or gauge field wave function renormalisation.<sup>5</sup><sup>5</sup>5A more detailed discussion of both $`R_k`$ and the dimensionless functions $`r(q^2/k^2)`$ is postponed until Sect. VI B.
We can distinguish three different contributions to the running of the potential (9) which are, from the left to the right, related to the massless scalar, massive scalar, and the gauge field fluctuations, respectively. Not all the three of them are of the same order of magnitude, though. Indeed, it was already noted that the gauge field fluctuations dominate (9) if the quartic scalar coupling $`\lambda `$ is much smaller than the gauge coupling squared, $`\lambda /e^21`$. This is the case for the physically relevant initial conditions, that is, for the starting point of the flow equation (9). Therefore, we can make a further approximation and neglect the contributions from the scalar field fluctuations compared to those from the gauge field. The flow equation thus takes the form
$$\frac{2\pi ^2}{k^2}\frac{d}{dk}U_k(\overline{\rho })=\mathrm{}_0^3\left(\frac{2\overline{e}_3^2(k)\overline{\rho }}{k^2}\right).$$
(12)
Integrating the approximated flow equation allows to control self-consistently whether the effects from scalar fluctuations remain negligible or not. It suffices to evaluate the right-hand side of (9) with $`U_k`$ from the solution of (12) to compare the contribution of the neglected terms to the running of, say, $`U_k^{\prime \prime }`$ with the leading contributions. It is well known that the scalar fluctuations are important for the inner part of the effective potential which becomes convex in the limit $`k0`$ . Therefore it is to be expected that this approximation becomes unreliable, within the non-convex part of the potential, at some scale $`k_{\mathrm{flat}}`$.
The solution to (12) is the first step of a systematic iteration to compute the solution to (9). The next step would be to replace $`U_k`$ on the r.h.s. of (9) by the solution to (12). Proceeding to the next iteration step the scalar fluctuations are eventually taken into account. Solving (9) with $`U_k`$ on the right-hand side replaced by the explicit solution of (12) is much easier than solving (9) directly, because the former becomes an ordinary differential equation, while the later is a partial one. This procedure can be interpreted as an expansion in terms of scalar loops around the gauge field sector. We will mainly use the first step in the sequel. In order to estimate the integrated contribution of the scalar fluctuations, we will in addition discuss the solution of (9) with $`U_k`$ on the right-hand side replaced by $`U_\mathrm{\Lambda }`$ (see Appendix C).
### F The coarse-grained free energy
The coarse grained free energy obtains as the solution to the coupled set of flow equations (4) and (9). In the present case, a solution can be written as
$$U_k(\overline{\rho })=U_\mathrm{\Lambda }(\overline{\rho })+\mathrm{\Delta }_k(\overline{\rho }).$$
(14)
Here, the term $`\mathrm{\Delta }(\overline{\rho })`$ stems from integrating out the $`3d`$ fluctuations between the scales $`\mathrm{\Lambda }`$ and $`k`$. With $`e_3^2(k)`$ from (7) and $`dU_k/dk`$ from (12), it reads
$$\mathrm{\Delta }_k(\overline{\rho })=\frac{1}{2\pi ^2}_k^\mathrm{\Lambda }𝑑\overline{k}_0^{\mathrm{}}𝑑y\frac{r^{}(y)y^{5/2}\overline{k}^3(1+\overline{k}/k_{\mathrm{cr}})}{y\overline{k}(1+r)(1+\overline{k}/k_{\mathrm{cr}})+2e_{}^2\overline{\rho }}+\mathrm{const}.$$
(15)
The constant is fixed by requiring that $`\mathrm{\Delta }_k(0)=0`$. In Eq. (15) we see that the term resulting form integrating-out $`3d`$ effective modes depends on the RS through the regulator function $`r(y)`$ and its first derivative. (Explicit expressions are given in the Appendix B).
The above expressions are enough to study all properties of the phase transitions as functions of the parameters of the potential $`U_\mathrm{\Lambda }`$.
We are aiming to use an initial condition at $`k=\mathrm{\Lambda }`$ obtained from perturbation theory in $`4d`$. This requires that the parametrisation of the $`3d`$ potential $`U_\mathrm{\Lambda }`$ is such that the matching equates the right parameters. In the universal limit $`\mathrm{\Lambda }\mathrm{}`$, the effective mass term contained via $`U_\mathrm{\Lambda }`$ is renormalised to $`U_\mathrm{\Lambda }U_\mathrm{\Lambda }C_\mathrm{\Lambda }\overline{\rho }`$. For a sharp cut-off, we find explicitly
$$C_\mathrm{\Lambda }(e)=\frac{e_{}^2}{\pi ^2}\left(\mathrm{\Lambda }k_{\mathrm{cr}}k_{\mathrm{cr}}^2\mathrm{ln}(\mathrm{\Lambda }/\mathrm{\Lambda }_0)\right).$$
(16)
For finite $`\mathrm{\Lambda }`$, this corresponds to a finite renormalisation of the parameters of the theory, i.e. the mass term, or, equivalently, a finite shift of the v.e.v. at the matching scale.<sup>6</sup><sup>6</sup>6This shift corresponds to the finite renormalisation as employed in . This finite renormalisation has its origin simply in the way how the flow equation integrates-out the $`3d`$ momentum scales. Only after this transformation it will be appropriate to identify the potential $`U_\mathrm{\Lambda }`$ at the scale of dimensional reduction with the renormalised effective potential obtained from a perturbative calculation.
## III Thermal initial conditions
We now specify in concrete terms the initial conditions for the effective $`3d`$ theory. The task is to relate the $`3d`$ renormalised parameters of the effective potential to those of the $`T=0`$ $`4d`$ theory. The initial conditions for the $`3d`$ running are the potential $`U_\mathrm{\Lambda }(\overline{\rho })`$ and the gauge coupling $`\overline{e}_3^2(\mathrm{\Lambda })`$. The effective perturbative $`3d`$ Lagrangean has been derived in . We start with the $`4d`$ effective action,
$$\mathrm{\Gamma }[\varphi ,A]=d^4x\left\{\frac{1}{4}F_{\mu \nu }F_{\mu \nu }+(𝒟_\mu \varphi )^{}(𝒟_\mu \varphi )\frac{m_\mathrm{H}^2}{2}\varphi ^{}\varphi +\frac{\lambda }{2}(\varphi ^{}\varphi )^2\right\},$$
(17)
where $`\varphi `$ is a single component $`4d`$ complex scalar field. The mass parameter $`m_\mathrm{H}`$ entering (17) denotes the $`T=0`$ Higgs boson mass. It is related to the other zero temperature parameters of the theory by
$$\frac{\lambda }{e^2}=\frac{m_\mathrm{H}^2}{M_\mathrm{W}^2}$$
(18)
with $`M_\mathrm{W}`$ the photon mass. In the phase with spontaneous symmetry breaking, $`m_\mathrm{H}^2>0`$, we have $`\varphi ^{}\varphi v^2/2=M_\mathrm{W}^2/2e^2`$. The effective action for the 3$`d`$ theory obtains as
$`\mathrm{\Gamma }_\mathrm{\Lambda }[\phi ,A]`$ $`=`$ $`{\displaystyle d^3x\left\{\frac{1}{4}F_{ij}F_{ij}+(𝒟_i\phi )^{}(𝒟_i\phi )+V_\mathrm{\Lambda }(\overline{\rho })\right\}},`$ (20)
$`V_\mathrm{\Lambda }(\overline{\rho })`$ $`=`$ $`m_3^2\phi ^{}\phi +{\displaystyle \frac{\overline{\lambda }_3}{2}}(\phi ^{}\phi )^2`$ (21)
where $`\phi `$ is the static component of $`\varphi `$ and $`i,j`$ the spatial components of $`\mu ,\nu `$. The electric component of the gauge field has been fully integrated out because it acquires a thermal (Debye) mass $`m_D`$. The effects of the fluctuation of this mode are suppressed by inverse powers of $`T`$ as $`m_\mathrm{D}T`$, like the nonstatic modes. Following , the matching conditions read to 1-loop accuracy
$`\overline{e}_3^2(\mathrm{\Lambda })`$ $`=`$ $`e^2T`$ (23)
$`\overline{\lambda }_3(\mathrm{\Lambda })`$ $`=`$ $`\left(\lambda +{\displaystyle \frac{e^4}{4\pi ^2}}\right)T{\displaystyle \frac{e^4}{4\pi }}{\displaystyle \frac{T^2}{m_\mathrm{D}(\mathrm{\Lambda })}}`$ (24)
$`m_3^2(\mathrm{\Lambda })`$ $`=`$ $`\left({\displaystyle \frac{1}{4}}e^2+{\displaystyle \frac{1}{6}}\lambda \right)T^2{\displaystyle \frac{1}{2}}m_\mathrm{H}^2{\displaystyle \frac{e^2}{4\pi }}Tm_\mathrm{D}(\mathrm{\Lambda })`$ (25)
$`m_\mathrm{D}^2(\mathrm{\Lambda })`$ $`=`$ $`{\displaystyle \frac{1}{3}}e^2T^2.`$ (26)
Using the above, and taking into account the finite renormalisation (16) as explained in Sect. II F, the renormalised effective initial potential $`U_\mathrm{\Lambda }(\overline{\rho })`$ entering (14) can be expressed in terms of the $`T=0`$ parameters and (III) as
$$U_\mathrm{\Lambda }(\overline{\rho })=m_\mathrm{R}^2\overline{\rho }+\frac{1}{2}\overline{\lambda }_\mathrm{R}\overline{\rho }^2$$
(28)
with
$`m_\mathrm{R}^2(\mathrm{\Lambda })`$ $`=`$ $`{\displaystyle \frac{1}{2}}m_\mathrm{H}^2\left({\displaystyle \frac{e^2}{4}}+{\displaystyle \frac{\lambda }{6}}{\displaystyle \frac{e^3}{4\sqrt{3}\pi }}\right)T^2+C_\mathrm{\Lambda }(e),`$ (29)
$`\overline{\lambda }_\mathrm{R}(\mathrm{\Lambda })`$ $`=`$ $`\left(\lambda +{\displaystyle \frac{e^4}{4\pi ^2}}{\displaystyle \frac{\sqrt{3}e^3}{4\pi }}\right)T,`$ (30)
and the dimensionless renormalised quartic coupling reads $`\lambda _\mathrm{R}=\overline{\lambda }_\mathrm{R}/\mathrm{\Lambda }`$. The renormalised v.e.v. at the scale of dimensional reduction follows as
$$\overline{\rho }_\mathrm{R}(\mathrm{\Lambda })=m_\mathrm{R}^2(\mathrm{\Lambda })/\overline{\lambda }_\mathrm{R}(\mathrm{\Lambda }).$$
(31)
All the $`3d`$ parameters are now defined at the reduction scale $`\mathrm{\Lambda }`$, which is on dimensional grounds linearly related to the temperature,
$$\mathrm{\Lambda }=\xi T.$$
(32)
Using eqs. (8), (23) and (32) it follows, that the cross-over scale $`k_{\mathrm{cr}}`$ is also related to $`T`$ as
$$k_{\mathrm{cr}}=\frac{\xi e^2}{\xi e_{}^2e^2}T.$$
(33)
Let us finally comment on the matching parameter $`\xi `$. On one hand, $`\xi `$ has to be smaller than $`2\pi `$, because elsewise the assumption that all heavy modes have been integrated out can no longer be maintained. On the other hand, a too small value for $`\xi `$, say $`\xi <1`$, would tend to neglect contributions from modes roughly within the window $`2\pi T`$ and $`T`$. For the problem under consideration $`\xi 1`$ turns out to be a good choice. This choice shall be adopted throughout. Our results do depend very little on a variation of this matching scale (see also the comment in Sect. V D below).
## IV The phase diagram at finite temperature
We have now all the ingredients to study in detail the phase diagram and the phase transition of scalar electrodynamics. In this section, we discuss the main characteristics of the phase diagram as well as some properties of the critical line. The following section collects our results for the thermodynamical quantities related to the first-order phase transition and a discussion of the characteristic scales of the problem.
### A The phase diagram
The ‘phases’ of scalar electrodynamics are distinguished by the location of the global minimum of the effective potential. Above the critical temperature, the ground state corresponds to vanishing field $`\overline{\rho }_0=0`$, that is, to the symmetric phase (SYM). Below the critical temperature, the ground state corresponds to $`\overline{\rho }_00`$, the phase with spontaneous symmetry breaking (SSB).<sup>7</sup><sup>7</sup>7It is sensible to speak of two distinct phases only for $`N>1`$ complex scalar fields. For $`N=1`$, the symmetry is never broken in the strict sense. However, we will stick to the usual – albeit slightly incorrect – terminology even for $`N=1`$. The corresponding phase diagram in the $`(T,m_\mathrm{H})`$-plane is displayed in Fig. 2. The phase transition between these two phases is first order for small $`\overline{\lambda }_3/\overline{e}_3^2`$, that is for small values of the Higgs field mass. In the context of superconductivity this region corresponds to the strongly type-I systems. For very large Higgs field mass, the phase transition turns second or higher order .<sup>8</sup><sup>8</sup>8The strongly type-II region has been studied using flow equations within a local polynomial approximation in . See also .
In Fig. 3, we have displayed the coarse grained free energy within the type-I region of parameters for $`m_\mathrm{H}=60`$ GeV for different scales and temperatures. At the critical temperature (left panel), it is realised that a barrier is building up for decreasing scale $`k`$, eventually creating a second minima at vanishing field. The minima are degenerate in the infra-red limit $`kk_{\mathrm{stable}}`$ (which corresponds roughly to $`k0`$ in the present approximation). Notice that the flattening of the inner part of the potential is not observed because the scalar fluctuations have been neglected at the present state. Rather, the effective potential reaches the degenerate shape already at some scale $`k_{\mathrm{stable}}`$, which should be larger than the scale where the flattening sets in.<sup>9</sup><sup>9</sup>9A quantitative discussion of these scales is given in Sect. V D below.
The temperature dependence of the coarse grained free energy at $`kk_{\mathrm{stable}}`$ is shown in the right panel. The metastability range $`\mathrm{\Delta }T=T_sT_b`$ between the barrier temperature $`T_b`$, where the potential develops a second minimum at the origin (lowest dashed curve) and the spinodal temperature $`T_s`$, where the asymmetric minimum disappears (upper-most dashed curve), is very small.
The physical quantities that characterise a first-order phase transition (except the metastability range) are defined at the critical temperature $`T_c`$, when the potential has two degenerate minima, the trivial one at $`\overline{\rho }=0`$ and a non-trivial one at $`\overline{\rho }=\overline{\rho }_00`$. The critical line of the phase diagram as depicted in Fig. 2 is obtained solving the criticality conditions
$$0=\frac{dU_k}{d\overline{\rho }}|_{\overline{\rho }=\overline{\rho }_0}$$
(35)
$$U_k(0)=U_k(\overline{\rho }_0).$$
(36)
Here we kept $`k`$ arbitrary though strictly only for $`k=0`$ are these conditions required physically. They establish a relationship between the parameters of the theory, and thereby define the critical line between the symmetric and the SSB phase in Fig. 2. It is helpful to rewrite the conditions (9) into
$`F_1\left(\overline{\rho }/T\right)`$ $`=`$ $`\lambda _\mathrm{R}`$ (38)
$`F_2\left(\overline{\rho }/T\right)`$ $`=`$ $`2{\displaystyle \frac{m_\mathrm{R}^2}{T^2}}.`$ (39)
The functions $`F_1`$ and $`F_2`$ are related to the fluctuation integral through
$`F_1(x)`$ $`=`$ $`{\displaystyle \frac{2}{x^2}}\left[\stackrel{~}{\mathrm{\Delta }}(x)x\stackrel{~}{\mathrm{\Delta }}^{}(x)\right]`$ (41)
$`F_2(x)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\hspace{0.17em}2}}{x}}\left[2\stackrel{~}{\mathrm{\Delta }}(x)x\stackrel{~}{\mathrm{\Delta }}^{}(x)\right],`$ (42)
with
$$\stackrel{~}{\mathrm{\Delta }}(\overline{\rho }/T)=\mathrm{\Delta }(\overline{\rho })/T^3.$$
(43)
The first condition determines the ratio $`x=\overline{\rho }/T_c`$ of the discontinuity to critical temperature in dependence on the $`4d`$ parameters as given through $`\lambda _\mathrm{R}(e,\lambda )`$ from (30). The second one relates the solution of (38) to the ratio of the Higgs boson mass to critical temperature and (29), and eventually to the critical temperature and the discontinuity itself.
Explicit expressions for the scale-dependent effective potential and the function $`\mathrm{\Delta }(\overline{\rho })`$ are given in the Appendix B.
### B Endpoint of the critical line
Some simple properties of the solutions to (36) can be deduced directly from the functions $`F_{1,2}`$. For $`x>0`$, these functions \[with $`\mathrm{\Delta }_k`$ from (B.9)\] are positive, finite, monotonically decreasing and vanishing for $`x\mathrm{}`$. They reach their respective maxima at $`x=0`$, with (for $`k=0`$)
$`F_1(0)`$ $`=`$ $`{\displaystyle \frac{2}{\pi ^2}}e_{}^2e^2,`$ (45)
$`F_2(0)`$ $`=`$ $`{\displaystyle \frac{2}{\pi ^2}}{\displaystyle \frac{\xi ^2e_{}^2e^2}{\xi e_{}^2e^2}}\left(1{\displaystyle \frac{e^2}{\xi e_{}^2e^2}}\mathrm{ln}\left(\xi {\displaystyle \frac{e_{}^2}{e^2}}\right)\right).`$ (46)
The renormalised $`3d`$ quartic coupling $`\lambda _\mathrm{R}`$, as given by (30) and fixed through the $`4d`$ parameters of the theory, is positive in the domain under consideration. Given the monotony property of $`F_1`$ it follows that a solution to (38) is unique (if it exists). There exists no solution for too large values of $`\lambda _\mathrm{R}`$. Its largest possible value corresponds to vanishing v.e.v. i.e. to $`x=0`$. Using eqs. (18) and (30) gives an upper bound on the scalar mass for the phase transition being first order. It reads
$$\frac{m_\mathrm{H}^2}{M_\mathrm{W}^2}\frac{2e_{}^2}{\pi ^2}.$$
(47)
For any finite value of $`e_{}^2`$ (47) predicts an upper limit for the mass of the Higgs particle. This is an immediate consequence of the existence of an effective fixed point for the running gauge coupling (4). Indeed, as the limit $`e_{}^2\mathrm{}`$ corresponds to perturbation theory we recover the standard perturbative prediction of a first-order phase transition for all Higgs boson mass. This endpoint is usually interpreted as the tri-critical point of the model, above which the phase transition turns from a first-order transition to a second-order one. However, the endpoint of the first-order transition line is within the domain of validity of the present computation only for sufficiently small values of $`e_{}^2`$.<sup>10</sup><sup>10</sup>10The endpoint presented in Fig. 2 corresponds to $`e_{}^25`$. For larger values of the Abelian fixed point, we expect that the precise location of the endpoint is also determined by the scalar field fluctuations.
In the opposite case, the smallest possible value for $`\lambda _\mathrm{R}`$ corresponds to $`x\mathrm{}`$, thus to $`\lambda _\mathrm{R}=0`$. This gives a lower bound on the mass of the Higgs particle according to
$$\frac{\sqrt{3}e}{4\pi }<\frac{m_\mathrm{H}^2}{M_\mathrm{W}^2}.$$
(48)
For $`M_W=80.6`$ GeV and $`e=0.3`$ the bound is satisfied at about $`m_\mathrm{H}16`$ GeV. This bound stems entirely from the initial conditions employed. This indicates that the dimensional reduction scenario is no longer appropriate for small $`m_\mathrm{H}`$. In the present work, we are also not interested in the region of parameter space where the Coleman-Weinberg mechanism already takes place within the original $`4d`$ theory, which happens at even smaller values for $`m_\mathrm{H}`$ (typically for $`\lambda /e^4`$ at about $`3/8\pi ^2`$ or smaller ).
## V Thermodynamics of the first-order phase transition
Here we present our results for the coarse grained free energy and related physical quantities close to the critical temperature of the first-order phase transition as a function of the effective Abelian fixed point. The initial conditions are specified through the gauge coupling at vanishing temperature $`e=0.3`$ and the photon mass $`M_\mathrm{W}=80.6`$ GeV. The ratio $`\lambda /e^2`$ of the $`4d`$ couplings ranges between $`0.060.75`$ for a Higgs field mass between $`2070`$ GeV. Relevant information is given by the critical temperature $`T_c`$, the discontinuity at the phase transition $`\overline{\rho }_0`$, the latent heat $`L`$ and the surface tension $`\sigma `$.<sup>11</sup><sup>11</sup>11A comment concerning the dimensions is in order: $`U,\sigma ,L`$ and $`\overline{\rho }`$ will be given in $`3d`$ units, unless otherwise stated. Their $`4d`$ counterparts are simply obtained by multiplying with $`T`$. We compare our findings to perturbation theory, and to lattice simulations (for the critical temperature). All our results are obtained as functions of the effective fixed point of the Abelian charge. Due to the approximations performed, they depend also on the regularisation scheme. We use a sharp cut-off regulator throughout the present section. The regularisation scheme dependence is discussed in the following section.
### A Discontinuity and critical temperature
We begin with the discontinuity and the critical temperature, which follow directly from solving the criticality conditions (36).
The critical temperature as a function of the Abelian fixed point is given in Fig. 4 for $`m_\mathrm{H}=30,50,70`$ GeV. It turns out that $`T_c`$ is rather insensitive against $`e_{}^2`$. We observe an effect of a few percent only for very small values of $`e_{}^2`$ (see also Fig. 15). This is not a feature of the Higgs mass being relatively small, as similar results are obtained for all $`m_\mathrm{H}`$.
Before continuing, let us briefly compare our findings for the critical temperature to existing lattice data. Lattice results have been reported for $`e=\frac{1}{3}`$, $`m_\mathrm{W}=80.6`$ GeV and $`m_\mathrm{H}=30`$ GeV for the non-compact $`U(1)`$-Higgs model in , and for the compact one in . The result reported in is $`T_c=131.18`$ GeV for a finite lattice spacing. The continuum limit gives the slightly lower value $`T_c=130.86`$ GeV . This is consistent with $`T_c=131.28`$ GeV, the result for the compact case . Here, for $`e_{}^2=6\pi ^2`$, we find $`T_c=128.11`$ GeV. As follows from Fig. 4, the critical temperature is essentially independent of the effective Abelian fixed point. The perturbative value is $`T_c=132.64`$ GeV . These results are in good numerical agreement.
We now turn to the discussion of the discontinuity. In Fig. 5 we compare the logarithm of the v.e.v. in $`4d`$ units (normalised to the v.e.v. at $`T=0`$) at different scales. The renormalisation of $`\overline{\rho }_0`$ between the $`T=0`$ and the $`k=\mathrm{\Lambda }`$ lines results from the integration of the heavy and super-heavy modes, given by (31). The scale $`k_{\mathrm{vev}}`$ is defined as the scale where the running of the potential minimum stops. This scale is related to the scale $`k_M`$, where the photon mass in the SSB regime is becoming larger than the coarse graining scale, and thus decouples. Indeed, in the present approximation, the flow equation for the potential minimum reads
$$\frac{d\overline{\rho }_0}{dk}=\frac{1}{\pi ^2}\frac{\overline{e}^2(k)}{\overline{\lambda }(k)}\mathrm{}_1^3(M^2(k)/k^2).$$
(49)
Here, $`\overline{\lambda }(k)=U_k^{\prime \prime }(\overline{\rho }_0(k))`$ denotes the quartic coupling at the minimum, and
$$M^2(k)=2\overline{e}^2(k)\overline{\rho }_0(k)$$
(50)
the photon mass squared. The running of the v.e.v. decouples at $`kk_{\mathrm{vev}}`$, which happens as soon as the $`3d`$ photon mass $`M`$ is sufficiently larger than the scale $`k`$ (roughly at $`M^2/k^210`$) such that the threshold function in (49) suppresses any further renormalisation ($`k_{\mathrm{vev}}/T`$ is displayed in Fig. 12).
From Fig. 5 we conclude that the main part of the actual running of the potential minimum comes from integrating-out the $`3d`$ fluctuations, as can be inferred from the wide separation of the $`k=\mathrm{\Lambda }`$ and the $`kk_{\mathrm{vev}}`$ lines as opposed to the comparatively narrow separation of the $`T=0`$ and the $`k=\mathrm{\Lambda }`$ lines.
Fig. 6 shows the v.e.v. $`\overline{\rho }_0`$ as a function of the Higgs field mass and the Abelian fixed point. The shaded region covers the region $`2e_{}\mathrm{}`$ for the Abelian fixed point value. For small $`m_\mathrm{H}`$, the effect is clearly negligible. With increasing $`m_\mathrm{H}`$, however, the influence of the running gauge coupling is increasing drastically, leading to a strong weakening of the phase transition (see also Fig.15).
Finally we compare in Fig. 7 the ratio of the $`4d`$ v.e.v. $`\varphi _0`$ to the initial $`T=0`$ v.e.v. $`v/\sqrt{2}`$ for different Abelian fixed point values with the findings of perturbation theory.<sup>12</sup><sup>12</sup>12We thank A. Hebecker for providing his data from for comparison in Figs. 7, 9, 10 and 11. Again, the shaded region covers the region $`2e_{}\mathrm{}`$ for the Abelian fixed point. We observe that the v.e.v. shows a small dependence on the Abelian fixed point for sufficiently small Higgs field mass. For larger values of $`m_\mathrm{H}`$, the v.e.v. approaches the perturbative two-loop result. It follows that the v.e.v. is rather stable against effects from the running Abelian charge, say a least for $`e_{}^2>20`$. Only for $`e_{}^24`$ the running becomes strong enough to result in a significant decrease of $`\overline{\rho }_0`$.
### B The critical potential
The critical potential is shown in Figs. 8 and 9 for $`m_\mathrm{H}=38`$ GeV. Fig. 8 gives the critical potential in units of the critical temperature for different values of $`e_{}^2`$ as functions of $`\overline{\rho }/T`$.<sup>13</sup><sup>13</sup>13Notice that comparing critical potentials (or other relevant quantities) in units of $`T`$ for different values of $`e_{}^2`$ is sensible due to the very weak dependence of $`T_c`$ on the effective fixed point (see Figs. 4 and 15). We note that for large $`e_{}^2>6\pi ^2`$, the shape of the potential is rather insensitive against a change in $`e_{}^2`$. Here, the additional scale dependence induced through the gauge coupling is quite small (a few percent). For small values of $`e_{}^2`$, the height of the barrier is reduced significantly, up to a factor of 3 at $`e_{}^2=4`$. The strong scaling of $`\overline{e}^2`$ thus weakens the phase transition considerably for small $`e_{}^26\pi ^2`$. Again, the quantitative change depends strongly on the value for the effective Abelian fixed point, if $`e_{}^26\pi ^2`$. The non-trivial running of $`\overline{e}^2(k)`$ has a stronger effect on the small $`\overline{\rho }`$ region of the potential. Here, the decoupling of the gauge field sets in only at smaller scales, which in turn results in a stronger quantitative effect due to the running gauge coupling.
Fig. 9 gives the critical potential in units of the $`4d`$ v.e.v. $`v/\sqrt{2}`$, and compares the solution of (12) with those obtained within perturbation theory (PT). Line (a) corresponds to PT to order $`(e^3,\lambda ^{3/2})`$ , line (b) to our result with $`e_{}^2=6\pi ^2`$, line (c) to PT at order $`(e^4,\lambda ^2)`$ and line (d) to our result with $`e_{}^2=4`$. For $`e_{}^2=6\pi ^2`$, the critical potential is situated half way between the one- and two-loop perturbative results. For decreasing $`e_{}^2`$, the critical potential approaches quickly the two-loop result, and becomes even smaller at about $`e_{}^24`$. It is interesting to note that a value for $`e_{}^2`$ can be found for which the two-loop perturbative result is matched perfectly.
### C Surface tension and latent heat
The interface tension for a planar interface separating the two degenerate vacua follows from (2) as
$$\sigma =2_0^{\phi _+}𝑑\phi \sqrt{Z_\phi U_{\mathrm{crit}}(\overline{\rho })}.$$
(51)
It is sensible to the actual shape of the critical potential and yields additional information regarding the strength of the phase transition. In Fig. 10 the surface tension is shown as a function of $`m_\mathrm{H}`$ and in comparison with perturbation theory. The shaded region covers the results for $`2e_{}\sqrt{6}\pi `$. We again note that the effect of the running coupling is negligible for small Higgs boson mass. In contrast to the v.e.v., the surface tension depends rather strongly on $`e_{}^2`$ already for moderate values of $`m_\mathrm{H}`$. An even stronger running of $`e_3^2`$ would lead to a dramatic decrease of the surface tension, up to several orders of magnitude.
Finally, we consider the latent heat $`L`$, defined at the critical temperature as
$$L=T\left(\frac{dU(\overline{\rho }_0)}{dT}\frac{dU(0)}{dT}\right)|_{T=T_c}$$
(52)
Using eqs. (36), (IV A) and (III) we obtain
$$L=\left(m_\mathrm{H}^22m_\mathrm{R}^2\right)\overline{\rho }_0+\frac{1}{2}\lambda _\mathrm{R}T\overline{\rho }_0^2+3\mathrm{\Delta }(\overline{\rho }_0)\overline{\rho }_0\mathrm{\Delta }^{}(\overline{\rho }_0)$$
(53)
The latent heat is related to the discontinuity and the mass of the scalar particle. Using (36), it can be shown that
$$L=\overline{\rho }_0m_\mathrm{H}^2,$$
(54)
which is also known as the Clausius-Clapeyron equation. This relation was shown to be fulfilled within an explicit gauge-invariant perturbative calculation . However, it holds not true within standard perturbation theory: the perturbative values for the latent heat as found in are all below the value given through the Clausius-Clapeyron relation (54). The deviation varies between a few percent up to 15-20$`\%`$ for $`m_\mathrm{H}`$ between 20 GeV and 70 GeV, and is larger at order $`(e^4,\lambda ^2)`$ than at order $`(e^3,\lambda ^{3/2})`$.
The latent heat in units of the critical temperature is displayed in Fig. 11 for various values of the effective Abelian fixed point, and in comparison with perturbation theory to order $`(e^3,\lambda ^{3/2})`$ and $`(e^4,\lambda ^2)`$. The shaded region covers the interval $`2e_{}\sqrt{6}\pi `$. We again observe a sharp decrease for small $`e_{}`$ and large Higgs boson mass as in Fig. 7. It is interesting to note that the curve for $`e_{}=4`$ roughly agrees with the two-loop perturbative result for all $`m_\mathrm{H}`$ above 30 GeV. This is not the case for the surface tension. Comparing Fig. 11 with Fig. 10, we notice that the effect of the running gauge coupling is more pronounced for the surface tension, because the entire region for $`\overline{\rho }\overline{\rho }_0`$ enters (51), while the latent heat is only affected by $`\overline{\rho }_0`$.
### D Characteristic scales
We discuss the results obtained so far in terms of the characteristic scales relevant for the phase transition. Most of the qualitative (and even quantitative) features can be understood once these scales are known.
In Fig. 12, we have depicted the relevant momentum scales as a function of the Higgs mass. The top line at $`k=\mathrm{\Lambda }`$ corresponds to the scale of dimensional reduction, that is, the starting point of the flow in $`3d`$. The scales $`k_s,k_{\mathrm{vev}}`$ and $`k_{\mathrm{stable}}`$ (full lines) describe characteristics of the potential, the scale $`k_{\mathrm{cr}}`$ (dashed lines, for two values of the Abelian fixed point) the characteristics of the gauge sector, and $`k_{\mathrm{flat}}`$ (dashed-dotted line) the scale where scalar fluctuations can no longer be neglected within the non-convex part of the potential. All these scales are now discussed in detail.
At $`k=k_s`$, the origin of the effective potential stabilizes, $`U^{}(\overline{\rho }=0)=0`$, as the mass term squared at vanishing field changes sign. The free energy has two local minima for scales below $`k_s`$. This scale is therefore a good estimate for the scale of discontinuity. In , an estimate for this scale has been given, based on a local polynomial approximation for the potential. Within our conventions, it reads $`k_{\mathrm{dis}}0.18e^4(T)/\lambda _\mathrm{R}(T)`$ for a sharp cut-off, and roughly coincides with $`k_s`$ as presented here ($`k_{\mathrm{dis}}/k_s`$ ranges between 1 to 3).
The scale $`kk_{\mathrm{vev}}`$ indicates when the v.e.v. $`\overline{\rho }_0`$ is within $`1\%`$ of its final value, eventually reached for $`k0`$. However, this is not yet the scale where the critical potential has reached a stable shape, which actually happens only at about $`kk_{\mathrm{stable}}`$. This results from the fact that the effective photon mass squared $`2\overline{e}^2(k)\overline{\rho }`$ (within the non-convex part of the potential) is smaller than the photon mass at the minimum in the SSB regime (50), and the decoupling takes place only at smaller scales. Here, we have obtained $`k_{\mathrm{stable}}`$ comparing the depth of the potential $`U(0)U(\overline{\rho }_0)`$ at $`\overline{\rho }_0`$ with the height of the barrier $`U(\overline{\rho }_{\mathrm{max}})U(0)`$, demanding this ratio to be below $``$ 5%. At $`k=k_{\mathrm{stable}}`$, the v.e.v. is as close as 0.1% to its final value. <sup>14</sup><sup>14</sup>14Remember that the critical potential at $`k_{\mathrm{stable}}`$, within the present approximations, is about the same as at $`k=0`$, as no substantial running takes place below $`k_{\mathrm{stable}}`$.
The cross-over scale $`k_{\mathrm{cr}}`$ characterises the cross-over from the Gaussian to the Abelian fixed point. For $`e_{}^2=6\pi ^2`$, we see that $`k_{\mathrm{cr}}`$ is about 1-2 orders of magnitude smaller than the scale $`k_s`$, which explains why the running gauge coupling has, in this case, only a small numerical effect on the properties of the phase transition. From the fact that the scales $`k_{\mathrm{vev}}`$ and $`k_{\mathrm{stable}}`$ are separated by an order of magnitude ($`k_{\mathrm{vev}}/k_{\mathrm{stable}}5`$), we can conclude that the running of the gauge coupling has a stronger effect on physical observables based on the entire effective potential (like the surface tension), than those related only to the v.e.v. (like the latent heat). This is quantitatively confirmed by the findings displayed in the Figs. 6, 7, 10 and 11. For $`e_{}^2=4`$, we realize that the corresponding cross-over scale is of the same order of magnitude as the scales $`k_s,k_{\mathrm{vev}}`$ and $`k_{\mathrm{stable}}`$.<sup>15</sup><sup>15</sup>15In Fig. 12, the scales $`k_s,k_{\mathrm{vev}}`$, $`k_{\mathrm{stable}}`$ and $`k_{\mathrm{flat}}`$ have been obtained for $`e_{}^2=6\pi ^2`$. The corresponding results for $`e_{}^2=4`$ do deviate (for larger Higgs mass) only slightly from the curves as presented here. This minor difference is of no relevance for the present discussion. This is the region where the running of the gauge coupling has a strong quantitative effect on the properties of the phase transition, leading to a significant decrease of the strength of the transition.
Finally, we have also indicated the scale $`k_{\mathrm{flat}}`$ (dashed-dotted line), which is an estimate for the scale where the flattening of the inner part of the effective potential sets in. We obtained $`k_{\mathrm{flat}}`$ from solving $`k^2+U_k^{}(\overline{\rho })0`$ numerically for $`k`$ in the non-convex part of the potential, with $`U_k`$ the leading order solution for the free energy.<sup>16</sup><sup>16</sup>16A similar though slightly shifted curve for $`k_{\mathrm{flat}}`$ is obtained from solving $`k^2+U_k^{}(\overline{\rho })+2\overline{\rho }U_k^{\prime \prime }(\overline{\rho })0`$.
In an estimate for the ratio of $`k_{\mathrm{flat}}/k_{\mathrm{stable}}`$ has been obtained, based on an investigation of the surface tension of the $`3d`$ Abelian Higgs model in the universal limit $`\mathrm{\Lambda }\mathrm{}`$. There, it was found that $`k_{\mathrm{flat}}^2/k_{\mathrm{stable}}^2\overline{e}^2/M`$, with $`M`$ being the $`3d`$ photon mass. The boundary $`k_{\mathrm{flat}}^2/k_{\mathrm{stable}}^21`$ yields the relation $`k_{\mathrm{flat}}(e^2T/2\overline{\rho }_0)^{1/4}k_{\mathrm{stable}}`$, which, using the data for $`k_{\mathrm{stable}}`$ as in Fig. 12, coincides within a few percent with the line for $`k_{\mathrm{flat}}`$ as obtained above. Corrections to the universal limit can be expanded as a series in $`M^2/\mathrm{\Lambda }^2`$ . In the present case, we start at a finite scale $`\mathrm{\Lambda }=\xi T`$, but the smallness of $`M^2/\mathrm{\Lambda }^2`$ (ranging from 0.2 to 0.001 for $`20\mathrm{GeV}m_\mathrm{H}70`$ GeV) is responsible for the small corrections with respect to the universal limit $`\mathrm{\Lambda }\mathrm{}`$. Being close to the universal limit of the effective $`3d`$ theory also explains why the dependence on the matching parameter $`\xi `$ is rather small.
We now come back to the discussion of $`k_{\mathrm{flat}}`$ from Fig. 12, which, by definition, sets the scale below which the scalar fluctuations trigger the flattening within the non-convex part of the potential, and hence the scale below which these fluctuations should no longer be neglected. First notice, that the scale of discontinuity $`k_s`$ is bigger than $`k_{\mathrm{flat}}`$ by an order of magnitude. We can thus expect that the scale of discontinuity is only weakly affected by the scalar fluctuations. Also, $`k_{\mathrm{vev}}>k_{\mathrm{flat}}`$ by a factor of $`5`$. Finally, for small Higgs field mass, $`k_{\mathrm{flat}}`$ is also smaller than $`k_{\mathrm{stable}}`$. In this region, only small quantitative changes are expected if the scalar fluctuations are taken into account. This is no longer the case for large Higgs field mass, where $`k_{\mathrm{flat}}k_{\mathrm{stable}}`$. However, as these effects concerns mainly the non-convex part of the potential, and thus quantities like the surface tension, we can still expect that the latent heat and the v.e.v. are only moderately affected.
These last observations are also relevant for the applicability of Langer’s theory of bubble nucleation. The concept of an interface tension, as defined in (51), is based on the implicit assumption that the scale $`k_{\mathrm{stable}}`$ can indeed be identified. A criterion for this being the case is the smallness of the perturbative expansion parameter. From our consideration we can conclude that this will become more and more difficult for increasing $`e^2T_c/2\overline{\rho }_01`$, that is, for very weakly first-order phase transitions.<sup>17</sup><sup>17</sup>17The treatment of very weakly first-order transitions based on coarse grained potentials has been considered in .
### E Higher order corrections
Finally, we comment on the higher order corrections which are expected from operators neglected within the present approximation. Clearly, the results presented here are affected by the approximations performed, most notably through $`(i)`$ the derivative expansion, $`(ii)`$ neglecting the scalar field fluctuations as opposed to the gauge field ones, $`(iii)`$ approximating the infra-red regime of the Abelian charge by an effective fixed point, and $`(iv)`$ computing the initial conditions perturbatively. We discuss these approximations now one by one.
$`(i)`$ The leading order terms of the derivative expansion are known to correctly describe critical equations of state and scaling solutions for a variety of $`O(N)`$-symmetric scalar models in $`3d`$. Although little is known about the convergence of such an expansion, it appeared that the smallness of the anomalous dimensions controls the influence of higher order derivative operators in the effective action. Therefore, an a posteriori consistency check for the reliability of the derivative expansion consists in computing the corresponding scalar and gauge field anomalous dimension $`\eta _\phi `$ and $`\eta _F`$. In the present case, this involves more complicated higher order threshold functions (for their definitions and further details, see ). At the scale $`kk_{\mathrm{stable}}`$, we can compute the scalar anomalous dimension self-consistently from the explicit solution for the effective potential, obtained while neglecting $`\eta _\phi `$. We find that $`|\eta _\phi |0.005`$ in the interval considered, which is consistent with our initial approximation $`\eta _\phi =0`$ and justifies the derivative expansion within the scalar sector. For $`N=1`$, the gauge field anomalous dimension $`\eta _F`$ can be estimated in a similar way. It becomes of order one only when the non-trivial fixed point is approached. We find that $`\eta _F`$ ranges from $`0.03`$ to $`0.4`$ within the range of Higgs field masses considered here and for $`e_{}^26\pi ^2`$. A main difference between the scalar and the gauge field sector is that the gauge field anomalous dimension grows large ($`\eta _F=1`$) at a scaling solution. Therefore, one expects that higher order corrections within a derivative expansion (or the momentum dependence of the gauge coupling) can become important at a scaling solution and should not be neglected. In the present case, however, the scales relevant for the first-order phase transition have been reached before the Abelian charge finally runs into its non-trivial fixed point, that is before $`\eta _F=1`$. Therefore we can expect that the derivative expansion behaves reasonably well even for the gauge field sector.
$`(ii)`$ In the same way, we can check the validity of neglecting scalar fluctuations within the non-convex part of the effective potential. It is found that the self-consistent inclusion of scalar fluctuations to leading order results in corrections of the order of a few percent, increasing with increasing Higgs field mass (see Appendix C). This agrees also with the discussion of the preceding section, where it was argued that scalar fluctuations should no longer be neglected as soon as $`k_{\mathrm{flat}}`$ is of the order of $`k_{\mathrm{stable}}`$. Clearly, the weaker the first-order phase transition the more scalar fluctuations will become relevant at the phase transition. For a quantitatively more reliable computation of thermodynamical quantities in the weakly type-I region, one has to go beyond the present approximation and include scalar fluctuations. All the present approximations can be improved in a systematic way, as has been emphasized earlier. This can be done either along the lines outlined in Sect. II E, or by a straightforward numerical integration of the flow equation as in .
$`(iii)`$ The main uncertainty in the present understanding of the $`U(1)`$-Higgs theory is linked to the gauge sector of the theory i.e. the precise infra-red behaviour of the Abelian gauge coupling. Here we have effectively parametrised this uncertainty in terms of an Abelian fixed point motivated by previous work based on large -$`N`$ extrapolations and Wilsonian RG techniques. A precise determination of the correct fixed point requires the study of the momentum and of the field dependence of the Abelian charge. Our approximation assumes that the field gradients of the function $`e_{}^2(k,\overline{\rho })`$ remain sufficiently small within the non-convex part of the potential at scales above $`kk_{\mathrm{stop}}`$. In the large -$`N`$ limit, where this fixed point is well understood, the results in the present approximation are in very good agreement with the result found within a fixed dimension computation.
$`(iv)`$ The points $`(i)(iii)`$ concerned the approximations on the level of the flow equation. These are the most important ones, because they act back on $`\mathrm{\Gamma }_k`$ upon integration of the flow. An additional approximation concerns the initial conditions to the flow. Here, they have been obtained from the dimensional reduction scenario within a perturbative loop computation. For the present purposes, it was sufficient to use a 1-loop perturbative matching as given in Sect. III, although the 2-loop matching has been reported as well . These higher order effects can be taken into account in principle; in practice, this shall not be necessary because their quantitative influence is smaller than the contributions from the scalar fluctuations for larger Higgs field mass, which have already been neglected. In any case, a small change of the initial condition cannot change the main effect reported here. Except for small Higgs field masses, the dominant contributions come from integrating-out modes in $`3d`$. This follows directly from Fig. 5, which shows that the main running of the v.e.v. takes place below the scale of dimensional reduction.
Finally, we remark that the quality of a given approximation can also be assessed by studying the dependence on the coarse graining scheme. This discussion will be the subject of the following section.
## VI Scheme dependence
All quantitative results present up to now have been obtained for a sharp cut-off regulator. It is a straightforward consequence of the Wilsonian renormalisation group approach that physical observables obtained from a solution to a Wilsonian flow equation will not depend on the precise form of the coarse-graining. Unfortunately, this conclusion holds only if the $`\mathrm{𝑓𝑢𝑙𝑙}`$ effective action is computed. On a technical level, this is barely possible, and truncations of the effective action have to be employed. It is precisely this truncation that can introduce a spurious coarse graining scheme dependence for physical observables. In this section we address the question as to what extend the physical observables as obtained in the preceding section do (or do not) depend on the precise form of the coarse graining. In doing so, we are able to present quantitative ‘error bars’ related to the scheme dependence. We also present evidence for an intimate quantitative link between the scheme dependence and the truncations employed.
### A Scheme dependence vs. truncations
Consider the case of computing some physical observable from the solution to a (truncated) Wilsonian flow. It goes without saying that a strong dependence of this observable on the coarse graining employed is not acceptable as it would cast serious doubts on the truncations performed so far. With ‘strong’ we mean ‘inducing large quantitative’, or even ‘qualitative’ changes. On the other hand, a weak scheme dependence of certain physical observables is a sign for the viability of the approximation employed. In fact, if we were able to solve the flow equations without truncating the effective action $`\mathrm{\Gamma }_k`$, the final result in the physical limit $`k0`$, which is by construction nothing else but the full quantum effective action $`\mathrm{\Gamma }`$, should not depend on the details of the particular coarse graining employed. There is little hope for this holding true for any truncation of the effective action $`\mathrm{\Gamma }_k`$ as any truncation necessarily neglects infinitely many operators.
The coarse graining procedure is implemented through the momentum-dependent operator $`R_k(q^2)`$. It couples to all the operators present in $`\mathrm{\Gamma }_k`$ in a well-defined way, that is, according to the flow equation (1). Replacing a coarse-graining by another coarse graining implies that the effective coupling of $`R_k(q^2)`$ to the operators contained in the effective action changes accordingly. A truncation of the effective action amounts to neglecting infinitely many operators to which the coarse graining, in principle, is sensitive. Therefore, studying the scheme dependence will probe whether some relevant operators (for the problem under investigation) have been neglected, or not. In this light, the indirect feed-back of some relevant operators should manifest itself through some strong eigenmode with respect to a change of the coarse graining procedure.
Although these arguments, as presented so far, are of a purely qualitative nature, we will show in the sequel that they can indeed be given a quantitative meaning.
### B Coarse grainings
Before studying in detail the scheme dependence of our results, we will briefly review the main requirements for a viable coarse graining procedure. There are basically three key points to be considered. The first one concerns the possible zero-modes of the propagators, which typically cause strong infra-red problems within perturbative loop expansions in $`d<4`$ dimensions. These are properly regularised, if
$$\underset{q^20}{lim}R_k(q^2)>0$$
(56)
holds true. This way, the effective inverse propagator for a massless mode reads $`q^2+R_k(q^2)`$, and has a well-defined infra-red limit. The second point concerns the infra-red limit of the effective action $`\mathrm{\Gamma }_k`$, which should coincide with the usual effective action for $`k0`$. This is the case, if
$$\underset{k0}{lim}R_k(q^2)=0.$$
(57)
Finally, we have to make sure that the correct initial effective action in the ultra-violet limit is approached which is guaranteed by
$$\underset{k\mathrm{}}{lim}R_k(q^2)\mathrm{}.$$
(58)
Any function $`R_k(q^2)`$ with the above properties can be considered as a coarse graining . It is convenient to re-write $`R_k`$ in terms of dimensionless functions $`r(q^2/k^2)`$ as
$$R_k(q)=Zq^2r(q^2/k^2),$$
(59)
where $`Z`$ corresponds to a possible wave-function renormalisation ($`Z_\varphi =1`$ in our approximation).
Let us introduce two classes of regulator functions which are commonly used in the literature. The first one is a class of power-like regularisation schemes given by the coarse-graining function
$$r_p(y)=y^n,$$
(60)
and $`yq^2/k^2`$. The particular case $`n=1`$ corresponds to a mass-like regulator $`R_kk^2`$, and $`n=2`$ to a quartic regulator $`R_kk^4/q^2`$. These algebraic regulators are often used because the related threshold functions can be computed analytically. On the other hand, these regulators decay only algebraically for large momenta, which can in principle lead to an insufficiency in the integrating-out of the hard UV modes.
A second convenient class of regulators consists of exponential ones, parametrised as
$$r_e(y)=\frac{1}{\mathrm{exp}(cy^n)1},$$
(61)
where $`c`$ is a constant. The exponential regulator with $`n=c=1`$ has been used previously in various numerical investigations . The suppression of large momentum modes $`q^2k^2`$ to the flow is now exponential and thus much stronger than in the case of algebraic regulators. It is expected that this property is at the basis for a good convergence of approximate solutions.
Both classes of regulator functions depend on the parameter $`n`$, with $`1n\mathrm{}`$. In the limit $`n\mathrm{}`$, they both approach to what is known as the sharp cut-off regulator, given by
$$r_s(y)=\frac{1}{\theta (y1)}1.$$
(62)
We will now consider the dependence of certain physical observables on particular choices of these regulators.
### C Tri-critical point and large -$`N`$ limit
We have given an estimate for the endpoint of the critical line in (47). Its mere existence is closely linked to the presence of an Abelian fixed point, although it will be within the domain of validity only for small values of the latter. Both functions $`F_1`$ and $`F_2`$ depend explicitly on the RS, and so does the solution to eqs. (36). In the general case, the endpoint of the critical line also depends on the RS. Instead of (47), which is the result for a sharp cutoff, we find for the general case
$$\frac{m_\mathrm{H}^2}{M^2}=\frac{8a_1}{3\pi ^2}e_{}^2,$$
(63)
where terms $`𝒪(e)`$ have been dropped. The entire scheme dependence is now encoded in the coefficient $`a_1`$, given by
$$a_1=\frac{3}{2}_0^{\mathrm{}}𝑑y\frac{r^{}(y)y^{\frac{1}{2}}}{[\mathrm{\hspace{0.17em}1}+r(y)]^3}$$
(64)
in $`d=3`$ dimensions. This coefficient belongs to a set of expansion coefficients $`a_k`$ characterising a coarse graining scheme (see Appendix A for their general definition and more details). For each of the two classes of regulators the coefficient $`a_1`$ can be calculated as a function of the parameter $`n`$. In Fig. 13, the dashed line corresponds to the power-like, and the full line to the exponential regulator class with $`c=\mathrm{ln}2`$. For this choice of $`c`$ both set of regulators are normalised to $`r(1)=1`$. For a power-like regulator, we find explicitly $`a_1=\frac{3}{4}\mathrm{\Gamma }[1+\frac{1}{2n}]\mathrm{\Gamma }[2\frac{1}{2n}]`$, and for the exponential one $`a_1=\frac{3}{8}n^1c^{1/2n}(2^{1/2n}2)\mathrm{\Gamma }[\frac{1}{2n}]`$. It is interesting to note that although these classes of regulators do have strong qualitative differences, the coefficient $`a_1`$, which only involves a folding of $`r(y)`$ over all momenta, is rather stable (i.e. $`\pm 10\%`$ about the mean value).
We shall compare the numerical value of the tri-critical point with results obtained in the large -$`N`$ limit via the $`ϵ`$-expansion or a fixed dimension computation in $`(d=3)`$ . As argued in Sect. II B, the Abelian fixed point reads $`e_{}^2=6\pi ^2/N`$ in the large -$`N`$ limit, and our above result therefore becomes
$$\frac{\lambda _3}{e_3^2}=16a_1\frac{1}{N}(9.412.0)\frac{1}{N}$$
(65)
The $`ϵ`$-expansion, to leading order, yields
$$\frac{\lambda _3}{e_3^2}=(54136ϵ)\frac{1}{N}$$
(66)
This is to be compared to the result of , which reads
$$\frac{\lambda _3}{e_3^2}=\frac{96}{\pi ^2}\frac{1}{N}9.9\frac{1}{N}$$
(67)
While (66) fails to give a reliable answer at $`ϵ=1`$, we observe that our result (65) is in good numerical agreement with (67).
### D Scheme dependence of the critical potential
Here, we consider the task of computing the critical potential for coarse grainings other than the sharp cut-off. First, we have to obtain the corresponding fluctuation integrals. The most general expression (for arbitrary scheme) has been given in Appendix B. This expression still contains an integral over momenta to be performed, which is how the scheme dependence enters into the expression for the fluctuation integral $`\mathrm{\Delta }_k`$. Then, the criticality conditions (36) have to be solved to find $`T_c`$ and $`\overline{\rho }_0`$. The sharp cut-off allowed an analytical computation of $`\mathrm{\Delta }_k`$, (B.9), and thus of the functions $`F_{1,2}`$ in (IV A).
Below, in addition to the sharp cut-off, we consider the classes of power-like regulators (60) and exponential regulators (61). From the power-like regulators, we consider the limiting cases $`n=1`$ (i.e. a mass-like regulator $`R_k=k^2`$) and $`n=\mathrm{}`$ (the sharp cutoff). As an intermediate case we consider also the case $`n=2`$ (i.e. the quartic regulator $`R_k=k^4/q^2`$). The exponential regulators are represented for $`n=1`$ (i.e. $`R_k=q^2/(\mathrm{exp}q^2/k^21)`$), and $`n=\mathrm{}`$ (the sharp cutoff). A continuity argument suggests that the critical potentials for intermediate values of the coarse graining parameter $`n`$ should appear within those limits set by $`n=1,2`$ and $`n=\mathrm{}`$.
No explicit analytical expressions for the coarse grained free energy have been found in these cases. For the mass-like and the quartic regulator we used the integrals (B.11) and (B.12), respectively, while (B.2) is used for the exponential regulator. Then, the problem of solving the criticality conditions reduces to the optimization of two integral equations.
We find that the critical temperature $`T_c`$ depends very weakly on the different schemes. Indeed, plotting $`T_c`$ as a function of the Higgs field mass we find that the lines corresponding to different schemes are almost on top of each other, inducing a relative error well below the $`1\%`$ level (and thus below the error already present due to other approximations). A similar situation holds for the v.e.v., where we find a relative error below a few percent.
In Fig. 14, the entire critical potential (in units of the $`4d`$ v.e.v.) is displayed for different coarse grainings at $`m_\mathrm{H}=38`$ GeV (left panel) and at $`m_\mathrm{H}=70`$ GeV (right panel). The labels $`s,q,m`$ and $`e`$ denote respectively the $`s`$harp cut-off, the $`q`$uartic/$`m`$ass-like regulator, and the $`e`$xponential cut-off from (61) for $`n=1`$ and $`c=1`$.
We first consider $`m_\mathrm{H}=38`$ GeV, and notice that the $`s`$ and $`q`$ lines turn out to be on top of each other. Furthermore, it is realised that the v.e.v. is nearly independent on the RS, as is the shape of the potential close to the minima. The main dependence concerns the local maximum of the critical potential. This dependence will therefore affect integrated quantities like the surface tension, but not those related to the v.e.v., like the latent heat. The error for the surface tension in the present case is about a few percent.
For $`m_\mathrm{H}=70`$ GeV the dependence on the scheme is more pronounced than in the previous case. Furthermore, the v.e.v. receives – for the mass-like regulator – a sizeable shift towards smaller values. Again, the variance is strongest around the maximum of the critical potential, and dominant in the non-convex region of the critical potential. The additional shift in the value of the v.e.v. entails a corresponding shift for the outer region of the effective potential, as opposed to the case for smaller Higgs field mass.
It is interesting to make contact with the qualitative considerations presented at the beginning of this section, and to compare the scheme dependence observed in Fig. 14 with the reliability of the coarse grained potential in its different regions, due to the approximations employed. Recall that the present computation is based on neglecting the scalar fluctuations. This approximation is more reliable for the outer part of the potential than for the non-convex part of it (more precisely, around a small region of the maximum of the inner part of the potential). Here, scalar fluctuations cause ultimately the flattening of the potential in the IR limit. While we have seen in Sect. V D that this approximation is still reliable for $`m_\mathrm{H}=38`$ GeV, we certainly expect larger corrections for $`m_\mathrm{H}=70`$ GeV (see the discussion of Sects. V D and V E). It is quite remarkable that the scheme dependence indeed seems to reflect the weakness of the approximation for this region of the potential. Our computation thus turns the qualitative statement into a quantitative one.
Finally, we briefly comment on the different regulators used. It is well-known that the mass-like regulator is marginal in the sense that it has a poor UV behavior which makes its use for certain applications questionable (a more refined discussion has been given in ). From Fig. 14, we learn that the critical potential as obtained for the mass-like regulator deviates the most from the results for the other regulators employed. Considering the class of power-like regulators, we see from Fig. 14 that the width between the quartic and the sharp cut-off limit is significantly smaller than the deviation for the mass-like regulator. This observation strongly suggests that the mass-like regulator should be discarded for quantitative considerations, although it remains, in the present example, a useful regulator for studying the main qualitative features of the problem.<sup>18</sup><sup>18</sup>18This conclusion coincides with those of based on more formal considerations regarding mass-like regulators. Discarding the mass-like regulator from our discussion, we end up with the observation that the error induced through the scheme dependence is of the same order of magnitude for algebraic as for exponential regulators. For the present case, and at this level of accuracy, no further qualitative differences are observed between the exponential regulators (61) and the power-like ones (60) for $`n2`$.
In summary, we conclude that a quantitative analysis of the scheme dependence indeed yields non-trivial information regarding the accuracy of the approximations or truncations employed, as suggested by the qualitative argument presented in Sect. VI A. In addition, we have found some evidence for why a mass-like regulator, as opposed to exponential or higher order power-like regulators, should be discarded for accurate quantitative considerations. However, as the qualitative features are still well described by a mass-like regulator, and as the quantitative deviation is not too big, this also suggests that a mass term regulator could be very useful for an error estimate.<sup>19</sup><sup>19</sup>19An error estimate based on the mass-like regulator is rather conservative as it seems to overestimate the scheme dependence. Typically, analytical computations are largely simplified for mass-like regulators, allowing for a simple cross-check of the results.
## VII Summary and outlook
We have studied in detail the first-order phase transition of Abelian Higgs models in 3+1 dimensions at finite temperature. Properties of the transition are determined by the underlying fixed point structure of the $`3d`$ theory such as the cross-over of the Abelian charge from the Gaussian to the Abelian fixed point. We computed all physical observables at the phase transition, the phase diagram in the domain of first-order transitions and the tri-critical point. The analysis has been restricted to the region of parameter space where the dimensional reduction scenario applies and a perturbative matching of the $`4d`$ parameters to the corresponding $`3d`$ ones is possible. The main contribution to the free energy (and thus to the physical observables at criticality) stem from the remaining effective $`3d`$ running for which we have used a Wilsonian renormalisation group to leading order in the derivative expansion, neglecting the scalar, but not the gauge field anomalous dimension. The latter is related to the non-trivial running of the Abelian gauge coupling, which is described by an effective fixed point. While this fixed point is well understood in the large -$`N`$ limit where the tri-critical fixed point is known, its precise form is not yet established for the relevant case of $`N=1`$. We therefore studied the parametric dependence of physical observables on the fixed point value. A quantitative discussion of the relevant physical scales, which are easily accessible within a Wilsonian framework, has also been given.
The main effect on physical observables due to the presence of a non-trivial fixed point depends on the ratio between the cross-over scale $`k_{\mathrm{cr}}`$ (which defines the cross-over to the Abelian fixed point) and the typical scales characterising the first-order phase transition (like the discontinuity scale $`k_{\mathrm{dis}}`$, or $`k_{\mathrm{stable}}`$). For $`k_{\mathrm{cr}}`$ small as compared to $`k_{\mathrm{stable}}`$ the observed dependence is weak. The sizeable deviations from the perturbative $`\overline{e}^2(k)\overline{e}^2(\mathrm{\Lambda })`$ \- behaviour only set in at very small scales below $`k_{\mathrm{stable}}`$ and are no longer relevant for the phase transition itself in this situation. The main effects are restricted to alterations in the far infra-red region, like the details of the flattening of the inner part of the potential. On the other hand, a strong dependence emerges for $`k_{\mathrm{cr}}`$ larger than $`k_{\mathrm{stable}}`$.
Most of our results for the physical observables can be summarised as in Fig. 15. Here, the reference values $`T_{\mathrm{ref}}`$ and $`\overline{\rho }_{\mathrm{ref}}`$ are given for $`e_{}=\sqrt{6}\pi `$ (which corresponds roughly to $`k_{\mathrm{cr}}k_{\mathrm{stable}}`$), and for the sharp cut-off regulator. In the present approximation, the critical temperature is insensitive to the running gauge coupling. On the other hand, the v.e.v. appears to be quite sensitive to the actual fixed point value, in particular for larger Higgs field mass. The phase transition weakens significantly for small fixed point values. The reason is that the gauge coupling is decreasing strongly for small fixed point values at scales larger than the scale where the critical potential reaches its degenerate shape, that is above the scale of decoupling. These results compare well with perturbation theory, except for very large or very small values for the Abelian fixed point. Corrections due to the non-trivial scaling of $`\overline{e}^2(k)`$ remain below $`10\%`$ for $`e_{}^2>6\pi ^2`$ and $`m_\mathrm{H}`$ below 70 GeV, but do grow large as soon as $`e_{}^2`$ is a below $`6\pi ^2`$. We conclude that $`e_{}^26\pi ^2`$ is a good leading order approximation for small Higgs field mass as higher order corrections are small. For $`m_\mathrm{H}=30`$ GeV, we also compared the value for the critical temperature with lattice simulations and found agreement below $`4\%`$. The sensitivity on $`e_{}^2<6\pi ^2`$ for larger Higgs mass, in turn, requires a better determination of the fixed point in this domain. This concerns in particular physical observables like the critical exponents at the endpoint of the line of first order phase transitions.
For generic regulator function the free energy in the type-I regime has been given as an integral (one remaining integration). For the case of a sharp cut-off regulator, we obtained an explicit analytical solution for the free energy, given the non-trivial scale dependence of the Abelian charge. In the present article, we evaluated all relevant quantities for initial conditions obtained from a perturbative dimensional reduction scenario relevant for a high temperature (cosmological) phase transition.
The explicit result for the effective potential can also be of use for applications to the superconducting phase transition, or for the nematic to smectic-A phase transition in certain liquid crystals. The main change would concern the initial potential for the effective $`3d`$ flow of the potential, and the numerical value of the Abelian charge at that scale. These changes affect in particular the ratio $`k_{\mathrm{cr}}/k_{\mathrm{stable}}`$, and therefore the above discussion, as both scales depend in a qualitatively different manner on $`e^2(\mathrm{\Lambda })`$ and $`U_\mathrm{\Lambda }`$.
In addition, we studied the dependence of our results on the coarse graining procedure employed. We have seen that the physical observables do depend only very weakly on the coarse graining. This is encouraging, as a strong dependence would have cast serious doubts on the approximations used. Furthermore, we employed a variety of qualitatively different coarse grainings ranging from the mass-like and other polynomial regulators over exponential ones to the sharp cut-off regulator. Therefore, our result can be seen as an important consistency check of the method. The weak variation w.r.t. the coarse graining which is to be interpreted as an ‘error bar’ for the observables, is smaller or of about the same size as the error expected from higher order operators for the coarse grainings studied. This ‘error bar’ would vanish only if no truncation to the effective action would have to be performed. We also observed an intimate relationship between the truncation of the effective action, and the error bar introduced through the scheme dependence. More precisely, it is observed that the scheme dependence is largest in regions where a similarly large effect due to the neglecting of the scalar fluctuations in the non-convex region of the potential is expected. While this result is not entirely unexpected, a quantitative evidence for it has never been presented before. It would be useful if further quantitative results in this direction could be established. This concerns in particular the cross-dependences between an optimal coarse-graining that minimizes the scheme dependence, and an optimized convergence of systematic truncations and approximations .
An important open question for future work concerns the precise IR behaviour of the Abelian charge. This, of course, is an intrinsic problem of the $`3d`$ theory. As argued, our current understanding is mainly limited due to an insufficient understanding of the field and/or momentum dependence of the Abelian charge. It might be fruitful to consider alternatively a thermal renormalisation group to improve the situation . At the same time, the inclusion of higher order corrections due to scalar fluctuations will also become important – close to the critical points – for a reliable determination of critical exponents and other universal quantities. It would also be interesting to consider the $`SU(2)`$-Higgs theory, where a non-trivial endpoint of the line of first-order phase transitions has been established recently. A field theoretical understanding of this endpoint is still missing, and a derivation of the related critical indices from field theory would be desirable. Again, one expects that the IR behaviour of the gauge coupling, in competition with the scalar fluctuations, is responsible for the existence of the endpoint.
## A RS dependence and threshold functions
The solution of the flow equation (and the related physical observables) can be written as momentum integrals over a measure, which depends on the precise implementation of the coarse graining. We employ the notation of , where a scheme dependent measure has been given (in $`d`$ dimensions) as
$$I_r[f]=\frac{d}{2}_0^{\mathrm{}}𝑑y\frac{r^{}(y)}{(1+r(y))^{1+d/2}}f(y)$$
(A.1)
for momentum-dependent functions $`f(y)`$, where $`y=q^2/k^2`$, and $`q`$ is the loop momenta. As a consequence of the conditions (VI B) on the regularisation function $`r(y)`$ it follows that the momentum measure $`r^{}(y)/(1+r)^{1+d/2}`$ is peaked. The measure is normalised to one,
$$I_r[1]=1.$$
(A.2)
This implies that $`I_r[f]`$ depends on the coarse graining as soon as $`f`$ displays a non-trivial dependence on momenta.
As an example, let’s consider the threshold functions $`\mathrm{}_n^d(\omega )`$, defined as
$$\mathrm{}_n^d(\omega )=\left(\delta _{n,0}+n\right)_0^{\mathrm{}}𝑑y\frac{r^{}(y)y^{1+d/2}}{[y(1+r)+\omega ]^{n+1}}.$$
(A.3)
They are related to the above measure through
$$\mathrm{}_n^d(\omega )=\frac{2}{d}\left(\delta _{n,0}+n\right)I_r\left[\frac{P^{d+2}}{(P^2+\omega )^{n+1}}\right].$$
(A.4)
Here, we also introduced the dimensionless effective (regularised) inverse propagator
$$P^2(y)=y+yr(y).$$
(A.5)
The threshold functions can always be expanded as a Taylor series in powers of $`\omega `$. Let us define the corresponding RS dependent expansion coefficients
$$a_k=I_r[P^k],$$
(A.6)
which are the $`k^{th}`$ moments of $`1/P`$ w.r.t. the measure $`I_r`$. These coefficients appear in the computation of the endpoint of the critical line (63), which is proportional to the coefficient $`a_1`$. For a power-like regulator $`r(y)=y^n`$ \[see eq. (60)\] we find for arbitrary dimension $`d`$
$$a_k=\frac{d}{2}\mathrm{\Gamma }[1+\frac{k}{2n}]\frac{\mathrm{\Gamma }[\frac{d}{2}+\frac{k}{2}(1\frac{1}{n})]}{\mathrm{\Gamma }[1+\frac{d}{2}+\frac{k}{2}]}.$$
(A.7)
A more detailed discussion of these coefficients and a related discussion of the convergence of amplitude expansions and optimised coarse-graining parameters is given in .
## B The fluctuation integral
The fluctuation integral reads
$$\mathrm{\Delta }_k(\overline{\rho })=\frac{1}{2\pi ^2}_k^\mathrm{\Lambda }𝑑\overline{k}_0^{\mathrm{}}𝑑y\frac{\overline{k}^2}{P^2}\frac{2e_{}^2\overline{\rho }r^{}(y)y^{5/2}}{2e_{}^2\overline{\rho }+P^2\overline{k}(1+\overline{k}/k_{\mathrm{cr}})}$$
(B.1)
Note that we have normalised $`\mathrm{\Delta }(0)=0`$ in the above definition. The remaining integrals in (B.1) can be solved in different ways, either first performing the momentum integration or the scale integration. Integrating first w.r.t. $`\overline{k}`$ yields (for the notation see Appendix A)
$$\mathrm{\Delta }_k(\overline{\rho })=I_r\left[𝒰(\overline{\rho },P)\right]$$
(B.2)
where
$`3\pi ^2𝒰(\overline{\rho },P)`$ $`=`$ $`2e_{}^2\overline{\rho }{\displaystyle _k^\mathrm{\Lambda }}𝑑\overline{k}{\displaystyle \frac{P^3\overline{k}^2}{P^2\overline{k}(1+\overline{k}/k_{\mathrm{cr}})+2e_{}^2\overline{\rho }}}`$ (B.3)
$`=`$ $`2e_{}^2\overline{\rho }Pk_{\mathrm{cr}}(k\mathrm{\Lambda })+e_{}^2\overline{\rho }Pk_{\mathrm{cr}}^2\mathrm{ln}\left({\displaystyle \frac{2e_{}^2\overline{\rho }/P^2+k+k^2/k_{\mathrm{cr}}}{2e_{}^2\overline{\rho }/P^2+\mathrm{\Lambda }+\mathrm{\Lambda }^2/k_{\mathrm{cr}}}}\right)`$ (B.5)
$`+2e_{}^2\overline{\rho }Pk_{\mathrm{cr}}(4e_{}^2\overline{\rho }/P^2k_{\mathrm{cr}})G_{k,\mathrm{\Lambda }}(18e_{}^2\overline{\rho }/P^2k_{\mathrm{cr}}),`$
with $`I_r`$ defined in (A.1) and $`P(y)`$ in (A.5). The function $`G(\mathrm{\Omega })`$ reads
$$G_{k,\mathrm{\Lambda }}(\mathrm{\Omega })=\{\begin{array}{cc}\hfill \frac{1}{2\sqrt{\mathrm{\Omega }}}\mathrm{ln}\left(\frac{1+2k/k_{\mathrm{cr}}\sqrt{\mathrm{\Omega }}}{1+2k/k_{\mathrm{cr}}+\sqrt{\mathrm{\Omega }}}\frac{1+2\mathrm{\Lambda }/k_{\mathrm{cr}}+\sqrt{\mathrm{\Omega }}}{1+2\mathrm{\Lambda }/k_{\mathrm{cr}}\sqrt{\mathrm{\Omega }}}\right)\text{for}& \mathrm{\Omega }>0\hfill \\ \hfill \frac{1}{\sqrt{\mathrm{\Omega }}}\left[\mathrm{arctan}\left(\frac{\sqrt{\mathrm{\Omega }}}{1+2k/k_{\mathrm{cr}}}\right)\mathrm{arctan}\left(\frac{\sqrt{\mathrm{\Omega }}}{1+2\mathrm{\Lambda }/k_{\mathrm{cr}}}\right)\right]\text{for}& \mathrm{\Omega }<0\hfill \\ \hfill \frac{2k_{\mathrm{cr}}(k\mathrm{\Lambda })}{(k_{\mathrm{cr}}+2k)(k_{\mathrm{cr}}+2\mathrm{\Lambda })}\text{for}& \mathrm{\Omega }=0.\hfill \end{array}$$
(B.6)
For a sharp cut-off regulator (62), the remaining momentum integration can be performed analytically to give
$`2\pi ^2\mathrm{\Delta }_k^{(s)}`$ $`=`$ $`\frac{1}{3}\mathrm{\Lambda }^3\mathrm{ln}\left(1+{\displaystyle \frac{2e_{}^2\overline{\rho }k_{\mathrm{cr}}}{\mathrm{\Lambda }(\mathrm{\Lambda }+k_{\mathrm{cr}})}}\right)\frac{1}{3}k^3\mathrm{ln}\left(1+{\displaystyle \frac{2e_{}^2\overline{\rho }k_{\mathrm{cr}}}{k(k+k_{\mathrm{cr}})}}\right)\frac{1}{3}k_{\mathrm{cr}}^3\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Lambda }+k_{\mathrm{cr}}}{k+k_{\mathrm{cr}}}}\right)`$ (B.9)
$`+\frac{1}{6}(k_{\mathrm{cr}}^36e_{}^2\overline{\rho }k_{\mathrm{cr}}^2)\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Lambda }(\mathrm{\Lambda }+k_{\mathrm{cr}})+2e_{}^2\overline{\rho }k_{\mathrm{cr}}}{k(k+k_{\mathrm{cr}})+2e_{}^2\overline{\rho }k_{\mathrm{cr}}}}\right)+\frac{4}{3}e_{}^2\overline{\rho }k_{\mathrm{cr}}(\mathrm{\Lambda }k)`$
$`\frac{1}{3}k_{\mathrm{cr}}(16e_{}^4\overline{\rho }^2+10e_{}^2\overline{\rho }k_{\mathrm{cr}}k_{\mathrm{cr}}^2)G_{k,\mathrm{\Lambda }}\left(18e_{}^2\overline{\rho }/k_{\mathrm{cr}}\right).`$
We have normalised $`\mathrm{\Delta }(\overline{\rho })`$ such that $`\mathrm{\Delta }(0)=0`$.
On the other hand, performing first the scheme dependent momentum integration leaves us with the following remaining integrals,
$`\mathrm{\Delta }^{(s)}(\overline{\rho })`$ $`=`$ $`{\displaystyle \frac{1}{2\pi ^2}}{\displaystyle _k^\mathrm{\Lambda }}𝑑\overline{k}\overline{k}^2\mathrm{ln}\left(1+{\displaystyle \frac{2e_{}^2\overline{\rho }}{\overline{k}(1+\overline{k}/k_{\mathrm{cr}})}}\right)`$ (B.10)
$`\mathrm{\Delta }^{(m)}(\overline{\rho })`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _k^\mathrm{\Lambda }}𝑑\overline{k}\overline{k}^2\left(\sqrt{1+{\displaystyle \frac{2e_{}^2\overline{\rho }}{\overline{k}(1+\overline{k}/k_{\mathrm{cr}})}}}1\right)`$ (B.11)
$`\mathrm{\Delta }^{(q)}(\overline{\rho })`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}\pi }}{\displaystyle _k^\mathrm{\Lambda }}𝑑\overline{k}\overline{k}^2\left(1\left(1+{\displaystyle \frac{e_{}^2\overline{\rho }}{\overline{k}(1+\overline{k}/k_{\mathrm{cr}})}}\right)^{1/2}\right)`$ (B.12)
Here, the indices refer to the sharp ($`s`$), the mass-like ($`m`$) and the quartic ($`q`$) cut-off function, as defined in Sect. VI.
## C Including scalar fluctuations
In order to obtain an estimate of the effect of the scalar fluctuations we will solve eq. (9) with $`U_k`$ on the r.h.s. replaced by $`U_\mathrm{\Lambda }`$, for a sharp cut-off regulator. The flow equation becomes
$$4\pi ^2\frac{dU_k(\overline{\rho })}{k^2dk}=\mathrm{}_0^3\left(\frac{m_{1,\mathrm{\Lambda }}^2(\overline{\rho })}{k^2}\right)+\mathrm{}_0^3\left(\frac{m_{2,\mathrm{\Lambda }}^2(\overline{\rho })}{k^2}\right)+2\mathrm{}_0^3\left(\frac{2\overline{e}_3^2(k)\overline{\rho }}{k^2}\right).$$
(C.1)
with the masses $`m_i^2`$ given through
$$m_1^2(\overline{\rho })=m_\mathrm{R}^2+\overline{\lambda }_\mathrm{R}\overline{\rho },m_2^2(\overline{\rho })=m_\mathrm{R}^2+3\overline{\lambda }_\mathrm{R}\overline{\rho }.$$
(C.2)
We introduce the functions
$`K(m^2)`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^2}}{\displaystyle _\mathrm{\Lambda }^k}𝑑yy^2\mathrm{ln}\left(1+{\displaystyle \frac{m^2}{y^2}}\right)`$ (C.3)
$`=`$ $`{\displaystyle \frac{1}{12\pi ^2}}[2m^2k2m^3\mathrm{arctan}\left({\displaystyle \frac{k}{m}}\right)+k^3\mathrm{ln}(1+{\displaystyle \frac{m^2}{k^2}})(k\mathrm{\Lambda })]`$ (C.4)
$`J_i(\overline{\rho })`$ $`=`$ $`K[m_i^2(\overline{\rho })]K[m_i^2(0)].`$ (C.5)
The solution to the flow (C.1) then obtains, using also $`\mathrm{\Delta }^{(s)}`$ from (B.9), as
$$U_k(\overline{\rho })=U_\mathrm{\Lambda }(\overline{\rho })+\mathrm{\Delta }^{(s)}(\overline{\rho })+J_1(\overline{\rho })+J_2(\overline{\rho }).$$
(C.6)
The effect of the additional terms on the shape of the critical potential is about a few percent, increasing towards higher values for $`m_\mathrm{H}`$.
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# Invariant Measures and Orbit Closures on Homogeneous Spaces for Actions of Subgroups Generated by Unipotent Elements
## 1 Introduction
In \[Ratner:measure, Ratner:distribution\] Ratner showed the validity of Raghunathan’s conjecture describing orbit closures for actions of unipotent subgroups on homogeneous spaces of Lie groups, and its analogous conjecture, due to Dani \[Dani:invent2\], describing ergodic invariant measures for such actions. Earlier in \[Margulis:varna, Margulis:selberg\] Margulis had conjectured that the conclusions of the orbit closure and the ergodic invariant measure conjectures should hold also for the actions of subgroups generated by unipotent elements, as compared to the subgroups themselves being unipotent. In fact for actions of connected subgroups generated by unipotent elements, this conjecture was also verified to be true in Ratner’s above mentioned papers. Using Ratner’s theorems for actions of unipotent one-parameter subgroups, in this article we show the validity of the generalized conjecture. This also answers a question raised by Ratner in \[Ratner:ICM, End of Section 4\].
Notation. Let $`G`$ be a Lie group, $`𝔤`$ its Lie algebra, and $`\mathrm{Ad}_G:G\mathrm{GL}(𝔤)`$ denote the Adjoint representation of $`G`$ on $`𝔤`$. An element $`uG`$ is called $`\mathrm{Ad}_G`$-unipotent, if $`\mathrm{Ad}_G(u)`$ is a unipotent linear transformation. A subgroup of $`G`$ consisting of $`\mathrm{Ad}_G`$-unipotent elements is called an $`\mathrm{Ad}_G`$-unipotent subgroup.
Let $`𝒰`$ denote the subgroup generated by a subset $`𝒰`$ in $`G`$. Let $`\mathrm{Zcl}(X)`$ denote the Zariski closure of a subset $`X`$ in $`\mathrm{GL}(𝔤)`$. For a subgroup $`F`$ of $`G`$, let $`F^0`$ denote the connected component of $`F`$ containing the identity element.
For a Borel measure $`\mu `$ on a second countable topological space, we denote by $`\mathrm{supp}(\mu )`$ the closed subset which is the complement of the union of all open sets with zero $`\mu `$-measure.
###### Theorem 1.1
Let $`G`$ be a Lie group and $`\mathrm{\Gamma }`$ a closed subgroup of $`G`$. Let $`W`$ be a subgroup of $`G`$ and $`𝒰W`$ such that $`𝒰`$ consists of $`\mathrm{Ad}_G`$-unipotent elements and $`\mathrm{Ad}_G(W)\mathrm{Zcl}(\mathrm{Ad}_G(𝒰))`$. Let $`\mu `$ be a finite $`W`$-invariant $`W`$-ergodic Borel measure on $`G/\mathrm{\Gamma }`$. Then there exists a closed subgroup $`H`$ of $`G`$ containing $`W`$ such that $`\mu `$ is $`H`$-invariant and $`\mathrm{supp}(\mu )`$ is a closed $`H`$-orbit.
A Borel measure on a locally compact second countable topological space is called locally finite, if it is finite on compact sets.
###### Theorem 1.2
Let $`G`$, $`\mathrm{\Gamma }`$, and $`W`$ be as in Theorem 1.1. Suppose that $`G/\mathrm{\Gamma }`$ has a finite $`G`$-invariant measure. Let $`\mu `$ be a locally finite $`W`$-invariant $`W`$-ergodic measure on $`G/\mathrm{\Gamma }`$. Then there exists a closed subgroup $`H`$ of $`G`$ containing $`W`$ such that $`\mu `$ is $`H`$-invariant and $`\mathrm{supp}(\mu )`$ is a closed $`H`$-orbit.
###### Theorem 1.3
Let $`G`$, $`\mathrm{\Gamma }`$ and $`W`$ be as in Theorem 1.1. Suppose that $`G/\mathrm{\Gamma }`$ has a finite $`G`$-invariant measure. Then for any $`xG/\mathrm{\Gamma }`$, there exists a closed subgroup $`F`$ of $`G`$ containing $`W`$ such that
$$\overline{Wx}=Fx.$$
Moreover, $`F^0x`$ has a finite $`F^0`$-invariant measure (cf. Conjectures 1.1 and 1 below). Also the action of $`W`$ is ergodic with respect to a locally finite $`F`$-invariant measure on $`Fx`$.
We may note that Theorem 1.1 and Theorem 1.3 have already been proved in the above mentioned papers of Ratner in the following special case: $`G`$ is connected, $`W`$ is of the form $`W=_{i=1}^{\mathrm{}}w_iW^0`$, where $`w_i`$ is $`\mathrm{Ad}_G`$-unipotent, $`i=1,2,\mathrm{}`$, $`W/W^0`$ is finitely generated, and $`W^0`$ is generated by one-parameter $`\mathrm{Ad}_G`$-unipotent subgroups contained in $`W^0`$. In the case when $`G`$ is not connected and $`W`$ is a nilpotent $`\mathrm{Ad}_G`$-unipotent subgroup of $`G`$, Theorem 1.1 was proved by Witte \[Witte:quotients, Theorem 1.2\]. In the case when $`G`$ is connected and $`W`$ is a $`\mathrm{Ad}_G`$-unipotent subgroup, it was shown by Dani \[Dani:rk=1, Theorem 4.3\] that if $`G/\mathrm{\Gamma }`$ has a finite invariant measure then any locally finite $`W`$-invariant $`W`$-ergodic measure is finite. In \[Margulis:varna, Remarks 3.12\], Margulis observed that the same holds for connected $`W`$. Thus for connected $`W`$, Theorem 1.2 reduces to Theorem 1.1, which was proved by Ratner (for connected $`W`$).
The following result is deduced from that above results using the ‘suspension techniques’ (cf. Witte \[Witte:quotients, Corollary 5.8\]).
###### Corollary 1.4
Let $`G`$ and $`W`$ be as in any one of the theorems stated above. Assume that $`W`$ is closed and let $`\mathrm{\Lambda }`$ be a closed subgroup of $`W`$ such that $`W/\mathrm{\Lambda }`$ has a finite $`W`$-invariant measure. Then all the theorems stated above are true for $`\mathrm{\Lambda }`$ in place of $`W`$.
Further, if $`G/\mathrm{\Gamma }`$ admits a finite $`G`$-invariant measure and $`W`$ is connected, then we have the following additional information:
1. Any locally finite $`\mathrm{\Lambda }`$-invariant $`\mathrm{\Lambda }`$-ergodic measure on $`G/\mathrm{\Gamma }`$ is finite.
2. For $`xG/\mathrm{\Gamma }`$, if $`\overline{\mathrm{\Lambda }x}=Fx`$ for a closed subgroup $`F`$ of $`G`$ then $`Fx`$ has a finite $`F`$-invariant measure.
From this corollary, we deduce the following.
###### Corollary 1.5
Let $`G`$ be a connected semisimple Lie group without nontrivial compact factors. Let $`\mathrm{\Gamma }`$ and $`\mathrm{\Lambda }`$ be lattices in $`G`$ such that at least one of them is irreducible in $`G`$; (see \[Raghunathan:book, Sect. 5.20\] for definition). Then either $`\mathrm{\Lambda }\mathrm{\Gamma }`$ is dense in $`G`$ or $`\mathrm{\Lambda }\mathrm{\Gamma }`$ is a subgroup of finite index in $`\mathrm{\Gamma }`$, as well as $`\mathrm{\Lambda }`$.
In view of the above results we may ask if the following is true.
###### Conjecture 1.1
Let the notation be as in Theorem 1.1. Suppose further that $`G/\mathrm{\Gamma }`$ has a finite $`G`$-invariant measure. Then the following statements hold:
1. Any locally finite $`W`$-invariant $`W`$-ergodic measure on $`G/\mathrm{\Gamma }`$ is finite.
2. For any $`xG/\mathrm{\Gamma }`$, if $`\overline{Wx}=Fx`$ for a closed subgroup $`F`$ of $`G`$, then $`Fx`$ has a finite $`F`$-invariant measure.
3. The closure of any $`W`$-orbit has finitely many connected components.
Note that by the above stated theorems and by Hedlund’s lemma 2.1, the three statements in the above conjecture are equivalent.
###### Remark 1.6
If the above conjecture is valid for the diagonal action of $`W`$ on $`W/\mathrm{\Lambda }\times G/\mathrm{\Gamma }`$, then it holds for the action of $`\mathrm{\Lambda }`$ on $`G/\mathrm{\Gamma }`$, where $`W`$ and $`\mathrm{\Lambda }`$ are as in corollary 1.4.
It seems that the generalized Raghunathan conjecture due to Margulis already includes Conjecture 1.1. Using some standard arguments, as in proof of Theorem LABEL:nimish:thm:closed:finite, one can reduce this conjecture to the case of $`G`$ being a semisimple group with no nontrivial compact factors and trivial center. Then one can express $`G`$ as a product of semisimple subgroups each intersecting $`\mathrm{\Gamma }`$ in an irreducible lattice. Using the structure of the cusps in the quotient of the $``$-rank one factors, one can take care of those factors. Thus the conjecture remains to be proved for higher rank semisimple groups $`G`$. We use the arithmeticity theorem of Margulis, and reduce the conjecture to its following typical case.
###### Conjecture 1.2
<sup>1</sup><sup>1</sup>1Recently Alex Eskin and G. A. Margulis informed the author that they can prove this conjecture. Let $`G=\mathrm{SL}_n()`$, $`\mathrm{\Gamma }=\mathrm{SL}_n()`$, and $`W\mathrm{SL}_n()`$ a closed subgroup of $`G`$ such that $`W`$ is contained in the Zariski closure of a subgroup generated by $`\mathrm{Ad}_G`$-unipotent elements of $`W`$. If $`W\mathrm{\Gamma }`$ is discrete, then $`W\mathrm{\Gamma }`$ is of finite index in $`W`$.
Finally, a challenging question is to describe invariant measures and orbit closures for actions of subgroups $`H`$ whose Zariski closure is generated by unipotent elements, which are not necessarily contained in $`H`$. For example, let $`H`$ be a Zariski dense subgroup of $`\mathrm{SL}_2()`$ not containing any unipotent elements, and consider the action of $`H`$ on $`\mathrm{SL}_2()/\mathrm{SL}_2()`$.
## 2 Preliminary Results
In this section we recall some standard results or their modifications about orbit closures on homogeneous spaces, and Zariski density of certain discrete subgroups as in Borel’s density theorem.
###### Lemma 2.1 (Hedlund’s Lemma)
Let $`X`$ be a second countable topological space and $`W`$ be a group of homeomorphisms of $`X`$. Let $`\mu `$ be a $`W`$-invariant $`W`$-ergodic Borel measure on $`X`$. Then $`\overline{Wx}=\mathrm{supp}\mu `$ for $`\mu `$-almost all $`xX`$.
###### Lemma 2.2
Let $`G`$ be a locally compact second countable group, $`\mathrm{\Gamma }`$ a discrete subgroup of $`G`$, and $`\pi :GG/\mathrm{\Gamma }`$ the natural quotient map. Let $`F`$ be a Borel measurable subgroup of $`G`$. Suppose there exists a locally finite $`F`$-invariant Borel measure concentrated on $`\pi (F)`$. Then $`F`$ is closed, $`\pi (F)`$ is closed, and $`\mathrm{supp}(\mu )=\pi (F)`$.
###### Proof
Since $`\mu `$ is locally finite, by dominated convergence theorem (see \[Ratner:solvable, Proposition 1.4\]), $`\mu `$ is invariant under the closure of $`F`$ in $`G`$, say $`H`$. We have a natural inclusion $`H/H\mathrm{\Gamma }G/\mathrm{\Gamma }`$, which is $`H`$-equivariant. Since $`\mu `$ is concentrated on $`\pi (F)\pi (H)`$, we can treat $`\mu `$ as a locally finite $`H`$-invariant Borel measure on $`H/H\mathrm{\Gamma }`$. Since $`\mu `$ is concentrated on an orbit of $`F`$, we conclude that a Haar measure on $`H`$ is strictly positive on $`F`$. Since $`F=FF^^1`$, we have that $`F`$ is an open, and hence a closed subgroup of $`H`$. Thus $`F=H`$.
Now since $`F`$ is closed, the result follows from the proof of \[Raghunathan:book, Theorems 1.12-1.13\] (cf. \[Ratner:measure, Proposition 1.4\]). Although these references assume that $`\mu `$ is finite, the local nature of the conclusion requires only the assumption that $`\mu `$ is locally finite. ∎
###### Lemma 2.3
Let $`G`$ be a Lie group and $`\mathrm{\Gamma }`$ be a discrete subgroup of $`G`$. Let $`\pi :GG/\mathrm{\Gamma }`$ be the quotient map. Let $`F`$ be a subgroup of $`G`$ such that
$$\mathrm{Ad}_G(F)\mathrm{Zcl}(\mathrm{Ad}_G(F\mathrm{\Gamma })).$$
Then $`\pi (Z_G(F))`$ is closed.
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# Screening of a charged particle by multivalent counterions in salty water: Strong charge inversion
## I Introduction
Charge inversion is a phenomenon in which a charged particle (a macroion) strongly binds so many counterions in a water solution that its net charge changes sign. As shown below the binding energy of a counterion with large charge $`Z`$ is larger than $`k_BT`$, so that this net charge is easily observable; for instance, it is the net charge that determines linear transport properties, such as particle drift in a weak field electrophoresis. Charge inversion is possible for a variety of macroions, ranging from the charged surface of mica or other solids to charged lipid membranes, DNA or actin. Multivalent metallic ions, small colloidal particles, charged micelles, short or long polyelectrolytes can play the role of multivalent counterions. Recently, charge inversion has attracted significant attention .
Charge inversion is of special interest for the delivery of genes to the living cell for the purpose of the gene therapy. The problem is that both bare DNA and a cell surface are negatively charged and repel each other, so that DNA does not approach the cell surface. The goal is to screen DNA in such a way that the resulting complex is positive . Multivalent counterions can be used for this purpose. The charge inversion depends on the surface charge density, so the cell surface charge can still be negative when DNA charge is inverted.
Charge inversion can be also thought of as an over-screening. Indeed, the simplest screening atmosphere, familiar from linear Debye-Hückel theory, compensates at any finite distance only a part of the macroion charge. It can be proven that this property holds also in non-linear Poisson-Boltzmann (PB) theory. The statement that the net charge preserves sign of the bare charge agrees with the common sense. One can think that this statement is even more universal than results of PB equation. It was shown , however, that this presumption of common sense fails for screening by $`Z`$-valent counterions ($`Z`$-ions) with large $`Z`$, such as charged colloidal particles, micelles or rigid polyelectrolytes, because there are strong repulsive correlations between them when they are bound to the surface of a macroion. As a result, $`Z`$-ions form strongly correlated liquid with properties resembling a Wigner crystal (WC) at the macroion surface. The negative chemical potential of this liquid leads to an additional ”correlation ” attraction of $`Z`$-ions to the surface. This effect is beyond the mean field PB theory, and charge inversion is its most spectacular manifestation.
Let us demonstrate fundamental role of lateral correlations between $`Z`$-ions for a simple model. Imagine a hard-core sphere with radius $`b`$ and with negative charge $`Q`$ screened by two spherical positive $`Z`$-ions with radius $`a`$. One can see that if Coulomb repulsion between $`Z`$-ions is much larger than $`k_BT`$ they are situated on opposite sides of the negative sphere (Fig. 1a).
If $`Q>Ze/2`$, each $`Z`$-ion is bound because the energy required to remove it to infinity $`QZe/(a+b)Z^2e^2/2(a+b)`$ is positive. Thus, the charge of the whole complex $`Q^{}=Q+2Ze`$ can be positive. For example, $`Q^{}=3Ze/2=3Q`$ at $`Q=Ze/2`$. This example demonstrates the possibility of an almost 300% charge inversion. It is obviously a result of the correlation between $`Z`$-ions which avoid each other and reside on opposite sides of the negative charge. On the other hand, the description of screening of the central sphere in the PB approximation smears the positive charge, as shown on Fig. 1b and does not lead to the charge inversion. Indeed, in this case charge accumulates in spherically symmetric screening atmosphere only until the point of neutrality at which electric field reverses its sign and attraction is replaced by repulsion.
Weak charge inversion can be also obtained as a trivial result of $`Z`$-ions discreteness without correlations. Indeed, discrete $`Z`$-ions can over-screen by a fraction of the ”charge quantum” $`Ze`$. For example, if central charge $`Q=Ze/2`$ binds one $`Z`$-ion, the net charge of the complex is $`Q^{}=Ze/2`$. This charge is, however, three times smaller than the charge $`3Ze/2`$ which we obtained above for screening of the same charge $`Ze/2`$ by two correlated $`Z`$-ions, so that for the same $`Q`$ and $`Z`$ correlations lead to stronger charge inversion.
Difference between charge inversion, obtained with and without correlations becomes dramatic for a large sphere with a macroscopic charge $`QZe`$. In this case, discreteness by itself can lead to inverted charge limited by $`Ze`$. On the other hand, it was predicted and confirmed by numerical simulations that due to correlation between $`Z`$-ions which leads to their WC-like short range order on the surface of the sphere, the net inverted charge can reach
$$Q^{}=0.84\sqrt{QZe},$$
(1)
i. e. can be much larger than the charge quantum $`Ze`$. This charge is still smaller than $`Q`$ because of limitations imposed by the very large charging energy of the macroscopic net charge.
In this paper, we consider systems in which inverted charge can be even larger than what Eq. (1) predicts. Specifically, we consider the problem of screening by $`Z`$-ions in the presence of monovalent salt, such as NaCl, in solution. This is a more practical situation than the salt-free one considered in Ref. . Monovalent salt screens long range Coulomb interactions stronger than short range lateral correlations between adsorbed $`Z`$-ions. Therefore, screening diminishes the charging energy of the macroion much stronger than the correlation energy of $`Z`$-ions. As a results, the inverted charge $`Q^{}`$ becomes larger than that predicted by Eq. (1) and scales linearly with $`Q`$. The amount of charge inversion at strong screening is limited only by the fact that the binding energy of $`Z`$-ions becomes eventually lower than $`k_BT`$, in which case it is no longer meaningful to speak about binding or adsorption. Nevertheless, remaining within the strong binding regime, we demonstrate on many examples throughout this work, that the inverted charge, in terms of its absolute value, can be larger than the original bare charge, sometimes even by a factor up to 3. We call this phenomenon strong or giant charge inversion and its prediction and theory are the main results of our paper (A brief preliminary version of this paper is given in Ref. ).
Since, in the presence of a sufficient concentration of salt, the macroion is screened at the distance smaller than its size, the macroion can be thought of as an overscreened surface, with inverted charge $`Q^{}`$ proportional to the surface area. In this sense, overall shape of the macroion and its surface is irrelevant, at least to a first approximation. Therefore, we consider screening of a planar macroion surface with a negative surface charge density $`\sigma `$ by finite concentration, $`N`$, of positive $`Z`$-ions, and concentration $`ZN`$ of neutralizing monovalent coions, and a large concentration $`N_1`$ of a monovalent salt. Correspondingly, we assume that all interactions are screened with Debye-Hückel screening length $`r_s=\left(8\pi l_BN_1\right)^{1/2}`$, where $`l_B=e^2/(Dk_BT)`$ is the Bjerrum length, $`e`$ is the charge of a proton, $`D80`$ is the dielectric constant of water. At small enough $`r_s`$, the method of a new boundary condition for the PB equation suggested in Ref. becomes less convenient and in this paper we develop more universal and direct theoretical approach to charge inversion problem.
Our goal is to calculate the two-dimensional concentration $`n`$ of $`Z`$-ions at the plane as a function of $`r_s`$ and $`N`$. In other words, we want to find the net charge density of the plane
$$\sigma ^{}=\sigma +Zen.$$
(2)
In particular, we are interested in the maximal value of the ”inversion ratio”, $`\sigma ^{}/\sigma `$, which can be reached at large enough $`N`$. The subtle physical meaning of $`\sigma ^{}`$ should be clearly explained. Indeed, the entire system, macroion plus overcharging $`Z`$-ions, is, of course, neutralized by the monovalent ions. One can ask then, what is the meaning of charge inversion? In other words, what is the justification of definition of Eq. (2) which disregards monovalent ions?
To answer we note that under realistic conditions, every $`Z`$-ion, when on the macroion surface, is attached to the macroion with energy well in excess of $`k_BT`$. At the same time, monovalent ions, maintaining electroneutrality over the distances of order $`r_s`$, interact with the macroion with energies less than $`k_BT`$ each. It is this very distinction that led us to define the net charge of the macroion including adsorbed $`Z`$-ions and excluding monovalent ions. Our definition is physically justified, it has direct experimental relevance. Indeed, it is conceivable that the strongly adsorbed $`Z`$-ions can withstand perturbation caused by the atomic force microscopy (AFM) experiment, while the neutralizing atmosphere of monovalent ions cannot. Therefore, one can, at least in principle, count the adsorbed $`Z`$-ions, thus directly measuring $`\sigma ^{}`$. To give a practical example, when $`Z`$-ions are the DNA chains, one can realistically measure the distance between neighboring DNAs adsorbed on the surface. In most cases, similar logic applies to an electrophoresis experiment in a weak external electric field such that the current is linear in applied field. Sufficiently weak field does not affect the strong (above $`k_BT`$) attachment of $`Z`$-ions to the macroion. In other words, macroion coated with bound $`Z`$-ions drifts in the field as a single body. On the other hand, the surrounding atmosphere of monovalent ions, smeared over the distances about $`r_s`$, drifts with respect to the macroion. Presenting linear electrophoretic mobility of a macroion as a ratio of effective charge to effective friction, we conclude that only $`Z`$-ions contribute to the former, while monovalent ions contribute only to the latter. In particular, and most importantly, the sign of the effect - in which direction the macroion moves, along the field or against the field - is determined by the net charge $`\sigma ^{}`$ which, once again, includes $`Z`$-ions and does not include monovalent ones. Furthermore, for a macroion with simple (e.g., spherical) shape, the absolute value of the net macroion charge can be also found using the mobility measurements and the standard theory of friction in electrolytes . This logic fails only for the regime which we call strongly non-linear. In this regime, majority of monovalent ions form a bound Gouy-Chapman atmosphere of the inverted charge, and, while surface charge as counted by AFM remains equal $`\sigma ^{}`$, the electrophoretic measurement yields universal value $`e/2\pi l_Br_s`$, which is inverted but is smaller than $`\sigma ^{}`$. For a macroion of the size smaller than $`r_s`$, its size determines the maximum inverted charge.
Now, as we have formulated major goal of the paper, let us describe briefly its structure and main results. In Sec. II - IV we consider screening of a charged surface by compact $`Z`$-ions such as charged colloidal particles, micelles or short polyelectrolytes, which can be modeled as a sphere with radius $`a`$. We call such $`Z`$-ions ”spherical”. Spherical ions form correlated liquid with properties similar to two-dimensional WC (Fig. 2).
In Sec. II we begin with screening of the simplest macroion which is a thin charged sheet immersed in water solution (Fig. 3a). This lets us to postpone the complication related to image potential which appears for a more realistic macroion which is a thick insulator charged at the surface (Fig. 3b). We calculate analytically the dependence of the inversion ratio, $`\sigma ^{}/\sigma `$, on $`r_s`$ in two limiting cases $`r_sR_0`$ and $`r_sR_0`$, where $`R_0=(\pi \sigma /Ze)^{1/2}`$ is the radius of a Wigner-Seitz cell at the neutral point $`n=\sigma /Ze`$ (we approximate the hexagon by a disk). We find that at $`r_sR_0`$
$$\sigma ^{}/\sigma =0.83(R_0/r_s)=0.83\zeta ^{1/2},(\zeta 1)$$
(3)
where $`\zeta =Ze/\pi \sigma r_s^2=(R_0/r_s)^2`$. At $`r_sR_0`$
$$\frac{\sigma ^{}}{\sigma }=\frac{2\pi \zeta }{\sqrt{3}\mathrm{ln}^2\zeta },(\zeta 1).$$
(4)
Thus $`\sigma ^{}/\sigma `$ grows with decreasing $`r_s`$ and can become larger than 100%. We also present numerical calculation of the full dependence of the inversion ratio on $`\zeta `$.
In Sec. III we discuss effects related to finite size of $`Z`$-ion. It is well known that monovalent ions can condense on the surface of a small and strongly charged spherical $`Z`$-ion. As a result, instead of the bare charge of $`Z`$-ions in Eqs. (3) and (4) one should use the net charge of $`Z`$-ions, which is substantially smaller. Thus, condensation puts a limit for the inversion ratio. The net charge grows with the radius $`a`$ of the $`Z`$-ion. Therefore, we study in this section the case when $`r_saR_0`$ and showed that the largest inversion ratio for spherical ions can reach a few hundred percent.
Sec. IV is devoted to more realistic macroions which have a thick insulating body with dielectric constant much smaller than that of water. In this case each $`Z`$-ion has an image charge of the same sign and magnitude. Image charge repels $`Z`$-ion and pushes WC away from the surface. In this case charge inversion is studied numerically in all the range of $`r_s`$ or $`\zeta `$. The result turns out to be remarkably simple: at $`\zeta <100`$, the inversion ratio is twice smaller than for the case of the charged sheet immersed in water. A simple interpretation of this result will be given in Sec. IV.
In Sec. V and VI we study adsorption of long rod-like $`Z`$-ions with negative linear charge bare density $`\eta _0`$ on a surface with a positive charge density $`\sigma `$. (We changed the signs of both surface and $`Z`$-ion charges to be closer to the practical case when DNA double helices are adsorbed on a positive surface.) Due to the strong lateral repulsion, charged rods tend to be parallel to each other and have a short range order of an one-dimensional WC (Fig. 4). In the Ref. one can find beautiful atomic force microscopy pictures of almost perfect one-dimensional WC of DNA double helices on a positive membrane. The adsorption of another rigid polyelectrolyte, PDDA, was studied in Ref. . Here we concentrate on the case of DNA.
It is well known that for DNA, the bare charge density, $`\eta _0`$ is four times larger than the critical density $`\eta _c=Dk_BT/e`$ of the Onsager-Manning condensation . According to the solution of nonlinear PB equation, most of the bare charge of an isolated DNA is compensated by positive monovalent ions residing at its surface so that the net charge of DNA is equal to $`\eta _c`$. The net charge of DNA adsorbed on a charged surface may differ from $`\eta _c`$ due to the repulsion of positive monovalent ions condensed on DNA from the charged surface. We, however, show that in the case of strong screening, $`r_sA_0`$ ($`A_0=\eta _c/\sigma `$), the potential of the surface is so weak that the net charge, $`\eta `$, of each adsorbed DNA is still equal to $`\eta _c`$. Simultaneously, at $`r_sA_0`$ the Debye-Hückel approximation can be used to describe screening of the charged surface by monovalent salt. In Sec.V, these simplifications are used to study the case of strong screening. We show that the competition between the attraction of DNA to the surface and the repulsion of the neighbouring DNAs results in the negative net surface charge density $`\sigma ^{}`$ and the charge inversion ratio, similar to Eq. (4):
$$\frac{\sigma ^{}}{\sigma }=\frac{\eta _c/\sigma r_s}{\mathrm{ln}(\eta _c/\sigma r_s)},(\eta _c\sigma /r_s1)$$
(5)
Thus the inversion ratio grows with decreasing $`r_s`$ as in the spherical $`Z`$-ion case. At small enough $`r_s`$ and $`\sigma `$, the inversion ratio can reach 400%. This is larger than for spherical ions because in this case, due to the large persistence length of DNA, the correlation energy remains large and WC-like short range order is preserved at smaller $`\sigma r_s`$. An expression similar Eq. (5) has been recently derived for the case of polyelectrolyte with small absolute value of the linear charge density, $`\eta _0\eta _c`$, and strong screening ($`r_sA`$) when screening of both the charged surface and the polyelectrolyte can be treated in Debye-Hückel approximation . The result of Ref. can be obtained if we replace the net charge $`\eta _c`$ by the bare charge $`\eta _0`$ in Eq. (5) .
In Sec. VI we study the adsorption of DNA rods in the case of weak screening by monovalent salt, $`r_sA_0`$. In this case, screening of the overcharged plane by monovalent salt becomes strongly nonlinear, with the Gouy-Chapman screening length $`\lambda =e/(\pi l_B\sigma ^{})`$ much smaller than $`r_s`$. Simultaneously, the charge of macroion repels monovalent coions so that some of them are released from DNA. As a result the absolute value of the net linear charge density of a rod, $`\eta `$, is larger than $`\eta _c`$. We derived two nonlinear equations for unknown $`\sigma ^{}`$ and $`\eta `$. Their solution at $`r_sA_0`$ gives:
$$\frac{\sigma ^{}}{\sigma }=\frac{\eta _c}{\pi a\sigma }\mathrm{exp}\left(\sqrt{\mathrm{ln}\frac{r_s}{a}\mathrm{ln}\frac{A_0}{2\pi a}}\right),$$
(6)
$$\eta =\eta _c\sqrt{\frac{\mathrm{ln}(r_s/a)}{\mathrm{ln}(A_0/2\pi a)}}.$$
(7)
At $`r_sA_0`$ we get $`\eta \eta _c`$, $`\lambda r_s`$ and $`\sigma ^{}/\sigma 1`$ so that Eq. (6) matches the strong screening result of Eq. (5). Since $`\eta `$ can not be smaller than $`\eta _c`$, the fact that $`\eta \eta _c`$ already at $`r_sA_0`$ proves that at $`r_sA_0`$, indeed, $`\eta \eta _c`$
In Sec. VII we return to spherical $`Z`$-ions and derive the system of nonlinear equations which is similar to one derived in Sec. VI for rod-like ones. This system lets us justify the use of Debye-Hückel approximation for screening of overcharged surface ( Sec. II) at $`r_s`$ smaller than $`r_m`$, where $`r_m=a\mathrm{exp}(R_0/1.65a)`$ is an exponentially large length. We show that even at $`r_sr_m`$ nonlinear equations lead only to a small correction to the power of $`r_s`$ in Eq. (3).
In Sec. I-VII we assume that the surface charges of a macroion are frozen and can not move. In Sec. VIII we explore the role of the mobility of these charges. Surface charge can be mobile, for example, on charged liquid membrane where hydrophilic heads can move along the surface. If a membrane surface has heads with two different charges, for example, 0 and -$`e`$, the negative ones can replace the neutral ones near the positive $`Z`$-ion, thus accumulating around it and binding it stronger to the surface. We show that this effect enhances charge inversion substantially. We conclude in Sec. IX.
## II Screening of charged sheet by spherical $`Z`$-ions
Assume that a plane with the charge density $`\sigma `$ is immersed in water (Fig. 3a) and is covered by $`Z`$-ions with two-dimensional concentration $`n`$. Integrating out all the monovalent ion degrees of freedom, or, equivalently, considering all interactions screened at the distance $`r_s`$, we can write down the free energy per unit area in the form
$$=\pi \sigma ^2r_s/D2\pi \sigma r_sZen/D+F_{ZZ}+F_{id},$$
(8)
where the four terms are responsible, respectively, for the self interaction of the charged plane, for the interaction between $`Z`$-ions and the plane, for pair interactions between $`Z`$-ions and for the entropy of ideal two-dimensional gas of $`Z`$-ions. Using Eq. (2) one can rewrite Eq. (8) as
$$=\pi (\sigma ^{})^2r_s/D+F_{OCP},$$
(9)
where $`F_{OCP}=F_c+F_{id}`$ is the free energy of the same system of $`Z`$-ions residing on a neutralizing background with surface charge density $`Zen`$, which is conventionally referred to as one component plasma (OCP), and
$$F_c=\pi (Zen)^2r_s/D+F_{ZZ}$$
(10)
is the correlation part of $`F_{OCP}`$. The transformation from Eq. (8) to Eq. (9) can be simply interpreted as the addition of uniform charge densities $`\sigma ^{}`$ and $`\sigma ^{}`$ to the plane. The first addition makes a neutral OCP on the plane. The second addition creates two planar capacitors with negative charges on both sides of the plane which screen the inverted charge of the plane at the distance $`r_s`$ (Fig. 3a). The first term of Eq. (9) is nothing but the energy of these two capacitors. There is no cross term corresponding to the interactions between the OCP and the capacitors because each planar capacitor creates a constant potential, $`\psi (0)=2\pi \sigma ^{}r_s/D`$, at the neutral OCP.
Using Eq. (10), the electrochemical potential of $`Z`$-ions at the plane can be written as $`\mu =Ze\psi (0)+\mu _{id}+\mu _c`$, where $`\mu _{id}`$ and $`\mu _c=F_c/n`$ are the ideal and the correlation parts of the chemical potential of OCP. In equilibrium, $`\mu `$ is equal to the chemical potential, $`\mu _b`$, of the ideal bulk solution, because in the bulk electrostatic potential $`\psi =0`$. Using Eq. (9), we have:
$$2\pi \sigma ^{}r_sZe/D=\mu _c+(\mu _b\mu _{id}).$$
(11)
As we show below, in most practical cases the correlation effect is rather strong, so that $`\mu _c`$ is negative and $`|\mu _c|k_BT`$. Furthermore, strong correlations imply that short range order of $`Z`$-ions on the surface should be similar to that of triangular Wigner crystal (WC) since it delivers the lowest energy to OCP. Thus one can substitute the chemical potential of Wigner crystal, $`\mu _{WC}`$, for $`\mu _c`$. One can also write the difference of ideal parts of the bulk and the surface chemical potentials of $`Z`$-ions as
$$\mu _b\mu _{id}=k_BT\mathrm{ln}(N_s/N),$$
(12)
where $`N_sn/a`$ is the bulk concentration of $`Z`$-ions at the plane. Then Eq. (11) can be rewritten as
$$2\pi \sigma ^{}r_sZe/D=k_BT\mathrm{ln}(N/N_0),$$
(13)
where $`N_0=N_s\mathrm{exp}(|\mu _{WC}|/k_BT)`$ is the concentration of $`Z`$-ions in the solution next to the charged plane. which plays the role of boundary condition for $`N(x)`$ when $`x0`$ . It is clear that when $`N>N_0`$, the net charge density $`\sigma ^{}`$ is positive, i.e. has the sign opposite to the bare charge density $`\sigma `$. The concentration $`N_0`$ is very small because $`|\mu _{WC}|/k_BT1`$. Therefore, it is easy to achieve charge inversion. According to Eq. (12) at large enough $`N`$ one can neglect second term of the right side of Eq. (11). This gives for the maximal inverted charge density
$$\sigma ^{}=\frac{D}{2\pi r_s}\frac{\left|\mu _{WC}\right|}{Ze}.$$
(14)
Eq. (14) has a very simple meaning: $`|\mu _{WC}|/Ze`$ is the ”correlation” voltage which charges two above mentioned parallel capacitors with ”distance between plates” $`r_s`$ and total capacitance per unit area $`D/(2\pi r_s)`$.
To calculate the correlation voltage $`\left|\mu _{WC}\right|/Ze`$, we start from the case of weak screening when $`r_s`$ is larger than the average distance between $`Z`$-ions. In this case, screening does not affect thermodynamic properties of WC. The energy per $`Z`$-ion $`\epsilon (n)`$ of such Coulomb WC at $`T=0`$ can be estimated as the energy of a Wigner-Seitz cell, because quadrupole-quadrupole interaction between neigbouring neutral Wigner-Seitz cells is very small. This gives
$$\epsilon (n)=(28/3\pi )Z^2e^2/RD1.15Z^2e^2/RD,$$
(15)
where $`R=(\pi n)^{1/2}`$ is the radius of a Wigner-Seitz cell. A more accurate calculation gives slightly higher energy:
$$\epsilon (n)1.11Z^2e^2/RD=1.96n^{1/2}Z^2e^2/D.$$
(16)
One can discuss the role of a finite temperature on WC in terms of the inverse dimensionless temperature $`\mathrm{\Gamma }=Z^2e^2/(RDk_BT)`$. We are interested in the case of large $`\mathrm{\Gamma }`$. For example, at a typical $`Zen=\sigma =1.0e/`$nm<sup>2</sup> and at room temperature, $`\mathrm{\Gamma }=10`$ for $`Z=4`$. Wigner crystal melts at $`\mathrm{\Gamma }=130`$, so that for $`\mathrm{\Gamma }<130`$ we deal with a strongly correlated liquid. Numerical calculations, however, confirm that at $`\mathrm{\Gamma }1`$ thermodynamic properties of strongly correlated liquid are close to that of WC . Therefore, for an estimate of $`\mu _c`$ we can still write $`F_c=n\epsilon (n)`$ and use
$$\mu _{WC}=\frac{\left[n\epsilon (n)\right]}{n}=1.65\mathrm{\Gamma }k_BT=1.65\frac{Z^2e^2}{RD}.$$
(17)
We see now that $`\mu _{WC}`$ is negative and $`|\mu _{WC}|k_BT`$, so that Eq. (14) is justified. Substituting Eq. (17) into Eq. (14), we get $`\sigma ^{}=0.83Ze/(\pi r_sR)`$. At $`r_sR`$, charge density $`\sigma ^{}\sigma `$, and $`Zen\sigma `$, one can replace $`R`$ by $`R_0=(\sigma \pi /Ze)^{1/2}`$. This gives
$$\sigma ^{}/\sigma =0.83\zeta ^{1/2},(\zeta 1),$$
(18)
where $`\zeta =Ze/\pi \sigma r_s^2`$ is the dimensionless charge of a $`Z`$-ion. Thus, at $`r_sR`$ or $`\zeta 1`$, inverted charge density grows with decreasing $`r_s`$. Extrapolating to $`r_s=2R_0`$ where screening starts to modify the interaction between $`Z`$-ions substantially, we obtain $`\sigma ^{}=0.4\sigma `$.
Now we switch to the case of strong screening, $`r_sR`$, or $`\zeta 1`$. It seems that in this case $`\sigma ^{}`$ should decrease with decreasing $`r_s`$, because screening reduces the energy of WC and leads to its melting. In fact, this is what eventually happens. However, there is a range of $`r_sR`$ where the energy of WC is still large. In this range, as $`r_s`$ decreases, the repulsion between $`Z`$-ions becomes weaker, what in turn makes it easier to pack more of them on the plane. Therefore, $`\sigma ^{}`$ continues to grow with decreasing $`r_s`$.
Although we can continue to use the capacitor model to deal with the problem, this model loses its physical transparency when $`r_sR`$, because there is no obvious spatial separation between the inverted charge $`\sigma ^{}`$ and its screening atmosphere. Therefore, at $`r_sR`$, we deal directly with the original free energy (8). The requirement that the chemical potential of $`Z`$-ion in the bulk solution equals that of $`Z`$-ions at the surface now reads
$$\frac{F}{n}=\mu _{id}\mu _b,$$
(19)
where
$$F=\frac{2\pi \sigma r_sZen}{D}+F_{ZZ}$$
(20)
is the interaction part of the total free energy (8) apart from the constant self-energy term $`\pi \sigma ^2r_s/D`$. According to Eq. (12), at large $`N`$ when
$$\mu _b\mu _{id}=k_BT\mathrm{ln}(N_s/N)2\pi \sigma r_sZe/D,$$
(21)
we can neglect the difference in the ideal part of the free energy of $`Z`$-ion at the surface and in the bulk. Therefore, the condition of equilibrium (19) can be reduced to the problem of minimization of the free energy (20) with respect to $`n`$. This direct minimization has a very simple meaning: new $`Z`$-ions are attracted to the surface, but $`n`$ saturates when the increase in the repulsion energy between $`Z`$-ions compensates this gain. Since this minimization balances the attraction to the surface with the repulsion between $`Z`$-ions, the inequality (21) also guarantees that thermal fluctuations of $`Z`$-ions around their WC positions are small. Therefore, $`F_{ZZ}`$ can be written as
$$F_{ZZ}=\underset{𝒓_i0}{}\frac{(Ze)^2}{Dr_i}e^{r_i/r_s},$$
(22)
where the sum is taken over all vectors of WC lattice. At $`r_sR`$, one needs to keep only interactions with the 6 nearest neighbours in Eq. (22). This gives
$$F=\frac{2\pi \sigma r_sZen}{D}+3n\frac{(Ze)^2}{DA}\mathrm{exp}(A/r_s),$$
(23)
where $`A=(2/\sqrt{3})^{1/2}n^{1/2}`$ is the lattice constant of this WC. Minimizing this free energy with respect to $`n`$ one gets $`Ar_s\mathrm{ln}\zeta `$, $`R(2\pi /\sqrt{3})^{1/2}r_s\mathrm{ln}\zeta `$ and
$$\frac{\sigma ^{}}{\sigma }=\frac{2\pi \zeta }{\sqrt{3}\mathrm{ln}^2\zeta },(\zeta 1).$$
(24)
It is clear from Eq. (24) that at $`r_sR`$, or $`\zeta 1`$ the distance $`R`$ decreases and inverted charge continues to grow with decreasing $`r_s`$. This result could be anticipated for the toy model of Fig. 1a if the Coulomb interaction between the spheres is replaced by a strongly screened one. Screening obviously affects repulsion between positive spheres stronger than their attraction to the negative one and, therefore, makes it possible to keep two $`Z`$-ions even at $`QZe`$.
Above we studied analytically two extremes, $`r_sR`$ and $`r_sR`$. In the case of arbitrary $`r_s`$ we can find $`\sigma ^{}`$ numerically. Indeed, minimizing the free energy (20) with the help of Eq. (22) one gets
$$\frac{1}{\zeta }=\underset{𝒓_i0}{}\frac{3+r_i/r_s}{8r_i/r_s}e^{r_i/r_s},$$
(25)
where the sum over all vectors of WC lattice can be evaluated numerically. Using Eq. (25) one can find the equilibrium concentration $`n`$ for any given value of $`\zeta `$. The resulting ratio $`\sigma ^{}/\sigma `$ is plotted by the solid curve in Fig. 5.
## III Condensation of monovalent coions on $`Z`$-ion. Role of finite size of $`Z`$-ion.
We are prepared now to address the question of maximal possible charge inversion. How far can a macroion be overcharged, and what should one do to achieve that? We see below that to answer this questions one should take into account the finite size of $`Z`$-ions.
Fig. 5 and Eq. (24) suggest that the ratio $`\sigma ^{}/\sigma `$ continues to grow with growing $`\zeta `$. However, the possibilities to increase $`\zeta `$ are limited along with the assumptions of the presented theory. Indeed, there are two ways to increase $`\zeta =Ze/\pi \sigma r_s^2`$, namely to choose a surface with a small $`\sigma `$ or to choose $`Z`$-ions with a large $`Z`$. The former way is restricted because, according to Eq. (21), $`Z`$-ion remains strongly bound to the charged plane only as long as $`2\pi r_s\sigma Ze/Dk_BTs`$ where
$$s=\mathrm{ln}(N_s/N)$$
(26)
is the entropy loss (in units of $`k_B`$) per Z-ion due to its adsorption to the surface. This gives for $`\zeta `$:
$$\zeta \zeta _{max}=2Z^2l_B/sr_s.$$
(27)
Therefore, the latter way, which is to increase $`Z`$, is really the most important one. The net charge $`Z`$ of a $`Z`$-ion is, however, restricted because at large charge $`Z_0`$ of the bare counterion monovalent coions of the charged plane (which have the sign opposite to $`Z`$-ions) condense on the $`Z`$-ion surface . Assuming that $`Z`$-ions are spheres of the radius $`a`$, their net charge, $`Z`$, at large $`Z_0`$ can be found from the equation
$$Ze^2/aD=k_BT\mathrm{ln}(N_{1,s}/N_1),$$
(28)
where $`Ze^2/aD`$ is the potential energy of a monovalent coion at the external boundary of the condensation atmosphere (”surface”) of $`Z`$-ion and $`k_BT\mathrm{ln}(N_{1,s}/N_1)`$ is the difference between the chemical potentials of monovalent coions in the bulk and at the $`Z`$-ion’s surface, $`N_{1,s}Z/a^3`$ is the concentration of coions at the surface layer. Eq. (28) gives
$$Z=\left(2a/l_B\right)\mathrm{ln}\left(r_s/a\right).$$
(29)
Using Eq. (29) and Eq. (27), we arrive at
$$\zeta _{max}=\frac{8a^2}{sl_Br_s}\left[\mathrm{ln}\frac{r_s}{a}\right]^2,(r_sa).$$
(30)
In the theory presented in Sec. II, the radius of $`Z`$-ion, $`a`$, was the smallest length, even smaller than $`r_s`$. Therefore, the largest $`a`$ we can substitute in Eq. (30) is $`a=r_s`$. For $`r_s=a=10\AA `$ and $`s=3`$ we get $`\zeta _{max}4`$ so that the inversion ratio can be as large as 2.
Since charge inversion grows with increasing $`a`$ we are tempted to explore the case $`r_saR_0`$. To address this situation, our theory needs a couple of modifications. Specifically, in the first term of Eq. (23) we must take into account the fact that only a part of a $`Z`$-ion interacts with the surface, namely the segment which is within the distance $`r_s`$ from the surface. Therefore, at $`r_sa`$ results depend on the shape of ions and distribution of charge. If the bare charge of $`Z`$-ion is uniformly distributed on the surface of a spherical ion this adds small factor $`r_s/2a`$ to $`\mu _{WC}`$ and the right side of Eq. (27). This gives
$$\zeta _{max}=Z^2l_B/sa.$$
(31)
One should also take into account that at $`ar_s`$ Eq. (29) should be replaced by
$$Z=a^2/r_sl_B,$$
(32)
which follows from the condition that potential at the surface of $`Z`$-ion $`Ze^2/aDZe^2/(a+r_s)D`$ is equal to $`k_BT\mathrm{ln}(N_{1,s}/N_1)`$. Substituting Eq. (32) to Eq. (31) we find that $`\zeta _{max}`$ is larger than that given by Eq. (30), namely
$$\zeta _{max}=\frac{2a^3}{sl_Br_s^2},(r_sa).$$
(33)
For $`a=20\AA `$, $`r_s=10\AA `$, $`l_B=7\AA `$ and $`s=3`$ we get $`\zeta _{max}8`$ so that the inversion ratio can be as large as 3.
Let us consider now a special case of the compact $`Z`$-ion when it is a short rod-like polyelectrolyte of length $`L<R`$ and radius $`a<r_s`$. Such rods lay at the surface of macroion and form strongly correlated liquid reminding WC, so that one can still start from Eq. (27). In this case, however, Eqs. (29) and (32) should be replaced by $`ZL\eta _c/e=L/l_B`$. Thus, $`\zeta _{max}=2R^2/sl_Br_s`$ and can be achieved at $`LR`$.
We conclude this section going back to spherical $`Z`$-ions and relatively weak screening. Until now we used everywhere the Debye-Hückel approximation for description of screening of surface charge density $`\sigma ^{}`$ by monovalent salt. Now we want to verify its validity. Theory of Sec. II requires that the correlation voltage applied to capacitors $`|\mu _{WC}|/Ze`$ is smaller than $`k_BT/e`$. Using Eqs. (14) and (17) one can rewrite this condition as $`Z<R/1.65l_B`$. Substituting $`Z`$ from Eq. (29) we find that one can use linear theory only when $`r_s<r_m`$, where
$$r_m=a\mathrm{exp}(R/3.3a).$$
(34)
For a large $`R/2a`$, the maximal screening radius of linear theory, $`r_m`$, is exponentially large. Nonlinear theory for $`r_s>r_m`$ is given in Sec. VII.
## IV Screening of a thick insulating macroion by spherical $`Z`$-ions: Role of images.
In Sec. II and III we studied a charged plane immersed in water so that screening charges are on both sides of the plane (Fig. 3a). In reality charged plane is typically a surface of a rather thick membrane whose (organic) material has the dielectric constant $`D_1`$ much less than that of water $`D_1D`$. It is well known in electrostatics that when a charge approaches the interface separating two dielectrics, it induces surface charge on interface. The potential created by these induced charges can be described as the potential of an image charge sitting on the opposite site of the interface (Fig. 3b). At $`D_1D`$, this image charge has the same sign and magnitude as the original charge. Due to repulsion from images, $`Z`$-ions are pushed off the surface to some distance, $`d`$. One can easily find $`d`$ in the case of a single $`Z`$-ion near the charged macroion in the absence of screening ($`r_s=\mathrm{}`$). The $`d`$-dependent part of the free energy of this system is
$$F=4\pi \sigma Zed/D+(Ze)^2/4Dd.$$
(35)
Here the first term is the work needed to move $`Z`$-ion from the surface to the distance $`d`$, and the second term is the energy of image repulsion. The coefficient $`4\pi `$ (instead of $`2\pi `$) in the first term accounts for the doubling of the plane charge due to the image of the plane. The ion sits at distance $`d=d_0`$ which minimizes the free energy of Eq. (35). Solving $`F/d=0`$, one gets
$$d_0=\frac{1}{4}\sqrt{\frac{Ze}{\pi \sigma }}=\frac{R_0}{4}.$$
(36)
In the presence of other counterions on the surface, the repulsive force is stronger, therefore one expects that $`d_0`$ is a little larger than $`R_0/4`$.
To consider the role of all images and finite $`r_s`$, let us start from the free energy per unit area describing the system:
$`F`$ $`=`$ $`{\displaystyle \frac{4\pi \sigma r_sZen}{D}}e^{d/r_s}+{\displaystyle \frac{n}{2}}{\displaystyle \underset{𝒓_i0}{}}{\displaystyle \frac{(Ze)^2}{Dr_i}}e^{r_i/r_s}`$ (38)
$`+{\displaystyle \frac{n}{2}}{\displaystyle \underset{𝒓_i}{}}{\displaystyle \frac{(Ze)^2}{D\sqrt{r_i^2+4d^2}}}e^{\sqrt{r_i^2+4d^2}/r_s},`$
where, as in Eq. (22), the sums are taken over all vectors of the WC lattice. The four terms in Eq. (38) are correspondingly the self energy of the plane, the interaction between the plane and the $`Z`$-ions, the interaction between $`Z`$-ions (the factor $`1/2`$ accounts for the double counting), and the repulsion between $`Z`$-ions and the image charges (the factor $`1/2`$ accounts for the fact that electric field occupies only half of the space).
At large concentration of $`Z`$-ions in the bulk, the difference in the ideal parts of the free energy of $`Z`$-ion in solution and at the surface can be neglected, therefore, one can directly minimize the free energy (38) to find the concentration of $`Z`$-ions, $`n`$, at the surface and the optimum distance $`d`$. The system of equations
$$\frac{F}{d}=0,\frac{F}{n}=0,$$
(39)
can be solved numerically and the results are plotted in Fig. 5. A remarkable feature of this plot is that, within 2% accuracy, the ratio $`\sigma ^{}/\sigma `$ for the image problem is equal to a half of the same ratio for the charged plane immersed in water (for which there are no images). If we try to interpret this result using Eq. (14) of the capacitor model (Sec. II) we can say that image charges do not modify the ”correlation” voltage $`|\mu _{WC}|/Ze`$. The only substantial difference between two cases is that for the thick macroion, instead of charging two capacitors, one has to charge only one capacitor (on one side of the surface) with capacitance per unit area $`D/4\pi r_s`$
The fact that image charges do not modify the ”correlation voltage” can be explained quite simply in the case of weak screening $`r_sR_0`$. In this limit, expanding the free energy (38) to the first order in $`d/r_s`$, we get
$$F=n\epsilon (n)+\frac{n}{2}Ze\varphi _{WC}(n,2d)+\frac{2\pi \sigma ^2d}{D}+\frac{2\pi (\sigma ^{})^2r_s}{D}.$$
(40)
The physical meaning of this equation is quite clear. The first two terms are energies of the WC and of its interaction with the image WC ($`\varphi _{WC}(n,2d)`$ is the potential of a WC with charge density $`Zen`$ at the location of an image of $`Z`$-ion.) The third term is the capacitor energy created by the charge of WC and the plane charge. And the final term is the usual energy of a capacitor made by the WC and the screening atmosphere.
At $`\sigma ^{}/\sigma 1`$ minimization of Eq. (40) with respect to $`d`$ gives the optimum distance $`d_0=0.3R_0`$, which is a little larger than the estimate (36). Minimization with respect to $`n`$ gives an equation similar to Eq. (14)
$$\sigma ^{}=\frac{D}{4\pi r_s}\frac{|\mu _{WC}|}{Ze},$$
(41)
where $`\mu _{WC}`$ differs from the corresponding value in the case of immersed plane (Eq. (17)) only by:
$$\delta \mu _{WC}=\frac{}{n}\left[\frac{n}{2}Ze\psi _{WC}(n,2d_0)\right].$$
(42)
It is known that $`\psi _{WC}(x)`$ decreases exponentially with $`x`$ when $`x>A/2\pi `$. Since $`2d_0/(A/2\pi )1.8`$, the potential $`\psi _{WC}(n,2d_0)\mathrm{exp}(2d_02\pi /A)`$ and $`\delta \mu _{WC}/|\mu _{WC}|(1d_0\frac{2\pi }{A})\mathrm{exp}(2d_0/(A/2\pi ))0.02`$. Thus, at $`r_sR_0`$ the chemical potential $`\mu _{WC}`$ remains practically unchanged by image charges.
In the opposite limit $`r_sR_0`$ one can calculate the ratio $`\sigma ^{}/\sigma `$ by direct minimization of the free energy, without the use of the capacitor model. Keeping only the nearest neighbour interactions in Eq. (38) one finds
$$d_0=r_s\mathrm{ln}\frac{\zeta }{8},$$
$$\frac{\sigma ^{}}{\sigma }\frac{2\pi \zeta }{\sqrt{3}\mathrm{ln}^2(\zeta ^2/10(d/r_s))}\frac{\pi \zeta }{2\sqrt{3}\mathrm{ln}^2\zeta }.$$
(43)
Comparing this result with Eq. (24) for the case of immersed plane (no image charges), one gets
$$\frac{(\sigma ^{}/\sigma )_{\mathrm{𝑖𝑚𝑎𝑔𝑒}}}{(\sigma ^{}/\sigma )_{\mathrm{𝑛𝑜}\mathrm{𝑖𝑚𝑎𝑔𝑒}}}=\frac{1}{4}\left[1+\frac{\mathrm{ln}10}{\mathrm{ln}\zeta }\right].$$
(44)
Eq. (44) shows that in the limit $`\zeta \mathrm{}`$, the ratio $`\sigma ^{}/\sigma `$ for the image problem actually approaches 1/4 of that for the problem without image. However, due to the logarithmic functions, it approaches this limit very slowly. Detailed numerical calculations show that even at $`\zeta =1000`$, the ratio (44) is still close to 0.5. In practice, $`\zeta `$ can hardly exceed $`20`$, and this ratio is always close to 0.5 as Fig. 5 suggested.
Although at a given $`\zeta `$, image charges do not change the results qualitatively, they, as we show below, reduce the value of $`\zeta _{max}`$ substantially. As in Sec. III, we find $`\zeta _{max}`$ from the condition that the bulk electrochemical potential of $`Z`$-ions can be neglected. When images are present, according to Eq. (38), one need to replace the right hand side of Eq. (21) by $`2\pi \sigma r_sZe\mathrm{exp}(d_0/r_s)`$. Using Eq. (43), this condition now reads
$$\zeta \zeta _{max}=4\sqrt{Z^2l_B/sr_s}$$
(45)
Using Eq. (45) instead of Eq. (27) and using Eq. (29) for $`Z`$ we get $`\zeta _{max}5`$ at $`r_s=a=10\AA `$ and $`s=3`$. Therefore, according to the dotted curve of Fig. 5 which was calculated for the case of image charges, the inversion ratio for a thick macroion can be as large as 100%.
## V Long charged rods as $`Z`$-ions. Strong screening by monovalent salt.
As we mentioned in Introduction the adsorption of long rod-like $`Z`$-ions such as DNA double helix on an oppositely charged surface leads to the strong charge inversion. In this case, correlations between rods cause parallel ordering of rods in a strongly correlated nematic liquid. In other words, in the direction perpendicular to the rods we deal with short range order of one-dimensional WC (Fig. 4).
Consider the problem of screening of a positive plane with surface charge density, $`\sigma `$, by negative DNA double helices with the net linear charge density $`\eta `$ and the length $`L`$ smaller than the DNA persistence length $`L_p`$ so that they can be considered straight rods. For simplicity, the charged plane is assumed to be thin and immersed in water so that we can neglect image charges. Modification of the results due to image charges is given later. Here, the strong screening case $`r_sA`$ is considered ($`A`$ is the WC lattice constant). The weak screening case, $`r_sA`$, is the topic of the next section.
We show below that at $`r_sA`$ screening radius $`r_s`$ is smaller than the Gouy-Chapman length for the bare plane. Therefore, one can use Debye-Hückel formula, $`\psi (0)=2\pi \sigma r_s/D`$, for the potential of the plane. On the other hand, the ”bare” surface charge of DNA is very large, and its corresponding Gouy-Chapman length is much smaller than $`r_s`$. As the result, one needs nonlinear theory for description of the net charge of DNA. It leads to Onsager-Manning conclusion that positive monovalent ions condense on the surface of DNA reducing its net charge, $`\eta `$, to $`\eta _c=Dk_BT/e`$. Far away from DNA, the linear theory can be used. When DNA rods condense on the plane, we can still use $`\eta _c`$ as the net charge density of DNA, because as we will see later, the strongly screened potential of plane only weakly affects condensation of monovalent ions on DNA.
Therefore, we can write the free energy per DNA as
f=2πσrsLηcD+12
=i-
i0
2Lηc2DK0(iArs),𝑓2𝜋𝜎subscript𝑟𝑠𝐿subscript𝜂𝑐𝐷12superscriptsubscript
=i-
i0
2𝐿superscriptsubscript𝜂𝑐2𝐷subscript𝐾0𝑖𝐴subscript𝑟𝑠f=-\frac{2\pi\sigma r\!_{s}L\eta_{c}}{D}+\frac{1}{2}\sum_{\parbox{28.18524pt}{$i=-\infty$\\
$~{}~{}~{}i\neq 0$}}^{\infty}\frac{2L\eta_{c}^{2}}{D}K_{0}\left(\frac{iA}{r\!_{s}}\right), (46)
where $`K_0(x)`$ is the modified Bessel function of $`0`$-th order. The first term of Eq. (46) describes the interaction energy of DNA rods with the charged plane, the second term describes the interaction between DNA rods arranged in one-dimensional WC, the factor $`1/2`$ accounts for the double counting of the interactions in the sum.
Since the function $`K_0(x)`$ exponentially decays at large $`x`$, at $`r_sA`$ one can keep only the nearest neighbour interactions in Eq. (46). This gives
$$f\frac{2\pi \sigma r_sL\eta _c}{D}+\frac{2L\eta _c^2}{D}\sqrt{\frac{\pi r_s}{2A}}\mathrm{exp}(A/r_s),$$
(47)
which is similar to Eq. (23). To find $`A`$, we minimize the free energy per unit area, $`F=nf`$, with respect to $`n`$, where $`n=1/LA`$ is the concentration of DNA helices at the charged plane. This yields:
$$\frac{\sqrt{2\pi }\sigma r_s}{\eta _c}=\sqrt{A/r_s}\mathrm{exp}(A/r_s).$$
(48)
Calculating the net negative surface charge density, $`\sigma ^{}=\eta _c/A+\sigma `$, we obtain for the inversion ratio
$$\frac{\sigma ^{}}{\sigma }\frac{\eta _c/\sigma r_s}{\mathrm{ln}(\eta _c/\sigma r_s)}(r_sA).$$
(49)
As we see from Eq. (48), the lattice constant $`A`$ of WC decreases with decreasing $`r_s`$ and charge inversion becomes stronger.
Let us now address the question of the maximal charge inversion in the case of screening by DNA. Similar to what was done in Sec. III, the charge inversion ratio is limited by the condition that the electrochemical potential of DNA in the bulk solution can be neglected and therefore, DNA is strongly bound to the surface. Using Eq. (47) and (48), this condition can be written by an equation similar to Eq. (21)
$$k_BTs2\pi \sigma r_sL\eta _c/D\text{or }\eta _c/\sigma r_s2\pi L/sl_B,$$
(50)
where $`s=\mathrm{ln}(N_{s,DNA}/N_{DNA})`$ is the entropy loss (in units of $`k_B`$) per DNA due to its adsorption to the surface. $`N_{s,DNA}`$ and $`N_{DNA}`$ are correspondingly the three-dimensional concentration of DNA at the charged surface and in the bulk. Inequality (50) also guarantees that WC-like short range order of DNA helices is preserved. To show this, let us assume that the left and right nearest neighbour rods at the surface are parallel to each other and discuss the amplitude of the thermal fluctuations of the central DNA along the axis $`x`$ perpendicular to DNA direction (in the limit $`r_sA`$, we need to deal only with two nearest neighbours of the central DNA). At $`x=0`$, the free energy of the rod is given by Eq. (47). At $`x0`$ the free energy of the central DNA is
$$f(x)\frac{2\pi \sigma r_s\eta _cL}{D}+\frac{2L\eta _c^2}{D}\sqrt{\frac{\pi r_s}{2A}}\mathrm{cosh}\left(\frac{x}{r_s}\right)e^{A/r_s}.$$
(51)
Using Eqs. (51) and (48), we find the average amplitude, $`x_0`$, of the fluctuations of $`x`$ from the condition $`f(x_0)f(0)k_BT`$. This gives $`x_0r_s\mathrm{ln}(Ae/2\pi \sigma r_s^2L)`$. The inequality (50) then gives:
$$x_0<r_s\mathrm{ln}\frac{A}{r_s}r_s\mathrm{ln}\mathrm{ln}\frac{\eta _c}{\sigma r_s}Ar_s\mathrm{ln}\frac{\eta _c}{\sigma r_s}.$$
(52)
Thus, DNA helices preserve WC-like short range order when the condition (50) is met.
This condition obviously puts only a weak restriction on maximum value of $`\sigma ^{}/\sigma `$. At $`L=L_p=50`$ nm and $`s=3`$, the parameter $`\eta _c/\sigma r_s`$ can be as large as 75 and, according to Eq. (49) the ratio $`\sigma ^{}/\sigma `$ can reach 15. Therefore, we can call this phenomenon strong charge inversion.
This limit can be easily reached at a very small $`\sigma `$. On the other hand, if we want to reach it making $`r_s`$ very small we have to modify this theory for the case when $`r_s`$ is smaller than the radius of DNA. In a way, this is similar to what was done in Sec. 3 for spherical $`Z`$-ions. At $`r_sa`$ one replaces the net charge of DNA, $`\eta _c`$ by $`\eta _ca/r_s`$ and adds small factor $`(r_s/\pi ^2a)^{1/2}`$ to the first term of Eq. (47). This modification changes only logarithmic term of Eq. (49) and does not change our conclusion about strong charge inversion.
One can numerically minimize the free energy (46) at all $`r_sA`$ to find $`\sigma ^{}/\sigma `$. The result is plotted by the solid curve in Fig. 6.
Let us now move to the more realistic case of a thick macroion, so that repulsion from image charges must be taken into consideration. As in the spherical $`Z`$-ion case, image charges push the WC off the surface to some distance $`d`$. The free energy per DNA rod can be written as
$`f`$ $`=`$ 4πσrsLηcDed/rs+12
=i-
i0
2Lηc2DK0(iArs)4𝜋𝜎subscript𝑟𝑠𝐿subscript𝜂𝑐𝐷superscript𝑒𝑑subscript𝑟𝑠12superscriptsubscript
=i-
i0
2𝐿superscriptsubscript𝜂𝑐2𝐷subscript𝐾0𝑖𝐴subscript𝑟𝑠\displaystyle-\frac{4\pi\sigma r\!_{s}L\eta_{c}}{D}e^{-d/r\!_{s}}+\frac{1}{2}\sum_{\parbox{28.18524pt}{$i=-\infty$\\
$~{}~{}~{}i\neq 0$}}^{\infty}\frac{2L\eta_{c}^{2}}{D}K_{0}\left(\frac{iA}{r\!_{s}}\right) (54)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{2L\eta _c^2}{D}}K_0\left({\displaystyle \frac{\sqrt{(iA)^2+4d^2}}{r_s}}\right),`$
where the three terms on the right hand side are correspondingly the interaction between the plane and the DNA, between the different DNAs and between the DNAs and their images.
The equilibrium distance $`d_0`$ and $`A`$ can be obtained by minimizing the free energy per unit area $`F=nf`$ with respect to $`d`$ and $`n=1/LA`$:
$$\frac{F}{d}=0,\frac{F}{n}=0,$$
(55)
This system of equations is solved numerically. The result for $`\sigma ^{}/\sigma `$ is plotted by the dotted curve in Fig. 6. It is clear that in the case of DNA, at a given value of $`\eta _c/\sigma r_s`$, image charges play even smaller role than for spherical $`Z`$-ions. The ratio $`\sigma ^{}/\sigma `$ in the case of a thick macroion is close to 70% of $`\sigma ^{}/\sigma `$ for the charged plane immersed in water, instead of 50% as in Fig. 5 for spherical $`Z`$-ions.
However, like in the case of spherical $`Z`$-ions, image charges modify the maximal possible value of $`\eta _c/\sigma r_s`$ significantly. When images are present, according to Eq. (54), one need to replace in Eq. (50) $`2\pi \sigma r_sL\eta _c/D`$ by $`2\pi \sigma r_sL\eta _c\mathrm{exp}(d_0/r_s)`$. Therefore, the condition that the bulk ideal chemical potential can be neglected and, therefore, DNA is strongly bound at the surface has the form
$$k_BTs2\pi \sigma r_sL\eta _c\mathrm{exp}(d_0/r_s).$$
(56)
Similarly to what was done above for the problem of charged plane immersed in water one can show that Eq. (56) guarantees, also, WC-like short range order of DNA helices. In the limit $`\eta _c\sigma r_s`$, keeping only the nearest neighbour interactions in the free energy (54) and minimizing with respect to $`d`$ one gets $`d_0r_s\mathrm{ln}(\eta _c/4\sigma r_s)`$. Substituting $`d_0`$ into Eq. (56) we arrive at the final form for the condition of Eq. (50):
$$\eta _c/\sigma r_s\sqrt{8\pi L/sl_B},(LL_p).$$
(57)
It is clear that the maximal $`\eta _c/\sigma r_s`$ and maximal inversion ratio grow with $`L`$. For $`L=L_p=50`$ nm and $`s=3`$, the maximal $`\eta _c/\sigma r_s`$ = 25. Therefore, according to the dotted curve in Fig. 6, the inversion ratio for a thick macroion $`\sigma ^{}/\sigma `$ can reach 4. Such inversion can still be considered as strong.
Until now we talked about relatively short DNA, $`LL_p`$, which can be considered as a rod. For DNA double-helices of a larger length ($`LL_p`$) the maximum inversion ratio saturates at the value obtained above at $`L=L_p`$. This happens because even a long DNA can not be adsorbed at the surface if for $`L=L_p`$ inequality (57) is violated. (See the theory of adsorption-desorption phase transition, for example, in Ref. ).
On the other hand, if inequality (57) holds at $`L=L_p`$, i. e. at $`\eta _c/\sigma r_s\sqrt{8\pi L_p/sl_B}`$, the adsorption of a long DNA is so strong that DNA lays flat on the charged surface. Since repulsion between neighbouring parallel DNA is balanced with attraction to the surface, interactions between parallel DNA helices are so strong that the same inequality guarantees WC-like short range order at the length scale $`L_p`$, even though DNA length is much larger than $`L_p`$. One can verify this statement studying lateral fluctuations of a DNA segment with length $`L_p`$ similarly to the calculation presented above for the problem of charged plane immersed in water (See Eqs. (51) and (52)). Thus, our theory and the plots of Fig. 5 are applicable for a long DNA and, therefore, for any flexible polyelectrolyte.
To conclude this section, we would like to mention another charge inversion problem similar to the problem we considered here. Giant charge inversion can be also achieved if a single very long DNA double helix screens a long and wide positively charged cylinder with radius greater or about the double helix DNA persistence length (Fig. 7).
In this case, an DNA double helix spirals around the cylinder. Neighbouring turns repel each other so that DNA forms an almost perfect coil which locally resembles one-dimensional WC. As a result, the cylinder charge inverts its sign: density of DNA charge per unit length of the cylinder becomes larger than the bare linear charge density of the cylinder. At small $`r_s`$ this charge inversion can be as strong as we discussed above. If cylinder diameter is smaller than DNA persistent length one should add elastic energy to the minimization problem. This, of course, will make charge inversion weaker than for wider cylinders, but still it can be quite large. We leave open the possibility to speculate on the relevance of these model systems to the fact that DNA overcharges a nucleosome by about 15% .
A similar problem of wrapping of a weakly charged polyelectrolyte around oppositely charged sphere was recently studied in the Debye-Hückel approximation in Ref. . A strong charge inversion was found in this case as well. Charge inversion for a charged sphere screened by an oppositely charged flexible polyelectrolyte was previously observed in experiment and numerical simulations .
## VI Long charged rods as $`Z`$-ions. Weak screening by monovalent salt
In this section, we consider screening of a positively charged plane by DNA rods in the case of weak screening, when $`r_sA`$. We saw in Sec. II that when the screening radius is larger than the lattice constant of WC, the capacitor model provides a transparent description of the charge inversion. Here we adopt this model, too. However, we find out that in the case of rods, the inversion charge $`\sigma ^{}`$ is so large that its screening by monovalent salt is nonlinear. In other words, at $`r_sA`$, the capacitors described in Sec. II becomes nonlinear. Correspondingly in this case one has to use the solution of the nonlinear PB equation for the plane potential:
$$\psi (0)(2k_BT/De)\mathrm{ln}(r_s/\lambda ).$$
(58)
where $`\lambda =e/\pi \sigma ^{}l_B`$ is the Gouy-Chapman length. It is shown below that $`A\lambda r_s`$ so that the use of Eq. (58) is justified.
The weak screening of the plane potential has also another important consequence. The net charge density of DNA, $`\eta `$, ceases to be equal to to the Onsager-Manning critical density $`\eta _c`$. The charge of the plane forces DNA to release some of monovalent coions condensed on it, so that $`\eta `$ becomes larger than $`\eta _c`$. Thus, in this case, we have to deal with a nonlinear problem with two unknowns, $`\eta `$ and $`\sigma ^{}`$.
One can find these unknowns from the two following physical conditions of equilibrium. The first one requires that the chemical potential of positive monovalent ions (coions) in the bulk of solution is equal to the chemical potential of coions condensed on the surface of DNA rods which, in turn, are adsorbed on the plane. The second condition requires that the chemical potentials for DNA rods in the bulk solution and DNA rods of the surface WC are equal. Let us write the first condition as
$$k_BT\mathrm{ln}\frac{N_{1,s}}{N_1}=e\psi (0)+\frac{2e\eta }{D}\mathrm{ln}(A/2\pi a),$$
(59)
where $`N_1`$ and $`N_{1,s}`$ are the concentrations of monovalent coion in the bulk and at the DNA surface respectively. The left-hand side of Eq. (59) is the entropy loss and the right-hand side is the potential energy gain when monovalent salt condenses on the DNA surface (the potential at the surface of DNA is the sum of $`\psi (0)`$, of the nonlinear plane capacitor made and the potential of the DNA charged cylinder with radius $`a`$ and the linear charge density $`\eta `$, screened at the distance $`A/2\pi `$, by neigbouring DNA). Far from the charged plane, DNA net charge regains its value $`\eta _c`$, the condition of equilibrium of condensed monovalent coions on isolated DNA rod with those in the bulk can be written in a way similar to Eq. (59):
$$k_BT\mathrm{ln}\frac{N_{1,s}}{N_1}=\frac{2e\eta _c}{D}\mathrm{ln}\frac{r_s}{a}.$$
(60)
Excluding $`\mathrm{ln}(N_{1,s}/N_1)`$ from Eqs. (59) and (60) and using Eq. (58) we can write the first equation for $`\lambda `$ (which represents $`\sigma ^{}`$) and $`\eta `$ as
$$\eta _c\mathrm{ln}\frac{r_s}{a}=\eta _c\mathrm{ln}\frac{r_s}{\lambda }+\eta \mathrm{ln}\frac{A}{2\pi a}.$$
(61)
The equality of the chemical potential of DNA in the bulk and of DNA condensed on the plane can be written in the form similar to Eq. (14)
$`L\eta \psi (0)`$ $`=`$ $`|\mu _{WC}|+{\displaystyle \frac{L\eta L\eta _c}{e}}k_BT\mathrm{ln}{\displaystyle \frac{N_{1,s}}{N_1}}`$ (63)
$`\left({\displaystyle \frac{L\eta ^2}{D}}\mathrm{ln}{\displaystyle \frac{\lambda }{a}}{\displaystyle \frac{L\eta _c^2}{D}}\mathrm{ln}{\displaystyle \frac{r_s}{a}}\right).`$
As in Eq. (14), we see that a ”correlation voltage”, $`|\mu _{WC}|/L\eta `$, charges two capacitors consisting of the overcharged plane and its screening atmosphere to a finite voltage $`\psi (0)`$. The new second and third terms on the right hand side are due to the change in the net charge of DNA, when it condenses on the plane. Specifically, the second term is the gain in the entropy of monovalent salt released and the third term is the loss in the self energy of DNA when its net charge changes from $`\eta _c`$ in the bulk solution to $`\eta `$ at the plane surface. Here $`\lambda `$ is the screening length near the plane surface. (This can be seen from the fact that the three-dimensional concentration of monovalent salt at the surface is of the order $`N_{1,s}\sigma ^{}/2e\lambda `$ and the corresponding screening length $`r_{s,surf}=(4\pi N_{1,s}l_B)^{1/2}(2\lambda e/\pi \sigma ^{}l_B)^{1/2}\lambda `$.)
A formal derivation of Eq. (63) is given in the end of this section.
The free energy per DNA of the one-dimensional WC of DNA rods at the surface can be written similarly to Eq. (46) with the screening length $`r_s`$ replaced by $`\lambda `$,
$`f`$ $`=`$ $`{\displaystyle \frac{2\pi (\eta /A)\lambda }{D}}L\eta +{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=\mathrm{},i0}{\overset{\mathrm{}}{}}}L\eta {\displaystyle \frac{2\eta }{D}}K_0\left({\displaystyle \frac{iA}{\lambda }}\right)`$ (64)
$``$ $`{\displaystyle \frac{L\eta ^2}{D}}\mathrm{ln}{\displaystyle \frac{2\pi \lambda }{A}}.`$ (65)
This result can be interpreted as the interaction of DNA with its Wigner-Seitz cell (a stripe with length $`L`$, width $`A`$ and charge density $`\eta /A`$).
The chemical potential $`\mu _{WC}`$ can be easily calculated:
$$\mu _{WC}=\frac{[nf]}{n}\frac{L\eta ^2}{D}\mathrm{ln}\frac{2\pi \lambda }{A},$$
(66)
where $`n=1/LA`$ is the concentration of DNA at the charged surface.
Substituting Eqs. (58), (60), and (66) into Eq. (63), we arrive at the second equation for $`\eta `$ and $`\lambda `$
$$2\eta \eta _c\mathrm{ln}\frac{r_s}{\lambda }=\eta _c^2\mathrm{ln}\frac{r_s}{a}\eta ^2\mathrm{ln}\frac{A}{2\pi a}+2\eta _c\eta \mathrm{ln}\frac{r_s}{a}.$$
(67)
Solving Eqs. (61) and (67) together with $`A=\eta /(\sigma +\sigma ^{})`$, we get
$$\eta \eta _c\sqrt{\frac{\mathrm{ln}(r_s/a)}{\mathrm{ln}(A_0/2\pi a)}},$$
(68)
$$\mathrm{ln}\frac{\lambda }{a}\sqrt{\mathrm{ln}\frac{r_s}{a}\mathrm{ln}\frac{A_0}{2\pi a}},$$
(69)
where $`A_0=\eta _c/\sigma `$.
Eq. (69) shows that the theory is self consistent: when $`r_sA_0`$, one has $`r_s\lambda A_0`$. This justifies the use of nonlinear potential for the plane. Eq. (68) demonstrates that $`\eta \eta _c`$ as we anticipated. Eq. (68), of course, is valid only if $`\eta \eta _0`$, where $`\eta _0`$ is bare linear charge density of DNA.
The ratio $`\sigma ^{}/\sigma `$ can now be easily calculated by substituting $`\lambda =e/\pi \sigma ^{}l_B`$ into Eq. (69). One arrives at Eq. (6) which shows that the ratio $`\sigma ^{}/\sigma `$ increases as $`r_s`$ decreases, but remains smaller than unity. When $`r_sA_0`$ one finds from Eqs. (68) and (69) that $`\eta \eta _c`$, $`\lambda r_sA_0`$, and $`\sigma ^{}/\sigma 1`$, what matches the Eq. (49) obtained for the strong screening limit ($`r_sA`$).
Let us now present a derivation of Eq. (63). To calculate the free energy of the system we use the standard charging procedure described, for example, in Ref. and used for DNA in Ref. . First, let us start by calculating the electrostatic free energy of a DNA dissolved in solution, which can be written as the work needed to charge the DNA up to the bare value $`\eta _0`$ per unit length
$$f=L_0^{\eta _0}\varphi (\eta ^{})𝑑\eta ^{},$$
(70)
where $`\varphi (\eta ^{})`$ is the self consistent surface potential of DNA when its charge is $`\eta ^{}`$ per unit length. Following Ref. , let us divide this charging process in two steps. First, the DNA is charged from $`0`$ up to $`\eta _c`$. In this step, one can use for $`\varphi (\eta ^{})`$ the linear (Debye-Hückel) potential
$$\varphi (\eta ^{})=\frac{2\eta ^{}}{D}\frac{K_0(a/r_s)}{K_1(a/r_s)a/r_s}\frac{2\eta ^{}}{D}\mathrm{ln}(r_s/a),(r_sa).$$
(71)
In the next step, DNA is charged from $`\eta _c`$ to $`\eta `$. In this step, one has to use nonlinear potential for $`\varphi (\eta ^{})`$. It can be written as a sum
$$\varphi (\eta ^{})=2\frac{k_BT}{De}\mathrm{ln}\frac{a}{\mathrm{\Lambda }(\eta ^{})}+\frac{2\eta _c}{D}\mathrm{ln}\frac{r_s}{a},$$
(72)
where the first term is the contribution of the interval $`2a>r>a`$ of the distances $`r`$ from the DNA axis. In this interval potential can be approximated by that of the charged plane with charge density $`\eta ^{}/2\pi a`$. It has Gouy-Chapman form with the corresponding Gouy-Chapman length $`\mathrm{\Lambda }(\eta ^{})=a\eta _c/\eta ^{}<a`$. The second term in Eq. (72) is the contribution of interval $`\mathrm{}>r>2a`$, where we deal with a cylinder of radius $`a`$ and linear net charge density $`\eta _c`$. Now, we can calculate the free energy of a DNA rod (which is also the chemical potential of DNA in the bulk solution, apart from an ideal part):
$`f`$ $`=`$ $`L{\displaystyle _0^{\eta _c}}\varphi (\eta ^{})𝑑\eta ^{}+L{\displaystyle _{\eta _c}^{\eta _0}}\varphi (\eta ^{})𝑑\eta ^{}`$ (73)
$`=`$ $`{\displaystyle \frac{L\eta _c^2}{D}}\mathrm{ln}{\displaystyle \frac{r_s}{a}}+{\displaystyle \frac{L2\eta _c\eta _0}{D}}\mathrm{ln}{\displaystyle \frac{a}{\mathrm{\Lambda }(\eta _0)}}+{\displaystyle \frac{L2\eta _c}{D}}(\eta _0\eta _c)\mathrm{ln}{\displaystyle \frac{r_s}{a}}.`$ (74)
In the Onsager-Manning condensation theory, one can think of the last two terms in the above expression as the free energy of the condensation layer.
When DNA rods are adsorbed on the surface of the macroion, the role of $`\eta _c`$ as a border between the linear and nonlinear charging regime is played by the net charge $`\eta `$. We calculate the total free energy of the system by first charging the plane surface to $`\sigma `$ and DNA to $`\eta `$ respectively, then continue charging the DNA from $`\eta `$ to the final value $`\eta _0`$. The first charging process leads to the standard contribution
$$L\eta \psi (0)+\mu _{WC}+\frac{L\eta ^2}{D}\mathrm{ln}\frac{\lambda }{a},$$
to the chemical potential of DNA, where the three terms result from, correspondingly, the capacitor energy of the screening atmosphere, the correlation energy of DNA and the self energy of DNA. The second charging process builds up the condensation layer around each DNA and gives a contribution
$$_\eta ^{\eta _0}\varphi (\eta )𝑑\eta =\frac{2\eta _c\eta _0}{D}\mathrm{ln}\frac{a}{\mathrm{\Lambda }(\eta _0)}+\frac{2\eta _c}{D}(\eta _0\eta )\mathrm{ln}\frac{r_s}{a}$$
where the nonlinear potential of Eq. (72) was used.
The chemical potential of DNA on the charged surface is the sum of the two above contributions:
$`L\eta \psi (0)+\mu _{WC}+{\displaystyle \frac{L\eta ^2}{D}}\mathrm{ln}{\displaystyle \frac{\lambda }{a}}+{\displaystyle \frac{L2\eta _c\eta _0}{D}}\mathrm{ln}{\displaystyle \frac{a}{\mathrm{\Lambda }(\eta _0)}}`$ (76)
$`+{\displaystyle \frac{L2\eta _c}{D}}(\eta _0\eta )\mathrm{ln}{\displaystyle \frac{r_s}{a}}.`$ (77)
Equating this expression to the chemical potential of DNA in the bulk (Eq. (LABEL:DNAmufullb)) one gets the desired Eq. (63).
So far, we have dealt only with the screening of charged surface by DNA double helices which are highly charged polyelectrolytes. The situation is simpler if one deals with weakly charged polyelectrolytes whose bare charge density $`\eta _0`$ is much smaller than $`\eta _c`$. In this case, there is no condensation on the polyelectrolyte. Therefore $`\eta _0`$ plays the role of the net charge $`\eta _c`$. In the weak screening case, $`r_s\eta _0/\sigma `$, this brings about small changes in Eq. (63), which now reads:
$$L\eta _0\psi (0)=|\mu _{WC}|\left(\frac{L\eta _0^2}{D}\mathrm{ln}\frac{\lambda }{a}\frac{L\eta _0^2}{D}\mathrm{ln}\frac{r_s}{a}\right).$$
(78)
Substituting Eq. (58) and (66) into Eq. (78), and solving for $`\lambda `$, we get
$$\lambda r_s\mathrm{exp}\left(\frac{\eta _0}{\eta _c}\mathrm{ln}\frac{\eta _0}{\sigma r_s}\right)(r_s\eta _0/\sigma ).$$
(79)
Nonlinear effects are important when $`\lambda r_s`$, or when the exponent in the above expression becomes less than $`1`$. This gives the minimal $`r_s`$ at which nonlinear effects are still important.
$$r_m=(\eta _0/\sigma )\mathrm{exp}(\eta _c/\eta _0).$$
(80)
As we see, $`r_m`$ is exponentially large at $`\eta _c/\eta _01`$. This makes this weak screening case practically unimportant. At smaller, more realistic value of $`r_s`$, one can use Debye-Hückel linear theory to describe the potential of the plane. For $`r_s<\eta _0/\sigma `$, this has been done in Ref. . The result is an expression similar to Eq. (49) with the net charge $`\eta _c`$ replaced by the bare charge $`\eta _0`$.
## VII Nonlinear screening of a charged surface by spherical $`Z`$-ions.
Let us now return to the screening of the charged plane by spherical $`Z`$-ions in the case when screening by monovalent salt is very weak. Our goal is to understand what happens when screening radius is larger than $`r_m`$ (see Eq. (34)), so that Debye-Hückel approximation of Sec. II for the description of screening of surface charge density $`\sigma ^{}`$ by monovalent salt fails and a nonlinear description is necessary.
The nonlinearity of screening leads to two important changes in the theory in Sec. II. First, the monovalent coions condense on the surface of the $`Z`$-ion and reduce its apparent charge. We discussed this condensation in Sec. III, but used for the net charge of $`Z`$-ion the value obtained for isolated $`Z`$-ion in the bulk solution (Eq. (28)). In this section we call this charge $`Z_c`$ (this quantity plays a similar role as $`\eta _c`$ in previous section) and save notation $`Z`$ for the net charge of $`Z`$-ion absorbed at the charged surface as a part of the WC. When positive $`Z`$-ions condense on the negative surface, a fraction of monovalent negative ions, condensed on the $`Z`$-ions is released. Therefore, strictly speaking, $`Z>Z_c`$. The charge $`Z_c`$ can be found from Eq. (28), which in the revised notation reads
$$e\frac{Z_ce}{aD}=k_BT\mathrm{ln}\frac{N_{1,s}}{N_1}.$$
(81)
Here, as in Sec. III, $`N_{1,s}`$ is the concentration of monovalent negative ions at the external boundary of the condensation atmosphere of the isolated spherical $`Z`$-ion. The net charge of a $`Z`$-ion in WC, $`Z`$, can be found from the condition of equilibrium of monovalent negative ions condensed on a $`Z`$-ion of the WC and those in the bulk solution
$$\frac{Ze^2}{aD}e\left(\frac{2.2Ze}{RD}+\psi (0)\right)=k_BT\mathrm{ln}\frac{N_{1,s}}{N_1}.$$
(82)
The term in the parentheses is the total potential of the plane and other adsorbed $`Z`$-ions at the considered $`Z`$-ion. This potential is the sum of the negative potential of WC and the potential due to the positive net charge $`\sigma ^{}`$ of the plane given by Eq. (58). Excluding $`k_BT\mathrm{ln}(N_{1,s}/N_1)`$ from Eqs. (81) and (83) we obtain the first equation for two unknowns $`Z`$ and $`\lambda `$, which is similar to Eq. (61):
$$\frac{Ze^2}{aD}\frac{2.2Ze^2}{RD}+2k_BT\mathrm{ln}(r_s/\lambda )=\frac{Z_be^2}{aD}.$$
(83)
To write the second equation for these unknowns we start from the condition that the chemical potentials of $`Z`$-ion at the charged surface and in the bulk solution are equal. In the close analogy with Eq. (63) of Sec. VI, we can write
$`Ze\psi (0)`$ $`=`$ $`|\mu _{WC}|+(ZZ_c)k_BT\mathrm{ln}{\displaystyle \frac{N_{1s}}{N_1}}`$ (85)
$`+{\displaystyle \frac{(Z_b^2Z^2)e^2}{2aD}}.`$
The second and third terms on the right-hand side account for the fact that monovalent ions are released when $`Z`$-ions condense on the plane surface (so that their entropy is gained) and simultaneously the self energy of the $`Z`$-ion is reduced. Using Eqs. (17), (58) and (81) we obtain the second equation for $`Z`$ and $`\lambda `$
$`2k_BTZ\mathrm{ln}{\displaystyle \frac{r_s}{\lambda }}`$ $`=`$ $`{\displaystyle \frac{1.65(Ze)^2}{RD}}+{\displaystyle \frac{(ZZ_c)Z_be^2}{aD}}`$ (87)
$`+{\displaystyle \frac{(Z_b^2Z^2)e^2}{2aD}}.`$
Solving Eqs. (83) and (87) we get
$$\frac{Z}{Z_c}1+\frac{0.56a}{R}$$
(88)
and
$$\lambda =r_s\mathrm{exp}\left(\frac{1.65a}{2R}\mathrm{ln}\frac{N_{1s}}{N_1}\right).$$
(89)
Approximating $`N_{1,s}`$ as $`N_{1,s}Z/a^3`$, we get
$$\mathrm{ln}\frac{N_{1s}}{N}=2\mathrm{ln}\frac{r_s}{a}.$$
Therefore
$$\lambda =r_s\mathrm{exp}\left(\frac{1.65a}{R}\mathrm{ln}\frac{r_s}{a}\right).$$
(90)
Nonlinear effects are important when $`\lambda r_s`$, or when the exponent in the above expression becomes less than $`1`$. This gives the minimal $`r_s`$ at which nonlinear effects are still important
$$r_m=a\mathrm{exp}(R/1.65a),$$
(91)
which matches the estimate Eq. (34) obtained from the side of the linear regime.
The ratio $`\sigma ^{}/\sigma `$ can be easily calculated from Eq. (89)
$`{\displaystyle \frac{\sigma ^{}}{\sigma }}`$ $`=`$ $`{\displaystyle \frac{e}{\pi \sigma l_Br_s}}\mathrm{exp}\left({\displaystyle \frac{1.65a}{R}}\mathrm{ln}{\displaystyle \frac{r_s}{a}}\right)`$ (92)
$`=`$ $`{\displaystyle \frac{e}{\pi \sigma l_Br_s}}\left({\displaystyle \frac{r_s}{a}}\right)^{1.65a/R}r_s^{(1+1.65a/R)}.`$ (93)
Once again, this ratio increases as $`r_s`$ decreases . Comparing Eq. (93) to Eq. (18), we see that nonlinear effects change the exponent in the dependence of $`\sigma ^{}/\sigma `$ on $`r_s`$ by $`1.65a/R1`$. Taking into account the fact that it is important only when $`r_s`$ is greater than an exponentially large critical value $`r_m`$ (see Eq. (91)), one can conclude from this section that, in practical situation, non-linear effects in the problem of screening of a charged surface by spherical $`Z`$-ions are not important.
## VIII Screening of a macroion with a mobile surface charge.
So far we have assumed that the bare surface charges of the macroion are fixed and can not move. For solid or glassy surfaces, colloidal particles and even rigid polyelectrolytes, such as double helix DNA and actin, this approximation seems to work well. On the other hand, for charged lipid membranes it can be violated. The membrane can have a mixture of neutral and, for example, negatively charged hydrophilic heads. In a liquid membrane heads are mobile so that negative ones can accumulate near the positive $`Z`$-ion and push the neutral heads outside (see Fig. 8). Since the background charges are now closer to the counterion, one can immediately predict that the energy of the WC is lower and charge inversion is stronger than that for the case of an uniform distribution of negative heads.
To simplify the calculation of the free energy, and gain more physical insight in the problem, let us use the same transformation as in the beginning of section II, namely we simultaneously add uniform planar charge densities $`\sigma ^{}`$ and $`\sigma ^{}`$ to the plane. The first addition makes a neutral WC on the plane. While the second addition creates the two planar capacitors. The free energy can be written as the sum of the energy of WC and two capacitors, in the same way as Eq. (9). Therefore, $`\sigma ^{}`$ is given by Eq. (14).
We use below the Wigner-Seitz approximation to calculate $`\mu _{WC}`$. This approximation gives the energy per ion of WC as the energy of one Wigner-Seitz cell and neglects the quadrupole-quadrupole interaction between Wigner-Seitz cells. It provides 5% accuracy for the energy of the standard WC on an uniform immobile background (see Eq. (16)). In the case of mobile charges, as one sees from Fig. 7, the quadrupole moment of the Wigner-Seitz cell is even smaller than that for WC on an uniform background with the same average charge density $`\sigma `$. Therefore, in the case of mobile charge, the accuracy of the Wigner-Seitz approximation is even better.
For simplicity, we assume the Wigner-Seitz cell of a counterion is a disk with radius $`R=(\pi /n)^{1/2}`$. The negative heads concentrate around the counterion and make a negative disk with radius $`R_{}<R`$ and charge density $`\sigma _{}`$ where $`\sigma _{}=\sigma /n\pi R_{}^2\sigma `$. The rest of the cell is occupied by neutral heads (Fig. 5). The fraction of negative heads $`f^2=R_{}^2/R^2`$ is fixed for each membrane. The uniform charge case is recovered when there are no neutral heads so that $`R_{}=R`$ and $`f=1`$.
Let us consider the weak screening case $`r_sR`$. Under the transformation mentioned above, we add a disk with radius $`R`$, density $`\sigma ^{}`$ to the Wigner-Seitz cell to neutralize it. Now, the total energy of a Wigner-Seitz cell is the sum of the interactions of the $`Z`$-ion with two disks of radiuses $`R_{}`$ and $`R`$, the self energy of the two disks and the interaction between the disks:
$`\epsilon (n)`$ $`=`$ $`{\displaystyle \frac{2\pi Ze\sigma _{}R_{}}{D}}{\displaystyle \frac{2\pi Ze\sigma ^{}R}{D}}+{\displaystyle \frac{8\pi }{3}}{\displaystyle \frac{(\sigma _{})^2R_{}^3}{D}}`$ (95)
$`+{\displaystyle \frac{8\pi }{3}}{\displaystyle \frac{(\sigma ^{})^2R^3}{D}}+{\displaystyle _{(R_{})}}{\displaystyle _{(R)}}𝑑𝒓𝑑𝒓^{\mathbf{}}{\displaystyle \frac{\sigma _{}\sigma ^{}}{D|𝒓𝒓^{\mathbf{}}|}}.`$
The integrations in the last term are taken over the disks with radius $`R_{}`$ and $`R`$ respectively. This last term can be written as
$$_{(R_{})}_{(R)}𝑑𝒓𝑑𝒓^{\mathbf{}}\frac{\sigma _{}\sigma ^{}}{D|𝒓𝒓^{\mathbf{}}|}=\frac{2\pi \sigma _{}\sigma ^{}R_{}^3}{D}𝒢(f),$$
(96)
where $`𝒢(f)`$ is a function of $`f`$ only and can be evaluated numerically for each value of $`f`$ (it decreases monotonically from $`8/3`$ at $`f=1`$ to $`0`$ at $`f=0`$). Using $`Zen=\sigma +\sigma ^{}`$ and Eq. (96), one gets from Eq. (95):
$`\epsilon (n)`$ $`=`$ $`{\displaystyle \frac{(Ze)^2}{RD}}\left(2{\displaystyle \frac{8}{3\pi }}\right)+{\displaystyle \frac{2\pi \sigma ^2R^3}{D}}\left({\displaystyle \frac{4}{3f}}+{\displaystyle \frac{4}{3}}f𝒢(f)\right)`$ (98)
$`+{\displaystyle \frac{2\pi \sigma ZeR}{D}}\left(1{\displaystyle \frac{8}{3\pi }}{\displaystyle \frac{1}{f}}+{\displaystyle \frac{f𝒢(f)}{\pi }}\right).`$
The last two terms is the correction to $`\epsilon (n)`$ due to the mobility of the surface charge. In the uniform limit, $`f=1`$, $`𝒢(f)=8/3`$, these two terms vanish and one gets back the usual formula for the energy per ion of WC in Wigner-Seitz cell approximation, Eq. (15).
The chemical potential for a counterion in the mobile charge case, $`\mu _{WC,m}`$, can be easily calculated as
$$\mu _{WC,m}=\frac{[n\epsilon (n)]}{n}\frac{(Ze)^2}{RD}\left(2+\frac{1}{f}+\frac{4}{3\pi f}\frac{2f𝒢(f)}{\pi }\right).$$
(99)
Here $`\sigma `$ is approximated by $`Zen`$, because at $`r_sR`$, the ratio $`\sigma ^{}/\sigma 1`$.
The ratio between chemical potential $`\mu _{WC,m}`$ for the mobile charges and the chemical potential $`\mu _{WC}`$ for the immobile charges has been evaluated numerically as function of the fraction $`f^2`$ of the negative heads. The result is plotted in Fig. 9.
Obviously, as $`f`$ decreases, $`\mu _{WC,m}`$ grows as expected. According to Eq. (14) this means that the inversion ratio $`\sigma ^{}/\sigma `$ grows with decreasing $`f`$, too. We do not continue the plot in Fig. 9 to very small $`f^2`$ because in this case, the entropy of negative heads plays important role and screening by negative heads of the membrane can be described in Debye-Hückel approximation . We do not consider this regime here.
Let us now move to the limit of strong screening, $`r_sR`$. In this limit, it is more convenient to directly minimize the free energy, instead of using the capacitor model. Since $`r_sR`$, one needs to keep only the nearest neighbour interactions in the free energy. Assuming $`r_sR_{}`$, one can write the free energy per unit area as
$$F=2\pi \sigma _{}r_sZen+3n\left(Ze\frac{\sigma }{n}\right)^2\frac{\mathrm{exp}(A/r_s)}{A},$$
(100)
where $`Ze\sigma /n=\sigma ^{}/n`$ is the charge of one Wigner-Seitz cell. In Eq. (100), the first term is the interaction of $`Z`$-ion with the negative background (the disk with charge density $`\sigma _{}`$), the second term is the interaction between neighbouring Wigner-Seitz cells. As usual, the quadrupole-quadrupole interaction between Wigner-Seitz cells is neglected.
Minimizing the free energy (100) with respect to $`n`$, one gets $`Ar_s\mathrm{ln}(f^2\zeta )`$ and
$$\frac{\sigma ^{}}{\sigma }=\frac{2\pi \zeta }{\sqrt{3}\mathrm{ln}^2(f^2\zeta )},(\zeta 1).$$
(101)
where $`\zeta =Ze/\pi \sigma r_s^2`$. Comparing to Eq. (24), one can see that, as in the weak screening case, the inversion ratio increases due to the mobility of the surface charge.
Theory of this section is based on the assumption that the charge of $`Z`$-ion is so large that it is screened nonlinearly by the disk of opposite charge. One can easily generalize this calculations to rod-like polyelectrolytes and study the role of a similar stripe of positive hydrofilic heads attracted by strongly negative DNA. Note that the idea of nonlinear concentration of charge in membranes with two types of heads has been used recently in a theory of DNA-cationic lipid complexes .
## IX Conclusion
We would like to conclude with another general physical interpretation of the origin of charge inversion. To do so, let us begin with brief discussion of a separate physical problem, namely, let us imagine that, instead of a macroion, a neutral macroscopic metallic particle is suspended in water with $`Z`$\- and mono-valent ions. In this case, each ion creates an image charge of opposite sign inside the metal and thus attracts to the metal. Obviously, this effect is by a factor $`Z^2`$ stronger for $`Z`$-ions than for monovalent ones. While directly at the metal surface, energy of interaction of $`Z`$-ion with image, $`(Ze)^2/4a`$, is much larger than $`k_BT`$. Therefore $`Z`$-ions are strongly bound to the metallic surface, making it effectively charged, while monovalent ions are loosely correlated with the surface, providing for its screening over the distances of the order of $`r_s`$. We can determine the net charge of metallic particle with bound $`Z`$-ions using the ”capacitor model” discussed above. Namely, the attraction of the $`Z`$-ions to their images plays the role of correlation part of the chemical potential $`\mu _c`$ and provides for the voltage $`Ze/4a`$ which charges a ”capacitor” with the width $`r_s`$ between metal surface and the bulk solution. This leads to the result that metal surface is charged with the net charge density $`\sigma ^{}=Zen=Ze/(16\pi ar_s)`$. Note that metallic particle becomes charged due to interactions, or correlations, between $`Z`$-ions and their images, even though the particle itself was neutral in the first place.
Major results of the present paper can be now interpreted using a similar language of images . Although now we consider a macroion with an insulating body, it has some bare charge $`\sigma `$ on its surface, which leads to adsorption of certain amount of $`Z`$-ions. The layer of adsorbed $`Z`$-ions plays the role of a metal. Indeed, consider bringing a new $`Z`$-ion to the macroion surface which has already some bound $`Z`$-ions. New $`Z`$-ion repels nearest adsorbed ones, creating a correlation hole for himself. In other words, it creates an image with the opposite charge behind the surface. Image attracts the $`Z`$-ion, thus providing for the negative $`\mu _c`$ in Eq. (11) and therefore leading to the charge inversion.
The analogy between the adsorbed layer of $`Z`$-ions and a metal surface holds only at length scales larger than some characteristic length. In WC this latter scale is equal to Wigner-Seitz cell radius $`R`$. This is why for WC $`\mu _c(Ze)^2/R`$ (see Eq. (17)). To make $`|\mu _c|k_BT`$, small enough radius $`R`$ is needed. This explains why a significant bare charge $`\sigma `$ is necessary to initiate adsorbtion of $`Z`$-ions and to create a metallic layer with images which can lead to charge inversion. From formal point of view, charge inversion in this case can be characterized by the ratio $`\sigma ^{}/\sigma `$, as we did throughout the paper, while for a neutral metallic particle such ratio is infinite.
In this paper, we considered adsorption of rigid $`Z`$-ions with the shapes of either small spheres or thin rods. The concept of effective metallic surface and image based language is perfectly applicable in both cases. It appears also applicable to the other problem, not considered in this paper, namely, that of adsorption of a flexible polyelectrolyte on an oppositely charged dielectric macroion surface . To our mind, this idea was already implicitly used in Ref. , which assumes that Coulomb self-energy of a polyelectrolyte molecule in the adsorbed layer is negligible. This means that charge of the polyelectrolyte molecule is compensated by the correlation hole, or image. It is the image charge that attracts a flexible polyelectrolyte molecule to the surface. Interestingly, conformations of both the polymer molecule and its image change when the molecule approaches the surface.
A similar role of images and correlations is actually well known in the physics of metals. In the Thomas-Fermi approximation (which is similar to PB one) the work function of a metal is zero (the work function is an analog of $`\mu _c`$). The finite value of the work function is known to result from the exchange and correlation between electrons. For a leaving electron it can be interpreted as interaction with its image charge in the metal .
We believe that interaction with image or, in other words, lateral correlations of $`Z`$-ions in the adsorbed layer is the only possible reason for a charge inversion exceeding one $`Z`$-ion charge (of course, we mean here purely Coulomb systems and do not speak about cases when charge inversion is driven by other forces, such as, e.g., hydrophobicity). In the Poisson-Boltzmann approximation, when charge is smeared uniformly along the surface, no charging of neutral metal or overcharging of charged insulating plane is possible.
###### Acknowledgements.
We are grateful to R. Podgornik and I. Rouzina for useful discussions. This work was supported by NSF DMR-9985985.
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# Static magnetization induced by time-periodic fields with zero mean
## Abstract
We consider a single spin in a constant magnetic field or an anisotropy field. We show that additional external time-periodic fields with zero mean may generate nonzero time-averaged spin components which are vanishing for the time-averaged Hamiltonian. The reason is a lowering of the dynamical symmetry of the system. A harmonic signal with proper orientation is enough to display the effect. We analyze the problem both with and without dissipation, both for quantum spins ($`s=1/2,1`$) and classical spins. The results are of importance for controlling the system’s state using high or low frequency fields and for using new resonance techniques which probe internal system parameters, to name a few.
Usually nonzero averages of observables, which would be expected to be zero by symmetry considerations, are generated either by constant external fields, or by internal interactions which may lead to phase transitions. However as we will show below such a situation is also possible if we use time-periodic fields with zero mean. The general idea behind the following results is purely symmetry related, and thus it seems to be worthwhile to understand the mechanisms which may lead to nonzero averages if such fields are applied. This work is motivated by a recent paper where similar ideas have been used to explain the phenomenon of directed currents in driven systems. The essence of the present paper is that we can lower the symmetry of a given dynamical system by applying time-periodic fields with zero mean, i.e. that the time-averaged Hamiltonian displays symmetries which would imply zero averages for corresponding observables. It will be the symmetry breaking in the temporal evolution which induces nonzero averages.
Let us start our considerations with a model describing an $`s=1/2`$ spin in a constant field $`h_z=2`$ directed along the $`z`$-direction and a time-periodic field $`2h_x(t)`$ with period $`T`$ and zero mean directed along the $`x`$-direction. The Hamiltonian is given by $`H=h_zS_z+2h_x(t)S_x`$ (here $`S_{x,y,z}`$ are the spin component operators related to the corresponding Pauli matrices, e.g. ). For the moment we assume that $`|h_x(t)|2`$ and the frequency $`\omega =2\pi /T2`$. In that case we can use the adiabatic approximation and neglect Zener transitions. The two eigenvalues of $`H`$ for a given value of $`h_x`$ are $`\lambda _\pm =\pm \sqrt{1+h_x^2}`$. The expectation value for $`S_x`$ in these states is given by
$$S_x=\frac{h_x}{2\sqrt{1+h_x^2}}.$$
(1)
Now we assume that the spin is in any of the two states. Slow variation of $`h_x`$ in time will keep the system in that state. Let us average $`S_x`$ over one period of oscillation. Because $`S_x`$ is odd in $`h_x`$, we will obtain nonzero time averages for the $`x`$-component of the spin if e.g. $`_0^Th_x^3dt0`$. This is possible if $`h_x(t)`$ contains several harmonics (SH), e.g. $`h_x(t)=h_1\mathrm{cos}(\omega t)+h_2\mathrm{cos}(2\omega t+\xi )`$ (see also ). In that case in lowest order in $`h_1,h_2`$ we obtain $`S_x=\frac{3}{16}h_1^2h_2\mathrm{cos}\xi `$. We conclude this example with stating that it is possible to generate a nonzero average $`S_x`$ spin component by applying a permanent field in $`z`$-direction and a time-periodic field with SH and zero average in $`x`$-direction.
Let us relate the results from the example given above to symmetry considerations. The Hamiltonian $`H`$ should be a periodic function of time $`H(t)=H(t+T)`$. Instead of solving the time-dependent Schrödinger equation, which would bring us to the analysis of unitary Floquet matrices , we follow the density matrix approach, which is suitable since we want to average over different initial conditions and are thus facing the dynamics of mixed states. The density matrix $`\rho `$ satisfies the quantum Liouville equation
$$\frac{\rho }{t}=\mathrm{i}[H,\rho ]\nu (\rho \rho _\beta )$$
(2)
where $`[A,B]=ABBA`$, $`\rho _\beta `$ is some equilibrium density matrix parametrized by the inverse temperature $`\beta `$ and $`\nu `$ is a phenomenological parameter measuring the coupling strength of the system described by $`H`$ to some environment. Note that $`\nu `$ is the characteristic inverse relaxation time of $`H`$ in the environmental bath.
Let us further define $`H_0=1/T_0^TH(t)dt`$ and $`H_1(t)H(t)H_0`$. Note that $`_0^TH_1(t)dt=0`$. Then we may choose $`\rho _\beta =\frac{1}{Z}\mathrm{e}^{\beta H_0}`$ with $`Z=\mathrm{Tr}(\mathrm{e}^{\beta H_0})`$. We define the value $`\overline{A}(t)`$ of an observable characterized by the operator $`A`$ as $`\overline{A}(t)=\mathrm{Tr}(A\rho (t))`$. The time average of $`\overline{A}(t)`$ shall be defined as $`\stackrel{~}{A}=lim_t^{}\mathrm{}\frac{1}{t^{}}_0^t^{}\overline{A}(t)dt`$. The averaged attenuation power (the rate of energy transfer from the time-periodic field to the heat bath) is given by $`W=\nu (\stackrel{~}{H}_0\mathrm{Tr}(H_0\rho _\beta ))`$.
We chose the relaxation term in (2) in an oversimplified form. There are many theories which exploit different concrete relaxation mechanisms (e.g. and references therein). The reason for choosing (2) instead is that it allows to discuss the following symmetry breaking without entering the details of the concrete dissipation mechanism. In other words, we deliberately choose the simplest dissipation term which conserves all symmetries of our dynamical system except time reversal.
Equation (2) is a linear equation for the matrix coefficients of $`\rho `$ with inhomogeneous terms due to $`\rho _\beta `$. The general solution is given by a sum of the general solution of the homogeneous equation (put $`\rho _\beta =0`$ in (2)) and a particular solution of the full equation. Since the homogeneous solution for $`\nu =0`$ is given by some unitary time evolution, $`\nu >0`$ will cause all solutions of the homogeneous equation to decay to zero for infinite time. For $`t1/\nu `$ any particular solution of the inhomogeneous equation trends to a unique time-periodic solution - the attractor of (2). This allows us to choose any (reasonable) initial condition $`\rho (t=0)`$. If $`H`$, $`\rho (t=0)`$ and $`\rho _\beta `$ are invariant under certain unitary transformations, it immediately follows that $`\rho (t)`$ keeps those symmetries, and consequently the attractor will have the same symmetries too. For large temperatures $`\rho _\beta `$ is approaching the unity matrix (up to some factor). Consequently in that limit, whatever the time dependence of $`H(t)`$, the solution of (2) will approach $`\rho _\beta `$. Finally we note that due to $`\mathrm{Tr}\rho _\beta =1`$ any choice of $`\rho (t=0)`$ with $`\mathrm{Tr}\rho (t=0)=1`$ implies $`\mathrm{Tr}\rho (t)=1`$ for all $`t`$.
Let us consider (2) for
$$H=h_0S_z+h(t)(\alpha S_x+\gamma S_z)$$
(3)
where $`\alpha =\mathrm{sin}(\varphi )`$ and $`\gamma =\mathrm{cos}(\varphi )`$. This model describes a spin in a constant magnetic field pointing in the $`z`$-direction, under the influence of an additional time-periodic field $`h(t)=h(t+T)`$. This oscillating field should have zero mean: $`_0^Th(t)dt=0`$. Let us define $`h(t)`$ having $`T_a`$ symmetry if $`h(t)=h(t)h_a(t)`$, $`T_s`$ symmetry if $`h(t)=h(t)h_s(t)`$, and $`T_{sh}`$ symmetry if $`h(t)=h(t+T/2)h_{sh}(t)`$ (note that in the two first cases any argument shift is allowed, so that e.g. $`h(t)=\mathrm{cos}(t+\mu )`$ posesses all three symmetries). For a monochromatic field (MCF) $`h(t)`$ and $`\varphi =\pi /2`$ (3) is the classical setup for performing magnetic resonance (MR) experiments ,.
For the $`s=\frac{1}{2}`$ case the spin component operators are given by the Pauli matrices: $`S_{x,y,z}=\frac{1}{2}\sigma _{x,y,z}`$. The density matrix $`\rho `$ has three independent real variables. Using the variables $`\overline{S}_{x,y,z}`$ we find
$`\dot{\overline{S}}_x=(h_0+\gamma h(t))\overline{S}_y\nu \overline{S}_x`$ (4)
$`\dot{\overline{S}}_y=\alpha h(t)\overline{S}_z(h_0+\gamma h(t))\overline{S}_x\nu \overline{S}_y`$ (5)
$`\dot{\overline{S}}_z=\alpha h(t)\overline{S}_y\nu (\overline{S}_zC)`$ (6)
where $`C=1/2\mathrm{tanh}(h_0\beta /2)`$. Note that the obtained set of equations for $`\nu =0`$ is equivalent to the Heisenberg equations for the operators $`S_{x,y,z}`$ and thus also to the equations of motion for a classical spin. In fact (4)-(6) is a particular case of the Bloch equations ,.
Let us discuss the symmetries of (4)-(6) which conserve $`H_0`$, i.e. $`\overline{S}_zS_z`$. Consider the case $`\gamma =0`$: if $`h(t)h_{sh}(t)`$ then a symmetry operation $`Q_1`$ is $`\overline{S}_x\overline{S}_x,\overline{S}_y\overline{S}_y,\overline{S}_z\overline{S}_z,tt+T/2`$. If $`Q_1`$ holds we conclude that $`\stackrel{~}{S}_x=\stackrel{~}{S}_y=0`$, while $`\stackrel{~}{S}_z`$ may be nonzero. Consider $`\gamma =0`$ and $`\nu =0`$: if $`h(t)h_a(t)`$ then a symmetry operation $`Q_2`$ is $`\overline{S}_x\overline{S}_x,\overline{S}_y\overline{S}_y,\overline{S}_z\overline{S}_z,tt`$. If $`Q_2`$ holds it follows $`\stackrel{~}{S}_x=0`$, while $`\stackrel{~}{S}_{y,z}`$ may be nonzero. Finally for $`\nu =0`$ and $`h(t)h_s(t)`$ a symmetry operation $`Q_3`$ is $`\overline{S}_x\overline{S}_x,\overline{S}_y\overline{S}_y,\overline{S}_z\overline{S}_z,tt`$. If $`Q_3`$ holds it follows $`\stackrel{~}{S}_y=0`$, while $`\stackrel{~}{S}_{x,z}`$ may be nonzero.
Let us note some consequences. If we choose $`h(t)=h_1\mathrm{cos}(\omega t)`$, then the classical MR setup with $`\gamma =0`$ ($`Q_1`$) yields nonzero values for $`\stackrel{~}{S}_z`$ only . If the probing field is not perpendicular to the $`z`$-axis ($`\gamma 0`$), nonzero values for $`\stackrel{~}{S}_x`$ and $`\stackrel{~}{S}_y`$ appear as well. $`\stackrel{~}{S}_y`$ will vanish in the limit of zero coupling to the environment $`\nu 0`$ ($`Q_3`$), so that this average can be used to measure the coupling strength. Applying e.g. $`h(t)=h_1\mathrm{sin}(\omega t)+h_2\mathrm{sin}(2\omega t)`$ (having $`h_a`$ symmetry but not $`h_{sh}`$ and $`h_s`$ one) we can suppress the value of $`\stackrel{~}{S}_x`$ relatively to $`\stackrel{~}{S}_y`$ for $`\gamma 0`$ and $`\nu 0`$ keeping $`\stackrel{~}{S}_y`$ finite ($`Q_2`$)!
Analytical solutions to (4)-(6) can be found e.g. for large $`\nu 1`$. Expanding in $`1/\nu `$ and averaging over time we find in lowest orders
$`\stackrel{~}{S}_x=C\alpha \gamma h^2{\displaystyle \frac{1}{\nu ^2}}C\alpha (\gamma h\ddot{h}+3\gamma h_0^2h^2+`$ (7)
$`(\alpha ^2+3\gamma ^2)h_0h^3+\gamma (\gamma ^2+\alpha ^2)h^4){\displaystyle \frac{1}{\nu ^4}}+O({\displaystyle \frac{1}{\nu ^5}})`$ (8)
$`\stackrel{~}{S}_y=C\alpha \left[2\gamma h_0h^2+(\gamma ^2+\alpha ^2)h^3\right]{\displaystyle \frac{1}{\nu ^3}}+O({\displaystyle \frac{1}{\nu ^5}})`$ (9)
where $`f(t)=\frac{1}{T}_0^Tf(t)dt`$. It is easy to cross check that all symmetry statements from above are correct. Nonzero values for $`h^3`$ can be obtained e.g. with $`h(t)=h_1\mathrm{sin}(\omega t)+h_2\mathrm{sin}(2\omega t+\xi )`$ for $`\xi 0,\pi `$ (see also ).
In Fig.1 we show the dependence of $`\stackrel{~}{S}_{x,y,z}`$ on $`\omega `$ for $`h(t)=\sqrt{2}\mathrm{cos}\omega t`$, $`\varphi =\pi /4`$, $`h_0=3`$, $`\nu =0.1`$ and $`\beta =10`$. The time-periodic field has a large amplitude compared to typical MR setups . This causes the $`\stackrel{~}{S}_z`$ curve to show a rather broad peak at $`\omega h_0`$ \- the position of the expected MR resonance. However we also observe sattelite peaks at lower frequencies which are clearly related to the variations of nonzero $`\stackrel{~}{S}_{x,y}`$ (for convenience these averages are scaled by a factor of 10 in Fig.1). In fact the positions of the sattelite peaks are subharmonics of the main resonance. The dependence of $`\stackrel{~}{S}_x`$ and $`\stackrel{~}{S}_y`$ on $`\omega `$ shows rather complex structures. We find that typically the dependence of these averages on $`\omega `$ becomes oscillatory for small $`\omega h_0`$, whereas large $`\omega `$ values yield smooth decay curves. Note also that these averages stay nonzero down to small frequencies in accord with the adiabatic example from above. Also important is to notice that the fluctuations of $`\overline{S}_x`$ and $`\overline{S}_y`$ around their mean values may happen with amplitudes being one order of magnitude larger than the mean values (see inset in Fig.1 ).
The above results hold also for larger spins. To show that they also hold for internal anisotropy fields rather than external fields, we consider a spin with $`s=1`$ and the Hamiltonian
$$H=S_z^2+h(t)(\alpha S_x+\gamma S_z)$$
(10)
which describes a spin with an anisotropy along the $`z`$-axis ($`S_z^2`$) under the influence of an external magnetic field $`h(t)`$ parallel to the $`xz`$ plane. The magnetic field is again time-periodic with period $`T`$ and has zero mean. The $`3\times 3`$ hermitian density matrix $`\rho `$ has 8 independent real parameters. Since $`H`$ in (10) is a real symmetric matrix, we can define $`\rho =R+\mathrm{i}I`$ where $`R`$ is a real symmetric matrix and $`I`$ a real antisymmetric one. Noting that also $`\rho _\beta `$ is a real diagonal matrix, (2) can be rewritten as
$`{\displaystyle \frac{R}{t}}=[H,I]\nu (R\rho _\beta )`$ (11)
$`{\displaystyle \frac{I}{t}}=[H,R]\nu I`$ (12)
It follows $`\overline{S}_x=\sqrt{2}(R(1,2)+R(2,3))`$, $`\overline{S}_y=\sqrt{2}(I(1,2)+I(2,3))`$ and $`\overline{S}_z=R(1,1)R(3,3)`$ . Using the abbrevations $`P_x=\sqrt{2}(R(1,2)R(2,3))`$, $`P_y=\sqrt{2}(I(1,2)I(2,3))`$, $`P_z=R(1,1)+R(3,3)`$, $`R_{13}=\sqrt{2}R(1,3)`$, $`I_{13}=\sqrt{2}I(1,3)`$, $`R_{22}=\sqrt{2}R(2,2)`$, $`D^1=1+2\mathrm{e}^\beta `$ and $`F^1=2+\mathrm{e}^\beta `$ the equations of motion become
$`\dot{\overline{S}}_x=P_y+\gamma h\overline{S}_y\nu \overline{S}_x`$ (13)
$`\dot{P}_x=\overline{S}_y\gamma hP_y+\sqrt{2}\alpha hI_{13}\nu P_x`$ (14)
$`\dot{\overline{S}}_y=P_x\gamma h\overline{S}_x+\alpha h\overline{S}_z\nu \overline{S}_y`$ (15)
$`\dot{P}_y=\overline{S}_x+\gamma hP_x+\alpha h\left[\sqrt{2}R_{22}P_z\sqrt{2}R_{13}\right]\nu P_y`$ (16)
$`\dot{\overline{S}}_z=\alpha h\overline{S}_y\nu \overline{S}_z`$ (17)
$`\dot{P}_z=\alpha hP_y\nu \left[P_z2F\right]`$ (18)
$`\dot{R}_{13}=2\gamma hI_{13}+\sqrt{2}\alpha hP_y\nu R_{13}`$ (19)
$`\dot{I}_{13}=2\gamma hR_{13}\sqrt{2}\alpha hP_x\nu I_{13}`$ (20)
$`\dot{R}_{22}=\sqrt{2}\alpha hP_y\nu \left[R_{22}\sqrt{2}D\right]`$ (21)
These equations conserve the trace $`\mathrm{Tr}\rho P_z+R_{22}/\sqrt{2}=1`$.
Now we can discuss the symmetries of (17) which change the sign of $`\overline{S}`$. Two of them hold only for $`\nu =0`$. First, if $`h(t)h_a(t)`$, then the equations are invariant under change of sign of the variables $`t,\overline{S}_x,\overline{S}_y,\overline{S}_z`$ (leaving all other variables unchanged). A second case takes place if $`h(t)h_s(t)`$. Then changing the sign of $`t,\overline{S}_y,P_y,I_{13}`$ (leaving all other variables unchanged) is an operation which keeps equations (17) invariant. These two cases imply that if $`h(t)`$ is antisymmetric, then for vanishing dissipation $`\nu 0`$ $`\stackrel{~}{S}_{x,y,z}0`$, while for symmetric $`h(t)`$ the same limit provides a vanishing of the $`y`$-component only $`\stackrel{~}{S}_y0`$.
For the general case $`\nu 0`$ two more symmetries may take place. If $`\gamma =0`$ (the field $`h(t)`$ acts perpendicularly to the anisotropy axis $`z`$), changing the sign of $`\overline{S}_y,\overline{S}_z,P_x,I_{13}`$ (and keeping all others) leaves (17) invariant. Finally if $`h(t)h_{sh}(t)`$, the shift $`tt+T/2`$ and simultaneous change of sign of the variables $`\overline{S}_x,\overline{S}_z,P_y,I_{13}`$ do not change the equations. It follows that $`\stackrel{~}{S}_y=\stackrel{~}{S}_z=0`$ for $`\gamma =0`$ and $`\stackrel{~}{S}_x=\stackrel{~}{S}_z=0`$ for $`h(t)`$ having shift symmetry.
It is interesting to note that for a MCF $`h(t)=\mathrm{cos}\omega t`$ and $`\nu 0`$, $`\gamma 0`$ the spin will point on average in $`y`$ direction, i.e. perpendicular to the plane spanned by the driving field and the local anisotropy axis! In Fig.2 we plot the dependence of $`\stackrel{~}{S}_y`$ on $`\omega `$ for this case ($`\beta =10`$, $`\nu =0.1`$, $`\gamma =\alpha =1`$), which confirms the symmetry considerations. Note that $`\stackrel{~}{S}_x`$ and $`\stackrel{~}{S}_z`$ are less than $`10^8`$ as found in the numerical studies.
To conclude this case we remark that it is again an easy task to perform expansions in $`1/\nu `$ for large $`\nu `$ values as shown above for the $`s=1/2`$ case. The resulting expressions also confirm the symmetry considerations.
So far we have discussed the results for quantum spin systems. It is also possible to analyze corresponding classical systems. E.g. the classical equations for (10) are given by
$`\dot{s}_x=2s_zs_y+\gamma hs_y`$ (22)
$`\dot{s}_y=2s_zs_x+h\left(\alpha s_z\gamma s_x\right)`$ (23)
$`\dot{s}_z=\alpha hs_y`$ (24)
Let us discuss the symmetry properties of (22)-(24). We denote on the left part the condition and on the right part the symmetry operations which leave the equations of motion invariant (note that we list only those variables which have to be changed):
$`\gamma =0`$ $`:`$ $`(s_y,s_z)(s_y,s_z)`$ (25)
$`T_a`$ $`:`$ $`tt,(s_x,s_y,s_z)(s_x,s_y,s_z)`$ (26)
$`T_s`$ $`:`$ $`tt,s_ys_y`$ (27)
$`T_{sh}`$ $`:`$ $`tt+{\displaystyle \frac{T}{2}},(s_x,s_z)(s_x,s_z)`$ (28)
If we add dissipation terms, these terms will break time reversal symmetry, and we are left only with (25) and (28). All of the above statements for the quantum system can be recovered. Especially nonzero dissipation and $`\gamma ,\alpha 0`$ lead to nonvanishing magnetization along the $`y`$-axis, even for MCF.
Let us summarize the presented results. We have shown that time-periodic magnetic fields with zero mean may induce nonzero averages of spin components which would be strictly zero in the absence of these fields. The spin is simultaneously experiencing some local anistropy field or simply an external constant field. In addition the spin is coupled to some thermal environment characterized by some finite temperature and a characteristic relaxation time . The reasoning follows symmetry considerations of the dynamical equations. In the case of a classical spin these equations formally coincide with the Heisenberg equations for the quantum spin operators. In the quantum case we instead solve the (purely linear!) equations of motion for the independent components of the density matrix. Remarkably the symmetry properties obtained from both approaches coincide.
The quantum approach shows that for infinite temperatures all spin component averages will vanish. This follows from $`\rho (t\mathrm{};\beta 0)=\rho _{\beta 0}`$ and $`\mathrm{Tr}S_{x,y,z}=0`$.
For the spin $`1/2`$ case we proposed a MR experiment to observe the effect. One should choose the time-periodic magnetic field to be not perpendicular to the static magnetic field. Further the amplitude of the time-periodic field should be not too small such that the generated $`\stackrel{~}{S}_x`$ and $`\stackrel{~}{S}_y`$ components are measurable. The attenuation spectrum should show resonances located at subharmonics of the original resonance. The intensity of the sattelite peaks is a function of both the angle between both fields and the inverse relaxation time $`\nu `$.
Experiments which probe directly the nonzero spin components can be performed by adding yet another probing field to the system, and varying its frequency while keeping the frequency of the original probing field. This will be studied in detail in future work.
Acknowledgements.
We thank A. Bartl, P. Fulde, D. A. Garanin and A. Latz for useful discussions.
Figure captions.
Fig.1. $`10\stackrel{~}{S}_x`$ (solid), $`10\stackrel{~}{S}_y`$ (dashed) and $`\stackrel{~}{S}_z`$ (dotted) as functions of $`\omega `$ (see text for parameters).
Inset: $`\overline{S}_{x,y,z}`$ versus time for one period of $`h(t)`$ at $`\omega =1.5`$ (same line codes as in Fig.1). Note that functions are not scaled here!
Fig.2.
$`\stackrel{~}{S}_y`$ as a function of $`\omega `$ (see text for parameters).
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# Off-Shell Formulation of Simple Supersymmetric Yang-Mills
## 1 Introduction
Day after day, supersymmetry consolidates its position in theoretical physics. Even if it was introduced more than 25 years ago, there are still problems with the geometric basis of extended $`\left(N>1\right)`$ supersymmetry. The situation of the extended superspace is far less satisfactory than the original N=1 superspace. At the level of the algebra the on-shell formalism closes up to modulo of the classical equations of motion. This fact seems odd at the quantum level since the equations of motion receive loop corrections<sup>1</sup><sup>1</sup>1Also, the supersymmetry transformations receive corrections and one should test the closure of the algebra order by order in perturbation theory.. The superspace introduces an elegant supermanifold with different enlarged superconnections, where some are truly integrable in the sense of having zero supercurvature. In principle, the extended superspace should be a very powerful tool for quantum calculations.
Before starting, we feel obliged to mention something about the history of the following conjecture: Ring Division Algebras $`𝕂`${ real $`,`$complex $`,`$quaternions $`,`$octonions $`𝕆`$ } are relevant to simple supersymmetric Yang-Mills. The first hint, as mentioned by Schwarz comes from the number of propagating Bose and Fermi degrees of freedom which is one for $`d=3`$, two for $`d=4`$, four for $`d=6`$ and eight for $`d=10`$ suggesting a correspondence with real $``$, complex $``$, quaternions $``$ and octonions $`𝕆`$. Kugo and Townsend investigated in detail the relationship between $`𝕂`$ and the irreducible spinorial representation of the Lorentz group in $`d=3,4,6,10`$, building upon the following chain of isomorphisms
$`so(2,1)`$ $``$ $`sl(2,)`$
$`so(3,1)`$ $``$ $`sl(2,)`$
$`so(5,1)`$ $``$ $`sl(2,).`$
They conjectured that $`so(9,1)sl(2,𝕆)`$, the correct relation turned out to be
$$so(9,1)sl(2,𝕆)G_2$$
as has been shown by Chung and Sudbery , i.e. the dimension of $`Sl(2,𝕆)`$ is 31. Also in , a quaternionic treatment of the $`d=6`$ case is presented. Later, Evans made a systematic investigation of the relationship between SSYM and ring division algebra in a couple of papers. In the first , he simplified the construction of SSYM by proving a very important identity between gamma matrices by using the intrinsic triality of ring division algebra instead of the “tour de force” used originally by Brink, Scherck and Schwarz via Fierz identities generalized to $`d>4`$ dimensions. Then, in the second paper , Evans made the connection even clearer by showing how the auxiliary fields are really related to ring division algebras. For $`d=3,4,6,10`$ we need $`k=0,1,3,7`$ auxiliary fields respectively. An alternative approach for the octonionic case was introduced by Berkovits who invented a larger supersymmetric transformation called generalized supersymmetry in . There has also been a twistor attempt by Bengtsson and Cederwall . For more references about the octonionic case and ten dimensional physics one may consult references in and its extension to p-branes by Belecowe and Duff . The early work of Nilsson may be relevant too.
As a first step towards an extended superspace, we address the point of the algebraic auxiliary fields for simple N=1 supersymmetric Yang-Mills (SSYM) definable only in $`d=3,4,6`$ and $`10`$ dimensions . The important point is: While the physical fields couple to ring division left action the auxiliary ones couple to right action (or vice versa). To admit a closed off-shell supersymmetric algebra, left and right action must commute i.e. we should have a parallelizable associative algebra. For d = 6, quaternions work fine but for d = 10, the only associative seven dimensional algebra that is known is the soft seven sphere. We shall show below how this works. In this work, we use the same symbols $`\left(\text{left action}𝔼_j\text{, right action}1|𝔼_j\right)`$ for either complex, quaternionic or octonionic numbers and each case should be distinguished by the range of the indices $`j`$ which run from $`1`$ to $`(1,3,7)`$ for complex, quaternions and octonions respectively.
In the second section, we review the relation between hypercomplex structure and Clifford algebra. The auxliary fields problem in 6 dimensions is presented in the third section. While section four is devoted to the ten dimensional case. The last section contains some superspace hints.
## 2 Hypercomplex Structure and Pure Spinors
Everything starts from Clifford algebra, so let’s review quickly the connection between hypercomplex structures and our gamma matrices in $`d=3,4,6,10`$. The solution is encoded completely in our $`\mathrm{\Gamma }_M`$. Algebraically, we can construct a Clifford algebra directly from complex, quaternions and octonions over $`1,3,7`$ Euclidean space which can be extended easily to a representation of the minimal irreducible spinorial subspaces in $`d=3,4,6`$ and $`10`$ Minkowskian space-time. Consider the set of matrices $`\left\{𝔼_j\right\}`$ for the following three different cases:
* the canonical complex structure over $`^2`$ is just $`2\times 2`$ matrix $`𝔼_1`$
$$e_1𝔼_1=(\delta _{0\mu }\delta _{1\nu }\delta _{0\nu }\delta _{1\mu });\left(𝔼_1\right)=\mathrm{𝟏}_\mu ;\mu =0,1,$$
(1)
by $`\mathrm{𝟏}_\mu `$ we always mean an ($`\mu \times \mu `$) unit matrix.
* the canonical left quaternionic structures over $`^4`$ are
$$(𝔼_j)_{\mu \nu }=(\delta _{0\mu }\delta _{j\nu }\delta _{0\nu }\delta _{j\mu }ϵ_{j\mu \nu })\mu ,\nu =\mathrm{0..3};j,k,h=\mathrm{1..3},$$
(2)
and
$$𝔼_j𝔼_k=(\delta _{jk}\mathrm{𝟏}_\mu +ϵ_{jkh}𝔼_h),$$
(3)
where $`ϵ_{jkh}`$ is the standard Levi-Civita symbol. Using Rotelli’s notation for right action, the canonical right quaternionic structures are,
$$(1|𝔼_j)_{\mu \nu }=(\delta _{0\mu }\delta _{j\nu }\delta _{0\nu }\delta _{j\mu }+ϵ_{j\mu \nu }),$$
(4)
and
$$1|𝔼_j1|𝔼_k=(\delta _{jk}\mathrm{𝟏}_\mu ϵ_{jkh}1|𝔼_h).$$
(5)
Let’s put these quaternionic structures into a form that can be recognized by physicsits
$$\left(𝔼_j\right)_{\mu \nu }=(1|𝔼_j)_{\mu \nu }=ϵ_{j\mu \nu }if\mu ,\nu =1,2,3.\left(𝔼_j\right)_{00}=(1|𝔼_j)_{00}=0.$$
$$\left(𝔼_j\right)_{0\nu }=\left(1|𝔼_j\right)_{0\nu }=\delta _{j\nu },\left(𝔼_j\right)_{\mu 0}=\left(1|𝔼_j\right)_{\mu 0}=\delta _{j\mu },$$
(6)
such mathematical quaternionic structures are well known in physics as the ’t Hooft eta symbols. We can check that
$$\left\{𝔼_{i,}𝔼_j\right\}=\left\{1\right|𝔼_i,1|𝔼_i\}=2\delta _{ij}\mathrm{𝟏}_4,$$
$`[𝔼_j,𝔼_k]`$ $`=`$ $`ϵ_{jkh}𝔼_h,`$
$`[1|𝔼_j,1|𝔼_k]`$ $`=`$ $`ϵ_{jkh}1|𝔼_h,`$ (7)
and the very important formula
$$[𝔼_i,1|𝔼_j]=0,$$
(8)
i.e. left and right quaternionic actions commute.
For octonions, the story is quite different, as they are non-associative. But as it is well known, for any Lie algebra the structure constants are proportional to the constant torsion over the group manifold whereas the torsion over the seven sphere $`S^7`$ varies from one point to another . The only way to solve these problems is to use the $`S^7`$ as an associative soft Lie algebra<sup>2</sup><sup>2</sup>2Soft Lie algebra is an algebra with structure functions instead of structure constants . as had been proposed by Englert, Servin, Troost, Van Proeyen and Spindel which can be derived from octonions (Look to for a full algebraic investigation of the soft seven sphere). For a generic octonionic number,
$$\phi =\phi ^\mu e_\mu =\phi _0e_0+\phi _ie_{i,}\left(\begin{array}{c}\phi _0\\ \phi _1\\ \phi _2\\ \phi _3\\ \phi _4\\ \phi _5\\ \phi _6\\ \phi _7\end{array}\right),\begin{array}{ccc}\mu ,\nu =\mathrm{0..7}\hfill & ,\hfill & j,k,h=\mathrm{1..7}\hfill \\ & & \\ & \phi ^\mu \hfill & \end{array},$$
(9)
such that $`e_0=1`$ and the other seven imaginary units satisfy $`e_ie_j=\delta _{ij}+f_{ijk}e_k[e_i,e_j]=2f_{ijk}e_k`$ where $`f_{ijk}`$ is completely antisymmetric and equals one for any of the following three-cycles (123), (145), (246), (347), (176), (257), (365). To construct the soft seven sphere Lie algebra, we just have to define the direction of action, for left and right action, we have
$`\delta _i\phi `$ $`=`$ $`e_i\phi ,`$
$`1|\delta _i\phi `$ $`=`$ $`\phi e_i,`$ (10)
then after simple calculations, we find
$`[\delta _j,\delta _k]`$ $`=`$ $`2f_{jkh}\delta _h2[\delta _j,1|\delta _k],`$
$`[1|\delta _j,1|\delta _k]`$ $`=`$ $`2f_{jkh}1|\delta _h2[\delta _j,1|\delta _k],`$
$`\{\delta _j,\delta _k\}`$ $`=`$ $`2\delta _{jk},`$
$`\{1|\delta _j,1|\delta _k\}`$ $`=`$ $`2\delta _{jk},`$ (11)
which are isomorphic to the following set$`\{𝔼_j,1|𝔼_j\}`$ of $`8\times 8`$ matrices,
$$\begin{array}{ccccc}\delta _j& & (𝔼_j)_{\mu \nu }& =& \delta _{0\mu }\delta _{j\nu }\delta _{0\nu }\delta _{j\mu }f_{j\mu \nu },\\ 1|\delta _j& & (1|𝔼_j)_{\mu \nu }& =& \delta _{0\mu }\delta _{j\nu }\delta _{0\nu }\delta _{j\mu }+f_{j\mu \nu },\end{array}$$
(12)
satisfying the algebra
$`[𝔼_j,𝔼_k]`$ $`=`$ $`2f_{jkh}𝔼_h2[𝔼_j,1|𝔼_k],`$
$`[1|𝔼_j,1|𝔼_k]`$ $`=`$ $`2f_{jkh}1|𝔼_h2[𝔼_j,1|𝔼_k],`$
$`\{𝔼_j,𝔼_k\}`$ $`=`$ $`2\delta _{jk},`$
$`\{1|𝔼_j,1|𝔼_k\}`$ $`=`$ $`2\delta _{jk},`$ (13)
they don’t close a Lie algebra but they close a soft Lie algebra defined by
$$[\delta _j,\delta _k]\phi 2f_{jkh}^{(+)}(\phi )e_h\phi [𝔼_j,𝔼_k]\phi =2f_{jkh}^{(+)}𝔼_h\phi ,$$
$$[1|\delta _j,1|\delta _k]\phi 2f_{jkh}^{()}(\phi )\phi e_h[1|𝔼_j,1|𝔼_k]\phi =2f_{jkh}^{()}1|𝔼_h\phi ,$$
(14)
where $`f_{jkh}^{(\pm )}\left(\phi \right)`$ are the left and right parallelizable torsion. One can check that our $`𝔼_i`$ defines what Cartan calls pure spinors ,
$$\phi ^t𝔼_i\phi =0$$
(15)
thus
$$f_{ijk}^{(+)}\left(\phi \right)=\frac{\phi ^t\left(𝔼_k𝔼_i𝔼_j\right)\phi }{r^2}.$$
(16)
and
$$f_{ijk}^{()}\left(\phi \right)=\frac{\phi ^t\left(1|𝔼_k\mathrm{\hspace{0.33em}\hspace{0.33em}1}|𝔼_i\mathrm{\hspace{0.33em}\hspace{0.33em}1}|𝔼_j\right)\phi }{r^2}.$$
(17)
where
$$\phi ^t\phi =r^2.$$
(18)
There is another interesting and very important property to note
$$\phi ^t[𝔼_i,1|𝔼_j]\phi =0$$
(19)
which may be the generalization of the standard Lie algebra relation, left and right action commute everywhere over the group manifold.
We close the algebra pointwisely using structure functions $`f_{ijk}(\phi )`$ instead of structure constants $`f_{ijk}`$ where $`\phi `$ may be considered as a coordinate system for an internal $`S^7`$ manifold not the space-time $`x`$ and they don’t mix
$$\frac{x}{\phi }=\frac{\phi }{x}=0.$$
(20)
Apart from the commutation of left and right actions, there are some other useful identities satisfied by our $`(𝔼_j,1|𝔼_j)`$ quaternionic or octonionic structures, they are
$`(𝔼_k)_{\mu \nu }(𝔼_j)_{\lambda \nu }+(𝔼_j)_{\mu \nu }(𝔼_k)_{\lambda \nu }`$ $`=`$ $`2\delta _{kj}\delta _{\mu \lambda },`$
$`(𝔼_k)_{\mu \nu }(𝔼_j)_{\mu \lambda }+(𝔼_j)_{\mu \nu }(𝔼_k)_{\mu \lambda }`$ $`=`$ $`2\delta _{kj}\delta _{\nu \lambda },`$
$`(𝔼_k)_{\mu \nu }(𝔼_k)_{\lambda \zeta }+(𝔼_k)_{\lambda \nu }(𝔼_k)_{\mu \zeta }`$ $`=`$ $`2\delta _{\mu \lambda }\delta _{\nu \zeta },`$ (21)
and the same holds equally well for $`(1|𝔼_j)`$, as had been noticed by Evans , they are direct consequences of the ring division triality.
Now, we have all the needed ingredients to construct our real universal $`\left(\mathrm{\Gamma }_M\right)_{ab}`$ matrices with spinorial lower indices $`a,b`$ of range the double of the $`\mu `$. For Minkowskian metric of signature $`\eta (,+,\mathrm{},+)`$, in $`d=3,4,6`$ and $`10`$, $`a,b=\mathrm{0..2}\mu +1`$, for simplicity, we use symmetric $`\mathrm{\Gamma }_M`$,
$$\begin{array}{ccc}_{\left(\mathrm{\Gamma }_j\right)_{ab}=\left(\begin{array}{cc}0& 𝔼_j\\ 𝔼_j& 0\end{array}\right)}\hfill & & _{\left(1|\mathrm{\Gamma }_j\right)_{ab}=\left(\begin{array}{cc}0& 1|𝔼_\mu \\ 1|𝔼_j& 0\end{array}\right),}\hfill \\ _{\left(\mathrm{\Gamma }_0\right)_{ab}=\left(\begin{array}{cc}\mathrm{𝟏}_\mu & 0\\ 0& \mathrm{𝟏}_\mu \end{array}\right);}\hfill & _{\left(\mathrm{\Gamma }_{d2}\right)=\left(\begin{array}{cc}0& \mathrm{𝟏}_\mu \\ \mathrm{𝟏}_\mu & 0\end{array}\right)}\hfill & _{\left(\mathrm{\Gamma }_{d1}\right)_{ab}=\left(\begin{array}{cc}\mathrm{𝟏}_\mu & 0\\ 0& \mathrm{𝟏}_\mu \end{array}\right),}\hfill \end{array}$$
(22)
The corresponding higher indices $`\left(\stackrel{~}{\mathrm{\Gamma }}\right)^{ab}`$’s are
$$\left(\stackrel{~}{\mathrm{\Gamma }}_0\right)^{ab}=\left(\mathrm{\Gamma }_0\right)_{ab}and\left(\stackrel{~}{\mathrm{\Gamma }}\right)^{ab}=\left(\mathrm{\Gamma }\right)_{ab}.$$
(23)
As a result, we find
$$\mathrm{\Gamma }^M\stackrel{~}{\mathrm{\Gamma }}^N+\mathrm{\Gamma }^N\stackrel{~}{\mathrm{\Gamma }}^M=1|\mathrm{\Gamma }^M1|\stackrel{~}{\mathrm{\Gamma }}^N+1|\mathrm{\Gamma }^N1|\stackrel{~}{\mathrm{\Gamma }}^M=2\eta ^{MN}$$
or in terms of components
$$_{\left(\mathrm{\Gamma }^M\right)_{ab}\left(\stackrel{~}{\mathrm{\Gamma }}^N\right)^{bc}+\left(\mathrm{\Gamma }^N\right)_{ab}\left(\stackrel{~}{\mathrm{\Gamma }}^M\right)^{bc}=\left(1|\mathrm{\Gamma }^M\right)_{ab}\left(1|\stackrel{~}{\mathrm{\Gamma }}^N\right)^{bc}+\left(1|\mathrm{\Gamma }^N\right)_{ab}\left(1|\stackrel{~}{\mathrm{\Gamma }}^M\right)^{bc}=2\eta ^{MN}\delta _a^c.}$$
(24)
Our $`\mathrm{\Gamma }`$’s satisfy the very important identity
$$\mathrm{\Gamma }_{Ma(b}\mathrm{\Gamma }_{cd)}^M=1|\mathrm{\Gamma }_{Ma(b}1|\mathrm{\Gamma }_{cd)}^M=0.$$
(25)
## 3 The SSYM’s Auxliary Fields
Using Evans ansatz , SSYM are composed of: Gauge fields $`A_M`$, spinors $`\mathrm{\Psi }^a`$, $`j(=1..d3)`$ algebraic auxiliary fields $`K^j`$. The gauge group indices will be suppressed in the following. The Lagrangian density is
$$=\frac{1}{4}F_{MN}F^{MN}+\frac{i}{2}\mathrm{\Psi }^t\mathrm{\Gamma }^M_M\mathrm{\Psi }+\frac{1}{2}\delta _{ij}K^iK^j,$$
(26)
where $`_M_M+A_M;`$ $`F_{MN}[_M,_N]`$ and the $`\mathrm{\Gamma }`$ are given in (22). The Lagrangian is invariant up to a total derivative iff (25) holds. Our supersymmetry transformations are<sup>3</sup><sup>3</sup>3Contrary to , we set $`\mathrm{\Lambda }_j=\stackrel{~}{\mathrm{\Lambda }}^j`$ from the strart.
$`\delta _\eta A_M`$ $`=`$ $`i\eta \mathrm{\Gamma }_M\mathrm{\Psi },`$
$`\delta _\eta \mathrm{\Psi }^\alpha `$ $`=`$ $`{\displaystyle \frac{1}{2}}F_{MN}\left(\mathrm{\Gamma }_{MN}\eta \right)^\alpha +K^j\left(\mathrm{\Lambda }_j\right)_\beta ^\alpha \eta ^\beta ,`$
$`\delta _\eta K_j`$ $`=`$ $`i\left(\mathrm{\Gamma }^M_M\mathrm{\Psi }\right)_\alpha \left(\mathrm{\Lambda }_j\right)_\beta ^\alpha \eta ^\beta ,`$
where $`\mathrm{\Lambda }_P`$ are some real matrices and Lorentz transformations are generated by $`\mathrm{\Gamma }_{MN}\stackrel{~}{\mathrm{\Gamma }}_{[M}\mathrm{\Gamma }_{N]}`$. Imposing the closure of the supersymmetry infinitesimal transformations
$$[\delta _ϵ,\delta _\eta ]=2iϵ^t\mathrm{\Gamma }^M\eta _M.$$
(28)
The closure on $`A_M`$ yields
$$\mathrm{\Gamma }_M\mathrm{\Lambda }_j+\left(\mathrm{\Lambda }_j\right)^t\mathrm{\Gamma }_M=0.$$
(29)
In addition to this condition the closure on $`K^j`$ also requires
$$\mathrm{\Lambda }_j\mathrm{\Lambda }_h+\mathrm{\Lambda }_h\mathrm{\Lambda }_j=2\delta _{jh}.$$
(30)
While closure on the fermionic field $`\mathrm{\Psi }^\alpha `$ holds iff
$$\left(\mathrm{\Gamma }^M\right)_{\alpha \beta }\left(\stackrel{~}{\mathrm{\Gamma }}_M\right)^{\gamma \delta }=2\delta _{(\alpha }^\gamma \delta _{\beta )}^\delta +2\left(\mathrm{\Lambda }_j\right)_{(\alpha }^\gamma \left(\mathrm{\Lambda }_j\right)_{\beta )}^\delta .$$
Now, we continue in a different way to Evans. To construct $`\mathrm{\Lambda }_j`$, we first impose the additional condition
$$\left(\mathrm{\Lambda }\right)^t=\left(\mathrm{\Lambda }\right),$$
(31)
we notice from (30) that the $`\mathrm{\Lambda }_j`$ form a real Clifford algebra, and from (29)
$$\mathrm{\Gamma }_M\mathrm{\Lambda }_j\mathrm{\Lambda }_j\mathrm{\Gamma }_M=0.$$
(32)
that they commute with our space-time $`\mathrm{\Gamma }_M`$ Clifford algebra. The solution of the auxiliary field problem for $`d=3,4,6`$ dimensions, using (22) is then simply
$$\mathrm{\Lambda }_j=\left(\begin{array}{cc}1|𝔼_j& 0\\ 0& 1|𝔼_j\end{array}\right),$$
(33)
because
$$\left\{1\right|𝔼_j,1|𝔼_h\}=2\delta _{jh},$$
(34)
and
$$[𝔼_j,1|𝔼_h]=0.$$
(35)
Of course this solution is not unique. For example, if someone had started with $`1|\mathrm{\Gamma }_M`$, he would have found $`\mathrm{\Lambda }_j=\left(\begin{array}{cc}𝔼_j\hfill & 0\hfill \\ 0\hfill & 𝔼_j\hfill \end{array}\right)`$.
Now, we can relax the conditions (22) and (31). In general, one replaces left/right action used for the gamma matrices by right/left action for the $`\mathrm{\Lambda }_j`$ e.g.
$$_{\left(\mathrm{\Gamma }_j\right)_{ab}=\left(\begin{array}{cc}0& 𝔼_j|𝔼_{j+1}\\ 𝔼_j|𝔼_{j+1}& 0\end{array}\right)\left(\mathrm{\Lambda }_j\right)_{ab}=\left(\begin{array}{cc}𝔼_{j+1}|𝔼_j& 0\\ 0& 𝔼_{j+1}|𝔼_j\end{array}\right)}$$
(36)
One writes any $`\mathrm{\Gamma }`$ and expand it in terms left/right action $`(𝔼_{i,}1|𝔼_j,𝔼_m|𝔼_n)`$ then the $`\mathrm{\Lambda }`$ will be given in terms of suitable $`(1|𝔼_{i,}𝔼_j,𝔼_n|𝔼_m)`$ taking into account that daigonal elements should be replaced by non-diagonal one and interchanging left/right actions simultaneously.
## 4 The Ten Dimensions Case
For $`d=10`$, working with octonions the situation is different. We know that octonionic left and right action commutes only when applied to $`\phi `$,
$$\phi ^t[𝔼_j,1|𝔼_h]\phi =0,$$
(37)
and $`\phi `$ is just an 8 dimensional column matrix. Up to now, we have not restricted $`\phi `$ by any other conditions. With two different $`\phi `$, $`(\phi ^{\left(1\right)},\phi ^{\left(2\right)})`$, we impose now the conditions that $`\phi ^{\left(i\right)}`$ be fermionic fields. We express our 16 dimensional Grassmanian variables $`ϵ,\eta `$ of eqn.(28) in terms of $`\phi `$,
$$\begin{array}{ccc}& ϵ=\eta ^t\hfill & \\ & \hfill & \\ ϵ=\left(\begin{array}{cc}\phi ^{\left(1\right)}\hfill & \phi ^{\left(2\right)}\hfill \end{array}\right);\hfill & & \eta =\left(\begin{array}{c}\phi ^{\left(1\right)}\hfill \\ \phi ^{\left(2\right)}\hfill \end{array}\right)\hfill \end{array}$$
(38)
We now rederive (28) for the octonions. The closure conditions of our algebra, without omitting the Grassmanian variables are
$`\eta ^t\left(\mathrm{\Gamma }_M\mathrm{\Lambda }_j\mathrm{\Lambda }_j\mathrm{\Gamma }_M\right)\eta `$ $`=`$ $`0,`$
$`\eta ^t\left(\mathrm{\Lambda }_j\mathrm{\Lambda }_h+\mathrm{\Lambda }_h\mathrm{\Lambda }_j\right)\eta `$ $`=`$ $`\eta ^t\left(2\delta _{jh}\right)\eta ,`$
$`\eta ^t\left(\left(\mathrm{\Gamma }^M\right)_{\alpha \beta }\left(\stackrel{~}{\mathrm{\Gamma }}_M\right)^{\gamma \delta }\right)\eta `$ $`=`$ $`\eta ^t\left(2\delta _{(\alpha }^\gamma \delta _{\beta )}^\delta +2\left(\mathrm{\Lambda }_j\right)_{(\alpha }^\gamma \left(\mathrm{\Lambda }_j\right)_{\beta )}^\delta \right)\eta ,`$ (39)
which are satisfied for the octonionic representation
$$\left(\mathrm{\Gamma }_j\right)_{ab}=\left(\begin{array}{cc}0& 𝔼_j\\ 𝔼_j& 0\end{array}\right),\mathrm{\Lambda }_j=\left(\begin{array}{cc}1|𝔼_j& 0\\ 0& 1|𝔼_j\end{array}\right).$$
(40)
By interchanging left/right action, we have different solutions as in the quaternionic case. In summary, while the fermionic fields couple to left/right action through the gamma matrices, the auxiliary fields couple to right/left action through the $`\mathrm{\Lambda }`$. For the octonionic case the presence of the Grassmanian variables is essential. Contrary to the standard supersymmetry transformation, our Grassman variables are the same ($`ϵ=\eta ^t`$), which is identical to the result obtained by Berkovits in . According to Evans , the attractive feature of this scheme is that the Lagrangian (26) and the transformation (LABEL:dddd) are manifestly invariant under the generalized Lorentz group $`SO(1,9)`$. In our formulation, we can show some additional characteristic. In some cases, the (38) condition may be relaxed, for equal $`j`$ or $`h`$ (no summation)
$$\begin{array}{c}\phi ^t𝔼_j[𝔼_j,1|𝔼_h]\phi \\ \phi ^t\mathrm{\hspace{0.33em}\hspace{0.33em}1}|𝔼_i[𝔼_j,1|𝔼_h]\phi \\ \phi ^tE_h[𝔼_j,1|𝔼_h]\phi \\ \phi ^t\mathrm{\hspace{0.33em}\hspace{0.33em}1}|E_h[𝔼_j,1|𝔼_h]\phi \end{array}\}=0.$$
(41)
i.e. relating $`ϵ`$ and $`\eta `$ by an $`S^7`$ is also allowed.
Now, Let us show what will happen to $`spin(1,9)`$ when we transform it to $`softspin(1,9)`$
$`softspin(1,9)`$ $``$ $`[\mathrm{\Gamma }_i,\mathrm{\Gamma }_j]\eta `$ (48)
$`=`$ $`[\left(\begin{array}{cc}0& 𝔼_i\\ 𝔼_i& 0\end{array}\right),\left(\begin{array}{cc}0& 𝔼_j\\ 𝔼_j& 0\end{array}\right)]\left(\begin{array}{c}\phi ^{\left(1\right)}\hfill \\ \phi ^{\left(2\right)}\hfill \end{array}\right)`$
$`=`$ $`\left(\begin{array}{cc}0& [𝔼_i,𝔼_j]\\ [𝔼_i,𝔼_j]& 0\end{array}\right)\left(\begin{array}{c}\phi ^{\left(1\right)}\hfill \\ \phi ^{\left(2\right)}\hfill \end{array}\right)`$ (53)
$`=`$ $`\left(\begin{array}{cc}0\hfill & f_{ijk}^{(+)}\left(\phi ^{\left(2\right)}\right)𝔼_k\hfill \\ f_{ijk}^{(+)}\left(\phi ^{\left(1\right)}\right)𝔼_k\hfill & 0\hfill \end{array}\right)\left(\begin{array}{c}\phi ^{\left(1\right)}\hfill \\ \phi ^{\left(2\right)}\hfill \end{array}\right).`$ (58)
## 5 Some Superspace Hints
Lastly, let us make some comments about a possible superspace. It seems that the best way to find the $`d=6,10`$ superspace for SSYM is by defining some quaternionic and octonionic Grassmann variables that decompose the corresponding spinors into an $`SL(2,H)`$ and an $`SL(2,softS^7)`$ respectively
$$\{\theta _\alpha ,\theta _\beta \}=\{\overline{\theta }_{\dot{\alpha }},\overline{\theta }_{\dot{\beta }}\}=\{\theta _\alpha ,\overline{\theta }_{\dot{\beta }}\}=0,$$
(60)
where $`\alpha =1,2`$ over quaternions or octonions. We know that the supersymmetry generators $`Q_\alpha `$ are derived from right multiplication
$`Q_\alpha `$ $`=`$ $`\left(_\alpha 1|\mathrm{\Gamma }_{\alpha \dot{\beta }}^\mu \overline{\theta }^{\dot{\beta }}P_\mu \right)`$ (61)
$`Q^\alpha `$ $`=`$ $`\left(^\alpha +\overline{\theta }_{\dot{\beta }}1|\stackrel{~}{\mathrm{\Gamma }}^{\mu \dot{\beta }\alpha }P_\mu \right)`$ (62)
also
$$\overline{Q}^{\dot{\alpha }}=\left(^{\dot{\alpha }}1|\stackrel{~}{\mathrm{\Gamma }}^{\mu \dot{\alpha }\alpha }\theta _\alpha P_\mu \right)$$
(63)
$$\overline{Q}_{\dot{\alpha }}=\left(_{\dot{\alpha }}+\theta ^\alpha 1|\mathrm{\Gamma }_{\alpha \dot{\alpha }}P_\mu \right)$$
(64)
whereas the covariant derivative $`D_\alpha `$ are obtained by left action
$`D_\alpha `$ $`=`$ $`\left(_\alpha +\mathrm{\Gamma }_{\alpha \dot{\beta }}^\mu \overline{\theta }^{\dot{\beta }}P_\mu \right)`$ (65)
$`D^\alpha `$ $`=`$ $`\left(^\alpha \overline{\theta }_{\dot{\beta }}\stackrel{~}{\mathrm{\Gamma }}^{\mu \dot{\beta }\alpha }P_\mu \right)`$ (66)
also
$$\overline{D}^{\dot{\alpha }}=\left(^{\dot{\alpha }}+\stackrel{~}{\mathrm{\Gamma }}^{\mu \dot{\alpha }\alpha }\theta _\alpha P_\mu \right)$$
(67)
$$\overline{D}_{\dot{\alpha }}=\left(_{\dot{\alpha }}\theta ^\alpha \mathrm{\Gamma }_{\alpha \dot{\alpha }}P_\mu \right)$$
(68)
Leading to a result acceptable but different from the standard $`N=1`$, $`d=4`$ superspace,
$`\{Q_\alpha ,\overline{Q}_{\dot{\alpha }}\}`$ $`=`$ $`2\left(1|\mathrm{\Gamma }_{\alpha \dot{\alpha }}^\mu \right)P_\mu ,`$
$`\{Q_\alpha ,Q_\beta \}`$ $`=`$ $`\{\overline{Q}_{\dot{\alpha }},\overline{Q}_{\dot{\beta }}\}=\mathrm{\hspace{0.33em}0},`$
$`\{D_\alpha ,\overline{D}_{\dot{\alpha }}\}`$ $`=`$ $`2\mathrm{\Gamma }_{\alpha \dot{\alpha }}^\mu P_\mu ,`$
$`\{D_\alpha ,D_\beta \}`$ $`=`$ $`\{\overline{D}_{\dot{\alpha }},\overline{D}_{\dot{\beta }}\}=\mathrm{\hspace{0.33em}0},`$
and iff left and right action commute, we restore
$`\{Q_\alpha ,\overline{D}_{\dot{\alpha }}\}`$ $`=`$ $`\{D_\alpha ,\overline{Q}_{\dot{\alpha }}\}=0,`$
$`\{Q_\alpha ,D_\beta \}`$ $`=`$ $`\{\overline{D}_{\dot{\alpha }},\overline{Q}_{\dot{\beta }}\}=\mathrm{\hspace{0.33em}0}.`$
On the other hand for octonions we would have the weaker conditions,
$`\left(\begin{array}{cc}\phi ^{\left(1\right)}\hfill & \phi ^{\left(2\right)}\hfill \end{array}\right)\{Q_\alpha ,\overline{D}_{\dot{\alpha }}\}\left(\begin{array}{c}\phi ^{\left(1\right)}\hfill \\ \phi ^{\left(2\right)}\hfill \end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}\phi ^{\left(1\right)}\hfill & \phi ^{\left(2\right)}\hfill \end{array}\right)\{D_\alpha ,\overline{Q}_{\dot{\alpha }}\}\left(\begin{array}{c}\phi ^{\left(1\right)}\hfill \\ \phi ^{\left(2\right)}\hfill \end{array}\right)=0,`$ (75)
$`\left(\begin{array}{cc}\phi ^{\left(1\right)}\hfill & \phi ^{\left(2\right)}\hfill \end{array}\right)\{Q_\alpha ,D_\beta \}\left(\begin{array}{c}\phi ^{\left(1\right)}\hfill \\ \phi ^{\left(2\right)}\hfill \end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}\phi ^{\left(1\right)}\hfill & \phi ^{\left(2\right)}\hfill \end{array}\right)\{\overline{D}_{\dot{\alpha }},\overline{Q}_{\dot{\beta }}\}\left(\begin{array}{c}\phi ^{\left(1\right)}\hfill \\ \phi ^{\left(2\right)}\hfill \end{array}\right)=\mathrm{\hspace{0.33em}0}.`$ (82)
The commutation of left and right actions is not just needed for associativity but for the invariance under supersymmetry transformation
$$\delta _\xi \xi Q+\overline{\xi }\overline{Q}$$
(83)
because only the associativity ensures
$$\left(\begin{array}{cc}\phi ^{\left(1\right)}\hfill & \phi ^{\left(2\right)}\hfill \end{array}\right)[\delta _\xi ,D_\alpha ]\left(\begin{array}{c}\phi _1\hfill \\ \phi _2\hfill \end{array}\right)=\left(\begin{array}{cc}\phi ^{\left(1\right)}\hfill & \phi ^{\left(2\right)}\hfill \end{array}\right)[\delta _\xi ,\overline{D}_{\dot{\alpha }}]\left(\begin{array}{c}\phi _1\hfill \\ \phi _2\hfill \end{array}\right)=0,$$
(84)
since $`\delta _\xi `$ is left action and $`D_\alpha `$ is right action which is a very important relation in the standard $`N=1`$ superspace for the invariance of the Lagrangian under supersymmetry transformation. We hope to return to this point in a future work.
I am grateful to C. Imbimbo, P. Rotelli and A. Van Proeyen for useful comments.
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# Koornwinder polynomials and affine Hecke algebras
## 1. Introduction
Cherednik and Macdonald clarified the structure of Macdonald polynomials using certain representations of affine Hecke algebras in terms of difference-reflection operators. The underlying data stem from a fixed, reduced, irreducible root system $`\mathrm{\Sigma }`$. The degrees of freedom are, besides the deformation parameter $`q`$, the number of different root length occurring in $`\mathrm{\Sigma }`$ (so at most two). In fact, their work shows that the Macdonald polynomials are naturally attached to the reduced, affine root system $`\stackrel{~}{\mathrm{\Sigma }}`$ associated with $`\mathrm{\Sigma }`$.
As announced by Macdonald \[4, sect. 8\], and partially carried out by Noumi , Sahi , , Noumi & Stokman and Nishino et al , the theory naturally extends to the setting of arbitrary (not necessarily reduced) irreducible affine root systems. This in particular allows to incorporate the very general six parameter family of Koornwinder polynomials in the theory.
The affine root system underlying the Koornwinder polynomials is the non-reduced, irreducible affine root system $`S`$ of type $`C^{}C_n`$, which was introduced by Macdonald in . The affine root system $`S`$ contains all (possibly non-reduced) irreducible affine root systems of classical type as an affine root sub-system. On the polynomial level, this property is reflected by van Diejen’s observation that the families of Macdonald polynomials associated with classical root systems are special cases or limit cases of the Koornwinder polynomials. In fact, there is a large class of interesting families of multivariable orthogonal polynomials associated with classical root systems which are degenerate cases of the Koornwinder polynomials, see e.g. and . This supports the idea that the role of the Koornwinder polynomials in the theory of multivariable orthogonal polynomials associated with classical root systems is similar to the important and dominant role of the Askey-Wilson polynomials in the theory of one variable (basic) hypergeometric orthogonal polynomials.
In this paper we continue the affine Hecke algebra approach to the theory of Koornwinder polynomials. In particular, Sahi’s bi-orthogonality relations for the non-symmetric Koornwinder polynomials are extended to the case of continuous parameter values, and the corresponding diagonal terms are evaluated explicitly. We furthermore derive the evaluation formulas for the non-symmetric Koornwinder polynomials. Anti-symmetric Koornwinder polynomials are defined, and the explicit connection between the non-symmetric and the (anti-)symmetric theory is established. This leads to new derivations of Koornwinder’s , van Diejen’s and Sahi’s results on the orthogonality relations, quadratic norm evaluations and evaluation formulas for the symmetric Koornwinder polynomials. We also shortly discuss the analogue of Weyl’s character formula for Koornwinder polynomials, and the (closely related) shift operators.
Instead of using shift operators to evaluate the diagonal terms for the non-symmetric Koornwinder polynomials, we use a method which is motivated by Cherednik’s beautiful approach for proving Opdam’s inversion formula of the non-symmetric Harish-Chandra transform (the so-called Cherednik-Opdam transform). This method amounts to evaluating the diagonal terms using an explicit description of the action of the double affine Hecke algebra on non-symmetric Koornwinder polynomials in terms of operators acting on the spectral parameter. This approach reveals interesting new structures, such as an expression of the diagonal terms as residues of the bi-orthogonality weight function. This method is also expected to be an important tool for obtaining a better understanding of $`q`$-analogues of the Cherednik-Opdam transform, see e.g. and for some preliminary considerations in the rank one setting.
The results of this paper extend the results of Macdonald and Cherednik on the bi-orthogonality relations, diagonal term evaluations and evaluation formulas for non-symmetric Macdonald polynomials associated with classical reduced root systems, as well as the results of Noumi & Stokman , in which the rank one setting was treated in detail.
Acknowledgements: The author is supported by a fellowship from the Royal Netherlands Academy of Arts and Sciences (KNAW). The research was done during the author’s stay at Université Pierre et Marie Curie (Paris VI) and Institut de Recherche Mathématique Avancée (Strasbourg) in France, supported by a NWO-TALENT stipendium of the Netherlands Organization for Scientific Research (NWO) and by the EC TMR network “Algebraic Lie Representations”, grant no. ERB FMRX-CT97-0100. The author thanks Masatoshi Noumi for fruitful discussions, and Siddhartha Sahi and Akinori Nishino for drawing his attention to the papers and , respectively.
## 2. The affine root system of type $`C^{}C_n`$
In this section we discuss the affine root system of type $`C^{}C_n`$, which was introduced by Macdonald in . Let $`n`$ be a positive integer $`2`$. Let $`V=(^n,.,.)`$ be Euclidean $`n`$-space with orthonormal basis $`\{ϵ_i\}_{i=1}^n`$. We write $`\widehat{V}`$ for the affine linear transformations from $`V`$ to $``$. As a vector space, $`\widehat{V}`$ can be identified with $`V\delta `$, where vectors in $`V`$ are considered as linear functionals on $`V`$ via the scalar product $`.,.`$, and where $`\delta `$ is the function identically equal to one on $`V`$. We extend the scalar product $`.,.`$ to a positive semi-definite form on $`\widehat{V}`$ by requiring that the constant function $`\delta `$ is in the radical of $`.,.`$.
Let $`S\widehat{V}`$ be the subset
$$\begin{array}{cc}\hfill S=& \{\pm ϵ_i+\frac{m}{2}\delta ,\pm 2ϵ_i+m\delta |m,i=1,\mathrm{},n\}\hfill \\ & \{\pm ϵ_i\pm ϵ_j+m\delta |m,\mathrm{\hspace{0.17em}1}i<jn\},\hfill \end{array}$$
(2.1)
where all the sign combinations occur. Let $`𝒲=𝒲(S)`$ be the sub-group of $`\text{GL}_{}(\widehat{V})`$ generated by the reflections $`s_\beta `$ ($`\beta S`$), where
$$s_f(g)=gg,f^{}f,f\widehat{V}\delta ,g\widehat{V},$$
and where $`f^{}=2f/f,f`$ is the co-root of $`f`$. Observe that $`s_f(g)=g\stackrel{~}{s}_f^1`$, with $`\stackrel{~}{s}_f:VV`$ the orthogonal reflection in the affine hyperplane $`f^1(0)`$.
By , $`S\widehat{V}`$ is an irreducible, affine root system. In particular, $`\alpha ,\beta ^{}`$ for all $`\alpha ,\beta S`$, and $`S`$ is stable under the action of $`𝒲`$. The sub-group $`𝒲\text{GL}(\widehat{V})`$ is called the affine Weyl group of $`S`$.
Let $`R`$ be the inmultiplyable roots in $`S`$ and $`R^{}S`$ the corresponding co-root system. Then $`R`$ and $`R^{}`$ are irreducible, reduced affine root systems in $`\widehat{V}`$, with affine Weyl group $`𝒲`$. The projection $`\mathrm{\Sigma }R`$ of $`R`$ on $`V`$ along the direct sum decomposition $`\widehat{V}=V\delta `$ is an irreducible root system of type $`C_n`$ with Weyl group $`W=S_n(\pm 1)^n𝒲`$ given by permutations and sign changes of the fixed basis $`\{ϵ_i\}_{i=1}^n`$ of $`V`$ (here $`S_n`$ denotes the symmetric group in $`n`$ letters). Due to the (non-disjoint) union $`S=RR^{}`$ of $`S`$ into the reduced affine root sub-system $`R`$ of type $`\stackrel{~}{C}_n`$ and its co-root system, we call $`S`$ of type $`C^{}C_n`$, cf. .
Let $`Q^{}`$ be the co-root lattice of $`\mathrm{\Sigma }`$, which coincides with the weight lattice $`\mathrm{\Lambda }`$ of $`\mathrm{\Sigma }`$. In fact, $`Q^{}=\mathrm{\Lambda }`$ is the full $``$-lattice in $`V`$ with basis $`\{ϵ_i\}_{i=1}^n`$. Then
$$𝒲=W\tau (Q^{}),$$
where $`\tau (v)\text{GL}(\widehat{V})`$ ($`vV`$) is the translation operator defined by $`\tau (v)f=f+v,f\delta `$ for $`f\widehat{V}`$. Observe that $`\tau (v)f=f\stackrel{~}{\tau }_v^1`$ with $`\stackrel{~}{\tau }_v:VV`$ given by $`\stackrel{~}{\tau }_v(\lambda )=\lambda v`$.
###### Remark 2.1.
In Macdonald’s and Cherednik’s work the translation operator $`\tau (v)`$ in fact corresponds to $`\tau (v)=\tau (v)^1`$. Later on, this change of convention (which is related to conjugation with the largest Weyl group element $`\sigma W`$) causes certain changes of signs compared with Cherednik’s and Macdonald’s theory, see e.g. remark 4.10(ii).
We fix a basis $`\{a_i\}_{i=0}^n`$ of $`R`$ by
$$a_0=\delta 2ϵ_1,a_i=ϵ_iϵ_{i+1}(i=1,\mathrm{},n1),a_n=2ϵ_n.$$
Observe that $`\{a_0^{}=a_0/2,a_1,\mathrm{},a_{n1},a_n^{}=a_n/2\}`$ is a basis of $`R^{}`$, as well as of $`S`$. Furthermore, $`\{a_i\}_{i=1}^n`$ is a basis of the gradient root system $`\mathrm{\Sigma }`$. We write $`\mathrm{\Sigma }^+`$ (respectively $`\mathrm{\Sigma }^{}`$) for the corresponding positive (respectively negative) roots in $`\mathrm{\Sigma }`$, and $`\mathrm{\Lambda }^+=_{i=1}^n_+\omega _i`$ for the corresponding cone of dominant weights of $`\mathrm{\Sigma }`$. Here $`\omega _i=ϵ_1+\mathrm{}+ϵ_i`$ $`(i=1,\mathrm{},n)`$ are the fundamental weights of $`\mathrm{\Lambda }`$, i.e. $`\omega _i,a_j^{}=\delta _{i,j}`$ for all $`i,j=1,\mathrm{},n`$, where $`\delta _{i,j}`$ is the Kronecker delta. We furthermore write $`Q^{,+}`$ for the positive span of the simple co-roots $`a_i^{}`$ ($`i=1,\mathrm{},n`$).
Let $`R^+`$ (respectively $`R^{}`$) be the positive (respectively negative) roots of $`R`$ with respect to the basis of the previous paragraph. In particular, $`R^+=\mathrm{\Sigma }^+\{\beta R|\beta (0)>0\}`$.
The affine Weyl group $`𝒲`$ is generated by the simple reflections $`s_i=s_{a_i}`$ ($`i=0,\mathrm{},n`$). In fact, $`𝒲`$ is isomorphic to the Coxeter group with generators $`s_i`$ ($`i=0,\mathrm{},n`$) satisfying $`s_i^2=1`$ and the braid relations $`s_is_{i+1}s_is_{i+1}=s_{i+1}s_is_{i+1}s_i`$ ($`i=0`$, $`i=n1`$), $`s_is_{i+1}s_i=s_{i+1}s_is_{i+1}`$ ($`i=1,\mathrm{},n2`$) and $`s_is_j=s_js_i`$ for $`|ij|2`$.
With our present conventions, Lusztig’s formula \[15, 1.4(a)\] for the length of an element in $`𝒲`$ is given by
$$l(\tau (\lambda )w)=\underset{\alpha \mathrm{\Sigma }^+}{}|\lambda ,w\alpha +\chi (w\alpha )|,\lambda \mathrm{\Lambda },wW,$$
(2.2)
where $`\chi (\alpha )=1`$ if $`\alpha \mathrm{\Sigma }^{}`$ and $`=0`$ otherwise.
We write $`\mathrm{\Sigma }=\mathrm{\Sigma }_m\mathrm{\Sigma }_l`$ for the decomposition of $`\mathrm{\Sigma }`$ into $`W`$-orbits, where $`\mathrm{\Sigma }_m`$ (respectively $`\mathrm{\Sigma }_l`$) is the set of roots of length two (respectively four). We furthermore set $`\mathrm{\Sigma }_s=\frac{1}{2}\mathrm{\Sigma }_l`$. There are five $`𝒲`$-orbits in $`S`$, namely
$$\begin{array}{cc}\hfill 𝒲a_0^{}& =\left(\frac{1}{2}+\right)\delta +\mathrm{\Sigma }_s,𝒲a_0=\left(1+2\right)\delta +\mathrm{\Sigma }_l,\hfill \\ \hfill 𝒲a_i& =\delta +\mathrm{\Sigma }_m(i\{1,\mathrm{},n1\}\text{ arbitrary}),\hfill \\ \hfill 𝒲a_n^{}& =\delta +\mathrm{\Sigma }_s,𝒲a_n=2\delta +\mathrm{\Sigma }_l.\hfill \end{array}$$
(2.3)
Observe that $`R`$ (respectively $`R^{}`$) has three $`𝒲`$-orbits, namely $`𝒲a_0`$, $`𝒲a_i`$ and $`𝒲a_n`$ (respectively $`𝒲a_0^{}`$, $`𝒲a_i`$ and $`𝒲a_n^{}`$), where $`i\{1,\mathrm{},n1\}`$ is arbitrary.
For later purposes, we define an action of $`𝒲`$ on $`V`$ which extends the canonical $`W`$-action on $`V`$. It suffices to specify the action of the simple reflection $`s_0`$ on $`V`$, which we take to be
$$s_0.x=(1x_1,x_2,\mathrm{},x_n),$$
where $`x_i=x,ϵ_i`$. Observe that $`\mathrm{\Lambda }V`$ is $`𝒲`$-stable, and that $`\tau (\lambda ).x=x+\lambda `$ for $`\lambda \mathrm{\Lambda }`$. We denote this action of $`𝒲`$ on $`V`$ with a dot and we call it the dot-action, in order to avoid confusion with the canonical action of $`𝒲`$ on $`\widehat{V}`$ and its induced dual action on $`V`$.
## 3. The (double) affine Hecke algebra
Let $`𝒜`$ be the group algebra of the weight lattice $`\mathrm{\Lambda }`$. We write $`x^\lambda `$ ($`\lambda \mathrm{\Lambda }`$) for the canonical basis of $`𝒜`$, so that $`x^0=1`$ is the unit element in $`𝒜`$ and $`x^\lambda x^\mu =x^{\lambda +\mu }`$ for all $`\lambda ,\mu \mathrm{\Lambda }`$. The group algebra $`𝒜`$ is isomorphic to the Laurent polynomials in the $`n`$ independent indeterminates $`x_i=x^{ϵ_i}`$ ($`i=1,\mathrm{},n`$). Let $`q\{0\}`$ be generic complex (in particular, not a root of unity) and let $`q^{1/2}`$ be a fixed square root of $`q`$. We write $`x^{\mu +c\delta }=q^cx^\mu `$ for $`\mu \mathrm{\Lambda }`$ and $`c\frac{1}{2}`$. Then the assignment $`w(x^\mu )=x^{w\mu }`$ for $`w𝒲`$ and $`\mu \mathrm{\Lambda }\widehat{V}`$ extends by linearity to an action of $`𝒲`$ on $`𝒜`$. In particular, the action of the simple reflections $`s_i`$ ($`i=0,\mathrm{},n`$) on $`𝒜`$ is given by
$$\begin{array}{cc}\hfill (s_0f)(x)& =f(qx_1^1,x_2,\mathrm{},x_n),\hfill \\ \hfill (s_if)(x)& =f(x_1,\mathrm{},x_{i1},x_{i+1},x_i,x_{i+2},\mathrm{},x_n)(i=1,\mathrm{},n1),\hfill \\ \hfill (s_nf)(x)& =f(x_1,\mathrm{},x_{n1},x_n^1),\hfill \end{array}$$
(3.1)
where $`f𝒜`$ and $`x=(x_1,\mathrm{},x_n)`$. In particular, the translation operators $`\tau (\mu )`$ ($`\mu \mathrm{\Lambda }`$) act as $`q`$-difference operators: $`\tau (\mu )(x^\lambda )=q^{\mu ,\lambda }x^\lambda `$ for all $`\lambda ,\mu \mathrm{\Lambda }`$.
The Noumi representation is a five($`=\mathrm{\#}\{𝒲\text{orbits in }S\}`$) parameter deformation of the above action of $`𝒲`$ on the group algebra $`𝒜`$. We incorporate the five extra degrees of freedom in a so-called multiplicity function $`𝐭=(t_\beta )_{\beta S}`$ of $`S`$, which is a $`𝒲`$-invariant map from $`S`$ to $`\{0\}`$ (so $`t_{w\beta }=t_\beta `$ for all $`w𝒲`$ and all $`\beta S`$). We furthermore set $`t_f=1`$ if $`f\widehat{V}S`$. A multiplicity function is thus uniquely determined by the five values $`t_{a_0^{}}`$, $`t_{a_0}`$, $`t_{a_i}`$ ($`i\{1,\mathrm{},n1\}`$ arbitrary), $`t_{a_n}`$ and $`t_{a_n^{}}`$. In order to avoid cumbersome notations, we sometimes write $`t_i`$ (respectively $`t_i^{}`$) for $`t_{a_i}`$ (respectively $`t_{a_i^{}}`$) and we write $`t`$ for the value of $`t_j=t_j^{}`$ with $`j\{1,\mathrm{},n1\}`$. Furthermore, we use the short-hand notation $`𝐤=(t_\beta )_{\beta R}(t_0,t,t_n)`$ and $`𝐤^{}=(t_\beta )_{\beta R^{}}(t_0^{},t,t_n^{})`$ for the corresponding multiplicity functions of $`R`$ and $`R^{}`$, respectively. We assume throughout the paper that the values of the multiplicity function $`𝐭`$ are generically complex.
In the Noumi representation, the role of the affine Weyl group $`𝒲`$ is replaced by the affine Hecke algebra of type $`\stackrel{~}{C}_n`$, which is defined as follows.
###### Definition 3.1.
The affine Hecke algebra $`H=H(R;𝐤)`$ of type $`\stackrel{~}{C}_n`$ is the unital, associative algebra with generators $`T_0,\mathrm{},T_n`$ and relations
$$(T_it_i)(T_i+t_i^1)=0,(i=0,\mathrm{},n),$$
(3.2)
and the braid relations
$$\begin{array}{cc}\hfill T_iT_{i+1}T_iT_{i+1}& =T_{i+1}T_iT_{i+1}T_i,(i=0,i=n1),\hfill \\ \hfill T_iT_{i+1}T_i& =T_{i+1}T_iT_{i+1},(i=1,\mathrm{},n2),\hfill \\ \hfill T_iT_j& =T_jT_i,|ij|2.\hfill \end{array}$$
(3.3)
We call (3.2) and (3.3) the $`H(R;𝐤)`$-relations for the $`(n+1)`$-tuple $`(T_0,\mathrm{},T_n)`$. Furthermore, we write $`H(R^{};𝐤^{})`$ for the affine Hecke algebra $`H`$ in which the parameter $`t_i`$ is replaced by $`t_i^{}`$ for $`i=0,\mathrm{},n`$.
We recall here some of the basic properties of the affine Hecke algebra $`H`$, see Lusztig for details and for a general discussion on affine Hecke algebras.
For a reduced expression $`w=s_{i_1}\mathrm{}s_{i_r}`$ of $`w𝒲`$ we set $`T_w=T_{i_1}\mathrm{}T_{i_r}`$. This is independent of the choice of reduced expression by the braid relations (3.3) for the $`T_i`$, and $`\{T_w\}_{w𝒲}`$ is a linear basis of $`H`$.
For $`\lambda \mathrm{\Lambda }^+`$ we set $`Y^\lambda =T_{\tau (\lambda )}`$, and for $`\lambda =\mu \nu \mathrm{\Lambda }`$ with $`\mu ,\nu \mathrm{\Lambda }^+`$ we set $`Y^\lambda =Y^\mu (Y^\mu )^1`$. The length identity (2.2) implies that the $`Y^\lambda `$ ($`\lambda \mathrm{\Lambda }`$) are well-defined (i.e. independent of the choice of decomposition $`\lambda =\mu \nu `$). Furthermore, the sub-space $`𝒜_Y=\text{span}\{Y^\lambda |\lambda \mathrm{\Lambda }\}`$ is a commutative subalgebra of $`H`$ isomorphic to $`𝒜`$ (in particular, $`Y^0=1`$ and $`Y^\lambda Y^\mu =Y^{\lambda +\mu }`$ for all $`\lambda ,\mu \mathrm{\Lambda }`$). We identify $`f(x)=_\lambda c_\lambda x^\lambda 𝒜`$ with $`f(Y)=_\lambda c_\lambda Y^\lambda 𝒜_Y`$ in the remainder of the paper. We write $`Y_i=Y^{ϵ_i}`$, which corresponds with $`x_i`$ under the identification of $`𝒜_Y`$ with $`𝒜`$.
The Hecke algebra $`H_0=H_0(\mathrm{\Sigma };t,t_n)`$ of the finite Weyl group $`W`$ can be identified with the subalgebra of $`H`$ generated by $`T_i`$ ($`i=1,\mathrm{},n`$). Then $`\{T_w\}_{wW}`$ is a linear basis of $`H_0`$ and
$$HH_0𝒜_Y𝒜_YH_0$$
as vector spaces by multiplication. The commutation relations between $`T_iH_0`$ ($`i=1,\mathrm{},n`$) and $`f(Y)𝒜_Y`$ are given by the formulas
$$\begin{array}{cc}\hfill T_if(Y)(s_if)(Y)T_i& =(tt^1)\left(\frac{f(Y)(s_if)(Y)}{1Y^{a_i}}\right),\hfill \\ \hfill T_nf(Y)(s_nf)(Y)T_n& =\left((t_nt_n^1)+(t_0t_0^1)Y_n^1\right)\left(\frac{f(Y)(s_nf)(Y)}{1Y_n^2}\right)\hfill \end{array}$$
(3.4)
for $`i=1,\mathrm{},n1`$, see \[15, prop. 3.6\]. These commutation relations can be used to prove inductively that
$$Y_i=T_i\mathrm{}T_{n1}T_nT_{n1}\mathrm{}T_1T_0T_1^1T_2^1\mathrm{}T_{i1}^1,i\{1,\mathrm{},n\}$$
(3.5)
in $`H`$, see Noumi or Sahi \[24, (11)\]. The expression (3.5) is the analogue in $`H`$ of the reduced expression
$$\tau (ϵ_i)=s_i\mathrm{}s_{n1}s_ns_{n1}\mathrm{}s_1s_0s_1\mathrm{}s_{i1},i\{1,\mathrm{},n\}$$
(3.6)
in $`𝒲`$. We define a rational function $`v_\beta (x)=v_\beta (x;𝐭;q)(x)=\text{Quot}(𝒜)`$ by
$$v_\beta (x;𝐭;q)=\frac{\left(1t_\beta t_{\beta /2}x^{\beta /2}\right)\left(1+t_\beta t_{\beta /2}^1x^{\beta /2}\right)}{\left(1x^\beta \right)},\beta R.$$
(3.7)
Observe that for $`\beta R`$ with $`\beta /2S`$, the expression (3.7) reduces to $`v_\beta (x)=(1t_\beta ^2x^\beta )/(1x^\beta )`$ since $`t_{\beta /2}=1`$. The following crucial theorem was proved by Noumi .
###### Theorem 3.2 (The Noumi representation).
The assignment
$$T_it_i+t_i^1v_{a_i}(x;𝐭;q)\left(s_i\text{id}\right)\text{End}_{}(𝒜)$$
for $`i=0,\mathrm{},n`$ uniquely extends to a representation $`\pi _{𝐭,q}:H(R;𝐤)\text{End}_{}(𝒜)`$.
The commutation relations (3.4) play a crucial role in the proof of theorem 3.2, compare with the argument in and \[18, (4.6)\] in case of reduced root systems.
We write $`T_i`$ for the image of $`T_iH(R;𝐤)`$ under the Noumi representation $`\pi _{𝐭,q}`$ for $`i=0,\mathrm{},n`$ if no confusion is possible, and we call them the difference-reflection operators associated with $`S`$.
We are now in a position to recall Sahi’s definition of the double affine Hecke algebra.
###### Definition 3.3.
The double affine Hecke algebra $`=(S;𝐭;q)`$ is the sub-algebra of $`\text{End}_{}(𝒜)`$ generated by $`\pi _{𝐭,q}(H(R;𝐤))`$ and $`𝒜`$, where the elements in $`𝒜`$ are considered as multiplication operators in $`\text{End}_{}(𝒜)`$.
We end this section by giving an alternative presentation of $``$ (also different from Sahi’s \[24, sect. 3\] presentation), which emphasizes its close connection with the affine root system $`S`$.
We write $`f(z)=_\lambda c_\lambda z^\lambda \text{End}_{}(𝒜)`$ for the multiplication operator associated with the Laurent polynomial $`f(x)=_\lambda c_\lambda x^\lambda 𝒜`$. In particular, $`z^{\lambda +m\delta }=q^mz^\lambda `$ ($`\lambda \mathrm{\Lambda }`$, $`m\frac{1}{2}`$) is the multiplication operator associated with $`x^{\lambda +m\delta }=q^mx^\lambda 𝒜`$. Then in $``$ we have the commutation relations
$$f(z)T_iT_i(s_if)(z)=\frac{(t_{a_i}t_{a_i}^1)+(t_{a_i/2}t_{a_i/2}^1)z^{a_i/2}}{1z^{a_i}}\left(f(z)(s_if)(z)\right)$$
(3.8)
for $`i=0,\mathrm{},n`$ and $`f𝒜`$. This follows from the fact that the difference-reflection operator $`T_i`$ can be rewritten as
$$T_i=t_{a_i}s_i+\frac{(t_{a_i}t_{a_i}^1)+(t_{a_i/2}t_{a_i/2}^1)x^{a_i/2}}{1x^{a_i}}(\text{id}s_i)$$
for $`i=0,\mathrm{},n`$.
###### Theorem 3.4.
The double affine Hecke algebra $`(S;𝐭;q)`$ is isomorphic to the unital, associative algebra $`(𝐭;q)`$ with generators $`V_0^{},V_0`$, $`V_i`$ ($`i=1,\mathrm{},n`$) and $`V_n^{}`$, satisfying
1. The $`H(R;𝐤)`$-relations for $`(V_0,V_1,\mathrm{},V_{n1},V_n)`$.
2. The $`H(R^{};𝐤^{})`$-relations for $`(V_0^{},V_1,\mathrm{},V_{n1},V_n^{})`$.
3. (Compatibility conditions). $`V_n^{}V_nV_{n1}\mathrm{}V_1V_0V_0^{}V_1V_2\mathrm{}V_{n1}=q^{1/2}`$ and $`[V_0,V_n^{}]=0=[V_0^{},V_n]`$.
The algebra isomorphism $`\varphi :(𝐭;q)(S;𝐭;q)`$ is explicitly given by $`\varphi (V_i)=T_i`$ ($`i=0,\mathrm{},n`$), $`\varphi (V_0^{})=T_0^1z^{a_0^{}}=q^{1/2}T_0^1z_1`$ and $`\varphi (V_n^{})=z^{a_n^{}}T_n^1=z_n^1T_n^1`$.
###### Proof.
For the existence of $`\varphi `$, we need to check that the elements $`T_0^{}=T_0^1z^{a_0^{}}`$, $`T_j`$ ($`j=0,\mathrm{},n`$) and $`T_n^{}=z^{a_n^{}}T_n^1`$ in $``$ respect the defining relations of the generators $`V_0^{}`$, $`V_j`$ ($`j=0,\mathrm{},n`$), $`V_n^{}`$ in $``$. The $`H(R;𝐤)`$-relations for $`(T_0,\mathrm{},T_n)`$ is precisely the content of theorem 3.2. The $`H(R^{};𝐤^{})`$-relations for the $`(n+1)`$-tuple $`(T_0^{},T_1,\mathrm{},T_{n1},T_n^{})`$ follows easily from (3.8) (see also \[24, sect. 3\]). By (3.8) it follows inductively that
$$\begin{array}{cc}\hfill z_i& =T_i^1T_{i+1}^1\mathrm{}T_n^1(T_n^{})^1T_{n1}^1\mathrm{}T_i^1\hfill \\ & =q^{1/2}T_{i1}\mathrm{}T_1T_0T_0^{}T_1\mathrm{}T_{i1}\hfill \end{array}$$
for $`i=1,\mathrm{},n`$ in $``$. In particular, the identity $`z_n^1z_n=1`$ in $``$ shows that the generators $`T_0^{}`$, $`T_j`$ ($`j=0,\mathrm{},n`$) and $`T_n^{}`$ satisfy the compatibility condition in $``$. Hence the algebra homomorphism $`\varphi `$ exists.
We write $`w_i=q^{1/2}V_{i1}\mathrm{}V_1V_0V_0^{}V_1\mathrm{}V_{i1}`$ for $`i\{1,\mathrm{},n\}`$, so that $`\varphi (w_i)=z_i`$ for $`i=1,\mathrm{},n`$. Since the $`z_i`$ ($`i=1,\mathrm{},n`$) and the $`T_j=\varphi (V_j)`$ ($`j=0,\mathrm{},n`$) generate $``$ as an algebra, we see that $`\varphi `$ is surjective. On the other hand, all fundamental relations of $`(S;𝐭;q)`$ as given by Sahi \[26, sect. 3\] can be easily checked for the generators $`w_i`$ ($`i=1,\mathrm{},n`$) and $`V_j`$ ($`j=0,\mathrm{},n`$) of $``$. This implies the injectivity of $`\varphi `$. ∎
###### Remark 3.5.
The presentation of $`(S;𝐭;q)`$ as given in theorem 3.4 clearly reflects the structure of the underlying non-reduced affine root system $`S`$. In particular, the first part of the compatibility condition can be recovered from the root data as follows. We put the simple roots for the indivisible roots $`R^{}`$ (respectively for the inmultiplyable roots $`R`$) above (respectively below) the corresponding vertices of the extended Dynkin diagram:
$$\begin{array}{ccccc}\underset{a_0}{\stackrel{a_0^{}}{}}& =& \underset{a_1}{\stackrel{a_1}{}}\frac{}{}\underset{a_2}{\stackrel{a_2}{}}\frac{}{}\mathrm{}\mathrm{}\frac{}{}\underset{a_{n1}}{\stackrel{a_{n1}}{}}& =& \underset{a_n}{\stackrel{a_n^{}}{}}\end{array}$$
Now we attach $`V_i(=T_i)`$ to the simple roots $`a_i`$, $`V_0^{}(=T_0^{}=T_0^1z^{a_0^{}})`$ to the co-root $`a_0^{}`$ and $`V_n^{}(=T_n^{}=z^{a_n^{}}T_n^1)`$ to the co-root $`a_n^{}`$ in the above diagram. We call them the simple generators of $``$. Then the compatibility condition amounts to the following rule: multiplying simple generators in the order of appearance of a single walk around the diagram in clockwise direction, gives $`q^{1/2}`$. The point of departure for the walk is irrelevant, since the compatibility condition is equivalent to the compatibility condition in which the factors in its left-hand side are permuted cyclically.
## 4. Non-symmetric Koornwinder polynomials and triangularity
In this section we show that the $`Y`$-operators $`Y^\lambda `$ ($`\lambda \mathrm{\Lambda }`$) act as triangular operators under the Noumi representation $`\pi _{𝐭,q}`$. We use this triangularity property to redefine Sahi’s non-symmetric Koornwinder polynomials. The advantage of this method is that triangularity properties of the non-symmetric Koornwinder polynomials are automatically incorporated in their definition, in contrast with Sahi’s , approach.
Let $`\lambda ^+\mathrm{\Lambda }^+`$ for $`\lambda \mathrm{\Lambda }`$ be the unique dominant weight in the orbit $`W\lambda `$. We will be needing the following two partial orders on the weight lattice $`\mathrm{\Lambda }`$.
###### Definition 4.1.
Let $`\lambda ,\mu \mathrm{\Lambda }`$.
(i) We write $`\lambda \mu `$ if $`\mu \lambda Q^{,+}`$ (and $`\lambda <\mu `$ if $`\lambda \mu `$ and $`\lambda \mu `$).
(ii) We write $`\lambda \mu `$ if $`\lambda ^+<\mu ^+`$, or if $`\lambda ^+=\mu ^+`$ and $`\lambda \mu `$ (and $`\lambda \mu `$ if $`\lambda \mu `$ and $`\lambda \mu `$).
###### Lemma 4.2.
Let $`\mu \mathrm{\Lambda }`$ and $`\alpha \mathrm{\Sigma }^+`$.
If $`\mu ,\alpha 2`$, then $`\mu r\alpha ^{}\mu `$ for $`r=1,\mathrm{},\mu ,\alpha 1`$.
If $`\mu ,\alpha 2`$, then $`\mu +r\alpha ^{}\mu `$ for $`r=1,\mathrm{},\mu ,\alpha 1`$.
###### Proof.
We write $`m_\alpha =\mu ,\alpha `$. Suppose that $`m_\alpha 2`$ and write $`\mu _r=\mu r\alpha ^{}`$ with $`r\{1,\mathrm{},m_\alpha 1\}`$. We show that $`\mu _r^+<\mu ^+`$.
Let $`wW`$ such that $`\mu _r^+=w\mu _r`$. If $`w\alpha ^{}Q^{,+}`$, then $`\mu _r^+=w\mu rw\alpha ^{}<w\mu \mu ^+`$. On the other hand, if $`w\alpha ^{}Q^{,+}`$, then $`\mu _r^+=w\mu rw\alpha ^{}<w\mu m_\alpha w\alpha ^{}=(ws_\alpha )\mu \mu ^+`$. This proves the assertion for $`m_\alpha 2`$. The case $`m_\alpha 2`$ can be obtained by applying the previous case to $`s_\alpha \mu `$. ∎
For $`\beta R`$ we define
$$(\beta )=t_\beta s_\beta +t_\beta ^1v_\beta (x)\left(1s_\beta \right)\text{End}_{}(𝒜),$$
(4.1)
where $`v_\beta ()`$ is given by (3.7). Let $`ϵ:\{\pm 1\}`$ be the function which maps a positive integer to $`1`$ and a strictly negative integer to $`1`$.
###### Lemma 4.3.
Let $`\lambda \mathrm{\Lambda }`$. For $`\beta =\alpha +m\delta R^+`$ with $`\alpha \mathrm{\Sigma }^+`$ we have
$$(\beta )(x^\lambda )=t_\beta ^{ϵ(\lambda ,\beta )}x^\lambda +\underset{\mu \lambda }{}c_{\lambda ,\mu }x^\mu $$
for certain constants $`c_{\lambda ,\mu }`$.
###### Proof.
Let $`D_\beta \text{End}_{}(𝒜)`$ for $`\beta R`$ be the divided difference-reflection operator defined by
$$D_\beta f=\frac{fs_\beta f}{1x^\beta },f𝒜.$$
Then for all $`\lambda \mathrm{\Lambda }`$ we have
$$D_\beta (x^\lambda )=\{\begin{array}{cc}x^{\lambda \beta }x^{\lambda 2\beta }\mathrm{}x^{\lambda \lambda ,\beta ^{}\beta }\hfill & \text{if }\lambda ,\beta ^{}>0,\hfill \\ 0\hfill & \text{if }\lambda ,\beta ^{}=0,\hfill \\ x^\lambda +x^{\lambda +\beta }+\mathrm{}+x^{\lambda (1+\lambda ,\beta ^{})\beta }\hfill & \text{if }\lambda ,\beta ^{}<0.\hfill \end{array}$$
The proof follows now easily from lemma 4.2 and from the definition (4.1) of $`(\beta )`$. ∎
Observe that $`(a_i)=T_is_i`$ for $`i=0,\mathrm{},n`$ and that $`(w(\beta ))=w(\beta )w^1`$ for all $`w𝒲`$ and all $`\beta R`$. Combined with (3.5) and (3.6), we obtain
$$\begin{array}{cc}\hfill Y_i=& (ϵ_iϵ_{i+1})(ϵ_iϵ_{i+2})\mathrm{}(ϵ_iϵ_n)(2ϵ_i)\hfill \\ & (ϵ_i+ϵ_n)\mathrm{}(ϵ_i+ϵ_{i+1})(ϵ_i+ϵ_{i1})\mathrm{}(ϵ_i+ϵ_1)\hfill \\ & (\delta +2ϵ_i)\tau (ϵ_i)(ϵ_1ϵ_i)^1\mathrm{}(ϵ_{i1}ϵ_i)^1\hfill \end{array}$$
(4.2)
for $`i=1,\mathrm{},n`$, cf. . Hence the triangularity of the factors $`()`$ in (4.2) (see lemma 4.3) implies the triangularity of $`Y_i`$ for $`i=1,\mathrm{},n`$, and hence of $`Y^\lambda `$ for all $`\lambda \mathrm{\Lambda }`$. For the explicit description of the diagonal terms of the $`Y`$-operators, we need to introduce some additional notations first.
We write $`f(y)`$ for the value of $`f𝒜`$ at $`y=(y_1,\mathrm{},y_n)(\{0\})^n`$. In particular, $`(y)^{\lambda +m\delta }=q^m(y)^\lambda `$ ($`m\frac{1}{2}`$, $`\lambda \mathrm{\Lambda }`$) is the value of $`x^{\lambda +m\delta }=q^mx^\lambda 𝒜`$ at $`y`$. Conversely, we let $`c^\lambda (\{0\})^n`$ for $`c\{0\}`$ and $`\lambda \mathrm{\Lambda }`$ be the vector $`c^\lambda =(c^{\lambda _1},\mathrm{},c^{\lambda _n})`$, where $`\lambda _i=\lambda ,ϵ_i`$.
###### Remark 4.4.
The brackets in the notation $`(y)^{\lambda +m\delta }`$ for the value of $`x^{\lambda +m\delta }`$ at $`y(\{0\})^n`$ will occasionally be omitted when $`m=0`$. For $`m0`$ the brackets are needed to distinguish the value of $`x^{\lambda +m\delta }`$ at $`y^1=(y_1^1,\mathrm{},y_n^1)`$ from the value of $`x^{\lambda m\delta }`$ at $`y`$.
Let now $`\mathrm{\Sigma }_m^+`$ (respectively $`\mathrm{\Sigma }_l^+`$) be the positive roots in $`\mathrm{\Sigma }`$ of squared length $`2`$ (respectively $`4`$). Let
$$\rho _m=2\underset{i=1}{\overset{n}{}}(ni)ϵ_i,\rho _l=\underset{i=1}{\overset{n}{}}ϵ_i$$
be the sum of co-roots $`\alpha ^{}`$ with $`\alpha \mathrm{\Sigma }_m^+`$ and $`\alpha \mathrm{\Sigma }_l^+`$, respectively. Let $`\lambda \mathrm{\Lambda }`$ and set
$$\rho _m(\lambda )=\underset{\alpha \mathrm{\Sigma }_m^+}{}ϵ(\lambda ,\alpha )\alpha ^{},\rho _l(\lambda )=\underset{\alpha \mathrm{\Sigma }_l^+}{}ϵ(\lambda ,\alpha )\alpha ^{}.$$
Then $`\rho _m(\lambda )=\rho _m`$ and $`\rho _l(\lambda )=\rho _l`$ for all dominant weights $`\lambda \mathrm{\Lambda }^+`$. We define now $`\gamma _\lambda =\gamma _\lambda (𝐤,q)^n`$ by
$$\gamma _\lambda =t_0^{\rho _l(\lambda )}.t_n^{\rho _l(\lambda )}.t^{\rho _m(\lambda )}.q^\lambda ,\lambda \mathrm{\Lambda },$$
(4.3)
where the product is the usual dot-product in $`^n`$. In particular,
$$\gamma _\lambda =(t_0t_nt^{2(n1)}q^{\lambda _1},t_0t_nt^{2(n2)}q^{\lambda _2},\mathrm{},t_0t_nq^{\lambda _n}),\lambda \mathrm{\Lambda }^+.$$
(4.4)
Then lemma 4.3, (4.2) and (2.3) lead to the following proposition.
###### Proposition 4.5.
For $`\lambda \mathrm{\Lambda }`$ and $`f(Y)𝒜_Y`$, we have
$$f(Y)(x^\lambda )=f(\gamma _\lambda )x^\lambda +\underset{\mu \lambda }{}c_{\lambda ,\mu }x^\mu $$
for certain constants $`c_{\lambda ,\mu }`$.
We give now some properties of the diagonal terms $`\gamma _\lambda `$ ($`\lambda \mathrm{\Lambda }`$) which will be used frequently in the remainder of the paper. First of all, the diagonal terms of the $`Y`$-operators can be related to the spectrum of the $`Y`$-operators as described in \[24, def. 2.5\] by observing that
$$\rho _m(\lambda )=w_\lambda \rho _m,\rho _l(\lambda )=w_\lambda \rho _l,\lambda \mathrm{\Lambda },$$
(4.5)
where $`w_\lambda W`$ is the unique element of minimal length such that $`\lambda ^+=w_\lambda ^1\lambda `$, see \[23, prop. 2.10\].
Secondly, the action of $`𝒲`$ on the diagonal terms $`\gamma _\lambda `$ ($`\lambda \mathrm{\Lambda }`$), induced from the dot-action of $`𝒲`$ on $`\mathrm{\Lambda }`$, is compatible with the action of $`𝒲`$ on $`𝒜`$ in the following way. We refer to \[24, thm. 5.3\] for the proof.
###### Lemma 4.6.
Let $`f𝒜`$. Let $`\lambda \mathrm{\Lambda }`$ and $`i\{0,1,\mathrm{},n\}`$ such that $`s_i.\lambda \lambda `$. Then $`f(\gamma _{s_i.\lambda }^1)=(s_if)(\gamma _\lambda ^1)`$. If furthermore $`i1`$, then also $`f(\gamma _{s_i.\lambda })=f(\gamma _{s_i\lambda })=(s_if)(\gamma _\lambda )`$.
###### Remark 4.7.
Observe that the condition $`s_i.\lambda \lambda `$ in lemma 4.6 is always met for $`i=0`$. Let now $`i\{1,\mathrm{},n\}`$ and $`\lambda \mathrm{\Lambda }`$ with $`s_i.\lambda (=s_i\lambda )=\lambda `$. Then we have $`\gamma _\lambda ^{a_i}=\gamma _0^{a_i}`$ ($`=t^2`$ if $`i<n`$ and $`=t_0^2t_n^2`$ if $`i=n`$). Hence $`(s_if)(\gamma _\lambda ^{\pm 1})=f(\gamma _{\lambda ,i}^{\pm 1})`$ for $`f𝒜`$, with $`\gamma _{\lambda ,i}=\gamma _\lambda .(\gamma _0^{a_i})^{a_i^{}}`$. Observe that $`\gamma _{\lambda ,i}\gamma _\mu `$ for all $`\mu \mathrm{\Lambda }`$ by the generic conditions on the parameters.
By lemma 4.6 we have $`f(\gamma _\lambda )=\left(w_\lambda ^1f\right)(\gamma _{\lambda ^+})`$ for $`f𝒜`$ and $`\lambda \mathrm{\Lambda }`$. Combined with (4.4) we conclude that the diagonal terms $`\gamma _\lambda `$ ($`\lambda \mathrm{\Lambda }`$) are mutually different for generic values of $`q`$ and $`𝐤`$. This leads to the following main result of this section.
###### Theorem 4.8.
There exists a unique basis $`\{P_\lambda \}_{\lambda \mathrm{\Lambda }}`$ of $`𝒜`$ such that
1. $`P_\lambda (x)=x^\lambda +_{\mu \lambda }c_{\lambda ,\mu }x^\mu `$ for certain constants $`c_{\lambda ,\mu }`$,
2. $`f(Y)P_\lambda =f(\gamma _\lambda )P_\lambda `$ for all $`f(Y)𝒜_Y`$,
for all $`\lambda \mathrm{\Lambda }`$.
###### Definition 4.9.
The Laurent polynomial $`P_\lambda ()=P_\lambda (;𝐭;q)`$ ($`\lambda \mathrm{\Lambda }`$) is called the monic, non-symmetric Koornwinder polynomial of degree $`\lambda `$.
The terminology introduced in definition 4.9 stems from the close connection between the Laurent polynomials $`P_\lambda `$ ($`\lambda \mathrm{\Lambda }`$) and Koornwinder’s multivariable analogues of the Askey-Wilson polynomials, see and , as well as section 5 and section 6.
###### Remark 4.10.
(i) The second property of theorem 4.8 already characterizes the non-symmetric Koornwinder polynomial $`P_\lambda `$ up to a constant. This characterizing property was used by Sahi \[24, def. 6.1\] to introduce the non-symmetric Koornwinder polynomials. The triangularity of the non-symmetric Koornwinder polynomials was derived by Sahi \[25, sect. 6\] using recursion formulas.
(ii) If one uses Macdonald’s and Cherednik’s convention for the translation operator $`\tau `$ (see remark 2.1), then the role of $`Y^\lambda =T_{\tau (\lambda )}`$ is taken over by $`T_{\tau (\sigma \lambda )}=T_\sigma ^1Y^\lambda T_\sigma `$ for $`\lambda \mathrm{\Lambda }^+`$, where $`\sigma W`$ is the longest Weyl group element. The common eigenfunctions then become $`T_\sigma ^1P_\lambda `$ ($`\lambda \mathrm{\Lambda }`$). From lemma 4.3 and theorem 4.8 it follows that the $`T_\sigma ^1P_\lambda `$ are triangular with respect to the partial order on $`\mathrm{\Lambda }`$ in which $`\mu `$ is less than $`\nu `$ iff $`\sigma \mu \sigma \nu `$ (i.e. the anti-dominant weight is highest in each $`W`$-orbit). This is in accordance with the triangular structure of non-symmetric Macdonald polynomials, see and .
## 5. Symmetric Koornwinder polynomials
In this section we recall Noumi’s results on the affine Hecke algebraic characterization of Koornwinder’s multivariable analogues of the Askey-Wilson polynomials. We present Noumi’s results here in a different order by making use of the triangularity of the $`Y`$-operators, see proposition 4.5.
We write $`𝒜^W=\{f𝒜|wf=fwW\}`$, and similarly $`𝒜_Y^W`$, where the action is given by $`w(Y^\lambda )=Y^{w\lambda }`$ for $`wW`$ and $`\lambda \mathrm{\Lambda }`$. A linear basis of $`𝒜^W`$ and $`𝒜_Y^W`$ is given by the monomials $`m_\lambda (x)=_{\mu W\lambda }x^\mu `$, respectively $`m_\lambda (Y)=_{\mu W\lambda }Y^\mu `$ ($`\lambda \mathrm{\Lambda }^+`$).
It follows from (3.4) that $`𝒜_Y^W`$ lies in the center $`𝒵(H)`$ of $`H`$. In fact, by \[15, prop. 3.11\] we know that $`𝒵(H)=𝒜_Y^W`$ (which also follows from results in section 6). The action of $`𝒜_Y^W`$ on $`𝒜`$ through the Noumi representation $`\pi _{𝐭,q}`$ preserves $`𝒜^W`$, compare e.g. with \[18, 4.8\].
We extend the action of $`𝒲`$ on $`𝒜`$ to an action on the quotient field $`(x)=\text{Quot}(𝒜)`$ by requiring $`w𝒲`$ to be an automorphism of $`(x)`$. Let $`(x)[𝒲]\text{End}_{}((x))`$ be the subalgebra generated by $`(x)`$ (acting as multiplication operators) and by $`𝒲`$. Observe that $`(x)[𝒲]`$, and that
$$(x)[𝒲]=\underset{w𝒲}{}(x)w=\underset{wW,\lambda \mathrm{\Lambda }}{}(x)\tau (\lambda )w$$
as a $`(x)`$-submodule of $`\text{End}_{}((x))`$, see the proof of \[24, thm. 3.2\]. Furthermore, $`(x)[\tau (\mathrm{\Lambda })]=_{\lambda \mathrm{\Lambda }}(x)\tau (\lambda )`$ is the subalgebra of $`(x)[𝒲]`$ consisting of $`q`$-difference operators with coefficients in $`(x)`$.
With $`D(x)[𝒲]`$, say
$$D=\underset{wW}{}D(x,w)w,D(x,w)(x)[\tau (\mathrm{\Lambda })],$$
we associate a $`q`$-difference operator by
$$D_{sym}=\underset{wW}{}D(x,w)(x)[\tau (\mathrm{\Lambda })].$$
Observe that $`Df=D_{sym}f`$ if $`f(x)`$ is $`W`$-invariant. Proposition 4.5 and lemma 4.6 imply that the $`q`$-difference operators $`f(Y)_{sym}`$ ($`f𝒜^W`$) are triangular endomorphisms of $`𝒜^W`$:
$$f(Y)_{sym}m_\lambda =f(\gamma _\lambda )m_\lambda +\underset{\mu \mathrm{\Lambda }^+:\mu <\lambda }{}c_{\lambda ,\mu }m_\mu ,f𝒜^W,\lambda \mathrm{\Lambda }^+$$
for certain constants $`c_{\lambda ,\mu }`$. This immediately implies the following result.
###### Theorem 5.1.
There exists a unique basis $`\{P_\lambda ^+\}_{\lambda \mathrm{\Lambda }^+}`$ of $`𝒜^W`$ such that
1. $`P_\lambda ^+=m_\lambda +_{\mu \mathrm{\Lambda }^+:\mu <\lambda }c_{\lambda ,\mu }m_\mu `$ for certain constants $`c_{\lambda ,\mu }`$,
2. $`f(Y)_{sym}P_\lambda ^+=f(\gamma _\lambda )P_\lambda ^+`$ for all $`f(Y)𝒜_Y^W`$,
for all $`\lambda \mathrm{\Lambda }^+`$.
Noumi identified the $`q`$-difference operator
$$m_{ϵ_1}(Y)_{sym}=\left(Y_1+\mathrm{}+Y_n+Y_1^1+\mathrm{}+Y_n^1\right)_{sym}(x)[\tau (\mathrm{\Lambda })]$$
with Koornwinder’s second order $`q`$-difference operator. Explicitly, Noumi showed that the $`q`$-difference operator $`L=m_{ϵ_1}(Y)_{sym}m_{ϵ_1}(\gamma _0)`$ is given by
$$L=\underset{j=1}{\overset{n}{}}\left(\varphi _j^+(x)(\tau (ϵ_j)1)+\varphi _j^{}(x)(\tau (ϵ_j)1)\right)$$
(5.1)
with $`\varphi _j^{}(x)=\varphi _j^+(x_1^1,\mathrm{},x_n^1)`$ and
$$\begin{array}{cc}\hfill \varphi _j^+(x)=& (t_0t_n)^1t^{2(1n)}\frac{(1ax_j)(1bx_j)(1cx_j)(1dx_j)}{(1x_j^2)(1qx_j^2)}\hfill \\ & .\underset{ij}{}\frac{(1t^2x_ix_j)(1t^2x_i^1x_j)}{(1x_ix_j)(1x_i^1x_j)}.\hfill \end{array}$$
Here $`\{a,b,c,d\}`$ is related to the multiplicity function $`𝐭`$ by
$$\{a,b,c,d\}=\{t_0t_0^{}q^{1/2},t_0(t_0^{})^1q^{1/2},t_nt_n^{},t_n(t_n^{})^1\}.$$
(5.2)
Since the spectrum of $`L\text{End}_𝐂(𝒜^W)`$ is already simple (see ), this result implies that the $`W`$-invariant Laurent polynomials $`P_\lambda ^+`$ ($`\lambda \mathrm{\Lambda }^+`$) coincide with Koornwinder’s multivariable analogues of the Askey-Wilson polynomials.
###### Definition 5.2.
The $`W`$-invariant Laurent polynomial $`P_\lambda ^+()=P_\lambda ^+(;𝐭;q)`$ ($`\lambda \mathrm{\Lambda }^+`$) is called the monic, symmetric Koornwinder polynomial of degree $`\lambda \mathrm{\Lambda }^+`$.
## 6. The action of the Hecke algebra of type $`C_n`$
We associate a dual multiplicity function $`\stackrel{~}{𝐭}`$ with the multiplicity function $`𝐭`$ by interchanging the value of $`𝐭`$ on the $`𝒲`$-orbit $`𝒲a_0`$ with its value on the $`𝒲`$-orbit $`𝒲a_n^{}`$. In other words, $`\stackrel{~}{𝐭}`$ is the unique multiplicity function of $`S`$ satisfying
$$\stackrel{~}{t}_0=t_n^{},\stackrel{~}{t}_0^{}=t_0^{},\stackrel{~}{t}=t,\stackrel{~}{t}_n^{}=t_0,\stackrel{~}{t}_n=t_n.$$
We write $`\stackrel{~}{𝐤}`$ and $`\stackrel{~}{𝐤}^{}`$ for the associated multiplicity functions of $`R`$ and $`R^{}`$ respectively, and $`\stackrel{~}{v}_\beta (x)=v_\beta (x;\stackrel{~}{𝐭};q)`$ for the function $`v_\beta ()`$ (3.7) with respect to dual parameters. Observe that Lusztig’s formulas (3.4) can now be written in a uniform way:
$$T_if(Y)(s_if)(Y)T_i=\left((\stackrel{~}{t}_{a_i}\stackrel{~}{t}_{a_i}^1)+(\stackrel{~}{t}_{a_i/2}\stackrel{~}{t}_{a_i/2}^1)Y^{a_i/2}\right)\left(\frac{f(Y)(s_if)(Y)}{1Y^{a_i}}\right)$$
(6.1)
for $`i=1,\mathrm{},n`$ and $`f𝒜`$. In the following proposition we expand $`T_iP_\lambda `$ as a linear combination of non-symmetric Koornwinder polynomials.
###### Proposition 6.1.
Let $`i\{1,\mathrm{},n\}`$ and $`\lambda \mathrm{\Lambda }`$. Then
$$T_iP_\lambda =\xi _i(\gamma _\lambda )P_\lambda +\eta _i(\gamma _\lambda )P_{s_i\lambda }$$
(6.2)
with
$$\xi _i(x)=\stackrel{~}{t}_i\stackrel{~}{t}_i^1\stackrel{~}{v}_{a_i}(x)=\frac{\left(\stackrel{~}{t}_{a_i}^1\stackrel{~}{t}_{a_i}\right)x^{a_i}+\left(\stackrel{~}{t}_{a_i/2}^1\stackrel{~}{t}_{a_i/2}\right)x^{a_i/2}}{1x^{a_i}}$$
(6.3)
and
$$\eta _i(\gamma _\lambda )=\{\begin{array}{cc}\stackrel{~}{t}_i,\hfill & \text{ if }\lambda ,a_i<0,\hfill \\ \stackrel{~}{t}_i^3\stackrel{~}{v}_{a_i}(\gamma _\lambda )\stackrel{~}{v}_{a_i}(\gamma _\lambda ),\hfill & \text{ if }\lambda ,a_i0.\hfill \end{array}$$
(6.4)
###### Proof.
The proof is based on the following consequence of Lusztig’s formula (6.1) and theorem 4.8: let $`\lambda \mathrm{\Lambda }`$ and $`i\{1,\mathrm{},n\}`$, then
$$\left(f(Y)(s_if)(\gamma _\lambda )\right)T_iP_\lambda =\left(f(\gamma _\lambda )(s_if)(\gamma _\lambda )\right)\xi _i(\gamma _\lambda )P_\lambda ,f𝒜,$$
(6.5)
with $`\xi _i`$ given by (6.3).
Suppose now first that $`\lambda ,a_i=0`$, i.e. that $`s_i\lambda =\lambda `$. Using remark 4.7, we see that $`\xi _i(\gamma _\lambda )+\eta _i(\gamma _\lambda )=t_i`$, so we have to show that $`T_iP_\lambda =t_iP_\lambda `$. Now (6.5), theorem 4.8 and remark 4.7 imply that $`T_iP_\lambda `$ is a constant multiple of $`P_\lambda `$. The constant multiple can be determined by computing the leading coefficient of $`T_iP_\lambda `$ using the identity $`T_i=s_i(a_i)^1+t_it_i^1`$ and using lemma 4.3.
Suppose now that $`\lambda ,a_i0`$, then (6.5), lemma 4.6 and theorem 4.8 imply that $`T_iP_\lambda `$ is of the form (6.2) for some constant $`\eta _i(\gamma _\lambda )`$, with $`\xi _i`$ given by (6.3). If $`\lambda ,a_i<0`$, then leading term considerations using lemma 4.3 and theorem 4.8 show that $`\eta _i(\gamma _\lambda )=t_i=\stackrel{~}{t}_i`$. The expression for $`\eta _i(\gamma _\lambda )`$ when $`\lambda ,a_i>0`$ follows now easily by applying $`T_i`$ on both sides of (6.2) and using the quadratic relation $`(T_it_i)(T_i+t_i^1)=0`$, compare with \[22, prop. 4.1\] for the proof in the rank one setting. ∎
We write $`S_i=[T_i,Y^{a_i}]=T_iY^{a_i}Y^{a_i}T_i`$ ($`i=1,\mathrm{},n`$). It follows from (6.1) that the $`S_i`$ satisfy the fundamental commutation relations $`f(Y)S_i=S_i(s_if)(Y)`$ for all $`f𝒜`$, cf. \[24, sect. 5\]. We call $`S_i`$ the (non-affine) intertwiner associated with the simple reflection $`s_i`$.
We use here a slightly different definition for the intertwiners $`S_i`$ compared with Sahi’s , intertwiners. The advantage of the present definition is that the $`S_i`$ ($`i=1,\mathrm{},n`$) satisfy the $`C_n`$-braid relations, see remark 7.6. In particular, we may write $`S_w=S_{i_1}\mathrm{}S_{i_r}`$ for a reduced expression $`w=s_{i_1}\mathrm{}s_{i_r}W`$, and $`S_w`$ satisfies the intertwining property $`S_wf(Y)=(wf)(Y)S_w`$ for all $`f𝒜`$. See also the paper , in which yet another definition for the intertwiners $`S_i`$ ($`i=0,\mathrm{},n`$) is used (including a non-affine intertwiner $`S_0`$). The intertwiners in , which satisfy the $`\stackrel{~}{C}_n`$-braid relations, are used to prove a Rodrigues type formula for non-symmetric Koornwinder polynomials.
With our present conventions, the action of the intertwiners on the non-symmetric Koornwinder polynomials is easily determined from proposition 6.1. We give here only the action of the intertwiner $`S_i`$ corresponding to a simple reflection $`s_i`$ of $`W`$.
###### Corollary 6.2.
$`S_iP_\lambda =\left(\gamma _\lambda ^{a_i}\gamma _{s_i\lambda }^{a_i}\right)\eta _i(\gamma _\lambda )P_{s_i\lambda }`$ for $`i=1,\mathrm{},n`$ and $`\lambda \mathrm{\Lambda }`$.
###### Remark 6.3.
Proposition 6.1 refines the non-affine part of Sahi’s recursion formula \[25, thm. 18\], while corollary 6.2 can be seen as a refinement of the non-affine part of \[24, thm. 5.3\]. Indeed, in \[25, thm. 18\] and \[24, thm. 5.3\] the formulas are given up to an unknown multiple constant. The unknown constants in the affine part of \[25, thm. 18\] and \[24, thm. 5.3\] can also be computed, but this requires the duality properties and the evaluation formulas for the non-symmetric Koornwinder polynomials, see proposition 7.8 and theorem 9.3.
By theorem 4.8, proposition 6.1 and corollary 6.2 we have a complete description of the action of $`H(R;𝐤)`$ on the non-symmetric Koornwinder polynomials under the Noumi representation $`\pi _{𝐭,q}`$. From this the $`H`$-module structure of $`𝒜`$ can be described in detail, see also Sahi . The result is as follows. We write
$$𝒜=\underset{\lambda \mathrm{\Lambda }^+}{}𝒜(\lambda ),𝒜(\lambda )=\text{span}\{P_\mu |\mu W\lambda \}.$$
(6.6)
Recall that the parameters $`𝐭`$ and $`q`$ are assumed to be generic.
###### Theorem 6.4.
The direct sum decomposition (6.6) is the multiplicity-free, irreducible decomposition of $`𝒜`$ as a $`(\pi _{𝐭,q},H(R;𝐤))`$-module. Furthermore, (6.6) is the decomposition of $`𝒜`$ into isotypical components under the action of the center $`𝒵(H(R;𝐤))=𝒜_Y^W`$. The central character $`\chi _\lambda `$ of $`𝒜(\lambda )`$ is given by $`\chi _\lambda (f)=f(\gamma _\lambda )`$ for $`f𝒵(H(R;𝐤))`$.
We end this section by defining anti-symmetric Koornwinder polynomials and by expanding (anti-)symmetric Koornwinder polynomials in terms of non-symmetric Koornwinder polynomials.
We associate a function $`\{t_w\}_{w𝒲}`$ with the multiplicity function $`𝐭=\{t_\beta \}_{\beta S}`$ by defining $`t_w=t_{i_1}\mathrm{}t_{i_r}`$ for a reduced expression $`w=s_{i_1}\mathrm{}s_{i_r}𝒲`$.
###### Remark 6.5.
Recall the well-known fact that $`R^+w^1R^{}=\{\beta _1,\mathrm{},\beta _r\}`$ with the $`r`$ distinct positive roots $`\beta _j`$ given by
$$\beta _j=s_{i_r}\mathrm{}s_{i_{j+1}}a_{i_j}(j=1,\mathrm{},r1),\beta _r=a_{i_r}.$$
In particular, it follows that
$$t_w=\underset{\beta R^+w^1R^{}}{}t_\beta ,w𝒲.$$
Restricted to the finite Weyl group $`W`$, the expression for $`t_w`$ reduces to
$$t_w=\underset{\alpha \mathrm{\Sigma }^+w^1\mathrm{\Sigma }^{}}{}t_\alpha ,wW.$$
Observe in particular that $`\stackrel{~}{t}_w=t_w`$ when $`wW`$, where $`\{\stackrel{~}{t}_w\}_{w𝒲}`$ is the function associated with the dual multiplicity function $`\stackrel{~}{𝐭}`$.
Let $`\chi _\pm :H_0`$ be the trivial and alternating character of the Hecke algebra $`H_0`$, i.e. $`\chi _\pm (T_i)=\pm t_i^{\pm 1}`$ for $`i=1,\mathrm{},n`$. Then the corresponding mutually orthogonal, primitive idempotents are given by
$$C_\pm =\frac{1}{_{wW}t_w^{\pm 2}}\underset{wW}{}(\pm 1)^{l(w)}t_w^{\pm 1}T_w.$$
(6.7)
We define for $`\lambda \mathrm{\Lambda }^+`$,
$$\begin{array}{cc}\hfill 𝒜_\pm (\lambda )& =\{f𝒜(\lambda )|C_\pm f=f\}\hfill \\ & =\{f𝒜(\lambda )|(T_it_i^{\pm 1})f=0i=1,\mathrm{},n\}.\hfill \end{array}$$
(6.8)
Observe in particular that $`𝒜_+(\lambda )=𝒜(\lambda )𝒜^W`$. Let $`\mathrm{\Lambda }^{++}=\kappa +\mathrm{\Lambda }^+`$ be the cone of regular dominant weights, where
$$\kappa =\frac{1}{2}\underset{\alpha \mathrm{\Sigma }^+}{}\alpha =\underset{i=1}{\overset{n}{}}\omega _i.$$
(6.9)
###### Theorem 6.6.
(i) $`𝒜_+(\lambda )`$ is spanned by the symmetric Koornwinder polynomial $`P_\lambda ^+`$ for all $`\lambda \mathrm{\Lambda }^+`$.
(ii) $`𝒜_{}(\lambda )`$ is one-dimensional if $`\lambda \mathrm{\Lambda }^{++}`$ and zero dimensional otherwise. For $`\lambda \mathrm{\Lambda }^{++}`$ there exists a unique $`P_\lambda ^{}𝒜_{}(\lambda )`$ of the form $`P_\lambda ^{}=x^\lambda +_{\mu \lambda }e_\mu x^\mu `$ for certain constants $`e_\mu `$.
(iii) The expressions for $`P_\lambda ^+`$ ($`\lambda \mathrm{\Lambda }^+`$) and for $`P_\lambda ^{}`$ ($`\lambda \mathrm{\Lambda }^{++}`$) as linear combinations of the non-symmetric Koornwinder polynomials $`P_\mu `$ ($`\mu W\lambda `$) are given by
$$P_\lambda ^\pm =\underset{\mu W\lambda }{}c_{\lambda ,\mu }^\pm P_\mu $$
with the coefficients $`c_{\lambda ,\mu }^\pm `$ ($`\mu W\lambda `$) given by
$$\begin{array}{cc}\hfill c_{\lambda ,\mu }^\pm & =(\pm 1)^{l(w_\mu )}\stackrel{~}{t}_{w_\mu }^2\underset{\stackrel{\alpha \mathrm{\Sigma }^+}{\mu ,\alpha <0}}{}\stackrel{~}{v}_{\pm \alpha }(\gamma _\mu )\hfill \\ & =(\pm 1)^{l(w_\mu )}\stackrel{~}{t}_{w_\mu }^2\underset{\alpha \mathrm{\Sigma }^+w_\mu ^1\mathrm{\Sigma }^{}}{}\stackrel{~}{v}_{\pm \alpha }(\gamma _\lambda ^1),\hfill \end{array}$$
(6.10)
where $`w_\mu `$ is the element of minimal length in $`W`$ such that $`w_\mu \lambda (=w_\mu \mu ^+)=\mu `$.
###### Proof.
We first introduce some notations and deduce some preliminary results which we will need for the proof.
Let $`\lambda \mathrm{\Lambda }^+`$ and let $`Q_\lambda ^\pm =_{\mu W\lambda }d_{\lambda ,\mu }^\pm P_\mu `$ be an element in $`𝒜_\pm (\lambda )`$. By proposition 6.1 we have for $`\mu \mathrm{\Lambda }`$ and $`i\{1,\mathrm{},n\}`$ that
$$\left(T_it_i^{\pm 1}\right)P_\mu =\xi _i^\pm (\gamma _\mu )P_\mu +\eta _i(\gamma _\mu )P_{s_i\mu }$$
with
$$\xi _i^\pm (x)=\xi _i(x)t_i^{\pm 1}=\stackrel{~}{t}_i^1\stackrel{~}{v}_{a_i}(x).$$
(6.11)
It follows now from the relations $`(T_it_i^{\pm 1})Q_\lambda ^\pm =0`$ that the coefficients $`d_{\lambda ,\mu }^\pm `$ satisfy the recurrence relations
$$d_{\lambda ,s_i\mu }^\pm =\left(\frac{\xi _i^\pm (\gamma _\mu )}{\eta _i(\gamma _{s_i\mu })}\right)d_{\lambda ,\mu }^\pm $$
(6.12)
for $`\mu W\lambda `$ and $`i\{1,\mathrm{},n\}`$ such that $`\mu ,a_i0`$. Let now $`\mu W\lambda `$ and choose a reduced expression $`w_\mu =s_{i_1}\mathrm{}s_{i_r}`$. We set
$$\mu _j=s_{i_{j+1}}\mathrm{}s_{i_r}\lambda (j=0,\mathrm{},r1),\mu _r=\lambda .$$
Then $`\mu _j,a_{i_j}>0`$ for $`j=1,\mathrm{},r`$ by remark 6.5. Iterating (6.12), we thus obtain
$$d_{\lambda ,\mu }^\pm =d_{\lambda ,\lambda }^\pm \underset{j=1}{\overset{r}{}}\left(\frac{\xi _{i_j}^\pm (\gamma _{\mu _j})}{\eta _{i_j}(\gamma _{\mu _{j1}})}\right).$$
(6.13)
Proof of (i) The recurrence formula (6.13) implies that $`\text{dim}(𝒜_+(\lambda ))1`$ for all $`\lambda \mathrm{\Lambda }^+`$. On the other hand, theorem 5.1 implies that the symmetric Koornwinder polynomial $`P_\lambda ^+`$ is a non-zero element in $`𝒜_+(\lambda )`$ for all $`\lambda \mathrm{\Lambda }^+`$.
Proof of (ii) Again by (6.13), we have $`\text{dim}(𝒜_{}(\lambda ))1`$. Let $`\lambda \mathrm{\Lambda }^+\mathrm{\Lambda }^{++}`$. Let $`s_iW`$ be a simple reflection in $`W`$ which stabilizes $`\lambda `$. Then proposition 6.1 implies that $`T_iP_\lambda =t_iP_\lambda `$. Hence the coefficient of $`P_\lambda `$ in the expansion of $`(T_i+t_i^1)Q_\lambda ^{}`$ as a linear combination of the $`P_\nu `$ ($`\nu W\lambda `$), is $`(t_i+t_i^1)d_{\lambda ,\lambda }^{}`$. On the other hand, $`(T_i+t_i^1)Q_\lambda ^{}=0`$, hence we conclude that $`d_{\lambda ,\lambda }^{}=0`$. Then (6.13) implies $`d_{\lambda ,\mu }^{}=0`$ for all $`\mu W\lambda `$, hence $`Q_\lambda ^{}=0`$. This proves that $`𝒜_{}(\lambda )=\{0\}`$ if $`\lambda \mathrm{\Lambda }^+\mathrm{\Lambda }^{++}`$.
Let now $`\lambda \mathrm{\Lambda }^{++}`$. Then $`C_{}P_{\sigma \lambda }𝒜_{}(\lambda )`$, where $`\sigma W`$ is the longest Weyl group element. Furthermore, by proposition 6.1 we have $`C_{}P_{\sigma \lambda }=_{\mu W\lambda }d_\mu P_\mu `$ with $`d_\lambda =(_{wW}t_w^2)^1(1)^{l(\sigma )}0`$, so $`C_{}P_{\sigma \lambda }`$ is a non-zero element in $`𝒜_{}(\lambda )`$. Hence $`\text{dim}(𝒜_{}(\lambda ))=1`$ if $`\lambda \mathrm{\Lambda }^{++}`$. The triangularity statement follows from theorem 4.8 and from the fact that the coefficient $`d_\lambda `$ in the above expansion of $`C_{}P_{\sigma \lambda }`$ is non-zero.
Proof of (iii) We have to show that
$$d_{\lambda ,\mu }^\pm =d_{\lambda ,\lambda }^\pm c_{\lambda ,\mu }^\pm ,\mu W\lambda ,$$
(6.14)
where the $`d_{\lambda ,\mu }^\pm `$ are the expansion coefficients of $`Q_\lambda ^\pm `$ in terms of non-symmetric Koornwinder polynomials $`P_\mu `$ ($`\mu W\lambda `$), and where $`c_{\lambda ,\mu }^\pm `$ is given by (6.10). We use again the recurrence formula (6.13) for the coefficients $`d_{\lambda ,\mu }^\pm `$. We set
$$\alpha _1=a_{i_1},\alpha _j=s_{i_1}\mathrm{}s_{i_{j1}}a_{i_j}(j=2,\mathrm{},r),$$
then the $`\alpha _j`$ ($`j=1,\mathrm{},r`$) are mutually different and
$$\{\alpha _1,\mathrm{},\alpha _r\}=\mathrm{\Sigma }^{}w_\mu \mathrm{\Sigma }^+=\{\alpha \mathrm{\Sigma }^{}|\mu ,\alpha >0\},$$
(6.15)
see remark 6.5. Now observe that $`\mu _{j1},a_{i_j}=\mu _j,a_{i_j}<0`$ for all $`j=1,\mathrm{},r`$, so that
$$\eta _{i_j}(\gamma _{\mu _{j1}})=\stackrel{~}{t}_{\alpha _j}(j=1,\mathrm{},r).$$
(6.16)
Furthermore, observe that
$$(\gamma _{\mu _j})^{a_{i_j}}=(\gamma _\mu )^{\alpha _j}=(\gamma _\lambda )^{w_\mu ^1\alpha _j},j=1,\mathrm{},r$$
(6.17)
by lemma 4.6. Substituting (6.16) and (6.11) in (6.13) and using (6.17) and the characterization of the roots $`\{\alpha _j\}_{j=1}^r`$ (see (6.15)), we obtain (6.14). ∎
###### Definition 6.7.
The Laurent polynomial $`P_\lambda ^{}𝒜_{}(\lambda )`$ ($`\lambda \mathrm{\Lambda }^{++}`$) is called the anti-symmetric Koornwinder polynomial of degree $`\lambda `$.
###### Remark 6.8.
(i) Part (i) of theorem 6.6 was also observed by Sahi \[24, cor. 6.6\].
(ii) Theorem 6.6 extends Macdonald’s \[18, sect. 6\] explicit expansion formulas for the (anti-)symmetric Macdonald polynomials associated with root systems of classical type.
## 7. Spectral difference-reflection operators and duality
Let $`x_\lambda =\gamma _\lambda (\stackrel{~}{𝐤},q)`$ ($`\lambda \mathrm{\Lambda }`$) be the spectrum (4.3) of the $`\stackrel{~}{Y}`$-operators, and denote $`\stackrel{~}{}=(S;\stackrel{~}{𝐭};q)`$ for the double affine Hecke algebra with respect to dual parameters. We define evaluation mappings $`\text{Ev}:`$ and $`\stackrel{~}{\text{Ev}}:\stackrel{~}{}`$ by
$$\text{Ev}(X)=\left(X(1)\right)(x_0^1),\stackrel{~}{\text{Ev}}(\stackrel{~}{X})=\left(\stackrel{~}{X}(1)\right)(\gamma _0^1)$$
for $`X`$ and $`\stackrel{~}{X}\stackrel{~}{}`$, where $`1𝒜`$ is the Laurent polynomial identically equal to one.
The evaluation $`\text{Ev}(P_\lambda (z))=P_\lambda (x_0^1)`$ of the non-symmetric monic Koornwinder polynomial $`P_\lambda ()=P_\lambda (;𝐭;q)`$ is generically non-zero by the analytic dependence of $`P_\lambda `$ on $`𝐭`$ and $`q`$ (we use here that $`P_\lambda (x)=x^\lambda `$ when $`t_a=1`$ for all $`aS`$). Similarly, $`\text{Ev}(P_\lambda ^+(z))=P_\lambda ^+(x_0^{\pm 1})`$ is non-zero for generic parameter values $`𝐭`$ and $`q`$. In section 9 we explicitly evaluate $`P_\lambda (x_0^1)`$ and $`P_\lambda ^+(x_0)=P_\lambda ^+(x_0^1)`$, so that the generic conditions on the parameters can be made completely explicit.
###### Definition 7.1.
(i) Let $`E(\gamma _\lambda ;)=E(\gamma _\lambda ;;𝐭;q)`$ be the constant multiple of the non-symmetric Koornwinder polynomial $`P_\lambda ()`$ of degree $`\lambda \mathrm{\Lambda }`$ which takes the value one at $`x=x_0^1`$.
(ii) Let $`E^+(\gamma _\lambda ;)=E^+(\gamma _\lambda ;;𝐭;q)`$ be the constant multiple of the symmetric Koornwinder polynomial $`P_\lambda ^+()`$ of degree $`\lambda \mathrm{\Lambda }^+`$ which takes the value one at $`x=x_0`$.
Sahi showed that the role of the geometric parameter $`x=x_\mu `$ and of the spectral parameter $`\gamma =\gamma _\lambda `$ are (in a suitable sense) interchangeable for the renormalized Koornwinder polynomials $`E(\gamma ;x^1)`$ and $`E^+(\gamma ;x)`$, see also van Diejen for a sub-class of the symmetric Koornwinder polynomials. These duality properties stem from a particular anti-algebra isomorphism of the double affine Hecke algebra $``$, which we define now first.
Recall the notations $`T_0^{}=T_0^1z^{a_0^{}}`$ and $`T_n^{}=z^{a_n^{}}T_n^1`$ for the simple generators associated with $`a_0^{}`$ and $`a_n^{}`$ respectively, see remark 3.5. We set
$$U_n=T_1T_2\mathrm{}T_{n1}T_n^{}T_{n1}^1\mathrm{}T_2^1T_1^1,$$
which is a conjugate of $`T_n^{}`$ in $`(S;𝐭;q)`$.
Set $`𝐭^1=(t_\beta ^1)_{\beta S}`$ for the inverse of the multiplicity function $`𝐭`$. We write $`T_i^{},T_j^{}`$, $`U_n^{}`$, $`Y^\lambda `$ $`z^\lambda `$ for the elements $`T_i,T_j^{}`$, $`U_n`$, $`Y^\lambda `$ and $`z^\lambda `$ in the double affine Hecke algebra $`^{}=(S;𝐭^1;q^1)`$. Similarly, we write $`\stackrel{~}{T}_i,\mathrm{}`$ (respectively $`\stackrel{~}{T}_i^{},\mathrm{}`$) for the elements $`T_i,\mathrm{}`$ in the double affine Hecke algebra $`\stackrel{~}{}`$ (respectively $`\stackrel{~}{}^{}=(S;\stackrel{~}{𝐭}^1;q^1)`$). The following theorem was proved by Sahi \[24, thm. 4.2\].
###### Theorem 7.2.
There exists a unique algebra isomorphism $`ϵ=ϵ_{𝐭,q}:\stackrel{~}{}^{}`$ satisfying $`ϵ(T_0)=(\stackrel{~}{U}_n^{})^1`$, $`ϵ(z_i)=\stackrel{~}{Y}_i^{}`$ and $`ϵ(T_i)=(\stackrel{~}{T}_i^{})^1`$ for $`i=1,\mathrm{},n`$. Furthermore, $`ϵ_{𝐭,q}^1=ϵ_{\stackrel{~}{𝐭}^1,q^1}`$.
The isomorphism $`ϵ`$ is a crucial building block for Sahi’s duality anti-isomorphism of the double affine Hecke algebra $``$. In fact, the duality anti-isomorphism is obtained by composing $`ϵ`$ with the anti-isomorphism $``$ defined in the following lemma.
###### Lemma 7.3.
There exists a unique algebra isomorphism $`=_{𝐭,q}:^{}`$ (respectively anti-algebra isomorphism $`=_{𝐭,q}:^{}`$) satisfying $`T_i(T_i^{})^1`$ ($`i=0,\mathrm{},n`$) and $`z_j(z_j^{})^1`$ ($`j=1,\mathrm{},n`$).
###### Proof.
This follows directly from the presentation of $``$ as given by Sahi \[24, sect. 3\]. ∎
###### Remark 7.4.
(i) Lemma 7.3 for $``$ was observed by Sahi \[24, prop. 7.1\].
(ii) In proposition 8.3 we interpret the anti-algebra isomorphism $``$ as a $``$-structure on $`\text{End}_{}(𝒜)`$ induced from a suitable non-degenerate bilinear form on $`𝒜`$.
We write $`\stackrel{~}{}^{}`$ (respectively $`\stackrel{~}{}^{}`$) for $``$ (respectively $``$) with respect to the parameters $`(\stackrel{~}{𝐭}^1,q^1)`$.
###### Definition 7.5.
(i) The algebra isomorphism $`\mathrm{\Phi }=\mathrm{\Phi }_{𝐭,q}=\stackrel{~}{}^{}ϵ:\stackrel{~}{}`$ is called the duality isomorphism of $``$.
(ii) The anti-algebra isomorphism $`\mathrm{\Psi }=\mathrm{\Psi }_{𝐭,q}=\stackrel{~}{}^{}ϵ:\stackrel{~}{}`$ is called the duality anti-isomorphism of $``$.
Observe that $`\mathrm{\Phi }`$ (respectively $`\mathrm{\Psi }`$) is uniquely characterized as the (anti-)algebra homomorphism $`\stackrel{~}{}`$ which maps $`U_n`$ to $`\stackrel{~}{T}_0`$, $`T_i`$ to $`\stackrel{~}{T}_i`$ and $`Y_i`$ to $`\stackrel{~}{z}_i^1`$ for $`i=1,\mathrm{},n`$. By \[24, sect. 7\], the inverse of $`\mathrm{\Psi }=\mathrm{\Psi }_{𝐭,q}`$ is given by $`\stackrel{~}{\mathrm{\Psi }}=\mathrm{\Psi }_{\stackrel{~}{𝐭},q}`$.
###### Remark 7.6.
Observe that the image of the non-affine intertwiners $`S_i=[T_i,Y^{a_i}]`$ ($`i=1,\mathrm{},n`$) under $`\mathrm{\Psi }`$ is given by
$$\mathrm{\Psi }(S_i)=\stackrel{~}{t}_i^1(\stackrel{~}{z}^{a_i}\stackrel{~}{z}^{a_i})\stackrel{~}{v}_{a_i}(\stackrel{~}{z})s_i\stackrel{~}{}$$
in view of the explicit expression for $`\stackrel{~}{T}_i`$ (see theorem 3.2). In particular, $`\mathrm{\Psi }(S_i)`$ is of the form $`f_i(\stackrel{~}{z}^{a_i^{}})s_i`$, with $`f_i`$ a Laurent polynomial in one variable and with $`f_1=f_2=\mathrm{}=f_{n1}`$. From these facts it is easy to prove that $`(S_1,\mathrm{},S_n)`$ satisfies the $`C_n`$-braid relations in $``$.
The two evaluation mappings Ev and $`\stackrel{~}{\text{Ev}}`$ are related via the duality anti-isomorphism:
$$\stackrel{~}{\text{Ev}}\left(\mathrm{\Psi }(X)\right)=\text{Ev}\left(X\right),X,$$
see \[24, thm. 7.3\]. This implies that the two pairings $`B:\times \stackrel{~}{}`$ and $`\stackrel{~}{B}:\stackrel{~}{}\times `$ defined by $`B(X,\stackrel{~}{X})=\text{Ev}\left(\stackrel{~}{\mathrm{\Psi }}(\stackrel{~}{X})X\right)`$ and $`\stackrel{~}{B}(\stackrel{~}{X},X)=\stackrel{~}{\text{Ev}}\left(\mathrm{\Psi }(X)\stackrel{~}{X}\right)`$ for $`X`$ and $`\stackrel{~}{X}\stackrel{~}{}`$ satisfy the duality property
$$B(X,\stackrel{~}{X})=\stackrel{~}{B}(\stackrel{~}{X},X),X,\stackrel{~}{X}\stackrel{~}{}.$$
(7.1)
Before we recall how (7.1) implies the duality properties of the Koornwinder polynomials, we first collect some elementary identities for the bilinear form $`B`$. The proof of the lemma is similar to the proof in the rank one setting, see \[22, lem. 10.5\].
###### Lemma 7.7.
Let $`f𝒜`$. Let $`X,X_1,X_2`$ and $`\stackrel{~}{X},\stackrel{~}{X}_1,\stackrel{~}{X}_2\stackrel{~}{}`$.
(i) $`B(X_1X_2,\stackrel{~}{X})=B(X_2,\mathrm{\Psi }(X_1)\stackrel{~}{X})`$ and $`B(X,\stackrel{~}{X}_1\stackrel{~}{X}_2)=B(\stackrel{~}{\mathrm{\Psi }}(\stackrel{~}{X}_1)X,\stackrel{~}{X}_2)`$.
(ii) $`B(XT_i,\stackrel{~}{X})=t_iB(X,\stackrel{~}{X})`$ for $`i=0,\mathrm{},n`$.
(iii) $`B((X(f))(z),\stackrel{~}{X})=B(Xf(z),\stackrel{~}{X})`$ and $`B(X,(\stackrel{~}{X}(f))(\stackrel{~}{z}))=B(X,\stackrel{~}{X}f(\stackrel{~}{z}))`$, where $`(X(f))(z)`$ is the multiplication operator in $``$ corresponding to the Laurent polynomial $`X(f)𝒜`$, and $`Xf(z)`$ is the product of the elements $`X`$ and $`f(z)`$ in $``$.
We write $`\stackrel{~}{E}(x_\lambda ;)`$ for the renormalized non-symmetric Koornwinder polynomial $`E(x_\lambda ;;\stackrel{~}{𝐭};q)`$, and similarly for $`\stackrel{~}{E}^+(x_\lambda ;)`$. Observe now that by lemma 7.7 and theorem 4.8,
$$f(\gamma _\lambda ^1)=\stackrel{~}{B}(f(\stackrel{~}{z}),E(\gamma _\lambda ;z)),g(x_\mu ^1)=B(g(z),\stackrel{~}{E}(x_\mu ;\stackrel{~}{z}))$$
(7.2)
for $`f,g𝒜`$ and $`\lambda ,\mu \mathrm{\Lambda }`$. Taking $`f=\stackrel{~}{E}(x_\mu ;)`$ and $`g=E(\gamma _\lambda ;)`$ and using the duality (7.1) for the pairing, we arrive at
$$E(\gamma _\lambda ;x_\mu ^1)=\stackrel{~}{E}(x_\mu ;\gamma _\lambda ^1),\lambda ,\mu \mathrm{\Lambda }$$
(7.3)
which is the duality for the renormalized Koornwinder polynomials, see \[24, thm. 7.4\]. Similarly, we derive from theorem 5.1 and (7.1) that
$$E^+(\gamma _\lambda ;x_\mu )=\stackrel{~}{E}^+(x_\mu ;\gamma _\lambda ),\lambda ,\mu \mathrm{\Lambda }^+,$$
(7.4)
see \[24, cor. 7.5\]. Using the duality (7.3), we can rewrite the action of $`T_i`$ ($`i=1,\mathrm{},n`$) and $`U_n`$ on the renormalized Koornwinder polynomials $`E(\gamma ;)`$ in terms of difference-reflection operators acting on the spectral parameter $`\gamma \text{Spec}(Y)=\{\gamma _\lambda |\lambda \mathrm{\Lambda }\}`$. Define an action of $`𝒲`$ on $`\text{Spec}(Y)`$ by $`w\gamma _\lambda =\gamma _{w.\lambda }`$ ($`\lambda \mathrm{\Lambda }`$, $`w𝒲`$).
###### Proposition 7.8.
(i) For $`\gamma \text{Spec}(Y)`$ we have
$$\left(U_nE(\gamma ;)\right)(x)=\stackrel{~}{t}_0E(\gamma ;x)+\stackrel{~}{t}_0^1\stackrel{~}{v}_{a_0}(\gamma ^1)\left(E(s_0\gamma ;x)E(\gamma ;x)\right).$$
(ii) For $`i=1,\mathrm{},n`$ and $`\gamma \text{Spec}(Y)`$ we have
$$\left(T_iE(\gamma ;)\right)(x)=\stackrel{~}{t}_iE(\gamma ;x)+\stackrel{~}{t}_i^1\stackrel{~}{v}_{a_i}(\gamma ^1)\left(E(s_i\gamma ;x)E(\gamma ;x)\right).$$
###### Proof.
By (7.2) and lemma 7.7 we have
$$B(E(\gamma _\lambda ;z),\stackrel{~}{T}_i\stackrel{~}{E}(x_\mu ;\stackrel{~}{z}))=\{\begin{array}{cc}\left(U_nE(\gamma _\lambda ;)\right)(x_\mu ^1)\hfill & \text{ if }i=0\hfill \\ \left(T_iE(\gamma _\lambda ;)\right)(x_\mu ^1)\hfill & \text{ if }i=1,\mathrm{},n\hfill \end{array}$$
for all $`\lambda ,\mu \mathrm{\Lambda }`$. So it suffices to prove that
$$\begin{array}{cc}\hfill B(E(\gamma _\lambda ;z),\stackrel{~}{T}_i\stackrel{~}{E}(x_\mu ;\stackrel{~}{z}))& =\stackrel{~}{t}_iE(\gamma _\lambda ;x_\mu ^1)\hfill \\ & +\stackrel{~}{t}_i^1\stackrel{~}{v}_{a_i}(\gamma _\lambda ^1)\left(E(\gamma _{s_i.\lambda };x_\mu ^1)E(\gamma _\lambda ;x_\mu ^1)\right)\hfill \end{array}$$
(7.5)
for all $`\lambda ,\mu \mathrm{\Lambda }`$ and all $`i=0,\mathrm{},n`$.
Formula (7.5) is easy when $`\lambda `$ is stabilized by $`s_i`$ since then we have
$$B(E(\gamma _\lambda ;z),\stackrel{~}{T}_i\stackrel{~}{E}(x_\mu ;\stackrel{~}{z}))=B((T_i(E(\gamma _\lambda ;.))(z),\stackrel{~}{E}(x_\mu ;\stackrel{~}{z}))=\stackrel{~}{t}_iE(\gamma _\lambda ;x_\mu ^1)$$
where the last equality follows from (the proof of) proposition 6.1 and (7.2). So we assume for the remainder of the proof that $`s_i.\lambda \lambda `$. We can use now (3.8) to commute $`\stackrel{~}{T}_i`$ and $`\stackrel{~}{E}(x_\mu ;\stackrel{~}{z})`$ in the left-hand side of (7.5). Combined with lemma 7.7 we then derive that
$$\begin{array}{cc}\hfill B(E(\gamma _\lambda ;z),\stackrel{~}{T}_i\stackrel{~}{E}(x_\mu ;\stackrel{~}{z}))& =\stackrel{~}{t}_i\left(s_i\stackrel{~}{E}(x_\mu ;)\right)(\gamma _\lambda ^1)\hfill \\ & +\psi _i(\gamma _\lambda ^1)\left(\stackrel{~}{E}(x_\mu ;\gamma _\lambda ^1)(s_i\stackrel{~}{E}(x_\mu ;))(\gamma _\lambda ^1)\right),\hfill \end{array}$$
where $`\psi _i(x)`$ is given by
$$\psi _i(x)=\frac{(\stackrel{~}{t}_{a_i}\stackrel{~}{t}_{a_i}^1)+(\stackrel{~}{t}_{a_i/2}\stackrel{~}{t}_{a_i/2}^1)x^{a_i/2}}{1x^{a_i}}=\stackrel{~}{t}_i\stackrel{~}{t}_i^1\stackrel{~}{v}_{a_i}(x).$$
Since $`s_i.\lambda \lambda `$, we can apply lemma 4.6 together with the duality (7.3) of the non-symmetric Koornwinder polynomials to obtain the desired formula (7.5). ∎
## 8. (Bi-)orthogonality relations and quadratic norms
From now on we assume that $`0<q,t<1`$ and that the parameters $`a,b,c,d`$ (see (5.2)) have moduli less than one. We define $`\mathrm{\Delta }()=\mathrm{\Delta }(;𝐭;q)`$ and $`\mathrm{\Delta }_+()=\mathrm{\Delta }_+(;𝐭;q)`$ by
$$\mathrm{\Delta }(x;𝐭;q)=\underset{\beta R^+}{}\frac{1}{v_\beta (x;𝐭;q)}$$
(8.1)
and
$$\begin{array}{cc}\hfill \mathrm{\Delta }_+(x;𝐭;q)=& \underset{\stackrel{\beta R}{\beta (0)0}}{}\frac{1}{v_\beta (x;𝐭;q)}\hfill \\ \hfill =& \underset{1i<jn}{}\frac{(x_ix_j,x_ix_j^1,x_i^1x_j,x_i^1x_j^1;q)_{\mathrm{}}}{(t^2x_ix_j,t^2x_ix_j^1,t^2x_i^1x_j,t^2x_i^1x_j^1;q)_{\mathrm{}}}\hfill \\ & .\underset{i=1}{\overset{n}{}}\frac{(x_i^2,x_i^2;q)_{\mathrm{}}}{(ax_i,ax_i^1,bx_i,bx_i^1,cx_i,cx_i^1,dx_i,dx_i^1;q)_{\mathrm{}}},\hfill \end{array}$$
(8.2)
where $`(y_1,\mathrm{},y_m;q)_{\mathrm{}}=_{j=1}^m(y_j;q)_{\mathrm{}}`$ with $`(y;q)_{\mathrm{}}=_{j=0}^{\mathrm{}}(1yq^j)`$ the $`q`$-shifted factorial. The second equality in (8.2) follows from the $`𝒲`$-orbit structure of the reduced affine root system $`R`$ (cf. (2.3)), together with (5.2). Observe that $`\mathrm{\Delta }_+()`$ is $`W`$-invariant, where $`W`$ acts by permutations and inversions of the coordinates $`x=(x_1,\mathrm{},x_n)`$. Furthermore,
$$\mathrm{\Delta }(x;𝐭;q)=𝒞(x;𝐭;q)\mathrm{\Delta }_+(x;𝐭;q)$$
(8.3)
with $`𝒞(x)=𝒞(x;𝐭;q)`$ given by
$$𝒞(x;𝐭;q)=\underset{\alpha \mathrm{\Sigma }^{}}{}v_\alpha (x;𝐭;q).$$
(8.4)
We define now bilinear forms $`.,.=.,._{𝐭,q}`$ and $`.,._+=.,._{+,𝐭,q}`$ on $`𝒜`$ by
$$\begin{array}{cc}\hfill f,g& =\frac{1}{(2\pi i)^n}_{x𝕋^n}f(x)(\sigma g)(x)\mathrm{\Delta }(x)\frac{dx}{x},\hfill \\ \hfill f,g_+& =\frac{1}{(2\pi i)^n}_{x𝕋^n}f(x)(\sigma g)(x)\mathrm{\Delta }_+(x)\frac{dx}{x}\hfill \end{array}$$
(8.5)
where $`\frac{dx}{x}=\frac{dx_1}{x_1}\mathrm{}\frac{dx_n}{x_n}`$ and $`𝕋`$ is the (positively oriented) unit circle. Recall here that $`\sigma `$ is the longest Weyl group element in $`W`$. Observe that the bilinear forms $`.,.`$ and $`.,._+`$ are non-degenerate in both factors.
The bilinear form $`.,._+`$ coincides with Koornwinder’s pairing for the symmetric Koornwinder polynomials, see also for an extension to more general parameter values. If on the other hand the parameter values are such that $`\mathrm{\Delta }()𝒜`$, then $`f,g`$ equals the constant term of the Laurent polynomial $`f(x)(\sigma g)(x)\mathrm{\Delta }(x)`$ and $`.,.`$ then coincides with Sahi’s pairing for the non-symmetric Koornwinder polynomials.
###### Lemma 8.1.
We have
$$\underset{wW}{}w𝒞(;𝐭;q)=K_{𝐭,q}$$
(8.6)
in $`(x)`$ for some constant $`K=K_{𝐭,q}`$. In particular,
$$f,g=\frac{K}{|W|}f,g_+,f,g𝒜^W$$
(8.7)
where $`|W|=2^nn!`$ is the cardinality of the finite Weyl group $`W`$.
###### Proof.
The first statement follows from \[17, (2.8 n.r)\] with the indeterminates in \[17, (2.8 n.r)\] specialized to $`u_\alpha ^{1/2}=t_\alpha t_{\alpha /2}`$ for $`\alpha \mathrm{\Sigma }_l^+`$, $`u_\alpha =t_\alpha ^2`$ for $`\alpha \mathrm{\Sigma }_m^+`$ and $`u_\alpha =t_\alpha ^2`$ for $`\alpha \frac{1}{2}\mathrm{\Sigma }_l^+`$. The identity (8.7) follows then from (8.3) and the invariance of the measure $`(𝕋^n,\frac{dx}{x})`$ under the action of $`W`$. ∎
A product form for the constant $`K`$ can be obtained by specializing the left hand side of (8.6) at $`x_0^1`$, see \[17, (2.4 n.r)\]. In fact, we have the following more general result.
###### Lemma 8.2.
Let $`\lambda \mathrm{\Lambda }^+`$, and write $`W_\lambda `$ (respectively $`W^\lambda `$) for the stabilizer sub-group of $`\lambda `$ in $`W`$ (respectively the minimal coset representatives of $`W/W_\lambda `$). Then $`K=_{wW^\lambda }𝒞(x_{w\lambda }^1)`$. In particular, $`K=𝒞(x_0^1)`$.
###### Proof.
By the definition (8.6) of $`K`$ we have
$$K=\underset{wW^\lambda ,uW_\lambda }{}(u^1w^1𝒞)(x_\lambda ^1).$$
(8.8)
We consider a term $`(u^1w^1𝒞)(x_\lambda ^1)`$ in this sum with $`u1`$. Then there exists a simple root $`a_i`$ ($`i\{1,\mathrm{},n\}`$) which is orthogonal to $`\lambda `$, and which is mapped to a negative root $`\alpha `$ by $`wu`$. Now remark 4.7 implies that the factor $`v_{u^1w^1\alpha }(x_\lambda ^1)=v_{a_i}(x_\lambda ^1)`$ of $`(u^1w^1𝒞)(x_\lambda ^1)`$ is zero. Hence the contribution in the sum (8.8) is zero unless $`u=1`$. The lemma follows now from lemma 4.6. ∎
###### Proposition 8.3.
For $`X`$ we have
$$X(f),g=f,X^{}(g),f,g𝒜,$$
where $`:^{}`$ is the anti-algebra isomorphism defined in lemma 7.3.
###### Proof.
The proposition is obviously correct for $`X=z^\lambda `$ ($`\lambda \mathrm{\Lambda }`$), so it suffices to prove it for $`X=T_i`$ ($`i=0,\mathrm{},n`$). Let $`f,g𝒜`$. It follows by direct computations that
$$\left(T_if\right)(x)\left(\sigma g\right)(x)f(x)\left(\sigma \left((T_i^{})^1g\right)\right)(x)=t_i^1h_i(x)v_{a_i}(x)$$
(8.9)
for $`i=0,\mathrm{},n`$, with
$$h_i(x)=\left(s_if\right)(x)\left(\sigma g\right)(x)f(x)\left(s_i(\sigma g)\right)(x)$$
and with the action of $`s_i`$ as defined in (3.1). Now observe that $`h_i`$ is $`s_i`$-alternating, i.e. $`s_ih_i=h_i`$ for $`i=0,\mathrm{},n`$. On the other hand,
$$v_{a_i}(x)\mathrm{\Delta }(x)=\underset{\beta R^+\{a_i\}}{}\frac{1}{v_\beta (x)}$$
(8.10)
is invariant under the action of $`s_i`$ for $`i=0,\mathrm{},n`$, where the action of $`s_i`$ is extended from $`𝒜`$ to (suitably nice) functions $`f`$ in the $`n`$ variables $`x=(x_1,\mathrm{},x_n)`$ via the formulas (3.1). This is an immediate consequence of the well-known fact that the roots $`R^+\{a_i\}`$ are permuted by the simple reflection $`s_i`$. Hence $`T_if,gf,(T_i^{})^1g`$ can be rewritten as an integral over $`(𝕋^n,\frac{dx}{x})`$ with $`s_i`$-alternating integrand for all $`i\{0,\mathrm{},n\}`$.
Now $`T_if,gf,(T_i^{})^1g=0`$ for $`i=1,\mathrm{},n`$ follows from the fact that the measure $`(𝕋^n,\frac{dx}{x})`$ is $`W`$-invariant. The case $`i=0`$ is more subtle. The behaviour of the measure $`(𝕋^n,\frac{dx}{x})`$ under the action of $`s_0`$ is given by
$$_{x𝕋^n}(s_0h)(x)\frac{dx}{x}=_{y_1q𝕋}_{y𝕋^{n1}}h(y_1,y)\frac{dy_1}{y_1}\frac{dy}{y},$$
which now implies that
$$\begin{array}{cc}\hfill T_0f,g& f,(T_0^{})^1g=\hfill \\ \hfill =& \frac{1}{2(2\pi i)^n}_{y_1𝕋q𝕋}_{y𝕋^{n1}}t_0^1h_0(y_1,y)v_{a_0}(y_1,y)\mathrm{\Delta }(y_1,y)\frac{dy_1}{y_1}\frac{dy}{y}.\hfill \end{array}$$
(8.11)
For fixed $`y𝕋^{n1}`$, the integrand in the right-hand side of (8.11) depends analytically on $`y_1\{v|q|v|1\}`$. Indeed, by a direct computation using (5.2) and the second expression of $`\mathrm{\Delta }_+(x)`$ in (8.2), we see that the $`y_1`$-dependent factor of $`v_{a_0}(y_1,y)\mathrm{\Delta }(y_1,y)`$ is given by
$$\begin{array}{cc}& \frac{(y_1^2,q^2y_1^2;q)_{\mathrm{}}}{(ay_1,by_1,cy_1,dy_1,qay_1^1,qby_1^1,qcy_1^1,qdy_1^1;q)_{\mathrm{}}}\hfill \\ & .\underset{j=2}{\overset{n}{}}\frac{(y_1y_j,y_1y_j^1,qy_1^1y_j,qy_1^1y_j^1;q)_{\mathrm{}}}{(t^2y_1y_j,t^2y_1y_j^1,qt^2y_1^1y_j,qt^2y_1^1y_j^1;q)_{\mathrm{}}},\hfill \end{array}$$
which has the desired analytic behaviour due to the conditions on the parameters $`q`$ and $`𝐭`$. Thus by Cauchy’s theorem we conclude that $`T_0f,gf,(T_0^{})^1g=0`$. This completes the proof of the proposition. ∎
###### Remark 8.4.
An algebraic proof of proposition 8.3 was given by Sahi \[25, thm. 16\] for those (discrete) values of $`𝐭`$ such that $`\mathrm{\Delta }(;𝐭;q)𝒜`$.
We write $`E^{}(\gamma _\lambda ^1;)`$ for the renormalized Koornwinder polynomial of degree $`\lambda \mathrm{\Lambda }`$ with respect to inverse parameters $`(𝐭^1,q^1)`$. Since $`(Y^\lambda )^{}=(Y^\lambda )^1`$ for $`\lambda \mathrm{\Lambda }`$, we obtain the following extension of \[25, cor. 17\] from theorem 4.8 and proposition 8.3.
###### Corollary 8.5 (Bi-orthogonality relations).
For $`\lambda ,\mu \mathrm{\Lambda }`$ with $`\lambda \mu `$ we have $`E(\gamma _\lambda ;),E^{}(\gamma _\mu ^1;)=0`$.
Recall that $`E^+(\gamma _\lambda ;)`$ ($`\lambda \mathrm{\Lambda }^+`$) can be characterized as the unique solution of the eigenvalue equation
$$Lf=\left(m_{ϵ_1}(\gamma _\lambda )m_{ϵ_1}(\gamma _0)\right)f,f𝒜^W$$
which takes the value one at $`x_0^{\pm 1}`$, where $`L`$ (5.1) is Koornwinder’s second order $`q`$-difference operator. Since $`L`$ and the eigenvalue $`m_{ϵ_1}(\gamma _\lambda )m_{ϵ_1}(\gamma _0)`$ are invariant under replacement of the parameters $`(𝐭,q)`$ by their inverses $`(𝐭^1,q^1)`$, we derive that
$$E^+(\gamma _\lambda ;x;𝐭;q)=E^+(\gamma _\lambda ^1;x;𝐭^1;q^1),\lambda \mathrm{\Lambda }^+.$$
(8.12)
Combined with (8.7), corollary 8.5 and theorem 6.6, we re-obtain Koornwinder’s orthogonality relations for the symmetric Koornwinder polynomials:
###### Corollary 8.6 (Orthogonality relations).
For $`\lambda ,\mu \mathrm{\Lambda }^+`$ with $`\lambda \mu `$, we have $`E^+(\gamma _\lambda ;),E^+(\gamma _\mu ;)_+=0`$.
We write $`\text{Spec}(Y^{})=\{\gamma _\lambda ^1|\lambda \mathrm{\Lambda }\}`$ for the spectrum of the $`Y^{}`$-operators, and $`F=F_{𝐭,q}`$ for the linear space of functions $`g:\text{Spec}(Y^{})`$ with finite support. We define a $`𝒲`$-module structure on $`F`$ by
$$(wg)(\gamma _\lambda ^1)=g(\gamma _{w^1.\lambda }^1),gF,w𝒲,\lambda \mathrm{\Lambda }.$$
###### Definition 8.7.
We call the linear map $`=_{𝐭,q}:𝒜F`$ defined by
$$(f)(\gamma )=f,E^{}(\gamma ;),f𝒜,\gamma \text{Spec}(Y^{})$$
(8.13)
the non-symmetric Koornwinder transform.
Observe that $``$ is injective since $`.,.`$ is non-degenerate, and that $``$ is surjective by corollary 8.5.
In the next proposition we present an action of the double affine Hecke algebra $`\stackrel{~}{}`$ on $`F`$ in terms of spectral difference-reflection operators, and we relate it to the action of $``$ on $`𝒜`$ via the non-symmetric Koornwinder transform $``$ and the duality isomorphism $`\mathrm{\Phi }`$.
###### Proposition 8.8.
The applications
$$\begin{array}{cc}\hfill (\stackrel{~}{T}_ig)(\gamma )& =\stackrel{~}{t}_ig(\gamma )+\stackrel{~}{t}_i^1\stackrel{~}{v}_{a_i}(\gamma )((s_ig)(\gamma )g(\gamma )),i\{0,\mathrm{},n\},\hfill \\ \hfill \left(f(\stackrel{~}{z})g\right)(\gamma )& =f(\gamma )g(\gamma ),f𝒜\hfill \end{array}$$
(8.14)
where $`gF`$ and $`\gamma \text{Spec}(Y^{})`$, uniquely extend to an action of $`\stackrel{~}{}`$ on $`F`$. Furthermore,
$$(X(f))=\mathrm{\Phi }(X)(f),X,f𝒜,$$
(8.15)
where $`\mathrm{\Phi }`$ is the duality isomorphism.
###### Proof.
The intertwining property (8.15) can be proved by checking it for the algebraic generators $`U_n`$, $`T_i`$ and $`Y_i`$ ($`i=1,\mathrm{},n`$) of $``$ using proposition 8.3 and proposition 7.8. The fact that (8.14) defines an action of $`\stackrel{~}{}`$ on $`F`$ follows then from (8.15) since $`\mathrm{\Phi }`$ is an algebra isomorphism and $``$ is bijective. ∎
Next we determine the inverse of the non-symmetric Koornwinder transform $``$. We let $`𝒢=𝒢_{𝐭,q}:F𝒜`$ be the linear map defined by
$$(𝒢g)(x)=\underset{\lambda \mathrm{\Lambda }}{}g(\gamma _\lambda ^1)E(\gamma _\lambda ;x;𝐭;q)w(\gamma _\lambda ^1;\stackrel{~}{𝐭};q),gF,$$
(8.16)
where the discrete weights $`\stackrel{~}{w}(\gamma _\lambda ^1)=w(\gamma _\lambda ^1;\stackrel{~}{𝐭};q)`$ ($`\lambda \mathrm{\Lambda }`$) are defined as follows:
$$w(\gamma _\lambda ^1;\stackrel{~}{𝐭};q)=𝒞(\gamma _\lambda ^1;\stackrel{~}{𝐭};q)w_+(\gamma _{\lambda ^+}^1;\stackrel{~}{𝐭};q),$$
(8.17)
with $`\stackrel{~}{w}_+(\gamma _\mu ^1)=w_+(\gamma _\mu ^1;\stackrel{~}{𝐭};q)`$ for $`\mu \mathrm{\Lambda }^+`$ given by the multiple residue
$$w_+(\gamma _\mu ^1;\stackrel{~}{𝐭};q)=\underset{x_1=\gamma _\mu ^{ϵ_1}}{\text{Res}}\left(\underset{x_2=\gamma _\mu ^{ϵ_2}}{\text{Res}}\left(\mathrm{}\underset{x_n=\gamma _\mu ^{ϵ_n}}{\text{Res}}\left(\frac{\mathrm{\Delta }_+(x;\stackrel{~}{𝐭};q)}{x_1\mathrm{}x_n}\right)\mathrm{}\right)\right).$$
(8.18)
Using the second expression in (8.2) together with (4.4), it is easily verified that the discrete weights $`\stackrel{~}{w}(\gamma )`$ and $`\stackrel{~}{w}_+(\gamma )`$ ($`\gamma \text{Spec}(Y^{})`$) are well defined and non-zero for generic parameters $`𝐭`$ and $`q`$ (in fact, all residues in (8.18) are taken at simple poles).
###### Proposition 8.9.
We have
$$𝒢(\stackrel{~}{X}g)=\mathrm{\Phi }^1(\stackrel{~}{X})𝒢(g),\stackrel{~}{X}\stackrel{~}{},gF.$$
###### Proof.
The proof for $`\stackrel{~}{X}=\stackrel{~}{z}^\lambda `$ with $`\lambda \mathrm{\Lambda }`$ is immediate. Hence it suffices to check the intertwining property for $`\stackrel{~}{X}=\stackrel{~}{T}_i`$ ($`i=0,\mathrm{},n`$). Let $`gF`$. By proposition 7.8 we have for $`i=0,\mathrm{},n`$,
$$𝒢\left(\stackrel{~}{T}_ig\right)\mathrm{\Phi }^1(\stackrel{~}{T}_i)\left(𝒢g\right)=\stackrel{~}{t}_i^1\underset{\lambda \mathrm{\Lambda }}{}h_i(\gamma _\lambda ;)\stackrel{~}{v}_{a_i}(\gamma _\lambda ^1)\stackrel{~}{w}(\gamma _\lambda ^1)$$
with $`h_i(\gamma _\lambda ;)𝒜`$ given by
$$h_i(\gamma _\lambda ;)=g(\gamma _{s_i.\lambda }^1)E(\gamma _\lambda ;)g(\gamma _\lambda ^1)E(\gamma _{s_i.\lambda };).$$
Since $`h_i(\gamma _{s_i.\lambda };)=h_i(\gamma _\lambda ;)`$ for $`i=0,\mathrm{},n`$ and $`\lambda \mathrm{\Lambda }`$, it thus suffices to prove that
$$\stackrel{~}{v}_{a_i}(\gamma _\lambda ^1)\stackrel{~}{w}(\gamma _\lambda ^1)=\stackrel{~}{w}_+(\gamma _{\lambda ^+}^1)\underset{\alpha \mathrm{\Sigma }^{}\{a_i\}}{}\stackrel{~}{v}_\alpha (\gamma _\lambda ^1)$$
(8.19)
is invariant under replacement of $`\lambda \mathrm{\Lambda }`$ by $`s_i.\lambda `$ for all $`i\{0,\mathrm{},n\}`$ and all $`\lambda \mathrm{\Lambda }`$. For $`i\{1,\mathrm{},n\}`$ this is immediate by lemma 4.6.
As usual, the proof for the affine part of the statement (the case $`i=0`$) is more subtle. We begin by rewriting $`\stackrel{~}{w}(\gamma _\lambda ^1)`$ as a (kind of) multiple residue of $`\stackrel{~}{\mathrm{\Delta }}(x)`$ at $`x=\gamma _\lambda ^1`$, where $`\stackrel{~}{\mathrm{\Delta }}(x)=\mathrm{\Delta }(x;\stackrel{~}{𝐭};q)`$. This can be done using the $`W`$-invariance of the weight function $`\mathrm{\Delta }_+(;\stackrel{~}{𝐭};q)`$. The result is as follows.
Let $`u_\lambda S_n`$ be the component in $`S_n`$ of the minimal coset representative $`w_\lambda W`$ with respect to the semi-direct product structure $`W=S_n(\pm 1)^n`$, and let $`n_\lambda =\mathrm{\#}\{i\{1,\mathrm{},n\}|\lambda _i<0\}`$. Then we have
$$\stackrel{~}{w}(\gamma _\lambda ^1)=\underset{x=\gamma _\lambda ^1}{\text{Res}}\left(\frac{\stackrel{~}{\mathrm{\Delta }}(x)}{x_1\mathrm{}x_n}\right)$$
(8.20)
for all $`\lambda \mathrm{\Lambda }`$, where the multiple residue at $`x=\gamma _\lambda ^1`$ is defined by
$$\underset{x=\gamma _\lambda ^1}{\text{Res}}()=(1)^{n_\lambda }\underset{x_{u_\lambda (1)}=\gamma _\lambda ^{ϵ_{u_\lambda (1)}}}{\text{Res}}\left(\underset{x_{u_\lambda (2)}=\gamma _\lambda ^{ϵ_{u_\lambda (2)}}}{\text{Res}}\left(\mathrm{}\underset{x_{u_\lambda (n)}=\gamma _\lambda ^{ϵ_{u_\lambda (n)}}}{\text{Res}}()\mathrm{}\right)\right).$$
In particular, we obtain
$$\stackrel{~}{v}_{a_0}(\gamma _\lambda ^1)\stackrel{~}{w}(\gamma _\lambda ^1)=\underset{x=\gamma _\lambda ^1}{\text{Res}}\left(\frac{\stackrel{~}{v}_{a_0}(x)\stackrel{~}{\mathrm{\Delta }}(x)}{x_1\mathrm{}x_n}\right)$$
(8.21)
for all $`\lambda \mathrm{\Lambda }`$. Now we consider (8.21) with $`\lambda `$ replaced by $`s_0.\lambda `$. We first consider the changes in the multiple residue. By the proof of \[24, thm. 5.3\], we have $`w_{s_0.\lambda }=s_{ϵ_1}w_\lambda `$ for all $`\lambda \mathrm{\Lambda }`$, i.e. $`n_{s_0.\lambda }=n_\lambda \pm 1`$ and $`u_{s_0.\lambda }=u_\lambda `$. Secondly,
$$(\gamma _{s_0.\lambda }^1)^{ϵ_i}=\{\begin{array}{cc}q\gamma _\lambda ^{ϵ_1}\hfill & \text{ if }i=1,\hfill \\ \gamma _\lambda ^{ϵ_i}\hfill & \text{ if }i=2,\mathrm{},n\hfill \end{array}$$
by lemma 4.6. So if we replace the residue at $`x_1=\gamma _\lambda ^{ϵ_1}`$ by the residue at $`x_1=q\gamma _\lambda ^{ϵ_1}`$ in the definition of the multiple residue at $`x=\gamma _\lambda ^1`$, then we obtain minus the multiple residue at $`x=\gamma _{s_0.\lambda }^1`$. On the other hand, we know by the proof of proposition 8.3 that $`v_{a_0}(x)\mathrm{\Delta }(x)`$ is invariant under the action of $`s_0`$. So the invariance of (8.21) under replacement of $`\lambda `$ by $`s_0.\lambda `$ follows from the simple observation that
$$\underset{y=y_0}{\text{Res}}\left(\frac{h(y)}{y}\right)=\underset{y=qy_0^1}{\text{Res}}\left(\frac{h(y)}{y}\right)$$
when $`h(y)`$ is a function depending on a single variable $`y`$, having a simple pole at $`y=y_0`$, and satisfying the invariance condition $`h(qy^1)=h(y)`$. ∎
###### Theorem 8.10.
We have $`𝒢=k\text{Id}_𝒜`$ and $`𝒢=k\text{Id}_F`$ with $`k=k_{𝐭,q}=w(\gamma _0^1;\stackrel{~}{𝐭};q)1,1_{𝐭,q}`$. In particular, we have for $`\lambda ,\mu \mathrm{\Lambda }`$,
$$\frac{E(\gamma _\lambda ;),E^{}(\gamma _\mu ^1;)}{1,1}=\delta _{\lambda ,\mu }\frac{\stackrel{~}{w}(\gamma _0^1)}{\stackrel{~}{w}(\gamma _\lambda ^1)}$$
(8.22)
where $`\delta _{\lambda ,\mu }`$ is the Kronecker delta.
###### Proof.
By proposition 8.8 and proposition 8.9 we have
$$𝒢((f))=𝒢((f(z)1))=f(z)𝒢((1)),f𝒜.$$
(8.23)
Furthermore, it follows from corollary 8.5 that
$$𝒢\left((E(\gamma ;))\right)=E(\gamma ;),E^{}(\gamma ^1;)\stackrel{~}{w}(\gamma ^1)E(\gamma ;)$$
(8.24)
for $`\gamma \text{Spec}(Y)`$. Formula (8.24) reduces to $`𝒢((1))=k\mathrm{\hspace{0.17em}1}`$ when $`\gamma =\gamma _0`$, with the constant $`k`$ as given in the statement of the theorem. Combined with (8.23) it follows that $`𝒢=k\text{Id}_𝒜`$. Since $``$ is bijective, we then also have $`𝒢=k\text{Id}_F`$.
It remains to prove (8.22). By corollary 8.5, we only have to prove (8.22) when $`\mu =\lambda `$. We fix $`\gamma =\gamma _\lambda \text{Spec}(Y)`$, $`\lambda \mathrm{\Lambda }`$. Since $`𝒢=k\text{Id}_𝒜`$, it follows that $`𝒢((E(\gamma ;)))=kE(\gamma ;)`$. Comparing this outcome with the right-hand side of (8.24), we obtain
$$E(\gamma ;),E^{}(\gamma ^1;)\stackrel{~}{w}(\gamma ^1)=k=1,1\stackrel{~}{w}(\gamma _0^1)$$
which yields the desired result.
###### Corollary 8.11.
For all $`\lambda ,\mu \mathrm{\Lambda }^+`$, we have
$$\frac{E^+(\gamma _\lambda ;),E^+(\gamma _\mu ;)_+}{1,1_+}=\delta _{\lambda ,\mu }\frac{\stackrel{~}{w}_+(\gamma _0^1)}{\stackrel{~}{w}_+(\gamma _\lambda ^1)}.$$
###### Proof.
First of all, observe that $`E^+(\gamma _{\lambda ^+};)=C_+E(\gamma _\lambda ;)`$ for all $`\lambda \mathrm{\Lambda }`$ and that $`(C_+)^{}=C_+^{}`$, where $`C_+`$ (respectively $`C_+^{}^{}`$, $`\stackrel{~}{C}_+\stackrel{~}{}`$) is the idempotent corresponding to the trivial representation of the underlying finite Hecke algebra of type $`C_n`$. Combined with proposition 8.3, (8.7) and (8.12) we can rewrite $`_+=|_{𝒜^W}`$ as
$$\left(_+f\right)(\gamma _\lambda ^1)=\frac{K}{|W|}f,E^+(\gamma _{\lambda ^+};)_+,f𝒜^W,\lambda \mathrm{\Lambda }.$$
In particular, the symmetric Koornwinder transform $`_+`$ maps into
$$F^W=\{gF|wg=gw𝒲\}=\{gF|\stackrel{~}{C}_+g=g\}.$$
Similarly, since $`\mathrm{\Phi }(C_+)=\stackrel{~}{C}_+`$, we derive from proposition 8.9, (8.17) and lemma 8.2 that $`𝒢_+=𝒢|_{F^W}`$ can be rewritten as
$$\left(𝒢_+g\right)(x)=\stackrel{~}{K}\underset{\lambda \mathrm{\Lambda }^+}{}g(\gamma _\lambda ^1)E^+(\gamma _\lambda ;x)\stackrel{~}{w}_+(\gamma _\lambda ^1),$$
where $`\stackrel{~}{K}=K_{\stackrel{~}{𝐭},q}`$. Using these alternative descriptions for $`_+`$ and $`𝒢_+`$ together with the orthogonality relations for the symmetric Koornwinder polynomials (see corollary 8.6), we obtain
$$𝒢_+(_+(E^+(\gamma _\lambda ;)))=\frac{K\stackrel{~}{K}}{|W|}E^+(\gamma _\lambda ;),E^+(\gamma _\lambda ;)_+\stackrel{~}{w}_+(\gamma _\lambda ^1)E^+(\gamma _\lambda ;)$$
for all $`\lambda \mathrm{\Lambda }^+`$. On the other hand, by theorem 8.10, we have $`𝒢_+(_+(E^+(\gamma _\lambda ;)))=kE^+(\gamma _\lambda ;)`$ for $`\lambda \mathrm{\Lambda }^+`$ with
$$k=1,1\stackrel{~}{w}(\gamma _0^1)=\frac{K\stackrel{~}{K}}{|W|}1,1_+\stackrel{~}{w}_+(\gamma _0^1).$$
Comparing the two different outcomes for $`𝒢_+(_+(E^+(\gamma _\lambda ;)))`$, we obtain the desired result in case $`\mu =\lambda `$. The off-diagonal case is covered by the orthogonality relations for the symmetric Koornwinder polynomials, see corollary 8.6. ∎
###### Remark 8.12.
The discrete weights $`\stackrel{~}{w}_+(\gamma _\lambda ^1)=w_+(\gamma _\lambda ^1;\stackrel{~}{𝐭};q)`$ ($`\lambda \mathrm{\Lambda }^+`$) also appear as weights in a partly discrete orthogonality measure for the symmetric Koornwinder polynomials, see . In particular, using the expression \[26, prop. 4.1\] for the discrete weights, it can be shown that the ratio $`E^+(\gamma _\lambda ;),E^+(\gamma _\lambda ;)_+/1,1_+=\stackrel{~}{w}_+(\gamma _0^1)/\stackrel{~}{w}_+(\gamma _\lambda ^1)`$ is equal to
$$\begin{array}{cc}& \underset{i=1}{\overset{n}{}}\{\frac{(q^1abcdt^{4(ni)};q)_{2\lambda _i}\left(c^2t^{4(ni)}\right)^{\lambda _i}}{(abcdt^{4(ni)};q)_{2\lambda _i}}\hfill \\ & .\frac{(qt^{2(ni)},abt^{2(ni)},adt^{2(ni)},bdt^{2(ni)};q)_{\lambda _i}}{(act^{2(ni)},bct^{2(ni)},cdt^{2(ni)},q^1abcdt^{2(ni)};q)_{\lambda _i}}\}\hfill \\ & .\underset{1i<jn}{}\frac{(abcdt^{2(2nij1)},q^1abcdt^{2(2nij)};q)_{\lambda _i+\lambda _j}}{(abcdt^{2(2nij)},q^1abcdt^{2(2nij+1)};q)_{\lambda _i+\lambda _j}}\frac{(qt^{2(ji1)},t^{2(ji)};q)_{\lambda _i\lambda _j}}{(qt^{2(ji)},t^{2(ji+1)};q)_{\lambda _i\lambda _j}}\hfill \end{array}$$
for all $`\lambda \mathrm{\Lambda }^+`$, where $`(y_1,\mathrm{},y_m;q)_k=_{j=1}^m(y_j;q)_k`$ and $`(y;q)_k=_{j=0}^{k1}(1yq^j)`$ for $`k_+`$.
## 9. Evaluation formulas
In this section we give an explicit expression for the value of the non-symmetric Koornwinder polynomial $`P_\lambda ()=P_\lambda (;𝐭;q)`$ at $`x_0^1`$. We start with two preliminary lemmas.
###### Lemma 9.1.
For $`\lambda \mathrm{\Lambda }`$, we have
$$\frac{P_\lambda ,E^{}(\gamma _\lambda ^1;)}{P_{\lambda ^+},E^{}(\gamma _{\lambda ^+}^1;)}=\stackrel{~}{t}_{w_\lambda }^2\underset{\alpha \mathrm{\Sigma }^+w_\lambda ^1\mathrm{\Sigma }^{}}{}\frac{1}{\stackrel{~}{v}_\alpha (\gamma _{\lambda ^+}^1)}.$$
###### Proof.
We prove the lemma using the non-affine intertwiners $`S_i`$ ($`i=1,\mathrm{},n`$). For the moment, we fix $`\lambda \mathrm{\Lambda }`$ and $`i\{1,\mathrm{},n\}`$ such that $`\mu =s_i\lambda \lambda `$. We first give some additional properties of $`S_i`$ which we will need for the proof.
The intertwiner $`S_i`$ is self-adjoint, i.e. $`(S_i)^{}=S_i^{}`$, where $`S_i^{}=[T_i^{},(Y^{})^{a_i}]^{}`$. This can be checked most easily in the image of the duality anti-isomorphism $`\mathrm{\Psi }_{𝐭^1,q^1}`$. It follows from proposition 7.8(ii) that $`S_i^{}E^{}(\gamma _\lambda ^1;)=L_{i,\lambda }E^{}(\gamma _\mu ^1;)`$ with $`L_{i,\lambda }=\left(\gamma _\lambda ^{a_i}\gamma _\lambda ^{a_i}\right)\stackrel{~}{t}_i^1\stackrel{~}{v}_{a_i}(\gamma _\lambda ^1)`$. On the other hand, $`S_iP_\lambda =K_{i,\lambda }P_\mu `$ with $`K_{i,\lambda }=\left(\gamma _\lambda ^{a_i}\gamma _\lambda ^{a_i}\right)\eta _i(\gamma _\lambda )`$ by corollary 6.2. Combined with proposition 8.3 we obtain
$$P_\mu ,E^{}(\gamma _\mu ^1;)=\frac{1}{L_{i,\lambda }}S_iP_\mu ,E^{}(\gamma _\lambda ^1;)=\frac{K_{i,\mu }}{L_{i,\lambda }}P_\lambda ,E^{}(\gamma _\lambda ^1;).$$
In particular, if $`\lambda ,a_i<0`$, then
$$P_\lambda ,E^{}(\gamma _\lambda ^1;)=\frac{\stackrel{~}{t}_i^2}{\stackrel{~}{v}_{a_i}(\gamma _\mu ^1)}P_\mu ,E^{}(\gamma _\mu ^1;).$$
The ratio $`P_\lambda ,E^{}(\gamma _\lambda ^1;)/P_{\lambda ^+},E^{}(\gamma _{\lambda ^+}^1;)`$ can now be evaluated inductively using similar techniques as in the proof of theorem 6.6. This gives the desired result. ∎
Let $`\delta _{\gamma _\lambda ^1}F`$ for $`\lambda \mathrm{\Lambda }`$ be the function which is equal to one at $`\gamma _\lambda ^1`$ and zero otherwise.
###### Lemma 9.2.
For $`\lambda \mathrm{\Lambda }^+`$ we have
$$\left(\mathrm{\Phi }(z^\lambda )\delta _{\gamma _0^1}\right)(\gamma _\lambda ^1)=\stackrel{~}{t}_{\tau (\lambda )}^1\underset{\beta R^+\tau (\lambda )R^{}}{}\stackrel{~}{v}_\beta (\gamma _\lambda ^1).$$
###### Proof.
Let $`\lambda \mathrm{\Lambda }^+`$. It follows from (2.2) that $`\tau (\lambda )`$ is the unique element of minimal length in $`𝒲`$ which maps $`\lambda `$ to $`0\mathrm{\Lambda }`$ under the dot-action. In particular, any element $`w𝒲`$ which is smaller than $`\tau (\lambda )`$ with respect to the Bruhat order, maps $`\lambda `$ to a non-zero element in $`\mathrm{\Lambda }`$.
Since $`\lambda \mathrm{\Lambda }^+`$, we have $`Y^\lambda =T_{\tau (\lambda )}`$, hence
$$\mathrm{\Phi }(z^\lambda )=(ϵ(z^\lambda ))^\stackrel{~}{}^{}=(\stackrel{~}{T}_{\tau (\lambda )}^{})^\stackrel{~}{}^{}=\stackrel{~}{T}_{\tau (\lambda )}^1.$$
Let now $`\tau (\lambda )=s_{i_1}s_{i_2}\mathrm{}s_{i_r}`$ be a reduced expression of $`\tau (\lambda )`$ in $`𝒲`$, then we obtain from proposition 8.8 and from the previous paragraph that
$$\left(\mathrm{\Phi }(z^\lambda )\delta _{\gamma _0^1}\right)(\gamma _\lambda ^1)=\underset{m=1}{\overset{r}{}}\stackrel{~}{t}_{a_{i_m}}^1\stackrel{~}{v}_{a_{i_m}}\left(\gamma _{(s_{i_{m+1}}\mathrm{}s_{i_{r1}}s_{i_r}).\lambda }^1\right),$$
where $`(s_{i_{m+1}}\mathrm{}s_{i_{r1}}s_{i_r}).\lambda `$ should be read as $`\lambda `$ when $`m=r`$. The lemma follows now easily from lemma 4.6 and remark 6.5. ∎
###### Theorem 9.3.
Let $`\lambda \mathrm{\Lambda }`$, then
$$P_\lambda (x_0^1)=\frac{\stackrel{~}{w}(\gamma _\lambda ^1)}{\stackrel{~}{w}(\gamma _0^1)}\stackrel{~}{t}_{w_\lambda }^2\stackrel{~}{t}_{\tau (\lambda ^+)}^1\underset{\beta }{}\stackrel{~}{v}_\beta (\gamma _{\lambda ^+}^1),$$
where the product is taken over $`\beta (R^+\tau (\lambda ^+)R^{})(\mathrm{\Sigma }^+w_\lambda ^1\mathrm{\Sigma }^{})`$.
###### Proof.
First of all, observe that
$$\begin{array}{cc}\hfill P_\lambda (x_0^1)E(\gamma _\lambda ;),E^{}(\gamma _\lambda ^1;)& =P_\lambda ,E^{}(\gamma _\lambda ^1;),\hfill \\ \hfill P_\lambda ,E^{}(\gamma _\lambda ^1;)& =x^\lambda ,E^{}(\gamma _\lambda ^1;)=1,1\left(\mathrm{\Phi }(z^\lambda )\delta _{\gamma _0^1}\right)(\gamma _\lambda ^1)\hfill \end{array}$$
(9.1)
for all $`\lambda \mathrm{\Lambda }`$. Indeed, the first formula is trivial, while the second formula is a direct consequence of corollary 8.5, the definition of the non-symmetric Koornwinder transform $``$ and its intertwining properties given in proposition 8.8.
We use now successively the first formula of (9.1), then lemma 9.1 and finally the second formula of (9.1) to arrive at
$$P_\lambda (x_0^1)=\frac{\stackrel{~}{t}_{w_\lambda }^21,1}{E(\gamma _\lambda ;),E^{}(\gamma _\lambda ^1;)}\left(\mathrm{\Phi }(z^{\lambda ^+})\delta _{\gamma _0^1}\right)(\gamma _{\lambda ^+}^1)\underset{\alpha \mathrm{\Sigma }^+w_\lambda ^1\mathrm{\Sigma }^{}}{}\frac{1}{\stackrel{~}{v}_\alpha (\gamma _{\lambda ^+}^1)}.$$
The theorem follows now from theorem 8.10 and lemma 9.2, combined with the inclusion $`\mathrm{\Sigma }^+w_\lambda ^1\mathrm{\Sigma }^{}R^+\tau (\lambda ^+)R^{}`$. This inclusion is a direct consequence of the inequality $`\lambda ^+,\alpha >0`$ for all $`\alpha \mathrm{\Sigma }^+w_\lambda ^1\mathrm{\Sigma }^{}`$, see (6.15). ∎
###### Corollary 9.4.
$$P_\lambda ^+(x_0^{\pm 1})=\frac{\stackrel{~}{w}_+(\gamma _\lambda ^1)}{\stackrel{~}{w}_+(\gamma _0^1)}\stackrel{~}{t}_{\tau (\lambda )}^1\underset{\beta R^+\tau (\lambda )R^{}}{}\stackrel{~}{v}_\beta (\gamma _\lambda ^1),\lambda \mathrm{\Lambda }^+.$$
(9.2)
###### Proof.
This follows from the evaluation of $`P_\mu (x_0^1)`$ ($`\mu \mathrm{\Lambda }`$) (see theorem 9.3), theorem 6.6(iii), (8.17) and lemma 8.2. ∎
###### Remark 9.5.
It is straightforward to explicitly write down the roots $`R^+\tau (\lambda )R^{}`$ for $`\lambda \mathrm{\Lambda }^+`$. Together with (4.4), (3.7), remark 8.12 and the $`𝒲`$-orbit structure of $`R`$, one can reformulate (9.2) now as
$$\begin{array}{cc}\hfill P_\lambda ^+(x_0)=& \underset{i=1}{\overset{n}{}}\frac{(act^{2(ni)},bct^{2(ni)},cdt^{2(ni)},q^1abcdt^{2(ni)};q)_{\lambda _i}}{(q^1abcdt^{4(ni)};q)_{2\lambda _i}\left(ct^{2(ni)}\right)^{\lambda _i}}\hfill \\ & .\underset{1i<jn}{}\frac{(q^1abcdt^{2(2nij+1)};q)_{\lambda _i+\lambda _j}(t^{2(ji+1)};q)_{\lambda _i\lambda _j}}{(q^1abcdt^{2(2nij)};q)_{\lambda _i+\lambda _j}(t^{2(ji)};q)_{\lambda _i\lambda _j}}\hfill \end{array}$$
for $`\lambda \mathrm{\Lambda }^+`$.
###### Remark 9.6.
Van Diejen \[9, thm. 5.1\] proved the evaluation formula (9.2) for a five parameter sub-family of the symmetric Koornwinder polynomials. This result was indirectly extended to the complete six parameter family of symmetric Koornwinder polynomials by Sahi’s duality results.
## 10. The generalized Weyl character formula and the constant term
In this section we discuss several results involving the anti-symmetric Koornwinder polynomials $`P_\lambda ^{}()=P_\lambda ^{}(;𝐭;q)`$ ($`\lambda \mathrm{\Lambda }^{++}`$) in an essential way. The most crucial result is the analogue of the Weyl character formula for the Koornwinder polynomials.
In order to state the generalized Weyl character formula, we need to introduce some notations first. Let $`𝐪=\{q_\beta \}_{\beta S}`$ be the multiplicity function satisfying $`q_{a_0}=q_{a_0^{}}=q_{a_n^{}}=1`$, $`q_{a_i}=q^{1/2}`$ ($`i\{1,\mathrm{},n1\}`$) and $`q_{a_n}=q`$. For two multiplicity functions $`𝐭`$ and $`𝐭^{}`$, we write $`\mathrm{𝐭𝐭}^{}`$ for the multiplicity function which takes the value $`t_\beta t_\beta ^{}`$ at $`\beta S`$. Then the generalized Weyl character formula is given by
$$P_{\lambda +\kappa }^{}(x;𝐭;q)=\chi (x;𝐭;q)P_\lambda ^+(x;\mathrm{𝐪𝐭};q),\lambda \mathrm{\Lambda }^+$$
(10.1)
with $`\kappa `$ given by (6.9) and with $`\chi (;𝐭;q)𝒜`$ given by
$$\chi (x;𝐭;q)=x^\kappa \underset{\alpha \mathrm{\Sigma }^{}}{}(1x^\alpha )v_\alpha (x;𝐭^1;q^1).$$
The proof of (10.1) is similar to the proof of the generalized Weyl character formula for Macdonald polynomials, see e.g. \[18, 7.3\].
The generalized Weyl character formula (10.1), together with the results of section 8, readily implies the norm relations
$$\begin{array}{cc}\hfill \frac{P_\lambda ^+(;𝐭;q),P_\lambda ^+(;𝐭;q)_{𝐭,q}}{P_\lambda ^{}(;𝐭;q),P_\lambda ^{}(;𝐭^1;q^1)_{𝐭,q}}& =t_\sigma ^2\frac{P_\lambda ^+(;𝐭;q),P_\lambda ^+(;𝐭;q)_{+,𝐭,q}}{P_{\lambda \kappa }^+(;\mathrm{𝐪𝐭};q),P_{\lambda \kappa }^+(;\mathrm{𝐪𝐭};q)_{+,\mathrm{𝐪𝐭},q}}\hfill \\ & =\underset{\alpha \mathrm{\Sigma }^+}{}\frac{v_\alpha (\gamma _\lambda ^1;\stackrel{~}{𝐭};q)}{v_\alpha (\gamma _\lambda ;\stackrel{~}{𝐭};q)}\hfill \end{array}$$
(10.2)
for $`\lambda \mathrm{\Lambda }^{++}`$. The norm relations (10.2) give an explicit description of the diagonal terms $`P_\lambda ^{}(;𝐭;q),P_\lambda ^{}(;𝐭^1;q^1)_{𝐭,q}`$ ($`\lambda \mathrm{\Lambda }^{++}`$) corresponding to the bi-orthogonality relations
$$P_\lambda ^{}(;𝐭;q),P_\mu ^{}(;𝐭^1;q^1)_{𝐭,q}=0\lambda ,\mu \mathrm{\Lambda }^{++}:\lambda \mu $$
(10.3)
in terms of the quadratic norms of the symmetric Koornwinder polynomials. Furthermore, the norm relations (10.2) enables one to obtain a new proof for Gustafson’s evaluation of the constant term $`1,1_+`$. In fact, by (10.2), one can express $`1,1_{+,𝐪^m𝐯,q}`$ in terms of
$$P_{m\kappa }^+(;𝐯;q),P_{m\kappa }^+(;𝐯;q)_{+,𝐯,q}$$
(10.4)
for all positive integers $`m`$. Let now $`𝐯`$ be a multiplicity function with $`v_\beta =1`$ for all $`\beta S`$ of length two, then the corresponding orthogonality measure reduces to the (coordinate-wise) product measure of the one-variable Askey-Wilson polynomials. In particular, (10.4) can be expressed in terms of the quadratic norms of the Askey-Wilson polynomials, which were evaluated in (see for an affine Hecke algebraic approach). This yields an evaluation of (10.4), and hence of $`1,1_{+,𝐪^m𝐯,q}`$. By analytic continuation, we arrive at Gustafson’s result that the constant term $`|W|^11,1_{+,𝐭,q}`$ is equal to
$$\underset{j=1}{\overset{n}{}}\frac{(t^2,t^{2(2nj1)}abcd;q)_{\mathrm{}}}{(q,t^{2(nj+1)},t^{2(nj)}ab,t^{2(nj)}ac,t^{2(nj)}ad,t^{2(nj)}bc,t^{2(nj)}bd,t^{2(nj)}cd;q)_{\mathrm{}}}.$$
###### Remark 10.1.
In view of theorem 8.10, corollary 8.11, theorem 9.3, corollary 9.4 and Gustafson’s constant term evaluation, we have arrived now at the stage that the quadratic norms (respectively diagonal terms) of the (non-)symmetric Koornwinder polynomials are completely explicit. In particular, the explicit evaluation of $`P_\lambda ^+,P_\lambda ^+_+`$ ($`\lambda \mathrm{\Lambda }^+`$) which we thus obtain, can be seen to coincide with van Diejen’s \[9, thm. 5.2\] explicit expression for $`P_\lambda ^+,P_\lambda ^+_+`$.
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# References
In the past few years there has been much development in our understanding of the dynamics of supersymmetric gauge theories and superstring theories. Among these it has been discovered that the noncommutative gauge theories naturally appear when the D-branes with constant $`B`$ fields is considered.
Recently Seiberg and Witten have argued that the noncommutative gauge theories realized as effective theories on D-branes are equivalent to some ordinary gauge theories . In a single D-brane case, they have shown that the effective action for the D-brane is consistent with the equivalence if all derivative terms are neglected. Furthermore, it has been shown that the D-brane action, including derivative terms, computed in the string theory is consistent with the equivalence if we keep the two derivative terms but neglect the fourth and higher order derivative terms .
For deeper understanding for these phenomena, it is natural to investigate the field theories on the noncommutative geometry by the field theoretical approaches. In particular by the perturbative analysis it was found in that the IR effects and UV effects are mixed in the noncommutative field theory.
To proceed further it may be important to study the noncommutative field theories with supersymmetry since their actions are highly constrained and we may understand the dynamics of these theories. To obtain the supersymmetric action, superfields on the noncommutative geometry may be desired. Explicit two-dimensional $`N=1`$ noncommutative superspace was obtained in .
However, in this note, instead of investigating the noncommutative superspace formalism, we consider the ordinary superspace and superfields and represent the noncommutative field theory using these notions. From the commutative supersymmetric action written by the superfields, we can obtain the noncommutative supersymmetric action by replacing the ordinary product between superfields to the $``$ product defined by the formula
$$f(x)g(x)=e^{\frac{i}{2}\mathrm{\Theta }^{ij}\frac{}{\xi ^i}\frac{}{\zeta ^j}}f(x+\xi )g(x+\zeta )|_{\xi =\zeta =0}.$$
(1)
Here we regard the additional terms depend on $`\mathrm{\Theta }^{ij}`$ as the interaction terms with derivatives although we do not expand the $``$ product explicitly. We can do so because if $`\mathrm{\Phi }`$ is superfield then $`^n\mathrm{\Phi }`$ is also superfields. Thus it is obvious that this action has the supersymmetry and the R symmetry when it exists in the action with $`\mathrm{\Theta }=0`$ because $`_\mu \theta =0`$ where $`\theta `$ is the fermionic coordinate. This treatment of the superfields is similar to the notions of superspace and superfields in noncommutative geometry , in which only the chiral superfields has been considered.
In this paper, we consider the $`N=1`$ supersymmetric theories on the noncommutative $`𝐑^4`$ although the above observation does not depend on the dimension of the spacetime and the number of the supersymmetries. We construct the $`N=1`$ supersymmetric action for the $`U(N)`$ vector multiplets and chiral multiplets of the fundamental, anti-fundamental and adjoint representations of the gauge group. The actions for gauge fields of the products gauge groups and its bi-fundamental matters are also obtained.
It is argued that even if we do not require the supersymmetry, only these gauge groups and the matters are possible for the noncommutative gauge theories. We also find that the scalar potentials have some characteristic forms and discuss the problem of the derivative terms of the auxiliary fields.
The convention and notation taken in this paper are same as in except for the spacetime indices, which are denote by $`\mu ,\nu ,\mathrm{}`$ in this paper, and the gauge field $`A_\mu `$.
First we consider the chiral superfields which satisfy $`\overline{D}_{\dot{\alpha }}\mathrm{\Phi }=0`$. Using the coordinate $`y^m=x^m+i\theta \sigma ^m\overline{\theta }`$, these are written as $`\mathrm{\Phi }(y,\theta ,\overline{\theta })=A(y)+\sqrt{2}\theta \psi (y)+\theta \theta F(y)`$. The supersymmetry transformations are identical for commutative counterparts
$`\delta _\xi A`$ $`=`$ $`\sqrt{2}\xi \psi ,`$
$`\delta _\xi \psi `$ $`=`$ $`i\sqrt{2}\sigma ^m\overline{\xi }_mA+\sqrt{2}\xi F,`$
$`\delta _\xi F`$ $`=`$ $`i\sqrt{2}\overline{\xi }\overline{\sigma }^m_m\psi ,`$ (2)
because we simply consider the ordinary superfields.
Defining
$$(\underset{i=1}{\overset{n}{}}f_i)_{}=f_1f_2\mathrm{}f_n,$$
(3)
the most generic action which can be constructed from the chiral superfields $`\mathrm{\Phi }^i`$ takes the form
$$S=d^4x(d^2\theta d^2\overline{\theta }K(\mathrm{\Phi }^i,\mathrm{\Phi }^j)_{}+[d^2\theta W(\mathrm{\Phi }^i)_{}+h.c.]),$$
(4)
where $`d^2\theta \theta ^2=1`$ and $`d^2\overline{\theta }\overline{\theta }^2=1`$. This is invariant under $`K(\mathrm{\Phi }^i,\mathrm{\Phi }^j)_{}K(\mathrm{\Phi }^i,\mathrm{\Phi }^j)_{}+F(\mathrm{\Phi })_{}+F(\mathrm{\Phi }^{})_{}^{}`$. Although the action of the component fields can be obtained straightforwardly, we will only give some examples below.
First we consider the action with $`K=\mathrm{\Phi }^{}\mathrm{\Phi }+a\mathrm{\Phi }\mathrm{\Phi }(\mathrm{\Phi }^{})+a^{}\mathrm{\Phi }(\mathrm{\Phi }^{})(\mathrm{\Phi }^{})`$ and $`W=0`$, where $`a`$ is some numerical constant. Note that
$$(AB)^{}=B^{}A^{}.$$
(5)
The part of the action which depends on $`F`$ becomes
$`S|_F`$ $`=`$ $`{\displaystyle }d^4x(F^{}F+(aAFF^{}+aFAF^{}+aFFA^{}+h.c))`$ (6)
$`=`$ $`{\displaystyle }d^4x(F^{}F+(aF(F^{}A)+aF(AF^{})+aF(FA^{})+h.c)).`$
Here we have used that
$$d^4xAB=d^4AB=d^4xBA,$$
(7)
which implies that the integral of the product of fields with $``$ product are unchanged by the cyclic rotation of the fields. Thus the action clearly contains the derivative of the auxiliary field $`F`$ and it is difficult to eliminate it from the action using the equation of motion. Moreover $`F`$ may become the propagating field if the noncommutative parameter $`\mathrm{\Theta }^{0\mu }0`$ for some $`\mu `$. To avoid these problems, we only consider the canonical Kähler potential $`K=_i\mathrm{\Phi }_i^{}\mathrm{\Phi }_i`$ below.
With this $`K`$, the action with non vanishing superpotential does not have the derivative of $`F`$ then $`F`$ can be eliminated. This can be seen from the fact that the terms which depend on $`F`$ in the superpotential are linear in $`F`$. For example, the $`F`$ dependent parts of the action with $`W=a\mathrm{\Phi }^n`$ become
$$d^4x(\frac{1}{2}F^{}F+a\underset{i=1}{\overset{n}{}}F(A^{ni})_{}(A^{i1})_{}+h.c)).$$
(8)
Next we consider the noncommutative Wess-Zumino model
$$S_{WZ}=d^4x(d^2\theta d^2\overline{\theta }\mathrm{\Phi }_i^{}\mathrm{\Phi }_i+[d^2\theta (\frac{1}{2}m_{ij}\mathrm{\Phi }_i\mathrm{\Phi }_j+\frac{1}{3}g_{ijk}\mathrm{\Phi }_i\mathrm{\Phi }_j\mathrm{\Phi }_k+g_i\mathrm{\Phi }_i)+h.c.]),$$
(9)
where the mass $`m_{ij}`$ is symmetric in their indices, however, the coupling $`g_{ijk}`$ is not necessarily symmetric. By trancing the procedure for the $`\mathrm{\Theta }=0`$ case, we can easily find that
$`S_{WZ}`$ $`=`$ $`{\displaystyle d^4x\left(_\mu A_i^{}^\mu A_i+i_\mu \psi _i^{}\overline{\sigma }^\mu \psi _i+F_i^{}F_i\right)}`$ (10)
$`+{\displaystyle }d^4x[{\displaystyle \frac{1}{3}}g_{ijk}(F_iA_jA_k+F_jA_kA_i+F_kA_iA_jA_i\psi _j\psi _kA_j\psi _k\psi _iA_k\psi _i\psi _j)`$
$`+g_iF_i+m_{ij}(A_iF_j{\displaystyle \frac{1}{2}}\psi _i\psi _j)+h.c].`$
The equation of motions of $`F_i`$ is
$$F_i^{}=g_i+m_{ij}A_j+\frac{1}{3}\left(g_{ijk}+g_{kij}+g_{jki}\right)A_jA_k,$$
(11)
and the supersymmetry transformation becomes (2) with this $`F_i`$. We note that the typical scalar potential has the form $`A^{}A^{}AA`$ and the notion of holomorphy is still valid at $`\mathrm{\Theta }0`$.
Now we consider the vector superfields $`V=V^{}`$ ,
$`V(x,\theta ,\overline{\theta })`$ $`=`$ $`C(x)+i\theta \chi (x)i\overline{\theta }\overline{\chi }(x)`$ (12)
$`+{\displaystyle \frac{i}{2}}\theta \theta \left[M(x)+iN(x)\right]{\displaystyle \frac{i}{2}}\overline{\theta }\overline{\theta }\left[M(x)iN(x)\right]`$
$`\theta \sigma ^\mu \overline{\theta }A_\mu (x)+i\theta \theta \overline{\theta }\left[\overline{\lambda }(x)+{\displaystyle \frac{i}{2}}\overline{\sigma }^\mu _\mu \chi (x)\right]`$
$`i\overline{\theta }\overline{\theta }\theta \left[\lambda (x)+{\displaystyle \frac{i}{2}}\sigma ^\mu _\mu \overline{\chi }(x)\right]+{\displaystyle \frac{1}{2}}\theta \theta \overline{\theta }\overline{\theta }\left[D(x)+{\displaystyle \frac{1}{2}}\mathrm{}C(x)\right].`$
Since the $``$ product contain derivative, the ordinary gauge invariant action for the vector superfields can not be generalized to $`\mathrm{\Theta }0`$ case. Then we should introduce noncommutative gauge field $`A_\mu =T^aA_\mu ^a`$, where $`T^a`$ is the matrix for a representation of the gauge group $`G`$ and satisfies that $`(T^a)^{}=T^a`$ and $`\mathrm{Tr}(T^aT^b)=k`$.
Hereafter we briefly discuss the some properties of the noncommutative gauge field without requiring supersymmetry in order to prepare to treat the vector superfield. We assume that the noncommutative gauge transformation is the naive generalization of the ordinary Non-Abelian gauge transformation,
$$A_\mu UA_\mu U^1+iU_\mu U^1,$$
(13)
where $`U=(e^{i\lambda })_{}`$, $`U^1=(e^{i\lambda })_{}=U^{}`$ and $`\lambda =T^a\lambda ^a=\lambda ^{}`$. The infinitesimal version of this is
$`\delta _\lambda A_\mu `$ $`=`$ $`_\mu \lambda +i\lambda A_\mu iA_\mu \lambda `$ (14)
$`=`$ $`T^a(_\mu \lambda ^a)+{\displaystyle \frac{i}{2}}[T^a,T^b](\lambda ^aA_\mu ^b+A_\mu ^b\lambda ^a)+{\displaystyle \frac{i}{2}}\{T^a,T^b\}(\lambda ^aA_\mu ^bA_\mu ^b\lambda ^a).`$
From this if $`\{T^a,T^b\}`$ is not a linear combination of $`T^d`$ for some $`a,b`$, the gauge transformation is not closed. Thus the noncommutative gauge transformation is consistent only for unitary group $`G=U(N)`$ or its direct product $`G=_aU(N_a)^{(a)}`$. In addition to this restriction, we should take $`T^a`$ as the matrix for the fundamental or anti-fundamental representation by the requirement of the closure of (14). However $`T^a`$ and $`\stackrel{~}{T}^a={}_{}{}^{t}T`$, which represent the fundamental and anti-fundamental representations respectively, give the different gauge transformations via (14). Then we take $`(T^a)_j^i`$ to be the matrix for the fundamental representation, $`N\times N`$ Hermitian matrix. Of course we can choose the matrix of the anti-fundamental representation instead of the one for fundamental representation.
We can see that the representation of the gauge group $`G`$ of the matter is also restricted to fundamental ($`𝐍`$), anti-fundamental ($`\overline{𝐍}`$), adjoint ($`𝐍\times \overline{𝐍}`$) or bi-fundamental ($`𝐍\times \overline{𝐌}`$) from the consideration of the possible form of the covariant derivative. In this restriction has been shown for $`G=U(1)`$ case, where $`T^1=k^{\frac{1}{2}}`$. For the fundamental matter represented as the column vector $`(\psi )^i=\psi ^i`$, the gauge transformation and the covariant derivative are given by $`\psi U\psi `$ and
$$𝒟_\mu \psi =_\mu \psi iA_\mu \psi ,$$
(15)
respectively. We can easily check $`𝒟_\mu \psi U(𝒟_\mu \psi )`$ under the gauge transformation. Noting $`𝒟_\mu \psi ^{}(𝒟_\mu \psi )^{}=_\mu \psi +i\psi ^{}A_\mu `$, the covariant derivative for the anti-fundamental matter $`(\stackrel{~}{\psi })_i=\stackrel{~}{\psi }_i`$, which is transformed as $`\stackrel{~}{\psi }\stackrel{~}{\psi }U^1`$, is
$$𝒟_\mu \stackrel{~}{\psi }=_\mu \stackrel{~}{\psi }+i\stackrel{~}{\psi }A_\mu .$$
(16)
The adjoint matter $`(\psi _{adj})_i^j`$, which is transformed as $`\psi _{adj}U\psi _{adj}U^1`$, can have the covariant derivative
$$𝒟_\mu \psi _{adj}=_\mu \psi _{adj}iA_\mu \psi _{adj}+i\psi _{adj}A_\mu .$$
(17)
This is also seen from the covariant derivative for the bi-fundamental matter $`(\psi _{N\overline{M}})_i^j`$, where $`j=1\mathrm{}N`$ and $`i=1\mathrm{}M`$, is
$$𝒟_\mu \psi _{N\overline{M}}=_\mu \psi _{N\overline{M}}iA_\mu ^{(1)}\psi _{N\overline{M}}+i\psi _{N\overline{M}}A_\mu ^{(2)}.$$
(18)
Here $`A_\mu ^{(1)}`$ and $`A_\mu ^{(2)}`$ are the gauge fields for $`U(N)`$ and $`U(M)`$ respectively and the gauge transformation for it is given by $`\psi _{N\overline{M}}U^{(1)}\psi _{N\overline{M}}U_{}^{(2)}{}_{}{}^{1}`$.
The interaction terms are severely constrained by the gauge symmetry and the possible forms of the terms are the polynomials of $`\stackrel{~}{\psi }(\psi _{adj}^{}{}_{}{}^{n})_{}\psi `$, and $`\mathrm{Tr}(\psi _{adj}^{}{}_{}{}^{n})_{}`$ for $`G=U(N)`$. For the product group case, there are other terms which are allowed by the symmetry.
On the basis of this observation, we return to consider the vector superfields $`V=T^aV^a`$. We define the noncommutative super gauge transformation as
$$(e^{2V})_{}(e^{2V^{}})_{}=(e^{i\mathrm{\Lambda }^{}})_{}(e^{2V})_{}(e^{i\mathrm{\Lambda }})_{}.$$
(19)
The chiral superfield
$$W_\alpha =\frac{1}{8}\overline{D}\overline{D}\left((e^{2V})_{}D_\alpha (e^{2V})_{}\right),$$
(20)
is transformed as
$$W_\alpha W_{\alpha }^{}{}_{}{}^{}=(e^{i\mathrm{\Lambda }})_{}W_\alpha (e^{i\mathrm{\Lambda }})_{}.$$
(21)
Because of
$$V^{}=V+i(\mathrm{\Lambda }\mathrm{\Lambda }^{})+\mathrm{},$$
(22)
we can choose the Wess-Zumino gauge in which $`C,\chi ,M,N`$ are eliminated. In the Wess-Zumino gauge we see
$$W_\alpha (y)=i\lambda _\alpha (y)+\theta _\alpha D(y)\frac{i}{2}(\sigma ^\mu \overline{\sigma }^\nu \theta )_\alpha F_{\mu \nu }(y)+\theta ^2(\sigma ^\mu 𝒟_\mu \overline{\lambda }(y))_\alpha ,$$
(23)
where
$$F_{\mu \nu }=_\mu A_\nu _\nu A_\mu iA_\mu A_\nu +iA_\nu A_\mu ,$$
(24)
and
$$𝒟_\mu \overline{\lambda }=_\mu \overline{\lambda }iA_\mu \overline{\lambda }+i\overline{\lambda }A_\mu .$$
(25)
Defining the complex gauge coupling $`\tau =\frac{\stackrel{~}{\theta }}{2\pi }+\frac{4\pi i}{g^2}`$, we obtain the action of the noncommutative supersymmetric $`U(N)`$ gauge field theory, Although it is straightforward to obtain the action for the case of general Kähler potential, we only consider here the canonical Kähler potential in order to avoid the problem with the derivative terms of the auxiliary fields.
$`S_V`$ $`=`$ $`{\displaystyle \frac{1}{16\pi k}}{\displaystyle }d^4xd\theta ^2\mathrm{Tr}(i\tau W^\alpha W_\alpha +h.c.)`$ (26)
$`=`$ $`{\displaystyle d^4x\mathrm{Tr}\left(\frac{1}{4g^2}F^{\mu \nu }F_{\mu \nu }\frac{i}{g^2}\lambda \sigma ^\mu (𝒟_\mu \overline{\lambda })+\frac{1}{2g^2}DD\frac{\stackrel{~}{\theta }}{64\pi ^2}ϵ_{\mu \nu \alpha \beta }F^{\mu \nu }F^{\alpha \beta }\right)}.`$
Note that the definitions of $`F_{\mu \nu }`$ and $`𝒟_\mu `$ depend on $`\mathrm{\Theta }`$.
Next we consider the chiral superfields coupled to the vector superfields. The gauge transformations for the fundamental, anti-fundamental, adjoint and bi-fundamental chiral superfields are given by
$`\mathrm{\Phi }`$ $``$ $`(e^{i\mathrm{\Lambda }})_{}\mathrm{\Phi },`$
$`\stackrel{~}{\mathrm{\Phi }}`$ $``$ $`\stackrel{~}{\mathrm{\Phi }}(e^{i\mathrm{\Lambda }})_{},`$
$`\mathrm{\Phi }_{adj}`$ $``$ $`(e^{i\mathrm{\Lambda }})_{}\mathrm{\Phi }_{adj}(e^{i\mathrm{\Lambda }})_{},`$
$`\mathrm{\Phi }_{N\overline{M}}`$ $``$ $`(e^{i\mathrm{\Lambda }^{(1)}})_{}\mathrm{\Phi }_{N\overline{M}}(e^{i\mathrm{\Lambda }^{(2)}})_{},`$ (27)
respectively.
We can see that the supersymmetric gauge invariant actions including kinetic terms are
$`S_\mathrm{\Phi }`$ $`=`$ $`{\displaystyle d^4xd^2\theta d^2\overline{\theta }\left(\mathrm{\Phi }^{}(e^{2V})_{}\mathrm{\Phi }\right)}`$
$`=`$ $`{\displaystyle }d^4x((𝒟_\mu A^{})(𝒟^\mu A)i\psi ^{}\overline{\sigma }^\mu (𝒟_\mu \psi )A^{}DA`$
$`i\sqrt{2}A^{}\lambda \psi +i\sqrt{2}\psi ^{}\lambda ^{}A+F^{}F),`$
$`S_{\stackrel{~}{\mathrm{\Phi }}}`$ $`=`$ $`{\displaystyle d^4xd^2\theta d^2\overline{\theta }\left(\stackrel{~}{\mathrm{\Phi }}(e^{2V})_{}\stackrel{~}{\mathrm{\Phi }}^{}\right)}`$
$`=`$ $`{\displaystyle }d^4x((𝒟_\mu \stackrel{~}{A})(𝒟^\mu \stackrel{~}{A}^{})i\stackrel{~}{\psi }\sigma ^\mu (𝒟_\mu \stackrel{~}{\psi }^{})+\stackrel{~}{A}D\stackrel{~}{A}^{}`$
$`i\sqrt{2}\stackrel{~}{A}\lambda ^{}\stackrel{~}{\psi }^{}+i\sqrt{2}\stackrel{~}{\psi }\lambda \stackrel{~}{A}^{}+\stackrel{~}{F}\stackrel{~}{F}^{}),`$
$`S_{\mathrm{\Phi }_{adj}}`$ $`=`$ $`{\displaystyle d^4xd^2\theta d^2\overline{\theta }\frac{1}{k}\mathrm{Tr}\left((e^{2V})_{}\mathrm{\Phi }_{adj}^{}(e^{2V})_{}\mathrm{\Phi }_{adj}\right)}`$
$`=`$ $`{\displaystyle }d^4x{\displaystyle \frac{1}{k}}\mathrm{Tr}((𝒟_\mu A_{adj}^{})(𝒟^\mu A_{adj})i\psi _{adj}^{}\overline{\sigma }^\mu (𝒟_\mu \psi _{adj})`$
$`A_{adj}^{}DA_{adj}+A_{adj}DA_{adj}^{}`$
$`i\sqrt{2}A_{adj}^{}\lambda \psi _{adj}+i\sqrt{2}\psi _{adj}^{}\lambda ^{}A_{adj}`$
$`i\sqrt{2}A_{adj}\lambda ^{}\psi _{adj}^{}+i\sqrt{2}\psi _{adj}\lambda A_{adj}^{}+F_{adj}^{}F_{adj}),`$
$`S_{\mathrm{\Phi }_{N\overline{M}}}`$ $`=`$ $`{\displaystyle d^4xd^2\theta d^2\overline{\theta }\mathrm{Tr}\left((e^{2V^{(2)}})_{}\mathrm{\Phi }_{N\overline{M}}^{}(e^{2V^{(1)}})_{}\mathrm{\Phi }_{N\overline{M}}\right)}`$ (28)
$`=`$ $`{\displaystyle }d^4x\mathrm{Tr}((𝒟_\mu A_{N\overline{M}}^{})(𝒟^\mu A_{N\overline{M}})i\psi _{N\overline{M}}^{}\overline{\sigma }^\mu (𝒟_\mu \psi _{N\overline{M}})`$
$`A_{N\overline{M}}^{}D^{(1)}A_{N\overline{M}}+A_{N\overline{M}}D^{(2)}A_{N\overline{M}}^{}`$
$`i\sqrt{2}A_{N\overline{M}}^{}\lambda ^{(1)}\psi _{N\overline{M}}+i\sqrt{2}\psi _{N\overline{M}}^{}(\lambda ^{(1)})^{}A_{N\overline{M}}`$
$`i\sqrt{2}A_{N\overline{M}}(\lambda ^{(2)})^{}\psi _{N\overline{M}}^{}+i\sqrt{2}\psi _{N\overline{M}}\lambda ^{(2)}A_{N\overline{M}}^{}+F_{N\overline{M}}^{}F_{N\overline{M}}).`$
In the $`\mathrm{\Theta }=0`$ case $`\mathrm{Tr}(e^{2V}\mathrm{\Phi }_{adj}^{}e^{2V}\mathrm{\Phi }_{adj})/k`$ is equivalent to
$$\underset{a,b=1}{\overset{N^2}{}}\mathrm{\Phi }_{adj}^{a}{}_{}{}^{}e^{2_{c=1}^{N^2}V^c(T_{adj}^c)_{ab}}\mathrm{\Phi }_{adj}^b,$$
(29)
where $`T_{adj}^c`$ is the matrix of the adjoint representation and we have used $`e^YXe^Y=X+[X,Y]+\frac{1}{2}[Y[Y,X]]+\mathrm{}`$. However the generalization of (29) to the $`\mathrm{\Theta }0`$ case is not noncommutative gauge invariant. For the anti-fundamental chiral superfield the similar phenomena can be shown and in general we should use the matrix of the fundamental representation $`T^a`$ only.
We note that there are no derivative terms of the auxiliary fields $`D`$ in the actions (28) and typical scalar potential are the form of $`A^{}AA^{}A`$, which is different from the one from the superpotential. In it has been shown that the noncommutative complex scalar field theory with the interaction $`A^{}AA^{}A`$ does not suffer from IR divergences at one-loop insertions level. It is also seen that the classical moduli space of vacua is unchanged by varying $`\mathrm{\Theta }`$.
Finally as in the commutative case we can obtain the transformation of the supersymmetry in the Wess-Zumino gauge
$`\delta _\xi A`$ $`=`$ $`\sqrt{2}\xi \psi ,`$
$`\delta _\xi \psi `$ $`=`$ $`i\sqrt{2}\sigma ^\mu \overline{\xi }(𝒟_\mu A)+\sqrt{2}\xi F,`$
$`\delta _\xi F`$ $`=`$ $`i\sqrt{2}\overline{\xi }\overline{\sigma }^\mu (𝒟_\mu \psi )2i\overline{\xi }\overline{\lambda }A,`$
$`\delta _\xi A_\mu `$ $`=`$ $`i\overline{\lambda }\overline{\sigma }^\mu \xi +i\overline{\xi }\overline{\sigma }^\mu \lambda ,`$
$`\delta _\xi \lambda `$ $`=`$ $`\sigma ^{\mu \nu }\xi F_{\mu \nu }+i\xi D,`$
$`\delta _\xi D`$ $`=`$ $`\xi \sigma ^\mu (𝒟_\mu \overline{\lambda })(𝒟_\mu \lambda ^{(a)})\sigma ^\mu \overline{\xi }.`$ (30)
This formula is valid for the chiral superfields of any representation of $`G`$.
The form of the noncommutative gauge invariant superpotential are constrained as stated for the component fields. Using (10), we can easily write down the action of the component fields for any superpotential which is renormalizable at $`\mathrm{\Theta }=0`$.
It is possible to generalize these considerations to the extended superspace. On the other hand, we can construct the action with the extended supersymmetry by the $`N=1`$ superfields for commutative case. We can easily construct the action with non-vanishing $`\mathrm{\Theta }`$ corresponding to the commutative action with the extended supersymmetry. Even at $`\mathrm{\Theta }0`$, these action has the extended supersymmetry because of the existence of the R symmetry which rotates the generators of the supersymmetry. In fact the $`U(N)`$ noncommutative gauge theory with one adjoint chiral superfield, $`N_f`$ fundamental and $`N_f`$ anti-fundamental chiral superfields and $`W=\sqrt{2}_{i=1}^{N_f}\stackrel{~}{\mathrm{\Phi }}_{(i)}\mathrm{\Phi }_{adj}\mathrm{\Phi }_{(i)}`$ has $`N=2`$ supersymmetry. We can also obtain the noncommutative $`N=4`$ supersymmetric action with $`W=\mathrm{Tr}(\mathrm{\Phi }_{adj}^{(1)}(\mathrm{\Phi }_{adj}^{(2)}\mathrm{\Phi }_{adj}^{(3)}\mathrm{\Phi }_{adj}^{(3)}\mathrm{\Phi }_{adj}^{(2)}))`$, where $`\mathrm{\Phi }_{adj}^{(i)}`$ are three adjoint chiral superfields.
The effective theories of the D-branes on the orbifold are the quiver gauge theories or the elliptic models which have bi-fundamental matter. Thus it is interesting that the action for the supersymmetric gauge theories with bi-fundamental matters can be constructed.
In this paper, we have considered the $`N=1`$ supersymmetric theories on the noncommutative $`𝐑^4`$. We have constructed the $`N=1`$ supersymmetric action for the $`U(N)`$ vector multiplets and chiral multiplets of the fundamental, anti-fundamental and adjoint representations of the gauge group. The actions for gauge fields of the products gauge groups and its bi-fundamental matters have also been obtained. We have been argued that only these gauge groups and the matters are possible for the noncommutative gauge theories. We have also found that the scalar potentials have some characteristic forms and discussed the problem of the derivative terms of the auxiliary fields.
It is interesting to generalize the results obtained in this paper to the nonlinearly realized supersymmetry. This is important because the supersymmetric DBI action which is the effective theory on a D-brane has this symmetry and has been used for the instanton in the D-brane with the $`B`$ field which is related to the noncommutative instanton .
Acknowledgements
I would like to thank T. Kawano, T. Yanagida and S-K. Yang for useful conversations. I would also like to thank K. Kurosawa and Y. Nomura for discussions. This work was supported in part by JSPS Research Fellowships for Young Scientists.
Note added:
As this article was being completed, we received the preprint which substantially overlap the present work.
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# Many-body CPA for the Holstein-DE model
## I Introduction
In this paper we study the Holstein-DE (double exchange) model
$$H=\underset{ij\sigma }{}t_{ij}^\text{ }c_{i\sigma }^{}c_{j\sigma }^\text{ }J\underset{i}{}𝐒_i^\text{ }\stackrel{}{\sigma }_i^\text{ }h\underset{i}{}L_i^zg\underset{i}{}n_i^\text{ }\left(b_i^{}+b_i^\text{ }\right)+\omega \underset{i}{}n_i^b$$
(1)
where $`i`$ and $`j`$ are site indices, $`c_{i\sigma }^{}`$ and $`b_i^{}`$ ($`c_{i\sigma }^\text{ }`$ and $`b_i^\text{ }`$) create (annihilate) an electron of spin $`\sigma `$ and a phonon respectively, $`𝐒_i`$ is a local spin, $`\stackrel{}{\sigma }_i=(1/2)_{\sigma \sigma ^{}}c_{i\sigma }^{}\stackrel{}{\sigma }_{\sigma \sigma ^{}}^\text{ }c_{i\sigma ^{}}^\text{ }`$ is the electron spin operator ($`\stackrel{}{\sigma }_{\sigma \sigma ^{}}`$ being the Pauli matrix), $`L_i^z=S_i^z+\sigma _i^z`$ is the $`z`$-component of angular momentum, $`n_i^\text{ }=_\sigma c_{i\sigma }^{}c_{i\sigma }^\text{ }`$ and $`n_i^b=b_i^{}b_i^\text{ }`$. The parameter $`t_{ij}`$ is the hopping integral, $`J>0`$ is the Hund coupling, $`h`$ is the Zeeman energy, $`g`$ is the electron-phonon coupling strength and $`\omega `$ is the Einstein phonon energy. $`H`$ is a model for the colossal magnetoresistance (CMR) manganite compounds with the double degeneracy of the conduction band neglected and a simplified form assumed for the electron-phonon coupling and phonon dispersion. The electron-phonon coupling in Eq. (1) is of the breathing-mode form, i.e. $`g^{}_in_ix_i`$ in the classical limit (where $`x_i`$ is the phonon displacement), but we regard it as an effective Jahn-Teller coupling.
Hamiltonian (1) was first studied by Röder et al , who treated the Hund coupling using a mean-field approximation and the electron-phonon coupling using a variational Lang-Firsov approximation. The same authors later used a similar method to study a more realistic model for CMR systems . In this paper we treat both the Hund and the electron-phonon coupling using an extension of a many-body coherent potential approximation (CPA) previously derived for the DE model . The CPA treats the Hund coupling better than mean-field theory and has the advantage over Lang-Firsov variational methods that the whole of the electronic spectrum can be studied, not just the coherent polaron band near the Fermi energy. In the limit of classical spins and phonons Millis et al used dynamical mean-field theory (DMFT) to study another more realistic model for CMR materials . Here however we concentrate on the effects of quantisation of the phonons.
Our approach has many similarities with studies of the Holstein model using DMFT, which have been carried out for the classical phonon and empty-band limits in which the model is a one-electron problem. Indeed the standard dynamical CPA is equivalent to DMFT for one-electron problems such as the binary alloy , the DE and Holstein models in the empty-band limit , and the DE model with classical local spins . DMFT should be regarded as the correct extension of the CPA to many-body problems . For the current many-body problem we regard our CPA as an approximate solution of DMFT, or as an extrapolation from the one-electron case. The CPA has the advantage of relative analytic simplicity, but does not treat the many-body dynamics as well as DMFT, retaining too much one-electron character.
The many-body CPA derived for the finite $`S`$ DE model in Ref. and Ref. was based on Hubbard’s scattering correction approximation for the Hubbard model . Hubbard’s approximation was derived by decoupling Green function equations of motion (EOM) according to an alloy analogy in which electrons of one spin are frozen while the propagation of those of the opposite spin is considered (within the CPA). Although more modern formulations of the CPA exist Hubbard’s EOM approach was found to be particularly suitable for extension to the DE model, where the possibility of electrons exchanging (spin) angular momentum with local spins complicates the problem. The resulting many-body CPA was exact in the atomic limit $`t_{ij}0`$ and recovered the one-electron CPA/DMFT in the empty-band and classical spin ($`S\mathrm{}`$) limits.
In Sec. II we solve the atomic limit of Hamiltonian (1), and in Sec. III we extend our many-body CPA to the Holstein-DE model. The properties of the CPA solution are discussed in Sec. IV, and the special case ($`J=h=0`$) of the Holstein model is considered in Sec. V. We give a summary in Sec. VI.
## II The atomic limit
In the atomic limit $`t_{ij}0`$ $`H`$ is exactly solvable using the canonical transformation $`H\stackrel{~}{H}=e^\nu ^{}He^\nu `$ where $`\nu =\frac{g}{\omega }n(b^{}b)`$ (Ref. ) (we drop site indices). In the presence of electron-phonon coupling the phonon potential is of the displaced harmonic oscillator form, and the effect of the canonical transformation is to shift the operators to take account of this: $`bb+\frac{g}{\omega }n`$ and $`c_\sigma ^\text{ }Xc_\sigma ^\text{ }`$ where $`X=\mathrm{exp}(g/\omega (b^{}b))`$. This transformation decouples the Hamiltonian $`\stackrel{~}{H}=H_b+H_f`$ into a bosonic component $`H_b=\omega n^b`$ and a fermionic component $`H_f=J𝐒\stackrel{}{\sigma }hL^zg^2/\omega n^2`$ where $`g^2/\omega =:\lambda \omega `$ is the binding energy of a polaron.
The one-electron Green function $`G_\sigma (t)=i\theta (t)\{c_\sigma (t),c_\sigma ^{}\}`$ can be separated into fermionic and bosonic traces using the invariance of the trace under cyclic permutations and $`e^\nu ^{}e^\nu =1`$,
$`G_\sigma (t)=i\theta (t)[{\displaystyle \frac{\mathrm{Tr}_f\{e^{\beta H_f}e^{iH_ft}c_\sigma e^{iH_ft}c_\sigma ^{}\}}{\mathrm{Tr}_f\{e^{\beta H_f}\}}}{\displaystyle \frac{\mathrm{Tr}_b\{e^{\beta H_b}e^{iH_bt}Xe^{iH_bt}X^{}\}}{\mathrm{Tr}_b\{e^{\beta H_b}\}}}`$ (2)
$`+{\displaystyle \frac{\mathrm{Tr}_f\{e^{\beta H_f}c_\sigma ^{}e^{iH_ft}c_\sigma e^{iH_ft}\}}{\mathrm{Tr}_f\{e^{\beta H_f}\}}}{\displaystyle \frac{\mathrm{Tr}_b\{e^{\beta H_b}X^{}e^{iH_bt}Xe^{iH_bt}\}}{\mathrm{Tr}_b\{e^{\beta H_b}\}}}].`$ (3)
We evaluate the bosonic traces directly and the fermionic traces using the equation of motion (EOM) method, and in the energy representation $`G_\sigma (ϵ)=_{\mathrm{}}^{\mathrm{}}dte^{iϵt}G_\sigma (t)`$ obtain
$`G_\sigma (ϵ)={\displaystyle \underset{r=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{I}_r\left(2\lambda \sqrt{b(\omega )(b(\omega )+1)}\right)}{(2S+1)\mathrm{exp}(\lambda (2b(\omega )+1))}}{\displaystyle \underset{\alpha =\pm }{}}[{\displaystyle \frac{e^{r\beta \omega /2}W_{\alpha \sigma }^{0+}+e^{r\beta \omega /2}W_{\alpha \sigma }^0}{ϵ\alpha JS/2+h\sigma /2+\lambda \omega (1+2\delta _{\alpha +})+\omega r}}`$ (4)
$`+{\displaystyle \frac{e^{r\beta \omega /2}W_{\alpha \sigma }^{1+}+e^{r\beta \omega /2}W_{\alpha \sigma }^1}{ϵ+\alpha J/2(S+1)+h\sigma /2+\lambda \omega (1+2\delta _{\alpha +})+\omega r}}].`$ (5)
Here the weight factors
$`W_{\alpha \sigma }^{0\gamma }=(S+1)n_\sigma ^\alpha n_\sigma ^\gamma \alpha \sigma S^zn_\sigma ^\alpha n_\sigma ^\gamma +\left(1\delta _{\alpha \gamma }\right)S^\sigma \sigma ^{+\sigma }`$ (7)
$`W_{\alpha \sigma }^{1\gamma }=Sn_\sigma ^\alpha n_\sigma ^\gamma +\alpha \sigma S^zn_\sigma ^\alpha n_\sigma ^\gamma \left(1\delta _{\alpha \gamma }\right)S^\sigma \sigma ^{+\sigma },`$ (8)
$`\alpha `$, $`\gamma =\pm `$, $`\delta _{\alpha +}=1`$ for $`\alpha =+`$ and $`0`$ for $`\alpha =`$, $`\mathrm{I}_r`$ is the modified Bessel function, $`b(\omega )=1/(\mathrm{exp}(\beta \omega )1)`$ is the Bose function, and we define $`n_\sigma ^+=n_\sigma `$, $`n_\sigma ^{}=1n_\sigma `$, and $`S^\sigma \sigma ^{+\sigma }=S^{}\sigma ^+`$ for $`\sigma =`$ and $`S^+\sigma ^{}`$ for $`\sigma =`$.
In the ($`J=h=0`$) case of the Holstein model Eq. (5) reduces to the formula
$`G_\sigma (ϵ)={\displaystyle \underset{r=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{I}_r\left(2\lambda \sqrt{b(\omega )(b(\omega )+1)}\right)}{\mathrm{exp}(\lambda (2b(\omega )+1))}}[{\displaystyle \frac{e^{r\beta \omega /2}n_\sigma ^{}n_\sigma +e^{r\beta \omega /2}n_\sigma ^{}n_\sigma ^{}}{ϵ+\lambda \omega +\omega r}}`$ (9)
$`+{\displaystyle \frac{e^{r\beta \omega /2}n_\sigma n_\sigma +e^{r\beta \omega /2}n_\sigma n_\sigma ^{}}{ϵ+3\lambda \omega +\omega r}}],`$ (10)
but we are mostly interested in the strong Hund-coupling limit, so we shift the energy to have the zero near the Fermi level, $`ϵϵJS/2`$, and let $`J\mathrm{}`$. In this limit we find
$$G_\sigma (ϵ)=\underset{r=\mathrm{}}{\overset{\mathrm{}}{}}\frac{\mathrm{I}_r\left(2\lambda \sqrt{b(\omega )(b(\omega )+1)}\right)}{(2S+1)\mathrm{exp}(\lambda (2b(\omega )+1))}\frac{e^{r\beta \omega /2}W_\sigma ^++e^{r\beta \omega /2}W_\sigma ^{}}{ϵ+h\sigma /2+\omega r+\lambda \omega }$$
(11)
where the weight factors
$`W_\sigma ^+=`$ $`(S+1+\sigma S^z)n_\sigma +S^\sigma \sigma ^{+\sigma }`$ $`=(2S+1)\left({\displaystyle \frac{n}{2}}+\sigma \sigma ^z\right)`$ (13)
$`W_\sigma ^{}=`$ $`(1n)(S+1+\sigma S^z)`$ $`=(S+1)(1n)+\sigma S^z2S\sigma \sigma ^z.`$ (14)
The paramagnetic state spectrum of Eq. (11) is plotted in Fig. 1 for the classical spin limit $`S\mathrm{}`$ at quarter-filling $`n=0.5`$. The spectrum consists of delta-function peaks separated in energy by $`\omega `$, and for clarity we include the peaks’ envelope curve in Fig. 1. Note that the symmetry of the spectrum about zero energy is due to choice of filling $`n=0.5`$; in general the lower and upper (zero temperature) ‘bands’ have weights $`n`$ and $`1n`$ respectively. By counting weights it may be seen that for any $`n`$ the zero temperature chemical potential $`\mu (T=0)=0`$ lies in the zero energy peak in the middle of the pseudogap.
## III The CPA Green function
We now derive a many-body CPA for the one-electron Green function $`G_\sigma ^{ij}(ϵ)`$ of the full Hamiltonian (1). As discussed in the introduction we proceed by decoupling equations of motion (EOM), adapting decoupling approximations previously used for the DE model . Recall that with the fermionic definition of Green functions, $`A;C_ϵ=i_0^{\mathrm{}}dt\mathrm{exp}(iϵt)\{A(t),C\}`$, the EOM is
$$ϵA;C_ϵ=\{A,C\}+[A,H];C_ϵ.$$
(15)
As in Ref. and in the original version of this method due to Hubbard we split the Green function into a low-energy component $`G_\sigma ^{ij}(ϵ)=n_i^{}c_{i\sigma }^\text{ };c_{j\sigma }^{}_ϵ`$ and a high-energy component $`G_\sigma ^{ij+}(ϵ)=n_i^+c_{i\sigma }^\text{ };c_{j\sigma }^{}_ϵ`$, i.e. $`G_\sigma ^{ij}(ϵ)=_{\alpha =\pm }G_\sigma ^{ij\alpha }(ϵ)`$. To close the system of EOM we in fact need to introduce the Green functions
$`S_\sigma ^{ij\alpha }(\stackrel{}{r},ϵ)`$ $`=`$ $`\mathrm{\Gamma }_i(\stackrel{}{r})n_i^\alpha c_{i\sigma }^\text{ };c_{j\sigma }^{}_ϵ`$ (17)
$`T_\sigma ^{ij\alpha }(\stackrel{}{r},ϵ)`$ $`=`$ $`\mathrm{\Gamma }_i(\stackrel{}{r})n_i^\alpha S_i^\sigma c_{i\sigma }^\text{ };c_{j\sigma }^{}_ϵ`$ (18)
where $`\stackrel{}{r}=(\rho ,\varphi ,\theta )`$ is a parameter and the operator $`\mathrm{\Gamma }_i(\stackrel{}{r})=\mathrm{exp}(\rho S_i^z)\mathrm{exp}(\varphi b_i^{})\mathrm{exp}(\theta b_i)`$. Note that $`\mathrm{\Gamma }`$ is a generating function for the operators $`S^z`$, $`b^{}`$ and $`b`$, so that $`^n/\varphi ^n\mathrm{\Gamma }_i(\stackrel{}{r})=(b_i^{})^n\mathrm{\Gamma }_i(\stackrel{}{r})`$ for instance. This is convenient in allowing us to close the system of EOM with a minimal number of equations.
When writing the EOM we use the convenient commutation identities $`[e^{\varphi b^{}},b]=\varphi e^{\varphi b^{}}`$ and $`[e^{\theta b},b^{}]=\theta e^{\theta b}`$, and the Feynman operator disentanglement relation $`e^{A+B}=e^Ae^Be^{1/2[A,B]}`$, which holds if $`[[A,B],A]=[[A,B],B]=0`$. We also work for $`\sigma =`$; the $`\sigma =`$ equations can be obtained using the symmetry of $`H`$. We introduce the operator
$$K^\alpha (\theta ,\varphi )=\omega \left(\varphi \frac{}{\varphi }\theta \frac{}{\theta }\right)+g\left(\frac{}{\varphi }+\frac{}{\theta }+\theta +\delta _{\alpha +}(\theta \varphi )\right),$$
(19)
and obtain the (exact) EOM
$`\left[ϵ+{\displaystyle \frac{h}{2}}+{\displaystyle \frac{J}{2}}{\displaystyle \frac{}{\rho }}+K^\alpha (\theta ,\varphi )\right]S_{}^{ij\alpha }(\stackrel{}{r},ϵ)+{\displaystyle \frac{J}{2}}e^{\rho \delta _{\alpha +}}T_{}^{ij\alpha }(\stackrel{}{r},ϵ)=\delta _{ij}\mathrm{\Gamma }(\stackrel{}{r})n_{}^\alpha `$ (20)
$`+\mathrm{\Gamma }_i(\stackrel{}{r})[n_i^\alpha ,H_0]c_i;c_j^{}_ϵ+{\displaystyle \underset{k}{}}t_{ik}\mathrm{\Gamma }_i(\stackrel{}{r})n_i^\alpha c_k;c_j^{}_ϵ`$ (21)
for $`S_{}^{ij\alpha }(\stackrel{}{r},ϵ)`$ and
$`\left[ϵ+{\displaystyle \frac{h}{2}}{\displaystyle \frac{J}{2}}\left(\delta _\alpha +{\displaystyle \frac{}{\rho }}\right)+K^\alpha (\theta ,\varphi )\right]T_{}^{ij\alpha }(\stackrel{}{r},ϵ)`$ (22)
$`+{\displaystyle \frac{J}{2}}e^{\rho \delta _{\alpha +}}\left[S(S+1)+\alpha {\displaystyle \frac{}{\rho }}{\displaystyle \frac{^2}{\rho ^2}}\right]S_{}^{ij\alpha }(\stackrel{}{r},ϵ)=\alpha \delta _{ij}\mathrm{\Gamma }(\stackrel{}{r})S^{}\sigma ^+`$ (23)
$`+\mathrm{\Gamma }_i(\stackrel{}{r})[n_i^\alpha ,H_0]S_i^{}c_i;c_j^{}_ϵ+{\displaystyle \underset{k}{}}t_{ik}\mathrm{\Gamma }_i(\stackrel{}{r})n_i^\alpha S_i^{}c_k;c_j^{}_ϵ`$ (24)
for $`T_{}^{ij\alpha }(\stackrel{}{r},ϵ)`$. Here $`H_0`$ is the electron hopping term of the Hamiltonian and we have dropped site indices in the expectations, assuming a homogeneous state. Equations (21) and (24) should be compared with their analogues for the DE model case: equations (4) and (5) in Ref. .
As usual we now neglect the penultimate Green functions in equations (21) and (24) (the ones containing $`[n_{i\sigma }^\alpha ,H_0]`$). This corresponds to making the alloy analogy. We treat the final Green functions in these equations using a CPA and treat all other terms exactly. In fact we use the approximations
$`{\displaystyle \underset{k}{}}t_{ik}\mathrm{\Gamma }_i(\stackrel{}{r})n_i^\alpha c_k^\text{ };c_j^{}_ϵ\mathrm{\Gamma }(\stackrel{}{r})n_{}^\alpha \left({\displaystyle \underset{k}{}}t_{ik}G_{}^{kj}(ϵ)J_{}(ϵ)G_{}^{ij}(ϵ)\right)+J_{}(ϵ)S_{}^{ij\alpha }(\stackrel{}{r},ϵ)`$ (26)
$`{\displaystyle \underset{k}{}}t_{ik}\mathrm{\Gamma }_i(\stackrel{}{r})n_i^\alpha S_i^{}c_k^\text{ };c_j^{}_ϵ\alpha \mathrm{\Gamma }(\stackrel{}{r})S^{}\sigma ^+\left({\displaystyle \underset{k}{}}t_{ik}G_{}^{kj}(ϵ)J_{}(ϵ)G_{}^{ij}(ϵ)\right)`$ (27)
$`+J_{}(ϵ+h)T_{}^{ij\alpha }(\stackrel{}{r},ϵ)`$ (28)
where $`J_\sigma (ϵ)=ϵ\mathrm{\Sigma }_\sigma (ϵ)G_\sigma (ϵ)^1`$, $`\mathrm{\Sigma }_\sigma (ϵ)`$ and $`G_\sigma (ϵ)`$ being the self-energy and local Green function respectively. The function $`E_\sigma (ϵ)=ϵJ_\sigma (ϵ)`$ is related to the Weiss function of DMFT, $`E_\sigma (i\omega _n)𝒢_\sigma ^1(i\omega _n)`$, $`i\omega _n`$ being a fermionic Matsubara frequency. Equations (26) and (28) are generalisations of Hubbard’s scattering correction approximation , and lead to the CPA equations in the case of the DE model. We make these particular approximations since Eq. (26) and Eq. (28) are of the usual CPA form, but do not give a formal justification.
We define $`E_\sigma ^h(ϵ)=E_\sigma (ϵ+h\delta _\sigma )+\sigma h/2`$ and $`\lambda _\sigma ^{ij}(ϵ)=\delta _{ij}+_kt_{ik}G_\sigma ^{kj}(ϵ)J_\sigma (ϵ)G_\sigma ^{ij}(ϵ)`$, and from equations (21), (24), (26) and (28) obtain
$`\left[\begin{array}{cc}E_{}^h(ϵ)+\frac{J}{2}\frac{}{\rho }+K^\alpha (\theta ,\varphi )& \frac{J}{2}e^{\rho \delta _{\alpha +}}\\ \frac{J}{2}e^{\rho \delta _{\alpha +}}\left(S(S+1)+\alpha \frac{}{\rho }\frac{^2}{\rho ^2}\right)& E_{}^h(ϵ)\frac{J}{2}\left(\delta _\alpha +\frac{}{\rho }\right)+K^\alpha (\theta ,\varphi )\end{array}\right]\left(\begin{array}{c}S_{}^{ij\alpha }(\stackrel{}{r},ϵ)\\ T_{}^{ij\alpha }(\stackrel{}{r},ϵ)\end{array}\right)`$ (33)
$`\lambda _{}^{ij}(ϵ)\left(\begin{array}{c}\mathrm{\Gamma }(\stackrel{}{r})n_{}^\alpha \\ \alpha \mathrm{\Gamma }(\stackrel{}{r})S^{}\sigma ^+\end{array}\right).`$ (36)
We make no further approximations. We use the top row of Eq. (36) to eliminate $`T_{}^{ij\alpha }(\stackrel{}{r},ϵ)`$, obtaining a second-order linear (parabolic) PDE for $`S_{}^{ij\alpha }(\stackrel{}{r},ϵ)`$. We take $`i=j`$ and use $`\lambda _\sigma ^{ii}(ϵ)=1`$ (Ref. ). In the strong Hund-coupling limit $`J\mathrm{}`$, which we take with the energy origin shifted as $`ϵϵJS/2`$, this second-order PDE simplifies to the first-order linear PDE
$`\left[{\displaystyle \frac{(1+S)E_{}^h(ϵ)+SE_{}^h(ϵ)}{2S+1}}+{\displaystyle \frac{E_{}^h(ϵ)E_{}^h(ϵ)}{2S+1}}{\displaystyle \frac{}{\rho }}+K(\theta ,\varphi )\right]S_{}(\stackrel{}{r},ϵ)=`$ (37)
$`\mathrm{\Gamma }(\stackrel{}{r}){\displaystyle \frac{(S+1+S^z)n_{}^{}+S^{}\sigma ^+}{2S+1}}.`$ (38)
Note that in this limit we may assume $`\alpha =`$. We change $`(\theta ,\varphi )`$ variables to $`\mathrm{\Phi }=(\theta g/\omega )(\varphi +g/\omega )`$ and $`\mathrm{\Theta }=1/\omega \mathrm{ln}(\varphi +g/\omega )`$, in terms of which
$$K(\mathrm{\Theta },\mathrm{\Phi })=g\left(\frac{g}{\omega }+\mathrm{\Phi }e^{\omega \mathrm{\Theta }}\right)+\frac{}{\mathrm{\Theta }}.$$
(39)
In the new $`(\rho ,\mathrm{\Theta },\mathrm{\Phi })`$ system of variables Eq. (38) contains derivatives with respect to $`\rho `$ and $`\mathrm{\Theta }`$ only, facilitating its solution. We find (see appendix) for $`\rho =\theta =\varphi =0`$
$`G_{}(ϵ)={\displaystyle \underset{m,n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^n\lambda ^{(m+n)/2}}{m!n!e^\lambda }}`$ (40)
$`\times {\displaystyle \frac{e^{g/\omega b^{}}(b^{})^m(bg/\omega )^ne^{g/\omega b}[(S+1+S^z)n_{}^{}+S^{}\sigma ^+]}{(S+1+S^z)E_{}^h(ϵ)+(SS^z)E_{}^h(ϵ)+(2S+1)(\lambda \omega +(mn)\omega )}}`$ (41)
where $``$ denotes quantum and statistical averaging and $`S^z`$ in the denominator acts on the left.
In principle the average in Eq. (41) should be determined self-consistently, but this is difficult to carry out. In previous many-body CPAs, e.g. Hubbard’s for the Hubbard model and ours for the DE model, it was found that the hopping does not affect the total weight in a band near a given atomic limit peak, at least when the bands are separated so that the weight associated with a given atomic limit peak is a meaningful quantity. For simplicity we therefore assume that all averages take their atomic limit values. Note that owing to the degeneracy of the atomic limit states this says nothing about the spin polarisation. This assumption means that we cannot take account of the effects of electron hopping on the phonon distribution.
The right-hand side of Eq. (41) then depends on the half-bandwidth $`W`$ only through $`E_\sigma ^h(ϵ)`$. We change summation variables $`(m,n)`$ to $`r=mn`$ and $`s=m+n`$ and use local spin projection operators $`P(S^z=m^z)`$ to pull the denominator of Eq. (41) out of the average. Since $`lim_{W0}[(S+1+m^z)E_{}^h(ϵ)+(Sm^z)E_{}^h(ϵ)]=(2S+1)(ϵ+h/2)`$ we can match our averages (summed over $`s`$) with atomic limit peak weights to see that
$`G_{}(ϵ)={\displaystyle \underset{r=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{I}_r\left(2\lambda \sqrt{b(\omega )(b(\omega )+1)}\right)}{\mathrm{exp}(\lambda (2b(\omega )+1))}}`$ (42)
$`\times {\displaystyle \frac{e^{r\beta \omega /2}(1n)(S+1+S^z)+e^{r\beta \omega /2}[(S+1+S^z)n_{}+S^{}\sigma ^+]}{(S+1+S^z)E_{}^h(ϵ)+(SS^z)E_{}^h(ϵ)+(2S+1)(\lambda \omega +r\omega )}}.`$ (43)
This should be compared with the atomic limit expression Eq. (11), to which Eq. (43) reduces as $`W0`$.
Now in the case of the empty-band limit of the Holstein model one is used to obtaining a CPA/DMFT expression for $`G_\sigma `$ in the form of a continued fraction. In Eq. (43) we have a simpler expression involving a sum over the atomic limit peaks, despite the more complex nature of the problem which we are considering (the many-electron case with both Holstein and DE interactions). One might suspect that our CPA is cruder than the one-electron CPA. In fact our expression for $`G_{}`$ for the Holstein model (given in Sec. V) in the limit $`n0`$ is equivalent to making the approximation
$$E_\sigma ^h(ϵ+r\omega )E_\sigma ^h(ϵ)+r\omega $$
(44)
in the one-electron CPA expression. We will mainly consider the case of an elliptic bare density of states (DOS), $`D(ϵ)=2/(\pi W^2)\sqrt{W^2ϵ^2}`$, for which it be shown that $`E(ϵ)=ϵW^2G(ϵ)/4`$. For the elliptic DOS approximation Eq. (44) is thus equivalent to neglecting energy shifts in the Green function on the right-hand side of the CPA equation. Since we do not recover the one-electron CPA/DMFT as $`n0`$, unlike in the case of the bare DE model , our CPA for the Holstein-DE model is not as good as our CPA for the DE model. We choose however to accept the increased crudeness of the approximation in return for the greatly increased simplicity; a CPA which correctly reduced to the one-electron CPA as $`n0`$ would probably be analytically intractable in the many-body case.
### A Calculation of Curie temperature
Using mean-field arguments Millis et al claimed that the bare DE model predicts a Curie temperature $`T_\mathrm{C}`$ at least an order of magnitude too large. However, subsequent more reliable treatments of the DE model taking into account quantum fluctuations showed that the DE model’s $`T_\mathrm{C}`$ is in fact in reasonable agreement with experiment. Now as discussed by Röder et al phonon coupling suppresses $`T_\mathrm{C}`$. We therefore calculate $`T_\mathrm{C}`$ to see if a phonon coupling strength $`g`$ exists that gives a much larger resistivity (than the $`g=0`$ case) without making $`T_\mathrm{C}`$ unphysically small.
For simplicity we work in the classical limit $`S=\mathrm{}`$, in which $`E_\sigma ^h(ϵ)=E_\sigma (ϵ)`$ since $`h1/S`$, and specialise to the case of an elliptic bare DOS, where as mentioned above $`E_\sigma (ϵ)=ϵW^2G_\sigma (ϵ)/4`$ is just a function of $`G_\sigma `$. We set $`h=0`$ in Eq. (43) and expand $`G_\sigma `$ about the paramagnetic state to first order (in $`\delta \sigma ^z`$ or $`\delta S^z`$), obtaining
$$\delta G_{}(ϵ)=\left(\underset{r}{}w_r(ϵ)e^{r\beta \omega /2}\right)\delta S^z+4\left(\underset{r}{}w_r(ϵ)\mathrm{sinh}(r\beta \omega /2)\right)\delta \sigma ^z$$
(45)
where
$$w_r(ϵ)=\frac{\mathrm{I}_r\left(2\lambda \sqrt{b(\omega )(1+b(\omega ))}\right)}{2\mathrm{exp}(\lambda (1+2b(\omega )))}\frac{(E(ϵ)+\omega r)^1}{1+(W^2/12)G^{}(ϵ)/E^{}(ϵ)},$$
(46)
$`A^{}(ϵ)=\mathrm{d}A(ϵ)/\mathrm{d}ϵ`$ and we drop spin indices on paramagnetic state quantities. From the spectral theorem $`\sigma ^z=[G_{}G_{}]/2`$ where
$$[A]=_{\mathrm{}}^{\mathrm{}}\frac{\mathrm{d}ϵ}{\pi }f(ϵ\mu )\mathrm{Im}\left[A(ϵi0)\right],$$
(47)
$`f(ϵ\mu )`$ being the Fermi function. Applying $``$ to Eq. (45) and rearranging leads to
$$\delta \sigma ^z=\frac{_r[w_r]\mathrm{exp}(r\beta \omega /2)}{14_r[w_r]\mathrm{sinh}(r\beta \omega /2)}\delta S^z.$$
(48)
In Ref. we showed (in the DE model case) that the CPA for electronic Green functions does not give a good estimate of local spin expectations. Fortunately for $`S=\mathrm{}`$ we can use DMFT to obtain an expression for $`\delta S^z`$ in terms of $`\delta G_{}`$ which is exact in the infinite-dimensional limit. After integrating out the bosonic degrees of freedom the DMFT effective action can be written in the Matsubara representation as
$`\stackrel{~}{S}={\displaystyle \underset{n}{}}\left(\begin{array}{cc}c_n^{}& c_n^{}\end{array}\right)\left(\begin{array}{cc}E_{}(i\omega _n)+(J/2)S^z& (J/2)S^{}\\ (J/2)S^+& E_{}(i\omega _n)(J/2)S^z\end{array}\right)\left(\begin{array}{c}c_n\\ c_n\end{array}\right)\beta hS^z`$ (54)
$`+{\displaystyle _0^\beta }d\tau {\displaystyle _0^\beta }d\tau ^{}n_i(\tau )\stackrel{~}{U}(\tau \tau ^{})n_i(\tau ^{})`$ (55)
where we work for $`S=\mathrm{}`$, $`h`$ and $`J`$ finite, and the $`n`$ subscripts refer to fermionic Matsubara frequencies. The last term in Eq. (55) is an attractive Hubbard-like term, with the Fourier transform of the interaction given by $`\stackrel{~}{U}(i\omega _n)=(1/2)g^2/(\omega ^2+\omega _n^2)`$. It is retarded in imaginary time and originates from the phonon coupling . We expand $`\stackrel{~}{S}`$ about the $`h=0`$ paramagnetic state with action $`\stackrel{~}{S}_0`$ and partition function $`Z_0`$, $`\stackrel{~}{S}=\stackrel{~}{S}_0+\delta \stackrel{~}{S}+\mathrm{}`$ where
$$\delta \stackrel{~}{S}=\beta hS^z\underset{n}{}\left(\delta E_{}(i\omega _n)c_n^{}c_n^\text{ }+\delta E_{}(i\omega _n)c_n^{}c_n^\text{ }\right),$$
(56)
and in terms of $`\delta \stackrel{~}{S}`$ we have
$$\delta S^z=\frac{1}{Z_0}\mathrm{d}^2S\left(\underset{n\sigma }{}\mathrm{d}c_{n\sigma }^{}\mathrm{d}c_{n\sigma }\right)S^ze^{\stackrel{~}{S}_0}\delta \stackrel{~}{S}$$
(57)
where $`\mathrm{d}^2S`$ is the integral over the surface of the unit sphere.
Now $`\delta E_\sigma (i\omega _n)=W^2\delta G_\sigma (i\omega _n)/4`$ for an elliptic DOS, and we use the relation $`\beta ^1_ng(i\omega _n)=[g]`$, which holds for functions $`g`$ analytic off the real axis, to write Eq. (57) as
$$\delta S^z=\beta \left\{\frac{h}{3}\frac{W^2}{2}\left[\left(\frac{S(\stackrel{}{r},ϵ)}{\rho }\right)_{\stackrel{}{r}=0}\delta G_{}(ϵ)\right]\right\}.$$
(58)
Note that the effects of phonon coupling enter only implicitly via Green functions. This is expected as the electron-phonon coupling is spin-symmetric. Setting $`h=0`$ in Eq. (58) and using (45), (48) and $`S(\stackrel{}{r},ϵ)/\rho |_{\stackrel{}{r}=0}=G(ϵ)/3`$ we obtain the Curie temperature equation
$$k_\mathrm{B}T_\mathrm{C}=\frac{W^2}{6}\left[\frac{_r[Gw_r]e^{r\beta \omega /2}+4_{rs}[Gw_r][w_s]\mathrm{sinh}(\beta \omega (rs)/2)}{14_r[w_r]\mathrm{sinh}(r\beta \omega /2)}\right]$$
(59)
upon dividing by $`\delta S^z`$. We discuss the value of $`T_\mathrm{C}`$ in the next section.
## IV Results
We now discuss numerical results obtained using our CPA, for simplicity using the elliptic bare DOS and working at $`J=S=\mathrm{}`$ and $`n=0.5`$. In the spin-saturated state the minority-spin weight at low energy is of order $`1/S`$, so the classical limit $`S\mathrm{}`$ is convenient as we do not need band-shifts, which are difficult to obtain within the CPA, for consistency of the saturated ferromagnetic state. Quarter-filling $`n=0.5`$ is used because owing to the symmetry of the spectrum about zero energy the chemical potential $`\mu (n=0.5)=0`$ for all $`T`$. For a homogeneous state the doping has a qualitative effect only as $`n0`$ or 1, and the results for a more physical value $`n0.7`$ are similar in form to those at $`n=0.5`$. Note however that the quantitative predictions of the model are very sensitive to the model parameters, especially the electron-phonon coupling strength but also the doping, so if the physical doping value is used the model parameters must be adjusted to retain quantitative agreement with experiment. We take $`\omega /k_\mathrm{B}=0.05W/k_\mathrm{B}600`$K for $`W1`$eV. Zhao et al report that $`\omega /k_\mathrm{B}100`$K for La<sub>1-x</sub>Ca<sub>x</sub>MnO<sub>3</sub> (Ref. ) so this may be a bit large.
The Curie temperature $`T_\mathrm{C}`$ obtained from equations (59) and (43) is plotted against electron-phonon coupling strength $`g`$ in Fig. 2. It will be found later that $`g0.16W`$ gives reasonable values for the resistivity. For values of $`g`$ in this range $`T_\mathrm{C}`$ is only suppressed by about a factor $`2`$, so for $`W`$1eV is still compatible with experiment.
The effects of phonon coupling on the (forced) $`h=T=0`$ paramagnetic state DOS are shown in Fig. 3. At $`g=0`$ we obtain the usual elliptic band . As the coupling $`g`$ is increased the DOS broadens, small subbands are split off from the band-edges, and a pseudogap appears near the Fermi energy. At a critical value $`g_c`$ the DOS splits near zero energy, leaving a small polaron band in the gap with low weight but very large mass. Increasing $`g`$ further causes more bands to be formed in the gap, with weights equal to the relevant atomic limit peak weights. The effect of increasing the temperature $`T`$ on the DOS in the pseudogap is shown in Fig. 4 for $`g=0.18W>g_c`$. With increasing $`T`$ the DOS at the Fermi surface increases rapidly and the polaron bands are smeared out.
For $`g>g_c`$ the majority of electrons are in fully occupied bands, and the itinerant electrons lie in a polaron band of very small weight near zero energy. At $`T=0`$ this band is equivalent to the one obtained in the standard strong-coupling theory of the Holstein model , where one averages the phonons out of the Hamiltonian $`\stackrel{~}{H}`$ considering only diagonal electron hopping processes in which the number of phonons in each state is conserved. In our approximation this polaron band is damped even at $`T=0`$ (i.e. Im$`\mathrm{\Sigma }(\mu )`$0), but it would be coherent (barring damping coming from the disordered local spins, i.e. in the saturated ferromagnetic state) in an approximation which took better account of the dynamics. In the usual strong-coupling treatment, which treats only the coherent polaron band, it is found that the DOS $`D(\mu )`$ at the Fermi surface decreases with increasing $`T`$. This is not inconsistent with our finding that $`D(\mu )`$ increases with $`T`$ since our DOS includes all the spectral weight, both (ideally) coherent and incoherent. For $`n0.5`$ the DOS is no longer symmetric about zero energy; the main lower and upper bands into which the DOS is split for $`g>g_c`$ have approximate weights $`n`$ and $`1n`$ respectively. Although these large features of the spectrum vary considerably with doping the zero-temperature chemical potential is always confined to the polaron band near zero energy, moving from the bottom at $`n=0`$ to the top at $`n=1`$ (so that we obtain a band insulator in these cases).
In our CPA we have no reliable means of calculating the probability distribution function $`P(S^z)`$, so to go below $`T_\mathrm{C}`$ we use the mean-field approximation for the ferromagnetic Heisenberg model with classical spins and nearest-neighbour coupling. Since we regard our CPA as an approximation to DMFT, which is also exact in the infinite-dimensional limit, this simple approximation may not be unreasonable. We obtain the coupling constant for the Heisenberg model by matching Curie temperatures. We take this coupling constant to be temperature-independent, but in a more systematic mapping onto the Heisenberg model one would expect the coupling constant to vary with temperature. Note also that in principle the Heisenberg model’s $`P(S^z)`$ is of a different form to the DE model’s . One effect of using a mean-field approximation for the magnetisation will be to obtain the mean-field magnetisation exponent of $`1/2`$; in three dimensions we expect the magnetisation to increase more rapidly below $`T_\mathrm{C}`$, but note that Schwartz et al find that the magnetisation exponent in La<sub>0.8</sub>Sr<sub>0.2</sub>MnO<sub>3</sub> is 0.45$`\pm `$0.05.
We plot the up- and down-spin DOSs for $`T=0.005W/k_\mathrm{B}T_\mathrm{C}`$ and $`g=0.16W>g_c`$ in Fig. 5, also showing the saturated ferromagnetic and paramagnetic state DOSs for comparison. The large difference between $`D_{\mathrm{ferro}}(\mu )`$ and $`D_{\mathrm{para}}(\mu )`$ mean that for given $`T`$ we expect the paramagnetic state to have a much higher resistivity than a magnetised state. Note that there are no separated polaron bands near $`\mu `$ in the up- and down-spin DOSs, even at this low temperature $`k_\mathrm{B}T=0.1\omega `$ where quantum effects might be expected to be important. The transfer of weight to the up-spin DOS has broadened the polaron bands enough to remove the gaps in the DOS, and mixing of the down-spins with the up-spins via the Hund coupling suffices to remove the gaps from the down-spin DOS too. It therefore appears that the development of magnetisation prevents quantum effects from becoming important for this coupling strength, at least as far as the DOS is concerned.
We calculate the resistivity $`\rho `$ using the formula obtained in Ref. , plotting it against temperature in Fig. 6 for various magnetic fields $`h`$. The form of the curve is broadly in agreement with experiment , with the resistivity peak (which occurs at $`T_\mathrm{C}`$) of the correct order of magnitude for La<sub>0.75</sub>Ca<sub>0.25</sub>MnO<sub>3</sub> (Ref. ) and we find a large negative magnetoresistance near the peak. Note that for $`W`$1eV the Curie temperature $`T_\mathrm{C}`$230K and the magnetic field $`B=0.004W/(g\mu _\mathrm{B})`$20T (for $`g7/2`$). The main differences between Fig. 6 and experiment are the large residual $`T=0`$ resistivity, due to the artificial incoherence of the CPA, and the less rapid drop in the resistivity below $`T_\mathrm{C}`$ and with $`h`$, possibly due to the mean-field form used for the magnetisation. The rise in the resistivity below $`T_\mathrm{C}`$ is due to the effects of the reduced spin polarisation on the DOS. Below $`T_\mathrm{C}`$ these effects dominate over the effects on the DOS of thermal smearing, which is responsible for the fall in the resistivity above $`T_\mathrm{C}`$.
In Fig. 7 we show the effect on the resistivity of increasing hydrostatic pressure, which we model as an increase in the bandwidth. Any change in the other terms of the Hamiltonian is neglected (note that the resistivity is proportional to the lattice constant $`a`$, so a decrease in $`a`$ will only reinforce the trend observed in Fig. 7). The strong suppression of the peak and the increase in Curie temperature is in agreement with the measurements of Neumeier et al on La<sub>0.67</sub>Ca<sub>0.33</sub>MnO<sub>3</sub>, where a drop in peak resistivity of a factor of $`2`$ is observed when a pressure of 1.62GPa is applied. This change comes mainly from a decease in the effective coupling constant $`g^2/(\omega W)`$.
## V The Holstein model
We now briefly consider the special case $`J=h=0`$ where Hamiltonian (1) reduces to the Holstein model. In this case CPA equation (36) takes the form
$$\left[E_\sigma (ϵ)+K^\alpha (\theta ,\varphi )\right]\mathrm{exp}\left(\varphi b_i^{}\right)\mathrm{exp}\left(\theta b_i\right)n_i^\alpha c_{i\sigma }^\text{ };c_{i\sigma }^{}_ϵ\mathrm{exp}\left(\varphi b^{}\right)\mathrm{exp}\left(\theta b\right)n_\sigma ^\alpha ,$$
(60)
and can be solved as in Sec. III to yield
$`G_\sigma (ϵ)={\displaystyle \underset{r=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{I}_r\left(2\lambda \sqrt{b(\omega )(b(\omega )+1)}\right)}{\mathrm{exp}(\lambda (2b(\omega )+1))}}[{\displaystyle \frac{e^{r\beta \omega /2}n_\sigma ^{}n_\sigma +e^{r\beta \omega /2}n_\sigma ^{}n_\sigma ^{}}{E_\sigma (ϵ)+\lambda \omega +\omega r}}`$ (61)
$`+{\displaystyle \frac{e^{r\beta \omega /2}n_\sigma n_\sigma +e^{r\beta \omega /2}n_\sigma n_\sigma ^{}}{E_\sigma (ϵ)+3\lambda \omega +\omega r}}].`$ (62)
There is now no mixing of up- and down-spins in the problem so our CPA takes the simple form $`G(ϵ)=G_{\mathrm{AL}}(E(ϵ))`$ (where $`G_{\mathrm{AL}}`$ is the atomic limit Green function) which one would guess for a many-body CPA: Eq. (62) is just the atomic limit result Eq. (10) with $`ϵE_\sigma (ϵ)`$. The first ($`G_\sigma ^{\alpha =}`$) term of Eq. (62) corresponds to polaron bands near energy $`\lambda \omega `$, and the second ($`G_\sigma ^{\alpha =+}`$) is the bipolaronic term, corresponding to bands near $`3\lambda \omega `$. Our approximation’s reliance on the atomic limit means that all bipolaron coupling takes place on-site.
In principle we can use the spectral theorem to determine all weights self-consistently in terms of the Green functions $`G_\sigma ^\alpha `$, but for low temperature and strong coupling $`g`$ most electrons are bound as bipolarons and we may set $`n_{}n_{}n_\sigma =n/2`$. The groundstate of the Holstein model is actually believed to be either superconducting (away from half-filling and at strong coupling) or a charge density wave (near half-filling and at weak coupling) . However, determining the weights self-consistently near the homogeneous state we do not find a (second-order) transition to a charge density wave state. This is reminiscent of the CPA for the Hubbard model, where no transition to ferromagnetism or antiferromagnetism exists. We are also unable to consider superconductivity within our approximation.
Note that the true CPA/DMFT result in the empty-band limit, first obtained by Sumi using the CPA and rederived by Ciuchi et al using DMFT, is
$$G(ϵ)=(1e^{\beta \omega })\underset{n=0}{\overset{\mathrm{}}{}}\frac{e^{n\beta \omega }}{E(ϵ)A_n(ϵ)B_n(ϵ)}$$
(63)
where $`A_n`$ is the finite continued fraction
$$A_n(ϵ)=\frac{ng^2}{E(ϵ+\omega )\frac{(n1)g^2}{E(ϵ+2\omega )\frac{(n2)g^2}{{}_{}{}^{}_{}^{}\frac{g^2}{E(ϵ+n\omega )}}}}$$
(65)
and $`B_n`$ is the infinite continued fraction
$$B_n(ϵ)=\frac{(n+1)g^2}{E(ϵ\omega )\frac{(n+2)g^2}{E(ϵ2\omega )\frac{(n+3)g^2}{E(ϵ3\omega )\mathrm{}}}}.$$
(66)
As mentioned in Sec. III our result in the empty-band limit is only equivalent to the true one-electron CPA result if the approximation $`E(ϵ+r\omega )E(ϵ)+r\omega `$ is made in equations (65) and (66).
## VI Summary
In this paper we have extended our many-body CPA, developed in references and for the DE model, to study the Holstein-DE model, which we regard as a simple model for CMR materials. We were interested in effects due to the quantisation of the phonons. Our CPA has the advantage over DMFT of being analytically relatively simple, although necessarily cruder, and over variational Lang-Firsov approaches of being able to study the whole of the spectrum, not just the low-energy coherent polaron band. We solved the Holstein-DE model exactly in the atomic limit in which the CPA becomes exact and solved the CPA equations in the strong Hund-coupling limit $`J\mathrm{}`$. Using a DMFT result for the local spin polarisation in terms of electronic Green functions we obtained an equation for the Curie temperature $`T_\mathrm{C}`$. For intermediate electron-phonon coupling strength we obtained reasonable agreement with experiment for most calculated quantities, including the Curie temperature and resistivity. It appears however that for this range of coupling the development of magnetisation below $`T_\mathrm{C}`$ prevents the quantisation of the phonons from affecting the DOS near the Fermi surface even at low temperatures.
## Acknowledgements
I am grateful to DM Edwards for helpful discussions and to the UK Engineering and Physical Sciences Research Council (EPSRC) for financial support.
##
We now solve Eq. (38) for $`S_{}`$ using the method of characteristics. For compactness of notation we define the operator $`U=\left[(S+1+S^z)n_{}^{}+S^{}\sigma ^+\right]/(2S+1)`$. In the $`(\rho ,\mathrm{\Theta },\mathrm{\Phi })`$ system of variables Eq. (38) takes the form
$`\left[{\displaystyle \frac{E_{}^h(ϵ)E_{}^h(ϵ)}{2S+1}}\right]{\displaystyle \frac{S_{}}{\rho }}+{\displaystyle \frac{S_{}}{\mathrm{\Theta }}}=e^{\rho S^z}e^{\varphi (\mathrm{\Theta })b^{}}e^{\theta (\mathrm{\Theta },\mathrm{\Phi })b}U`$ (67)
$`\left[{\displaystyle \frac{g^2}{\omega }}+g\mathrm{\Phi }e^{\omega \mathrm{\Theta }}+{\displaystyle \frac{(1+S)E_{}^h(ϵ)+SE_{}^h(ϵ)}{2S+1}}\right]S_{}`$ (68)
where $`\varphi (\mathrm{\Theta })=e^{\omega \mathrm{\Theta }}g/\omega `$ and $`\theta (\mathrm{\Theta },\mathrm{\Phi })=\mathrm{\Phi }e^{\omega \mathrm{\Theta }}+g/\omega `$. The characteristic equations are
$`{\displaystyle \frac{\mathrm{d}\rho }{\mathrm{d}s}}`$ $`=`$ $`\left[{\displaystyle \frac{E_{}^h(ϵ)E_{}^h(ϵ)}{2S+1}}\right],{\displaystyle \frac{\mathrm{d}\mathrm{\Theta }}{\mathrm{d}s}}=1`$ (70)
$`{\displaystyle \frac{\mathrm{d}S_{}}{\mathrm{d}s}}`$ $`=`$ $`e^{\rho S^z}e^{\varphi (\mathrm{\Theta })b^{}}e^{\theta (\mathrm{\Theta },\mathrm{\Phi })b}U\left[{\displaystyle \frac{g^2}{\omega }}+g\mathrm{\Phi }e^{\omega \mathrm{\Theta }}+{\displaystyle \frac{(1+S)E_{}^h(ϵ)+SE_{}^h(ϵ)}{2S+1}}\right]S_{}.`$ (71)
The first two are solved immediately as
$$\rho =\rho _0+\left[\frac{E_{}^h(ϵ)E_{}^h(ϵ)}{2S+1}\right]s,\mathrm{\Theta }=s$$
(72)
where $`\rho _0`$ is an arbitrary constant and we set the constant in the $`\mathrm{\Theta }`$ equation to zero without loss of generality. These solutions are substituted into Eq. (71), which may then be written as
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}s}}\left[e^{\left(\left[(1+S)E_{}^h(ϵ)+SE_{}^h(ϵ)\right]/(2S+1)+g^2/\omega \right)s\frac{g}{\omega }\mathrm{\Phi }e^{\omega s}}S_{}\right]=e^{\left(\left[(1+S)E_{}^h(ϵ)+SE_{}^h(ϵ)\right]/(2S+1)+g^2/\omega \right)s}`$ (73)
$`\times e^{\left(\rho _0+s\left[E_{}^hE_{}^h\right]/(2S+1)\right)S^z}e^{\frac{g}{\omega }b^{}}e^{e^{\omega s}b^{}}e^{\mathrm{\Phi }e^{\omega s}\left(b\frac{g}{\omega }\right)}e^{\frac{g}{\omega }b}U.`$ (74)
We expand $`\mathrm{exp}(e^{\omega s}b^{})`$ and $`\mathrm{exp}(\mathrm{\Phi }e^{\omega s}(bg/\omega ))`$ in Eq. (74) as series and integrate to find
$`S_{}\mathrm{exp}\{({\displaystyle \frac{(1+S)E_{}^h+SE_{}^h}{2S+1}}+{\displaystyle \frac{g^2}{\omega }})s{\displaystyle \frac{g}{\omega }}\mathrm{\Phi }e^{\omega s}\}=S_0+{\displaystyle \underset{m,n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Phi }^n}{m!n!}}e^{\rho _0S^z}e^{\frac{g}{\omega }b^{}}\left(b^{}\right)^m`$ (75)
$`\times {\displaystyle \frac{(2S+1)e^{\left([(1+S+S^z)E_{}^h+(SS^z)E_{}^h]/(2S+1)+g^2/\omega +(mm)\omega \right)s}}{(1+S+S^z)E_{}^h+(SS^z)E_{}^h+(2S+1)(g^2/\omega +(mm)\omega )}}(b{\displaystyle \frac{g}{\omega }})^ne^{\frac{g}{\omega }b}U.`$ (76)
We then write the characteristics as intersections of the surfaces $`S_0=S_0(\rho ,\mathrm{\Theta },\mathrm{\Phi })`$ and $`\rho _0=\rho _0(\rho ,\mathrm{\Theta },\mathrm{\Phi })`$, and the general solution of Eq. (38) is of the form $`S_0(\rho ,\mathrm{\Theta },\mathrm{\Phi })=F(\rho _0(\rho ,\mathrm{\Theta },\mathrm{\Phi }))`$ where $`F`$ is an arbitrary function. Rearranging we obtain
$`S_{}=\mathrm{exp}\left({\displaystyle \frac{g}{\omega }}\mathrm{\Phi }e^{\omega \mathrm{\Theta }}\right)\{\mathrm{exp}[({\displaystyle \frac{(1+S)E_{}^h+SE_{}^h}{2S+1}}+\lambda \omega )\mathrm{\Theta }]F(\rho {\displaystyle \frac{E_{}^hE_{}^h}{2S+1}}\mathrm{\Theta })`$ (77)
$`+{\displaystyle \underset{m,n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Phi }^n}{m!n!}}e^{(mn)\omega \mathrm{\Theta }}{\displaystyle \frac{(2S+1)e^{\rho S^z}e^{g/\omega b^{}}(b^{})^m(bg/\omega )^ne^{g/\omega b}U}{(S+1+S^z)E_{}^h+(SS^z)E_{}^h+(2S+1)(\lambda \omega +(mn)\omega )}}\}.`$ (78)
Now from definition Eq. (17) of $`S_{}`$ it may be seen that $`S_{}`$ is of the form
$$S_{}(\rho +\frac{E_{}^hE_{}^h}{2S+1}\mathrm{\Theta },\mathrm{\Theta },\mathrm{\Phi })=\underset{m=S}{\overset{S}{}}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}a_{mn}(\rho ,\mathrm{\Phi })\mathrm{exp}\left[\left(m\frac{E_{}^hE_{}^h}{2S+1}+n\omega \right)\mathrm{\Theta }\right].$$
(79)
The final term of Eq. (78) is compatible with this form but the term proportional to $`F`$ is not, so we must have $`F0`$. Finally, in our original $`(\rho ,\theta ,\varphi )`$ system of variables
$`S_{}(\rho ,\theta ,\varphi )=\mathrm{exp}\left({\displaystyle \frac{g}{\omega }}\left(\theta {\displaystyle \frac{g}{\omega }}\right)\right){\displaystyle \underset{m,n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\theta g/\omega )^n(\varphi +g/\omega )^m}{m!n!}}`$ (80)
$`\times {\displaystyle \frac{(2S+1)e^{\rho S^z}e^{g/\omega b^{}}\left(b^{}\right)^m\left(bg/\omega \right)^ne^{g/\omega b}}{(1+S+S^z)E_{}^h+(SS^z)E_{}^h+(2S+1)(g^2/\omega +(mn)\omega )}}U.`$ (81)
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# Optical off-nuclear spectra of quasar hosts and radio galaxies
## 1 Introduction
Understanding the host galaxies of active galactic nuclei (AGN) is now recognised as an important step on the path towards reaching an understanding of AGN themselves - how they form, how they are fuelled and how the differences between the various classes of object arise. In two areas in particular the nature of the host galaxy gives us a direct insight into the workings of the AGN: galaxy properties seem to play a role in determining the radio loudness of the central engine (powerful radio sources are almost never found in spiral or disc-dominated systems – though see McHardy et al. 1994); and amongst radio-loud objects the host galaxies offer a powerful, orientation-independent means of testing models which attempt to unify different types of AGN via beaming and viewing angle effects (eg Urry & Padovani 1995). The observations presented here attempt to address both of these issues by comparing the galaxies associated with the three main types of powerful AGN: radio-quiet quasars (RQQs), radio-loud quasars (RLQs) and radio galaxies (RGs) of Fanaroff & Riley Type II.
### 1.1 Quasar hosts
Only once the problem of separating the diffuse galaxy emission from the wings of the quasar point spread function (PSF) has been overcome can one can begin to describe and classify the morphologies, brightness profiles and interaction histories of the quasar hosts. Over the last decade improvements in ground-based techniques and the advent of the Hubble Space Telescope (HST) have revolutionised our understanding of quasar host galaxies.
Evidence for mergers or interactions in the form of morphological disturbances and close companions is a common feature of these images, but a significant number of quasars are also found in what appear to be undisturbed hosts. In addition, the idea that radio-loudness is a straightforward consequence of the host galaxy type has had to be abandoned. Although some radio-quiet quasars are found in spirals (eg Hutchings et al. 1994; Örndahl, Rönnback & van Groningen 1997) in general the hosts of both radio-loud and radio-quiet quasars tend to have properties consistent with early-type galaxies (eg Véron-Cetty & Woltjer 1990; Disney et al. 1995, Hutchings & Morris 1995, Bahcall et al. 1997, McLeod, Rieke & Storrie-Lombardi 1999) and typically have luminosities $`>L^{}`$ (eg Dunlop et al. 1993, Bahcall, Kirhakos & Schneider 1994, 1995ab, 1996; Hutchings et al. 1994, Boyce et al. 1998, Hooper, Impey & Foltz 1997) and in many cases are comparable in mass to brightest cluster galaxies (BCGs).
Meanwhile, McLeod & Rieke (1995a) find evidence for a lower limit on the $`H`$-band luminosities (and hence the mass of the red, established stellar populations) of galaxies hosting radio-quiet AGN, which appears to increase as the nuclear luminosity increases, implying that the nuclear activity is closely linked to the mass of the bulge component of the host. This impression has been reinforced by McLure et al. (1999), who find that all the RQQs with $`M_R24`$ in their HST sample occur in massive elliptical galaxies, with only the least luminous radio-quiet objects lying in disc-dominated hosts.
Many long-held views about the triggering of nuclear activity and the origins of radio loudness are currently being reassessed in the light of these imaging studies, but images cannot tell the whole story. A completely independent way of characterising the host galaxies of AGN is via analysis and classification of their stellar populations and starformation histories. The aim of the observations described in this paper was to obtain high signal-to-noise spectra of quasar hosts and radio galaxies for use in spectrophotometric modelling to determine the nature and history of their stelllar constituents.
### 1.2 Spectroscopy of quasar hosts
Previous off-nuclear spectroscopy of the host galaxies of quasars has produced mixed results. Boroson & Oke (1982) were the first to detect an unequivocably stellar continuum from the nebulosity surrounding the radio-loud quasar 3C48 and subsequent studies revealed stellar continua and emission/absorption features around several other quasars (Boroson, Oke & Green 1982; Boroson & Oke 1984; Boroson, Persson & Oke 1985; Hickson & Hutchings 1987; Hutchings & Crampton 1990). However, except for the general result that there appeared to be systematic differences between the spectra of radio-loud and radio-quiet quasar hosts, these studies produced little advance in our understanding of the relationship between RQQs, RLQs and RGs for the reasons outlined below.
In general the quasar targets were chosen virtually at random, often for the sole reason that they were ‘interesting’ and/or unusual. Until now a programme of off-nuclear spectroscopy for statistically useful and properly matched samples of RQQs and RLQs has never been carried out, nor has any attempt been made to compare the off-nuclear spectra of RLQs with equivalent off-nuclear spectra of radio galaxies.
The early work also focussed mostly on emission-line activity, with discussion of the stellar continuum being confined to classification as red or blue, and the identification of a few stellar features. Only limited attempts were made to use the form of the spectrum to investigate the composition and evolution of the stellar population.
Finally, few of the host galaxy spectra were taken sufficiently off-nucleus - eg the spectra taken by Boroson and collaborators were taken only 3<sup>′′</sup> from the quasar. This was done in the belief that any further off-nucleus the host galaxy would be too faint for a reasonable spectrum to be obtained, but inevitably resulted in significant contamination of the host galaxy spectrum by scattered light from the quasar nucleus. A scaled version of the quasar spectrum had therefore to be subtracted from the off-nuclear spectrum to reveal the spectrum of the underlying host, but the extra noise and systematic errors introduced by this process severely limited the quality of the final spectra.
Thus, the main obstacles in these previous attempts to classify the hosts of powerful AGN were the difficulty in separating the underlying starlight from the glare of the quasar and - less immediately, but still of some importance - the lack of well-defined samples of a sufficient size to carry out straightforward statistical comparisons.
As far as possible, the design of the current study was intended to circumvent both these problems with the goal of obtaining clean spectra of the stellar component of radio-loud and radio-quiet quasar hosts and radio galaxies, which could then be used to search for systematic differences and similarities between the stellar populations and of the galaxies hosting each type of activity. Despite that fact that the radio galaxies lack bright nuclear point sources we were careful to adopt exactly the same observing strategy as used for the quasars to ensure that the data would be directly comparable.
In this paper we describe the observations and present the spectra. A second paper (Nolan et al. 2000) will describe the results of spectrophotometric modelling to estimate the ages and starformation histories of the galaxies (preliminary results have already been reported by Kukula et al. 1997). The current paper is organised as follows. Section 2 describes the samples used and Section 3 details the observing and data reduction strategies chosen to optimise the amount of starlight collected by the various instruments employed throughout the study. Section 4 gives an overview of the data obtained and Sections 5 and 6 contain more detailed information on the individual spectra.
## 2 Sample selection
Dunlop et al. (1993) and Taylor et al. (1996) observed a sample of intermediate-redshift ($`0.1z0.3`$) radio-loud and -quiet quasars<sup>1</sup><sup>1</sup>1Dunlop et al. define ‘radio quiet’ objects as those with $`L_{5GHz}<10^{24}`$W Hz<sup>-1</sup>sr<sup>-1</sup>. and Fanaroff-Riley Type II radio galaxies (RGs) in the near infrared ($`K`$-band:2.2$`\mu `$m) in order to compare the luminosities and morphologies of the galaxies associated with these three main types of powerful active nucleus. The choice of waveband was informed by the low quasar:host ratio in the near infrared, which allowed accurate determination of the galaxy properties from the ground with the minimum amount of confusion from the point spread function of the central quasar. Their sample was carefully constructed to ensure that the different types of object could be compared directly with one another: the radio-loud and -quiet quasars both have the same distribution of optical luminosities ($`23M_V26`$) and both the radio-loud quasars and the radio galaxies have similar extended radio luminosities and morphologies and steep radio spectra. The sample is therefore ideal for investigating the influence of galaxy properties on the ‘radio loudness’ of otherwise very similar quasars, and also for testing unified models of RLQs and RGs which predict that the properties of their hosts should be identical. A substantial subset of the sample has also formed the basis for an $`R`$-band imaging study using HST (Kukula et al. 1999, McLure et al. 1999, Dunlop et al. 2000), in which the enhanced angular resolution of HST has allowed both unambiguous identification of the host morphology and identification of detailed substructure in the quasar hosts and radio galaxies.
This sample provides an ideal starting point for a spectroscopic study of host galaxies because, in addition to the careful selection criteria, the existence of deep near-infrared images of every object provides us with a unique opportunity to minimise the contamination of the galaxy spectra by quasar light. Armed with knowledge of the extent and orientation of the host galaxy on the sky one is able to choose a slit position which is far enough from the nucleus to avoid the worst excesses of scattered quasar light, but which simultaneously maximises the amount of galaxy light falling onto the slit.
Out of the the 40 objects in the original Taylor et al. sample a total of 26 objects were observed in the current study (9 RQQs, 10 RLQs and 7 RGs). Details are listed in Table 1. Note that the radio source 3C59 has been shown by Meurs & Unger (1991) to consist of three separate sources, of which only the weakest appears to be associated with the quasar 0204+292. As a result, the radio luminosity of 0204+292 places it below our dividing line of $`L_{5GHz}=10^{24}`$W Hz<sup>-1</sup>sr<sup>-1</sup> and we classify it as an RQQ.
## 3 Observations
In order to asssess the feasibility of our observing strategy initial observations of 11 of the nearest and brightest objects in the sample were carried out using the Mayall 4-m Telescope at Kitt Peak National Observatory. This was followed by a larger programme of observations using the 4.2-m William Herschel Telescope (WHT), part of the Isaac Newton Group of telescopes on La Palma.
Tables 2 and 3 list the observations made on the Mayall 4-m Telescope and the WHT respectively. Six objects (0054$`+`$144, 0157$`+`$001, 0736$`+`$017, 2141$`+`$279, 2247$`+`$140 and 2344$`+`$184) were observed with both telescopes in order to provide a check for consistency between the two sets of observations. The spectra for these objects are discussed in more detail in Section 6.
### 3.1 The Mayall 4-m Telescope
Observations of 11 objects were carried out with the R. C. Spectrograph on the Mayall 4-m Telescope at Kitt Peak in 1992 September and 1993 March (Table 1). The long-slit spectrograph uses a Tektronix $`2048\times 2048`$ chip with 24-$`\mu `$m pixels, designated T2KB. The slit width was set to 3<sup>′′</sup> and the instrument was configured to give a spectral resolution of 1.9Å/pixel and a spatial resolution of 0.69<sup>′′</sup>/pixel.
The slit was first centred on the quasar (or, if the object was a radio galaxy, on the galaxy centroid), before being rotated and offset to the desired position, usually 5 arcseconds off-nucleus.
To allow the removal of cosmic rays five 1800-second off-nuclear exposures were obtained for each object, along with a shorter exposure of the quasar itself. The spectra typically spanned wavelengths from 3500Å to 6000Å, although the precise range varied from object to object according to redshift; details for individual objects are given in Table 2. Data reduction was carried out using standard iraf routines.
### 3.2 Observations with ISIS on the WHT
After the two runs at Kitt Peak had demonstrated that galaxy light could indeed be separated from that of the quasar and that useful spectra could be obtained, further observations were made of 22 objects (6 of which had already been observed at Kitt Peak) with the Intermediate Dispersion Spectroscopic and Imaging System (ISIS) at the 4.2 m William Herschel Telescope (WHT) on La Palma. ISIS uses a dichroic mirror to split the incoming light into red and blue beams which are then treated separately in spectrographs which have been optimised for the appropriate wavelength ranges ($`30006000`$Å for the blue arm and $`500010000`$Å for the red). This arrangement allows a larger wavelength coverage than is possible with the single spectrograph on the Mayall 4-m, enabling us to extend our spectra further into the red and thus giving greater scope for constraining models of spectrophotometric evolution.
The data were obtained in four separate observing runs in 1993 November, 1994 May, 1994 September and 1995 March. The instrumental set-up differed slightly between the runs, resulting in variations in the wavength ranges obtained. For the two sessions in 1994 the detectors in both the blue and red arms were Tektronix (‘Tek’) CCDs, with 24-$`\mu `$m pixels. In November 1993 and March 1996 the Tektronix chip in the red arm was replaced by an EEV P88300 chip, with 22.5-$`\mu `$m pixels, which allowed a slightly larger wavelength coverage. R158 gratings were used in each arm, giving a spectral resolution of 2.88Å pixel<sup>-1</sup> for the blue Tek chip and 2.90 (Tek) or 2.72Å pixel<sup>-1</sup> (EEV) for the red chips.
As with the observations at Kitt Peak the slit was first centred on the quasar nucleus, or the optical peak of the radio galaxy, and then offset to the desired position, 5<sup>′′</sup> from the quasar, and rotated to be at right angles to the direction of the offset. The slit width was set to 2<sup>′′</sup>, placing the inner edge of the slit at least 4<sup>′′</sup> from the quasar position.
Once again exposures were limited to 1800 seconds duration. Whenever possible we aimed to obtain five such frames per object, giving a total on-source exposure time of 2.5 hours. On-nuclear quasar spectra were also taken when time permitted. Data reduction was carried out using the figaro package, part of the Starlink suite of astronomical software.
#### 3.2.1 Calibration and splicing of red and blue spectra
The galaxies observed with the WHT tend to be fainter and/or at greater redshifts than those observed on the Mayall 4-m at Kitt Peak, and so additional care needed to be taken during the reduction of the WHT spectra. In particular, the procedure employed to optimize the extraction of the off-nuclear spectrum from the CCD frame means that the final flux calibration is only relative and not absolute - comparisons with the absolute fluxes obtained on the Mayall 4-m are therefore not meaningful.
The use of the red and blue arms of ISIS, whilst extending the wavelength coverage significantly, also introduces its own special problems. The large wavelength range made available by ISIS means that the spectra are affected by several prominent atmospheric emission features (Figure 1) which must be removed. Of these the strongest are the two oxygen lines at 5577 and 6300Å, the sodium D line at 5890Å, and the series of OH bands at wavelengths $`>6500`$Å. Sky lines were removed by fitting a third order polynomial to two empty strips of sky on either side of the target spectrum but the removal process sometimes left a residual imprint of these features and this constitutes a major source of noise in some of the fainter spectra in our sample. Where this is the case the affected regions are mentioned in the description of the individual spectrum.
A second problem involves the splicing together of the two spectra from the red and blue arms of the instrument to give continuous wavelength coverage from 3200Å to 9000Å.
The reflection and transmission responses of the ISIS dichroic cross at $`6100`$Å and the instrument settings were designed to ensure a large overlap between the wavelength coverages of the two arms so that the full reflection/transmission curves could be followed in the blue and the red spectra respectively. However, the faintness of many of the spectra effectively causes the flux calibration to fail at the red end of the blue spectrum and the blue end of the red, producing an artificial ‘lump’ in the spliced spectrum at the crossover point (6100Å or 6050Å, depending on the exact instrumental settings at the time of the observations - see column 8 in Table 3).
Accordingly, in plotting the WHT spectra in Figure 2 we have blanked out the data points for 100Å on either side of the join region and replaced them with an averaged bridging section. We have masked out this section of the data in all subsequent modelling due to its inherrent unreliability (see Nolan et al. 2000).
Due to the luminosity of the quasars themselves the on-nuclear spectra obtained with the WHT do not suffer from this problem and the red and blue sections have been spliced together directly with no noticeable discontinuity in flux at the join.
## 4 Results
The galaxy spectra themselves are displayed in Figure 2 along with nuclear spectra for the same object, where such are available. All the data have been smoothed to 10-Å bins, and the WHT data are displayed with the region linking the individual red and blue spectra blanked out as described above. The expected position of various stellar absorption features, including the 4000Å break are indicated by dotted lines in the off-nuclear spectra and the wavelengths of redshifted \[Oiii\]$`\lambda 5007`$ and H$`\alpha `$ are also marked.
### 4.1 Data quality
A cursory examination of Figure 2 is enough to show that the quality of the off-nuclear spectra varies considerably from object to object. In the two initial observing runs at Kitt Peak priority was given to nearby objects with bright, prominent host galaxies, and the eleven spectra obtained with the Mayall 4-m enjoy a high signal-to-noise ratio. By contrast, the objects observed in later runs on the WHT tend to have larger redshifts and the signal-to-noise in many of these spectra is correspondingly lower due to the rapid reduction in surface brightness and increase in galactocentric radius with $`z`$.
### 4.2 Degree of quasar contamination
The primary goal of the observing program - to obtain spectra from the stellar component of the host whilst avoiding scattered emission from the active nucleus - appears to have been satisfied to a large extent. This can be demonstrated most easily by comparing the off-nuclear (galaxy) and nuclear spectra for the same object in Figure 2. The quasar spectra show prominent broad lines, notably those of H$`\alpha \lambda 6563`$ and H$`\beta \lambda 4861`$ , with an increase in flux towards shorter wavelengths (particularly when the ‘blue bump’ continuum feature begins to emerge at $`\lambda _{rest}5000`$Å).
Whilst emission lines do occur in the off-nuclear spectra, they tend to be relatively narrow forbidden lines such as \[Oiii\]$`\lambda \lambda \lambda 4363,4959,5007`$. Where permitted lines occur, they lack the extremely broad profiles seen in the quasar nuclei, and are thus unlikely to result from scattering-induced contamination by nuclear light. In many cases the line ratios in the off-nuclear spectra also differ from those measured in the quasars, indicating that the emission arises under different conditions than those prevailing in the active nucleus.
However, in spectra obtained under conditions of poor seeing, there is likely to be some degree of nuclear contamination and this will be exacerbated if the object was observed far from the zenith, where differential atmospheric refraction may also lead to signifcantly more contamination at the blue end of our wavelength range than in the red. Although the extent of such contamination is difficult to measure directly, Tables 2 and 3 list the atmospheric seeing and airmass at the time each spectrum was obtained, to allow a rough assessment of the problem to be made.
### 4.3 Spectral features
The characteristic shape of the stellar continuum, including features such as the 4000Å break and various stellar absorption lines, is easily recognisable in most cases. The 4000Å break is particularly important for comparison between data and models of the stellar population. Since it covers a large wavelength interval it is relatively insensitive to the effects of instrument resolution and noise (Hamilton 1985). The break amplitude (defined as the ratio of the average flux density between rest-frame 4050 and 4250Å to that between 3750 and 3950Å) is therefore widely used as a tracer of spectral evolution (eg Bruzual 1983). The discontinuity results from the combined effect, shortwards of 4000Å, of lines of several elements heavier than helium, in a variety of ionization states, along with the crowding of higher order Balmer lines. If a significant population of massive young stars is present the enhanced degree of ionization causes the feature to weaken; it is most prominent in the spectrum of a well-established stellar population in which the most massive stars have had time to evolve away from the main sequence. Hence the 4000Å break is sensitive to both spectral type and metallicity, although if we assume a constant metallicity (a reasonable assumption for a particular galaxy at a fixed radius) it becomes a good indicator of the mean age of the local stellar population. The strength of this feature, as measured in each of the current spectra, is listed in Table 4. Other stellar absorption features, such as G band (4300 to 4320Å) and the Mg Ib (5173Å) and Fe$`\lambda 5270`$ lines, are also clearly present in many of the spectra.
Several of the off-nuclear spectra (eg 0054+144, 2201+315), despite showing a clear break at 4000Å and a continuum longwards of this wavelength which can be fitted extremely well by a passively ageing stellar population (Nolan et al. 2000), also display a contribution from a component which rises steeply towards the blue. The slope of this feature closely resembles that of a quasar SED, and its presence is often (but not always) accompanied by emission lines characteristic of quasar nuclei, suggesting that it is in fact scattered light from the quasar itself, the result either of atmospheric scattering or (since the feature is not always correlated with poor observing conditions) of scattering within the ISM of the host galaxy. Another possibility is that the blue excess indicates the presence of a substantial population of young stars within the host galaxy. If this latter case were true then it would pose a serious problem for the unification of RLQs and RGs since none of the radio galaxies display such a component. The issue of excess blue continuum is raised on a case-by-case basis in the following section, but a full discussion in the light of detailed stellar population synthesis modelling is deferred to the companion paper by Nolan et al. (2000).
The average values of the 4000Å break strength for the three types of object in the sample are 1.4 (RQQ hosts), 1.5 (RLQ hosts) and 1.7 (RGs). These are all somewhat lower than the value of $`2.0`$ measured for local inactive elliptical galaxies by Hamilton (1985), but we note that there is a wide scatter in our sample and that the lowest values are all associated either with spectra in which the signal to noise is particularly poor (eg 0244+292, 1004+130, 1334+008) or those which clearly show an extra source of continuum emission at short wavelengths (eg the quasars 0054+144, 1217+023, 2141+175). The cleanest spectra generally have break strengths which are consistent with those seen in ‘normal’ well-established galaxies.
As a final caveat we note that these spectra tell us only about the stellar composition of the region of the galaxy covered by the slit - we cannot, for example, rule out the presence of a significantly different stellar population closer to the nucleus, or concentrated in clumps which the slit happens to avoid. However, the slit has a width of at least 2 arcsec and, as can be seen from the contour plots in Figure 2, its length cuts across a significant fraction of the galaxy in the transverse direction. The area covered often amounts to several square arcseconds (equivalent to several tens of square kiloparsecs at typical redshifts) and therefore represents a good general sample of the outer regions of the host.
## 5 Individual objects
Objects are listed under their IAU names (alternative names are listed in Table 1), in order of increasing right ascension. The classification of each object as either an RQQ, an RLQ or a radio galaxy is indicated in parentheses, along with the name(s) of the telescope(s) on which off-nuclear spectra were obtained. The form of the spectrum is described, noting any peculiar features as well as the presence or otherwise of a 4000Å break at the expected observed wavelength ($`\lambda _{obs}`$). We also note the morphology of the galaxy (disc or elliptical) based on its surface brightness profile in the $`K`$ or $`R`$-band continuum images by Taylor et al. (1996) ($`K`$-band; UKIRT) or McLure et al. (1999) and Dunlop et al. (2000) ($`R`$-band; HST). For a more detailed description of previous imaging studies of the host galaxies see Dunlop et al. (1993) (RLQs and RQQs) or Taylor et al. (1996) (RGs).
0007$`+`$106 (RQQ; M4M): the nuclear spectrum of this radio-quiet quasar shows prominent broad H$`\alpha \lambda 6583`$ and H$`\beta \lambda 4861`$ emission as well as narrower forbidden line emission from species including \[0iii\] and \[Fevii\]. The off-nuclear spectrum of the host galaxy has a high signal-to-noise and displays little sign of contamination from the quasar: the contribution from emission lines is very small, and the 4000Å break is clearly visible (redshifted to 4356Å) despite the rapid increase in quasar flux towards the blue end of the spectrum, which would tend to mask the break if scattering were significant. G band and Mg ib absorption are also present. We note that the slit crosses the optical arc-like structure to the north of the quasar which Hutchings et al. (1984) suggest may be a spiral arm. Previous spectroscopy of this region by Green, Williams & Morton (1978) showed narrow emission lines (with different line ratios from those in the nucleus) and a red continuum which they attribute to starlight. However, this region only constitutes a small fraction of the galaxy light intercepted by the slit in the current observations. Taylor et al. (1996) find that an exponential disc profile provides a good fit to the NIR surface brightness distribution of the galaxy. (Morphology: disc; Taylor et al. 1996.)
0054$`+`$144 (RQQ; M4M & WHT): quasar continuum emission dominates the nuclear spectrum of this RQQ, although H$`\beta \lambda 4861`$ and \[Oiii\]$`\lambda \lambda \lambda 4363,4959,5007`$ lines are visible. These lines are not prominent in either the Mayall 4-m or WHT off-nuclear spectra of the host (particularly the latter spectrum) but bluewards of the (relatively weak) 4000Å break (at $`\lambda _{obs}=4684`$Å) the galaxy spectrum displays a marked increase in flux, very similar in form to that displayed by the quasar itself. This component is seen in both off-nuclear spectra, which were taken at different times, under different seeing conditions, and used different slit positions, so it is not clear whether we are seeing quasar light which is being scattered into our line of sight either by the atmosphere or by the interstellar medium of the host galaxy, or whether the blue excess is due to a population of young stars. G band, Mg ib and Fe5270 absorption features appear in the Mayall 4-m spectrum. McLure et al. (1999) classify the galaxy as an elliptical: its light profile is very well described by an $`r^{1/4}`$ law and the galaxy itself is quite red ($`RK=3.14`$), but a tidal interaction with a nearby companion is suggested by the extension to the NW of the nucleus (Dunlop et al. 1993). (Morphology: elliptical; McLure et al. 1999.)
0137$`+`$012 (RLQ; M4M): the 4000Å break in this object is quite clear (at $`\lambda _{obs}=5032`$Å) and there is little evidence of nuclear emission lines. (Morphology: elliptical; McLure et al. 1999.)
0157$`+`$001 (RQQ; M4M & WHT): the off-nuclear spectrum obtained with the Mayall 4-m does not appear to be strongly contaminated by emission from the nucleus and the 4000Å break is visible at $`\lambda _{obs}=4656`$Å. The G band absorption feature is also present. However, in the spectrum taken with the WHT using a similar slit position, there appears to be a significant contribution from the quasar continuum bluewards of the break. The signal:noise ratio in the WHT spectrum is much reduced longwards of 7500Å due to residuals from the subtracted OH bands. In both cases the slit intercepts the prominent tidal arm which extends north and NW of the nucleus (MacKenty & Stockton 1984) and is known to contain several emission line regions (Stockton & MacKenty 1987), although emission lines are not strongly evident in the integrated spectra presented here. Previous spectroscopy of this structure by Heckman et al. (1984) showed a velocity difference of $`300\pm 200`$ km s<sup>-1</sup> between the arm and the quasar itself. McLure et al. (1999) find that the underlying smooth $`R`$-band continuum light is well described by an $`r^{1/4}`$-law, though conceivably this might be another result of the tidal interaction. (Morphology: elliptical; McLure et al. 1999.)
0204$`+`$292 (RQQ; WHT): the 4000Å break ($`\lambda _{obs}=4436`$Å) is quite strong in this WHT spectrum, but residuals from oxygen and OH features from the sky spectrum are also present. (Morphology: disc; Taylor et al. 1996)
0230$``$027 (RG; WHT): the off-nuclear spectrum of this small, faint radio galaxy suffers from low signal to noise and residual sky features are visible. However a break can be seen at $`\lambda _{obs}=4952`$Å. The nuclear spectrum of this object is also dominated by starlight, although narrow emission lines are present. (Morphology: elliptical; Dunlop et al. 2000.)
0244$`+`$194 (RQQ; WHT): the signal to noise in this off-nuclear spectrum is very low. Apart from residual sky features there is some evidence for a drop in the continuum level around $`\lambda _{obs}=4704`$Å (the expected position of the 4000Å break), but also for a blue component shortwards of this, perhaps indicative of nuclear contamination. (Morphology: elliptical; McLure et al. 1999.)
0345$`+`$337 (RG; WHT): the spectrum is noisy, but a break feature appears to be present at $`\lambda _{obs}=4936`$Å. The WHT slit intercepts a bright knot NW of the quasar which Taylor et al. (1996) suggest may be an embedded companion galaxy. (Morphology: elliptical; McLure et al. 1999.)
0736$`+`$017 (RLQ; M4M & WHT): the WHT spectrum, though noisier and suffering from OH-band residuals at long wavelengths, agrees well with the spectrum obtained previously at Kitt Peak. The 4000Å break is quite clear at $`\lambda _{obs}=4764`$Å and weak \[Oiii\], H$`\alpha `$ and H$`\beta `$ features can also be seen, though they appear to be narrower than those in the nuclear spectrum. The galaxy itself is highly disturbed, but McLure et al. (1999) fit an $`r^{1/4}`$-law profile to the smooth component of the $`R`$-band continuum. (Morphology: elliptical; McLure et al. 1999.)
0917$`+`$459 (RG; WHT): the off-nuclear spectrum is pure stellar continuum, with a strong break at $`\lambda _{obs}=4696`$Å. The galaxy isophotes are complex but McLure et al. (1999) find the underlying distribution to be well described by an $`r^{1/4}`$-law. (Morphology: elliptical; McLure et al. 1999.)
1004$`+`$130 (RLQ; WHT): OH-band residuals are the only prominent feature of this off-nuclear spectrum, with little evidence for either emission lines or a 4000Å break (at $`\lambda _{obs}=4960`$Å). The absence of a strong break may reflect the fact that the elliptical host galaxy of this quasar is known to possess unusual ‘spiral’ features close to the nucleus (McLure et al. 1999), perhaps indicating the presence of a significant population of young stars. Stockton & MacKenty (1987) note that there is no significant extended \[Oiii\] emission in this object. (Morphology: elliptical; McLure et al. 1999.)
1020$``$103 (RLQ; M4M): despite the presence of stellar absorption features such as G band and a 4000Å break at $`\lambda _{obs}=4788`$Å, the off-nuclear spectrum also contains many emission lines as well as a blue excess, which may indicate a significant contribution from scattered quasar light (the seeing was quite poor for much of the observations). Dunlop et al. (1993) note that the quasar is off-centre and that the galaxy isophotes appear to be swept back towards the SW, providing evidence for disturbance in this object. (Morphology: elliptical; Dunlop et al. 1999.)
1215$``$033 (RG; WHT): this spectrum displays only a weak break at $`\lambda _{obs}=4736`$Å . OH-band residuals dominate longwards of 7500Å. (Morphology: elliptical; Dunlop et al. 2000.)
1217$`+`$023 (RLQ; WHT): a weak break feature is present at $`\lambda _{obs}=4960`$Å but shortwards of this the spectrum is dominated by a component which rises towards the blue possibly indicating scattered nuclear continuum (there is however little evidence for accompaning nuclear line emission, and the seeing during the observations was excellent). At the red end, poor subtraction of OH bands has reduced the signal to noise ratio of the spectrum. (Morphology: elliptical; Dunlop et al. 2000.)
1330$`+`$022 (RG; M4M): weak \[Oiii\] lines occur in the off-nuclear spectrum of this radio galaxy and the stellar continuum shows G band absorption and a strong 4000Å break feature at $`\lambda _{obs}=4860`$Å. By contrast, the nuclear spectrum shows evidence for broad H$`\beta \lambda 4861`$ and the less prominent break may indicate a contribution from a quasar-type continuum or perhaps a nuclear starburst region. (Morphology: elliptical; Dunlop et al. 2000.)
1334$`+`$008 (RG; WHT): a very noisy spectrum obtained under poor conditions, the underlying continuum is confused by many residual sky features. The apparent increase in flux shortwards of 4500Å almost certainly reflects a failure of the flux calibration at very low light levels. However, there is some evidence for a break at the expected wavelength of $`\lambda _{obs}=5196`$Å. The slit intercepts one of the secondary nuclei reported by Taylor et al. (1996). (Morphology: elliptical; Taylor et al. 1996.)
1549$`+`$203 (RQQ; WHT): the extended nebulosity around this RQQ, though confused by a foreground galaxy cluster, shows a strong break feature at the expected wavelength of $`\lambda _{obs}=5000`$Å and little evidence for emission lines. OH-band residuals add to the noise levels at long wavelengths. (Morphology: elliptical; Dunlop et al. 2000.)
1635$`+`$119 (RQQ; WHT): the 4000Å break is clearly visible at $`\lambda _{obs}=4584`$Å, but residual sky features degrade the quality of the spectrum towards the red end. (Morphology: elliptical; McLure et al. 1999.)
2135$``$147 (RLQ; WHT): generally low signal to noise with prominent residuals due to sky features. However, there is weak evidence for a break at $`\lambda _{obs}=4800`$Å. The galaxy appears to be disturbed and possesses a secondary nucleus to the SE of the quasar (Stockton 1982) which may also be active (Hickson & Hutchings 1987). The WHT slit intercepts a region to the SW of the quasar at which Stockton & MacKenty (1987) report extended \[Oiii\] emission. Narrow H$`\alpha `$ is present in the current spectrum but the \[Oiii\] lines fall within the join region where the signal is unreliable. (Morphology: elliptical; Dunlop et al. 2000.)
2141$`+`$175 (RLQ; WHT): low signal to noise and the presence of a strong blue component may serve to mask any evidence of a 4000Å break in this object (expected at $`\lambda _{obs}=4852`$Å). The idea that this blue component originates as scattered quasar light is bolstered by the possible presence of the H$`\alpha `$ line, but the issue is confused by strong sky residuals (although the seeing was good, the airmass during the observations was relatively high, so differential refraction might explain the presence of quasar contamination at shorter wavelengths without requiring the presence of a strong H$`\alpha `$ line). The elongated appearance of the galaxy is apparently the result of an edge-on tidal arm consisting of old stars (Stockton & Farnham 1991). However, the WHT slit crosses the galaxy to the NE of the quasar, where the starlight appears to follow a bulge-dominated $`r^{1/4}`$-law (McLure et al. 1999a). (Morphology: elliptical; McLure et al. 1999.)
2141$`+`$279 (RG; M4M & WHT): the off-nuclear spectrum taken at Kitt Peak shows a weak \[Oiii\]$`\lambda 5007`$ line, a break feature at $`\lambda _{obs}=4860`$Åand Mg ib absorption. The WHT spectrum is much noisier, but generally consistent with the features in the earlier data. The nuclear spectrum of this radio galaxy is also dominated by starlight, although prominent narrow lines are present. (Morphology: elliptical; McLure et al. 1999.)
2201$`+`$315 (RLQ; M4M): this object shows prominent G band absorption but only a weak 4000Å break at $`\lambda _{obs}=5192`$Å, along with low-equivalent-width \[Oiii\] lines. A blue component shortwards of the break is consistent with scattering of the nuclear quasar continuum. (Morphology: ambiguous; Taylor et al. 1996.)
2215$``$037 (RQQ; WHT): the spectrum is relatively noisy with prominent residuals from all the bright sky features. The apparent ‘hump’ at $`6000`$Å is an artifact that results from the sodium D sky line coinciding with the beginning of the masking region that links the red and blue halves of the the spectra from ISIS. There is little evidence for a prominent break at the expected wavelength of $`\lambda _{obs}=4964`$Å. Both Hutchings et al. (1989) and McLure et al. (1999b) classify the galaxy as a non-interacting elliptical system. (Morphology: elliptical; Dunlop et al. 2000.)
2247$`+`$140 (RLQ; M4M & WHT): the galaxy’s 4000Å break is visible at $`\lambda _{obs}=4948`$Å in the Mayall 4-m spectrum, along with G band absorption. The data obtained on the WHT are consistent with this, although residual sky features from oxygen and sodium D lines and the OH bands are also present. (Morphology: elliptical; McLure et al. 1999.)
2344$`+`$184 (RQQ; M4M & WHT): both the Mayall 4-m and WHT spectra show a break at $`\lambda _{obs}=4552`$Å as well as G band absorption and are generally consistent with one another despite their differing slit positions. The nuclear spectrum of this object is also dominated by starlight; 2344$`+`$184 is one of the least powerful quasars in the current study and technically qualifies as a Type 1 Seyfert galaxy. The surface brightness profile is that of a disc galaxy, although a bulge dominates in the central regions (McLure et al. 1999) and there is also evidence for a bar (Hutchings, Janson & Neff 1989). (Morphology: disc; McLure et al. 1999.)
2349$``$014 (RLQ; WHT): there is a weak break feature at $`\lambda _{obs}=4692`$Å and also evidence for weak H$`\beta `$, \[Oiii\]$`\lambda 5007`$ and H$`\alpha `$ lines. The red end of the spectrum is dominated by residual atmospheric OH band emission. (Morphology: elliptical; McLure et al. 1999.)
## 6 Comparison of Mayall 4-m and WHT results
Six objects in the current sample (the RQQs 0054$`+`$144, 0157$`+`$001 and 2344$`+`$184, the RLQs 0736$`+`$017 and 2247$`+`$140, and the radio galaxy 2141$`+`$279) were observed with both the Mayall 4-m Telescope and the WHT. This duplication allows us to check for systematic differences between the spectra obtained with each instrument and to this end the are replotted one above the other in Figure 3. Since the observations were often made under different atmospheric conditions they also provide a means of assessing the degree to which any contamination by nuclear light can be attributed to either the airmass and seeing conditions at the time of the observations or scattering within the host galaxy itself. Where a different slit position was used, the two spectra allow us to examine the degree of homogeneity in the stellar composition of the galaxy.
It should be noted that due to the optimization method used to extract the WHT spectra, their flux calibration can only be considered as relative, not absolute. Different slit widths were also used on the two instruments. Comparisons of the flux densities obtained on the two telescopes are therefore not meaningful.
0054$`+`$144 (RQQ): different slit positions were used for the Mayall 4-m and WHT observations. The spectrum obtained with the Mayall 4-m shows a contribution from H$`\beta `$/\[Oiii\] at $`\lambda _{obs}5800`$Å, which is lacking in the WHT data. However, the underlying continuum is very similar in both spectra and the measured depth of the (weak) 4000Å break is very similar in each. This implies that the emission lines detected in the Mayall 4-m spectrum are a local feature, produced in situ rather than being due to scattered light from emission-line regions in the nucleus, since they do not appear to be accompanied by a corresponding amount of scattered nuclear continuum. The consistency in the strength of the 4000Å break at the two observing epochs, along with its relative weakness, suggests the presence of a significant population of young stars in the host galaxy. McLure et al. (1999) report a tidal feature visible in their $`R`$-band HST image of the elliptical host galaxy which may be linked to the origin of this young stellar population.
0157$`+`$001 (RQQ): the slit positions used in the two sets of observations were essentially the same. However, the WHT spectrum shows excess blue continuum shortwards of this feature. The two spectra were taken under what were ostensibly very similar atmospheric conditions, but we note that a very good approximation of the WHT spectrum can be obtained simply by adding a scaled version of the nuclear spectrum to the off-nuclear Mayall 4-m data. This suggests that atmospheric scattering is to blame for the blue excess in the WHT spectra.
0736$`+`$017 (RLQ): the two spectra were both taken under good atmospheric conditions and using the same slit position. The agreement between the two is excellent.
2141$`+`$279 (RG): the WHT spectrum has a much lower signal:noise ratio than the Mayall 4-m spectrum, and uses a different slit position. However, the overall continuum shape is in good agreement with the earlier data. Both slits intercept an extension to the NE of the nucleus which may be a tidal feature caused by interaction with a northern companion.
2247$`+`$140 (RLQ): the same slit position was used for both spectra, although the seeing during the WHT run was quite poor ($`1.6^{\prime \prime }`$). The two datasets are consistent with one another, with a relatively strong 4000Å break and little sign of an additional blue continuum component from either young stars or scattered quasar light.
2344$`+`$184 (RQQ): different slit positions were used on the two different instruments, and the WHT spectrum shows an excess of blue emission and a correspondingly weaker 4000Å break. The seeing conditions were actually worse during the M4M observations, making atmospheric scattering of quasar light an unlikely culprit. Moreover, we note that the nuclear spectrum of this low-luminosity quasar is not particularly blue and itself shows prominent stellar continuum features (see Figure 2). More likely is that the WHT slit crossed a region of the galaxy containing a large number of young stars. The host of 2344$`+`$184 is a disc galaxy, with a central bulge and prominent spiral arms (McLure et al. 1999). (Hutchings et al. (1989) also suggest the presence of a bar.)
## 7 Summary
We describe and present optical spectra of 26 galaxies hosting powerful nuclear activity. The sample contains carefully matched subsamples of all three types of powerful AGN - radio-quiet and radio-loud quasars, and FRii radio galaxies - enabling us to investigate the relationship between the host galaxy and the radio properties of the resident AGN and also to test unified models for radio-loud objects.
The spectra were taken 5 arcsec off-nucleus and, via a careful choice of slit position, aim to maximise the amount of galaxy light entering the instrument whilst avoiding contamination from the active nucleus. In the majority of cases this approach appears to have been successful; the continuum is clearly overwhelmingly stellar in origin and even when the presence of scattered nuclear emission is suspected, features such as the 4000Å break are still clearly discernable. In almost all objects we detect at least some stellar signatures.
A second paper (Nolan et al. 2000) will discuss spectrophotometric modelling of the galaxy spectra in order to determine the ages and starformation histories of their constituent stellar populations.
## Acknowledgments
The authors would like to thank the staff at KPNO and the ING for their assistance during the observations, and the anonymous referee whose comments and suggestions resulted in several improvements in the final paper. DHH and MJK acknowledge PPARC support. This research has made use of the NASA/IPAC Extragalactic Database (NED), which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration.
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# Casimir Dispersion Forces and Orientational Pairwise Additivity
## I Introduction and summary
External objects that are immersed in a fluctuating medium, and modify the fluctuations in their vicinity, experience induced interactions with one another . These interactions are most often independent of the structural details, and are in turn highly sensitive to the geometry of the objects and their mutual arrangements while immersed in the medium.
The strong dependence of these interactions on the shape of the objects raise the issue of pairwise additivity: Is it possible to express the fluctuation–induced interaction between two extended bodies as the sum of a pair potential, or the interaction between several bodies as the sum of two-body interactions?
It is well known that a pairwise summation of the van der Waals interaction gives the correct power law for the Casimir energy . Let us take a pair potential of the form $`A/r^n`$, with $`n=6`$ for the thermal case and $`n=7`$ for the quantum case , and $`A`$ being a constant to be determined. If one tries to fix the coefficient by summing the pair potential over two bodies and equating the result to the expression for the Casimir interaction between the bodies, one finds out that a different coefficient is needed for every geometry.
To understand this, one should note that the van der Waals interaction is due to dipolar fluctuations. When two extended bodies are at a close separation, one can show that the fluctuations of all the multipoles in fact contribute comparably to the Casimir energy, and thus summation of the contribution due to the dipolar fluctuations cannot by itself account for the interaction . When the bodies are at large separations (larger than their typical sizes), the contribution due to higher multipoles are in fact systematically weaker. However, there is still a discrepancy between the sum of a (van der Waals) pair potential, and the contribution of the dipolar fluctuations to the Casimir energy. In the spirit of a (second order) perturbation theory, the correct way of calculating the dipolar Casimir energy is to consider the pairwise sum of the dipole-dipole interactions over the two bodies, and then square it and take the average. This is clearly in contrast with the pairwise summation of the van der Waals interaction, which corresponds to taking the square of the local dipolar fluctuations and averaging, and then summing over the two bodies.
The same picture can help us answer the second question. Many-body interactions can be expressed as the sum of many-body interactions of the multipoles of different bodies in the medium. When extended bodies are at close separations, and all the multipoles have comparable contributions, many-body interactions of nontrivial forms result . On the other hand, for bodies at large separations the leading order contribution comes from the sum of two-body interactions of the lowest nonvanishing multipole .
In this article, we study the issue of orientational pairwise additivity , which is to determine whether the orientational dependence of the interactions could be obtained from the summation of a pair potential. A path integral formulation is used to study the fluctuation–induced interactions between manifolds of arbitrary shape at large separations, in the context of a multipole expansion. It is shown that the form of the interaction crucially depends on whether the manifolds are grounded or isolated in an electrostatic analogy. In the grounded case, the manifolds are connected to a charge reservoir to maintain a constant potential, and thus the leading fluctuations are monopolar. Isolated manifolds, however, are constrained to have fixed overall charges, and can only undergo dipolar fluctuations. The leading interaction between grounded manifolds is found to be of the form $`(\mathrm{monopole}\mathrm{monopole})^2`$, and is independent of their shapes and orientations. The leading shape dependent term comes from the $`(\mathrm{monopole}\mathrm{dipole})^2`$ term, which gives rise to orientational dependencies that are pairwise additive. The interaction between isolated manifolds, however, is dominated by the $`(\mathrm{dipole}\mathrm{dipole})^2`$ term to the leading order, which is not pairwise additive.
The rest of the paper is organized as follows. In Sec. II, the path integral formulation is developed and general expressions are derived for the fluctuation–induced interactions for different types of boundary conditions. In Sec. III, the interactions are examined for the specific examples of symmetric objects such as spheres, and also highly asymmetric objects such as rods and disks, where the above features can be manifestly understood. Critical fluids are examined in Sec. IV, as a special case, and a conclusion follows in Sec. V.
## II Path Integral Formulation
Consider a $`d`$-dimensional medium, in which a field $`\varphi `$ is undergoing scale-free (massless) thermal fluctuations, described by the Hamiltonian
$$[\varphi ]=\frac{K}{2}d^d𝐱(\varphi )^2.$$
(1)
The field could represent a component of the electromagnetic field in a dielectric medium , the electrostatic potential in charged fluids at very low salt concentrations , an order parameter field for a critical binary mixture or a magnetic system , or a massless Goldstone mode arising from a continuous symmetry breaking . For simplicity, in what follows we are going to think in the context of charged fluids, and later on comment on the specific effects in the case of critical fluids.
Let us assume that there are $`n`$ manifolds immersed in the medium, denoted by $`M_\alpha `$ ($`\alpha =1,\mathrm{},n`$), which modify the fluctuations. In the electrostatic context, one can view each manifold as a conductor that requires a constant value for the potential field in the whole volume that it encloses. A restricted partition function, which requires a value of $`\varphi _\alpha `$ for the potential field on the $`\alpha `$th manifold, can then be written as
$$𝒵[\varphi _\alpha ]=𝒟\varphi (𝐱)\underset{\alpha =1}{\overset{n}{}}\delta \left\{\varphi |_{M_\alpha }\varphi _\alpha \right\}\mathrm{e}^{[\varphi ]}.$$
(2)
Following Ref. , the functional delta functions can next be represented by introducing the Lagrange multiplier fields $`\rho _\alpha (𝐱)`$, as
$`𝒵[\varphi _\alpha ]`$ $`=`$ $`{\displaystyle 𝒟\varphi (𝐱)\underset{\alpha =1}{\overset{n}{}}_{M_\alpha }𝒟\rho _\alpha (𝐱)\mathrm{exp}\left\{\frac{K}{2}d^d𝐱(\varphi )^2+i\underset{\alpha }{}d^d𝐱\rho _\alpha (𝐱)\left[\varphi (𝐱)\varphi _\alpha \right]\right\}}`$ (3)
$`=`$ $`𝒵_0\times {\displaystyle \underset{\alpha =1}{\overset{n}{}}}{\displaystyle _{M_\alpha }}𝒟\rho _\alpha (𝐱)\mathrm{exp}\left\{{\displaystyle \frac{1}{2K}}{\displaystyle \underset{\alpha ,\beta }{}}{\displaystyle d^d𝐱d^d𝐱^{}\rho _\alpha (𝐱)G(𝐱𝐱^{})\rho _\beta (𝐱^{})}i{\displaystyle \underset{\alpha }{}}\varphi _\alpha {\displaystyle d^d𝐱\rho _\alpha (𝐱)}\right\},`$ (4)
in which
$$G(𝐱𝐱^{})=(^2)_{𝐱,𝐱^{}}^1=\frac{1}{S_d(d2)|𝐱𝐱^{}|^{d2}},$$
(5)
with $`S_d=\frac{2\pi ^{d/2}}{\mathrm{\Gamma }(d/2)}`$ (the surface area of the $`d`$-dimensional sphere), $`𝒵_0`$ is the free partition function, and $`_{M_\alpha }𝒟\rho _\alpha (𝐱)`$ implies a functional integration only in the region enclosed by $`M_\alpha `$. (In other words, the Lagrange multiplier field $`\rho _\alpha (𝐱)`$ is nonzero only within the volume of $`M_\alpha `$.) Note that one should view the $`\rho _\alpha (𝐱)`$ fields as fluctuating charge density fields, and Eq.(3) as the partition function of a set of interacting Coulomb plasmas, in the electrostatic context .
The fluctuation–induced interactions between the conductors can now be inferred from the above partition function. However, it is important to specify the boundary conditions for the conductors. One possibility is that the conductors are grounded, that is to say they are maintained at a constant fixed potential by being in contact with a large reservoir of charges; a so-called “ground.” In this case, the free energy of the system is obtained as
$$F_{\mathrm{gr}}=k_\mathrm{B}T\mathrm{ln}𝒵[\varphi _\alpha =0].$$
(6)
The other possibility is that the conductors are made isolated, and maintain constant amounts of net charges, which we assume to be zero. In this case, the potential field at the conductors can take any value to help maintain the neutrality, and thus the free energy is obtained as
$$F_{\mathrm{is}}=k_\mathrm{B}T\mathrm{ln}\left(_{\mathrm{}}^+\mathrm{}\underset{\alpha }{}d\varphi _\alpha 𝒵[\varphi _\alpha ]\right).$$
(7)
To further proceed, we focus on the situation in which the manifolds are far from each other, namely, they are at separations much larger than their typical sizes. In this case, we can perform a multipole expansion for the charge density distribution. For example, the Coulomb interaction between the $`\alpha `$th and the $`\beta `$th conductors can be written as ($`\alpha \beta `$)
$`h_{\alpha \beta }`$ $`=`$ $`{\displaystyle d^d𝐱d^d𝐱^{}\frac{\rho _\alpha (𝐱)\rho _\beta (𝐱^{})}{S_d(d2)|𝐱𝐱^{}+𝐑_{\alpha \beta }|^{d2}}},`$ (8)
$`=`$ $`{\displaystyle \frac{Q_\alpha Q_\beta }{S_d(d2)R_{\alpha \beta }^{d2}}}\left({\displaystyle \frac{Q_\beta 𝐏_\alpha \widehat{𝐑}_{\alpha \beta }Q_\alpha 𝐏_\beta \widehat{𝐑}_{\alpha \beta }}{S_dR_{\alpha \beta }^{d1}}}\right)+\left({\displaystyle \frac{𝐏_\alpha 𝐏_\beta d𝐏_\alpha \widehat{𝐑}_{\alpha \beta }𝐏_\beta \widehat{𝐑}_{\alpha \beta }}{S_dR_{\alpha \beta }^d}}\right)+\mathrm{},`$ (9)
in which $`R_{\alpha \beta }`$ is the distance between the two conductors, and the multipoles are defined as
$`Q_\alpha `$ $`=`$ $`{\displaystyle d^d𝐱\rho _\alpha (𝐱)}=\stackrel{~}{\rho }_\alpha (𝐤)|_{𝐤=\mathrm{𝟎}},`$ (10)
$`P_{\alpha ,i}`$ $`=`$ $`{\displaystyle d^d𝐱x_i\rho _\alpha (𝐱)}={\displaystyle \frac{1}{i}}{\displaystyle \frac{}{k_i}}\stackrel{~}{\rho }_\alpha (𝐤)|_{𝐤=\mathrm{𝟎}},`$ (11)
with
$$\stackrel{~}{\rho }_\alpha (𝐤)=d^d𝐱\rho _\alpha (𝐱)\mathrm{e}^{i𝐤𝐱}.$$
(12)
We also need to make a similar multipole expansion for the self energy terms at each manifold
$`h_{\alpha \alpha }`$ $`=`$ $`{\displaystyle d^d𝐱d^d𝐱^{}\frac{\rho _\alpha (𝐱)\rho _\alpha (𝐱^{})}{S_d(d2)|𝐱𝐱^{}|^{d2}}},`$ (13)
$`=`$ $`{\displaystyle \frac{d^d𝐤}{(2\pi )^d}\frac{1}{k^2}\stackrel{~}{\rho }_\alpha (𝐤)\stackrel{~}{\rho }_\alpha (𝐤)}.`$ (14)
We can introduce the multipoles, using the Taylor expansion of the charge density in Fourier space
$`\stackrel{~}{\rho }_\alpha (𝐤)`$ $`=`$ $`\stackrel{~}{\rho }_\alpha (\mathrm{𝟎})+{\displaystyle \frac{1}{i}}{\displaystyle \frac{\stackrel{~}{\rho }_\alpha (𝐤)}{k_i}}|_{𝐤=\mathrm{𝟎}}ik_i+\mathrm{},`$ (15)
$`=`$ $`Q_\alpha +P_{\alpha ,i}\times ik_i+\mathrm{},`$ (16)
The above expansion can be formally viewed as an expansion in powers of $`kL_\alpha `$, where $`L_\alpha `$ is a typical size of the manifold. The expansion is thus convergent only for sufficiently small values of $`k`$, corresponding to length scales larger than the size of the manifolds. Since the self energy integral in Eq.(14) involves contributions from higher wavevectors, a multipole expansion for the self energy will be divergent. However, the expansion in Eq.(16) indicates that all the information concerning the first few multipoles of the charge distribution is already contained in the low $`k`$ behavior of the function $`\stackrel{~}{\rho }_\alpha (𝐤)`$. Since we are only interested in the dependence of the partition function Eq.(3) on the distances $`R_{\alpha \beta }`$, all we need to know about the self energy is its dependence on the first few multipoles, which is in fact well behaved.
Let us denote the domain of convergence for the expansion in Eq.(16) in $`k`$-space by $`𝒟_\alpha `$. This domain contains the origin, and its shape is determined by the geometry of the conductor. Loosely speaking, its size in each direction is set by the inverse of the size of the conductor in that direction. Now we can restrict the $`k`$-integral in the self energy only to this domain, and neglect the contribution from the outside of $`𝒟_\alpha `$, because all the dependence on the first few multipoles is included in the domain $`𝒟_\alpha `$. We thus have
$`h_{\alpha \alpha }`$ $`=`$ $`{\displaystyle _{𝒟_\alpha }}{\displaystyle \frac{d^d𝐤}{(2\pi )^d}}{\displaystyle \frac{1}{k^2}}\stackrel{~}{\rho }_\alpha (𝐤)\stackrel{~}{\rho }_\alpha (𝐤)+\mathrm{},`$ (17)
$`=`$ $`\gamma _\alpha Q_\alpha ^2+\gamma _{\alpha ,ij}P_{\alpha ,i}P_{\alpha ,j}+\mathrm{},`$ (18)
in which
$`\gamma _\alpha `$ $`=`$ $`{\displaystyle _{𝒟_\alpha }}{\displaystyle \frac{d^d𝐤}{(2\pi )^d}}{\displaystyle \frac{1}{k^2}},`$ (19)
$`\gamma _{\alpha ,ij}`$ $`=`$ $`{\displaystyle _{𝒟_\alpha }}{\displaystyle \frac{d^d𝐤}{(2\pi )^d}}{\displaystyle \frac{k_ik_j}{k^2}},`$ (20)
and so forth.
Putting all the pieces together, the $`R_{\alpha \beta }`$-dependent part of the partition function can be written as
$`𝒵[\varphi _\alpha ]`$ $`=`$ $`{\displaystyle }{\displaystyle \underset{\alpha }{}}dQ_\alpha d𝐏_\alpha \mathrm{}\mathrm{e}^{i_\alpha \varphi _\alpha Q_\alpha }\times \mathrm{exp}\{{\displaystyle \frac{1}{2K}}{\displaystyle \underset{\alpha }{}}[\gamma _\alpha Q_\alpha ^2+\gamma _{\alpha ,ij}P_{\alpha ,i}P_{\alpha ,j}+\mathrm{}]`$ (22)
$`{\displaystyle \frac{1}{2K}}{\displaystyle \underset{\alpha \beta }{}}[{\displaystyle \frac{Q_\alpha Q_\beta }{S_d(d2)R_{\alpha \beta }^{d2}}}{\displaystyle \frac{\widehat{R}_{\alpha \beta ,i}(Q_\beta P_{\alpha ,i}Q_\alpha P_{\beta ,i})}{S_dR_{\alpha \beta }^{d1}}}+{\displaystyle \frac{(\delta _{ij}d\widehat{R}_{\alpha \beta ,i}\widehat{R}_{\alpha \beta ,j})P_{\alpha ,i}P_{\beta ,j}}{S_dR_{\alpha \beta }^d}}+\mathrm{}]\}.`$
Note that we have neglected a Jacobian in changing the measure of integration. However, since the transformation from the charge density distribution to the multipole description is linear, one can show that the Jacobian is just an uninteresting constant.
Finally, using Eqs.(6) and (22), the interaction free energy for grounded manifolds can be obtained as
$$F_{\mathrm{gr}}=\frac{k_\mathrm{B}T}{4S_d^2}\underset{\alpha \beta }{}\left[\frac{\gamma _\alpha ^1\gamma _\beta ^1}{(d2)^2R_{\alpha \beta }^{2(d2)}}+\frac{(\gamma _\alpha ^1\gamma _{\beta ,ij}^1+\gamma _{\alpha ,ij}^1\gamma _\beta ^1)\widehat{R}_{\alpha \beta ,i}\widehat{R}_{\alpha \beta ,j}}{R_{\alpha \beta }^{2(d1)}}\right]+O(1/R^{2d}).$$
(23)
The first term in Eq.(23) is a squared monopole–monopole interaction, and is independent of the relative orientations of the conductors in space. The second term, on the other hand, has the form of a squared monopole–dipole interaction, and does depend on the orientations through an effective dipole–dipole interaction, which is pairwise additive.
Similarly, for isolated manifolds, Eqs.(7) and (22) yield the interaction as
$$F_{\mathrm{is}}=\frac{k_\mathrm{B}T}{4S_d^2}\underset{\alpha \beta }{}\frac{\gamma _{\alpha ,ik}^1\gamma _{\beta ,jl}^1}{R_{\alpha \beta }^{2d}}(\delta _{ij}d\widehat{R}_{\alpha \beta ,i}\widehat{R}_{\alpha \beta ,j})(\delta _{kl}d\widehat{R}_{\alpha \beta ,k}\widehat{R}_{\alpha \beta ,l})+O(1/R^{2d+2}).$$
(24)
Note that the leading term in Eq.(24) is a squared dipole–dipole interaction, and thus it is not orientationally pairwise additive.
## III Application to Specific Geometries
The multipole expansion allowed us to calculate the general forms of the fluctuation–induced interactions between manifolds of arbitrary shape and with arbitrary orientations with respect to one another, for the two cases of isolated and grounded boundary conditions. All the specific informations about the shapes and the orientations of the manifolds are encoded in the $`\gamma `$-tensors defined above. These informations are in fact of three kinds: (i) the overall magnitude of the tensors which are set by the typical sizes of the manifolds, (ii) the orientational dependencies which make up the tensorial structure, and are dictated by the structure of the symmetry axes or “the principal axes” of the manifolds, and (iii) overall numerical prefactors of order unity. In this section, we try to use symmetry arguments to determine the $`\gamma `$-tensors for some simple geometries within the numerical prefactors, without actually specifying the exact shape of the integration domain $`𝒟`$. The final piece of information, which is the numerical prefactor, appears to be very sensitive to the exact geometry of the manifold (and thus to that of $`𝒟`$), and can in general be calculated using the techniques developed in Ref. .
### A Two Spheres
The $`\gamma `$-tensors for a sphere of radius $`L`$ can be easily estimated using symmetry: $`\gamma _s_0^{1/L}k^{d1}𝑑k/k^21/L^{d2}`$ and $`\gamma _{s,ij}\delta _{ij}_0^{1/L}k^{d1}𝑑k\delta _{ij}/L^d`$. Using Eqs.(23) and (24), the interaction between a sphere of radius $`L_1`$, and another sphere of radius $`L_2`$ which is at a distance $`R`$, reads
$$F_{\mathrm{gr}}^{\mathrm{sph}}k_\mathrm{B}T\times \frac{L_1^{d2}L_2^{d2}}{R^{2(d2)}},$$
(25)
for the grounded case, and
$$F_{\mathrm{is}}^{\mathrm{sph}}k_\mathrm{B}T\times \frac{L_1^dL_2^d}{R^{2d}},$$
(26)
for the isolated case, with no orientational dependence due to symmetry.
### B Two Rods
The calculation of the $`\gamma `$-tensors for a rod of length $`L`$ and thickness $`a`$ is more tricky. Using the cylindrical symmetry one obtains: $`\gamma _r_0^{1/L}𝑑k_z_0^{1/a}𝑑k_{}k_{}^{d2}/(k_z^2+k_{}^2)1/L^{d2}`$ for $`d3`$, and $`1/(La^{d3})`$ for $`d>3`$, where $`z`$-axis is parallel to the director of the cylinder, and $``$ denotes the remaining directions that are perpendicular to it. The second rank tensor $`\gamma _{r,ij}`$ is diagonal with only two independent components: $`\gamma _{r,zz}_0^{1/L}𝑑k_z_0^{1/a}𝑑k_{}k_{}^{d2}k_z^2/(k_z^2+k_{}^2)1/L^d`$ for $`d3`$, and $`1/(L^3a^{d3})`$ for $`d>3`$, and $`\gamma _{r,}_0^{1/L}𝑑k_z_0^{1/a}𝑑k_{}k_{}^d/(k_z^2+k_{}^2)1/(La^{d1})`$. If the unit vector $`\widehat{𝐝}`$ denotes the director of the rod, the inverse second rank $`\gamma `$-tensor which appears in the expression for the interaction can be written as $`\gamma _{r,ij}^1L^d\widehat{d}_i\widehat{d}_j`$ for $`d3`$, in the limit of small thickness. Note that in this limit, the inverse $`\gamma `$-tensors are vanishing for $`d>3`$, and thus rods do not interact in these high dimensions.
Using Eqs.(23) and (24), the orientation dependent part of the interaction between two rods of lengths $`L_1`$ and $`L_2`$, and directors $`\widehat{𝐝}_1`$ and $`\widehat{𝐝}_2`$, which are a distance $`R`$ apart, reads $`(d3)`$
$$F_{\mathrm{gr}}^{\mathrm{rod}}k_\mathrm{B}T\times \frac{L_1^{d1}L_2^{d1}}{R^{2(d1)}}\times \left[\frac{L_1}{L_2}(\widehat{𝐝}_1\widehat{𝐑}_{12})^2+\frac{L_2}{L_1}(\widehat{𝐝}_2\widehat{𝐑}_{12})^2\right],$$
(27)
for the grounded case, which is pairwise additive, and
$$F_{\mathrm{is}}^{\mathrm{rod}}k_\mathrm{B}T\times \frac{L_1^dL_2^d}{R^{2d}}\times \left[\widehat{𝐝}_1\widehat{𝐝}_2d(\widehat{𝐝}_1\widehat{𝐑}_{12})(\widehat{𝐝}_2\widehat{𝐑}_{12})\right]^2,$$
(28)
for the isolated case, which has a squared dipolar form and is not pairwise additive.
### C Two Disks
The $`\gamma `$-tensors for a disk of radius $`L`$ and thickness $`a`$ can be similarly calculated within numerical prefactors using symmetry. The zeroth rank tensor can be calculated as: $`\gamma _d_0^{1/a}𝑑k_z_0^{1/L}𝑑k_{}k_{}^{d2}/(k_z^2+k_{}^2)1/L^{d2}`$, where $`z`$-axis is normal to the disk, and $``$ denotes the remaining directions in the subspace of the disk. The second rank tensor $`\gamma _{d,ij}`$ is diagonal with only two independent components: $`\gamma _{d,zz}_0^{1/a}𝑑k_z_0^{1/L}𝑑k_{}k_{}^{d2}k_z^2/(k_z^2+k_{}^2)1/(aL^{d1})`$, and $`\gamma _{d,}_0^{1/a}𝑑k_z_0^{1/L}𝑑k_{}k_{}^d/(k_z^2+k_{}^2)1/L^d`$. If we denote the unit vector perpendicular to the disk by $`\widehat{𝐧}`$, the inverse second rank $`\gamma `$-tensor which appears in the expression for the interaction can be written as $`\gamma _{d,ij}^1L^d(\delta _{ij}\widehat{n}_i\widehat{n}_j)`$, in the limit of small thickness.
Using Eqs.(23) and (24), the orientation dependent part of the interaction between a disk of radius $`L_1`$ and normal vector $`\widehat{𝐧}_1`$, and another one with radius $`L_2`$ and normal vector $`\widehat{𝐧}_2`$ that is a distance $`R`$ apart, reads
$$F_{\mathrm{gr}}^{\mathrm{disk}}k_\mathrm{B}T\times \frac{L_1^{d1}L_2^{d1}}{R^{2(d1)}}\times \left\{\frac{L_1}{L_2}\left[1(\widehat{𝐧}_1\widehat{𝐑}_{12})^2\right]+\frac{L_2}{L_1}\left[1(\widehat{𝐧}_2\widehat{𝐑}_{12})^2\right]\right\},$$
(29)
for the grounded case, which is pairwise additive, and
$`F_{\mathrm{is}}^{\mathrm{disk}}k_\mathrm{B}T\times {\displaystyle \frac{L_1^dL_2^d}{R^{2d}}}`$ $`\times `$ $`[d^2d2+(\widehat{𝐧}_1\widehat{𝐧}_2)^2+(2dd^2)(\widehat{𝐧}_1\widehat{𝐑}_{12})^2+(2dd^2)(\widehat{𝐧}_2\widehat{𝐑}_{12})^2`$ (31)
$`2d(\widehat{𝐧}_1\widehat{𝐑}_{12})(\widehat{𝐧}_2\widehat{𝐑}_{12})(\widehat{𝐧}_1\widehat{𝐧}_2)+d^2(\widehat{𝐧}_1\widehat{𝐑}_{12})^2(\widehat{𝐧}_2\widehat{𝐑}_{12})^2],`$
for the isolated case, which has a squared dipolar form and is not pairwise additive.
## IV Critical Fluids
As mentioned above, interactions could be induced between objects that modify thermal fluctuations of an order parameter field for a critical binary mixture or a magnetic system . In this case, two kinds of boundary conditions are usually considered: (i) the ordinary boundary condition that suppresses the order parameter at the boundary, and thus does not break its symmetry, and (ii) the symmetry breaking boundary condition, which sets a nonvanishing value for the order parameter at the boundary. Note that the ordinary boundary condition is the same as the grounded boundary condition in the electrostatic terminology.
The fluctuations in a critical fluid are characterized by the universality class of the system. For the case when the fluid can be described by a Gaussian Hamiltonian as given in Eq.(1), all of the above results for the grounded manifolds hold for the case of ordinary boundary condition. The interaction between manifolds with symmetry breaking boundary conditions, where the value of the order parameter is set to $`\mathrm{\Phi }_\alpha `$ on the $`\alpha `$th manifold, is calculated as
$$F_{\mathrm{sb}}=k_\mathrm{B}T\mathrm{ln}𝒵[\varphi _\alpha =\mathrm{\Phi }_\alpha ],$$
(32)
where $`𝒵[\varphi _\alpha ]`$ is given by Eq.(22). One obtains
$$F_{\mathrm{sb}}^{\mathrm{Gauss}}=\frac{k_\mathrm{B}T}{2}\underset{\alpha \beta }{}\left[\frac{K\mathrm{\Phi }_\alpha \mathrm{\Phi }_\beta \gamma _\alpha ^1\gamma _\beta ^1}{S_d(d2)R_{\alpha \beta }^{d2}}\right]+O(1/R^{2d4}).$$
(33)
It is important to note that this interaction is independent of the orientations of the manifolds, and that the leading order orientation dependent term for the symmetry breaking boundary condition is the same as the case of ordinary boundary condition, and is given as in Eq.(23). This interaction is orientationally pairwise additive.
For a nontrivial universality class, one should make use of more complicated Hamiltonians with nonlinear terms. It is then possible to calculate the Casimir energy expressions using field theoretical techniques . The interaction between two spheres in an arbitrary critical system has in fact been calculated exactly in Ref. using conformal-invariance methods. The interaction for the case of symmetry breaking boundary conditions (on both spheres) is obtained as $`1/R^{d2+\eta }`$, while for the case of ordinary boundary conditions it is found as $`1/R^{2(d1/\nu )}`$, where $`\eta `$ and $`\nu `$ are critical exponents of the system . One can easily check that for the case of Gaussian universality class, where $`\eta =0`$ and $`\nu =1/2`$, they coincide with the results of Eq.(33) and (23).
It is interesting to note that the power law for the symmetry breaking case is given by the two-point correlation function of the field (the spin-spin correlation in magnetic terminology), while the one for the ordinary case is given by the four-point correlation function (the energy-energy correlation) . Guided by this, one can think of an effective Gaussian Hamiltonian of the form
$$_{\mathrm{cf}}[\varphi ]=\frac{K}{2}\frac{d^d𝐪}{(2\pi )^d}q^{2\eta }|\varphi (𝐪)|^2,$$
(34)
which yields a correct form for the two-point function, and calculate the fluctuation–induced interactions using the above methods . However, although it yields a correct result for the symmetry breaking case (almost by construction), it gives a corresponding form for the ordinary case as $`1/R^{2(d2+\eta )}`$ which is not correct. The reason is that the above effective Gaussian Hamiltonian does not give a correct four-point correlation function. However, it can be constructed to do so by using $`q^{1/\nu }`$ instead of $`q^{2\eta }`$ in Eq.(34).
## V Conclusion
The analysis that is presented here is aimed at emphasizing the crucial role of the type of boundary conditions on fluctuation–induced interactions. Unlike the case of extended objects at close separations, where different types of boundary conditions all lead to the same form of interaction, we found that for objects at large separations the type of boundary conditions determine the very form of the interaction, and whether or not it is pairwise additive.
## ACKNOWLEDGMENTS
I am grateful to M. Kardar for invaluable discussions and comments. This research was supported in part by the National Science Foundation under Grants No. PHY94-07194 and DMR-98-05833.
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# 1 Introduction
## 1 Introduction
Quantization of gravitation is still today, one of the most challenging problem in theoretical physic. A possible approach to have an idea of what happens when passing from classical theory to the quantum one consists in trying simpler models. The case of two commuting Killing vector reduction of source-free general relativity is probably one of the most interesting. This model has the particularity to exhibit an infinite dimensional symmetry group, the so-called Geroch group . It is known to be integrable since the works of Belinskii, Zakharov and Maison . One way to quantize is to apply methods used for standard integrable models. This quantization problem, which is equivalent to $`\mathrm{𝐒𝐋}(\mathrm{𝟐},𝐑)/\mathrm{𝐒𝐎}(\mathrm{𝟐})`$ coset space $`\sigma `$-models coupled to two-dimensional gravity and a dilaton (it can be generalized to $`𝐆/𝐇`$ coset space), has been of course the subject of several papers ( and references therein).
Before trying to quantize this system, it is necessary to study in details the Poisson algebra. The classical theory is also interesting for its own. An extensive work has already been done in this domain by Julia, Korotkin, Nicolai and Samtleben . A dynamical $`r`$matrix formulation of the model has been proposed (see e.g. ) where the dilaton is given a priori. A restriction of this approach is that the brackets are evaluated on the constraint surface, which prevents deduction of the associated Yang-Baxter equations. In an expression was proposed for the Lax connection based on twisted $`sl(2,R)`$ affine Kac-Moody and Virasoro algebras that reproduces the equations of motion. Using dressing transformations, it provides a rather elegant method to generate solutions . The aim of this article is to show that this form of the Lax connection can also provide a good basis to obtain a $`r`$matrix formulation of this problem. This means that all fields are considered as dynamical variables and pure $`c`$-number Yang-Baxter equations can be deduced. The structure we obtain is closed to Toda affine model’s one. We hope that we will be able to transpose what have been done in this domain to our problem. This will provide an alternative algebraic approach for the quantization of 2d reduced gravity and thus a complementary point of view to what has already been done.
This paper is organized as follow. Section 2 sums up some of the main results of , in particular we introduce the Lax pair. Section 3 deals with the Hamiltonian formulation of the theory. The calculation of the Poisson brackets of the Lax connection, the key point of this paper, and the deduction of the associated Yang-Baxter equations are described in section 4. We show that despite the $`r`$dynamical behavior of the model, we obtain pure $`c`$-number modified Yang-Baxter relations of kind
$`[r_{12}^{ϵ_1},r_{23}^{ϵ_2}]+[s_{23}^{ϵ_2},s_{31}^{ϵ_3}]+[s_{31}^{ϵ_3},r_{12}^{ϵ_1}]{\displaystyle \frac{1}{2}}k_2s_{31}^{ϵ_3}{\displaystyle \frac{1}{2}}k_3r_{12}^{ϵ_1}{\displaystyle \frac{1}{4}}[U_{23},c_{12}]`$ $`=`$ $`0`$
that can be interpreted as consistency conditions for a simpler static linear model (but still lack of a quadratic interpretation). As an application of the previous results, we determine the Poisson brackets of monodromy matrices in section 5, and we point out the problem of coincident points. Section 6 is an attempt to find classical observables. We show that if we impose reasonable boundaries conditions, it is possible to construct an infinite set of these objects. Finally, we have gathered in Appendix A and B, some expressions and sketches of demonstrations related to section 4.
## 2 Equation of motion and Lax connection
In this section, we will review some basics facts and results found in . We will also introduce notations used in this article. Let us recall the parameterization we choose for the metric :
$`ds^2=\rho ^{\frac{1}{2}}e^{2\widehat{\sigma }}(dt^2+dx^2)+\rho S_{ij}(x,t)dy^idy^j`$ (1)
where $`\rho `$ is called the dilaton and $`\widehat{\sigma }`$ the conformal factor. The symmetric $`2\times 2`$ matrix $`S`$, normalized by $`det(S)=1`$, can be written as $`S=𝒱^t𝒱`$, where $`𝒱`$ is an element of $`SL(2,R)`$. $`𝒱`$ is equivalent to internal zweibein, up to a $`\sqrt{\rho }`$ factor. There is a manifest local $`SO(2)`$ gauge symmetry when multiplying $`𝒱`$ to the left with any element of $`SO(2)`$.
We introduce the decomposition of $`sl(2,R)=hg`$ where $`h=so(2)`$ is the maximal compact subalgebra of $`sl(2,R)`$. We will use the following notation for the generators : $`T^\alpha `$ with a Greek index correspond to the generator of $`h`$ and $`T^a`$ with a Latin index correspond to the generators of $`g`$. We choose these generators such that they are orthogonal and normalized with respect to the Killing form. To formulate the vacuum Einstein’s equation (the so-called Ernst’s equations), we introduce the connection $`𝒱𝒱^1`$. We denote each component of this connection as $`P_x+Q_x=𝒱_x𝒱^1`$ and $`P_t+Q_t=𝒱_t𝒱^1`$, where $`P`$ is an element of $`g`$ and $`Q`$ belong to $`h`$. With these objects, Einstein’s equation can be brought to the form
$`_xQ_t_tQ_x+[Q_x,Q_t]+[P_x,P_t]`$ $`=`$ $`0`$ (2)
$`_xP_t+[Q_x,P_t]`$ $`=`$ $`_tP_x+[Q_t,P_x]`$ (3)
$`_x\left(\rho P_x\right)+[Q_x,\rho P_x]`$ $`=`$ $`_t\left(\rho P_t\right)+[Q_t,\rho P_t]`$ (4)
$`\left(_t^2_x^2\right)\rho `$ $`=`$ $`0`$ (5)
$`\left(\left(_t\pm _x\right)\rho \right)\left(_t\pm _x\right)\widehat{\sigma }`$ $`=`$ $`\rho {\displaystyle \frac{1}{2}}tr\left(\left(P_x\pm P_t\right)^2\right)`$ (6)
Before dealing with the Lax connection, we shall introduce the algebra we will use. Consider the $`sl(2,R)`$ affine Kac-Moody algebra defined by the commutation relations :
$`[X\lambda ^n,Y\lambda ^m]`$ $`=`$ $`[X,Y]\lambda ^{m+n}+n{\displaystyle \frac{k}{2}}tr\left(XY\right)\delta _{n+m,0}`$ (7)
We twist this algebra with the order two automorphism that leaves $`h`$ invariant. It means that for some element $`X\lambda ^n`$, if $`n`$ is even, then $`X`$ is an element of $`h`$, else $`X`$ is an element of $`g`$. In fact, we will use the semi direct product of this algebra with the Virasoro’s one. We recall that the commutation relations for the Virasoro algebra are
$`[L_n,L_m]`$ $`=`$ $`(nm)L_{m+n}+n(n^21){\displaystyle \frac{c}{12}}\delta _{n+m,0}`$
and the crossed Lie bracket is $`[L_n,X\lambda ^m]=\frac{m}{2}X\lambda ^{n+m}`$. For convenience, we introduce a particular notation for two elements of the Virasoro algebra $`E_\pm =L_0L_{\pm 1}`$ which verify the commutation relation $`[E_+,E_{}]=E_++E_{}`$.
As we said before, the model is integrable. It means that an auxiliary linear system
$`\left(_t+A_t\right)\mathrm{\Psi }=0`$ $`\mathrm{and}`$ $`\left(_x+A_x\right)\mathrm{\Psi }=0`$ (8)
can be found such that the zero curvature condition $`[_t+A_t,_x+A_x]=0`$ reproduces the equations of motion. The expression of the components of Lax connection that fulfills this requirement, is
$`A_x`$ $`=`$ $`{\displaystyle \frac{1}{2}}\rho ^1(\mathrm{\Pi }_{\widehat{\sigma }}_x\rho )E_+{\displaystyle \frac{1}{2}}\rho ^1(\mathrm{\Pi }_{\widehat{\sigma }}+_x\rho )E_{}+{\displaystyle \frac{1}{2}}\left(P_{xa}+P_{ta}\right)T^a\lambda `$ (9)
$`+{\displaystyle \frac{1}{2}}\left(P_{xa}P_{ta}\right)T^a\lambda ^1+Q_{x\alpha }T^\alpha +\mathrm{\Pi }_\rho {\displaystyle \frac{k}{2}}`$
$`A_t`$ $`=`$ $`{\displaystyle \frac{1}{2}}\rho ^1(\mathrm{\Pi }_{\widehat{\sigma }}_x\rho )E_++{\displaystyle \frac{1}{2}}\rho ^1(\mathrm{\Pi }_{\widehat{\sigma }}+_x\rho )E_{}+{\displaystyle \frac{1}{2}}\left(P_{xa}+P_{ta}\right)T^a\lambda `$ (10)
$`{\displaystyle \frac{1}{2}}\left(P_{xa}P_{ta}\right)T^a\lambda ^1+Q_{t\alpha }T^\alpha _x\widehat{\sigma }{\displaystyle \frac{k}{2}}`$
where $`\mathrm{\Pi }_{\widehat{\sigma }}=_t\rho `$ and $`\mathrm{\Pi }_\rho =_t\widehat{\sigma }`$.
We will often use the notation $`A`$ for $`A_x`$ and the term connection for the Lax connection (from now, we will deal no more with the connection $`𝒱𝒱^1`$, so there will be no misunderstanding). Notice that the zero curvature condition reproduces the equations of motion (2) to (5) and a second order equation for $`\widehat{\sigma }`$
$`\left(_t^2_x^2\right)\widehat{\sigma }`$ $`=`$ $`{\displaystyle \frac{1}{2}}tr\left(P_x^2P_t^2\right)`$ (11)
which is a consequence of the two linear equations for the conformal factor (6).
This connection was first used as a powerful method to generate solutions to Einstein’s equations. Readers who are interested, can found details in and . When comparing with other Lax connection used in the literature (see e.g.), we remark that all fields are considered on an equal footing (in general, the dilaton is supposed to be given and the conformal factor is deduced from the other variables). So, if we want to consider all of them at the same time, this connection seems to be a good candidate. We will see that it introduces no additional difficulty and, at the contrary, yields simpler and more compact expressions for the Poisson brackets for the connection.
## 3 Action and canonical brackets
The hamiltonian formulation of 2d reduced gravity, has already been studied in various papers . So, up to some minor changes, the formulation we shall use is identical to what can be found in the literature. Let us recall it.
First of all, let describes our phases space. It is defined by the canonical variables $`P_x,Q_x,\rho ,\widehat{\sigma }`$ and their associated momenta $`\mathrm{\Pi }_P,\mathrm{\Pi }_Q,\mathrm{\Pi }_\rho ,\mathrm{\Pi }_{\widehat{\sigma }}`$. Here we use the canonical Poisson brackets to define the symplectic structure ($`\{\rho (x),\mathrm{\Pi }_\rho (y)\}=\delta (xy)`$ and so on).
To make contact with the model, we have to express the quantities $`P_t`$ and $`Q_t`$ in terms of the canonical variables. We also need the generators of the transformations associated to the invariances of our system which are invariance under reparameterization and the local $`SO(2)`$ invariance. To solve this problem, we can either try to deduce these formulae from some mathematical procedures (see for example), or just give expressions as definition and verify if this choice is coherent. Here, we will adopt the second method. First we define the variable $`P_t`$ as
$`\rho P_t`$ $`=`$ $`_x\mathrm{\Pi }_P+[Q_x,\mathrm{\Pi }_P]+[P_x,\mathrm{\Pi }_Q]`$ (12)
$`Q_t`$ will be considered to have vanishing brackets with all variables.
This phase space is reduced by the three constraints arising from the gauge invariance mentioned above. First we have the Hamiltonian $``$ which can be written as
$``$ $`=`$ $`\mathrm{\Pi }_\rho \mathrm{\Pi }_{\widehat{\sigma }}_x\rho _x\widehat{\sigma }+{\displaystyle \frac{1}{2}}\rho tr\left(P_t^2+P_x^2\right)+tr\left(Q_t\mathrm{\Phi }\right)`$ (13)
With this expression for the Hamiltonian, the equations of motion (2-5 and 11) are correctly reproduced (the brackets needed for these calculations are given below). We still have two other generators of gauge transformations. We denote $`𝒫`$ the generator of diffeomorphisms in the spatial direction whose expression is
$`𝒫`$ $`=`$ $`\mathrm{\Pi }_\rho _x\rho +\mathrm{\Pi }_{\widehat{\sigma }}_x\widehat{\sigma }+\rho tr\left(P_tP_x\right)+tr\left(Q_x\mathrm{\Phi }\right)`$ (14)
The linear combinations $`𝒞_\pm =\pm 𝒫0`$ of these two constraints are equivalent to the two linear equations for the conformal factor (6). Finally, the generator $`\mathrm{\Phi }`$ of the $`SO(2)`$ gauge invariance takes the following form
$`\mathrm{\Phi }`$ $`=`$ $`_xQ_x+[Q_x,\mathrm{\Pi }_Q]+[P_x,\mathrm{\Pi }_P]`$ (15)
It could be easily verified that $`\mathrm{\Phi }`$ belongs to $`so(2)`$. Notice that all these constraints are first class constraints.
We will use the standard index-free tensor notation. For some element $`X`$, we define $`X_1XI`$ and $`X_2IX`$. We will also introduce the decomposition of the Casimir element $`𝒞_{12}`$ of $`sl(2,R)`$: $`𝒞_{12}=c_{12}+d_{12}`$ with $`c_{12}=T^\alpha T_\alpha `$ and $`d_{12}=T^aT_a`$. The validity of this decomposition is due to orthogonality of generators with respect to the Killing form.
We shall list all the basic Poisson brackets needed for further calculations.
$`\{P_{t1}(x),P_{t2}(y)\}`$ $`=`$ $`\rho ^2(x)\delta (xy)[d_{12},\mathrm{\Phi }_2(x)]`$
$`\{P_{t1}(x),P_{x2}(y)\}`$ $`=`$ $`\rho ^1(x)\delta ^{}(xy)d_{12}+\rho ^1(x)\delta (xy)[d_{12},Q_{x2}(x)]`$
$`\{P_{t1}(x),Q_{x2}(y)\}`$ $`=`$ $`\rho ^1(x)\delta (xy)[d_{12},P_{x2}(x)]`$
$`\{\mathrm{\Phi }_1(x),\mathrm{\Phi }_2(y)\}`$ $`=`$ $`\delta (xy)[c_{12},\mathrm{\Phi }_2(x)]`$
$`\{P_{t1}(x),\mathrm{\Phi }_2(y)\}`$ $`=`$ $`\delta (xy)[c_{12},P_{t1}(x)]`$
$`\{P_{x1}(x),\mathrm{\Phi }_2(y)\}`$ $`=`$ $`\delta (xy)[c_{12},P_{x1}(x)]`$
$`\{Q_{x1}(x),\mathrm{\Phi }_2(y)\}`$ $`=`$ $`\delta ^{}(xy)c_{12}+\delta (xy)[c_{12},Q_{x2}(x)]`$
The four last commutators show that $`\mathrm{\Phi }`$ is the generator of the local $`SO(2)`$ invariance. Equivalent calculations can be done to prove other expressions.
## 4 Poisson brackets for the Lax connection
### 4.1 Main result
Now that we have defined the Poisson algebra, we can deduce the Poisson bracket for the Lax connection (the so-called fundamental Poisson brackets). The raw formula derived from a direct calculation, is quite long and has no pedagogic interest. We won’t write its expression (readers who want to obtain it will encounter no difficulty). What is really interesting is that it can be put on a $`r`$matrix form. A brief survey of the method used to find this expression is described in Appendix A. The result we have obtained is
$`\{A_1(x),A_2(y)\}`$ $`=`$ $`{\displaystyle \frac{1}{\rho (x)}}\delta (xy)\left([r_{12}^ϵ,A_1(x)]+[s_{12}^ϵ,A_2(x)]\right)`$
$`+`$ $`\left({\displaystyle \frac{1}{\rho (x)}}s_{12}{\displaystyle \frac{1}{\rho (y)}}r_{12}\right)_x\delta (xy)`$
$``$ $`{\displaystyle \frac{1}{8}}{\displaystyle \frac{1}{\rho ^2(x)}}\delta (xy)[U_{12},\mathrm{\Phi }_1(x)\mathrm{\Phi }_2(x)]`$
where
$`U_{12}`$ $`=`$ $`d_{12}(\lambda _1\lambda _1^1)(\lambda _2\lambda _2^1)`$
This formula is the one obtained in the case of non-ultralocal theories with an additional term coming from the local $`SO(2)`$ invariance. This term has to be considered with caution when dealing with the Jacobi identity. Here are the expressions of the $`r`$ and $`s`$matrices :
$`r_{12}^\pm `$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{(1\lambda _1^2)(1\lambda _2^2)}{\lambda _1^2\lambda _2^2}}c_{12}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\lambda _1\lambda _2^1(1\lambda _2^2)^2}{\lambda _1^2\lambda _2^2}}d_{12}{\displaystyle \frac{1}{2}}\left(E_\pm k+{\displaystyle \frac{1}{2}}k\left(E_++E_{}\right)\right)`$ (17)
$`s_{12}^\pm `$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{(1\lambda _1^2)(1\lambda _2^2)}{\lambda _1^2\lambda _2^2}}c_{12}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\lambda _1^1\lambda _2(1\lambda _1^2)^2}{\lambda _1^2\lambda _2^2}}d_{12}{\displaystyle \frac{1}{2}}\left(kE_{}+{\displaystyle \frac{1}{2}}\left(E_++E_{}\right)k\right)`$ (18)
The matrices involved here are pure c-number, all coordinates dependencies have been factorized in the $`\rho ^1`$ factors. The rational functions that appear in (17) and (18), have only a meaning as formal power series. So, whether we choose $`|\lambda _1|<|\lambda _2|`$ or $`|\lambda _1|>|\lambda _2|`$ when developing, we obtain two different sets of matrices (here + convention refers to the case $`|\lambda _1|<|\lambda _2|`$). Fully developed formulas are given in appendix A.
Although our Lax connection is not exactly the same that the one used by Korotkin and Samtleben in , it however is possible to compare some pieces of (4.1) with their expression. The algebra we used is just a way to eliminate the coordinates dependence of the moving poles. In order to compare the two formulae we have to restore this dependence. It can be achieved by formally substituting $`\frac{1\gamma }{1+\gamma }`$ by $`\lambda `$ where $`\gamma `$ is the moving pole (a detailed explanation of this equivalence can be found in ). One can now see, their ultralocal part of the Poisson brackets for the connection is equivalent to ours if we just keep the loop part of the $`r`$matrices (remember that the introduction of the central extension $`k`$ and the Virasoro algebra is a method to take $`\rho `$ and $`\widehat{\sigma }`$ into account). The case of the non-ultralocal part is more difficult. No direct comparison can be done because the dilaton produces additional terms. Notice that their brackets is calculated on the constraint surface, thus they have no additional term involving $`\mathrm{\Phi }`$ and they can’t explicitly verify the Jacobi identity (which imply that Yang-Baxter equations can’t be found).
Now, the fundamental Poisson brackets have to satisfy the standard relations of Poisson brackets and to be independent of the convention we choose. We will focus here on antisymmetry and independence, dealing with the Jacobi identity in the next subsection. Proofs of these two properties lie on the same relations of the $`r`$matrices, that can be divided into two sets. On one hand, there is a set a pure numerical identities :
$`r_{12}^ϵ`$ $`=`$ $`s_{21}^ϵ`$
$`r_{12}^ϵr_{12}^ϵ`$ $`=`$ $`s_{12}^ϵs_{12}^ϵ`$ (19)
$`U_{12}`$ $`=`$ $`U_{21}`$
and on the other hand, we have got a relation involving connection and dilaton :
$`[r_{12}^ϵr_{12}^ϵ,A_1+A_2]`$ $`=`$ $`\left(\rho ^1_x\right)\rho (r_{12}^ϵr_{12}^ϵ)`$ (20)
Using these formulae, antisymmetry and independence can be easily shown. For example, from eq.(19), we have
$`\{A_1(x),A_2(y)\}|_{ϵ=+}\{A_1(x),A_2(y)\}|_{ϵ=}`$ $`=`$ $`\rho ^1(x)[r_{12}^+r_{12}^{},A_1(x)+A_2(x)]\delta (xy)`$
$`+\left(r_{12}^+r_{12}^{}\right)(\rho ^1(x)\rho ^1(y))_x\delta (xy)`$
Now with eq.(20) and the identity $`(\rho ^1(x)\rho ^1(y))_x\delta (xy)=\rho ^2(x)_x\rho \delta (xy)`$, we easily prove that the right hand-side of the above equation vanishes.
Notice that in more conventional cases like Toda field theories, only numerical relations are used. In particular the difference $`r_{12}^+r_{12}^{}`$ is generally proportional to the Casimir tensor which would simplify (20). This more complicated form is a consequence of the $`\rho ^1`$ terms in (4.1).
### 4.2 Jacobi identity and Yang-Baxter equations
Finally, we have to prove the validity of the Jacobi identity. Performing the calculation leads to the following formula :
$`\{A_1(x),\{A_2(y),A_3(z)\}\}+\mathrm{perm}.`$ $`=`$ $`\rho ^2(x)\delta (xy)\delta (yz)([A_1(x),A_{123}]+\mathrm{perm}.)`$
$`+`$ $`\left(_x\rho ^2\right)\delta (xy)\delta (yz)(B_{123}+\mathrm{perm}.)`$
$`+`$ $`\rho ^3(x)\delta (xy)\delta (yz)([\mathrm{\Phi }_1(x),C_{123}]+\mathrm{perm}.)`$
$`+`$ $`\rho ^2(x)\delta (yz)_x\delta (xy)D_{123}+\mathrm{perm}.`$
Explicit expressions of $`A_{123}`$, $`B_{123}`$, $`C_{123}`$ and $`D_{123}`$ are gathered in appendix B. The fact that eq.(4.2) is equal to zero, is a consequence of the properties of $`A_{123}`$ and $`C_{123}`$. One can shown they are invariant under cyclic permutations and equal to zero (we left details in appendix B). This leads to the modified Yang-Baxter equations
$`[r_{12}^{ϵ_1},r_{23}^{ϵ_2}]+[s_{23}^{ϵ_2},s_{31}^{ϵ_3}]+[s_{31}^{ϵ_3},r_{12}^{ϵ_1}]{\displaystyle \frac{1}{2}}k_2s_{31}^{ϵ_3}{\displaystyle \frac{1}{2}}k_3r_{12}^{ϵ_1}{\displaystyle \frac{1}{4}}[U_{23},c_{12}]`$ $`=`$ $`0`$ (22)
$`[r_{23}^{ϵ_2},U_{12}]+[s_{23}^{ϵ_2},U_{13}]+{\displaystyle \frac{1}{2}}k_3U_{12}{\displaystyle \frac{1}{2}}k_2U_{13}`$ $`=`$ $`0`$ (23)
Notice that the choice of the three conventions has to be consistent. This condition can be written as $`\left|ϵ_1+ϵ_2+ϵ_3\right|=1`$.
Thus, the validity conditions of the Jacobi identity are pure c-numbers equations. It could seem to be a miracle, especially when looking at first sight at equation (4.1) with its explicit $`\rho `$ dependence. Actually this dependence is encoded in the terms involving the central extension. The $`SO(2)`$ gauge invariance generates the term $`[U_{23},c_{12}]`$ and eq.(23). If we want to compare these Yang-Baxter equations with those obtained in simpler cases, we have to drop the central extension and the local $`SO(2)`$ constraint. In this case, we obtain the same results as in the case of non-ultralocal theories with constant $`r`$matrices (see ).
A possible interpretation of these Yang-Baxter equations is as consistency conditions for linear Poisson brackets involving three objects $`L^\pm `$ and $`\varphi `$. Defining the following algebra
$`\{L_1^\pm ,L_2^\pm \}`$ $`=`$ $`[r_{12}^ϵ,L_1^\pm ]+[s_{12}^ϵ,L_2^\pm ]+{\displaystyle \frac{1}{2}}k_2L_1^\pm {\displaystyle \frac{1}{2}}k_1L_2^\pm {\displaystyle \frac{1}{8}}[U_{12},\varphi _1\varphi _2]`$
$`\{L_1^+,L_2^{}\}`$ $`=`$ $`[r_{12}^{},L_1^+]+[s_{12}^{},L_2^{}]+{\displaystyle \frac{1}{2}}k_2L_1^+{\displaystyle \frac{1}{2}}k_1L_2^{}{\displaystyle \frac{1}{8}}[U_{12},\varphi _1\varphi _2]`$
$`\{L_1^{},L_2^+\}`$ $`=`$ $`[r_{12}^+,L_1^{}]+[s_{12}^+,L_2^+]+{\displaystyle \frac{1}{2}}k_2L_1^{}{\displaystyle \frac{1}{2}}k_1L_2^+{\displaystyle \frac{1}{8}}[U_{12},\varphi _1\varphi _2]`$
$`\{L_1^\pm ,\varphi _2\}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[L_1^\pm L_2^\pm ,c_{12}]`$
$`\{\varphi _1,\varphi _2\}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[\varphi _1\varphi _2,c_{12}]`$
We suppose that $`k`$ commutes with all other elements. $`U_{12}`$ and $`c_{12}`$ are considered to be symmetric when permuting the two spaces. $`r_{12}`$ and $`s_{12}`$ have to fulfilled (19) and (20) with $`\rho `$ constant (thus right-hand side of (20) vanishes). Under these assumptions, Yang-Baxter equations (22) and (23) can be deduced from Jacobi identity of this algebra.
Finally, remark that all these formulae are independent from the choice of $`SL(2,R)/SO(2)`$ coset model. The generalization to any $`G/H`$ coset is obvious : the tensors $`c_{12}`$ and $`d_{12}`$ have to be replaced by those of the new algebra.
## 5 Monodromy matrices
Now that we have determined the Poisson brackets of the Lax connection in the previous section, we can go further and calculate those of the monodromy matrices. So we want to determine $`\{\mathrm{\Psi }_1(x,x_0),\mathrm{\Psi }_2(y,y_0)\}`$ , where $`x>x_0`$, $`y>y_0`$ and the four points are distinct points. The way to achieve this calculation is as follows. We will use the Leibniz rule
$`\{\mathrm{\Psi }_{ij}(x),\mathrm{\Psi }_{kl}(y)\}`$ $`=`$ $`{\displaystyle 𝑑z𝑑z^{}\frac{\delta \mathrm{\Psi }_{ij}(x)}{\delta A_{mn}(z)}\{A_{mn}(z),A_{m^{}n^{}}(z^{})\}\frac{\delta \mathrm{\Psi }_{kl}(y)}{\delta A_{m^{}n^{}}(z^{})}}`$
The only difficulty is to find the functional derivative of $`\mathrm{\Psi }(x,x_0)`$ with respect to $`A(z)`$. It can be done by solving the differential equations
$`_x\delta \mathrm{\Psi }(x,x_0)+A(x)\delta \mathrm{\Psi }(x,x_0)+\delta A(x)\mathrm{\Psi }(x,x_0)=0`$
$`_{x_0}\delta \mathrm{\Psi }(x,x_0)\delta \mathrm{\Psi }(x,x_0)A(x_0)\mathrm{\Psi }(x,x_0)\delta A(x_0)=0`$
With the condition $`\delta \mathrm{\Psi }(x,x)=0`$, the solution is given by
$`\delta \mathrm{\Psi }(x,x_0)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑z\mathrm{\Theta }(xz)\mathrm{\Theta }(zx_0)\mathrm{\Psi }(x,z)\delta A(z)\mathrm{\Psi }(z,x_0)`$
where $`\mathrm{\Theta }(x)`$ is the Heaviside function ($`\mathrm{\Theta }(x)=1`$ for $`x>0`$, $`\mathrm{\Theta }(x)=0`$ elsewhere). The rest of the calculations is quite easy, consisting in putting together terms to form total derivatives. Thus, we obtain the following expression :
$`\{\mathrm{\Psi }_1(x,x_0),\mathrm{\Psi }_2(y,y_0)\}`$ $`=`$ $`\mathrm{\Theta }(y,x,y_0)\rho ^1(x)\mathrm{\Psi }_2(y,x)r_{12}^ϵ\mathrm{\Psi }_1(x,x_0)\mathrm{\Psi }_2(x,y_0)`$
$`\mathrm{\Theta }(x,y,x_0)\rho ^1(y)\mathrm{\Psi }_1(x,y)s_{12}^ϵ\mathrm{\Psi }_1(y,x_0)\mathrm{\Psi }_2(y,y_0)`$
$`+\mathrm{\Theta }(y,x_0,y_0)\rho ^1(x_0)\mathrm{\Psi }_1(x,x_0)\mathrm{\Psi }_2(y,x_0)r_{12}^ϵ\mathrm{\Psi }_2(x_0,y_0)`$
$`+\mathrm{\Theta }(x,y_0,x_0)\rho ^1(y_0)\mathrm{\Psi }_1(x,y_0)\mathrm{\Psi }_2(y,y_0)s_{12}^ϵ\mathrm{\Psi }_1(y_0,x_0)`$
$`{\displaystyle \frac{1}{8}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑z\mathrm{\Theta }(xz)\mathrm{\Theta }(zx_0)\mathrm{\Theta }(yz)\mathrm{\Theta }(zy_0)`$
$`\rho ^2(z)\mathrm{\Psi }_1(x,z)\mathrm{\Psi }_2(y,z)[U_{12},\mathrm{\Phi }_1(z)\mathrm{\Phi }_2(z)]\mathrm{\Psi }_1(z,x_0)\mathrm{\Psi }_2(z,y_0)`$
with $`\mathrm{\Theta }(x,y,z)`$ equal to 1 if $`x>y>z`$, and 0 for the other case.
We can easily verify that these brackets are consistent. It is also a consequence of relations (19), (22), (23) and (20). This last one gives the following relation for monodromy matrices
$`\rho ^1(x)\left(r_{12}^ϵr_{12}^ϵ\right)\mathrm{\Psi }_1(x,y)\mathrm{\Psi }_2(x,y)`$ $`=`$ $`\rho ^1(y)\mathrm{\Psi }_1(x,y)\mathrm{\Psi }_2(x,y)\left(r_{12}^ϵr_{12}^ϵ\right)`$ (25)
We need also to evalute the brackets between $`\mathrm{\Psi }`$ and $`\rho `$ that can be deduced form those of $`A`$ and $`\rho `$
$`\{A(x),\rho (y)\}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\delta (xy)k`$ (26)
$`\{\mathrm{\Psi }(x,y),\rho (z)\}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Theta }(xz)k\mathrm{\Psi }(x,y){\displaystyle \frac{1}{2}}\mathrm{\Theta }(yz)k\mathrm{\Psi }(x,y)`$ (27)
What happens if we let the two ending points (or the two starting points) tend to the same value? In this case, we have to face to a well-known problem of non-ultralocal theories (see for example) that the brackets are ill-defined (in particular, the Jacobi identity is no longer valid). In the case of non-linear sigma models, it has been shown that no regularization of the Poisson brackets is coherent. One way to go beyond this problem is to have additional informations about the boundaries conditions. For example, if we consider $`\rho `$ as the radial coordinate, we can choose a frame such that $`\rho (x)`$ tends to $`\mathrm{}`$ when $`x`$ goes to a given point $`x_{\mathrm{}}`$. If we take this naive limit in eq.(5), assuming that the $`\mathrm{\Psi }(x,x_{\mathrm{}})`$ terms have a good behavior compared to $`\rho ^1(x)`$, we obtain the following relation
$`\{\mathrm{\Psi }_1(x,x_{\mathrm{}}),\mathrm{\Psi }_2(y,x_{\mathrm{}})\}`$ $`=`$ $`\mathrm{\Theta }(yx)\rho ^1(x)\mathrm{\Psi }_2(y,x_{\mathrm{}})\mathrm{\Psi }_2^1(x,x_{\mathrm{}})r_{12}^ϵ\mathrm{\Psi }_1(x,x_{\mathrm{}})\mathrm{\Psi }_2(x,x_{\mathrm{}})`$
$`\mathrm{\Theta }(xy)\rho ^1(y)\mathrm{\Psi }_1(x,x_{\mathrm{}})\mathrm{\Psi }_1^1(y,x_{\mathrm{}})s_{12}^ϵ\mathrm{\Psi }_1(y,x_{\mathrm{}})\mathrm{\Psi }_2(y,x_{\mathrm{}})`$
$`{\displaystyle \frac{1}{8}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑z\mathrm{\Theta }(xz)\mathrm{\Theta }(yz)\rho ^2(z)\mathrm{\Psi }_1(x,x_{\mathrm{}})\mathrm{\Psi }_2(y,x_{\mathrm{}})`$
$`\mathrm{\Psi }_1^1(z,x_{\mathrm{}})\mathrm{\Psi }_2^1(z,x_{\mathrm{}})[U_{12},\mathrm{\Phi }_1(z)\mathrm{\Phi }_2(z)]\mathrm{\Psi }_1(z,x_{\mathrm{}})\mathrm{\Psi }_2(z,x_{\mathrm{}})`$
It can be shown that the above definition, these brackets are well-defined. We focus the readers attention on the consequences of this limit for relation (25). We obtain
$`\left(r_{12}^ϵr_{12}^ϵ\right)\left(\rho ^1(x)\mathrm{\Psi }_1(x,x_{\mathrm{}})\mathrm{\Psi }_2(x,x_{\mathrm{}})\right)=0`$
This boundaries condition implies that $`\rho ^1\mathrm{\Psi }_1\mathrm{\Psi }_2`$ has to be in the kernel of $`r_{12}^ϵr_{12}^ϵ`$, which imposes strong constraints on $`\mathrm{\Psi }`$.
Can we go further? If we keep in mind the equivalence of $`\rho `$ as the radial coordinate, it would seem interesting to consider the case $`\rho =0`$. It could be a way to define an algebra for spatially independent objects. Unfortunately, it leads to more difficult problems when trying to evaluate Poisson brackets of these objects. But as we shall see in the following section, such an approach is more successful when studying physical observables.
## 6 Classical observables
### 6.1 General Framework
By definition, classical observables are functionals of the phase space variables that commute with the constraints. Before finding these observables, we need some preliminary calculations. First, we have to determine the commutators between the constraints and the connection :
$`\{(x),A_x(y)\}`$ $`=`$ $`A_t(x)_x\delta (xy)[A_x(x),A_t(x)]\delta (xy)`$ (29)
$`\rho ^1(x)[\mathrm{\Phi }(x),P_t(x)]\delta (xy)`$
$`\{𝒫(x),A_x(y)\}`$ $`=`$ $`A_x(x)_x\delta (xy)`$ (30)
$`\{\mathrm{\Phi }_1(x),A_{x2}(y)\}`$ $`=`$ $`[c_{12},A_{x2}(x)]\delta (xy)+_x\delta (xy)c_{12}`$ (31)
The equivalent relations for the wave function can be found by solving the differential equation associated to one of the commutators, obtained when derivating with respect to the spatial parameter of the wave function. Thus, we deduce the following identities (on the constraint surface):
$`\{(x),\mathrm{\Psi }(y,z)\}`$ $`=`$ $`A_t(y)\mathrm{\Psi }(y,z)\delta (xy)\mathrm{\Psi }(y,z)A_t(z)\delta (xz)`$ (32)
$`\{𝒫(x),\mathrm{\Psi }(y,z)\}`$ $`=`$ $`A_x(y)\mathrm{\Psi }(y,z)\delta (xy)\mathrm{\Psi }(y,z)A_x(z)\delta (xz)`$ (33)
$`\{\mathrm{\Phi }_1(x),\mathrm{\Psi }_2(y,z)\}`$ $`=`$ $`c_{12}\mathrm{\Psi }_2(y,z)\delta (xy)\mathrm{\Psi }_2(y,z)c_{12}\delta (xz)`$ (34)
Finding quantities that have vanishing brackets with the local $`SO(2)`$ constraint is quite obvious. If we consider $`\zeta `$ defined by $`_\mu \zeta +Q_\mu \zeta =0`$, then it is straightforward to show that $`\zeta ^1(x)\mathrm{\Psi }(x,y)\zeta (y)`$ commutes with $`\mathrm{\Phi }`$. Difficulties arise when we try to find objects invariant under diffeomorphisms. If we keep in mind what was done in simpler cases, we should attempt to consider the monodromy matrix between the boundaries. We shall see that it can be achieved by using two particular values of $`\rho `$ and imposing physical boundary conditions to the solutions.
### 6.2 Vacuum solution and level one representations
We shall recall some formulae and results described in that will be helpful when dealing with boundary conditions. In particular, we shall introduce the level one representations for the affine algebra,
One of the most simple solution of the Einstein’s equations is the vacuum solution which corresponds to the case where where all P and Q fields are null (notice there is a slightly change comparing to , where $`\widehat{\sigma }`$ is also zero). The associated Lax connection belongs to the Virasoro algebra
$`A_x`$ $`=`$ $`{\displaystyle \frac{1}{2}}\rho ^1(\mathrm{\Pi }_{\widehat{\sigma }}_x\rho )E_+{\displaystyle \frac{1}{2}}\rho ^1(\mathrm{\Pi }_{\widehat{\sigma }}+_x\rho )E_{}+\mathrm{\Pi }_\rho {\displaystyle \frac{k}{2}}`$
$`A_t`$ $`=`$ $`{\displaystyle \frac{1}{2}}\rho ^1(\mathrm{\Pi }_{\widehat{\sigma }}_x\rho )E_++{\displaystyle \frac{1}{2}}\rho ^1(\mathrm{\Pi }_{\widehat{\sigma }}+_x\rho )E_{}_x\widehat{\sigma }{\displaystyle \frac{k}{2}}`$
Solving (8) with this connection, we deduce the vacuum wave function $`\mathrm{\Psi }_V`$
$`\mathrm{\Psi }_V`$ $`=`$ $`e^{\frac{1}{2}\widehat{\sigma }k}\left({\displaystyle \frac{\rho +\mathrm{\Pi }_{\widehat{\sigma }}+b}{2\rho }}\right)^{E+}\left({\displaystyle \frac{\rho +\mathrm{\Pi }_{\widehat{\sigma }}+b}{b^{}}}\right)^E`$
$`=`$ $`e^{\frac{1}{2}\widehat{\sigma }k}\left({\displaystyle \frac{\rho \mathrm{\Pi }_{\widehat{\sigma }}+c}{2\rho }}\right)^E\left({\displaystyle \frac{\rho \mathrm{\Pi }_{\widehat{\sigma }}+c}{c^{}}}\right)^{E+}`$
We shall introduce the level one representations as a tool to specify physical observables. First we introduce the two dimensional representation of $`sl(2,R)`$ involving Pauli matrices
$`T^x={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ $`T^y={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right)`$ $`T^z={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$ (42)
Notice that $`T^y`$ is the $`SO(2)`$ generator.
Let denote $`Z(\mu )`$ the field
$`Z(\mu )={\displaystyle \underset{n\mathrm{odd}}{}}p_n{\displaystyle \frac{\mu ^n}{n}}`$ $`\mathrm{with}`$ $`[p_n,p_m]=n\delta _{n+m,0}`$
and the vertex operator $`W_2(\mu )=:e^{2iZ(\mu )}:`$. We denote $`|0>`$ the vacuum of the Fock space generated by the operators $`p_n`$. The level one representations with highest weight $`\mathrm{\Lambda }_\pm `$ are defined by
$`i\mu {\displaystyle \frac{dZ(\mu )}{d\mu }}`$ $`=`$ $`{\displaystyle \underset{n\mathrm{odd}}{}}(T^z\lambda ^n)\mu ^n`$
$`\pm iW_2(\mu )`$ $`=`$ $`2{\displaystyle \underset{n\mathrm{even}}{}}(T^y\lambda ^n)\mu ^n2{\displaystyle \underset{n\mathrm{odd}}{}}(T^x\lambda ^n)\mu ^n`$
The heighest weight vectors $`|\mathrm{\Lambda }_\pm >`$ are identified with $`|0>`$. The Virasoro generators are represented by
$`L_n`$ $`=`$ $`{\displaystyle \frac{1}{8\pi i}}{\displaystyle _𝒞}𝑑\mu \mu ^{2n+1}:(i_\mu Z)^2:+{\displaystyle \frac{1}{16}}\delta _{n,0}`$
where $`𝒞`$ is a contour around zero.
In the next subsection, we shall need to conjugate elements of the algebra by the vacuum wave function. We gather here all needed formulae. The case of the Virasoro algebra elements $`E_+`$ and $`E_{}`$ is quite obvious
$`D^{E_+}E_{}D^{E_+}=DE_{}+(D1)E_+`$ $`\mathrm{and}`$ $`B^E_{}E_+B^E_{}=BE_++(B1)E_{}`$
Operators $`\mu ^1W_2(\mu )`$ and $`_\mu Z`$ are primary fields of weight 1. Under a diffeomorphism $`\mu F(\mu )`$, they transform as
$`\mu ^1W_2(\mu )`$ $``$ $`{\displaystyle \frac{\mu }{2F^2}}\left(_\mu F^2\right)W_2(F(\mu ))`$
$`_\mu Z`$ $``$ $`{\displaystyle \frac{\mu ^2}{2F^2}}\left(_\mu F^2\right)(_\mu Z)(F(\mu ))`$
It can be proved that conjugations $`AD^{E_+}\mu ^1W_2(\mu )D^{E+}`$ and $`B^E_{}\mu ^1W_2(\mu )B^E_{}`$, are associated to the following diffeomorphisms
$`F_+^2(\mu )`$ $`=`$ $`{\displaystyle \frac{\mu ^2}{\mu ^2+(1\mu ^2)D}}`$
$`F_{}^2(\mu )`$ $`=`$ $`1+(\mu ^21)B`$
### 6.3 Boundaries conditions and physical observables
Let give a brief sketch of the strategy we shall use. Following the idea proposed in , we shall consider $`\mathrm{\Psi }`$ between the points $`x_0`$ where $`\rho =0`$, and $`x_{\mathrm{}}`$ where $`\rho \mathrm{}`$ . We shall need to define our phase space by specifying the behavior of the other fields in these limits. As we will see, it is not possible to obtain physical observables by only imposing boundaries conditions on $`\mathrm{\Psi }`$. We shall construct physical observables in the form $`M_0^1(x_0)\mathrm{\Psi }(x_0,x_{\mathrm{}})M_{\mathrm{}}(x_{\mathrm{}})`$ in such a way the contribution of $`M_0`$ and $`M_{\mathrm{}}`$ to the Poisson brackets with the constraints eliminates the unwanted terms.
First we shall look at the case where $`\rho \mathrm{}`$, with the picture that in this limit $`\rho `$ becomes equivalent to the usual radial coordinate (assuming it happens when $`xx_{\mathrm{}}`$). It seems physically reasonable to restrict our phase space to solutions which tend asymptotically to the flat space solution when $`\rho \mathrm{}`$. This solution corresponds to the Minkowski metric element expressed in cylindrical coordinates $`ds^2=dt^2+d\rho ^2+\rho ^2d\theta ^2+dz^2`$. We shall take
$`\mathrm{\Pi }_\rho =\mathrm{\Pi }_{\widehat{\sigma }}=0,`$ $`\widehat{\sigma }={\displaystyle \frac{1}{4}}\mathrm{ln}\rho ,`$ $`P_x=\rho ^1\sqrt{2}T^z`$
and all other components of $`P`$ and $`Q`$ equal to zero. The components of the Lax connection are
$`(A_f)_x`$ $`=`$ $`\rho ^1\left({\displaystyle \frac{1}{2}}(E_+E_{})+{\displaystyle \frac{\sqrt{2}}{2}}T^z(\lambda +\lambda ^1)\right)`$ (43)
$`(A_f)_t`$ $`=`$ $`\rho ^1\left({\displaystyle \frac{1}{2}}(E_++E_{})+{\displaystyle \frac{\sqrt{2}}{2}}T^z(\lambda \lambda ^1)+{\displaystyle \frac{k}{8}}\right)`$ (44)
The wave function which is solution of (8) with the above connection can be written as
$`\mathrm{\Psi }_f`$ $`=`$ $`h_{}\left({\displaystyle \frac{2\rho }{\rho +1}}\right)^E_{}\left({\displaystyle \frac{1}{2}}(\rho +1)\right)^{E_+}\rho ^{\frac{k}{2}}h_{}^1`$
$`=`$ $`h_+\left({\displaystyle \frac{2\rho }{\rho +1}}\right)^{E_+}\left({\displaystyle \frac{1}{2}}(\rho +1)\right)^E_{}\rho ^{\frac{k}{2}}h_+^1`$
with $`h_\pm =\mathrm{exp}\left(\sqrt{2}T^z\mathrm{ln}\left(\frac{1+\lambda ^{\pm 1}}{1\lambda ^{\pm 1}}\right)\right)`$. If we look at the Lax connection, we see that it is proportional to $`\rho ^1`$. Thus, if we impose that the wave function $`\mathrm{\Psi }(x)`$ tends asymptotically toward $`\mathrm{\Psi }_f(x)`$ when $`xx_{\mathrm{}}`$, equations (32) and (33) become
$`\{(x),\mathrm{\Psi }(y,x_{\mathrm{}})\}`$ $`=`$ $`A_t(y)\mathrm{\Psi }(y,x_{\mathrm{}})\delta (xy)`$ (45)
$`\{𝒫(x),\mathrm{\Psi }(y,x_{\mathrm{}})\}`$ $`=`$ $`A_x(y)\mathrm{\Psi }(y,x_{\mathrm{}})\delta (xy)`$ (46)
Now let us consider the more tricky case of $`\rho =0`$. We shall suppose there is a point $`x_0`$ such that
$`\rho (x_0)=0`$ $`\mathrm{and}`$ $`\mathrm{\Pi }_{\widehat{\sigma }}(x_0)=0`$
In the picture of cylindrical symmetry, it means that close to the symmetry axis, $`\rho `$ can be identified with the usual radial coordinate. We assume $`P`$ behaves like $`\rho ^1`$ when $`xx_0`$. Actually, if we admit these quantities have a $`\rho `$ power law behavior in this limit (which seems physically reasonable), then, using equations of motion (2, 3, 4), we can show that it is the only solution. But the monodromy matrix $`\mathrm{\Psi }(x_0,x_{\mathrm{}})`$ is not an observable because the Lax connection diverges when $`\rho 0`$. To avoid this problem, we apply the technique explained at the beginning of this subsection: we multiply $`\mathrm{\Psi }(x,x_{\mathrm{}})`$ by $`M_0^1(x)`$ with $`M_0`$ equal to the vacuum wave function $`\mathrm{\Psi }_V`$ whose form in the limit $`xx_0`$ becomes
$`\mathrm{\Psi }_V`$ $`e^{\frac{1}{2}\widehat{\sigma }k}\left({\displaystyle \frac{b}{2\rho }}\right)^{E+}\left(b^{}\right)^E`$ $`e^{\frac{1}{2}\widehat{\sigma }k}\left({\displaystyle \frac{c}{2\rho }}\right)^E\left(c^{}\right)^{E+}`$ (47)
(relabeling the prime constants to simplify formulae).
The linear equations satisfied by $`\mathrm{\Psi }_V^1(x)\mathrm{\Psi }(x,x_{\mathrm{}})`$ are of type
$`\left(\mathrm{\Psi }_V^1(x)\mathrm{\Psi }(x,x_{\mathrm{}})\right)`$ $`=`$ $`\left(\mathrm{\Psi }_V(x)^1\stackrel{~}{A}(x)\mathrm{\Psi }_V(x)\right)\mathrm{\Psi }_V(x)^1\mathrm{\Psi }(x,x_{\mathrm{}})`$
where $`\stackrel{~}{A}`$ is the Kac-Moody algebra part of the connection (the contribution of $`\mathrm{\Psi }_V`$ has canceled Virasoro and central extension parts). Poisson brackets of $`\mathrm{\Psi }_V^1(x)\mathrm{\Psi }(x,x_{\mathrm{}})`$ with $``$ and $`𝒫`$ are obtained by substituting $`\mathrm{\Psi }_V^1\stackrel{~}{A}\mathrm{\Psi }_V`$ to $`A`$ in formulae (45) and (46). To evaluate $`\mathrm{\Psi }_V^1\stackrel{~}{A}\mathrm{\Psi }_V`$, we use the level one representation described in the previous subsection and the second form of (47). The diffeomorphism associated to this conjugation is $`F^2(\mu )=\frac{2\rho +(\mu ^21)c}{2\rho +(\mu ^21)c(1c^{})}`$. For example, for one of the components of $`A`$, we have
$`\mathrm{\Psi }_V^1(P_t)_xT^x\lambda \mathrm{\Psi }_V`$ $`=`$ $`{\displaystyle \frac{(P_t)_x}{4\pi }}\mathrm{\Psi }_V^1\left({\displaystyle _𝒞}𝑑\mu \mu W_2(\mu )\right)\mathrm{\Psi }_V`$
$`=`$ $`{\displaystyle \frac{(P_t)_x}{4\pi }}\left(\rho {\displaystyle \frac{c^{}}{c(c^{}1)^2}}W_2({\displaystyle \frac{1}{1c^{}}}){\displaystyle _𝒞}𝑑\mu {\displaystyle \frac{\mu ^3}{(\mu ^21)^2}}+𝒪(\rho ^2)\right)`$
$`=`$ $`(P_t)_x\left(0+𝒪(\rho ^2)\right)=𝒪(\rho )\mathrm{as}\rho 0`$
The last equation follows from the fact that $`P`$ is proportional to $`\rho ^1`$ as $`xx_0`$ and that the contour $`𝒞`$ around $`0`$ can be chosen as small as we want, such that there is no pole contribution to the integral. Doing these calculations in the level one representation provides a way to give meaning to quantities like $`W_2(\frac{1}{1c^{}})`$, which would be ambiguous in the abstract Kac-Moody algebra.
Identical results can be obtained for the other components. Thus $`\mathrm{\Psi }_V^1(x_0)\stackrel{~}{A}(x_0)\mathrm{\Psi }_V(x_0)`$ is equal to zero, whereas $`A(x)`$ was divergent when $`xx_0`$, and $`\mathrm{\Psi }_V^1(x_0)\mathrm{\Psi }(x_0,x_{\mathrm{}})`$ has null Poisson brackets with the generators of diffeomorphisms.
Defining $`SO(2)`$ gauge-invariant object from $`\mathrm{\Psi }_V^1(x_0)\mathrm{\Psi }(x_0,x_{\mathrm{}})`$ is quite straightforward . Using the technique presented in subsection 6.1 and the fact that $`so(2)`$ elements commute with $`\mathrm{\Psi }_V`$, we can show that the quantity
$`\stackrel{~}{\mathrm{\Psi }}(x_0,x_{\mathrm{}})`$ $`=`$ $`\zeta ^1(x_0)\mathrm{\Psi }_V^1(x_0)\mathrm{\Psi }(x_0,x_{\mathrm{}})\zeta (x_{\mathrm{}})`$ (48)
still has vanishing brackets with $``$ and $`𝒫`$, and is moreover $`SO(2)`$ gauge-invariant. It proves that $`\stackrel{~}{\mathrm{\Psi }}(x_0,x_{\mathrm{}})`$ is a physical observable. Note that these are operators (ie infinite dimensional matrices) acting on the level one representation of $`sl(2,R)`$. They thus provide an infinite set of physical observables.
To summarize: We have supposed the two following boundary conditions: 1) when $`\rho `$ goes to infinity, the wave function tends asymptotically to the flat space wave function; 2) we can find a point $`x_0`$ where $`\rho `$ and $`\mathrm{\Pi }_{\widehat{\sigma }}`$ are null. Notice that these conditions are fulfilled in concrete examples like cylindrical gravitational waves (see e.g. ). Using these hypotheses and the level one representations, we have shown that $`\zeta ^1(x_0)\mathrm{\Psi }_V^1(x_0)\mathrm{\Psi }(x_0,x_{\mathrm{}})\zeta (x_{\mathrm{}})`$ generates an infinite set of classical physical observables. It remains to decipher the Poisson bracket algebra they generate which should be closer to those of the Toda’s theories.
## 7 Appendix A:
The aim of this appendix is to give more details about the Poisson brackets of the connection. First, let us write developed formulae for the $`r`$ and $`s`$matrices :
* plus convention $`\left(|\lambda _1|<|\lambda _2|\right)`$:
$`r_{12}^+`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1\lambda _1^2)(1\lambda _2^2){\displaystyle \underset{n0}{}}\left({\displaystyle \frac{\lambda _1}{\lambda _2}}\right)^{2n}c_{12}{\displaystyle \frac{1}{2}}(\lambda _2\lambda _2^1)^2{\displaystyle \underset{n0}{}}\left({\displaystyle \frac{\lambda _1}{\lambda _2}}\right)^{2n+1}d_{12}`$
$`{\displaystyle \frac{1}{2}}E_+k{\displaystyle \frac{1}{4}}k\left(E_++E_{}\right)`$
$`s_{12}^+`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1\lambda _1^2)(1\lambda _2^2){\displaystyle \underset{n0}{}}\left({\displaystyle \frac{\lambda _1}{\lambda _2}}\right)^{2n}c_{12}{\displaystyle \frac{1}{2}}(\lambda _1\lambda _1^1)^2{\displaystyle \underset{n0}{}}\left({\displaystyle \frac{\lambda _1}{\lambda _2}}\right)^{2n+1}d_{12}`$
$`{\displaystyle \frac{1}{2}}kE_{}{\displaystyle \frac{1}{4}}\left(E_++E_{}\right)k`$
* minus convention $`\left(|\lambda _1|>|\lambda _2|\right)`$:
$`r_{12}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1\lambda _1^2)(1\lambda _2^2){\displaystyle \underset{n0}{}}\left({\displaystyle \frac{\lambda _2}{\lambda _1}}\right)^{2n}c_{12}+{\displaystyle \frac{1}{2}}(\lambda _2\lambda _2^1)^2{\displaystyle \underset{n0}{}}\left({\displaystyle \frac{\lambda _2}{\lambda _1}}\right)^{2n+1}d_{12}`$
$`+{\displaystyle \frac{1}{2}}E_{}k+{\displaystyle \frac{1}{4}}k\left(E_++E_{}\right)`$
$`s_{12}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1\lambda _1^2)(1\lambda _2^2){\displaystyle \underset{n0}{}}\left({\displaystyle \frac{\lambda _2}{\lambda _1}}\right)^{2n}c_{12}+{\displaystyle \frac{1}{2}}(\lambda _1\lambda _1^1)^2{\displaystyle \underset{n0}{}}\left({\displaystyle \frac{\lambda _2}{\lambda _1}}\right)^{2n+1}d_{12}`$
$`+{\displaystyle \frac{1}{2}}kE_++{\displaystyle \frac{1}{4}}\left(E_++E_{}\right)k`$
Notice that the differences $`r_{12}^+r_{12}^{}`$ and $`s_{12}^+s_{12}^{}`$ are equal. But in contrary to the usual cases, they are not proportional to the Casimir tensor (see equation (20)).
In order to demystify this formulae, we will sketch the way we have obtained them. First, we consider only the loop part of the algebra and the ultralocal contribution. The more general expression for the two $`r`$matrices we can take, is of type $`f(\lambda _1,\lambda _2)c_{12}+g(\lambda _1,\lambda _2)d_{12}`$. Comparing with the raw formula of the Poisson brackets, we deduce the loop part of (17) and (18)
$`r_{12}=f(\lambda _1,\lambda _2)c_{12}+g(\lambda _1,\lambda _2)d_{12}`$ $`\mathrm{and}`$ $`s_{12}=f(\lambda _1,\lambda _2)c_{12}g(\lambda _2,\lambda _1)d_{12}`$ (49)
with
$`f(\lambda _1,\lambda _2)={\displaystyle \frac{1}{2}}{\displaystyle \frac{(1\lambda _1^2)(1\lambda _2^2)}{\lambda _1^2\lambda _2^2}}`$ $`\mathrm{and}`$ $`g(\lambda _1,\lambda _2)={\displaystyle \frac{1}{2}}{\displaystyle \frac{\lambda _1\lambda _2^1(1\lambda _2^2)^2}{\lambda _1^2\lambda _2^2}}`$
These rational functions verify some non-trivial algebraic relations that are helpful when dealing with Jacobi identity
$`g(\lambda _1,\lambda _3)g(\lambda _3,\lambda _2)+f(\lambda _2,\lambda _3)g(\lambda _1,\lambda _2)+f(\lambda _3,\lambda _1)g(\lambda _1,\lambda _2)`$ $`=`$ $`0`$
$`g(\lambda _1,\lambda _2)g(\lambda _3,\lambda _2)f(\lambda _2,\lambda _3)g(\lambda _1,\lambda _3)f(\lambda _1,\lambda _2)g(\lambda _1,\lambda _3)`$ $`=`$ $`0`$
$`g(\lambda _1,\lambda _2)g(\lambda _1,\lambda _3)f(\lambda _1,\lambda _3)g(\lambda _3,\lambda _2)f(\lambda _1,\lambda _2)g(\lambda _2,\lambda _3)`$ $`=`$ $`(\lambda _2\lambda _2^1)(\lambda _3\lambda _3^1)`$
Now we have to add the central extension (which impose to choose a convention for the previous equations). We can easily see that we need also to introduce the Virasoro algebra. The only possible and non-trivial terms are $`E_\pm k`$ and $`kE_\pm `$. The goal we want to reach is to include all terms whose variables are different from the dilaton and its spatial derivative, into the two commutators $`[r_{12}^ϵ,A_1(x)]+[s_{12}^ϵ,A_2(x)]`$. The first miracle is that it can be achieved. Thus, we fix one part of the $`r`$matrices on the Virasoro algebra (remember that $`k`$ is in the center of the algebra, so $`[kE_\pm ,A_1(x)]=0`$ and so on). The last step consists in fixing the rest of the $`r`$matrices in such a way that we obtain a compact form for (4.1). And here we have a second miracle: it can be done for the ultralocal and the non-ultralocal part simultaneously.
We see this formula works due to non-trivial cancelations. In one sense, it’s a proof of the validity of our calculation and a hint of deeper algebraic structure of the problem.
## 8 Appendix B:
Here we will give some hints for the proof of the Jacobi identity. First of all, let recall the expression of eq.(4.2) :
$`\{A_1(x),\{A_2(y),A_3(z)\}\}+\mathrm{perm}.`$ $`=`$ $`\rho ^2(x)\delta (xy)\delta (yz)([A_1(x),A_{123}]+\mathrm{perm}.)`$
$`+`$ $`\left(_x\rho ^2\right)\delta (xy)\delta (yz)(B_{123}+\mathrm{perm}.)`$
$`+`$ $`\rho ^3(x)\delta (xy)\delta (yz)([\mathrm{\Phi }_1(x),C_{123}]+\mathrm{perm}.)`$
$`+`$ $`\rho ^2(x)\delta (yz)_x\delta (xy)D_{123}+\mathrm{perm}.`$
with the following values for the coefficients
$`A_{123}`$ $`=`$ $`[r_{12}^{ϵ_1},r_{23}^{ϵ_2}]+[s_{23}^{ϵ_2},s_{31}^{ϵ_3}]+[s_{31}^{ϵ_3},r_{12}^{ϵ_1}]{\displaystyle \frac{1}{2}}k_2s_{31}^{ϵ_3}{\displaystyle \frac{1}{2}}k_3r_{12}^{ϵ_1}{\displaystyle \frac{1}{4}}[U_{23},c_{12}]`$
$`B_{123}`$ $`=`$ $`[s_{23}^{ϵ_2},s_{31}^{ϵ_3}][r_{23}^{ϵ_2},r_{12}^{ϵ_1}]+{\displaystyle \frac{1}{2}}[r_{23}^{ϵ_2},s_{12}^{ϵ_1}]{\displaystyle \frac{1}{2}}[s_{23}^{ϵ_2},r_{31}^{ϵ_3}]{\displaystyle \frac{1}{4}}[U_{23},c_{12}]`$
$`C_{123}`$ $`=`$ $`{\displaystyle \frac{1}{4}}[r_{23}^{ϵ_2},U_{12}]+{\displaystyle \frac{1}{4}}[s_{23}^{ϵ_2},U_{31}]+{\displaystyle \frac{1}{8}}k_3U_{12}{\displaystyle \frac{1}{8}}k_2U_{13}`$
$`D_{123}`$ $`=`$ $`[s_{23}^{ϵ_2},s_{31}^{ϵ_3}][r_{23}^{ϵ_2},r_{12}^{ϵ_1}]+[r_{23}^{ϵ_2},s_{12}^{ϵ_1}][s_{23}^{ϵ_2},r_{31}^{ϵ_3}]{\displaystyle \frac{1}{4}}[U_{23},c_{12}]`$
$`+{\displaystyle \frac{1}{4}}k_3\left(s_{12}^{ϵ_1}r_{12}^{ϵ_1}\right){\displaystyle \frac{1}{4}}k_2\left(s_{31}^{ϵ_3}r_{31}^{ϵ_3}\right)+{\displaystyle \frac{1}{4}}k_1\left(s_{23}^{ϵ_2}+r_{23}^{ϵ_2}\right)`$
To obtain this expression, we need Poisson brackets of connection with the constraint $`\mathrm{\Phi }`$ and the dilaton (see eq.(31) and (26)).
Showing Jacobi identity with (4.2) is not obvious, because the three $`\delta \delta `$ distributions and $`\delta \delta `$ are not linearly independent. So, we have to express one of $`\delta \delta `$ with respect to the other distributions. It leads to the following formula :
$`\{A_1(x),\{A_2(y),A_3(z)\}\}+\mathrm{perm}.`$ $`=`$ $`\rho ^2(x)\delta (xy)\delta (yz)([A_1(x),A_{123}]+\mathrm{perm}.)`$
$`+`$ $`\left(_x\rho ^2\right)\delta (xy)\delta (yz)\left(B_{123}+B_{231}+B_{312}2D_{123}\right)`$
$`+`$ $`\rho ^3(x)\delta (xy)\delta (yz)([\mathrm{\Phi }_1(x),C_{123}]+\mathrm{perm}.)`$
$`+`$ $`\rho ^2(y)\delta (zx)_y\delta (yz)\left(D_{123}D_{123}\right)`$
$`+`$ $`\rho ^2(z)\delta (xy)_z\delta (yz)\left(D_{312}D_{123}\right)`$
It can be easily shown that we have the relations
$`B_{123}+B_{231}+B_{312}2D_{123}=A_{231}+A_{312}A_{123}`$
$`D_{123}D_{231}=A_{123}A_{231}`$
$`D_{123}D_{312}=A_{123}A_{312}`$
Thus, our problem is entirely expressed in terms of $`A_{123}`$ and $`C_{123}`$. As announced previously, these coefficients are invariant under cyclic permutations and equal to zero :
$`A_{123}=A_{231}=A_{312}=0`$ $`\mathrm{and}`$ $`C_{123}=C_{231}=C_{312}=0`$
If the calculations for the $`C`$ coefficients are rather easy, those for the $`A`$ ones are more tedious. In particular, when dealing with terms of type $`kd`$, the choice of the convention has to be coherent. Another point is the fundamental Poisson brackets are not to be taken on the constraint surface. Else the $`[U,c]`$ term disappears, and the equation for $`A`$ on the loop part of the algebra is no longer verified.
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# Superluminal Caustic is Just a Common Misconception: A Comment on astro-ph/0001199 by Zheng Zheng and Andrew Gould
## Acknowledgments
It is our pleasure to express our gratitude to Clara Bennett for the copy of “CRC Concise Encyclopedia of Mathematics” given to the author during the last winter solstice. It is a great gift to be snowed in with.
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# Warm inflation with coupled thermal quantum fluctuations: a new semiclassical approach*footnote **footnote *This essay received an “honorable mention” in the Annual Essay Competition of the Gravity Research Foundation for the year 2000.
## Abstract
I consider a new semiclassical expansion for the inflaton field in the framework of warm inflation scenario. The fluctuations of the matter field are considered as thermally coupled with the particles of the thermal bath. This coupling parameter depends on the temperature of the bath. The power spectrum remains invariant under this new semiclassical expansion for the inflaton. I find that the thermal component of the amplitude for the primordial field fluctuations should be very small at the end of inflation.
PACS number(s): 98.80.Cq.; 04.62.+v
Inflation is needed because it solves the horizon, flatness, and monopole problems of the very early universe. This theory also provides a mechanism for the creation of primordial density fluctuations. The differential microwave radiometer (DMR) on the Cosmic Background Explorer (COBE) has made the first direct probe of the initial density perturbations through detection of the temperature anisotropies in the cosmic background radiation (CBR). The results are consistent with the scaling spectrum given by the inflationary model. For inflation the simplest assumption is that there are two scales: a long - time, long - distance scale associated with the vacuum energy dynamics and a single short - time, short - distance scale associated with a random force component. The Hubble time during inflation, $`1/H`$, appropriately separates the two regimes.
Quantum fluctuations and thermal fluctuations of matter fields can play a prominent role in inflationary cosmology. During inflation, vacuum fluctuations on scales smaller than the size of the horizon are magnified into classical perturbations on scales bigger than the Hubble radius. The classical perturbations can lead to classical curvature of spacetime and energy density perturbations after inflation. These density perturbations should be responsible for the formation of large - scale structure of the universe, as well as, the anisotropies in the cosmic microwave background. Structure formation scenarios, can receive important restrictions based on the measured $`\delta T_r/T_r=1.1\times 10^5`$. According to the standard inflationary model, the formation of large - scale structure in the universe has its origin in the growth of primordial inhomogeneities in the matter distribution.
In a previous work I studied a stochastic approach for the backreaction of the metric produced by the fluctuations of the matter field in the framework of warm inflation scenario. In that work I considered a thermal coupling between the scalar field and the fields in the thermal bath. The aim of this letter is to consider a coupled thermal interaction between the matter field fluctuations and the fields in the thermal bath to study the matter field fluctuations in the warm inflation scenario.
Warm inflation takes into account separately, the matter and radiation energy fluctuations. In this scenario the field $`\phi `$ interacts with the particles of a thermal bath with a mean temperature ($`T_r`$) smaller than the Grand Unified Theories (GUT) critical temperature: $`T_r<T_{GUT}10^{15}`$ GeV. This scenario was introduced by A. Berera. In the warm inflation era, the kinetic component of energy, $`\rho _{kin}`$, must be smaller than the vacuum energy, which is given by the effective potential $`V(\phi )`$
$$\rho (\phi )\rho _mV(\phi )\rho _{kin},$$
(1)
where $`\rho _{kin}(\phi )=\rho _r(\phi )+\dot{\phi }^2/2`$, and the radiation energy density is $`\rho _r(\phi )=\frac{\tau (\phi )}{8H(\phi )}\dot{\phi }^2`$. Here, $`\phi `$ is a scalar field of matter, $`\tau (\phi )`$ is an effective friction parameter that represent the interaction of the matter field with other fields of the thermal bath.
In this work I consider a new semiclassical expansion for the inflaton field $`\phi `$
$$\phi (\stackrel{}{x},t)=\varphi _c(t)+\alpha (t)\psi (\stackrel{}{x},t),$$
(2)
where $`\varphi _c(t)=<0|\phi (\stackrel{}{x},t)|0>`$ and $`<0|\psi (\stackrel{}{x},t)|0>=<0|\dot{\psi }(\stackrel{}{x},t)|0>=0`$. Here, $`|0>`$ is the vacuum state. Here, $`\alpha (t)`$ is a dimensionless time - dependent function that characterize the thermal coupling between the inflaton field fluctuations and the fields in the thermal bath. Furthermore, the classical time - dependent field $`\varphi _c(t)`$ gives the instantaneous background. In principle, a permanent or temporary coupling of the scalar field with other fields might also lead to dissipative processis producing entropy at different eras of the cosmic evolution. We introduce the function $`\alpha (t)`$ with the aim to make an estimation of how important is the thermal amplitude of the fluctuations at the end of inflation. I will consider $`\alpha =\left(\frac{T_r(t)}{M}\right)^\beta `$, where $`T_r(t)`$ is the time - dependent radiation temperature and $`M10^{15}`$ GeV is the GUT mass.
The dynamics of the classical field $`\varphi _c(t)`$ was obtained in previous works
$$\ddot{\varphi }_c+\left[3H_c+\tau _c\right]\dot{\varphi }_c+V^{}(\varphi _c)=0,$$
(3)
where $`H_cH(\varphi _c)=\dot{a}/a`$, $`\tau _c\tau (\varphi _c)`$ and $`V^{}(\varphi _c)\frac{dV(\phi )}{d\phi }|_{\varphi _c}`$. Furthermore $`\dot{\varphi }_c=\frac{M_p^2}{4\pi }H_c^{}\left(1+\frac{\tau _c}{3H_c}\right)^1`$ and the classical effective potential \[see for example ref. ()\] can be obtainded from the homogeneous Friedmann equation $`H_c^2(\varphi _c)=\frac{4\pi }{3M_p^2}\left[\left(1+\tau _c/(4H_c)\right)\dot{\varphi }_c^2+2V(\varphi _c)\right]`$:
$$V(\varphi _c)=\frac{M_p^2}{8\pi }\left[H_c^2\frac{M_p^2}{12\pi }(H_c^{})^2\left(1+\frac{\tau _c}{4H_c}\right)\left(1+\frac{\tau _c}{3H_c}\right)^2\right].$$
(4)
The radiation energy density of the background is
$$\rho _r[\varphi _c(t)]\frac{\tau _c}{8H_c}\left(\frac{M_p^2}{4\pi }\right)^2(H_c^{})^2\left(1+\frac{\tau _c}{3H_c}\right)^2,$$
(5)
and the temperature of the bath is $`T_r\rho _r^{1/4}[\varphi _c(t)]`$. The spatially homogeneous matter energy density is given by
$$\rho _m(\varphi _c)=\frac{\dot{\varphi }_c^2}{2}+V(\varphi _c),$$
(6)
for $`V(\varphi _c)\frac{\dot{\varphi }_c^2}{2}`$.
I study the perturbations on a globally flat, homogeneous and isotropic spacetime, described by a flat Friedmann - Robertson - Walker (FRW) metric $`ds^2=dt^2+a^2d\stackrel{}{x}^2`$. The equation for the quantum perturbations $`\psi `$, with the semiclassical expansion (2), is
$`\ddot{\psi }`$ $`+`$ $`\left[2{\displaystyle \frac{\dot{\alpha }}{\alpha }}+(3H_c+\tau _c)\right]\dot{\psi }{\displaystyle \frac{1}{a^2}}^2\psi `$ (7)
$`+`$ $`\left[(3H_c+\tau _c){\displaystyle \frac{\dot{\alpha }}{\alpha }}+{\displaystyle \frac{\ddot{\alpha }}{\alpha }}+V^{\prime \prime }(\varphi _c)\right]\psi =0.`$ (8)
To simplify the structure of this equation, we can introduce the map $`\chi =e^{3/2{\scriptscriptstyle (H_c+\tau _c/3+{\scriptscriptstyle \frac{2\dot{\alpha }}{3\alpha }})𝑑t}}\psi `$
$$\ddot{\chi }\frac{1}{a^2}^2\chi \mu ^2\chi =0,$$
(9)
where $`\mu ^2(t)=\frac{k_o^2(t)}{a^2}`$ is a effective squared time - dependent parameter of mass
$$\mu ^2(t)=\frac{9}{4}\left(H_c+\tau _c/3\right)^2V^{\prime \prime }(\varphi _c)+\frac{3}{2}\left(\dot{H}_c+\frac{\dot{\tau }_c}{3}\right).$$
(10)
Note that this parameter does not depends on $`\alpha `$. For $`H_c+\tau _c/3+\frac{2\dot{\alpha }}{3\alpha }=0`$, one obtains the particular map $`\psi _{cg}=\chi _{cg}`$. Furthermore, for $`\tau _c/3+2\dot{\alpha }/(3\alpha )=0`$ the resulting map is $`\psi _{cg}=e^{3/2{\scriptscriptstyle H_c𝑑t}}\chi _{cg}`$, which coincides with the map developed in standard inflation with the standard semiclassical expansion $`\phi =\varphi _c+\varphi `$.
The backreaction of the metric with the quantum fluctuations of the matter field introduces an effective curvature ($`K`$) on the globally flat FRW background metric, which appears in the semiclassical Friedmann equation
$$H_c^2+\frac{K}{a^2}=\frac{8\pi }{3M_p^2}E\left|\rho _m(\phi )+\rho _r(\phi )\right|E,$$
(11)
where
$`{\displaystyle \frac{K}{a^2}}`$ $`=`$ $`{\displaystyle \frac{8\pi }{3M_p^2}}[(1+{\displaystyle \frac{\tau _c}{8H_c}})({\displaystyle \frac{\dot{\alpha }^2}{2}}<\psi ^2>+{\displaystyle \frac{\alpha ^2}{2}}<\dot{\psi }^2>`$ (12)
$`+`$ $`\alpha \dot{\alpha }<\psi \dot{\psi }>)+{\displaystyle \frac{\alpha ^2}{a^2}}\left(\stackrel{}{}\psi \right)^2+{\displaystyle \frac{V^{\prime \prime }}{2}}\alpha ^2\psi ^2].`$ (13)
To study the consequences of this approach, we can consider a power - law expansion of the universe. In this model the scale factor is $`a(t/t_o)^p`$, and $`H_c(t)=p/t`$. The effective classical potential and the radiation energy density are given by
$`V[\varphi _c(t)]`$ $`=`$ $`{\displaystyle \frac{3M_p^2}{8\pi }}t^2`$ (14)
$`\times `$ $`\left[p^2{\displaystyle \frac{M_p^2}{2\pi m^2}}p^2\left(1+{\displaystyle \frac{\tau _ct}{4p}}\right)\left(1+{\displaystyle \frac{\tau _ct}{3p}}\right)^2\right],`$ (15)
$`\rho _r[\varphi _c(t)]`$ $`=`$ $`\left({\displaystyle \frac{p\tau _ct}{32}}\right)\left({\displaystyle \frac{M_p^2}{\pi m}}\right)^2\left(1+{\displaystyle \frac{\tau _ct}{3p}}\right)^2t^2,`$ (16)
where $`p/t=H_oe^{\varphi _c(t)/m}`$.
The redefined matter field perturbations on the infrared sector \[$`k^2k_o^2(t)`$\] takes into account only the modes much bigger than the size of the horizon. It can be written as a Fourier expansion in terms of the modes $`\chi _k(\stackrel{}{x},t)=e^{i\stackrel{}{k}.\stackrel{}{x}}\xi _k(t)`$
$$\chi _{cg}=\frac{1}{(2\pi )^{3/2}}d^3k\theta (ϵk_ok)\left[a_k\chi _k+a_k^{}\chi _k^{}\right],$$
(17)
where $`a_k`$ and $`a_k^{}`$ are the annihilation and creation operators with commutation relations $`[a_k,a_k^{}^{}]=\delta ^{(3)}(\stackrel{}{k}\stackrel{}{k^{}})`$ and $`[a_k^{},a_k^{}^{}]`$ $`=`$ $`[a_k,a_k^{}]=0`$. The dimensionless constant $`ϵ=k/k_o1`$ is introduced to take into account only the modes with wavelengths much bigger than the size of the horizon. As was showed in a previous work, the redefined fluctuations $`\chi _{cg}`$ are classical in the infrared sector. This implies that $`\xi _k\dot{\xi }_k^{}\dot{\xi }_k\xi _k^{}0`$.
To make a calculation in the power - law inflation model I consider $`\alpha =[T_r(t)/M]^\beta `$ (with $`\beta 0`$ and $`\alpha <1`$), and $`\tau _c=\gamma (p/t)`$. Here, $`M10^{15}`$ GeV, is the GUT mass. However, the form chosen for the function $`\alpha `$ must come from the analysis of some quantum field theoretical description of the coupling between the scalar field and the radiation fields which would result in the introduction of an affective temperature dependent parameter in the expansion of the field. The effective squared parameter of mass (10), for $`\rho _r=\frac{\pi ^2}{30}N(T_r)T_r^4`$ — where $`N(T_r)`$ is the number of relativistic degrees of freedom at temperature $`T_r`$ — is given by
$`\mu ^2(t)`$ $`=`$ $`{\displaystyle \frac{t^2}{12}}[3p^2(9+6\gamma +\gamma ^2)6p(3+\gamma )+36\gamma +144`$ (18)
$``$ $`{\displaystyle \frac{m^2}{M_p^2}}(288\pi +192\pi \gamma +32\pi \gamma ^2)].`$ (19)
The effective potential and the homogeneous component of the radiation energy density are, in this case
$`V(\varphi _c)`$ $`=`$ $`{\displaystyle \frac{3M_p^2}{8\pi }}H_o^2e^{2\varphi _c/m}\left[1{\displaystyle \frac{M_p^2}{2\pi m^2}}\left(1+{\displaystyle \frac{\gamma }{4}}\right)\left(1+{\displaystyle \frac{\gamma }{3}}\right)^2\right],`$ (20)
$`\rho _r(\varphi _c)`$ $`=`$ $`{\displaystyle \frac{\gamma }{32}}\left(1+{\displaystyle \frac{\gamma }{3}}\right)^2\left({\displaystyle \frac{M_p^2}{\pi m}}\right)^2H_o^2e^{2\varphi _c/m},`$ (21)
which increase exponentially on $`\varphi _c`$. Note that $`\dot{\rho }_r<0`$, which means that $`\dot{\alpha }<0`$.
The equation of motion for the time - dependent modes $`\xi _k`$ is
$$\ddot{\xi }_k+\left[\frac{k^2}{a^2}\mu ^2(t)\right]\xi _k=0,$$
(22)
where the squared time - dependent parameter of mass is given by eq. (19). As $`\mu (t)`$ is independent of $`\alpha `$ \[see eq. (10)\], there are no changes in the dynamics of the functions $`\xi _k`$. The asymptotic solution for eq. (22) (for large $`p`$-values), is
$$\xi _k(t)\frac{\sqrt{t/t_o}}{2\sqrt{2\pi }}\mathrm{\Gamma }(\nu )\left[\frac{k(t/t_o)^{1p}}{2(p1)H_o}\right]^\nu ,$$
(23)
where $`\mathrm{\Gamma }(\nu )`$ is the gamma function, $`\nu =\frac{1}{2(p1)}\sqrt{1+4L^2}`$ and
$`L^2`$ $`=`$ $`{\displaystyle \frac{1}{12}}[3p^2(9+6\gamma +\gamma ^2)6p(3+\gamma )`$ (24)
$`+`$ $`144{\displaystyle \frac{m^2}{M_p^2}}(288\pi +192\pi \gamma +32\pi \gamma ^2)].`$ (25)
The squared fluctuations for $`\psi _{cg}`$ $`=`$ $`e^{3/2{\scriptscriptstyle (H_c+\tau _c/3+{\scriptscriptstyle \frac{2\dot{\alpha }}{3\alpha }})𝑑t}}\chi _{cg}`$ are given by
$$\psi _{cg}^2(\stackrel{}{x},t)_{IR}=e^{3{\scriptscriptstyle (H_c+\tau _c/3+{\scriptscriptstyle \frac{2\dot{\alpha }}{3\alpha }})𝑑t}}_0^{k_o}\frac{dkk^2}{(2\pi ^2)}\xi _k^2(t).$$
(26)
As the coupling parameter $`\alpha (t)`$ decreases during the rapid expansion of the universe \[$`\dot{\alpha }(t)<0`$\], the sign of the effective curvature $`K`$ will be dependent on the term $`\alpha \dot{\alpha }<\varphi \dot{\varphi }>`$ in eq. (13). However, during inflation the increasing rate of the scale factor is bigger than the $`K`$ one, so that at the end of inflation the ratio $`\frac{K}{a^2}`$ becomes nearly zero.
From the condition $`n1=2(1\nu )`$ for the spectral index $`n`$, one obtains the following condition for the constraint $`|n1|<0.3`$ obtained with the COBE data
$$\frac{\left[0.702(p1)\right]^21}{4}<L^2<\frac{\left[1.3(p1)\right]^21}{4}.$$
(27)
A spectral index $`n1`$, also has been found in a model with cosmic strings plus cold or hot dark matter. The standard choice of $`n=1`$ was first advocated by Harrison and Zel’dovich on the ground that it is scale invariant at the epoch of horizon entry. The condition (27) can be written as
$$\gamma _1(p)<\gamma <\gamma _2(p),$$
(28)
where $`\gamma _1(p)`$ and $`\gamma _2(p)`$ are given by the expressions
$`\gamma _1(p)={\displaystyle \frac{1}{f_1(p)}}[10^4(75pM_p^2225p^2M_p^2+2400m^2\pi )`$ (29)
$`+500M_p(23654592m^2\pi 106891191p^2M_p^2`$ (30)
$`2217618p^3M_p^2+1108809p^4M_p^211827296p^2m^2\pi `$ (31)
$`+1164172704m^2\pi )^{1/2}]`$ (32)
$`\gamma _2(p)={\displaystyle \frac{1}{f_2(p)}}[300([M_p^23M_p^2+32m^2\pi )`$ (33)
$`+10M_p(32448pm^2\pi 41679p^2M_p^23042p^3M_p^2`$ (34)
$`+1521p^4M_p^216224p^2m^2\pi +454176m^2\pi )^{1/2}],`$ (35)
with
$`f_1(p)`$ $`=`$ $`750000p^2M_p^28000000m^2\pi ,`$ (36)
$`f_2(p)`$ $`=`$ $`300p^2M_p^23200m^2\pi .`$ (37)
Finally, we can calculate the dependence in the amplitude of the fluctuations due to the semiclassical expansion here introduced. If $`\psi _{cg}^2_{IR}^{(\beta 0)}`$ is the amplitude for the fluctuations of the matter field within thermal coupling (i.e., with $`\beta 0`$) and $`\psi _{cg}^2_{IR}^{(\beta =0)}`$ is the amplitude for the matter field fluctuations without thermal coupling (i.e., for $`\beta =0`$), the only contribution in the amplitude for the fluctuations that arises from the thermal coupling will be
$$R(t)=\frac{\psi _{cg}^2_{IR}^{(\beta =0)}}{\psi _{cg}^2_{IR}^{(\beta 0)}}=\left(\frac{T_r(t)}{M}\right)^{2\beta },$$
(38)
where
$`\left({\displaystyle \frac{T_r(t)}{M}}\right)^2`$ $`=`$ $`{\displaystyle \frac{M_p^2}{8\pi M^2}}\sqrt{{\displaystyle \frac{15\gamma }{N}}}\left(1+{\displaystyle \frac{\gamma }{3}}\right)^1`$ (39)
$`=`$ $`{\displaystyle \frac{M_p^2}{8\pi M^2}}\sqrt{{\displaystyle \frac{15\gamma }{N}}}{\displaystyle \frac{p}{m}}\left(1+{\displaystyle \frac{\gamma }{3}}\right)^1t^1,`$ (40)
due to $`T_r(t)\rho _r^{1/4}(t)`$. The eq. (38) defines the ratio between the thermal and matter field fluctuations. For $`\beta >0`$, one obtains that this ratio decreases with time as $`t^\beta `$, and thus the thermal fluctuations should be more small than the matter fluctuations at the end of inflation. Furthermore, as $`\left(\frac{T_r(t)}{M}\right)^21`$, one obtains that $`\tau _c/H_c<10^{15}`$. Here, $`\beta `$ should be calculared from quantum field theory. (see for example ref. )).
To summarize, in this work I considered a new semiclassical expansion for the inflaton field in warm inflation. The matter field fluctuations are considered as thermally coupled with the particles of the thermal bath with temperature $`T_r<T_{GUT}`$. The coupling parameter $`\alpha =(T_r/M)^\beta `$, depends on the temperature of the bath. In this framework, I find that the equation of motion for the redefined matter field fluctuations $`\chi `$, remains invariant under the transformation $`\chi =e^{3/2{\scriptscriptstyle (H_c+\tau _c/3+{\scriptscriptstyle \frac{2\dot{\alpha }}{3\alpha }})𝑑t}}\psi `$. This implies that the chosen semiclassical expansion $`\phi =\varphi _c+\alpha (t)\psi `$ does not modifies the power spectrum — obtained with the standard semiclassical expansion $`\phi =\varphi _c+\varphi `$ — for the fluctuations of energy density. The thermal contribution to the fluctuations are given by the ratio $`R=\alpha ^2`$. I find — studying a power - law expansion for the universe — that the thermal component of the amplitude for these fluctuations in the infrared sector, should be very small at the end of inflation.
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# Localization Regions of Local Observables
## 1 Introduction
The algebraic approach to relativistic quantum physics aims at joining the structures familiar from nonrelativistic quantum mechanics to those of special relativity. The object of investigation is a net $`𝒜`$ of local observables that associates with every bounded open region $`𝒪`$ in the Minkowski spacetime $`𝐑^{1+s}`$ a unital C-algebra $`𝒜(𝒪)`$ of bounded operators in a Hilbert space $``$ in such a way that $`𝒪P𝐑^{1+s}`$ implies $`𝒜(𝒪)𝒜(P)`$ (isotony), and such that the elements of algebras associated with spacelike separated regions commute (locality). For every region $`𝒪`$, the elements of $`𝒜(𝒪)`$ are interpreted as the observables measurable in a lab located in $`𝒪`$.
This paper deals with the question whether, given a single element $`A𝒜(𝒪)`$ for some bounded $`𝒪𝐑^{1+s}`$, one can find a smallest nonempty region $`L(A)`$ within which $`A`$ can be measured. From a theorem of Landau it is well known that in at least 1+2 spacetime dimensions, one observable cannot be measured in two labs located in disjoint spacetime regions. A generalization of this theorem will be proved below, and for at least 1+2 spacetime dimensions, it will lead to several meaningful definitions of convex localization regions. Some additional technical assumptions then provide a strongest localization prescription with the property that observables with spacelike separated localization regions commute. But the fact that locality of the net alone (not even together with those assumptions that lead to the definition of a nonempty localization region) does not imply this version of locality, is a little surprising.
This article is structured as follows: Section 2.1 discusses the notation, concepts, and basic assumptions that play a role in this paper. Section 2.2 collects and completes the tools that will be used in what follows. Section 3 contains the main results of the paper. In Section 3.1 the theorem due to Landau mentioned above is discussed and generalized: it states that in 1+2 or more spacetime dimensions, the algebras of any two double cones with disjoint closures only have in common the complex multiples of the identity operator. Using the techniques introduced in the first sections, it is shown that this result can be generalized to the empty-intersection theorem: instead of considering two double cones, one can consider one double cone and any finite number of wedge regions; if the closures of the regions under consideration have an empty common intersection, then the corresponding algebras have in common exactly the c-numbers.
In Section 3.2, finally, it is discussed how one can, in 1+2 or more spacetime dimensions, use the empty-intersection theorem in order to associate a nonempty localization region with any given single local observable that is not a multiple of the identity. But even though locality of the net is assumed from the outset, the question whether local observables with spacelike separated localization regions commute turns out to be nontrivial. A necessary and sufficient criterion for the locality of such a localization prescription is provided by the nonempty-intersection theorem: the criterion requires that, given any finite family of wedges, all local observables contained in the algebras associated with all these wedges are contained in the algebra associated with any neighbourhood of the intersection of the wedges as well. It is shown that an additional additivity assumption, which is typically fulfilled by nets arising from Wightman fields, implies this criterion, so it should be quite difficult to find nonpathological counterexamples. On the other hand it is illustrated why the nonempty-intersection criterion, though looking quite natural, is far from self-evident. It does not follow from the locality property of the net.
In the Conclusion, some related results are discussed briefly.
## 2 Preliminaries
### 2.1 Notation and assumptions
In what follows, $``$ will be an infinite-dimensional (not necessarily separable) Hilbert space, and $`𝒜`$ will be a net of observables as defined above, i.e., a net which satisfies locality and isotony. The union of all algebras $`𝒜(𝒪)`$ associated with bounded open regions<sup>1</sup><sup>1</sup>1In this article we refer to arbitrary subsets of $`𝐑^{1+s}`$ as ‘regions’. $`𝒪𝐑^{1+s}`$, $`s1`$, is an involutive algebra $`𝒜_{\mathrm{loc}}`$ called the algebra of local observables. Throughout Section 3 it will be assumed that $`s2`$, but for the introductory part this assumption is not yet necessary.
The following (standard) assumptions on $`𝒜`$ will be made throughout this article:
> (A) Translation covariance. $`𝒜`$ is covariant under a strongly continuous unitary representation $`U`$ of the group $`(𝐑^{1+s},+)`$ of spacetime translations, i.e.,
>
> $$U(a)𝒜(𝒪)U(a)=𝒜(𝒪+a)$$
>
> for every bounded region $`𝒪`$ and every $`a𝐑^{1+s}`$.
>
> (B) Spectrum condition. The spectrum of the four-momentum operator generating $`U`$ is contained in the closure of the forward light cone.
>
> (C) Existence and uniqueness of the vacuum. The space of $`U`$-invariant vectors in $``$ is one-dimensional. $`\mathrm{\Omega }`$ will denote an arbitrary, but fixed unit vector $`\mathrm{\Omega }`$ in this space, the vacuum vector. $`\mathrm{\Omega }`$ is cyclic with respect to $`𝒜_{\mathrm{loc}}`$, i.e., $`\overline{𝒜_{\mathrm{loc}}\mathrm{\Omega }}=`$.
Throughout Section 3 the following additional assumption will be made:
> (D) Reeh-Schlieder property. For every nonempty bounded open region $`𝒪`$, the space $`𝒜(𝒪)\mathrm{\Omega }`$ is dense in $``$.
Assumptions (A) and (B) make sure that the system described by the net has a well-defined four-momentum whose spectrum ensures energetic stability of the system.
Given Assumptions (A) and (B), Assumption (C) holds if and only if $`\mathrm{\Omega }`$ induces a unique and pure vacuum state, which, in turn, holds if and only if the algebra $`𝒜_{\mathrm{loc}}`$ is irreducible. A sufficient condition for this uniqueness is that the bicommutant $`𝒜_{\mathrm{loc}}^{\prime \prime }`$ of $`𝒜_{\mathrm{loc}}`$ is a factor (Thm. III.3.2.6 in ). But as soon as $``$ is separable, this implies that the uniqueness part of Assumption (C) does not mean any loss of generality: every von Neumann algebra in a separable Hilbert space admits a direct-integral decomposition into factors, and since the unitaries representing the translations commute with the elements of the center of $`𝒜_{\mathrm{loc}}^{\prime \prime }`$, one can conclude that almost all factors of the central decomposition inherit Properties (A) and (B) (cf. also the remarks in , Sect. III.3.2, and references therein).
The Reeh-Schlieder property, Assumption (D), holds for all Wightman fields . If the region $`𝒪`$ contains an open cone, it is well known to follow from conditions (A) through (C) (cf., e.g., the Appendix in ). The Reeh-Schlieder property holds as soon as one has weak additivity (, cf. also Thm. 7.3.37 in ): if $`𝒪`$ is any bounded open region, then
$$\left(\underset{a𝐑^{1+s}}{}𝒜(𝒪+a)\right)^{\prime \prime }=𝒜_{\mathrm{loc}}^{\prime \prime }.$$
Conversely, it is well known that weak additivity can be derived from the Reeh-Schlieder property as well, provided Assumptions (A) through (C) hold (cf., e.g., Lemma 2.6 in ). For the reader’s convenience we include a proof of this fact in the appendix.
Some special classes of spacetime regions will be used below. The first one is the class $`𝒦`$ of double cones, i.e., all regions of the form $`(a+V_+)(bV_+)`$, $`a,b𝐑^{1+s}`$. The class $`𝒲`$ of wedges consists of the region $`W_1:=\{x𝐑^{1+s}:x_1>|x_0|\}`$ and its images under Poincaré transforms.
If $`M`$ is a region in $`𝐑^{1+s}`$, one denotes by $`M^c`$ the causal complement or spacelike complement, which is the region consisting of all points that are spacelike with respect to all points of $`M`$. The spacelike complement of the spacelike complement $`(M^c)^c=:M^{cc}M`$ is called the causal completion of $`M`$, and $`M`$ is called causally complete if $`M=M^{cc}`$. It is convenient to denote the interior of $`M^c`$ by $`M^{}`$.
$`𝒦`$ and $`𝒲`$ are subclasses of the class $`𝒞`$ of convex, causally complete and open proper subsets of $`𝐑^{1+s}`$. The wedges in $`𝒲`$ are maximal elements of $`𝒞`$ in the sense that for every wedge $`W𝒲`$, every element $`R𝒞`$ with $`RW`$ is a wedge. Every element $`R`$ of $`𝒞`$ is an intersection of wedges (cf. , Thm. 3.2)<sup>2</sup><sup>2</sup>2For the proof of this statement it is essential that $`𝒞`$ consists of open regions (the statement also implies to regions with a nonempty interior). As a counterexample, consider the lightlike half plane $`R:=\{x𝐑^{1+s}:x_1=x_0,x_2>0,\mathrm{}x_s>0\}`$. One checks that this region is causally complete and convex, while it is not an intersection of wedges.. The class of all wedges that contain a region $`R`$ will be denoted by $`𝒲_R`$.
In general, the causal complement of a region in $`𝒞`$ is not convex. If $`R𝒞`$, then $`R^{}`$ is a union of wedges (, Thm. 3.2), and $`𝒲^R^{}`$ will denote the class of all wedges that are subsets of $`R^{}`$. If $`𝒪`$ is an open convex region and if $`P`$ is a convex region that is spacelike separated from $`𝒪`$, there is a wedge $`X𝒲`$ such that $`𝒪X`$ and $`PW^c`$ (cf. , Prop. 3.1).
$``$ will denote the bounded elements of the class $`𝒞`$. Clearly, the double cones are in $``$. Every element of $``$ is contained in some double cone, and it is precisely the intersection of all such double cones (, Prop. 3.8). The class of all double cones which contain a region $`𝒪`$ will be denoted by $`𝒦_𝒪`$, and the class of all double cones contained in an arbitrary region $`R`$ will be called $`𝒦^R`$.
In Section 3.2, two more technical assumptions will occur:
> (E) Wedge duality. For all $`W𝒲`$, one has $`𝒜(W^{})^{}=𝒜(W)^{\prime \prime }`$.
>
> (F) Wedge additivity. For each wedge $`W𝒲`$ and each double cone $`𝒪𝒦`$ with $`WW+𝒪`$ one has
>
> $$𝒜(W)^{\prime \prime }\left(\underset{aW}{}𝒜(a+𝒪^{})^{}\right)^{\prime \prime }.$$
All nets arising from finite-component Wightman fields satisfy wedge duality . One checks that wedge duality implies the condition of essential duality known from the analysis of superselection sectors, since for any two spacelike separated regions in $``$, one can find a wedge which contains one of the two, whereas its spacelike complement contains the other one (see above, and cf. Lemma 3.2.2 below). The algebras $`𝒜(𝒪^{})^{}`$ will also occur below for other regions $`𝒪`$, e.g., for double cones. By locality, $`𝒜(𝒪^{})^{}𝒜(𝒪)`$, and with the above remarks, it is easy to show that as soon as wedge duality holds, one obtains that for spacelike separated double cones $`𝒪`$ and $`P`$, the elements of the algebras $`𝒜(𝒪^{})^{}`$ and $`𝒜(P^{})^{}`$ commute, a property which is called essential duality and which is used in the theory of localized superselection sectors (cf. and references given there).
Condition (F) strengthens Condition (D) slightly, but it is a standard property of all Wightman fields as well. It has been used extensively by Thomas and Wichmann in . The authors have obtained results in the spirit of Theorem 3.1.5 and Proposition 3.2.5 below, but their results do not imply ours.
Occasionally, terminology borrowed from PDEs and General Relativity will be used (timelike curves, Cauchy surfaces, etc.). These notions will not be defined in detail, but will be used as in .
### 2.2 Commutator functions and wave equation techniques
It is a classical result of the Wightman approach to quantum field theory that one can reconstruct a Wightman field from its vacuum expectation values . The following lemma shows how one can reconstruct commutation relations of a net of observables from the behaviour of its vacuum expectation values. Since these have some convenient properties, this will facilitate the subsequent investigations. $`𝒜`$ will be a local net of local observables satisfying the above Assumptions (A) through (C).
#### 2.2.1 Lemma
> For an arbitrary double cone $`𝒪𝒦`$, let $`A`$ be an element of $`𝒜(𝒪^{})^{}`$.
>
> > (i) If a region $`R𝐑^{1+s}`$ contains some open cone and has the property that $`\mathrm{\Omega },AB\mathrm{\Omega }=\mathrm{\Omega },BA\mathrm{\Omega }`$ for all $`B𝒜(R)`$, then $`A𝒜(R)^{}`$.
> >
> > (ii) Assume that $`𝒜`$ has the Reeh-Schlieder property, and suppose there is a double cone $`P𝒦`$ with the property that $`\mathrm{\Omega },AB\mathrm{\Omega }=\mathrm{\Omega },BA\mathrm{\Omega }`$ for all $`B𝒜(P)`$. If there is a double cone $`QP`$ with the property that $`A𝒜(Q)^{}`$, then $`A𝒜(P)^{}`$.
> >
> > (iii) Assume that $`𝒜`$ exhibits the Reeh-Schlieder property, and suppose there is a double cone $`P𝒦`$ with the property that $`\mathrm{\Omega },AB\mathrm{\Omega }=\mathrm{\Omega },BA\mathrm{\Omega }`$ for all $`B_{a𝐑^{1+s}}𝒜(P+a)`$. Then $`A`$ is a multiple of the identity.
Proof. (i): If $`S`$ is an open cone contained in $`R`$, there is a translation $`a𝐑^{1+s}`$ such that $`S+aR𝒪^{}`$. Choose $`C`$ and $`D`$ in $`𝒜(S+a)`$ and $`B`$ in $`𝒜(R)`$. Since $`A𝒜(𝒪^{})^{}`$, the operators $`A`$ and $`C^{}`$ commute:
$$C\mathrm{\Omega },ABD\mathrm{\Omega }=\mathrm{\Omega },C^{}ABD\mathrm{\Omega }=\mathrm{\Omega },AC^{}BD\mathrm{\Omega }.$$
Since $`C^{}BD`$ is in $`𝒜(R)`$, the assumption implies
$$\mathrm{\Omega },AC^{}BD\mathrm{\Omega }=\mathrm{\Omega },C^{}BDA\mathrm{\Omega },$$
and since $`D`$ and $`A`$, in turn, commute because of $`A𝒜(𝒪^{})^{}`$, one concludes
$$C\mathrm{\Omega },ABD\mathrm{\Omega }=\mathrm{\Omega },C^{}BDA\mathrm{\Omega }=C\mathrm{\Omega },BAD\mathrm{\Omega }.$$
But since $`C`$ and $`D`$ are arbitrary elements of $`𝒜(S+a)`$, and since $`\mathrm{\Omega }`$ is cyclic with respect to this algebra, it follows that $`AB=BA`$; since $`B𝒜(R)`$ was arbitrary, one obtains $`A𝒜(R)^{}`$, which is (i).
(ii) Choose $`C`$ and $`D`$ in $`𝒜(Q)`$ and $`B`$ in $`𝒜(P)`$. Since $`A`$ has been assumed to be in $`𝒜(Q)^{}`$, it commutes with $`C^{}`$, so
$$C\mathrm{\Omega },ABD\mathrm{\Omega }=\mathrm{\Omega },C^{}ABD\mathrm{\Omega }=\mathrm{\Omega },AC^{}BD\mathrm{\Omega }.$$
Since $`C^{}BD`$ is in $`𝒜(P)`$, the assumption implies
$$\mathrm{\Omega },AC^{}BD\mathrm{\Omega }=\mathrm{\Omega },C^{}BDA\mathrm{\Omega },$$
and since $`D`$ and $`A`$ commute by the assumption that $`A𝒜(Q)^{}`$, one concludes
$$C\mathrm{\Omega },ABD\mathrm{\Omega }=\mathrm{\Omega },C^{}BDA\mathrm{\Omega }=C\mathrm{\Omega },BAD\mathrm{\Omega }.$$
But since $`C`$ and $`D`$ are arbitrary elements of $`𝒜(Q)`$, and since by the Reeh-Schlieder property, $`\mathrm{\Omega }`$ is cyclic with respect to this algebra, it follows that $`AB=BA`$; since $`B𝒜(P)`$ was arbitrary, one obtains $`A𝒜(P)^{}`$, which is (ii).
(iii) There is a translation $`a𝐑^{1+s}`$ such that $`P+a𝒪^{}`$, so that $`A𝒜(𝒪^{})^{}𝒜(P+a)^{}`$. Now choose a $`b𝐑^{1+s}`$ such that $`P+b`$ intersects $`P+a`$, and let $`Q`$ be a double cone contained in $`(P+b)(P+a)`$. Isotony implies that $`A𝒜(Q)^{}`$. Since by assumption, $`\mathrm{\Omega },AB\mathrm{\Omega }=\mathrm{\Omega },BA\mathrm{\Omega }`$ for all $`B𝒜(P+b)`$, (ii) implies that $`A𝒜(P+b)^{}`$. Now one can iterate this procedure: choose an arbitrary $`c𝐑^{1+s}`$ such that $`(P+c)(P+b)`$ is nonempty, choose a new double cone $`Q`$ in this intersection, and conclude from (ii) that $`A𝒜(P+c)^{}`$. Note that only the double cone $`P+a`$ chosen in the first step needs to be spacelike separated from $`𝒪`$, since each step uses the result of the preceding one, so one finds that for every $`a𝐑^{1+s}`$, one proves that $`A𝒜(P+a)^{}`$ with a finite number of steps. The statement now follows from weak additivity, which follows from the Reeh-Schlieder property (see above), and from irreducibility.
$`\mathrm{}`$
Given any two local observables $`A,B𝒜_{\mathrm{loc}}`$, the commutator function $`f_{A,B}`$ will henceforth be defined by
$$𝐑^{1+s}x\mathrm{\Omega },[A,U(x)BU(x)]\mathrm{\Omega }=:f_{A,B}(x).$$
Due to Lemma 2.2.1, the analysis of the support of this function yields information on the structure of the net. Crucial for this analysis is the fact that $`f_{A,B}`$ is a boundary value of a solution of the wave equation, and a well-known lemma due to Asgeirsson concerning such solutions (cf., , Sect. 4.4.D in , or ) immediately implies the following lemma, which, for this reason, will be referred to as Asgeirsson’s Lemma. Another important consequence of the ‘wave nature’ of the function $`f_{A,B}`$ is a theorem due to Jost, Lehmann and Dyson which will also be recalled for the reader’s convenience.
#### 2.2.2 Lemma (Asgeirsson)
> If the commutator function $`f_{A,B}`$ and all its partial derivatives are zero along a timelike curve segment $`\gamma `$, $`f_{A,B}`$ vanishes in the entire double cone $`\gamma ^{cc}`$.
Proof. The Fourier transform of the operator valued function $`𝐑^{1+s}xU(x)`$ is the spectral measure of the four-momentum operator. It follows that the Fourier transform $`\widehat{f}_{A,B}`$ of the function $`f_{A,B}`$ is a finite (not necessarily positive) measure, and by the spectrum condition, one has $`\mathrm{supp}\widehat{\mathrm{f}}_{\mathrm{A},\mathrm{B}}\overline{\mathrm{V}}`$. It follows that the function
$$F(x,\sigma ):=(2\pi )^{\frac{1+s}{2}}\mathrm{cos}(\sigma \sqrt{k^2})e^{ikx}𝑑\widehat{f}_{A,B}(k)$$
is a continuous function with $`F(x,0)=f_{A,B}(x)`$ for all $`x𝐑^{1+s}`$. This $`F`$ is a solution of the 1+(s+1)-dimensional wave equation. This implies the statement by Asgeirsson’s result for solutions of the wave equation, see the references quoted above.
$`\mathrm{}`$
Evidently, the assumption of the lemma is satisfied as soon as $`f_{A,B}`$ vanishes in some open neighbourhood of $`\gamma `$. In the proof of Theorem 2.2.5 below, however, the function $`F`$ defined in the proof is analysed, and the information one has about $`f_{A,B}`$ from locality merely implies that $`F`$ vanishes in a null set of $`𝐑^{1+(s+1)}`$. In this case one makes use of the fact that $`F`$ has been constructed in such a way that all its partial derivatives, including the one in the $`\sigma `$-direction, are zero at all points of this null set; one may then use the above lemma to show that the region where $`F`$ vanishes also extends into the $`\sigma `$-direction.
#### 2.2.3 Definition
> Let $`R`$ be a region in Minkowski space.
>
> (i) $`R`$ will be called Asgeirsson complete if for every timelike curve segment $`\gamma R`$, the double cone $`\gamma ^{cc}`$ is a subset of $`R`$ as well. The smallest Asgeirsson complete extension of $`R`$ will be called the Asgeirsson hull of $`R`$.
>
> (ii) $`R`$ will be called timelike convex if it contains as subsets all double cones with tips in $`R`$, i.e., if $`(R+V_+)(RV_+)R`$.
>
> (iii) $`R`$ will be called a Jost-Lehmann-Dyson region if it is timelike convex and if every inexdentible timelike curve in $`𝐑^{1+s}`$ intersects $`RR^c`$.
Timelike convex regions contain all timelike path segments connecting two points in the region, so the terminology is in harmony with other notions of convexity. In the term ‘double cone complete’ was used instead of ‘timelike convex’, but the latter term was also used in (Par. IV) and will be used in what follows to facilitate reading. The following lemma collects some relations between these notions most of which will be used below.
#### 2.2.4 Lemma
> (i) Every causally complete region is timelike convex.
>
> (ii) Every timelike convex region is Asgeirsson complete.
>
> (iii) Every timelike convex and bounded open region is a Jost-Lehmann-Dyson region.
>
> (iv) The causal complement of a Jost-Lehmann-Dyson region is a Jost-Lehmann-Dyson region.
>
> (v) Let $`R`$ and $`S`$ be timelike convex regions, and assume that there exists a Cauchy surface $`T`$ which is a subset of both $`R`$ and $`S`$. Then the region $`RS`$ is timelike convex (and, like $`R`$ and $`S`$, trivially, a Jost-Lehmann-Dyson region).
>
> (vi) Let $`(R_\rho )_{\rho >0}`$ be an increasing family of Jost-Lehmann-Dyson regions. Then $`R:=_\rho R_\rho `$ is a Jost-Lehmann-Dyson region.
Before proving the lemma, we give some counterexamples to strengthened statements or converse implications. An example of a timelike convex region (and Jost-Lehmann-Dyson region) that is not causally complete (cf. (i)) is the time slice region $`\{x𝐑^{1+s}:\mathrm{\hspace{0.17em}0}x_01\}`$. An example of an Asgeirsson complete region that is not timelike convex (cf. (ii)) is the union of two disjoint double cones at a timelike distance; this shows that the classes of timelike convex regions and of Jost-Lehmann-Dyson regions, respectively, are not stable under taking unions, so Statement (v) is far from tautological. The same holds for the class of Asgeirsson complete regions: consider the regions $`R_+:=\{x𝐑^{1+s}:\rho x_1<x_0<\rho x_1+1\}`$ and $`R_{}:=\{x𝐑^{1+s}\rho x_1<x_0<1\rho x_1\}`$ for some $`\rho `$ with $`0<\rho 1`$. One easily checks that both regions are Asgeirsson complete, while their union is not: its Asgeirsson hull is $`𝐑^{1+s}`$. If $`\rho <1`$, the two regions are even Jost-Lehmann-Dyson regions, while their union evidently is not (cf. (v) and (vi)).
An example of a timelike convex region which is neither causally complete nor a Jost-Lehmann-Dyson region (cf. (iii))is the region
$$R:=\{x𝐑^{1+s}:\mathrm{\hspace{0.17em}1}<x^2<2,x_0>0\},$$
since there are timelike curves which do not intersect $`R`$, e.g., the curve $`𝐑t(\mathrm{sinh}t,\mathrm{cosh}t,0,\mathrm{},0)`$.
Proof of Lemma 2.2.4. (i) Let $`R`$ be a causally complete region, and pick two points $`xR`$ and $`yR(x+V_+)`$. The causal completion $`\{x,y\}^{cc}`$ of the set $`\{x,y\}`$ is the closure of the double cone $`(x+V_+)(yV_+)`$, and since $`\{x,y\}R`$ implies $`\{x,y\}^{cc}R^{cc}=R`$, this immediately implies (i).
Statement (ii) immediately follows from the definition.
(iii) Let $`R`$ be timelike convex, bounded and open. Since $`R`$ is open, a point $`x𝐑^{1+s}`$ is not contained in the spacelike complement $`R^c`$ if and only if it is timelike with respect to some point of $`R`$, i.e., $`𝐑^{1+s}\backslash R^c=R+V`$, where $`V`$ is the open light cone. Now let $`\gamma `$ be an inextendible timelike curve that does not intersect $`RR^c`$. Since $`\gamma `$ does not intersect $`R^c`$, it has to stay within the region $`R+V`$. But since $`\gamma `$ is an inextendible timelike curve, while $`R`$ is bounded, $`\gamma `$ cannot stay in the future $`R+V_+`$ or the past $`RV_+`$ of $`R`$, i.e., it has to pass from $`RV_+`$ to $`R+V_+`$. Since both these regions are open, while $`\gamma `$ is continuous, it follows that it has to hit the region $`(R+V_+)(RV_+)`$. But this region coincides with $`R`$ since $`R`$ is timelike convex and open, so $`\gamma `$ hits $`R`$, which is a contradiction and proves (iii).
(iv) The causal complement of any region is causally complete, by (i), this enhances timelike convexity of $`R^c`$. The condition that $`RR^c`$ is intersected by every inextendible timelike curve implies that $`R^cR^{cc}`$ ($`=R^cR`$) is intersected by every such curve. This proves (iv).
(v) Let $`x`$ and $`y`$ be points in $`RS`$ that are timelike with respect to each other. Since $`𝐑^{1+s}`$ is timelike convex, one finds an inextendible timelike curve $`\gamma `$ hitting both $`x`$ and $`y`$. Let $`z`$ be the unique point where $`\gamma `$ hits $`T`$. Since $`R`$ and $`S`$ are timelike convex, and since $`zTRS`$, the closed double cones with the tips $`z`$ and $`x`$ and the tips $`z`$ and $`y`$, respectively, are subsets in $`RS`$. If with respect to the time ordering along $`\gamma `$, $`z`$ is earlier or later than both $`x`$ and $`y`$, it follows that the double cone with tips $`x`$ and $`y`$ is contained in $`RS`$ as well. If $`z`$ is between $`x`$ and $`y`$, then, as before, we can conclude that the segments of $`\gamma `$ between $`z`$ and $`x`$ and between $`z`$ and $`y`$ is a subset of $`RS`$, and since $`zTRS`$, it follows that all of the segment of $`\gamma `$ joining $`x`$ to $`y`$ is a subset of $`RS`$. Since $`\gamma `$ can be any inextendible timelike curve hitting $`x`$ and $`y`$, one obtains that all timelike curve segments joining $`x`$ and $`y`$ are contained in $`RS`$, so the double cone with tips $`x`$ and $`y`$ is contained in $`RS`$, which completes the proof of (v).
(vi) Let $`x`$ and $`y`$ be two points in $`R`$ at a timelike distance. There are a $`\rho _x>0`$ and a $`\rho _y>0`$ such that $`xR_{\rho _x}`$ and $`yR_{\rho _y}`$, so it follows that both $`x`$ and $`y`$ are elements of $`R_{max\{\rho _x,\rho _y\}}`$. Since this region is timelike convex, it follows that the double cone with tips $`x`$ and $`y`$ is in $`R`$, proving that $`R`$ is timelike convex.
It remains to be shown that every inextendible timelike curve intersects $`RR^{}`$. Let $`\gamma `$ be an inextendible timelike curve that does not intersect $`R`$. Since all $`R_\rho `$ are Jost-Lehmann-Dyson regions, it follows that $`\gamma `$ has to intersect every $`R_\rho ^{}`$, so it has to intersect the region $`_{\rho >0}R_\rho ^{}=R^{}`$. This completes the proof.
$`\mathrm{}`$
The above statements and proofs can be extended in a straightforward manner to the spacetime one obtains by endowing the cyclinder
$$Z_\rho :=\{x=(x_0,\stackrel{}{x})𝐑^{1+s}:\stackrel{}{x}=\rho \}$$
with the spacetime structure it inherits from $`𝐑^{1+s}`$, provided $`s2`$. For $`s=1`$ this spacetime fails to be timelike convex, and the proof of part (v) does no longer work. For further results of the above kind, see . The useful property of Jost-Lehmann-Dyson regions (which is the reason to call them so) is established by the following theorem.
#### 2.2.5 Theorem (Jost, Lehmann, Dyson)
> Let $`A`$ and $`B`$ be local observables, and assume that the commutator function $`f_{A,B}`$ vanishes in a Jost-Lehmann-Dyson region $`R`$. Then the support of $`f_{A,B}`$ is contained not only in the complement of $`R`$, but even in the (in general, smaller) union of all admissible mass hyperboloids of $`R`$, i.e., the mass hyperboloids
>
> $$H_{a,\sigma }:=\{x𝐑^{1+s}:(xa)^2=\sigma ^2\},a𝐑^{1+s},\sigma 𝐑,$$
>
> which do not intersect $`R`$.
Sketch of proof. Define $`F`$ as in the proof of Lemma 2.2.2. Since $`F`$ is a solution of the wave equation, it is well-known that for every Cauchy surface $`\zeta `$ in $`𝐑^{1+(s+1)}`$, there exists a distribution $`F_\zeta `$ with support in $`\zeta `$ such that $`F=F_\zeta D_{1+(s+1)}`$, where $`D_{1+(s+1)}`$ denotes a fundamental solution of the 1+(s+1)-dimensional wave equation (see, e.g., , pp. 175-184). The support of $`D_{1+(s+1)}`$ is contained in the closed light cone $`\overline{\widehat{V}}`$ of $`𝐑^{1+(s+1)}`$ <sup>3</sup><sup>3</sup>3This notation is consistent since $`\widehat{V}`$ is, indeed, the 1+(s+1)-dimensional Asgeirsson hull of $`V`$. Note that $`\widehat{\overline{V}}\overline{\widehat{V}}`$.. Since $`R`$ is a Jost-Lehmann-Dyson region in $`𝐑^{1+s}`$, its 1+(s+1)-dimensional Asgeirsson hull $`\widehat{R}`$ is easily seen to be a Jost-Lehmann-Dyson region in $`𝐑^{1+(s+1)}`$. Provided this region is ‘well-behaved’, there is a Cauchy surface $`\zeta `$ in $`\widehat{R}\widehat{R}^c`$. This Cauchy surface has the property that for every point $`z\zeta `$, either both the forward and the backward part of $`\overline{\widehat{V}}+z`$ or neither of them intersects $`R`$. The former case occurs if and only if $`z\zeta \widehat{R}`$. The latter case occurs if and only if $`z\zeta \widehat{R}^c`$, the Asgeirsson hull $`\widehat{R}`$ of $`R`$ and the spacelike complement being taken in the spacetime $`𝐑^{1+(s+1)}`$. But since all partial derivatives of $`F`$ can be checked to vanish in all points in $`R`$, one obtains from Lemma 2.2.2 that $`F`$ vanishes in $`\widehat{R}`$, the support of $`F_\zeta `$ contains only points of the second kind, i.e., it is contained in $`\widehat{R}^c\zeta `$. This implies that the support of $`F`$ is contained in $`(\widehat{R}\zeta )+\overline{\widehat{V}}`$.
Since $`f_{A,B}`$ is a boundary value of $`F`$ and since the intersection of $`\overline{\widehat{V}}+c`$ with $`𝐑^{1+s}`$ is the convex hull of a shifted mass hyperboloid, the support of the boundary value $`f_{A,B}`$ of the function $`F`$ is contained in the union of admissible mass hyperboloids, as stated.
$`\mathrm{}`$
We conclude this section with another lemma to be used below that concerns the geometry of Minkowski space.
#### 2.2.6 Lemma
> Let $`P𝒦`$ be a double cone.
>
> (i) If $`𝒪`$ is a double cone, so is $`(𝒪+P)^{cc}`$.
>
> (ii) If $`W`$ is a wedge, so is $`(W+P)^{cc}`$.
Proof. Denote by $`a_𝒪`$ and $`a_P`$ the lower tips, and by $`b_𝒪`$ and $`b_P`$ the upper tips of $`𝒪`$ and $`P`$, respectively. Let $`x=a_𝒪+\xi `$ and $`y=a_P+\eta `$ be points in $`𝒪`$ and $`P`$, respectively. Then $`x+y=a_𝒪+a_P+\xi +\eta `$, and since $`\xi `$ and $`\eta `$ are elements of $`V_+`$, so is $`\xi +\eta `$, so $`x+y`$ is contained in $`a_𝒪+a_P+V_+`$. In the same way one proves that $`x+yb_𝒪+b_PV_+`$, so one has
$$𝒪+P(a_𝒪+a_P+V_+)(b_𝒪+b_PV_+).$$
Since the right hand side is a double cone and, hence, causally complete, it follows that $`(𝒪+P)^{cc}`$ is a subset of this double cone as well. On the other hand it is straightforward to check that the tips $`a_𝒪+a_P`$ and $`b_𝒪+b_P`$ of this double cone and the straight line joining them are contained in $`𝒪+P`$, whence the converse inclusion follows as well, so the proof of (i) is complete.
The region $`W+P`$ is a union of wedges that are images of $`W`$ under translations. Consequently, $`(W+P)^c`$ is the intersection of the corresponding translates of $`W^c`$ under translations. But this intersection is the closure of a wedge, so it follows that the causal complement of this region, $`(W+P)^{cc}`$, is a wedge. This proves (ii).
$`\mathrm{}`$
## 3 Results
By definition, a local net associates algebras with regions. In the sequel it will be discussed how to associate a localization region with a given algebra and even with a single local observable. The analysis is based on a theorem due to Landau . In order to localize single observables, a new generalization of Landau’s theorem will be used. It will be stated and proved below.
This section is structured as follows: in Section 3.1, the theorem of Landau and its consequences for the localization of algebras will be discussed, and the mentioned generalization will be proved. This generalization will be the basis for the analysis of localization regions for single local observables, which is presented in Section 3.2.
In what follows, Assumptions (A) through (D) will be made without further mentioning, and it will assumed in addition that $`s2`$; Landau’s theorem and all generalizations discussed below heavily depend on this assumption, and so do the consequences to be discussed later on.
### 3.1 Landau’s theorem and the empty-intersection theorem
Using the wave equation techniques discussed in the preceding section, Landau proved the following:
#### 3.1.1 Theorem (Landau)
> If the closures of two double cones $`𝒪`$ and $`P`$ are disjoint, then
>
> $$𝒜(𝒪^{})^{}𝒜(P^{})^{}=𝐂\mathrm{id}_{}.$$
This already implies that for an $`𝒪`$ satisfying the assumptions of the corollary, the region
$$L(𝒜(𝒪^{})^{}):=\{P𝒦:𝒜(P)𝒜(𝒪^{})^{}\},$$
which will be called the localization region of the algebra $`𝒜(𝒪^{})^{}`$, coincides with $`𝒪`$ (cf. ):
#### 3.1.2 Corollary
> Let $`𝒪𝐑^{1+s}`$ be a bounded, causally complete and convex open region.
>
> (i) For every open region $`M𝐑^{1+s}`$, one has $`𝒜(M)𝒜(𝒪^{})^{}`$ if and only if $`M𝒪`$.
>
> (ii) $`L(𝒜(𝒪))=𝒪`$.
Proof. By isotony and locality, the condition in statement (i) is sufficient. To prove that it is necessary, assume $`M𝒪`$. Then, since $`𝒦`$ is a topological base and since the region $`M\backslash \overline{𝒪}`$ has a nonempty interior, $`M\backslash \overline{𝒪}`$ contains a double cone $`P𝒦`$ whose closure is disjoint from $`\overline{𝒪}`$. Since $`\overline{𝒪}`$ is an intersection of closures of wedges, it follows from this that a wedge $`W`$ can be found whose closure is disjoint from $`\overline{P}`$ and contains $`\overline{𝒪}`$. Since $`\overline{P}`$ is compact, the distance between $`\overline{P}`$ and $`\overline{W}`$ is $`>0`$, so eventually shifting it a little bit, one can choose $`W`$ in such a way that $`\overline{W}`$ is a subset not only of $`\overline{W}`$, but also of $`W`$ itself.
By Proposition 3.8 (b) in , on can now conclude that there is a double cone $`Q`$ with $`QW`$ and $`Q𝒪`$ (note that $`𝒪`$ itself does not need to be a double cone). Landau’s theorem now implies that $`𝒜(P)𝒜(Q^{})^{}=𝐂\mathrm{id}_{}`$. It follows from the Reeh-Schlieder property that $`𝒜(P)𝐂\mathrm{id}_{}`$, so $`𝒜(P)𝒜(Q^{})^{}`$. Since $`𝒜(P)𝒜(M)`$ follows from isotony, $`𝒜(M)`$ cannot be a subset of $`𝒜(Q^{})^{}`$, and since $`𝒪Q`$, it cannot be a subset of $`𝒜(𝒪^{})^{}`$. This proves (i) and, trivially, implies (ii).
$`\mathrm{}`$
The proof of Corollary 3.1.2 can be made shorter as soon as one knows that Landau’s theorem still works if one of the two double cones is replaced by a wedge. That this, indeed, is possible, has been shown in the context of the proof of the P<sub>1</sub>CT-part of the first uniqueness theorem for modular symmetries (Theorem 2.1 in ).
#### 3.1.3 Theorem
> If the closures of a double cone $`𝒪`$ and a wedge $`W`$ are disjoint, then
>
> $$𝒜(𝒪^{})^{}𝒜(W^{})^{}=𝐂\mathrm{id}_{}.$$
Using this generalized version of Landau’s theorem, one concludes that in Lemma 3.1.2, the assumption that $`𝒪`$ is bounded may be omitted:
#### 3.1.4 Corollary
> Let $`R𝐑^{1+s}`$ be a causally complete convex open region.
>
> (i) For every open region $`M𝐑^{1+s}`$, one has $`𝒜(M)𝒜(R^{})^{}`$ if and only if $`MR`$.
>
> (ii) $`L(𝒜(R))=R`$.
Proof. By isotony and locality, the condition is sufficient. To prove that it is necessary, assume $`MR`$. Then, since $`𝒦`$ is a topological base and since the region $`M\backslash \overline{R}`$ has a nonempty interior, $`M\backslash \overline{R}`$ contains a double cone $`𝒪𝒦`$ whose closure is disjoint from $`\overline{R}`$. As in the proof of Corollary 3.1.2, it follows that a wedge $`W`$ can be found whose closure is disjoint from $`\overline{𝒪}`$ and whose interior contains $`\overline{R}`$. Landau’s theorem now implies that $`𝒜(𝒪)𝒜(W^{})^{}=𝐂\mathrm{id}_{}`$. It follows from the Reeh-Schlieder property that $`𝒜(𝒪)𝐂\mathrm{id}_{}`$, so $`𝒜(𝒪)𝒜(W^{})^{}`$. Since $`𝒜(𝒪)𝒜(M)`$ follows from isotony, $`𝒜(M)`$ cannot be a subset of $`𝒜(W^{})^{}`$, and since $`RW`$, it cannot be a subset of $`𝒜(R^{})^{}`$, proving both statements.
$`\mathrm{}`$
In order to investigate the localization behaviour of a single local observable, a further generalization of Landau’s theorem will be used. It is the main result of this section. It is a generalization of Theorem 2.1 in . $`𝒜_{\mathrm{loc}}^d`$ will denote the algebra of local observables of the dual net $`𝒜^d`$.
#### 3.1.5 Theorem (empty-intersection theorem)
> Let $`(W_\nu )_{1\nu n}`$ be a family of $`n`$ wedges in $`𝒲`$. If $`_\nu \overline{W}_\nu =\mathrm{}`$, then
>
> $$𝒜_{\mathrm{loc}}\underset{\nu }{}𝒜(W_\nu ^{})^{}=𝐂\mathrm{id}_{}.$$
Proof. Choose an $`A𝒜_{\mathrm{loc}}_\nu 𝒜(W_\nu ^{})^{}`$, and let $`𝒪`$ be a double cone with $`A𝒜(𝒪)`$ (or $`A𝒜(𝒪^{})^{}`$).
Since the wedges $`\overline{W}_\nu `$ have empty common intersection, so do the compact regions $`\overline{𝒪}\overline{W}_\nu `$. But if a finite family of compact regions have empty common intersection, there is an $`\epsilon >0`$ such that the family of the $`\epsilon `$-neighbourhoods of the regions still have empty common intersection. The proof of this is an elementary induction proof: any two disjoint compact regions have a positive distance, which implies the statement for two regions. Now assume the statement to hold for any family of $`n`$ compact sets, and let $`C_1,\mathrm{},C_{n+1}`$ be a family of n+1 regions. If $`n`$ of these regions already have empty common intersection, there is nothing more to prove. So consider the case that the set $`\mathrm{\Gamma }:=_{\nu =1}^nC_\nu `$ is nonempty. This region is compact and, as shown, has a finite distance $`\delta `$ from $`C_{n+1}`$, so the $`\delta /3`$ neighbourhoods of the two regions still have a finite distance. But the $`\delta /3`$-neighbourhood of $`\mathrm{\Gamma }`$ is the intersection of the $`\delta /3`$-neighbourhoods of $`C_1,\mathrm{},C_n`$, so the statement follows for $`\epsilon =\delta /3`$.
It follows from this that there is a double cone $`P`$ which is so small that the wedge $`\stackrel{~}{W}_\nu :=(W_\nu P)^{\prime \prime }=(W_\nu P)^{cc}`$, $`\nu n`$, and the double cone $`\stackrel{~}{𝒪}:=(𝒪P)^{\prime \prime }=(𝒪P)^{cc}`$ (cf. Lemma 2.2.6 above to see that these regions are a wedge and a double cone, respectively) still have empty common intersection. Choose any $`B𝒜(P)`$. By locality, the commutator function $`f_{A,B}`$ vanishes in the region $`R:=\stackrel{~}{𝒪}^{}_\nu \stackrel{~}{W}_\nu ^{}`$.
There is no admissible mass hyperboloid for this region. To see this, note that if a (shifted) mass hyperboloid is disjoint from a union of wedges, so is the unique shift $`x+\overline{V}`$, $`x𝐑^{1+s}`$, of the closure of the full light cone which contains the hyperboloid. Now choose $`x𝐑^{1+s}`$ such that $`x+\overline{V}`$ is disjoint from all $`\stackrel{~}{W}_\nu ^{}`$, $`\nu n`$, and from $`\stackrel{~}{𝒪}^{}`$, which is a union of wedges, too. This is equivalent to $`\{x\}^{}\stackrel{~}{𝒪}^{}_\nu \stackrel{~}{W}_\nu ^{}`$, i.e.,
$$x\stackrel{~}{𝒪}^{\prime \prime }\underset{\nu }{}\stackrel{~}{W}_\nu ^{\prime \prime }=\stackrel{~}{𝒪}\underset{\nu }{}\stackrel{~}{W}_\nu =\mathrm{}.$$
Hence there is no admissible mass hyperboloid for $`R`$.
If $`R`$ is a Jost-Lehmann-Dyson region, it follows from Theorem 2.2.5 that $`f_{A,B}(x)`$ vanishes for all $`x𝐑^{1+s}`$ and all $`B𝒜(P)`$, so using part (iii) of Lemma 2.2.1, one concludes that $`A𝐂\mathrm{id}_{}`$, and the proof is complete.
But since $`R`$ does not need to be a Jost-Lehmann-Dyson region, Asgeirsson’s lemma will be used to show that the function $`f_{A,B}`$ vanishes in a larger region $`NR`$ which is a Jost-Lehmann-Dyson region. Since there is no admissible hyperboloid for $`R`$, there is, a fortiori, no admissible hyperboloid for $`N`$, so the proof will be complete as soon as $`N`$ has been shown to exhibit the stated properties.
To this end, choose coordinates such that $`\stackrel{~}{𝒪}`$ is the double cone
$$(\rho _0e_0+V_+)(\rho _0e_0V_+),$$
where $`e_0`$ denotes the unit vector in the 0-direction, and $`\rho _0>0`$ is the radius of the double cone $`\stackrel{~}{𝒪}`$. Let $`Z_\rho =\{x=(x_0,\stackrel{}{x})𝐑^{1+s}:\stackrel{}{x}=\rho \}`$ be the boundary of the cylinder of radius $`\rho `$ around the time axis in $`𝐑^{1+s}`$, and define
$`R_{\rho ,0}`$ $`:=`$ $`\stackrel{~}{𝒪}^{}Z_\rho ,`$
$`R_{\rho ,\nu }`$ $`:=`$ $`\stackrel{~}{W}_\nu ^{}Z_\rho ,\nu n;`$
All these regions are bounded subsets of $`𝐑^{1+s}`$. Due to our choice of coordinates, the region $`R_{\rho ,0}`$ is a strip:
$$R_{\rho ,0}=\{xZ_\rho :|x_0|\rho \rho _0\}$$
(which is empty if $`\rho <\rho _0`$).
For $`1\nu n`$, the wedge $`\stackrel{~}{W}_\nu ^{}`$ is timelike convex in $`𝐑^{1+s}`$, so the region $`R_{\rho ,\nu }`$ is timelike convex with respect to the inherited spacetime structure of $`Z_\rho `$. We now show that there is a $`\rho _\nu >0`$ such that the union $`R_{\rho ,0}R_{\rho ,\nu }`$ is timelike convex as well for all $`\rho >\rho _\nu `$. To this end, let $`C`$ be a spacelike hypersurface in $`\stackrel{~}{W}_\nu \stackrel{~}{W}_\nu ^c`$. As a spacelike surface, it is a subset of $`\stackrel{~}{𝒪}^{}`$ up to a compact set. For $`\rho `$ so large that this compact set is enclosed by $`Z_\rho `$ one finds that $`CZ_\rho `$ is a subset of $`R_{\rho ,0}`$. Since $`CZ_\rho `$ is a Cauchy surface in the spacetime $`Z_\rho `$, it follows that $`R_{\rho ,\nu }C`$ and $`R_{\rho ,0}`$ are timelike convex regions in the spacetime $`Z_\rho `$ whose intersection contains a Cauchy surface, so part (v) of Lemma 2.2.4 implies that $`R_{\rho ,0}R_{\rho ,\nu }`$ is timelike convex. This proves that $`\rho _\nu `$ with the stated properties exists for $`1\nu n`$.
Now choose $`\rho >\widehat{\rho }:=\mathrm{max}_\nu \rho _\nu `$, and apply Lemma 2.2.4 (v) another $`n1`$ times to conclude that the region
$$R_\rho :=RZ_\rho =\underset{0\nu n}{}R_{\rho ,\nu }$$
is timelike convex in $`Z_\rho `$. Since the $`𝐑^{1+s}`$-Asgeirsson hull $`\widehat{R}_\rho `$ is open, bounded, and timelike convex, it is a Jost-Lehmann-Dyson region by Lemma 2.2.4 (iii).
On the other hand, the part of $`\stackrel{~}{𝒪}^{}`$ and $`\stackrel{~}{W}_\nu ^{}`$, respectively, which is enclosed by $`Z_\rho `$ is a subset of the $`𝐑^{1+s}`$-Asgeirsson hull $`\widehat{R}_{\rho ,\nu }`$ of $`R_{\rho ,\nu }`$. It follows that
$$RN:=\underset{\rho >\widehat{\rho }}{}\underset{\nu n}{}\widehat{R}_{\rho ,\nu }=\underset{\rho >\widehat{\rho }}{}\widehat{R}_\rho ,$$
and by Asgeirsson’s lemma, $`f_{A,B}`$ vanishes in $`N`$. Since the Jost-Lehmann-Dyson region $`\widehat{R}_\rho `$ increases with $`\rho `$, it follows from Lemma 2.2.4 (vi) that $`N`$ is a Jost-Lehmann-Dyson region. This is what remained to be shown, so the proof is complete.
$`\mathrm{}`$
Actually, the following, slightly stronger version has been established by the preceding proof:
#### 3.1.6 Corollary
> Let $`W_1,\mathrm{}W_n`$ be wedges in $`𝒲`$, and let $`𝒪𝒦`$ be a double cone. If $`\overline{𝒪}_{1\nu n}\overline{W}_\nu =\mathrm{}`$, then
>
> $$𝒜(𝒪^{})^{}\underset{\nu }{}𝒜(W_\nu ^{})^{}=𝐂\mathrm{id}.$$
After completing this article, it was brought to the author’s attention that in 1+3 dimensions, one can also use the results of Thomas and Wichmann for the above proof. As the region $`\stackrel{~}{𝒪}^{}_\nu \stackrel{~}{W}_\nu ^{}`$ is a union of wedges, one can apply Theorem 3.6 in to prove that the function $`f_{A,B}`$ vanishes in the causal closure of this region, which one can check to be all of $`𝐑^{1+s}`$. Their proof, which uses a completely different line of argument, has been written down for 1+3 dimensions only, and it is not evident whether it also holds in other dimensions; verifying this would require to check several hard proofs at the end of . On the other hand the Thomas-Wichmann analysis is more general in other aspects, so the reader interested in the above special problem may find the above argument more direct.
### 3.2 The localization region of a single local observable and the nonempty-intersection theorem
Theorem 3.1.5 prepares for the definition of a localization region for local observables. As the following proposition shows, there are several ‘natural’ ways how to define such a localization region, and it follows from the empty-intersection theorem that all of them yield nonempty localization regions.
#### 3.2.1 Proposition
> Let $`𝒳`$ be any of the classes $`𝒦`$, $``$, $`𝒲`$ and $`𝒞`$. For every $`A𝒜_{\mathrm{loc}}`$ which is not a multiple of the identity, the localization regions
>
> $`𝐋^𝒳(A)`$ $`:=`$ $`{\displaystyle \{\overline{𝒪}:𝒪𝒳:A𝒜(𝒪)^{\prime \prime }\}}`$
> $`L^𝒳(A)`$ $`:=`$ $`{\displaystyle \{\overline{𝒪}:𝒪𝒳:A𝒜(𝒪^{})^{}\}}`$
> are nonempty, causally complete, convex, and compact sets. Between them, one has the following equalities and inclusions:
>
> $$\begin{array}{ccccccc}𝐋^{}(A)& =& 𝐋^𝒦(A)& & L^𝒦(A)& =& L^{}(A)\\ & & & & & & \\ 𝐋^𝒞(A)& =& 𝐋^𝒲(A)& & L^𝒲(A)& =& L^𝒞(A)\end{array}$$
Proof. We start with the proof of the equalities and inclusions. The equalities immediately follow from the definitions, since on the one hand, $`𝒦`$ and $`𝒲𝒞`$, while on the other hand, every region in $``$ is an intersection of double cones in $`𝒦`$ and every region in $`𝒞`$ is an intersection of wedges in $`𝒲`$ (see Section 2.1). The inclusions in the upper and the lower row of the diagram immediately follow from locality. The inclusions in the two columns follow from the fact that every double cone is an intersection of wedges and that, by isotony, an observable contained in the algebra associated with a given double cone is contained in all algebras associated with wedges containing this double cone.
By these inclusions, it is sufficient to prove that $`L^𝒲(A)`$ is nonempty if $`A𝐂\mathrm{id}`$. It already follows from Theorem 3.1.5 that the intersection of the closures of any finite family of wedges whose algebras contain $`A`$ is nonempty. But the family of all wedges whose algebras contain $`A`$ is never finite.
Since $`A`$ is a local observable, there is a double cone $`𝒪`$ with $`A𝒜(𝒪)`$, and it follows from isotony, locality, and the above inclusions that $`L^𝒲(A)\overline{𝒪}`$. But this implies that
$$L^𝒲(A)=\{\overline{𝒪}\overline{W}:W𝒲,A𝒜(W^{})^{}\},$$
which is an intersection of subsets of the compact set $`\overline{𝒪}`$. But if for a class of closed subsets of a compact space, every finite subclass has a nonempty intersection, it follows from the Heine-Borel property that the whole class has a nonempty intersection. Now Corollary 3.1.6 implies the statement.
$`\mathrm{}`$
In the sequel the maps $`𝒜_{\mathrm{loc}}A𝐋^𝒳(A)`$ and $`𝒜_{\mathrm{loc}}AL^𝒳(A)`$ will be referred to as localization prescriptions. Clearly, the localization prescriptions $`𝐋^𝒦`$ and $`L^𝒦`$ coincide if the net satisfies Haag duality, i.e., if $`𝒜(𝒪^{})^{}=𝒜(𝒪)`$ for all $`𝒪𝒦`$, and the prescriptions $`𝐋^𝒲`$ and $`L^𝒲`$ coincide if the net satisfies wedge duality. Furthermore, wedge duality also makes $`L^𝒲`$ coincide with $`L^𝒦`$ by the following lemma (cf. also , Lemma 4.1).
#### 3.2.2 Lemma
> Assume the net $`𝒜`$ to satisfy wedge duality. For every region $`R𝒞`$, one has
>
> $$𝒜(R^{})^{}=\underset{W𝒲_R}{}𝒜(W)^{\prime \prime }=:(R),$$
>
> and the net $``$ satisfies locality.
Proof. We first show that the net $`(𝒜(R^{})^{})_{R𝒞}`$ satisfies locality. This immediately follows from the fact remarked above that if $`R`$ and $`S`$ are spacelike separated regions in $`𝒞`$, there is a wedge $`W𝒲`$ with $`RW`$ and $`SW^{}`$. For such a constellation one has
$$𝒜(R^{})^{}𝒜(W^{})^{}=𝒜(W)^{\prime \prime }𝒜(S^{})^{\prime \prime },$$
which is the stated locality for the net $`(𝒜(R^{})^{})_{R𝒞}`$.
One proves in the same way that the net $``$ satisfies locality with respect to $`𝒜`$, i.e., $`(R)𝒜(R^{})^{}`$ for all $`R𝒞`$. On the other hand,
$$𝒜(R^{})^{}\underset{W𝒲_R}{}𝒜(W^{})^{}=\underset{W𝒲_R}{}𝒜(W)^{\prime \prime }=(R)\text{for all}R𝒞,$$
and this completes the proof.
$`\mathrm{}`$
So if one assumes wedge duality, the localization prescriptions $`L^𝒲(A)`$, $`𝐋^𝒲(A)`$, and $`L^𝒦(A)`$ coincide and provide the smallest localization region out of the above suggestions. In what follows, wedge duality will be assumed, and for every local observable $`A𝒜_{\mathrm{loc}}`$, we simply write $`L^𝒲(A)=𝐋^𝒲(A)=L^𝒦(A)=:L(A)`$.
Now the question arises in how far $`L(A)`$ can be considered as the region where the observable $`A`$ can be measured. For this interpretation to be consistent it is important that the localization prescription $`L`$ satisfies locality in the sense that observables with spacelike separated localization regions commute. This does not follow from the locality assumption made for the net $`𝒜`$. To illustrate this, consider the wedge $`X:=W_1+e_1`$, where $`e_1`$ is the unit vector in the 1-direction, and its images $`Y`$ and $`Z`$ under rotations in the 1-2-plane by $`120^{}`$ and $`240^{}`$, respectively. Assume a local observable $`A`$ to be contained in $`𝒜(X)^{\prime \prime }`$ and in $`𝒜(Y)^{\prime \prime }`$, while another local observable $`B`$ is contained in $`𝒜(Y)^{\prime \prime }`$ and $`𝒜(Z)^{\prime \prime }`$. In this case, the localization regions $`L(A)`$ and $`L(B)`$ are spacelike with respect to each other, but locality of the net alone is not yet sufficient to conclude that $`A`$ and $`B`$ should commute, since not any two of the three wedges are spacelike separated.
Actually, this simplified sketch already points towards the sufficient and necessary condition for locality of $`L`$ provided by the following theorem.
#### 3.2.3 Theorem (nonempty-intersection theorem)
> Assume $`𝒜`$ to satisfy wedge duality.
>
> The localization prescription $`𝒜_{\mathrm{loc}}AL(A)`$ satisfies locality if and only if for every finite family $`W_1,\mathrm{},W_n`$ of wedges and for every causally complete and convex region $`R𝒞`$ with $`_\nu \overline{W}_\nu R`$, one has
>
> $$𝒜_{\mathrm{loc}}\underset{1\nu n}{}𝒜(W_\nu )^{\prime \prime }𝒜(R^{})^{}.$$
Proof. To prove that the condition is sufficient, let $`𝐁_\epsilon (L(A))`$ be the boundary of the open $`\epsilon `$-neighbourhood $`𝐁_\epsilon (L(A))`$ of $`L(A)`$ for $`\epsilon >0`$, and define
$$𝒲_A:=\{W𝒲:X𝒲:A𝒜(X)^{\prime \prime },\overline{X}W\}.$$
A class of closed subsets of the compact space $`𝐁_\epsilon (L(A))`$ is defined by
$$𝒳:=\{𝐁_\epsilon (L(A))\overline{W}:W𝒲_A\}.$$
$`𝒳`$ has empty intersection, and by the Heine-Borel property, there is a finite subclass of $`𝒳`$ whose intersection is still empty, i.e., there are wedges $`W_1,\mathrm{},W_n𝒲_A`$ such that
$$𝐁_\epsilon (L(A))\underset{\nu }{}\overline{W}_\nu =\mathrm{}.$$
Due to the convexity of $`L(A)`$ and of wedges it follows that the region
$$R:=\underset{\nu }{}W_\nu $$
is a subset of $`𝐁_\epsilon (L(A)),`$ and that $`R`$. By the definition of the class $`𝒲_A`$, there are wedges $`X_1,\mathrm{},X_n`$ in $`W_A`$ such that $`\overline{X}_\nu W_\nu `$ for $`1\nu n`$.
Since $`R𝒞`$, one now obtains from the condition that
$$A𝒜_{\mathrm{loc}}\underset{\nu }{}𝒜(X_\nu )^{\prime \prime }𝒜(R^{})^{}𝒜(𝐁_\epsilon (L(A))^{})^{},$$
as stated. This holds for each $`\epsilon >0`$, and evidently, the same reasoning proves that $`B𝒜(𝐁_\epsilon (L(B))^{})^{}`$.
Since $`L(A)`$ and $`L(B)`$ are compact, convex, and spacelike separated, the euclidean distance between these regions is positive, and one can choose $`\epsilon >0`$ so small that $`𝐁_\epsilon (L(A))`$ and $`𝐁_\epsilon (L(B))`$ still are spacelike separated. As remarked in Section 2.1, ther is a wedge $`X`$ such that $`𝐁_\epsilon (L(A))X`$ and $`𝐁_\epsilon (L(B))X^{}`$. Using wedge duality and Lemma 3.2.2, one concludes
$$A𝒜\left(𝐁_\epsilon (L(A))^{}\right)^{}𝒜(X)^{\prime \prime },$$
and
$$B𝒜\left(𝐁_\epsilon (L(B))^{}\right)^{}𝒜(X^{})^{\prime \prime }=𝒜(X)^{},$$
so $`AB=BA`$, proving that the condition is sufficient.
To prove that the condition is necessary, let $`W_1,\mathrm{},W_n`$ be a family of wedges, and choose an $`R𝒞`$ with $`_\nu \overline{W}_\nu R`$. Whenever $`A𝒜_{\mathrm{loc}}_\nu 𝒜(W_\nu )^{\prime \prime }`$ and $`B𝒜_{\mathrm{loc}}𝒜(X)^{\prime \prime }`$ for any $`X𝒲^R^{}`$, locality of $`L`$ implies that $`AB=BA`$, and one concludes that
$`A`$ $`{\displaystyle \underset{X𝒲^R^{}}{}}(𝒜_{\mathrm{loc}}𝒜(X)^{\prime \prime })^{}={\displaystyle \underset{X𝒲^R^{}}{}}𝒜(X)^{}={\displaystyle \underset{X𝒲^R^{}}{}}𝒜(X^{})^{\prime \prime }`$
$`={\displaystyle \underset{X𝒲_R}{}}𝒜(X)^{\prime \prime }=𝒜(R^{})^{},`$
where Lemma 3.2.2 has been used in the last step.
$`\mathrm{}`$
This theorem immediately implies the following corollary.
#### 3.2.4 Corollary
> Assume $`𝒜`$ to satisfy wedge duality, and suppose that the localization prescription $`L`$ satisfies locality. If $`A`$ is a local observable and $`R𝒞`$ is a causally complete convex open region in $`𝐑^{1+s}`$ such that $`L(A)R`$, then $`A𝒜(R^{})^{}`$.
As the following proposition shows, the additional assumption of wedge duality (Assumption (F) above) is sufficient to ensure locality of $`L`$.
#### 3.2.5 Proposition
> Assume $`𝒜`$ to satisfy wedge duality and wedge additivity. Then the localization prescription $`L`$ satisfies locality.
Proof. Let $`A`$ and $`B`$ be local observables with spacelike separated localization regions. There is a wedge $`W`$ such that $`L(A)W`$ and $`L(B)W^{}`$. So as soon as one proves that this implies $`A𝒜(W)^{\prime \prime }`$ and $`B𝒜(W^{})^{\prime \prime }`$, wedge duality implies the statement.
To this end, we consider any $`A𝒜_{\mathrm{loc}}`$ and any wedge $`W`$ whose closure is spacelike separated from $`L(A)`$, and show that $`A𝒜(W)^{}`$. This follows from wedge additivity as soon as one has found a double cone $`P`$ with the property that $`WW+P`$ and that $`f_{A,B}`$ vanishes in $`W`$ for all $`B𝒜(P^{})^{}`$.
So fix an $`\epsilon >0`$ such that the $`\epsilon `$-neighbourhood $`𝐁_\epsilon (L(A))`$ of $`L(A)`$ is still spacelike separated from $`\overline{W}`$. As in the proof of Theorem 3.2.3, we choose a finite number of wedges $`X_1,\mathrm{},X_n`$ in the class $`𝒲_A`$ such that
$$\underset{\nu }{}\overline{X}_\nu 𝐁_\epsilon (L(A)).$$
Now define
$$P:=(\rho e_0+V_+)(\rho e_0V_+)$$
for some $`\rho >0`$ (again, $`e_0`$ denotes the unit vector in the time direction), and note that $`WW+P`$. Fixing $`\rho >0`$ sufficiently small, one can make sure that
$$W\left(\underset{\nu }{}(X_\nu P)\right)^{}.$$
Choosing any $`B𝒜(P^{})^{}`$, one obtains from wedge duality that the commutator function $`f_{A,B}`$ defined above vanishes in the region
$$R:=\underset{\nu }{}(X_\nu P)^{},$$
which is a union of wedges.
As in the proof of Theorem 3.1.5, $`f_{A,B}`$ can be shown to vanish in a larger region $`NR`$ which is a Jost-Lehmann-Dyson region. This can be shown by mimicking the corresponding part of that proof, as it does not depend on the assumption that the intersection of the closed wedges under consideration is empty. So one can keep $`A`$, $`B`$, and the double cone $`P`$, choose some double cone $`𝒪`$ with $`A𝒜(𝒪)`$, replace $`X_1,\mathrm{},X_n`$ by $`W_1,\mathrm{},W_n`$, and proceed like above to construct $`N`$.
A mass hyperboloid $`H`$ is admissible with respect to $`N`$ only if it is admissible with respect to $`R`$, and as $`R`$ is a union of closed wedges, this is the case only if the whole unique shift of the open light cone which contains $`H`$ is disjoint from $`R`$. But by Theorem 2.2.5, this implies that in particular, $`f_{A,B}`$ vanishes in the region $`W`$, completing the proof.
$`\mathrm{}`$
Thomas and Wichmann have obtained a similar result for 1+3 dimensions from slightly stronger assumptions (Theorem 4.10 in ).
One may ask what is the difference between the intersection condition found in Theorem 3.2.3 and the ‘brute force’ condition that $`𝒜(𝒪)𝒜(P)=𝒜(𝒪P)`$ for all $`𝒪,P𝒞`$. Clearly, this condition is the stronger one of the two, and it would imply locality of $`L`$ in a straightforward fashion. Furthermore the property appears so natural that one could expect it to be a general feature of local nets.
In , Landau has given examples of theories which exhibit wedge duality, wedge additivity, and, hence, locality of $`L`$ and the equivalent condition given above, while $`𝒜(𝒪)𝒜(P)𝒜(𝒪P)`$ for $`𝒪,P𝒞`$.
To illustrate the geometrical trick of Landau’s example, start from some local net $``$ of observables in 1+(s+1) dimensions, and with every double cone $`𝒪=(a+V_+)(b+V_{})`$ in $`𝐑^{1+s}`$, associate the algebra
$$_0(𝒪):=((a+\widehat{V}_+)(b+\widehat{V}_{}))=:(\widehat{𝒪}),$$
where, as before, $`\widehat{V}_+`$ and $`\widehat{V}_{}`$ denote the 1+(s+1)-dimensional forward and backward light cone, respectively.
One easily checks that $`_0(𝒪)_0(P)`$ might not coincide with $`_0(𝒪P)`$, since the intersection of the 1+(s+1)-dimensional Asgeirsson hulls of $`𝒪`$ and $`P`$ differs from the 1+(s+1)-dimensional Asgeirsson hull of the intersection $`𝒪P`$, i.e., $`\widehat{𝒪}\widehat{P}\widehat{𝒪P}`$. Indeed, Landau has given examples for theories where the corresponding algebras differ. In particular, they differ if the ‘large’ net $``$ has the intersection property, i.e., if $`(𝒪)(P)=(𝒪P)`$ for all $`𝒪,P`$. This shows that the intersection property cannot be a general property of all local nets of observables.
While Landau’s examples do satisfy all of our above conditions, they illustrate that the sufficient and necessary condition for $`L`$ to be local is not self-evident, as it is similar to the intersection property violated by Landau’s examples. On the other hand, the fact that all our sufficient conditions for the locality of $`L`$ hold, gives some hope that locality of $`L`$ is a rather natural property of local nets.
## 4 Conclusion
Generalizing Landau’s result that the algebras associated with two strictly disjoint double cones have a trivial intersection, the empty-intersection theorem makes it possible to associate a nonempty causally complete, convex and compact localization region with every single local operator of a local net. If one makes the additional assumption of wedge duality, there is a natural way how to obtain a smallest localization region from the empty-intersection theorem. Even in this situation it is a nontrivial issue whether observables with spacelike separated localization regions commute. As a necessary and sufficient condition for this, the nonempty-intersection theorem establishes a special intersection property, and sufficient for this property is the additional condition of wedge additivity, a property typically shared by models arising from Wightman fields. As these results depend on very weak additional assumptions, locality of the localization prescription $`L`$ turns out to be a rather natural property of local nets.
The question what the intersection of two algebras of local observables contains has arised earlier, as, e.g., the remarks in Section III.4.2 of Haag’s monograph show. Haag’s ‘Tentative Postulate’ that the map $`𝒪𝒜(𝒪)`$ be a homomorphism from the orthocomplemented lattice of all causally complete regions (which, in general, are neither bounded nor convex) of Minkowski space into the orthocomplemented lattice of von Neumann algebras on a Hilbert space does not hold in general as it stands (cf. also Haag’s heuristic remarks which illustrate the physical limits of the postulate). But if a net satisfies wedge duality and strong additivity for wedges, the above results, indeed, imply parts of Haag’s conjecture: for arbitrary finite families of wedges, one obtains relations in the spirit of (III.4.7) through (III.4.11) in for the dual net.
The results of this article have been used for the analysis of the Unruh effect and related symmetries of quantum fields . Proceeding, so to speak, in the converse direction, Thomas and Wichmann have investigated the implications that the symmetries providing the Unruh effect exert on the localization behaviour of local observables. Assuming the theory to exhibit the Unruh-effect and a couple of (standard) technical properties including wedge additivity, they found that the localization region of an observable $`A`$ with respect to a minimal Poincaré covariant local net generated by $`A`$ is the smallest region $`𝒪_A`$ in $``$ with the property that for every $`(a,\mathrm{\Lambda })𝒫_+^{}`$, one has $`a+\mathrm{\Lambda }𝒪_A𝒪_A^{}`$ if and only if $`[A,U(a,\mathrm{\Lambda })AU(a,\lambda )^{}]=0`$ . This definition of a localization region no longer explicitly refers to any other operators of the net (while it does refer to the representation, which, by the Bisognano-Wichmann symmetries, is closely related to the net). While this interesting conclusion has been derived from the Unruh effect and other assumptions of relevance in the above discussion, all these assumptions have been avoided above since they are a goal rather than a starting point of the above analysis. In this sense, the results of Thomas and Wichmann are complementary to the above results.
## Appendix
For the reader’s convenience we include a proof that the Reeh-Schlieder property entails weak additivity (cf. also Lemma 2.6 in ).
### Lemma
> Let $`𝒜`$ be a local net of local observables satisfying Conditions (A) through (D) above, let $`𝒪𝐑^{1+s}`$ be a bounded open region, and let $`a𝐑^{1+s}`$ be some timelike vector. Then
>
> $$𝒞_{𝒪,a}:=\left(\underset{t𝐑}{}𝒜(𝒪+ta)\right)^{\prime \prime }=().$$
Proof. For any $`a`$ and $`𝒪`$ as above, let $`A`$ be any local observable commuting with all elements of $`𝒞_{𝒪,a}`$, and pick a $`B𝒜(𝒪)`$. Define $`f_+(t):=\mathrm{\Omega },A^{}U(ta)B\mathrm{\Omega }`$ and $`f_{}(t):=\mathrm{\Omega },BU(ta)A^{}\mathrm{\Omega }`$. By the spectral theorem and the spectrum condition, the Fourier transforms of these functions are (not necessarily positive, but bounded) measures one of which has its support in the closed positive half axis, while the other one has its support in the closed negative half axis. Since $`f_+`$ and $`f_{}`$ coincide by construction, it follows that the Fourier transform of $`f_+`$ (and of $`f_{}`$) is a measure with support $`\{0\}`$, i.e., some multiple of the Dirac measure, so that $`f_+`$ is a constant function. Using this, the spectral theorem, and uniqueness of the vacuum, one concludes
$$A\mathrm{\Omega },B\mathrm{\Omega }=f_+(0)=f_+(t)=\mathrm{\Omega },A\mathrm{\Omega }\mathrm{\Omega },B\mathrm{\Omega }=:\overline{\alpha }\mathrm{\Omega },B\mathrm{\Omega }=\alpha \mathrm{\Omega },B\mathrm{\Omega }.$$
Since by the Reeh-Schlieder property, $`B\mathrm{\Omega }`$ runs through a dense set, one concludes $`A\mathrm{\Omega }=\alpha \mathrm{\Omega }`$, and since $`A`$ is a local observable, one obtains $`A=\alpha \mathrm{id}`$ since $`\mathrm{\Omega }`$ is cyclic with respect to $`𝒜(𝒪^{})`$, which implies that it is separating with respect to $`𝒜(𝒪)`$ dut to locality. This proves the lemma.
$`\mathrm{}`$
## Acknowledgements
It was a very important help that D. Arlt, W. Kunhardt, and N. P. Landsman read preliminary versions of the manuscript very carefully. I would also like to thank K. Fredenhagen for helpful discussions and his encouraging interest. During a visit to Göttingen I obtained helpful hints from M. Lutz and M. Requardt.
The above research is part of a project that has been re-initiated at the Erwin Schrödinger Institute in Vienna, where I had the occasion to join the project ‘Local Quantum Physics’ in the autumn of 1997, and a ‘Nachlese’ meeting in the spring of 1999. I would like to thank the ESI for the kind invitations and the warm hospitality exhibited to me. At the ESI I profited from discussions with H.-J. Borchers, D. Guido, S. Trebels, and E. Wichmann.
This work has been supported by the Deutsche Forschungsgemeinschaft, the European Union’s TMR Network ‘Noncommutative Geometry’, a Feodor-Lynen grant of the Alexander von Humboldt foundation, a part of which has been funded by the University of Amsterdam, and a Casimir-Ziegler award of the Nordrhein-Westfälische Akademie der Wissenschaften.
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# Theorem 0.0.1 (Z)
## 1 Preliminary results on $`\overline{\rho }_{2,\alpha ,\beta }^N`$
Up until the last step, the analysis of $`\overline{\rho }_{2,\alpha ,\beta }^N`$ is analogous to the analysis of $`\rho _{2,\alpha ,\beta }^N`$ in \[Z\]. As in \[Z, Theorem (5.1.1)\] we have:
###### Theorem 1.0.1
Let $`\widehat{H}_{\alpha ,\beta }=\alpha \varphi (\widehat{I})+\beta \widehat{I}`$ where $`|\varphi ^{\prime \prime }|C_o>0`$ on $`[1,1].`$ Let $`\overline{\rho }_{2,\alpha ,\beta }^N`$ be as above. Then for any $`f`$ with supp$`\widehat{f}`$ compact:
$$_T^T_T^T|\overline{\rho }_{2,\alpha ,\beta }^N(f)\rho _2^{POISSON}(f)|^2𝑑\alpha 𝑑\beta =O(\frac{(\mathrm{log}N)^2}{N}).$$
###### Corollary 1.0.2
Let $`N_m=[m(\mathrm{log}m)^4]`$. Then for almost all $`(\alpha ,\beta )`$ in the Lebesgue sense,
$$\underset{m\mathrm{}}{lim}\overline{\rho }_{2,\alpha ,\beta }^{N_m}(f)=\rho _2^{POISSON}(f).$$
To fill in the gaps in the sparse susequence $`\{N_m\}`$, consider $`\overline{\rho }_{2,\alpha ,\beta }^M`$ for $`N_m<M<N_{m+1}.`$ Obviously,
$$\overline{\rho }_{2,\alpha ,\beta }^M(f)\overline{\rho }_{2,\alpha ,\beta }^{N_m}(f)=\frac{N_mM}{M}\overline{\rho }_{2,\alpha ,\beta }^{N_m}(f)+\frac{1}{M}\underset{n=N_m}{\overset{M}{}}\rho _{2,\alpha ,\beta }^n(f).$$
(2)
We have $`MN_m<<(N_{m+1}N_m)(m+1)(\mathrm{log}(m+1)^4)m(\mathrm{log}m)^4<<(\mathrm{log}m)^4.`$ So in the first sum $`\frac{N_mM}{M}<<m^{1+ϵ}.`$ In the second we have $`O((\mathrm{log}m)^4)`$ terms. Under the assumption supp$`\widehat{f}[1,1]`$ the trivial bound $`\rho _2^n(f)<<n`$ already gives
$$\frac{N_mM}{M}\overline{\rho }_2^{N_m}(f)+\frac{1}{M}\underset{n=N_m}{\overset{M}{}}\rho _{2,\alpha ,\beta }^n(f)<<(MN_m)<<(\mathrm{log}m)^4.$$
(3)
So we just need a tiny improvement on the trivial bound to prove that these terms tend to zero. In the following section we will prove that for almost all $`(\alpha ,\beta )`$, $`\rho _{2,\alpha ,\beta }^n(f)C(\alpha ,\beta )n^{1\frac{2}{K}+ϵ}`$ where $`K=2^{k1}`$ with $`k`$ the degree of $`\varphi `$. From this it also follows by standard density arguments that $`\overline{\rho }_{2,\alpha ,\beta }^N[a,b]\rho _2^{POISSON}[a,b]`$ for all intervals $`[a,b]`$. We refer to \[R.S\] for the details of the density argument.
## 2 The Main Lemma
The purpose of this section is to prove:
###### Lemma 2.0.1
Suppose that $`\varphi `$ is a polynomial of degree $`k`$ satisfying the hypotheses: (i) $`|\varphi ^{\prime \prime }|>0`$ and (ii) $`|\alpha \varphi ^{}+\beta |>0`$ on $`[1,1]`$, . Then for any $`\widehat{f}C_o(\mathrm{I}\mathrm{R})`$ and almost all $`(\alpha ,\beta )`$, we have: $`n^2\overline{\rho }_{2,\alpha ,\beta }^{(n)}(f)C(\alpha ,\beta )n^{1\frac{2}{K}+ϵ}`$, where $`K=2^{k1}`$.
Recall that the local PCF’s have the form
$$\rho _{2,\alpha ,\beta }^n=\underset{\mathrm{}ZZ}{}\widehat{f}(\frac{\mathrm{}}{n})|\underset{k=1}{\overset{n}{}}e(\alpha n\mathrm{}[\varphi (\frac{k}{n})+\beta \frac{k}{n}])|^2.$$
Since $`\widehat{f}`$ is compactly supported, the $`\mathrm{}`$-sum runs over an interval of integers of the form $`[Cn,Cn]`$ for some $`C>0.`$ For simplicity of notation, and with no loss of generality, we will assume the sum over $`\mathrm{}`$ runs over the interval $`[n,n].`$ Throughout we use the notation $`e(x)=e^{2\pi ix}.`$
### 2.1 The quadratic case
The case of quadratic polynomials is more elementary than that of polynomials of general degree and we can prove our main result without analysing continued fraction convergents to $`\alpha `$. Hence we begin by discussing this case. The relevant exponential sum is
$$|\underset{k=1}{\overset{n}{}}e(\alpha n\mathrm{}[\varphi (\frac{k}{n})+\beta \frac{k}{n}])|^2=\underset{h=n}{\overset{n}{}}\underset{x=1}{\overset{2n}{}}e(\mathrm{}h(\alpha \frac{x}{n}+\beta )).$$
For $`f`$ with supp$`\widehat{f}`$ in $`[1,1]`$ we have
$$n^2\rho _{2,\alpha ,\beta }^n(f)<<|\underset{|\mathrm{}|n}{}\underset{h=n}{\overset{n}{}}\underset{x=1}{\overset{2n}{}}e(\mathrm{}h(\alpha \frac{x}{n}+\beta ))|.$$
The following estimate is weaker than that claimed in the Main Lemma but is sufficient for the proof of the theorem.
###### Lemma 2.1.1
Let $`\alpha `$ be a diophantine number satisfying $`|\alpha \frac{a}{q}|\frac{K(\alpha )}{q^{2+ϵ}}`$ for any rational number $`\frac{a}{q}.`$ Then for all $`\beta `$, $`\rho _{2,\alpha ,\beta }^n(f)<<n^{\frac{1}{2}+ϵ}.`$
Proof:
We begin with the standard estimate (e.g. \[K, Lemma 1\])
$$|\underset{x=1}{\overset{2n}{}}e(\mathrm{}h(\alpha \frac{x}{n}+\beta ))|=|\underset{x=1}{\overset{2n}{}}e(\mathrm{}h(\alpha \frac{x}{n}))|\mathrm{min}(2n,\frac{1}{2\mathrm{}h\frac{\alpha }{n}})$$
where $`||||`$ denotes the distance to the nearest integer. This gives
$$n^2\rho _{2,\alpha ,\beta }^n(f)<<\underset{\mathrm{}n}{}\underset{h=n}{\overset{n}{}}\mathrm{min}(2n,\frac{1}{2\mathrm{}h\frac{\alpha }{n}}).$$
The variable $`x=h\mathrm{}`$ runs over $`[n^2,n^2]`$; when $`x0`$, the multiplicity $`c_x=\mathrm{\#}\{(h,\mathrm{}):h\mathrm{}=x\}`$ is well-known to have order $`n^ϵ`$ (e.g \[V, Lemma 2.5\]). Then there are $`2n`$ terms where $`h\mathrm{}=0`$, each contributing $`n`$ to the sum. Hence,
$$n^2\rho _{2,\alpha ,\beta }^n(f)<<n^2+n^ϵ\underset{x=n^2}{\overset{n^2}{}}\mathrm{min}(2n,\frac{1}{2x\frac{\alpha }{n}}).$$
(4)
At this point we are close to the well-known estimate ( e.g. Korobov \[K, Lemma 14\])
$$\underset{x=1}{\overset{Q}{}}\mathrm{min}(P,\frac{1}{\alpha x+\beta })<<(1+\frac{Q}{q})(P+q\mathrm{log}P)$$
where $`\alpha =\frac{a}{q}+\frac{\theta }{q^2}`$ with $`|\theta |<1`$ and with $`(a,q)=1.`$ In our situation $`Q=n^2,P=n`$, giving $`(1+\frac{n^2}{q})(n+q\mathrm{log}n)`$, but the estimate does not apply because our ‘$`\alpha `$’ is $`\frac{\alpha }{n}`$; the rational approximation $`\frac{a}{qn}`$ to $`\frac{\alpha }{n}`$ has a remainder of only $`\frac{1}{nq^2}`$ rather than $`\frac{1}{(nq)^2}.`$ This complicates the argument and worsens the resulting estimate.
Since we do not know the continued fraction expansion of $`\frac{\alpha }{n}`$, we use the rational approximation $`\frac{\alpha }{n}=\frac{a}{qn}+\frac{\theta }{nq^2}`$. It is not necessary that $`(a,n)=1`$ so we rewrite $`\frac{a}{qn}=\frac{a^{}}{qn^{}}`$ with $`(a^{},n^{})=1`$ (hence $`(a^{},n^{}q)=1`$). Then
$$\frac{\alpha }{n}=\frac{a^{}}{n^{}q}+\frac{\theta }{nq^2},(a^{},n^{})=1,|\theta |<1.$$
Now break up $`[n^2,n^2]`$ into blocks of length $`n^{}q`$. There are at most $`2[\frac{n^2}{n^{}q}]+1`$ such blocks. Hence
$$n^ϵ\underset{x=n^2}{\overset{n^2}{}}\mathrm{min}(2n,\frac{1}{2x\frac{\alpha }{n}})<<n^ϵ\underset{y=0}{\overset{[\frac{n^2}{n^{}q}]+1}{}}\underset{x=1}{\overset{n^{}q}{}}\mathrm{min}(2n,\frac{1}{2(x+yqn^{})\frac{\alpha }{n}}).$$
(5)
The above rational approximation brings
$$\frac{\alpha x}{n}+yqn^{}\frac{\alpha }{n}=\frac{a^{}x}{n^{}q}+\frac{x\theta }{nq^2}+ya^{}+\frac{yn^{}\theta }{nq}.$$
Hence
$$\frac{\alpha x}{n}+yqn^{}\frac{\alpha }{n}=\frac{a^{}x}{n^{}q}+\frac{x\theta }{nq^2}+\beta $$
where $`\beta =\{\frac{yn^{}\theta }{nq}\}`$. Write $`\beta =\frac{b(y)}{n^{}q}+\frac{\theta _1}{n^{}q}`$ with $`b(y)ZZ`$ and with $`|\theta _1|<1`$. Since $`|x|n^{}q`$ we have
$$\frac{a^{}x+b(y)}{n^{}q}=x\frac{\alpha }{n}+yqn^{}\frac{\alpha }{n}\frac{x\theta }{nq^2}\frac{\theta _1}{n^{}q}x\frac{\alpha }{n}+yqn^{}\frac{\alpha }{n}+\frac{1}{n^{}q}+\frac{n^{}}{nq}.$$
The remainder $`\frac{n^{}}{nq}`$ is much larger than occurs in the standard argument and since it is possible that $`n^{}=n`$ we can only be sure that the remainder is $`O(\frac{1}{q}).`$
Therefore we are only sure that our sum is
$$<<n^ϵ\underset{y=0}{\overset{[\frac{n^2}{n^{}q}]+1}{}}\underset{x=1}{\overset{n^{}q}{}}\mathrm{min}(2n,\frac{1}{2\frac{a^{}x+b(y)}{n^{}q}+O(\frac{1}{q})}).$$
Since $`(a^{},n^{}q)=1`$, the numbers $`a^{}x+b(y)`$ run thru a complete residue system modulo $`n^{}q`$ as $`x`$ runs thru $`1,\mathrm{}n^{}q`$. Hence, the $`x`$-sum is independent of $`a^{},b(y)`$ and we may rewrite it as
$$<<(\frac{n^{2+ϵ}}{n^{}q}+1)\underset{2xn^{}q1}{}\mathrm{min}(2n,\frac{2}{\frac{x}{n^{}q}+O(\frac{1}{q})}).$$
The distance $`\frac{x}{n^{}q}+O(\frac{1}{q})`$ can be less than $`\frac{1}{n}`$ over the range of terms $`x[0,Cn^{}]`$ and $`x[n^{}qCn^{},n^{}q]`$ where $`C`$ is the implicit constant in $`O(\frac{1}{q}).`$ For these we must take $`n`$ in the minimum. Since there are $`O(n)`$ such terms in the $`x`$-sum, their contribution to the entire sum is $`<<n^{2+ϵ}\frac{n^2}{n^{}q}.`$
For the remaining terms we use that $`\mathrm{min}(2n,\frac{2}{\frac{x}{n^{}q}})`$ is an even function of $`x`$ to put the $`x`$-sum in the form
$$\underset{Cn^{}x\frac{qn^{}}{2}}{}\mathrm{min}(2n,\frac{2}{\frac{x}{n^{}q}+O(\frac{1}{q})}).$$
The minimum is now surely attained by $`\frac{2}{\frac{x}{n^{}q}+O(\frac{1}{q})}`$ and since it stays in the left half of the interval we have
$$\frac{1}{\frac{x}{n^{}q}+O(\frac{1}{q})}=\frac{1}{\frac{x}{n^{}q}+O(\frac{1}{q})}.$$
Therefore
$$\underset{Cn^{}x\frac{qn^{}}{2}}{}\mathrm{min}(2n,\frac{2}{\frac{x}{n^{}q}+O(\frac{1}{q})})<<n^{}q\underset{Cn^{}x\frac{qn^{}}{2}}{}\frac{1}{xO(n^{})}<<n^{}q\mathrm{log}(n^{}q).$$
The whole $`x`$-sum is therefore $`<<(\frac{n^{2+ϵ}}{n^{}q}+1)[n^2+n^{}q\mathrm{log}(n^{}q)].`$
In sum, we have
$$n^ϵ\underset{x=n^2}{\overset{n^2}{}}\mathrm{min}(2n,\frac{1}{x\frac{\alpha }{n}})<<(\frac{n^{2+ϵ}}{n^{}q}+1)[n^2+qn^{}\mathrm{log}(n^{}q)].$$
Hence
$$\rho _{2,\alpha ,\beta }^n<<1+(\frac{n^ϵ}{n^{}q}+n^2)[n^2+qn^{}\mathrm{log}(n^{}q)].$$
The first parenthetical term is of size $`n^{1+ϵ}/q`$ when $`n^{}=n`$ while the trivial bound was $`n`$. It is at this point that we must restrict to diophantine numbers satisfying $`|\alpha \frac{a}{q}|\frac{K(\alpha )}{q^{2+ϵ}}`$ for all rational $`\frac{p}{q}.`$ By Dirichlet’s box principle there exists $`qn^r`$ and a rational $`\frac{a}{q}`$ with $`(a,q)=1`$ such that $`|\alpha \frac{a}{q}|\frac{1}{qn^r}.`$ It follows that $`q>n^{rϵ}.`$ Substituting into our estimate, we get
$$\rho _{2,\alpha ,\beta }^n<<1+(\frac{n^{r+ϵ}}{n^{}}+n^2)[n^2+n^rn^{}\mathrm{log}(n)]<<n^ϵ((a,n)n^{1r}+\frac{1}{(a,n)}n^{1+r}.$$
Since $`1(a,n)n`$ the final estimate is
$$<<n^ϵ(n^{2r}+n^{1+r}).$$
The terms balance when $`r=\frac{3}{2}`$ to give
$$\rho _{2,\alpha ,\beta }^n(f)<<n^{\frac{1}{2}+ϵ}.$$
Remark In the next section we will see that there are rational numbers $`\frac{a}{q}`$ satisfying the above requirements and also satisfying $`(a,n)C(\alpha )n^ϵ.`$ This changes the final estimate to $`<<n^ϵ(n^{1r}+n^{1+r})`$ and gives $`\rho _{2,\alpha ,\beta }^n(f)<<n^ϵ.`$
### 2.2 The general polynomial case
Now let $`\varphi (x)=\alpha _ox^k+\alpha _1x^{k1}\mathrm{}+\alpha _k`$ be a general polynomial. We would like to estimate
$$\rho _2^n(f)=\frac{1}{n^2}\underset{\mathrm{}ZZ}{}\widehat{f}(\frac{\mathrm{}}{n})|\underset{k=1}{\overset{n}{}}e(n\mathrm{}\varphi (\frac{x}{n}))|^2.$$
As in the classical Weyl inequality (cf. \[V, Lemma 2.4\]) we will estimate $`|_{k=1}^ne(n\mathrm{}\varphi (\frac{x}{n}))|^2`$ by squaring and differencing repeatedly until we reach the linear case. Let $`\mathrm{\Delta }_j`$ be the jth iterate of the forward difference operator, so that
$$\begin{array}{c}\mathrm{\Delta }_1\varphi (x;h)=\varphi (x+h)\varphi (x)\hfill \\ \mathrm{\Delta }_{j+1}\varphi (x;h_1,\mathrm{},h_{j+1})=\mathrm{\Delta }_1(\mathrm{\Delta }_j\varphi (x;h_1,\mathrm{},h_j;h_{j+1})).\hfill \end{array}$$
We recall (cf. \[V, Lemma 2.3\]):
###### Lemma 2.2.1
We have
$$|\underset{x=1}{\overset{n}{}}e(f(x))|^{2^j}(2n)^{2^jj1}\underset{|h_1|<n}{}\mathrm{}\underset{|h_j|<n}{}[\underset{xI_j}{}e(\mathrm{\Delta }_jf(x;h_1,\mathrm{},h_j))]$$
where the intervals $`I_j=I_j(h_1,\mathrm{},h_j)`$ satisfy $`I_1[1,n]`$, $`I_jI_{j1}.`$
Now let
$$T(\varphi ;n,\mathrm{})=\underset{x=1}{\overset{n}{}}e(n\mathrm{}\varphi (\frac{x}{n}))$$
with $`\varphi (x)=\alpha _ox^k+\mathrm{}+\alpha _o`$ and put $`K=2^{k1}.`$ Apply the previous lemma with $`j=k1`$ to get:
$$|T(\varphi ;n,\mathrm{})|^K<<n^{Kk}\times $$
$$\underset{h_1}{}\mathrm{}\underset{h_{k1}}{}\underset{xI_{k1}}{}e(h_1\mathrm{}h_{k1}\mathrm{}p_{k1}(x;h_1,\mathrm{},h_{k1};n,\mathrm{})).$$
Here, the sum runs over $`h_j`$ with $`|h_j|n`$ and
$$p_{k1}(x;h_1,\mathrm{},h_{k1};n)=k!n^{k+1}\alpha _o(x+\frac{1}{2}h_1+\mathrm{}+\frac{1}{2}h_{k1})+(k1)!n^{k+2}\alpha _1.$$
This is just as in the standard Weyl estimate (\[V\]\[D, §3\]) except for the powers of $`n`$ in the coefficients of $`p_{k1}.`$
Then write
$$\rho _2^n(f)=\frac{1}{n}\underset{\mathrm{}}{}\widehat{f}(\frac{\mathrm{}}{n})[\frac{1}{n}|T(\varphi ;n,\mathrm{})|^2]<<\frac{1}{n}\underset{\mathrm{}n}{}(\frac{1}{n}|T(\varphi ;n,\mathrm{})|^2)$$
(6)
Since the $`\mathrm{}`$-sum is an average, we may apply Holder’s inequality with exponent $`\frac{K}{2}`$ to get
$$\rho _2^n(f)<<[\frac{1}{n}\underset{\mathrm{}n}{}|\frac{1}{\sqrt{n}}T(\varphi ;n,\mathrm{})|^K]^{\frac{2}{K}}$$
(7)
Therefore
$$[\rho _2^n(f)]^{\frac{K}{2}}<<n^{Kk}n^{\frac{K}{2}1}\underset{\mathrm{}n}{}\underset{h_1}{}\mathrm{}\underset{h_{k1}}{}\underset{xI_{k1}}{}e(h_1\mathrm{}h_{k1}\mathrm{}p_{k1}(x;h_1,\mathrm{},h_{k1};n)).$$
There are $`n^{k1}`$ terms with $`h_1\mathrm{}h_{k1}\mathrm{}=0`$, each contributing $`n`$ to the $`x`$-sum. So the contributions of such terms to the total sum is $`O(n^k)`$, and we get
$$[\rho _2^n(f)]^{\frac{K}{2}}<<n^{\frac{K}{2}k1}[n^k+\underset{\mathrm{}n}{}\underset{h,x}{\overset{}{}}e(h_1\mathrm{}h_{k1}\mathrm{}p_{k1}(x;h_1,\mathrm{},h_{k1};n))]$$
(8)
where the primed sum runs only over non-zero values of $`h_1\mathrm{}h_{k1}\mathrm{}.`$
As in the case with $`k=2`$ above we sum over $`x`$ to get
$$[\rho _2^n(f)]^{\frac{K}{2}}<<n^{\frac{K}{2}k1}[n^k+\underset{\mathrm{}n}{}\underset{h}{\overset{}{}}min(n,\frac{1}{k!h_1\mathrm{}h_{k1}\mathrm{}n^{k+1}\alpha })]$$
(9)
and then rewrite the variable $`k!h_1\mathrm{}h_{k1}\mathrm{}`$ as a new variable $`x`$ ranging over $`[0,k!n^k].`$ As before, the number $`c_x`$ of ways of representing $`x0`$ as a product $`k!h_1\mathrm{}h_{k1}\mathrm{}`$ is $`O(n^ϵ)`$ so
$$[\rho _2^n(f)]^{\frac{K}{2}}<<n^{\frac{K}{2}k1+ϵ}[n^k+\underset{xk!n^k}{}min(n,\frac{1}{xn^{k+1}\alpha })].$$
(10)
We now repeat the steps of the quadratic case but with $`\frac{\alpha }{n^{k1}}`$ replacing $`\frac{\alpha }{n}.`$ Thus, the rational approximation $`\alpha =\frac{a}{q}+\frac{\theta }{q^2}`$ gives the approximation $`\frac{\alpha }{n^{k1}}=\frac{a}{n^{k1}q}+\frac{\theta }{n^{k1}q^2}`$ and hence requires us to break up the sum over $`[0,k!n^k]`$ into blocks of size $`n^{k1}q/(a,n^{k1})`$. Precisely the same argument (with $`n_k^{}=\frac{n^{k1}}{(a,n^{k1})}`$) then gives
$$\underset{xk!n^k}{}min(n,\frac{1}{xn^{k+1}\alpha })<<(\frac{n^k}{qn_k^{}}+1)(n^k+qn_k^{}\mathrm{log}(qn_k^{})).$$
Hence we get
$$[\rho _2^n(f)]^{\frac{K}{2}}<<n^{\frac{K}{2}k1+ϵ}[n^k+(\frac{n^k}{qn_k^{}}+1)(n^k+qn_k^{}\mathrm{log}(qn_k^{}))]<<n^{\frac{K}{2}k1+ϵ}[n^k+\frac{n^{2k}}{qn_k^{}}+qn_k^{}].$$
(11)
Recalling that $`n_k^{}=\frac{n^{k1}}{(a,n^{k1})}`$ the last expression is
$$<<n^{\frac{K}{2}1+ϵ}[1+\frac{n^k(a,n^{k1})}{qn^{k1}}+\frac{q}{n(a,n^{k1})}].$$
Thus,
$$[\rho _2^n(f)]<<n^{1\frac{2}{K}+ϵ}[1+\frac{n(a,n^{k1})}{q}+\frac{q}{n(a,n^{k1})}]^{\frac{2}{K}}$$
(12)
The exponent of the right side will be less than one if and only if the exponent of $`[1+\frac{n(a,n^{k1})}{q}+\frac{q}{n(a,n^{k1})}]`$ is less than one. Thus we are in very much the same situation as in the quadratic case (although the resulting exponent will be increasingly bad as $`K\mathrm{}`$). However, the estimate $`(a,n)n`$ used in the quadratic case does not generalize well to higher degree: In higher degree, the estimate $`(a,n^{k1})n^{k1}`$ leads to $`r=\frac{k+1}{2}`$ and an exponent larger than one. Therefore we need to choose a rational approximation satisfying $`(a,q)=1`$ and $`|\alpha \frac{a}{q}|<\frac{1}{q^2}`$ and with low value of $`(a,n^{k1})`$. The natural candidates for such numbers are the continued fraction convergents $`\frac{p_m}{q_m}=[a_o,a_1,\mathrm{},a_m]`$ to $`\alpha =[a_o,a_1,\mathrm{}].`$ Therefore we need to study the behaviour of
$$f_n(\alpha ):=\mathrm{min}\{\frac{n(p_m(\alpha ),n^{k1})}{q_m(\alpha )}+\frac{q_m(\alpha )}{n(p_m(\alpha ),n^{k1})}\}.$$
(13)
Since $`\frac{p_m}{q_m}=\alpha +O(\frac{1}{q_m^2})`$ we can (and will) replace the $`q_m`$ in this definition by $`p_m`$ Since it is presumably hard to arrange for $`(p_m(\alpha ),n^{k1})`$ to be large, we will require that $`p_m(\alpha )[n^{rϵ},n^r]`$ for some exponent $`r`$ to be determined later. Before proceeding let us recall how the index $`m`$ is related to $`n,r`$.
###### Proposition 2.2.2
For any $`r,ϵ>0`$, any $`M𝐍`$ and almost any $`\alpha \mathrm{I}\mathrm{R}`$, there exists $`n_o𝐍`$ with the following property: for $`nn_o`$ there exist at least $`M`$ consecutive convergents $`p_{mM}(\alpha ),p_{mM+1}(\alpha ),\mathrm{},p_m[n^{rϵ},n^r]`$ with $`mC(\alpha )\mathrm{log}n.`$
Proof: By a theorem of Khinchin and Levy \[Kh\], one knows that for almost all $`\alpha `$ the convergents satisfy
$$\underset{m\mathrm{}}{lim}q_m^{\frac{1}{m}}=\gamma ,\gamma :=\frac{\pi ^2}{12\mathrm{log}12}.$$
(14)
The first claim is equivalent to the statement that there exists $`m`$ such that, for $`0jM,`$
$$(rϵ)\mathrm{log}n<\mathrm{log}p_{mj}=m\mathrm{log}\gamma +o(m)<r\mathrm{log}n.$$
Evidently there exists $`C(\alpha )>0`$ such that $`mrC(\alpha )\mathrm{log}n`$, proving the second claim. The first claim is states that for sufficiently large $`n`$, there are at least $`k`$ consecutive solutions $`m`$ of
$$[\frac{(rϵ)}{\gamma }+o(1)]\mathrm{log}nm[\frac{r}{\gamma }+o(1)]\mathrm{log}n.$$
This is obvious since the width of the interval equals $`[\frac{ϵ}{\gamma }+o(1)]\mathrm{log}n`$, which is positive and unbounded.
We then have:
###### Proposition 2.2.3
Fix $`k,r,ϵ>0`$. Then for almost all $`\alpha \mathrm{I}\mathrm{R}`$ there exists a convergent $`\frac{p_m(\alpha )}{q_m(\alpha )}`$ with $`p_m(\alpha )[n^{rϵ},n^r]`$ and with $`(p_m(\alpha ),n^{k1})n^ϵ.`$
Proof By the previous proposition, for any $`M>0`$, there are at least $`M`$ consecutive $`p_m`$’s in $`[n^{rϵ},n^r]`$ for sufficiently large $`n`$. Our goal is to find one satisfying $`(p_m(\alpha ),n^{k1})n^{1+ϵ}.`$
To this end we recall \[Kh\] that
$$\{\begin{array}{c}p_m=a_mp_{m1}+p_{m2}\hfill \\ q_m=a_mq_{m1}+q_{m2}\hfill \end{array}$$
and hence that $`p_mq_{m1}p_{m1}q_m=\pm 1.`$ It follows that $`p_m(\alpha ),p_{m1}(\alpha )`$ are relatively prime. This pattern continues in a sufficiently useful way. By a simple induction we find that for $`k<m`$,
$$p_mq_{mk}p_{mk}q_k=\pm E_{k1}(a_m,a_{m1},\mathrm{},a_{mk+1})$$
(15)
where $`E_0=1,E_1(a_m)=a_m,E_2(a_m,a_{m1})=a_ma_{m_1}+1`$ and where
$$E_k(a_m,a_{m1},\mathrm{},a_{mk})=a_{mk}E_{k1}(a_m,a_{m1},\mathrm{},a_{mk+1})+E_{k2}(a_m,a_{m1},\mathrm{},a_{mk+2}).$$
Hence any common divisor of $`p_m,p_{m1},p_{m2}`$ is a divisor of $`a_m`$, and so on.
We now claim that for the $`M`$ consective $`p_m`$’s in $`[n^{rϵ},n^r]`$ we have:
$$(p_{mM},n^{k1})(p_{mM+1},n^{k1})\mathrm{}(p_m,n^{k1})n^{k1}\mathrm{\Pi }_{j=0}^M\mathrm{\Pi }_{\mathrm{}=1}^{Mj}E_{\mathrm{}}(a_{mj},a_{mj1},\mathrm{},a_{mj\mathrm{}+1})$$
(16)
The idea of the argument is that, were all the $`p_{mj}`$’s relatively prime, then each $`(p_{mj},n^{k1})`$ would contribute a distinct factor of $`n^{k1}`$ and hence the product would be $`n^{k1}.`$ The $`p_{mj}`$’s are of course not relatively prime but (15) gives an upper bound on the greatest common divisors of each pair.
Thus, let us start with $`p_m`$ and consider the degree to which factors in $`(p_m,n^{k1})`$ are replicated by the lower $`(p_{mj},n^{k1})`$’s. Since $`(p_m,p_{m1})=1`$ there is no duplication of factors due to the nearest neighbor. Since $`(p_m,p_{m2})|a_m`$ the greatest common factor of $`(p_{m2},n^{k1}),(p_m,n^{k1})`$ is less than $`(a_m,n^{k1})`$ and hence less than $`a_m.`$ Similarly the greatest common factor of $`(p_{m3},n^{k1}),(p_m,n^{k1})`$ is less than $`E_2(a_m,a_{m1}).`$ In all, the product $`(p_{mM},n^{k1})(p_{mM+1},n^{k1})\mathrm{}(p_m,n^{k1})`$ replicates factors of $`(p_m,n^{k1})`$ by at most $`E_1(a_m)\mathrm{}E_M(a_m,a_{m1},\mathrm{},a_{mM+1}).`$
Next, move on to $`(p_{m1},n^{k1}).`$ These factors of $`n^{k1}`$ can get duplicated in $`(p_{m3},n^{k1})`$ and so on down to $`(p_{mk},n^{k1}).`$ One gets a similar estimate as in the first case but with the indices lowered by one. Proceeding down to $`(p_{mM},n^{k1})`$ proves the claim.
To complete the proof of the proposition, we use another fact from the metric theory of continued fractions \[Kh, Theorem 30\]: For almost any $`\alpha \mathrm{I}\mathrm{R}`$, there exists $`C(\alpha )>0`$ such that $`a_m(\alpha )C(\alpha )m^{1+ϵ}`$. By Proposition (2.2.2), the relevant values of $`m`$ are of order $`\mathrm{log}(n).`$ Therefore, for the $`p_m,p_{m1},\mathrm{},p_{mM}`$ under consideration we have $`a_{mj}<<\mathrm{log}n.`$ Since $`E_{\mathrm{}}`$ is a polynomial in the $`a_{mj}`$’s of degree $`\mathrm{}`$, we have
$$E_{\mathrm{}}(a_{mj},a_{mj1},\mathrm{},a_{mj\mathrm{}+1})<<(\mathrm{log}n)^{\mathrm{}}.$$
Therefore
$$\mathrm{\Pi }_{j=0}^M\mathrm{\Pi }_{\mathrm{}=1}^{Mj}E_{\mathrm{}\mathrm{}}(a_{mj},a_{mj1},\mathrm{},a_{mj\mathrm{}+1})<<(\mathrm{log}n)^{M^3}.$$
(17)
It follows that
$$\mathrm{\Pi }_{j=0}^M(p_{mj},n^{k1})C(\alpha )n^{k1}(\mathrm{log}n)^{M^3}.$$
(18)
Hence at least one factor must be $`C(\alpha )^{1/M}n^{\frac{k1}{M}}(\mathrm{log}n)^{M^2}.`$ The proposition follows from the fact that $`M`$ can be arbitrarily large.
We now complete the proof of the lemma and of our main result. We have proved the existence of $`(p_m,q_m)`$ with all the necessary properties and such that $`q_m[n^{rϵ},n^r],(p_m,n^{k1})<<n^ϵ.`$ It follows that
$$\frac{n(p_m,n^{k1})}{q_m}+\frac{q_m}{n(p_m,n^{k1})}<<n^{1+ϵr}+n^{r1}.$$
(19)
The terms balance when $`r=1/2`$ and give the power $`n^\mathrm{}ϵ.`$ It follows from (12) that $`\rho _2^n(f)<<n^{1\frac{2}{K}+ϵ}.`$
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