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# On motives for Deligne-Mumford stacks
## 1 Preliminaries on Chow rings of Deligne-Mumford stacks
We start with the first definition of the Chow rings of Deligne-Mumford stacks.
For every $`DM`$-stack $`F`$, we consider $`𝒦_m`$, the sheaf of the $`m`$-th $`K`$-groups on $`F_{et}`$ associated with the abelian presheaf
$$\begin{array}{cccc}𝒦_m:& F_{et}& & Ab\\ & U& & K_m(U)\end{array}$$
###### Definition 1.1
(\[G\]) The codimension $`m`$ rationnal Chow group of $`F`$ is defined to be
$$A^m(F):=H^m(F_{et},𝒦_m).$$
We will note $`A^{}(F):=_mA^m(F)`$.
As it is shown in \[G\], the theory $`A^{}(F)`$ is a good Chow cohomology theory. Without recalling all the properties w recall three of them which will be usefull for us.
* (Product) For every $`DM`$-stack $`A^{}(F)`$ has a structure of a graded commutative ring.
* (Functoriality) For every morphism of $`DM`$-stacks, $`f:FF^{}`$, there exist an inverse image
$$f^{}:A^{}(F^{})A^{}(F)$$
which is a morphism of graded rings.
There is also a direct image
$$f_{}:A^{}(F)A^{}(F^{})$$
which is a morphism of $``$-vector spaces. This morphism is moreover graded of degree $`DimF^{}DimF`$ if $`F`$ and $`F^{}`$ are pure dimensional.
* (Projection formula) For every morphism $`f:FF^{}`$ between two $`DM`$-stacks we have
$$f_{}(x.f^{}(y))=f_{}(x).y$$
for every $`xA^{}(F)`$ and $`yA^{}(F^{})`$.
In particular, if $`F`$ and $`F^{}`$ are connected and of the same dimension, we have
$$f_{}f^{}=\times m$$
where $`m`$ is the generic degree of $`f`$ (in the stack sense \[V\]).
* (Compatibility) For every variety $`X`$, $`A^{}(X)CH^{}(X)_{}`$ is the usual Chow ring of $`X`$.
In order to introduce the second definition $`A_\chi ^{}`$ recall the main result of \[T3\]. It will not be used but explains the definition of $`A_\chi ^{}`$.
We begin with a $`DM`$-stack $`F`$, and consider $`C_F^t`$, the classifying stack of cyclic subgroups of automorphisms in $`F`$. More precisely it is defined in the following way.
A $`S`$-group scheme $`GS`$ is called cyclic (and finite), if locally for the etale topology on $`S`$ it is isomorphic (as a $`S`$-group scheme) to $`Spec{\displaystyle \frac{𝒪_S[T]}{T^m1}}`$. In other words, $`G`$ is a multiplicative type $`S`$-group scheme whose sheaf of characters has cyclic geometric fibres.
The stack $`C_F^t`$ is now defined by the following.
* For any $`k`$-scheme $`X`$, the objects in $`C_F^t(X)`$ are pairs $`(s,c)`$, where $`s`$ is an object in $`F(X)`$, and $`c`$ is a sub-group scheme of the $`X`$-group scheme of automorphisms of $`s`$, $`Aut_X(s)X`$, such that $`c`$ is a cyclic $`X`$-group scheme.
* An isomorphism between two objects in $`C_F^t(X)`$, $`(s,c)`$ and $`(s^{},c^{})`$, is an isomorphism $`u:ss^{}`$ in $`F(X)`$, such that $`u^1.c^{}.u=c`$.
The map which sends $`(s,c)`$ to $`s`$ gives a morphism $`\pi _F:C_F^tF`$, which is oubviously representable. Futhermore, we have the following local description of $`\pi _F`$. Locally (for the etale topology) on its moduli space $`F`$ is given by a quotient of a smooth scheme $`X`$ by a finite group $`H`$. So, to obtain a local descrition of $`\pi _F`$ it is enough to consider the case where $`F=[X/H]`$.
Let $`c(H)`$ be a set of representative of conjugacy classes of cyclic subgroups of $`H`$ whose orders are prime to the caracteristic of $`k`$. For each $`cc(H)`$, let $`X^c`$ the closed sub-scheme of fixed points of $`c`$ in $`X`$, and $`N_c`$ the normalisor of $`c`$ in $`H`$. Note that $`N_c`$ acts on $`X^c`$ by restriction. Then we have a natural equivalence
$$C_F^t\underset{cc(H)}{}[X^c/N_c].$$
As every $`c`$ is cyclic of order invertible on $`X`$, it is a diagonalisable group scheme, and so as $`X`$ is smooth, $`X^c`$ is also smooth.
From this local description we deduce that the stack $`C_F^t`$ is smooth, and the natural morphism $`\pi _F:C_F^tF`$ is representable, finite and unramified. In particular $`C_F^t`$ is again smooth and proper.
On the stack $`C_F^t`$ lives the universal cyclic group stack, $`q:𝒞_F^tC_F^t`$. It classifies triplets $`(s,c,h)`$, where $`(s,c)`$ is an object of $`C_F^t(X)`$ and $`h`$ a section of $`c`$ over $`X`$. Thus, for any morphism $`(s,c):UC_F^t`$, the pull-back of $`𝒞_F^t`$ on $`U`$ is isomorphic to the cyclic $`U`$-group scheme $`cU`$.
Let $`\chi _F`$ be the sheaf of characters of $`𝒞_F^t`$ on $`C_F^t`$. It is defined by $`\chi _F:=\underset{¯}{Hom}_{Gp}(𝒞_F^t,𝔾_m)`$. More explicitely, its restriction on the small etale site of $`C_F^t`$ is given by
$$\begin{array}{cccc}\chi _F:& (C_F^t)_{et}& & Ab\\ & ((s,c)C_F^t(U))& & Hom_{Gp}(c,𝔾_{m,U})\end{array}$$
As $`𝒞_F^t`$ is a cyclic group stack, $`\chi `$ is a locally constant sheaf on $`(C_F^t)_{et}`$, locally isomorphic to a constant finite cyclic group sheaf.
Let us consider the sheaf of group-algebras associated to $`\chi `$, $`[\chi _F]`$. It is a locally constant sheaf of $``$-algebras on $`(C_F^t)_{et}`$, which is locally isomorphic to the constant sheaf with fibre $`{\displaystyle \frac{[T]}{T^m1}}`$. As $`{\displaystyle \frac{[T]}{T^m1}}`$ is a product of cyclotomic fields with only one of maximal degree, namely $`(\zeta _m)`$, the kernels of the local quotients $`{\displaystyle \frac{[T]}{T^m1}}(\zeta _m)`$ glue together to give a well defined ideal sheaf $`_F[\chi _F]`$. We then define
$$\mathrm{\Lambda }_F:=\frac{[\chi _F]}{_F}.$$
Note that this is a well defined sheaf of $``$-algebras on $`C_F^t`$, locally isomorphic to the constant sheaf associated with a cyclotomic field.
We can now state the main result of \[T3\]. For a sketch of proof the reader can consult \[T3\], or \[T1, $`3.15`$\] for a particular case.
###### Theorem 1.2
There exist a functorial ring isomorphism
$$\varphi _F:K_{}(F)H^{}((C_F^t)_{et},\underset{¯}{K}\mathrm{\Lambda }_F).$$
Remark: Here $`K_{}(F)`$ is the ring of $`K`$-theory of perfect complexes on $`F`$, and $`\underset{¯}{K}`$ is the presheaf of $`K`$-theory spectrum on $`(C_F^t)_{et}`$.
The theorem justifies the following definition.
###### Definition 1.3
For any $`DM`$-stack $`F`$, the codimension $`m`$ rationnal Chow group with coefficients in the characters of $`F`$ is defined by
$$A_\chi ^m(F):=H^m((C_F^t)_{et},𝒦_m\mathrm{\Lambda }_F).$$
We will note $`A_\chi ^{}(F):=_mA_\chi ^m(F)`$.
There is a natural decomposition $`C_F^tFC_{F,+}^t`$ coming from the section $`FC_F^t`$ maping an object $`s`$ to $`(s,\{e\})`$. As the sheaf $`\chi _F`$ restricts to the constant sheaf $``$ on $`F_{et}`$, this induces a group decomposition
$$A_\chi ^{}(F)A^{}(F)\times A_{\chi 1}^{}(F).$$
###### Proposition 1.4
1. (Product) For every $`DM`$-stack, there is structure of graded commutative $``$-algebra on $`A_\chi ^{}(F)`$. Furthermore, the decomposition $`A_\chi ^{}(F)A^{}(F)\times A_{\chi 1}^{}(F)`$ becomes a $``$-algebra decomposition.
2. (Functoriality) For every morphism of $`DM`$-stacks $`f:FF^{}`$, there is an inverse image
$$f^{}:A_\chi ^{}(F^{})A_\chi ^{}(F)$$
which makes $`A_\chi ^{}`$ into a functor $`𝒟^o(gradedalgebras)`$. Furthermore the decomposition $`A_\chi ^{}A^{}\times A_{\chi 1}^{}`$ is compatible with these inverse images.
There exist a direct image
$$f_{}:A_\chi ^{}(F)A_\chi ^{}(F^{})$$
which makes $`A_\chi ^{}`$ into a functor $`𝒟Vect`$.
3. (Projection formula) For every morphism of $`DM`$-stacks $`f:FF^{}`$, we have
$$f_{}(x.f^{}(y))=f_{}(x).y$$
for every $`xA_\chi ^{}(F)`$ and $`yA_\chi ^{}(F^{})`$.
4. (Compatibility) For every variety $`X`$, $`A_\chi ^{}(X)CH^{}(X)_{}`$ is the usual Chow ring of $`X`$.
Proof: $`(1)`$ The product in $`K`$-theory gives morphisms of sheaves on $`(C_F^t)_{et}`$
$$𝒦_p𝒦_m𝒦_{p+m},$$
defining a graded ring structure on $`𝒦_{}:=_m𝒦_m`$. By tensoring with the sheaf of algebras $`\mathrm{\Lambda }_F`$ we obtain a sheaf of graded $``$-algebras $`𝒦_{}\mathrm{\Lambda }_F`$. It is then a general fact that the cohomology
$$A_\chi ^{}(F)H^{}((C_F^t)_{et},𝒦_{}\mathrm{\Lambda }_F)$$
is naturally a graded $``$-algebra.
$`(2)`$ Every morphism between two $`DM`$-stacks $`f:FF^{}`$ induces a morphism $`Cf:C_F^tC_F^{}^t`$. It sends an object $`(s,c)C_F^t(X)`$ to $`(f(s),f(c))C_F^{}^t`$. Furthermore there is a morphism of sheaves of groups $`Cf^1(\chi _F^{})\chi _F`$ given by restrictions of characters, giving a morphism of sheaves of algebras
$$Res_f:Cf^1(\mathrm{\Lambda }_F^{})\mathrm{\Lambda }_F.$$
On the other hand we have inverse image in $`K`$-theory, which gives a morphism of sheaves on graded algebras
$$Cf^{}:Cf^1(𝒦_{})𝒦_{}.$$
By tensorisation this gives
$$Cf^1(𝒦_{}\mathrm{\Lambda }_F^{})𝒦_{}\mathrm{\Lambda }_F$$
which allows to define inverse images
$$f^{}:H^{}((C_F^{}^t)_{et},𝒦_{}\mathrm{\Lambda }_F^{})H^{}((C_F^t)_{et},Cf^1(𝒦_{}\mathrm{\Lambda }_F^{}))H^{}((C_F^t)_{et},𝒦_{}\mathrm{\Lambda }_F).$$
To define the direct images we use the induction morphism of characters
$$Ind_f:\chi _FCf^1(\chi _F^{}).$$
This induces a morphism of sheaves of $``$-vector spaces
$$Ind_f:\mathrm{\Lambda }_FCf^1(\mathrm{\Lambda }_F^{}).$$
For every $`m`$, we use the Gersten resolution of the sheaf $`𝒦_m`$ (\[G2, $`7`$\])
$$𝒦_m_m^m_m^{m1}\mathrm{}_m^0.$$
Let $`_{}^{}:=_m_m^{}`$. Thinking of $`_m^0`$ in cohomological degree $`0`$, we have
$$A_\chi ^{}(F)H^0((C_F^t)_{et},_{}^{}\mathrm{\Lambda }_F).$$
The direct image is a morphism of complexes of sheaves on $`(C_F^{}^t)_{et}`$ (\[G2, $`7`$\])
$$Cf_{}(_{}^{})_{}^{}.$$
Tensoring with $`\mathrm{\Lambda }_F^{}`$ gives
$$Cf_{}(_{}^{}Cf^1(\mathrm{\Lambda }_F^{}))Cf_{}(_{}^{})\mathrm{\Lambda }_F^{}_{}^{}\mathrm{\Lambda }_F^{}.$$
We then compose with $`Ind_f`$ and take the cohomology to obtain
$$f_{}:A_\chi ^{}(F)A_\chi ^{}(F^{}).$$
$`(3)`$ Using the two previous explicits definitions of $`f_{}`$ and $`f^{}`$ the proof is exactly the same as for the case of scheme (\[G2, $`7`$\]).
$`(4)`$ If $`X`$ is a variety, then $`C_X^tX`$ and $`\mathrm{\Lambda }_X`$, so the isomorphism $`A_\chi ^{}(X)CH^{}(X)_{}`$ is given by the Bloch’s formula (\[G2, $`7`$\])
$$CH^p(X)_{}H^p(X_{zar},𝒦_p)H^p(X_{et},𝒦_p).$$
$`\mathrm{}`$
Remark: The Riemann-Roch formula of \[T1, $`4.11`$\] extends to a formula with values in $`A_\chi ^{}`$. Indeed, by using the construction of chern classes in \[G2\] and the theorem 1.2 one can define a Chern character
$$Ch^\chi :K_0(F)A_\chi ^{}(F).$$
The Todd class $`Td(F)`$ defined in \[T1, $`4.8`$\] can also be defined as $`Td^\chi (F)A_\chi ^{}(F)`$ in a very similar manner.
To prove the Riemann-Roch formula for $`Ch^\chi ().Td^\chi `$, we first use the projection formula to do galois descent and reduce the problem to the case where $`k`$ is algebraically closed. We choose an embending $`\mu _{\mathrm{}}(k)^{}`$. Then the formula follows from \[T2, $`3.36`$\] and the fact that (see 3.6 for a proof of this)
$$A_\chi ^{}(F)(\mu _{\mathrm{}}(k))A_{rep}^{}(F).$$
It is also true that the Chern character
$$Ch^\chi :K_0(F)_{}A_\chi ^{}(F)$$
is a ring isomorphism.
## 2 First construction
The construction of the category of Chow motives for $`DM`$-stacks using the theory $`A^{}`$ was done in \[B-M, $`8`$\]. We will denote it by $`^{DM}`$, and call its objects the $`DMC`$-motives, as suggested in \[B-M\]. We start by recalling briefly its construction.
For $`F,F^{}𝒟`$, we define the vector space of correspondences of degree $`m`$ between $`F`$ and $`F^{}`$
$$S^m(F,F^{}):=\{xA^{}(F\times F^{})/(p_2)_{}(x)A^m(F^{})\}$$
where $`p_2:F\times F^{}F^{}`$ is the second projection.
We have the usual composition
$$:S^m(F,F^{})S^n(F^{},F^{\prime \prime })S^{p+m}(F,F^{\prime \prime })$$
given by the formula
$$xy:=(p_{13})_{}(p_{12}^{}(x).p_{23}^{}(y)),$$
where the $`p_{ij}`$ are the natural projections of $`F\times F^{}\times F^{\prime \prime }`$ on two of the three factors.
Objects of $`^{DM}`$ are triplets $`(F,p,m)`$, with $`F𝒟`$, $`p`$ an idempotent in the ring of correspondences $`S^0(F,F)`$, and $`m`$. The morphisms between $`(F,p,m)`$ and $`(F^{},q,n)`$ are given by
$$Hom_{^{DM}}((F,p,m),(F^{},q,n)):=qS^{nm}(F,F^{})pS^{nm}(F,F^{}).$$
Recall also that for any morphism in $`𝒟`$, $`f:FF^{}`$, we have its graph
$$\mathrm{\Gamma }_f=f\times Id:FF^{}\times F,$$
and so a well defined element
$$[f^{}]:=(\mathrm{\Gamma }_f)_{}(1)S^0(F^{},F).$$
We can also consider its transposed
$$[f_{}]:=[f^{}]^tS^{}(F,F^{}).$$
This allows us to define a functor
$$\begin{array}{cccc}h:& 𝒟^o& & ^{DM}\\ & F& & (F,Id,0)\\ & f& & [f^{}]\end{array}$$
As for the case of varieties, the category $`^{DM}`$ is $``$-linear karoubian category (\[B-M, $`8.1`$\]). In particular, this implies that if a morphism in $`^{DM}`$ possesses a left inverse then it is a direct factor.
It is also symetric monoidal for the tensor product defined by
$$(F,p,m)(F^{},q,n):=(F\times F^{},pq,n+m).$$
As usual we shall write $`𝕃^m=(Speck,Id,m)`$ for the $`m`$-th power of the Lefschetz motive. Note that for every $`DMC`$-motive $`M`$, we have $`M(F,p,0)𝕃^m`$. As $`(F,p,0)`$ is a direct factor in $`h(F)`$, this shows that $`M`$ is a direct factor in some $`h(F)𝕃^m`$.
Finally, there is a natural fully faithfull tensorial functor
$$^{DM}$$
from the usual category of Chow motives of varieties to the category of $`DMC`$-motives. This functor fits into a commutative diagramm
The following theorem is a positive answer to the question \[B-M, $`8.2`$\].
###### Theorem 2.1
The previous functor
$$^{DM}$$
is an equivalence of $``$-tensorial categories.
Proof: By noticing that the essential image is closed by direct factors (because any direct factor of $`(X,p,m)`$ in $`^{DM}`$ is of the form $`(X,pq,m)`$), we only have to check that for each connected $`F𝒟`$, $`h(F)`$ is a direct factor of some $`h(X)`$ for $`X𝒱𝒜`$.
Let $`F𝒟`$, and by \[L-M, $`16.6`$\] choose an integral scheme $`X`$ and a finite and surjective morphism $`XF`$. Using \[J\] we can find $`YX`$ which is generically finite, with $`Y`$ a variety. We know consider the composed morphism $`f:YF`$, as well as
$$[f_{}]:h(Y)h(F)$$
$$[f^{}]:h(F)h(Y).$$
The indentity principle \[B-M, $`8.2`$\] and the projection formula implies that $`\frac{1}{m}.[f_{}]`$ is a left inverse to $`[f^{}]`$. This implies that $`[f^{}]`$ is a direct factor. More explicitely we have $`h(F)(X,\frac{1}{m}[f_{}][f^{}],0)`$. $`\mathrm{}`$
Inverting the equivalence $`^{DM}`$ gives a functor
$$h:𝒟^o.$$
As an inverse of a monoidal functor has a natural monoidal structure, $`h`$ is naturally a monoidal functor. We obtain this way natural isomorphisms
$$h(F\times F^{})h(F)h(F^{})$$
wich are associatives, commutatives and unitaries. In particular, the diagonal of a $`DM`$-stack $`F`$ gives a commutative algebra structure on the motive $`h(F)`$. This can be used for example to show that every good cohomology theory for varieties extends to $`DM`$-stacks.
## 3 Second construction
###### Definition 3.1
For two $`DM`$-stacks $`F`$ and $`F^{}`$, we define the vector space of $`\chi `$-correspondences of degree $`m`$ between $`F`$ ans $`F^{}`$ by
$$S_\chi ^m(F,F^{}):=\{xA_\chi ^{}(F\times F^{})/(p_2)_{}(x)A_\chi ^m(F^{})\}.$$
As in the previous case, we have a composition
$$:S_\chi ^m(F,F^{})S_\chi ^n(F^{},F^{\prime \prime })S_\chi ^{p+m}(F,F^{\prime \prime })$$
given by the formula
$$xy:=(p_{13})_{}(p_{12}^{}(x).p_{23}^{}(y)),$$
where the $`p_{ij}`$ are the natural projections of $`F\times F^{}\times F^{\prime \prime }`$ on two of the three factors.
###### Definition 3.2
We define the category of $`DMC_\chi `$-motives, $`_\chi ^{DM}`$ as follows.
* Objects of $`_\chi ^{DM}`$ are triplets $`(F,p,m)`$, where $`F`$ is a $`DM`$-stack, $`pS_\chi ^0(F,F)`$ an idempotent, and $`m`$.
* The set of morphisms between $`(F,p,m)`$ and $`(F^{},q,n)`$ is defined by
$$Hom_{_\chi ^{DM}}((F,p,m),(F^{},q,n)):=qS_\chi ^{nm}(F,F^{})pS_\chi ^{nm}(F,F^{}).$$
* The composition of morphisms in $`_\chi ^{DM}`$ is given by composition of $`\chi `$-correspondences.
For any morphism $`f:FF^{}`$ between two $`DM`$-stacks, we define as usual
$$[f^{}]:=(\mathrm{\Gamma }_f)_{}(1)S_\chi ^0(F^{},F),$$
as well as its transposed
$$[f_{}]:=[f^{}]^tS_\chi ^{}(F,F^{}).$$
Using this we define a natural functor
$$\begin{array}{cccc}h_\chi :& 𝒟^o& & _\chi ^{DM}\\ & F& & (F,Id,0)\\ & f& & [f^{}]\end{array}$$
The same arguments as for motives of varieties show that $`_\chi ^{DM}`$ is a $``$-tensorial karoubian category. There is alos a tensor product, given as usual by $`(F,p,m)(F^{},q,n):=(F\times F^{},pq,m+n)`$.
Note that the compatibility propety of 1.4 implies that there is a natural fully faithfull functor
$$_\chi ^{DM}.$$
###### Definition 3.3
For any $`DMC_\chi `$-motive $`M`$, we define its $`m`$-th Chow group by
$$A_\chi ^m(M):=Hom_{_\chi ^{DM}}(𝕃^m,M).$$
We will note $`A_\chi ^{}(M):=_mA_\chi ^m(M)`$.
Remark: Using the Chern character we have
$$Ch^\chi :K_0(F)_{}A_\chi ^{}(F).$$
Finally, the indentity principle says that the functor $`_\chi ^{DM}Hom(𝒟^o,Ab)`$, which sends $`M`$ to the functor $`FA_\chi ^{}(Mh(F))`$ is fully faithfull (it follows immediately from the Yoneda lemma and the fact that every $`DMC_\chi `$-motive is a direct factor of a $`h(F)𝕃^m`$).
###### Theorem 3.4
The natural functor
$$_\chi ^{DM}$$
is an equivalence of $``$-tensorial categories.
Proof: As for 2.1 it is enough to show that every $`h_\chi (F)`$ is a direct factor in some $`h_\chi (X)`$.
For the next lemma recall that for any stack $`F`$ we can define its inertia stack $`I_F`$ (\[V\]), whose objects are pairs $`(s,h)`$, with $`s`$ an object of $`F`$ and $`h`$ and automorphism of $`s`$. It can for example be defined by the formula
$$I_F:=F\times _{F\times F}F.$$
###### Lemma 3.5
For any Deligne-mumford stack proper over $`Speck`$ (non necerally smooth), there exist varieties $`Y_i`$ and finite groups $`H_i`$ together with a proper representable morphism
$$F_0:=\underset{i}{}Y_i\times BH_iF$$
such that the induced morphism
$$I_{F_0}I_F$$
is generically finite and surjective.
Proof: By \[L-M, $`16.6`$\] we can choose a finite and surjective morphism $`XF`$, with $`X`$ a normal scheme. Let $`FM`$ be the moduli space of $`F`$, and consider $`F_X`$, the normalization of the fibre product $`F\times _MX`$. By definition, the stack $`F_X`$ is normal and the projection to its moduli space $`F_XX`$ possesses a section. It follows from \[V, $`2.7`$\] that $`F_X`$ is a neutral gerb. By choosing a finite and etale morphism $`YX`$ and defining $`F^{}:=F_X\times _XY`$, we find a trivial gerb $`F^{}Y\times BH`$, together with a morphism $`F^{}F`$. By construction this morphism is generically obtained by a pull back of a etale morphism on $`M`$. This implies that there exists a dense open sub-stack $`U`$ of $`F`$, such that $`I_UI_F`$ is in the image of $`I_F^{}I_F`$. Proceding by noetherian induction we find reduced schemes $`X_i`$, and finite groups $`H_i`$, with a morphism $`F^{}:=_iX_i\times BH_iF`$ such that $`I_F^{}I_F`$ is finite and surjective.
We now apply \[J\] to each $`X_i`$ and choose generically finite morphism $`Y_iX_i`$, with $`Y_i`$ a variety. Let $`F_0:=_iY_i\times BH_i`$. As $`I_{F_0}I_F\times _FF_0`$, the induced morphism
$$I_{F_0}I_F$$
is still surjective and generically finite. $`\mathrm{}`$
###### Lemma 3.6
Suppose that $`k`$ is algebraically closed, and choose an embedding $`\mu _{\mathrm{}}(k)^{}`$. Let $`F`$ be a $`DM`$-stack, and denote by $`I_F^t`$ the open and closed sub-stack of $`I_F`$ whose objects are pairs $`(s,h)`$, such that the order of $`h`$ is prime to the characteristic of $`k`$. Then there exist an $`(\mu _{\mathrm{}}(k))`$-algebra isomorphism
$$A_\chi ^{}(F)(\mu _{\mathrm{}}(k))A^{}(I_F^t)(\mu _{\mathrm{}}(k)).$$
Furthermore this isomorphism is compatible with inverse and direct images.
Proof: Let $`u:I_F^tC_F^t`$ the morphism which sends an object $`(s,h)`$ to $`(s,<h>)`$, where $`<h>`$ is the subgroup generated by $`h`$ in $`Aut(s)`$. This is a representable finite et etale morphism. It is easy to see that there is an isomorphism of sheaves of graded $`(\mu _{\mathrm{}}(k))`$-algebras on $`(C_F^t)_{et}`$
$$u_{}(𝒦_{}(\mu _{\mathrm{}}(k)))𝒦_{}\mathrm{\Lambda }_F(\mu _{\mathrm{}}(k)).$$
This induces the required isomorphism
$$A^{}(I_F^t)(\mu _{\mathrm{}}(k))H^{}((I_F^t)_{et},𝒦_{}(\mu _{\mathrm{}}(k)))H^{}((C_F^t)_{et},u_{}(𝒦_{}(\mu _{\mathrm{}}(k))))$$
$$H^{}((C_F^t)_{et},𝒦_{}\mathrm{\Lambda }_F(\mu _{\mathrm{}}(k)))A_\chi ^{}(F)(\mu _{\mathrm{}}(k)).$$
The compatibility with inverse and direct images is clear by definitions. $`\mathrm{}`$
Let $`g:F_0:=_iY_i\times BH_iF`$ be a morphism as in 3.5.
###### Lemma 3.7
The element $`\beta :=g_{}(1)A_\chi ^{}(F)`$ is invertible.
Proof: We first use the projection formula 1.4 to show that for any finite extension $`k^{}/k`$, we have a natural isomorphism of algebras
$$A^{}(F)_\chi A_\chi ^{}(F\times _{Speck}Speck^{})^{Gal(k^{}/k)}.$$
This allows to assume that $`k`$ is algebraically closed.
Applying the lemma 3.6, it is enough to show that $`Ig_{}(1)A^{}(I_F^t)`$ is invertible. But, as $`Ig`$ is generically finite and surjective this is obvious. $`\mathrm{}`$
Consider $`\mathrm{\Delta }_{}(\beta )A_\chi ^{}(F\times F)=S_\chi ^{}(F,F)`$, where $`\mathrm{\Delta }:FF\times F`$ is the diagonal. By the previous lemma $`\beta `$ is invertible in the graded ring $`S^{}(F,F)`$. Let $`\alpha :=[g_{}]\beta ^1S^{}(F^{},F)`$. Then we have $`[g^{}]\alpha =1`$. This shows that the $`0`$-th component of $`\alpha `$ is a left inverse of $`[g^{}]`$. As the category $`_\chi ^{DM}`$ is karoubian, this implies that $`[g^{}]`$ is a direct factor, and so that $`h_\chi (F)`$ is a direct factor in $`h_\chi (F^{})`$.
As $`h_\chi (F^{})_ih_\chi (X_i)h_\chi (BH_i)`$ it remains to show that for any finite group $`H`$, $`h_\chi (BH)`$ is isomorphic to some power of the trivial motive $`h_\chi (Speck)`$.
Let $`Ch^\chi :K_0(BH)A_\chi ^0(BH)`$ the Chern character, $`\rho _1,\mathrm{},\rho _r`$ a set of representatives of irreducibles representations of $`H`$ over $`k`$, and $`\alpha _i:=Ch^\chi (\alpha _i)`$. These elements define morphisms of $`DMC_\chi `$-motives $`\alpha _i:h_\chi (Speck)h_\chi (BH)`$ Let us consider the sum
$$\underset{i}{}\alpha _i:h_\chi (Speck)^rh_\chi (BH),$$
and prove that it is an isomorphism. By the identity principle, we have to show that for every $`DM`$-stack $`F`$, the induced morphism
$$\underset{i}{}\alpha _i:(A_\chi ^{}(F))^rA_\chi ^{}(F\times BH)$$
is an isomorphism. But as $`Ch^\chi `$ is an isomorphism, the previous morphism is isomorphic to the Kunneth morphism
$$A_\chi ^{}(F)A_\chi ^{}(BH)A_\chi ^{}(F\times BH),$$
and so the theorem follows from the following lemma.
###### Lemma 3.8
For every $`DM`$-stack $`F`$ and every finite group $`H`$, the Kunneth morphism
$$A_\chi ^{}(F)A_\chi ^{}(BH)A_\chi ^{}(F\times BH)$$
is an isomorphism.
Proof: Using galois descent we can suppose that $`k`$ is algebraically closed. Then, using the lemma 3.6 we reduce the problem to show that the Kunneth morphism
$$A^{}(I_F^t)A^{}(I_{BH}^t)A^{}(I_{F\times BH}^t)$$
is an isomorphism.
Let $`A`$ be a set of representative of conjugacy classes of elements in $`H`$ with order prime to the characteristic of $`k`$. We have
$$I_{BH}^t\underset{hA}{}BZ_hI_{F\times BH}^tI_F^t\times I_{BH}^t,$$
where $`Z_h`$ is the centralisator of $`h`$ in $`H`$. So we only need to prove that the Kunneth morphism
$$A^{}(I_F^t)A^{}(BZ_h)A^{}(I_F^t\times BZ_h)$$
is an isomorphism. But this morphism fits into a commutative diagram
where $`v:I_F^t\times BZ_hI_F^t`$ is the first projection. Now, as $`A^{}(BZ_h)`$, the vertical morphism is an isomorphism. On the other hand, $`v`$ has a natural section $`u:I_F^tI_F^t\times BZ_h`$ and the projection formula shows that $`u^{}`$ is an isomorphism, which implies that $`v^{}`$ is an isomorphism. $`\mathrm{}`$
Inverting the equivalence $`_\chi ^{DM}`$ gives a functor
$$h_\chi :𝒟^o.$$
As for the case of the first construction this functor has a natural monoidal structure. This implies that for any $`DM`$-stack $`F`$, the motive $`h_\chi (F)`$ has a natural structure of a commutative algebra in $``$. In particular any good cohomology theory for varieties extends trough $`h_\chi `$ to a new theory for stacks.
###### Proposition 3.9
The functor $`h`$ is a direct factor of the functor $`h_\chi `$.
Proof: This follows immediately from the natural decomposition $`A_\chi ^{}A^{}\times A_{\chi 1}^{}`$. $`\mathrm{}`$
Remark: For any complex variety $`V`$ and $`\beta H_2(V,)`$, we can define the Gromov-Witten correspondence (\[B\])
$$I_{g,n}(V,\beta )S^{}(V^n,_{g,n}),$$
which is a morphism of graded $`DMC`$-motives (\[B-M, $`8`$\]). It seems natural to ask if this correspondence extends in a natural way to
$$I_{g,n}^\chi (V,\beta )S_\chi ^{}(V^n,_{g,n})$$
(i.e. as a morphism of graded $`DMC_\chi `$-motives). This question is of course linked to the question of constructing an extended virtual fundamental class $`_{g,n}^\chi (V,\beta )A_{}^\chi (_{g,n}(V,\beta ))`$.
## 4 Examples
We have seen that the two Chow coholomogy theories $`A^{}`$ and $`A_\chi ^{}`$ give natural functors
$$h,h_\chi :𝒟^o,$$
such that $`h`$ is a direct factor of $`h_\chi `$. In this last chapter we will give some examples of motives associated to certain stacks, and see some expicit relations between $`h_\chi `$ and $`h`$.
The proofs of the following three facts are left to the reader (they all follow from the indentity principle and the explicit description of the stacks $`C_F^t`$ and the sheaves $`\mathrm{\Lambda }_F`$). For the sake of simplicity we will suppose that $`k`$ contains the roots of unity.
If a finite group $`H`$ acts on a motive $`M`$ we will denote by $`M^H`$ the direct factor of $`M`$ corresponding to the projector $`{\displaystyle \frac{1}{m}}.{\displaystyle \underset{hH}{}}h`$.
1. Quotients stacks
Let $`H`$ be a finite group acting on a variety $`X`$. Let $`c(H)`$ be a set of representatives of conjugacy classes of cyclic sub-groups of $`H`$, whose orders are prime to the characteristic of $`k`$. For every $`cc(H)`$ let $`X^c`$ be the sub-variety of $`X`$ of fixed points of $`c`$, and $`N_c`$ the normalisator of $`c`$ in $`H`$. For any $`cc(H)`$, let $`s(c)`$ be the set of injectives characters $`ck^{}`$.
Then the group $`N_c`$ acts on $`X^c`$ and $`s(c)`$, and so on the product $`X^c\times s(c)`$, and there is an isomorphism
$$h_\chi ([X/H])\underset{cc(H)}{}h(X^c\times s(c))^{N_c}.$$
Furthermore, the motive $`h([X/H])`$ corresponds to the component of the trivial sub-group
$$h([X/H])h(X)^H.$$
For example if $`X=Speck`$ we obtain
$$h_\chi (BH)h(Speck)^r,$$
where $`r`$ is the number of irreducible representations of $`H`$ in $`k`$-vetor spaces. But notice that this isomorphism does not preserve the product structures (given on any $`h_\chi (F)`$ by the diagonal morphism). Indeed, if $`\rho _1,\mathrm{},\rho _r`$ are the irreducibles representations of $`H`$ over $`k`$, then we have the mutiplication rules
$$\rho _i\rho _j\underset{k}{}\rho _k^{n_k^{i,j}}.$$
Then, the product on $`h_\chi (BH)`$ corresponds on $`h(Speck)^r`$ to the morphism
$$h(Speck)^rh(Speck)^rh(Speck)^{r^2}h(Speck)^r$$
given by the $`r^2`$ by $`r`$ matrix $`(n_k^{i,j})_{i,j,k}`$.
2. Gerbs
Let $`F`$ be a connected $`DM`$-stack which is a gerb (i.e. the morphism $`C_F^tF`$ is etale), and $`FX`$ its projection to its moduli space. Recall that locally for the etale topology of $`X`$, $`F`$ is equivalent to $`X\times BH`$, for $`H`$ a finite group. This defines a locally constant sheaf of groups up to inner automorphisms on $`X_{et}`$, which is classified by its monodromy
$$\pi _1^{et}(X)Out(H).$$
Let $`cycl(H)`$ be the set of cyclic sub-groups of $`H`$ of order prime to the characteristic of $`k`$, and for any $`ccycl(H)`$, $`s(c)`$ the set the of injectives characters $`ck^{}`$. The group $`H`$ acts by conjugaison on $`{\displaystyle \underset{ccycl(H)}{}}s(c)`$, and let $`R(H):=({\displaystyle \underset{ccycl(H)}{}}s(c))/H`$ be the quotient. The group $`Aut(H)`$ acts naturally on $`R(H)`$ and any inner automorphisms of $`H`$ acts trivially, so we deduce a morphism
$$\pi _1^{et}(X)Aut(R(H)),$$
which it turns corresponds to a finite etale covering $`YX`$.
There is then an isomorphism
$$h_\chi (F)h(Y).$$
Note that the trivial subgroup with the trivial character induces a section $`XY`$, which gives a decomposition
$$h(Y)h(X)h(Y)^1.$$
Furthermore, $`h(F)`$ corresponds to the factor $`h(X)`$.
3. $`1`$-Dimensional complex orbifolds
Suppose that $`k=`$ is the field of complex number, and that $`F`$ is a $`1`$-dimensional $`DM`$-stack, which is generically a variety (i.e. $`C_F^tF`$ is birationnal). Let $`C`$ be the moduli space of $`F`$, which is a smooth projective curve, and note $`x_1,\mathrm{},x_r`$ the points of $`C`$ where $`F`$ is not a scheme. Locally for the analytic topology around each $`x_i`$, $`F`$ is a quotient stack of a disc by a cyclic group $`/n_i`$. There is then an isomorphism
$$h_\chi (F)h(C)\underset{i}{}h(Spec)^{n_i1},$$
where $`h(F)`$ corresponds to the factor $`h(C)`$.
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# Untitled Document
APPLICATION OF LIE TRANSFORMATION GROUP METHODS TO
CLASSICAL THEORIES OF PLATES AND RODS
V. VASSILEV<sup>1</sup><sup>1</sup>1E-mail: vassil@bgcict.acad.bg and P. DJONDJOROV<sup>2</sup><sup>2</sup>2E-mail: padjon@bgcict.acad.bg
Institute of Mechanics, Bulgarian Academy of Sciences
Acad. G. Bontchev St., Block 4, 1113 Sofia, Bulgaria
> Abstract—In the present paper, a class of partial differential equations related to various plate and rod problems is studied by Lie transformation group methods. A system of equations determining the generators of the admitted point Lie groups (symmetries) is derived. A general statement of the associated group-classification problem is given. A simple intrinsic relation is deduced allowing to recognize easily the variational symmetries among the ”ordinary” symmetries of a self-adjoint equation of the class examined. Explicit formulae for the conserved currents of the corresponding (via Noether’s theorem) conservation laws are suggested. Solutions of group-classification problems are presented for subclasses of equations of the foregoing type governing stability and vibration of plates, rods and fluid conveying pipes resting on variable elastic foundations and compressed by axial forces. The obtained group-classification results are used to derive conservation laws and group-invariant solutions readily applicable in plate statics or rod dynamics.
1. INTRODUCTION
A wide variety of classical theories of plates and rods<sup>3</sup><sup>3</sup>3In this work, following Antman (1984) we use ”rod” as a generic name for ”arch”, ”bar”, ”beam”, ”ring”, ”column”, ”tube”, ”pipe”, etc. We employ ”rod” in the intuitive sense of a slender solid body. rest on linear fourth-order partial differential equations in one dependent and two independent variables. Some of them, such as the Poisson-Kirchhoff type theories for small bending of plates, are developed within the framework of the linear elastostatics to determine the state of equilibrium of thin elastic plates in terms of the transversal displacement of the plate middle plane, the derived governing equations providing the background for solving problems concerning stability and vibration of such structural elements. Others (among them – the dynamic theory of Bernoulli-Euler beams, for instance) are deduced on the ground of the linear elastodynamics to describe the dynamic behaviour of rods in terms of the transversal displacement of the rod axis.
The aim of the present paper is to study the invariance properties (symmetries) of the equations of the foregoing type with respect to local Lie groups of point transformations of the involved independent and dependent variables. The work is motivated both by the aforesaid wide applicability of the equations in question in structural mechanics, and by the remarkable efficiency demonstrated by the symmetry methods, especially when applied to differential equations arising in physics and engineering.
Actually, once the invariance properties of a given differential equation are established, several important applications of its symmetries arise. First, it is possible to distinguish classes of solutions to this equation invariant under the transformations of symmetry groups admitted. The determination of such a group-invariant solution assumes solving a reduced equation involving less independent variables than the original one. Typical examples of group-invariant solutions are axisymmetric solutions, self-similar solutions, travelling waves, etc., which have proved to be quite useful in many branches of physics and engineering. For a self-adjoint differential equation another substantial application of its symmetries is available. As is well known, the self-adjoint equations are the Euler-Lagrange equations of a certain action functional. If a one-parameter symmetry group of such an equation turned out to be its variational symmetry as well, that is a symmetry of the associated action functional, then Noether’s theorem guarantees the existence of a conservation law for the smooth solutions of this equation. Needless to recall or discuss here the fundamental role that the conserved quantities and conservation laws (or the corresponding integral relations, i.e. the balance laws) have played in natural sciences, but it is worthy to point out that the available conservation laws (balance laws) should not be overlooked (as it is often done) in the examination of discontinuous solutions (acceleration waves, shock waves, etc.) or in the numerical analysis (when constructing finite difference schemes or verifying numerical results, for instance) of any system of differential equations of physical interest. It should be remarked also that the path-independent integrals (such as the well known $`J`$-, $`L`$\- and $`M`$-integrals) related to the conservation laws are basic tools in fracture analysis of solids and structures.
The aforementioned and many other applications of the symmetries of differential equations and variational problems as well as the foundations of the Lie transformation group methods, including the basic notions, statements and techniques, can be found in the books by Ovsiannikov (1982), Ibragimov (1985), Bluman and Kumei (1989) and Olver (1993) (see also the references given in these books). In the present paper, however, as fare as the application of the symmetry groups of the equations studied is concerned, our attention is restricted to the constructing of group-invariant solutions and conservation laws. Of course, the first task is to find these symmetry groups, and as here we do not deal with a single differential equation but with a class of differential equations, this means to solve a group-classification problem.
The layout of the paper is as follows. A detailed description of the differential equations to be studied as well as the variational statement for the self-adjoint equations among them are given in Section 2. Several examples of mechanical systems governed by such equations complete this Section. In Section 3, a system of equations determining the generators of the symmetry groups admitted by the equations of the class considered is derived, the invariance properties inherent to the whole class because of its linearity and homogeneity being taken into account, and then the general statement of the associated group-classification problem is given. After that, the variational symmetries of the self-adjoint equations of the examined class are investigated. A simple intrinsic relation allowing to recognize easily the variational symmetries among the ”ordinary” point Lie symmetries of such an equation is deduced, and explicit formulae for the conserved currents of the conservation laws corresponding to the variational symmetries via Noether’s theorem are suggested. Group-classification results, conservation laws and group-invariant solutions are presented in Section 4 for differential equations governing stability and vibration of plates of Poisson-Kirchhoff type. Similar results are displayed in Section 5 for equations governing vibration of rods on a variable elastic foundation and dynamic stability of fluid conveying pipes or rods compressed by axial forces. In the reminder of this Section, conservation laws for rods derived in the present contribution are compared to other ones reported in the literature.
2.BASIC EQUATIONS
Consider a fourth-order homogeneous linear partial differential equation
$$A^{\alpha \beta \gamma \delta }(x)w_{\alpha \beta \gamma \delta }+A^{\alpha \beta \gamma }(x)w_{\alpha \beta \gamma }+A^{\alpha \beta }(x)w_{\alpha \beta }+A^\alpha (x)w_\alpha +A(x)w=0,$$
(1)
in two independent variables $`x=(x^1,x^2)`$ and one dependent variable $`w(x)`$. Here and throughout: Greek indices have the range 1, 2, and the usual summation convention over a repeated index (one subscript and one superscript) is employed; $`w_{\alpha _1\alpha _2\mathrm{}\alpha _k}`$ $`(k=1,2,\mathrm{})`$ denote the $`k`$-th order partial derivatives of the dependent variable, i.e.
$$w_{\alpha _1\alpha _2\mathrm{}\alpha _k}=\frac{^kw}{x^{\alpha _1}x^{\alpha _2}\mathrm{}x^{\alpha _k}}(k=1,2,\mathrm{}).$$
Further, a similar notation will be used for the partial derivatives of any other function of the variables $`x^1,x^2`$ but, in this case, the indices indicating the differentiation will be preceded by a coma, e.g.
$$A_{,\alpha _1\alpha _2\mathrm{}\alpha _k}^{\alpha \beta \gamma \delta }=\frac{^kA^{\alpha \beta \gamma \delta }}{x^{\alpha _1}x^{\alpha _2}\mathrm{}x^{\alpha _k}}(k=1,2,\mathrm{}).$$
The coefficients of equation (1) are supposed to be smooth functions possessing as many derivatives as may be required on a certain domain of interest, and to be symmetric under any permutation of their indices, i.e.
$$A^{\alpha \beta \gamma \delta }=A^{\beta \gamma \delta \alpha }=A^{\gamma \delta \alpha \beta }=A^{\delta \alpha \beta \gamma },A^{\alpha \beta \gamma }=A^{\beta \gamma \alpha }=A^{\gamma \alpha \beta },A^{\alpha \beta }=A^{\beta \alpha }.$$
Using the total derivative operators
$$D_\alpha =\frac{}{x^\alpha }+w_\alpha \frac{}{w}+w_{\alpha \mu }\frac{}{w_\mu }+w_{\alpha \mu \nu }\frac{}{w_{\mu \nu }}+w_{\alpha \mu \nu \sigma }\frac{}{w_{\mu \nu \sigma }}+\mathrm{},$$
the equation (1) may be written in the form
$$𝒟[w]=0,$$
(2)
where $`𝒟`$ is the linear differential operator given by the expression
$$𝒟=A^{\alpha \beta \gamma \delta }D_\alpha D_\beta D_\gamma D_\delta +A^{\alpha \beta \gamma }D_\alpha D_\beta D_\gamma +A^{\alpha \beta }D_\alpha D_\beta +A^\alpha D_\alpha +A.$$
(3)
An equation of form (2) is the Euler-Lagrange equation associated with a certain variational problem involving only one dependent variable if and only if the differential operator $`𝒟`$ is self-adjoint, that is
$$𝒟=𝒟^{},$$
(4)
where $`𝒟^{}`$ is the (formal) adjoint operator of $`𝒟`$ (cf. Olver, 1993). The explicit form of $`𝒟^{}`$ is
$$𝒟^{}=D_\alpha D_\beta D_\gamma D_\delta A^{\alpha \beta \gamma \delta }D_\alpha D_\beta D_\gamma A^{\alpha \beta \gamma }+D_\alpha D_\beta A^{\alpha \beta }D_\alpha A^\alpha +A.$$
(5)
In such a case, (2) can be associated with the variational problem for the functional
$$A[w]=\frac{1}{2}w𝒟[w]𝑑x^1𝑑x^2,$$
since the application of the Euler operator
$$𝖤=\frac{}{w}D_\mu \frac{}{w_\mu }+D_\mu D_\nu \frac{}{w_{\mu \nu }}D_\mu D_\nu D_\sigma \frac{}{w_{\mu \nu \sigma }}+D_\mu D_\nu D_\sigma D_\tau \frac{}{w_{\mu \nu \sigma \tau }}\mathrm{}$$
on the Lagrangian density
$$L=\frac{1}{2}w𝒟[w],$$
(6)
yields
$$𝒟[w]=𝖤(L),$$
(7)
due to (4) and (5).
Let us give several examples of plate and rod structures whose governing equations are self-adjoint and belong to the class specified above. Henceforward, when regarding plates the variables $`x^1,x^2`$ will represent the coordinates of the plate middle plane. As for the rod problems, $`x^1`$ will be associated with the spatial variable along the rod axis, and $`x^2`$ – with the time. In both cases, the dependent variable $`w`$ will represent the transversal displacement field.
Example 1. Small bending of plates resting on elastic foundations. Consider a thin elastic plate of variable bending rigidity $`D(x)`$ resting on an elastic foundation of Winkler type with variable modulus $`k(x)`$ and subjected to an edge loading leading to the appearance of nonuniform membrane stresses $`N^{\alpha \beta }(x)`$. In this physical situation, the equation governing the small bending of the plate assumes the following form
$$\mathrm{\Delta }[D\mathrm{\Delta }w][(1\nu )\epsilon ^{\alpha \mu }\epsilon ^{\beta \nu }D_{,\alpha \beta }N^{\mu \nu }]w_{\mu \nu }+kw=0,$$
(8)
the membrane stress tensor $`N^{\alpha \beta }`$ being symmetric, $`N^{\alpha \beta }=N^{\beta \alpha }`$, and divergence free, i.e. $`N_{,\mu }^{\alpha \mu }=0`$. Here: $`\nu `$ is Poisson’s ratio; $`\mathrm{\Delta }`$ is the Laplace operator, that is $`\mathrm{\Delta }\delta ^{\alpha \beta }^2/x^\alpha x^\beta `$, where $`\delta ^{\alpha \beta }`$ is the Kronecker delta symbol ($`\delta ^{11}=\delta ^{22}=1,\delta ^{12}=\delta ^{21}=0`$) and $`\epsilon ^{\alpha \beta }`$ is the alternating symbol ($`\epsilon ^{11}=\epsilon ^{22}=0,\epsilon ^{12}=\epsilon ^{21}=1`$).
Example 2. Elastodynamics of Bernoulli-Euler beams. Consider a nonhomogeneous Bernoulli-Euler beam with bending rigidity $`B(x^1)`$ and inertia term $`H(x^1)`$. The differential equation governing the small vibration of such a beam is (Chien at al., 1993):
$$Bw_{1111}+2B_{,1}w_{111}+B_{,11}w_{11}+Hw_{22}=0.$$
Example 3. Elastic beams resting on elastic foundations. Consider an elastic beam of constant bending rigidity $`K`$ and constant mass density $`m`$ (mass of the beam per unit length), resting on an elastic foundation with variable modulus $`k(x)`$. Suppose it is subjected to a constant follower force $`p`$. Then, according to the Bernoulli-Euler theory, the differential equation for small transverse vibration of the beam is (see Smith and Herrmann, 1972):
$$Kw_{1111}+pw_{11}+k(x)w+mw_{22}=0.$$
(9)
Example 4. Pipes conveying fluid. Consider an elastic circular-cylindrical pipe of uniform outer radius of the pipe cross section, which is supposed to be small in comparison with certain characteristic pipe length (for a simply supported pipe say the length of the span). Let the pipe conveys inviscid incompressible fluid with a flow velocity $`U=const`$. Then, the equation of motion is (see Gregory and Paidoussis, 1966):
$$EJw_{1111}+MU^2w_{11}+2MUw_{12}+(m+M)w_{22}=0,$$
(10)
where $`E`$ is Young’s modulus of the pipe material, $`J`$ is the (axial) moment of inertia of the pipe cross section, $`m`$ is the mass (constant) of the pipe per unit length, $`M`$ is the mass (also constant) of the fluid per unit length.
Combining and generalizing equations (9) and (10) presented in Examples 3 and 4, in Section 5 we will pay particular attention to the differential equations of the form
$$\gamma w_{1111}+\chi ^{\alpha \beta }w_{\alpha \beta }+\kappa (x)w=0,$$
(11)
where $`\gamma 0`$ and $`\chi ^{\alpha \beta }`$ are real constants, while $`\kappa (x)`$ is a smooth function.
3. SYMMETRIES AND CONSERVATION LAWS
Consider a local one-parameter Lie group of point transformations acting on some open subset $`\mathrm{\Omega }`$ of the space $`𝐑^3`$ representing the independent and dependent variables $`x^1,x^2,w`$ involved in our basic equation (2). The infinitesimal generator of such a group is a vector field $`X`$ on $`𝐑^3`$,
$$X=\xi ^\mu (x,w)\frac{}{x^\mu }+\eta (x,w)\frac{}{w},$$
(12)
whose components $`\xi ^\mu (x,w)`$ and $`\eta `$$`(x,w)`$ are supposed to be functions of class $`C^{\mathrm{}}`$on $`\mathrm{\Omega }`$. By virtue of Theorem 2.31 (Olver, 1993), a vector field $`X`$ of form (12) generates a point Lie symmetry group of equation (2) if and only if there exists a function $`\lambda `$ depending on $`x`$, $`w`$ and derivatives of $`w`$ (that is a differential function) such that the following infinitesimal criterion of invariance,
$$\underset{4}{𝑋}\left(𝒟[w]\right)\lambda 𝒟[w]=0,$$
(13)
holds; here $`\underset{k}{𝑋}`$ denote the $`k`$-th prolongation of $`X`$ (Ovsiannikov, 1982).
The invariance criterion (13) leads, through the standard computational procedure (see, e.g. Ovsiannikov, 1982 or Olver, 1993), to the following results:
(i) each equation of form (2) being linear and homogeneous is invariant under the point Lie groups generated by the vector fields
$$X_0=w\frac{}{w},X_u=u\left(x\right)\frac{}{w},$$
(14)
where $`u\left(x\right)`$ is an arbitrary solution of the equation considered, the invariance criterion (13) being fulfilled with $`\lambda =1`$ for $`X_0`$, and $`\lambda =0`$ for the generators $`X_u`$;
(ii) an equation of form (2) admits other vector fields (12), in addition to the aforementioned (14), if and only if they have the special form
$$X=\xi ^\mu \left(x\right)\frac{}{x^\mu }+\sigma \left(x\right)w\frac{}{w},$$
(15)
the functions $`\xi ^\mu \left(x\right)`$ and $`\sigma \left(x\right)`$ being nontrivial solutions of the following system of determining equations (called further the DE system for easy reference):
$$\xi ^\mu A_{,\mu }^{\alpha \beta \gamma \delta }+(\sigma \lambda )A^{\alpha \beta \gamma \delta }A^{\alpha \beta \gamma \mu }\xi _{,\mu }^\delta A^{\alpha \beta \mu \delta }\xi _{,\mu }^\gamma A^{\alpha \mu \gamma \delta }\xi _{,\mu }^\beta A^{\mu \beta \gamma \delta }\xi _{,\mu }^\alpha =0,$$
(16)
| $`4A^{\alpha \beta \gamma \mu }\sigma _{,\mu }2A^{\alpha \beta \mu \nu }\xi _{,\mu \nu }^\gamma 2A^{\alpha \gamma \mu \nu }\xi _{,\mu \nu }^\beta 2A^{\beta \gamma \mu \nu }\xi _{,\mu \nu }^\alpha `$ |
| --- |
| $`+\xi ^\mu A_{,\mu }^{\alpha \beta \gamma }+(\sigma \lambda )A^{\alpha \beta \gamma }A^{\alpha \beta \mu }\xi _{,\mu }^\gamma A^{\alpha \mu \gamma }\xi _{,\mu }^\beta A^{\mu \beta \gamma }\xi _{,\mu }^\alpha =0,`$ |
(17)
| $`6A^{\alpha \beta \mu \nu }\sigma _{,\mu \nu }2A^{\alpha \mu \nu \sigma }\xi _{,\mu \nu \sigma }^\beta 2A^{\beta \mu \nu \sigma }\xi _{,\mu \nu \sigma }^\alpha `$ |
| --- |
| $`+3A^{\alpha \beta \mu }\sigma _{,\mu }(3/2)A^{\alpha \mu \nu }\xi _{,\mu \nu }^\beta (3/2)A^{\beta \mu \nu }\xi _{,\mu \nu }^\alpha `$ |
| $`+\xi ^\mu A_{,\mu }^{\alpha \beta }+(\sigma \lambda )A^{\alpha \beta }A^{\alpha \mu }\xi _{,\mu }^\beta A^{\mu \beta }\xi _{,\mu }^\alpha =0,`$ |
(18)
| $`4A^{\alpha \mu \nu \sigma }\sigma _{,\mu \nu \sigma }A^{\mu \nu \sigma \tau }\xi _{,\mu \nu \sigma \tau }^\alpha `$ |
| --- |
| $`+3A^{\alpha \mu \nu }\sigma _{,\mu \nu }A^{\mu \nu \sigma }\xi _{,\mu \nu \sigma }^\alpha `$ |
| $`+2A^{\alpha \mu }\sigma _{,\mu }A^{\mu \nu }\xi _{,\mu \nu }^\alpha `$ |
| $`+\xi ^\mu A_{,\mu }^\alpha +(\sigma \lambda )A^\alpha A^\mu \xi _{,\mu }^\alpha =0,`$ |
(19)
$$A^{\alpha \beta \gamma \delta }\sigma _{,\alpha \beta \gamma \delta }+A^{\alpha \beta \gamma }\sigma _{,\alpha \beta \gamma }+A^{\alpha \beta }\sigma _{,\alpha \beta }+A^\alpha \sigma _{,\alpha }+\xi ^\mu A_{,\mu }+(\sigma \lambda )A=0,$$
(20)
for a certain function $`\lambda `$ depending on $`x^1`$ and $`x^2`$ only. (Here, by a trivial solution we mean not only $`\xi ^\mu =0,\sigma =0,`$ but also $`\xi ^\mu =0,\sigma =c=const0`$, since the latter leads to the vector field $`cX_0`$ generating the same group as $`X_0`$ which is already identified to be admitted by each equation of the type considered.)
Thus, given an equation of form (2), the question is whether there exist vector fields $`XcX_0`$ of form (15) which leave it invariant, and the answer depends on whether the respective DE system has at least one nontrivial solution. In this context the coefficients of the equation are supposed to be known functions, and thereby (16) – (20) constitute an over-determined system of linear homogeneous partial differential equations with respect to the unknowns $`\xi ^\mu `$ and $`\sigma `$. Therefore, as a rule, it turns out possible to find in an explicit form some (or even all) nontrivial solutions of the DE system, and thus to determine several (all) additional point Lie symmetry groups inherent to the equation in question.
It should be remarked that various equations of form (2) admit only the point Lie groups generated by the vector fields (14) with $`u(x)`$ being any solution of the respective equation. For instance, it is easy to check that all equations of the form (11) such that $`\chi ^{\alpha \beta }=\delta ^{\alpha \beta }`$ and $`\kappa (x)=p(x)`$, where $`p(x)`$ is an arbitrary polynomial of $`x^1`$ and $`x^2`$, belong to this variety. Without too much difficulties one can ascertain that the same holds true for the equations of the form (8) with $`D=const`$, $`N^{\alpha \beta }=\delta ^{\alpha \beta }`$ and $`k(x)=p(x)`$.
On the other hand, there are equations of the foregoing type which are invariant under a larger group; an immediate example is the biharmonic equation, $`\mathrm{\Delta }^2w=0`$, which admits the seven-parameter group generated by the linear combinations of $`X_0`$ and the following six additional basic vector fields (cf. Ovsiannikov, 1972):
$`X_1={\displaystyle \frac{}{x^1}},X_3=x^2{\displaystyle \frac{}{x^1}}x^1{\displaystyle \frac{}{x^2}},X_5=2x^1x^2{\displaystyle \frac{}{x^1}}\left[\left(x^1\right)^2\left(x^2\right)^2\right]{\displaystyle \frac{}{x^2}}+2x^2w{\displaystyle \frac{}{w}},`$
$`X_2={\displaystyle \frac{}{x^2}},X_4=x^1{\displaystyle \frac{}{x^1}}+x^2{\displaystyle \frac{}{x^2}},X_6=\left[\left(x^1\right)^2\left(x^2\right)^2\right]{\displaystyle \frac{}{x^1}}+2x^1x^2{\displaystyle \frac{}{x^2}}+2x^1w{\displaystyle \frac{}{w}}.`$
An important problem naturally arises in the light of the above note. It may be placed in the category of the so-called group-classification problems (see Ovsiannikov, 1982) and consist in determination of all those equations of the type considered that admit a larger group together with this group itself. Its most general statement assumes all functions $`A^{\alpha \beta \gamma \delta }(x),A^{\alpha \beta \gamma }(x),A^{\alpha \beta }(x),A^\alpha (x),A(x),\xi ^\alpha (x)`$ and $`\sigma (x)`$ involved in the determining equations (16) – (20) to be regarded as unknown variables and to find all solutions of this system. Here we are not going to study this rather complicated nonlinear problem in general. However, in Sections 4 and 5, restricting our attention to the equations of form (8), $`D=const`$, and (11), respectively, we will examine the corresponding group-classification problems.
Let us now specialize to the case of self-adjoint equations of form (2). Suppose that
$$𝒟[w]=0,𝒟=𝒟^{},$$
(21)
is such an equation. Then, of particular interest are its variational symmetries – the Lie groups generated by the so-called infinitesimal divergence symmetries (see Definition 4.33 in Olver, 1993) of any variational functional with (21) as the associated Euler-Lagrange equation. (Note that if two functionals lead to the same Euler-Lagrange equation, then they have the same collection of infinitesimal divergence symmetries.) This interest is motivated by the fact that, in virtue of Noether’s theorem, each variational symmetry of a given self-adjoint equation corresponds to a conservation law admitted by the smooth solutions of the equation. Thus, if a vector field $`X`$ of form (12) is found to generate a variational symmetry of equation (21), then Noether’s theorem implies the existence of a conserved current, which, in the present case, is a couple of differential functions $`P^\alpha `$ such that
$$D_\alpha P^\alpha =Q𝒟[w],$$
(22)
where $`Q`$ is the characteristic of $`X`$; by definition
$$Q=\eta w_\mu \xi ^\mu .$$
(23)
The total divergence of the conserved current $`P^\alpha `$ vanishes on the smooth solutions of (21) and so we have the conservation law
$$D_\alpha P^\alpha =0,$$
(24)
(22) being its expression in characteristic form, and $`Q`$ – its characteristics. Therefore, to derive the conservation laws of the foregoing type, one can proceed by first determining the variational symmetries of equation (21), and than using their characteristics (23) to find, from (22), explicit expressions for the corresponding conserved currents.
Having analyzed earlier the invariance properties of the whole class of equations (2), it is convenient to base the determination of the variational symmetries of equation (21) on the following observation. A vector field $`X`$ of form (12) generates a variational symmetry of equation (21) if and only if $`X`$ is an infinitesimal symmetry of this equation, that is (13) holds, and
$$\underset{4}{𝑋}\left(𝒟[w]\right)+\left(\frac{\eta }{w}+D_\mu \xi ^\mu \right)𝒟[w]=0.$$
(25)
This is a consequence of Lemma 4.34 and Proposition 5.55 (Olver, 1993), see also Lemma 7.46 (Olver, 1995). Subtracting (13) from (25) we can replace the latter with
$$\left(\frac{\eta }{w}+D_\mu \xi ^\mu +\lambda \right)𝒟[w]=0,$$
and as $`𝒟[w]`$ is not supposed to vanish identically we arrive at the conclusion that
$$\frac{\eta }{w}+D_\mu \xi ^\mu +\lambda =0,$$
(26)
is a necessary and sufficient condition for an infinitesimal symmetry admitted by a self-adjoint equation of form (2) to be its infinitesimal variational symmetry as well. It should be remarked that the same holds true for any self-adjoint partial differential equation in one dependent variable $`w`$ and $`n`$ independent variables $`x=(x^1,\mathrm{},x^n)`$; of course, in the general case the summation index $`\mu `$ will take the values $`1,\mathrm{},n`$ in both formulae (12) and (26). For a vector field of form (15) the relation (26) simplifies, and reads
$$\sigma +\xi _{,\mu }^\mu +\lambda =0.$$
(27)
Thus to find the variational symmetries of an equation of form (21), it suffices to check which of its ”ordinary” symmetries satisfy the additional requirement (26). For instance, the result (i) implies that $`X_0=w/w`$ does not generate a variational symmetry of any equation of form (21), while a vector field $`X_u=u(x)/w`$ generates a variational symmetry of an equation of form (21) whenever $`u(x)`$ is its solution (this is a common property of all systems of linear homogeneous partial differential equations, see Section 5.3 in Olver, 1993).
Suppose one has established that a vector field $`X`$ with characteristic $`Q`$ generates a variational symmetry of a given equation of form (21), and now wishes to find the conserved current $`P^\alpha `$ of the corresponding conservation law (24). For this purpose one can use formulae (5.150) and (5.151) given by Olver (1993) which express (in an explicit form) a null Lagrangian as a divergence. Indeed, in this case the right hand side of (22) is a total divergence or, in other words, a null Lagrangian. However, bearing in mind the recommendation of Olver (1993) to use these formulae only as a last resort since ”the homotopy formula (5.151) can rapidly become unmanageable”, in the present paper we suggest another way for determination of the sought conserved currents.
Our starting point is the so-called Noether identity (cf. Ibragimov, 1985):
$$\underset{\mathrm{}}{𝑋}()+(D_\alpha \xi ^\alpha )=Q𝖤()+D_\alpha N^\alpha (),$$
(28)
which holds for any differential function $``$ and vector field $`X`$ of the types considered here. In (28), $`N^\alpha `$ are the differential operators given by the expressions
$`N^\alpha `$ $`=`$ $`\xi ^\alpha +Q\left\{{\displaystyle \frac{}{w_\alpha }}+{\displaystyle \underset{s1}{}}(1)^sD_{\nu _1}\mathrm{}D_{\nu _s}{\displaystyle \frac{}{w_{\alpha \nu _1\mathrm{}\nu _s}}}\right\}`$
$`+{\displaystyle \underset{r1}{}}\left(D_{\mu _1}\mathrm{}D_{\mu _r}Q\right)\left\{{\displaystyle \frac{}{w_{\alpha \mu _1\mathrm{}\mu _r}}}+{\displaystyle \underset{s1}{}}(1)^sD_{\nu _1}\mathrm{}D_{\nu _s}{\displaystyle \frac{}{w_{\alpha \mu _1\mathrm{}\mu _r\nu _1\mathrm{}\nu _s}}}\right\},`$
$`Q=\eta \xi ^\alpha w_\alpha `$ being the characteristic of the vector field $`X`$. Setting $`=L`$ in (28), and taking into account (6) and (7), after a little manipulation we obtain the identity
$$D_\mu N^\mu (w𝒟[w])=w\underset{4}{𝑋}(𝒟[w])\{\eta +(D_\mu \xi ^\mu )w2Q\}𝒟[w],$$
(30)
valid for any self-adjoint differential operator $`𝒟`$ of form (3) and vector field of form (12).
In particular, for $`X_v=v(x)/w`$, where $`v(x)`$ is an arbitrary smooth function, we have
$$\xi ^\alpha =0,Q=\eta =v,$$
(31)
and hence
$$\underset{4}{X_v}\left(𝒟[w]\right)=𝒟[v],$$
(32)
since $`𝒟`$ is a linear differential operator. Substituting (31) and (32) into (30) we obtain
$$D_\mu N^\mu (w𝒟[w])=v𝒟[w]w𝒟[v],$$
(33)
which is nothing but the reciprocity relation associated with the equation $`𝒟[w]=0`$. Under the additional assumption $`v=u(x)`$, where $`u(x)`$ is an arbitrary smooth solution of the latter equation, the reciprocity relation (33) becomes
$$D_\mu N^\mu (w𝒟[w])=u𝒟[w].$$
(34)
Taking into account (34) we can give now the following general formula for the conserved currents $`P^\alpha `$ of the conservation laws with characteristics $`Q=u`$ corresponding to the infinitesimal variational symmetries $`X_u=u/w`$ of equation (21):
$$P^\alpha =P_{(u)}^\alpha +G^\alpha ,$$
where
$$P_{(u)}^\alpha =N^\alpha (wD[w]),$$
(35)
and $`G^\alpha `$ is any null divergence. Of course,
$$D_\mu P_{(u)}^\mu =u𝒟[w],$$
and
$$D_\mu P_{(u)}^\mu =0,$$
(36)
on the smooth solutions of the equation (21).
Next, let $`X`$ be an infinitesimal variational symmetry of equation (21) with characteristic $`Q=w\sigma w_\mu \xi ^\mu `$. Then, on account of (25), (30) takes the form
$$D_\mu N^\mu \left(\frac{1}{2}w𝒟[w]\right)=Q𝒟[w],$$
and hence we can write down the following explicit formula for the conserved currents $`P^\alpha `$ of the conservation laws with characteristics $`Q=w\sigma w_\mu \xi ^\mu `$ corresponding to the aforementioned variational symmetries of equation (21), namely
$$P^\alpha =B^\alpha +G^\alpha ,$$
$$B^\alpha =N^\alpha (\frac{1}{2}wD[w])+\frac{1}{2}D_\mu \left(w\xi ^\alpha A^{\mu \beta \gamma \delta }D_\beta D_\gamma D_\delta ww\xi ^\mu A^{\alpha \beta \gamma \delta }D_\beta D_\gamma D_\delta w\right),$$
(37)
where, as before, $`G^\alpha `$ is any null divergence. Of course,
$$D_\mu B^\mu =Q𝒟[w],$$
and on the smooth solutions of respective equation (21) we have
$$D_\mu B^\mu =0,$$
Let us remark that the special null divergence
$$\frac{1}{2}D_\mu \left(w\xi ^\alpha A^{\mu \beta \gamma \delta }D_\beta D_\gamma D_\delta ww\xi ^\mu A^{\alpha \beta \gamma \delta }D_\beta D_\gamma D_\delta w\right),$$
is used in the expression (37) for the conserved current $`B^\alpha `$ to cut the fourth-order derivatives of the dependent variable $`w`$ away since in practice we are usually interested in conserved currents which involve derivatives of order not higher than $`k1`$, where $`k`$ is the order of the equation considered. Making use of (S0.Ex11) it is easy to check that the right-hand side of (37) incorporates derivatives of $`w`$ of order less than fourth. In the subsequent Sections just (35) and (37) will be referred to as the expressions for the conserved currents of the conservation laws with characteristics $`Q=u`$ and $`Q=w\sigma w_\mu \xi ^\mu `$, respectively, derived for equations of the form (21).
To summarize, given an equation of form (21), the crucial point on the way of deriving conservation laws admitted by its smooth solutions is to find vector fields of form (15) generating ”ordinary” point Lie symmetries of the given equation. For that purpose, we should look for solutions of the respective DE system (16) – (20). Once such vector fields are found, it is easy to check which of their linear combinations satisfy the requirement (26) and hence generate variational symmetries of the equation considered. Now, using the characteristics of these symmetries we first construct the operators $`N^\alpha `$ from formulae (S0.Ex11) and then calculate from (37) the conserved currents of the corresponding conservation laws.
4. SYMMETRIES, CONSERVATION LAWS AND GROUP-INVARIANT SOLUTIONS OF PLATE EQUATIONS
In Section 2 (Example 1), we have quoted the self-adjoint equation (8) describing the small bending of a plate resting on an elastic foundation. Many problems concerning stability and vibration of isotropic thin elastic plates are studied on the ground of this type of equations. Here, we analyze the invariance properties of a generic equation of this form under the assumption that the bending rigidity of the plates considered is uniform, that is $`D=const`$. In this case (8) may be written as follows
$$A^{\alpha \beta \gamma \delta }w_{\alpha \beta \gamma \delta }+A^{\alpha \beta }(x)w_{\alpha \beta }+A(x)w=0,$$
(38)
with
$$A^{\alpha \beta \gamma \delta }=\frac{1}{3}(\delta ^{\alpha \beta }\delta ^{\gamma \delta }+\delta ^{\alpha \gamma }\delta ^{\beta \delta }+\delta ^{\alpha \delta }\delta ^{\beta \gamma }),A_{,\mu }^{\alpha \mu }=0,$$
(39)
assuming that
$$A^{\alpha \beta }=\frac{1}{D}N^{\alpha \beta },A=\frac{1}{D}k.$$
In view of the general results of Section 3, it is clear that $`X_u=u(x)/w`$ generates a variational symmetry of any equation of form (38) whenever $`u(x)`$ is its solution, while $`X_0=w/w`$ alone could never generate a variational symmetry of an equation of form (38), though it always is its infinitesimal point Lie symmetry. Substituting (39) into the determining equations (16) – (20) and taking into account that $`A^{\alpha \beta \gamma }=0`$, $`A^\alpha =0`$, we obtain, after a straightforward computation, that an equation of form (38) is invariant under a point Lie group generated by a vector field $`X`$ of form (15), $`XcX_0`$, if and only if
$$\sigma =\frac{1}{2}\xi _{,\mu }^\mu ,$$
(40)
$$\delta ^{\alpha \mu }\xi _{,\mu }^\beta +\delta ^{\mu \beta }\xi _{,\mu }^\alpha \delta ^{\alpha \beta }\xi _{,\mu }^\mu =0,$$
(41)
$$\xi ^\mu A_{,\mu }^{\alpha \beta }A^{\alpha \mu }\xi _{,\mu }^\beta A^{\mu \beta }\xi _{,\mu }^\alpha +2\xi _{,\mu }^\mu A^{\alpha \beta }+2\delta ^{\alpha \tau }\delta ^{\beta \nu }\xi _{,\mu \tau \nu }^\mu =0,$$
(42)
$$A^{\alpha \nu }\xi _{,\mu \nu }^\mu A^{\mu \nu }\xi _{,\mu \nu }^\alpha =0,$$
(43)
$$2\xi ^\mu A_{,\mu }+4\xi _{,\mu }^\mu A+A^{\mu \nu }\xi _{,\tau \mu \nu }^\tau =0.$$
(44)
At that,
$$\lambda =\frac{3}{2}\xi _{,\mu }^\mu .$$
(45)
Substituting expression (40) into (15), and expressions (40) and (45) into condition (27), we immediately arrive at the conclusion that the generator of such a group is a vector field of form
$$X=\xi ^\mu \frac{}{x^\mu }+\frac{1}{2}\xi _{,\mu }^\mu w\frac{}{w},$$
(46)
each such symmetry of (38) being variational symmetry of the latter equation as well, and hence there exist a conservation law with characteristic $`Q=(1/2)\xi _{,\mu }^\mu ww_\mu \xi ^\mu `$ and conserved current $`B^\alpha `$ given by (37) admitted by the smooth solutions of the equation considered. Thus to derive the conservation laws, which correspond to the variational symmetries of an equation of form (38) it suffices to know the results of the group classification of the class of equations in question; of course, the same holds true for the derivation of group-invariant solutions to (38). This group-classification problem is studied in Vassilev (1988), (1991) and (1997). The classification results presented below are obtained in these works.
It is shown that the scalar fields
$$s_{\left(1\right)}=A^{\mu \nu }\delta _{\mu \nu },s_{\left(2\right)}=\left(8A\delta _{\alpha \mu }\delta _{\beta \nu }A^{\alpha \beta }A^{\mu \nu }\right)^{1/2},s_{\left(3\right)}=\left(\delta ^{\mu \nu }s_{\left(1\right),\mu }s_{\left(1\right),\nu }\right)^{1/3},$$
(47)
are of key importance for the group classification of the considered class of equations. These scalar fields are called the invariants of equation (38) since here they play a role similar to the role that Laplace’s and Cotton’s invariants play in the group classification of the second-order linear partial differential equations (see Ovsiannikov, 1982 and Ibragimov, 1985). The following two properties of the scalar fields (47) give us both an additional reason to call them invariants of (38) and explicit expressions for the invariants of groups admitted by (38). First, if an equation of form (38) admits a vector field of form (46), then
$$\xi _{,\mu }^\mu s_{(j)}+\xi ^\mu s_{(j),\mu }=0(j=1,2,3),$$
and hence $`U_{(j)}=w\sqrt{s_{(j)}}`$ are invariants of the corresponding Lie group whenever $`s_{(j)}0`$. Second, if an equation of form (38) admits a vector field of form (46) and is such that at least two of its invariants (47), say $`s_{(k)}`$ and $`s_{(l)}(kl3)`$ , are not identically equal to zero, then $`s_{(k)}/s_{(l)}`$ is an invariant of the corresponding symmetry group. Note, that the invariants $`s_{(k)}`$ and $`s_{(l)}`$of such an equation of form (38) provide two couples of functionally independent invariants, namely $`U_{(k)}=w\sqrt{s_{(k)}}`$ and $`s_{(k)}/s_{(l)}`$ as well as $`U_{(l)}=w\sqrt{s_{(l)}}`$ and $`s_{(k)}/s_{(l)}`$, of the admitted symmetry group, both couples being readily applicable for constructing group-invariant solutions to the respective equation. However, if even one of the invariants (47) of an equation of form (38) is not identically equal to zero, then this equation admits at most a 3-parameter group with generators of form (46). On the other hand, if all invariants (47) of an equation of form (38) are identically equal to zero, then this equation admits a 6-parameter group with generators of form (46). Below, the latter case is set out in detail.
Let $`\omega (z)const`$ be an analytic function of the complex variable $`z=x^1+ix^2`$, and let $`E_\omega `$ be the equation of the form (38) with coefficients
$$A^{11}=A^{22}=4\mathrm{R}e\left\{\varphi \right\},A^{12}=A^{21}=4\mathrm{I}m\left\{\varphi \right\},A=4\varphi \overline{\varphi },$$
(48)
where $`\varphi `$ is the Schwarzian derivative of the function $`\omega `$, that is
$$\varphi =\left(\frac{\omega ^{\prime \prime }}{\omega ^{}}\right)^{}\frac{1}{2}\left(\frac{\omega ^{\prime \prime }}{\omega ^{}}\right)^2,$$
(49)
$`\overline{\varphi }`$ is the complex conjugated of $`\varphi `$, and the prime is used to denote differentiation with respect to the variable $`z`$. Substituting (48) into (47) one can see that all invariants of $`E_\omega `$ are identically equal to zero. Then, taking into account the DE system (41)–(44), (48) and (49), one can verify by direct computing that $`E_\omega `$ admits the 6-parameter group generated by the vector fields
$$Z_{(j)}=\xi _{(j)}^\mu \frac{}{x^\mu }+\frac{1}{2}\xi _{(j),\mu }^\mu w\frac{}{w}(j=1,\mathrm{},6),$$
the functions $`\xi _{(j)}^\mu `$ being given by the expressions
| $`\xi _{(1)}^1=\mathrm{R}e\left\{\omega _1\right\},`$ | $`\xi _{(1)}^2=\mathrm{I}m\left\{\omega _1\right\},`$ |
| --- | --- |
| $`\xi _{(2)}^1=\mathrm{R}e\left\{i\omega _1\right\},`$ | $`\xi _{(2)}^2=\mathrm{I}m\left\{i\omega _1\right\},`$ |
| $`\xi _{(3)}^1=\mathrm{R}e\left\{\omega _2\right\},`$ | $`\xi _{(3)}^2=\mathrm{I}m\left\{\omega _2\right\},`$ |
| $`\xi _{(4)}^1=\mathrm{R}e\left\{i\omega _2\right\},`$ | $`\xi _{(4)}^2=\mathrm{I}m\left\{i\omega _2\right\},`$ |
| $`\xi _{(5)}^1=\mathrm{R}e\left\{\omega _3\right\},`$ | $`\xi _{(5)}^2=\mathrm{I}m\left\{\omega _3\right\},`$ |
| $`\xi _{(6)}^1=\mathrm{R}e\left\{i\omega _3\right\},`$ | $`\xi _{(6)}^2=\mathrm{I}m\left\{i\omega _3\right\}.`$ |
where
$$\omega _1=\frac{1}{\omega ^{}},\omega _2=\frac{\omega }{\omega ^{}},\omega _3=\frac{\omega ^2}{\omega ^{}}.$$
(50)
It should be remarked that each equation of form (38) which admits a 6-parameter group with generators of form (46) is of type $`E_\omega `$, meaning that it can be generated in the above manner using a suitable analytic function $`\omega `$. The coefficients of each equation of this type are of the form $`A^{\alpha \beta }=\delta ^{\alpha \mu }\delta ^{\beta \nu }\phi _{,\mu \nu }`$, $`A=(1/8)\delta ^{\alpha \mu }\delta ^{\beta \nu }\phi _{,\alpha \beta }\phi _{,\mu \nu }`$, where $`\phi `$ is a harmonic function, that is $`\delta ^{\alpha \beta }\phi _{,\alpha \beta }=0`$, and vice versa. It is noteworthy that each equation with variable coefficients of type $`E_\omega `$ can be mapped to an equation with constant coefficients belonging to the same family. It is easy to verify by direct computing that the equation $`E_\omega `$ corresponding to an analytic function $`\omega `$ whose Schwarzian derivative is not constant transforms to a constant coefficients one under the following change of the independent and dependent variables:
$$y^\alpha =f^\alpha (x^1,x^2),W=wU(x^1,x^2),$$
(51)
$$f^1(x^1,x^2)=\mathrm{R}e\{f^1dz\},f^2(x^1,x^2)=\mathrm{I}m\{f^1dz\},U(x^1,x^2)=\left(f\stackrel{}{f}\right)^{1/2},$$
where $`f`$ is any linear combination of the functions (50) such that $`f0`$, i.e.
$$f=k_1\omega _1+k_2\omega _2+k_3\omega _3,$$
(52)
where $`k_1`$, $`k_2`$, and $`k_3`$ are complex constants such that $`k_1^2+k_2^2+k_3^20`$.
Consider, as a simple example, the equation $`E_\omega `$ corresponding to $`\omega =z^{\sqrt{\varkappa /2}}`$, where $`\varkappa `$ is a positive real constant. In this case (49) gives $`\varphi =4\left(2\varkappa \right)z^2`$ and hence, according to (48), the coefficients of $`E_\omega `$ read
$$A^{11}=A^{22}=\left(2\varkappa \right)\frac{\left(x^1\right)^2\left(x^2\right)^2}{\left[\left(x^1\right)^2+\left(x^2\right)^2\right]^2},A^{12}=A^{21}=\left(2\varkappa \right)\frac{2x^1x^2}{\left[\left(x^1\right)^2+\left(x^2\right)^2\right]^2},$$
$$A=\left(2\varkappa \right)^2\frac{1}{4\left[\left(x^1\right)^2+\left(x^2\right)^2\right]^2}.$$
(53)
Using the function $`f=z`$ obtained from (52) for $`\omega =z^{\sqrt{\varkappa /2}}`$, $`k_1=k_3=0`$, $`k_2=1+\sqrt{\varkappa /2}`$, we introduce, according to (51), the new independent and dependent variables
$$y^1=\frac{1}{2}\mathrm{ln}\left[\left(x^1\right)^2+\left(x^2\right)^2\right],y^2=\mathrm{arctan}\left(\frac{x^2}{x^1}\right),W=w\left[\left(x^1\right)^2+\left(x^2\right)^2\right]^{1/2}.$$
(54)
Note that the inverse transformations are given by the expressions
$$x^1=e^{y^1}\mathrm{cos}y^2,x^2=e^{y^1}\mathrm{sin}y^2,w=W\left[\left(x^1\right)^2+\left(x^2\right)^2\right]^{1/2}.$$
(55)
Under the change of the variables according to (54), the considered equation $`E_\omega `$ transforms to the following one,
$$\delta ^{\alpha \beta }\delta ^{\mu \nu }\frac{^4W}{y^\alpha y^\beta y^\mu y^\nu }\varkappa \frac{^2W}{y^1y^1}+\varkappa \frac{^2W}{y^2y^2}+\frac{1}{4}\varkappa ^2W=0,$$
(56)
which belongs to the same class (since it corresponds to the analytic function $`\omega =e^{z\sqrt{\varkappa /2}}`$) but whose coefficients are constant.
Equation (56) admits the 6-parameter group of variational symmetries generated by the basic vector fields
$`V_\alpha `$ $`=`$ $`{\displaystyle \frac{}{y^\alpha }},`$
$`V_3`$ $`=`$ $`e^{\theta y^1}\mathrm{cos}\left(\theta y^2\right){\displaystyle \frac{}{y^1}}+e^{\theta y^1}\mathrm{sin}\left(\theta y^2\right){\displaystyle \frac{}{y^2}}+\theta e^{\theta y^1}w\mathrm{cos}\left(\theta y^2\right){\displaystyle \frac{}{w}},`$
$`V_4`$ $`=`$ $`e^{\theta y^1}\mathrm{cos}\left(\theta y^2\right){\displaystyle \frac{}{y^1}}+e^{\theta y^1}\mathrm{cos}\left(\theta y^2\right){\displaystyle \frac{}{y^2}}\theta e^{\theta y^1}w\mathrm{sin}\left(\theta y^2\right){\displaystyle \frac{}{w}},`$
$`V_5`$ $`=`$ $`e^{\theta y^1}\mathrm{cos}\left(\theta y^2\right){\displaystyle \frac{}{y^1}}e^{\theta y^1}\mathrm{sin}\left(\theta y^2\right){\displaystyle \frac{}{y^2}}\theta e^{\theta y^1}w\mathrm{cos}\left(\theta y^2\right){\displaystyle \frac{}{w}},`$
$`V_6`$ $`=`$ $`e^{\theta y^1}\mathrm{cos}\left(\theta y^2\right){\displaystyle \frac{}{y^1}}+e^{\theta y^1}\mathrm{sin}\left(\theta y^2\right){\displaystyle \frac{}{y^2}}\theta e^{\theta y^1}w\mathrm{cos}\left(\theta y^2\right){\displaystyle \frac{}{w}},`$
where $`\theta =\sqrt{\varkappa /2}`$. These vector fields give rise to six linearly independent conservation laws for equation (56). The characteristics of these conservation laws are
$$Q_{\left(j\right)}=\frac{1}{2}W\frac{}{y^\mu }V_j\left(y^\mu \right)W_\mu V_j\left(y^\mu \right)\left(j=1,\mathrm{},6\right).$$
Here, $`V_j`$ are regarded as operators acting on the functions $`\zeta :R^2R`$. The corresponding conserved currents can be easily calculated from (37).
Finally, let us remark that each one-parameter group generated by a linear combination of the basic vector fields $`V_j`$ can be used for constructing group-invariant solutions of equation (56). Consider, for instance, the group $`H\left(V_3+V_5\right)`$ generated by the vector field $`V_3+V_5`$. The functions $`s=\mathrm{sin}\left(\theta y^2\right)/\mathrm{cosh}\left(\theta y^1\right)`$ and $`u=W/\mathrm{cosh}\left(\theta y^1\right)`$ constitute a complete set of invariants for this group and hence, following the well known algorithm (Ovsiannikov, 1982; Olver, 1993), we seek the $`H\left(V_3+V_5\right)`$–invariant solutions of equation (56) in the form
$$W=u\left(s\right)\mathrm{cosh}\left(\theta y^1\right),s=\frac{\mathrm{sin}\left(\theta y^2\right)}{\mathrm{cosh}\left(\theta y^1\right)}.$$
(57)
Substituting (57) into (56), we get the reduced equation
$$\left(s^21\right)^2\frac{d^4u}{ds^4}+8s\left(s^21\right)\frac{d^3u}{ds^3}+4\left(3s^21\right)\frac{d^2u}{ds^2}=0.$$
The general solution to this ordinary differential equation is
$$u\left(s\right)=C_1+C_2\mathrm{ln}\left(\frac{s+1}{s1}\right)+C_3s+C_4s\mathrm{ln}\left(\frac{s+1}{s1}\right),$$
where $`C_1`$, $`C_2`$, $`C_3`$ and $`C_4`$ are real constants. Hence the $`H\left(V_3+V_5\right)`$–invariant solutions of equation (56) are given by the expression
$$W(y^1,y^2)=C_1\mathrm{cosh}\left(\theta y^1\right)+C_2\mathrm{cosh}\left(\theta y^1\right)\mathrm{ln}\left[\frac{\mathrm{sin}\left(\theta y^2\right)+\mathrm{cosh}\left(\theta y^1\right)}{\mathrm{sin}\left(\theta y^2\right)\mathrm{cosh}\left(\theta y^1\right)}\right]$$
$$+C_3\mathrm{sin}\left(\theta y^2\right)+C_4\frac{\mathrm{sin}\left(\theta y^2\right)}{\mathrm{cosh}\left(\theta y^1\right)}\mathrm{ln}\left[\frac{\mathrm{sin}\left(\theta y^2\right)+\mathrm{cosh}\left(\theta y^1\right)}{\mathrm{sin}\left(\theta y^2\right)\mathrm{cosh}\left(\theta y^1\right)}\right].$$
Using the inverse transformations (55) one can convert the above solutions of equation (56) into solutions of the equation $`E_\omega `$, $`\omega =z^{\sqrt{\varkappa /2}}`$, with variable coefficients (53).
4. SYMMETRIES, CONSERVATION LAWS AND GROUP-INVARIANT SOLUTIONS OF ROD EQUATIONS
In Section 2, combining and generalizing Examples 3 and 4 we have introduced the class of self-adjoint partial differential equations
$$A^{1111}w_{1111}+A^{\alpha \beta }w_{\alpha \beta }+Aw=0,$$
(58)
with coefficients
$$A^{1111}=\gamma ,A^{\alpha \beta }=\chi ^{\alpha \beta },A=\kappa (x),$$
(59)
where $`\gamma =const0`$, $`\chi ^{\alpha \beta }`$ are arbitrary constants (but $`\left(\chi ^{12}\right)^2+\left(\chi ^{22}\right)^20`$, otherwise (58) degenerates and becomes an ordinary differential equation), and $`\kappa (x)`$ is an arbitrary function. Equations of this special type are used by many authors to study applied engineering problems concerning dynamics and stability of both elastic beams resting on elastic foundations (see e.g. Smith and Herrmann, 1972) and pipes conveying fluid (see e.g. Gregory and Paidoussis, 1966). In the present Section we first examine the point Lie symmetries of (58) and solve the corresponding group-classification problem. Then we derive conservation laws and group-invariant solutions of various rod equations of form (58).
Consider the group-classification problem. In view of the results (i) and (ii) of Section 3, it is clear that each equation of form (58) is invariant under the point Lie groups generated by the vector fields $`X_0=w/w`$ and $`X_u=u\left(x\right)/w`$ where $`u\left(x\right)`$ is any smooth solution of the respective equation and the objective is to find those equations of the type considered which admit vector fields $`X`$ of form (15), $`XcX_0`$, $`c=const0`$.
Substituting (59) into (16) – (20), and taking into account that $`A^\alpha ,A^{\alpha \beta \gamma }`$ and all $`A^{\alpha \beta \gamma \delta }`$ except $`A^{1111}`$ are equal to zero, we obtain after a little manipulation the following system of determining equations for the functions $`\xi ^\mu \left(x\right)`$ and $`\sigma \left(x\right)`$ associated with the sought vector fields $`X`$ of form (15):
$$\xi _{,1}^2=0,2\sigma _{,1}3\xi _{,11}^1=0,$$
(60)
$$5\gamma \xi _{,111}^1+2\chi ^{11}\xi _{,1}^12\chi ^{12}\xi _{,2}^1=0,$$
(61)
$$\chi ^{22}(2\xi _{,1}^1\xi _{,2}^2)=0,\chi ^{12}(3\xi _{,1}^1\xi _{,2}^2)\chi ^{22}\xi _{,2}^1=0,$$
(62)
$$2\chi ^{12}\sigma _{,2}\chi ^{22}\xi _{,22}^1=0,2\chi ^{12}\sigma _{,1}+\chi ^{22}(2\sigma _{,2}\xi _{,22}^2)=0,$$
(63)
$$\sigma _{,1111}+\chi ^{\alpha \beta }\sigma _{,\alpha \beta }+\xi ^\mu \kappa _{,\mu }+4\xi _{,1}^1\kappa =0,$$
(64)
the auxiliary function $`\lambda `$ being expressed by
$$\lambda =\sigma 4\xi _{,1}^1.$$
(65)
We look for the equations of form (58) whose coefficients $`\gamma `$, $`\chi ^{\alpha \beta }`$ and $`\kappa \left(x\right)`$ are such that system (60) – (64) possesses solutions different from the trivial one $`\xi ^\mu =0`$, $`\sigma =const0`$.
Observing the system of determining equations (60) – (64) we see that for the purposes of the group-classification it is convenient to divide the equations of form (58) into three subclasses depending on whether the coefficients $`\chi ^{\alpha \beta }`$ of the given equation are such that: (A) $`\chi ^{22}0,det(\chi ^{\alpha \beta })0`$; (B) $`\chi ^{22}0,det(\chi ^{\alpha \beta })=0`$ or (C) $`\chi ^{22}=0,det(\chi ^{\alpha \beta })0`$. This covers all possibilities except for $`\chi ^{22}=0,det(\chi ^{\alpha \beta })=0`$ when (58) becomes an ordinary differential equation that falls outside our interest in the present paper.
For convenience we introduce the notation
$$Y_\alpha =\frac{}{x^\alpha },Y_3=\left(x^1+\frac{\chi ^{12}}{\chi ^{22}}x^2\right)\frac{}{x^1}+2x^2\frac{}{x^2},Y_4=\left(x^1+\frac{\chi ^{11}}{\chi ^{12}}x^2\right)\frac{}{x^1}+3x^2\frac{}{x^2}.$$
Let $`\chi ^{22}0`$. Then, the first equation (62) is equivalent to
$$2\xi _{,1}^1\xi _{,2}^2=0.$$
(66)
Differentiating (66) with respect to $`x^1`$ and taking into account the first equation (60) we obtain $`\xi _{,11}^1=0`$. Hence, (61) and the second equation (62) reduce to
$$\chi ^{11}\xi _{,1}^1\chi ^{12}\xi _{,2}^1=0,\chi ^{12}\xi _{,1}^1\chi ^{22}\xi _{,2}^1=0.$$
(67)
Consider the subclass (A): $`\chi ^{22}0,`$ $`det(\chi ^{\alpha \beta })0`$. Then, the first equation (60) together with (66) and (67) lead to $`\xi _{,1}^1=\xi _{,2}^1=\xi _{,1}^2=\xi _{,2}^2=0`$, i.e., $`\xi ^\alpha =c^\alpha =const`$. Consequently, the second equations of (60) and (63) imply $`\sigma =c=const`$. At this, the first equation (63) is satisfied and (64) becomes
$$c^\alpha \kappa _{,\alpha }=0.$$
(68)
All this means that when $`\chi ^{22}0`$ and $`det(\chi ^{\alpha \beta })0`$ the DE system has only the trivial solution unless
$$\kappa (x)=f(\beta ^2x^1\beta ^1x^2),$$
(69)
for a certain smooth function $`fconst`$ and certain constants $`\beta ^\alpha `$ such that $`\left(\beta ^1\right)^2+\left(\beta ^1\right)^20`$. In this latter case, the DE system has the nontrivial solution
$$\xi ^1=\beta ^1,\xi ^2=\beta ^2,\sigma =0,$$
(70)
and so the differential equations of that kind admit additionally the one-parameter symmetry group associated with the vector field $`\beta ^1Y_1+\beta ^2Y_2`$. In the case $`\kappa =const`$, (68) is satisfied for any couple of constants $`c^\alpha `$, and hence such equations admit the 2-parameter symmetry group with generators $`Y_1`$ and $`Y_2`$. This completes the analysis of subclass (A).
Consider now the subclass (B): $`\chi ^{22}0,det(\chi ^{\alpha \beta })=0`$. Differentiating the second equation (67) successively with respect to $`x^1`$ and $`x^2`$ we obtain
$$\chi ^{12}\xi _{,11}^1\chi ^{22}\xi _{,12}^1=0,\chi ^{12}\xi _{,12}^1\chi ^{22}\xi _{,22}^1=0,$$
which, on account of $`\xi _{,11}^1=0`$, leads to $`\xi _{,12}^1=\xi _{,22}^1=0`$. This result, together with (66) imply $`\xi _{,22}^2=0`$, and the nontrivial solution of (60) – (63) is now obvious:
$$\xi ^1=c^1+c^3\left(x^1+\frac{\chi ^{12}}{\chi ^{22}}x^2\right),\xi ^2=c^2+2c^3x^2,\sigma =0,$$
(71)
with $`c^1,c^2`$ and $`c^3`$ – arbitrary constants. Equation (64) reduces to
$$\xi ^\alpha \kappa _{,\alpha }+4\xi _{,1}^1\kappa =0.$$
(72)
For an arbitrary $`\kappa (x)`$ it leads to $`\xi ^\alpha =0`$; it is easily verified that for $`\kappa (x)=x^1+x^2x^2`$, $`\xi ^\alpha =0`$ is the only solution of (72).
If $`\kappa (x)`$ is a function of form (69), then (72) implies (70) as a nontrivial solution to the system (60) – (64). Therefore, a generic equation of this kind admits only the one-parameter group generated by $`\beta ^1Y_1+\beta ^2Y_2`$, unless $`\kappa (x)`$ is of one of the following two special forms. The first one is
$$\kappa (x)=\kappa _0\left(\beta +x^2\right)^2,\kappa _0=const0,\beta =const,$$
(73)
when the nontrivial solution of (60) – (64) is (71) with $`c^2=2\beta `$, $`c^3=1`$, $`c^1`$ – arbitrary, and hence the equations of subclass (B) with $`\kappa (x)`$ of form (73) admit the 2-parameter symmetry group generated by the vector fields $`Y_1`$ and $`2\beta Y_2+Y_3`$. The second special form of the function $`\kappa (x)`$ is
$$\kappa (x)=\kappa _0\left(\beta +x^1\frac{\chi ^{12}}{\chi ^{22}}x^2\right)^4,\kappa _0=const0,\beta =const,$$
(74)
when the respective differential equations admit the 2-parameter group spanned over the vector fields $`\beta Y_1+Y_3`$ and $`\left(\chi ^{12}/\chi ^{22}\right)Y_1+Y_2`$.
Another extension of the symmetry group is possible if there exist two constants $`\beta ^1`$ and $`\beta ^2`$, as well as a smooth function $`f0`$ such that
$$\kappa (x)=\left(\beta ^2+x^2\right)^2f(y),y=\left(\beta ^2+x^2\right)^{1/2}\left(\beta ^1+x^1\frac{\chi ^{12}}{\chi ^{22}}x^2\right).$$
(75)
If $`\kappa (x)`$ is of form (75), then the nontrivial solution of the determining equations (60) – (64) is (71) with $`c^1=\beta ^1+2\left(\chi ^{12}/\chi ^{22}\right)\beta ^2`$, $`c^2=2\beta ^2`$, $`c^3=1`$ and the differential equations of that sort admit the one-parameter group, generated by
$$\left(\beta ^1+2\frac{\chi ^{12}}{\chi ^{22}}\beta ^2\right)Y_1+2\beta ^2Y_2+Y_3$$
only, except for the cases $`f(y)=\kappa _0y^4`$ ($`\kappa _0=const`$), when $`\kappa (x)`$ takes the form (74), and $`f(y)=const0`$ when $`\kappa (x)`$ becomes (73).
The differential equations of form (58) with $`\chi ^{22}0`$, $`det(\chi ^{\alpha \beta })=0`$ admit the 2-parameter group generated by $`Y_1`$ and $`Y_2`$ when $`\kappa (x)=const0`$ or the 3-parameter group with generators $`Y_1`$, $`Y_2`$ and $`Y_3`$ when $`\kappa (x)=0`$.
Finally, consider the subclass (C): $`\chi ^{22}=0,det(\chi ^{\alpha \beta })0`$. Substituting $`\chi ^{22}=0`$ in the determining equations (60) – (63), the latter simplify to
$$\xi _{,1}^2=0,\xi _{,11}^1=0,\chi ^{11}\xi _{,1}^1\chi ^{12}\xi _{,2}^1=0,3\xi _{,1}^1\xi _{,2}^2=0,\sigma _{,1}=0,\sigma _{,2}=0,$$
and their nontrivial solution is easily obtained:
$$\xi ^1=c^1+c^3\left(x^1+\frac{\chi ^{11}}{\chi ^{12}}x^2\right),\xi ^2=c^2+3c^3x^2,\sigma =0,$$
(76)
where $`c^i`$ are arbitrary constants (note that if $`\chi ^{22}=0,`$ then $`det(\chi ^{\alpha \beta })0`$ assumes $`\chi ^{12}0`$). Equation (64) takes the form (72) and for an arbitrary $`\kappa (x)`$ it leads to $`\xi ^\alpha =0`$. If $`\kappa (x)`$ is of the form (69), such equations admit the one-parameter group generated by $`\beta ^1Y_1+\beta ^2Y_2`$ only, except for two special forms of $`\kappa (x)`$, namely
$$\kappa (x)=\kappa _0\left(\beta +x^2\right)^{4/3},$$
(77)
and
$$\kappa (x)=\kappa _0\left(\beta +2x^1\frac{\chi ^{11}}{\chi ^{12}}x^2\right)^4.$$
(78)
where $`\kappa _0`$ and $`\beta `$ are constants. In the case (77), the nontrivial solution of the determining equation (60) – (64) is (76) with $`c^2=3\beta `$, $`c^3=1`$, $`c^1`$ – arbitrary, and the differential equation considered admits the 2-parameter group spanned over the vector fields $`Y_1`$ and $`3\beta Y_2+Y_4`$. In the case (78) the group admitted is also a 2-parameter one, but generated by $`\beta Y_1+Y_4`$ and $`\left(\chi ^{11}/\chi ^{12}\right)Y_1+2Y_2`$.
Another extension of the symmetry group is possible if
$$\kappa (x)=\left(\beta ^2+x^2\right)^{4/3}f(y),y=\left(\beta ^2+x^2\right)^{1/3}\left(\beta ^1+2x^1\frac{\chi ^{11}}{\chi ^{12}}x^2\right),$$
(79)
where $`\beta ^\alpha `$ are constants and $`f\mathrm{\hspace{0.33em}0}`$ is a smooth function. In this case, the nontrivial solution of (60) – (64) is (76) with $`c^1=\beta ^1+3(\chi ^{11}/\chi ^{12})\beta ^2`$, $`c^2=6\beta ^2`$, $`c^3=1`$ and we conclude that the equations of subclass (C) with $`\kappa (x)`$ of form (79) admit only the one-parameter group generated by $`[\beta ^1+3(\chi ^{11}/\chi ^{12})\beta ^2]Y_1+6\beta ^2Y_2+Y_4`$, except for $`f(y)=\kappa _0y^4`$ ($`\kappa _0=const`$), when (79) coincides with (78) or $`f(y)=const0`$ when $`\kappa (x)`$ has the form (77).
Evidently, the differential equations of form (58) with $`\chi ^{22}=0`$, $`det(\chi ^{\alpha \beta })0`$ admit: the 2-parameter group with generators $`Y_1`$ and $`Y_2`$ when $`\kappa (x)=const0`$, and the 3-parameter group associated with $`Y_1`$, $`Y_2`$ and $`Y_4`$ – when $`\kappa (x)=0`$.
The results of the above group-classification analysis are summarized in Table 1, where the equations invariant under larger groups are given through their coefficients together with the generators of the associated symmetry groups.
Table 1. Equations of form (58) invariant under larger symmetry groups.
| # | Coefficients | Generators |
| --- | --- | --- |
| 1 | $`\kappa (x)=f(\beta ^2x^1\beta ^1x^2)`$ | $`\beta ^1Y_1+\beta ^2Y_2`$ |
| 2 | $`\begin{array}{c}\chi ^{22}0,det(\chi ^{\alpha \beta })=0,\kappa (x)=\left(\beta ^2+x^2\right)^2f(y),\hfill \\ y=\left(\beta ^2+x^2\right)^{1/2}[\beta ^1+x^1\left(\chi ^{12}/\chi ^{22}\right)x^2]\hfill \end{array}`$ | $`\begin{array}{c}[\beta ^1+2(\chi ^{12}/\chi ^{22})\beta ^2]Y_1\hfill \\ +2\beta ^2Y_2+Y_3\hfill \end{array}`$ |
| 3 | $`\begin{array}{c}\chi ^{22}=0,det(\chi ^{\alpha \beta })0,\kappa (x)=\left(\beta ^2+x^2\right)^{4/3}f(y),\hfill \\ y=\left(\beta ^2+x^2\right)^{1/3}[\beta ^1+2x^1\left(\chi ^{11}/\chi ^{12}\right)x^2]\hfill \end{array}`$ | $`\begin{array}{c}[\beta ^1+3(\chi ^{11}/\chi ^{12})\beta ^2]Y_1\hfill \\ +6\beta ^2Y_2+2Y_4\hfill \end{array}`$ |
| 4 | $`\chi ^{22}0,det(\chi ^{\alpha \beta })=0,\kappa (x)=\kappa _0\left(\beta +x^2\right)^2,`$ | $`Y_1`$, $`2\beta Y_2+Y_3`$ |
| 5 | $`\chi ^{22}=0,det(\chi ^{\alpha \beta })0,\kappa (x)=\kappa _0\left(\beta +x^2\right)^{4/3}`$ | $`Y_1`$, $`3\beta Y_2+Y_4`$ |
| 6 | $`\begin{array}{c}\chi ^{22}0,det(\chi ^{\alpha \beta })=0,\hfill \\ \kappa (x)=\kappa _0\left(\beta +x^1\left(\chi ^{12}/\chi ^{22}\right)x^2\right)^4\hfill \end{array}`$ | $`\begin{array}{c}\beta Y_1+Y_3,\hfill \\ (\chi ^{12}/\chi ^{22})Y_1+Y_2\hfill \end{array}`$ |
| 7 | $`\begin{array}{c}\chi ^{22}=0,det(\chi ^{\alpha \beta })0,\hfill \\ \kappa (x)=\kappa _0\left(\beta +2x^1\left(\chi ^{11}/\chi ^{12}\right)x^2\right)^4\hfill \end{array}`$ | $`\begin{array}{c}\beta Y_1+2Y_4,\hfill \\ (\chi ^{11}/\chi ^{12})Y_1+2Y_2\hfill \end{array}`$ |
| 8 | $`\chi ^{22}det(\chi ^{\alpha \beta })0,\kappa (x)=const`$ | $`Y_1`$, $`Y_2`$ |
| 9 | $`\chi ^{22}det(\chi ^{\alpha \beta })=0,\kappa (x)=const0`$ | $`Y_1`$, $`Y_2`$ |
| 10 | $`\chi ^{22}0,det(\chi ^{\alpha \beta })=0,\kappa (x)=0`$ | $`Y_1`$, $`Y_2`$, $`Y_3`$ |
| 11 | $`\chi ^{22}=0,det(\chi ^{\alpha \beta })0,\kappa (x)=0`$ | $`Y_1`$, $`Y_2`$, $`Y_4`$ |
Having completely solved the group-classification problem, our next step is to identify the variational symmetries of those equations of form (58) which have been found to admit larger symmetry groups. For this purpose, we are to apply the condition (26) to the linear combinations of (14) and the vector fields presented in Table 1, the respective functions $`\lambda `$ being given by (65). Thus, we found that all vector fields quoted under numbers 1, 3, 5, 7, 8, 9 and 11 generate variational symmetries as well. In case # 2 the variational symmetries are associated with $`[\beta ^1+2(\chi ^{12}/\chi ^{22})\beta ^2]Y_1+2\beta ^2Y_2+Y_3+(1/2)X_0`$. Similarly, in case # 4 the variational symmetries are generated by $`Y_1`$ and $`2\beta Y_2+Y_3+(1/2)X_0`$, in case # 6 – by $`(\chi ^{12}/\chi ^{22})Y_1+Y_2`$ and $`\beta Y_1+Y_3+(1/2)X_0`$, and in case # 10 – by $`Y_1`$, $`Y_2`$ and $`Y_3+(1/2)X_0`$.
Once the variational symmetries of the equations (58) are identified, we can derive the corresponding conservation laws. The conserved currents of the conservation laws for the equations given in Table 1 are computed using (37), the above notes concerning the corresponding variational symmetries being taken into account. The obtained conservation laws are listed in Table 2 (in the same order as in Table 1) in terms of the differential functions:
$`B_{(1)}^1`$ $`=`$ $`{\displaystyle \frac{1}{2}}[\gamma (2w_1w_{111}w_{11}^2)+\chi ^{11}w_1^2\chi ^{22}w_2^2+\kappa w^2]{\displaystyle \frac{1}{2}}(\chi ^{2\mu }ww_\mu ),_2,`$
$`B_{(1)}^2`$ $`=`$ $`\chi ^{2\mu }w_1w_\mu +{\displaystyle \frac{1}{2}}(\chi ^{2\mu }ww_\mu ),_1,`$
$`B_{(2)}^1`$ $`=`$ $`\chi ^{1\mu }w_2w_\mu +\gamma (w_{11}w_{12}w_2w_{111}){\displaystyle \frac{1}{2}}(\gamma w_1w_{11}\chi ^{1\mu }ww_\mu ),_2,`$
$`B_{(2)}^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\gamma w_{11}^2+\chi ^{22}w_2^2\chi ^{11}w_1^2+\kappa w^2)+{\displaystyle \frac{1}{2}}(\gamma w_1w_{11}\chi ^{1\mu }ww_\mu ),_1,`$
$`B_{(3)}^\alpha `$ $`=`$ $`\left(x^1+{\displaystyle \frac{\chi ^{12}}{\chi ^{22}}}x^2\right)B_{(1)}^\alpha +2x^2B_{(2)}^\alpha +\chi ^{\alpha \mu }ww_\mu +{\displaystyle \frac{1}{2}}\gamma \delta ^{1\alpha }(ww_{111}w_1w_{11}),`$
$`B_{(4)}^\alpha `$ $`=`$ $`\left(x^1+{\displaystyle \frac{\chi ^{11}}{\chi ^{12}}}x^2\right)B_{(1)}^\alpha +3x^2B_{(2)}^\alpha `$
$`+{\displaystyle \frac{1}{2}}[\chi ^{\alpha \mu }ww_\mu +\delta ^{1\alpha }(\chi ^{11}ww_1+2\chi ^{12}ww_2\gamma w_1w_{11})].`$
Table 2. Conservation laws for equations of form (58)
| # | Conservation laws |
| --- | --- |
| 1 | $`D_\alpha [\beta ^1B_{(1)}^\alpha +\beta ^2B_{(2)}^\alpha ]=0`$ |
| 2 | $`D_\alpha [\left(\beta ^1+2(\chi ^{12}/\chi ^{22})\beta ^2\right)B_{(1)}^\alpha +2\beta ^2B_{(2)}^\alpha +B_{(3)}^\alpha ]=0`$ |
| 3 | $`D_\alpha [\left(\beta ^1+3(\chi ^{11}/\chi ^{12})\beta ^2\right)B_{(1)}^\alpha +6\beta ^2B_{(2)}^\alpha +2B_{(4)}^\alpha ]=0`$ |
| 4 | $`D_\alpha B_{(1)}^\alpha =0`$, $`D_\alpha [2\beta B_{(2)}^\alpha +B_{(3)}^\alpha ]=0`$ |
| 5 | $`D_\alpha B_{(1)}^\alpha =0`$, $`D_\alpha [3\beta B_{(2)}^\alpha +B_{(4)}^\alpha ]=0`$ |
| 6 | $`D_\alpha [\beta B_{(1)}^\alpha +B_{(3)}^\alpha ]=0`$, $`D_\alpha [(\chi ^{12}/\chi ^{22})B_{(1)}^\alpha +B_{(2)}^\alpha ]=0`$ |
| 7 | $`D_\alpha [\beta B_{(1)}^\alpha +2B_{(4)}^\alpha ]=0`$, $`D_\alpha [(\chi ^{11}/\chi ^{12})B_{(1)}^\alpha +2B_{(2)}^\alpha ]=0`$ |
| 8 | $`D_\alpha B_{(1)}^\alpha =0`$, $`D_\alpha B_{(2)}^\alpha =0`$ |
| 9 | $`D_\alpha B_{(1)}^\alpha =0`$, $`D_\alpha B_{(2)}^\alpha =0`$ |
| 10 | $`D_\alpha B_{(1)}^\alpha =0`$, $`D_\alpha B_{(2)}^\alpha =0`$, $`D_\alpha B_{(3)}^\alpha =0`$ |
| 11 | $`D_\alpha B_{(1)}^\alpha =0`$, $`D_\alpha B_{(2)}^\alpha =0`$, $`D_\alpha B_{(4)}^\alpha =0`$ |
According to the general results of Section 3, each equation (58) admits conservation laws with characteristics $`Q=u(x)`$, where $`u(x)`$ is any smooth solution of the equation considered. These conservation laws are of the form (36), that is
$$D_\alpha P_{(u)}^\alpha =0,$$
the corresponding conserved currents $`P_{(u)}^\alpha `$ being given by the expression (35). Here, on account of (59), (35) simplifies and reads
$$P_{(u)}^\alpha =\chi ^{\alpha \mu }(uw_\mu u,_\mu w)+\delta ^{1\alpha }\gamma (uw_{111}+u,_{11}w_1u,_{111}wu,_1w_{11}).$$
(80)
Let us now specialize to the differential equation
$$EJw_{1111}+mw_{22}=0,$$
(81)
governing the dynamics of a classic homogeneous Bernoulli-Euler beam. Here $`EJ`$ is the bending rigidity of the beam and $`m`$ is the mass of the beam per unit length. According to the above analysis, (81) admits the following six linearly independent infinitesimal variational symmetries:
$$Y_1,Y_2,Y_3+\frac{1}{2}X_0,Y_5=\frac{}{w},Y_6=x^1\frac{}{w},Y_7=x^2\frac{}{w},$$
(82)
where $`Y_5`$, $`Y_6`$ and $`Y_7`$ are vector fields of the type $`X_u=u\left(x\right)/w`$ corresponding to the solutions $`u=1`$, $`u=x^1`$ and $`u=x^2`$ of (81), respectively. Here, the independent variables $`x^1`$ and $`x^2`$ are the spatial variable along the rod axis and the time, respectively, so that the conservation laws admitted by the smooth solutions of equation (81) may be written in the more familiar form
$$\frac{\Psi }{x^2}+\frac{P}{x^1}=0,$$
where $`\Psi `$ and $`P`$ denote the density and flux of the conservation law, respectively. The densities and fluxes of the conservation laws for (81) associated with the vector fields (82) together with their physical interpretation are presented in Table 3.
Table 3. Conservation laws for Bernoulli-Euler beams
| Generators | Conservation laws |
| --- | --- |
| $`\begin{array}{c}\text{space translations}\\ Y_1\end{array}`$ | $`\begin{array}{c}\text{wave momentum}\hfill \\ \mathrm{\Psi }_{(1)}=mw_1w_2\hfill \\ P_{(1)}=(1/2)[EJ(2w_1w_{111}w_{11}^2)mw_2^2]\hfill \end{array}`$ |
| $`\begin{array}{c}\text{time translations}\\ Y_2\end{array}`$ | $`\begin{array}{c}\text{energy}\hfill \\ \mathrm{\Psi }_{(2)}=(1/2)\left(EJw_{11}^2+mw_2^2\right)\hfill \\ P_{(2)}=EJ(w_{11}w_{12}w_2w_{111})\hfill \end{array}`$ |
| $`\begin{array}{c}\text{scaling}\\ Y_3+(1/2)X_0\end{array}`$ | $`\begin{array}{c}\mathrm{\Psi }_{(3)}=x^1\mathrm{\Psi }_{(1)}+2x^2\mathrm{\Psi }_{(2)}+mww_2\hfill \\ P_{(3)}=x^1P_{(1)}+2x^2P_{(2)}+(1/2)EJ(ww_{111}w_1w_{11})\end{array}`$ |
| $`Y_5`$ | $`\begin{array}{c}\text{linear momentum}\hfill \\ \mathrm{\Psi }_{(5)}=mw_2,P_{(5)}=EJw_{111}\hfill \end{array}`$ |
| $`Y_6`$ | $`\begin{array}{c}\text{similar to Eshelby energy-momentum tensor}\hfill \\ \mathrm{\Psi }_{(6)}=x^1mw_2,P_{(6)}=EJ(x^1w_{111}w_{11})\hfill \end{array}`$ |
| $`\begin{array}{c}\text{Galilean boost}\\ Y_7\end{array}`$ | $`\begin{array}{c}\text{center-of-mass theorem}\hfill \\ \mathrm{\Psi }_{(7)}=m(x^2w_2w),P_{(7)}=EJx^2w_{111}\hfill \end{array}`$ |
Conservation laws in the dynamics of rods are considered in many papers (see e.g. Antman, 1984; Kienzler, 1986; Chien et al., 1993; Maddocks and Dichmann, 1994; Tabarrok et al., 1994; Djondjorov, 1995). However, the particular form of the differential equations examined in the present study allows comparison with the results reported by Chien et al. (1993) and by Maddocks and Dichmann (1994) only.
Chien et al. (1993) derive conservation laws for the statics and dynamics of rods employing a technique called by the authors Neutral Action (NA) method. The conservation laws for rod equations established in this Section could be compared to their ones only for the differential equation (81) which coincides with the equation
$$Bw_{1111}+2B_{,1}w_{111}+B_{,11}w_{11}+Hw_{22}=0,$$
considered in Chien et al. (1993) when $`B=EJ`$ and $`H=m`$. The comparison shows that the conserved currents of the conservation laws for (81) with characteristics other than $`Q=u(x)`$ obtained by Chien et al. (1993) coincide with ours presented in Table 3. As for the conserved currents of the conservation laws for (81) with characteristics $`Q=u(x)`$, where $`u(x)`$ is any solution of (81), our general formula (80) implies
$$P_{(u)}^\alpha =\delta ^{2\alpha }m(uw_2u,_2w)+\delta ^{1\alpha }EJ(uw_{111}+u,_{11}w_1u,_{111}wu,_1w_{11}).$$
Only a part of these conservation laws are identified and presented in (Chien et al., 1993), namely those associated with the solutions to (81) of the form
$$u(x)=C_1(x^1)^3+C_2(x^1)^2+C_3x^1+C_4+C_5x^2,C_i=const\left(i=1,\mathrm{},5\right),$$
while, in fact, equation (81) has an infinite-dimensional space of solutions.
Five conservation laws in the dynamics of rods are reported in (Maddocks and Dichmann, 1994) within a general nonlinear direct theory. The restricted version of this theory describing small planar bending of an uniform inextensible unshearable isotropic elastic rod with a linear constitutive law, the rotatory inertia of the rod cross section being neglected, is exactly the classic Bernoulli-Euler theory for homogeneous beams whose governing equation is (81). Rewriting the conservation laws in (Maddocks and Dichmann, 1994) taking into account the aforementioned restrictions we observe that: (1) the conservation law for the total angular momentum (formula 2.14 in Maddocks and Dichmann, 1994) degenerates to the well known basic relation of Bernoulli-Euler theory $`Q=M/x^1`$ (here $`Q`$ and $`M`$ denote shear force and bending moment, respectively, see Washizu, 1982); (2) the density and flux of the conservation law associated with material isotropy (formula 4.5 in Maddocks and Dichmann, 1994) vanish identically; (3) the conservation law corresponding to material homogeneity (formula 3.2 in Maddocks and Dichmann, 1994) reduces to conservation of the wave momentum (see Table 3); (4) the expressions for the densities and fluxes of energy (formula 2.19 in Maddocks and Dichmann, 1994) and linear momentum (formula 2.12 in Maddocks and Dichmann, 1994) conservation laws coincide with the respective ones presented in Table 3. The set of conservation laws with characteristics $`Q=u(x)`$, where $`u(x)`$ is any solution of (81), as well as the conservation law associated with the variational scaling symmetry $`Y_3+(1/2)X_0`$ (see Table 3) have no analogues in (Maddocks and Dichmann, 1994).
Three interesting kinds of group-invariant solutions to certain equations of the class (58) are identified below. First of them corresponds to vector fields $`cY_1Y_2,`$ where $`c=const`$. These group-invariant solutions are travelling waves
$$w=U(s),s=x^1\pm cx^2,$$
admissible only for equations (58) with $`\kappa (x^1,x^2)=f(s)`$. The reduced equations determining such group-invariant solutions are
$$\gamma \frac{d^4U}{ds^4}+(\chi ^{11}\pm 2\chi ^{12}c+\chi ^{22}c^2)\frac{d^2U}{ds^2}+f(s)U=0.$$
The second one corresponds to the vector field $`Y_3`$ and is of the form
$$w=U(s),s=x^1(x^2)^{1/2}\frac{\chi ^{12}}{\chi ^{22}}(x^2)^{1/2}.$$
The vector field $`Y_3`$ is admitted only if $`\kappa (x^1,x^2)=(x^2)^2f(s)`$ (see cases # 2, 4, 6 and 10 in Table 1). The reduced equations for these invariant solutions are
$$4\gamma \frac{d^4U}{ds^4}+\chi ^{22}s^2\frac{d^2U}{ds^2}+3\chi ^{22}s\frac{dU}{ds}+4f(s)U=0.$$
The third kind of group-invariant solutions corresponds to the vector field $`Y_4`$:
$$w=U(s),s=2x^1(x^2)^{1/3}\frac{\chi ^{11}}{\chi ^{12}}(x^2)^{2/3}.$$
The vector field $`Y_4`$ is admitted only if $`\kappa (x^1,x^2)=(x^2)^{4/3}f(s)`$ (see cases # 3, 5, 7 and 11 in Table 1). The reduced equations for the invariant solution under consideration are
$$48\gamma \frac{d^4U}{ds^4}4\chi ^{12}s\frac{d^2U}{ds^2}4\chi ^{12}\frac{dU}{ds}+3f(s)U=0.$$
Obviously, the latter two kinds of group-invariant solutions could be reduced to self-similar solutions if $`\chi ^{12}=0`$ or $`\chi ^{11}=0`$, respectively.
6. CONCLUDING REMARKS
In this paper, Lie transformation group methods have been applied to the class of partial differential equations (1). This class is of interest for structural mechanics since the governing equations of various classical plate and rod theories belong to it; the examples given in Section 2 illustrate this fact. In the context of structural mechanics, the results of the group analysis of equations (1) give a number of attractive possibilities. Here, the established point Lie symmetries of (1) are used to construct group-invariant solution to the governing equations of several plate and rod models, to derive conservation laws revealing important features of such models and to find transformations simplifying the differential structure of equations associated with particular plate problems.
First of all, the well known computational procedure for finding the most general point Lie symmetry group has been applied to the foregoing class of equations. As a result, the system of equations (16) – (20) is derived determining the equations of the type considered that admit a larger group together with the generators of this group; naturally, all equations of this class being linear and homogeneous admit the point Lie groups generated by the vector fields (14). The system (16) – (20) allows the associated group-classification problem to be stated and examined.
In Section 4, this problem is solved for the plate equations (38) in terms of their invariants $`s_{\left(1\right)}`$, $`s_{\left(2\right)}`$ and $`s_{\left(3\right)}`$ defined by (47). The equations of form (38) with $`s_{\left(1\right)}s_{\left(2\right)}s_{\left(3\right)}0`$ are found to admit the largest symmetry groups. It is noteworthy that each such equation with variable coefficients can be transformed, using a suitable change of variables, to an equation with constant coefficients belonging to the same class. An example of such a transformation is given at the end of Section 4 where, in addition, a class of group-invariant solutions to the equation considered is presented.
The group-classification problem for the rod equations (58) is entirely solved in Section 5. All equations of that kind admitting point Lie symmetry groups, in addition to the ones generated by (14), are determined and presented in Table 1 together with the generators of the respective groups. The largest symmetry groups are admitted by the equations of form (58) whose coefficients are such that $`\chi ^{22}det(\chi ^{\alpha \beta })=0`$, $`\kappa (x)0`$. The most interesting group-invariant solutions for equations (58) are identified and the corresponding reduced equations are presented at the end of Section 5.
Once the ”ordinary” point Lie symmetries of an equation of form (1) are determined, one can easily find, using the general criterion (26), which of them are variational symmetries of this equation. Then, (35) and (37) provide explicit expressions for the conserved currents of the conservation laws associated through Noether’s theorem with the established variational symmetries. These expressions will involve derivatives of the dependent variable of lowest possible order, which is important in view of their application in structural mechanics. The reciprocity relation valid for each equation of form (1) is given explicitly by formula (33).
In Section 4, it is shown, using the consequence (27) of the general criterion (26), that each point Lie symmetry of a plate equation of form (38) generated by a vector field of form (46) is variational symmetry of this equation. Therefor, each such symmetry gives rise to a conservation law with characteristic $`Q=(1/2)\xi _{,\mu }^\mu ww_\mu \xi ^\mu `$ and conserved current given by (37) admitted by the smooth solutions of the respective equation.
The conservation laws for the rod equations listed in Table 1 are given in Table 2. Inspecting these results one can see that the equations for unsupported rods and rods on Winkler foundations admit two independent conservation laws associated with the wave momentum ($`D_\alpha B_{(1)}^\alpha =0`$) and energy ($`D_\alpha B_{(2)}^\alpha =0`$). Equations (9) and (10) governing the stability of unsupported axially compressed beams and fluid conveying pipes belong to this class. Rod equations with $`\kappa (x)=0`$ and $`det(\chi ^{\alpha \beta })=0`$ admit a supplementary conservation law $`D_\alpha B_{(3)}^\alpha =0`$ associated with the infinitesimal scaling symmetry $`Y_3`$. Such an equation is (81) governing the vibration of the classic Bernoulli-Euler beam. Table 3 contains several physically important conservation laws for this equation. A comparison between the conservation laws derived here for equation (81) and the relevant results in (Chien et al., 1993) and (Maddocks and Dichmann, 1994) is presented in Section 5.
Acknowledgements—This research was supported by Contract MM 517/1995 with the NSF, Bulgaria.
REFERENCES
Antman, S.S., 1984. The theory of rods. In: Mechanics of solids, Vol. II. Springer-Verlag, Berlin, pp. 641-703.
Bluman, G.W., Kumei, S., 1989. Symmetries and Differential Equations. Springer-Verlag, New York.
Chien, N., Honein, T., Herrmann, G., 1993. Conservation laws for nonhomogeneous Bernoulli-Euler beams. International Journal of Solids and Structures 30 (23), 3321-3335.
Djondjorov, P., 1995. Invariant properties of Timoshenko beams equations. International Journal of Engineering Science 33 (14), 2103-2114.
Gregory, R.W., Paidoussis, M.P., 1966. Unstable oscillation of tubular cantilevers conveying fluid. I: Theory. Proceedings of Royal Society London A–293, 512-527.
Ibragimov, N.H., 1985. Transformation Groups Applied to Mathematical Physics. Reidel, Boston.
Kienzler, R., 1986. On existence and completeness of conservation laws associated with elementary beam theory. International Journal of Solids and Structures 22 (7), 789-796.
Maddocks, J., Dichmann, D., 1994. Conservation laws in the dynamics of rods. Journal of Elasticity 34, 83-96.
Olver, P.J., 1993. Applications of Lie Groups to Differential Equations, Second Edition, Graduate Texts in Mathematics, Vol. 107. Springer-Verlag, New York.
Olver, P.J., 1995. Equivalence, Invariance and Symmetry. Cambridge University Press, Cambridge.
Ovsiannikov, L.V., 1972. Group properties of the equations of mechanics. In: Mechanics of Continuous Media and Relevant Problems of Analysis. Nauka, Moscow, pp. 381-393.
Ovsiannikov, L.V., 1982. Group Analysis of Differential Equations. Academic Press, New York.
Smith, T.E., Herrmann, G., 1972. Stability of a beam on an elastic foundation subjected to follower forces. Journal of Sound and Vibration 39, 628-629.
Tabarrok, B., Tezer, C., Stylianou, M., 1994. A note on conservation principles in classical mechanics. Acta Mechanica 107, 137-152.
Vassilev, V., 1988. Group properties of a class of fourth-order partial differential equations. Annals of the University of Sofia 82 (Part II - Mechanics), 163-178.
Vassilev, V., 1991. Group analysis of a class of equations of the plate and shell theory. Ph.D. thesis, Institute of Mechanics and Biomechanics, Bulgarian Academy of Sciences, Sofia.
Vassilev, V., 1997. Application of Lie groups to the theory of shells and rods. Nonlinear Analysis 30 (8), 4839-4848.
Washizu, K., 1982. Variational Methods in Elasticity and Plasticity. Pergamon Press, Oxford.
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# The finite temperature QCD phase transition with domain wall fermions
## I Introduction
Many of the properties of low energy QCD are a direct consequence of the breaking of chiral symmetry by the QCD vacuum. It is expected that this spontaneous chiral symmetry breaking will disappear as the temperature is increased. Both the nature of this symmetry restoration (abrupt phase transition or continuous cross-over) and the character of the high-temperature quark-gluon plasma phase are active areas of both theoretical and experimental research.
An especially promising approach to the theoretical study of equilibrium properties of both the QCD phase transition and the high-temperature plasma phase is direct numerical simulation of the Feynman path integral using the methods of lattice gauge theory. The quantum partition function is written as a Euclidean path integral that can be studied ab initio using the discrete, lattice formulation of Wilson. While the local color gauge symmetry of the theory remains exact at any lattice spacing in Wilson’s formulation, much of the theory’s flavor symmetry, and especially its chiral component, is explicitly broken.
This difficulty in representing the continuum flavor symmetries in a lattice fermion formulation is a serious problem that has persisted for more than two decades. When the fermion action is naively discretized the low-energy fermionic degrees of freedom increase by a factor of $`2^4`$. This well-known “doubling” problem can only be remedied by methods that explicitly break the chiral flavor symmetries for finite lattice spacing . The chiral symmetries are then recovered together with the Lorentz symmetry as the lattice spacing is sent to zero. The most popular of these methods are staggered and Wilson fermions.
Although, in principle these methods should be able to approximate the continuum theory in a controlled way, in practice this problem has been a formidable obstacle to lattice studies of the QCD phase transition. For example, the Wilson fermion formulation explicitly breaks all of the continuum chiral symmetries making phenomena driven by the spontaneous breakdown of chiral symmetry difficult to study. While staggered fermions do possess a one-dimensional continuous chiral symmetry at finite lattice spacing, this formulation explicitly breaks the vector flavor symmetry so instead of three light Goldstone pions with mass on the order of the critical temperature $`T_c160`$ MeV as found in Nature, present staggered simulations have masses for two of the three pions in the range 500-600 GeV, certainly too large.
In addition, the subtle effects of the continuum axial anomaly which are closely connected with the order of the transition are badly mutilated by both fermion formalisms at finite lattice spacing. While the anomalous $`U_A(1)`$ continuum chiral symmetry is explicitly broken by both formalisms, the fermion zero modes required by Atiyah-Singer index theorem are shifted away from zero by finite lattice spacing effects.
In principle, each of these difficulties can be addressed by simply working at smaller lattice spacing. However, present numerical methods scale very poorly as the lattice spacing is decreased, with the required numerical effort growing as $`1/a^{810}`$ for lattice spacing $`a`$.
Domain wall fermions (DWF) offer a new approach to the problem of including fermions in lattice gauge theory calculations. In this formulation, introduced by Kaplan , the fermionic fields are defined on a five-dimensional hyper-cubic lattice using a local action. The fifth direction can be thought of as an extra space-time dimension or as a new internal flavor space. The gauge fields are represented in the standard way in four dimensional space-time and are coupled to the extra fermion degrees of freedom in a diagonal fashion.
In this paper, we use a variant of Kaplan’s approach, developed by Shamir, in which the two four-dimensional faces orthogonal to the new fifth dimension are treated differently, with free boundary conditions imposed on the fermion fields. This key ingredient allows a system made up of naively massive fermions to develop chiral surface states on these boundaries (domain walls) with the positive chirality states bound exponentially to one wall and the negative chirality states bound to the other.
The two chiralities overlap only by an amount that is exponentially small in $`L_s`$, the number of lattice sites along the fifth direction. The resulting mixed state forms a Dirac 4-spinor that propagates in the four-dimensional space-time with an exponentially small mass. Therefore, the amount of chiral symmetry breaking that is artificially induced by the regulator can be controlled by the new parameter $`L_s`$. In the $`L_s\mathrm{}`$ limit the chiral symmetry is exact even at finite lattice spacing. Thus, the domain wall fermion method has succeeded in disentangling the chiral limit ($`L_s\mathrm{}`$) and the continuum limit ($`a0`$). Furthermore, the direct computing requirement grows only linearly with $`L_s`$.
Here we report the first full QCD simulations using domain wall fermions in four dimensions. The properties and parameter space of domain wall fermions appropriate for a study of QCD thermodynamics are explored in detail. Small lattices of size $`8^3\times 4`$ were used to perform numerical simulations of full, two-flavor QCD at finite temperature. Preliminary results of this work have appeared in . These studies have been carried out using the QCDSP supercomputer at Columbia. Based on the work reported here, results of physical interest have been obtained on larger lattices for a variety of observables. Preliminary results of these studies can be found in and will be presented in follow-on papers.
For a detailed introduction to the subject and relevant references the reader is referred to Refs. , and the reviews in Refs. . Earlier numerical work using domain wall fermions has explored the parameter space of a QCD-like, dynamical vector theory in two dimensions, the two flavor Schwinger model. For applications to quenched QCD see Refs. for applications to four-Fermi models see Ref. and for possible alternatives to domain wall fermion simulations see Refs. .
In Section II the action of the theory and a brief description of the numerical methods are presented. In Section III some important analytical facts are outlined in order to help guide the numerical investigation. In Section IV we study the chiral properties of the theory both below and above the chiral phase transition. Our numerical results suggest that domain wall fermions are able to sustain the desired chiral properties of QCD, even at finite lattice spacing. Both a low temperature phase where the $`SU(2)\times SU(2)`$ chiral symmetry is broken spontaneously to an $`SU(2)`$ vector symmetry and a high temperature phase where the full $`SU(2)\times SU(2)`$ chiral symmetry is intact can be recognized.
In Section IV the dependence on the two new regulator parameters, the number of sites in the fifth direction $`L_s`$, and the domain wall “height” $`m_0`$, is studied numerically. Finally, in Section VI conclusions and outlook are presented. Appendix A gives the explicit form of the gamma matrices used in this work while Appendix B describes the molecular dynamics equations of motion. Tables summarizing the numerical results are given at the end of the paper.
## II Hybrid Monte Carlo with domain wall fermions
In this section the action of QCD with domain wall fermions, its implementation for the Hybrid Monte Carlo (HMC) algorithm, and the parameters used in the simulations are described. In the following, we discuss the case of two degenerate flavors implemented using the HMC $`\mathrm{\Phi }`$ algorithm. (An odd number of flavors can be simulated using the HMC $`R`$ algorithm).
Domain wall fermions can be used in numerical simulations in a fashion similar to traditional Wilson fermions. In fact, if the fifth direction is thought of as an internal flavor direction then an HMC simulation with DWF is identical to a simulation of many flavors of Wilson fermions with a sophisticated mass matrix. We use the partition function of QCD with domain wall fermions proposed in but with a slightly different heavy flavor subtraction as in . In particular:
$$Z=[dU][d\overline{\mathrm{\Psi }}d\mathrm{\Psi }][d\mathrm{\Phi }_{PV}^{}d\mathrm{\Phi }_{PV}]e^S$$
(1)
$`U_\mu (x)`$ is the gauge field, $`\mathrm{\Psi }(x,s)`$ is the fermion field and $`\mathrm{\Phi }_{PV}(x,s)`$ is a bosonic, Pauli-Villars field. The variable $`x`$ specifies the coordinates in the four-dimensional space-time box with extent $`L`$ along each of the spatial directions and extent $`N_t`$ along the time direction while $`s=0,1,\mathrm{},L_s1`$ is the coordinate of the fifth direction, with $`L_s`$ assumed to be even. The action $`S`$ is given by:
$$S=S_G(U)+S_F(\overline{\mathrm{\Psi }},\mathrm{\Psi },U)+S_{PV}(\mathrm{\Phi }_{PV}^{},\mathrm{\Phi }_{PV},U)$$
(2)
where:
$$S_G=\beta \underset{p}{}(1\frac{1}{3}\mathrm{ReTr}[U_p])$$
(3)
is the standard plaquette action, $`\beta =6/g_0^2`$ and $`g_0`$ is the lattice gauge coupling. The fermion action for two flavors is:
$$S_F=\underset{x,x^{},s,s^{},f}{}\overline{\mathrm{\Psi }}_f(x,s)D_F(x,s;x^{},s^{})\mathrm{\Psi }_f(x^{},s^{})$$
(4)
with flavor index $`f=1,2`$ and Dirac operator:
$$D_F(x,s;x^{},s^{})=\delta _{s,s^{}}D^{}(x,x^{})+D^{}(s,s^{})\delta _{x,x^{}}$$
(5)
$`D^{}(x,x^{})`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mu =1}{\overset{4}{}}}\left[(1\gamma _\mu )U_\mu (x)\delta _{x+\widehat{\mu },x^{}}+(1+\gamma _\mu )U_\mu ^{}(x^{})\delta _{x\widehat{\mu },x^{}}\right]`$ (6)
$`+`$ $`(m_04)\delta _{x,x^{}}`$ (7)
$`D^{}(s,s^{})`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[(1\gamma _5)\delta _{s+1,s^{}}+(1+\gamma _5)\delta _{s1,s^{}}2\delta _{s,s^{}}\right]`$ (8)
$``$ $`{\displaystyle \frac{m_f}{2}}\left[(1\gamma _5)\delta _{L_s1,s}\delta _{s^{},0}+(1+\gamma _5)\delta _{s,0}\delta _{L_s1,s^{}}\right]`$ (9)
Here, $`s`$ and $`s^{}`$ lie in the range $`0s,s^{}L_s1`$. In the above equations $`m_0`$ is a five-dimensional mass representing the height of the domain wall in Kaplan’s original language. In order for the doubler species to be removed in the free theory one must choose $`0<m_0<2`$. The parameter $`m_f`$ explicitly mixes the two chiralities and, as a result, controls the bare fermion mass of the four-dimensional effective theory.
While the DWF Dirac operator defined above is not hermitian, it does obey the identity :
$$\gamma _5R_5D_F\gamma _5R_5=D_F^{}$$
(10)
with $`R_5`$ the reflection operator along the fifth direction. As a result the single-flavor Dirac determinant is real: $`detD_{F}^{}{}_{}{}^{}=detD_F^{}=det\gamma _5R_5D_F\gamma _5R_5=detD_F`$ and the two-flavor determinant which follows from integrating out the fermions in Eq. 1, $`detD_{F}^{}{}_{}{}^{2}`$, is positive. The gamma matrices used in this work are given in Appendix A. Also notice that $`D_F`$ is the same as the $`D_F^{}`$ of .
The Pauli-Villars action is designed to cancel the contribution of the heavy fermions in the large $`L_s`$ limit. Normally, such heavy fermions decouple from low energy physics and can be safely ignored. However, in the present situation the number of heavy fermions grows proportional to $`L_s`$ and can potentially overwhelm the effects of the fixed number of low energy degrees of freedom of interest. Specifically this difficulty will arise for the order of limits for which DWF are intended: first $`L_s\mathrm{}`$ followed by $`a0`$. .
There is some flexibility in the definition of the Pauli-Villars action since different actions can easily have the same $`L_s\mathrm{}`$ limit. However, the choice of the Pauli-Villars action may affect the approach to the $`L_s\mathrm{}`$ limit. A slightly different action than that proposed by Furman and Shamir is used here. This action is easier to implement numerically and, even for finite $`L_s`$, it exactly cancels the fermion action when $`m_f=1`$ resulting in a pure gauge theory. For two fermion flavors, the Pauli-Villars action we use is:
$$S_{PV}=\underset{x,x^{},s,s^{}}{}\mathrm{\Phi }_{PV}^{}(x,s)M_F(x,s;x^{},s^{})|_{m_f=1}\mathrm{\Phi }_{PV}(x^{},s^{})$$
(11)
where $`M_F=D_F^{}D_F`$.
The traditional HMC $`\mathrm{\Phi }`$ algorithm was constructed directly from the action of Eq. 2. In order to improve performance a standard even-odd preconditioning of the Dirac operator $`D_F`$ was employed. The even-odd preconditioning was done on the five dimensional space. All necessary matrix inversions were done using a standard conjugate gradient (CG) algorithm. As expected the even-odd preconditioning resulted in a reduction of the required number of conjugate gradient iterations and a consequent speed-up of a factor of approximately two.
The only new ingredient in our HMC algorithm is the appearance of the bosonic Pauli-Villars fields. The probability distribution of these fields is generated with a heat bath step at the beginning of each HMC “trajectory”: a field of Gaussian random numbers is generated with distribution $`e^{\eta _{PV}^{}\eta _{PV}}`$ and from it the Pauli-Villars fields $`\mathrm{\Phi }_{PV}(x,s)`$ are obtained by $`\mathrm{\Phi }_{PV}=[D_F(m_f=1)]^1\eta _{PV}`$ using the CG algorithm.
Since the Pauli-Villars action in Eq. 11 is polynomial in the domain wall operator $`D_F`$, its gradient with respect to the gauge fields, needed to evolve the gauge degrees of freedom, can be computed without performing any Dirac inversions. This contrasts favorably with the fermion contribution to the gauge force which requires one inversion per molecular dynamics step. As a result, the relative computational cost involved in calculating the Pauli-Villars force is negligible. Furthermore, because the Pauli-Villars fields are bosonic their molecular dynamics force term enters with an opposite sign that of the fermion force, resulting in a large, approximate cancellation. Because of this cancelation the HMC force term is approximately independent of $`L_s`$ and it is not necessary to decrease the HMC step size as $`L_s`$ is increased.
In the approach described above the presence of the Pauli-Villars fields increases the memory requirement. However, it should be noted that there is an alternative approach that does not involve Pauli-Villars fields. To see this consider the result after integration over both the Pauli-Villars and fermion fields. It is $`detM_F(m_f)/detM_F(m_f=1)=det[M_F(m_f)/M_F(m_f=1)]`$. Therefore, one could simulate the same action without Pauli-Villars fields by simply using as the fermion matrix $`M_F(m_f)/M_F(m_f=1)`$. Inversion of this matrix will involve inversion of $`M_F(m_f)`$ using the CG algorithm as in the previous method while the final result would have to be multiplied by the matrix $`M_F(m_f=1)`$. If, for example, the CG algorithm required $`100`$ iterations to converge, this extra matrix multiplication will increase the computing cost by only $`1\%`$. The only disadvantage of this approach is that the equations of motion become slightly more complicated.
Since this work is the first to implement DWF in dynamical QCD the approach with Pauli-Villars fields was adopted for simplicity and because it has been proven reliable in numerical simulations of the Schwinger model. For the convenience of the reader the molecular dynamics equations of motion with Pauli-Villars fields and an even-odd preconditioned DWF Dirac operator are given in Appendix B.
Fermionic Green’s functions were computed using the method described in Ref. . Standard fermion fields in the four-dimensional space–time are constructed from the five-dimensional fermion fields using the projection prescription:
$`\psi (x)`$ $`=`$ $`P_\mathrm{L}\mathrm{\Psi }(x,0)+P_\mathrm{R}\mathrm{\Psi }(x,L_s1)`$ (12)
$`\overline{\psi }(x)`$ $`=`$ $`\overline{\mathrm{\Psi }}(x,L_s1)P_\mathrm{L}+\overline{\mathrm{\Psi }}(x,0)P_\mathrm{R}`$ (13)
where $`P_{\mathrm{R}/\mathrm{L}}=\frac{1}{2}(1\pm \gamma ^5)`$. This somewhat arbitrary choice defines operators which should have a large overlap with the physical low energy fermion modes bound to the $`s=0`$ and $`s=L_s1`$ walls. The right- and left-handed components found on opposite walls are combined to assemble the desired physical 4-spinors.
Since these are the first simulations of DWF in dynamical QCD there are no previous results that would allow an independent check of the methods and code. Tests using the chiral condensate from the free field analytical results of were done in order to check the Dirac operator and inverter. The subtraction of Pauli-Villars fields was tested by performing simulations with $`m_f=1`$ and comparing with equivalent results from quenched simulations. Finally, two flavor dynamical simulations were done on $`2^4`$ lattices and the results were compared with simulations using the overlap formalism relevant for the DWF action for the same parameters. In particular for $`\beta =5.6`$, $`m_f=0.1`$, $`m_0=0.9`$ the overlap simulation gave $`\overline{\psi }\psi =1.672(2)\times 10^3`$ and average plaquette $`\mathrm{plaq}=5.765(79)\times 10^1`$ while the DWF simulation with $`L_s=18`$ gave $`\overline{\psi }\psi =1.653(33)\times 10^3`$ and average plaquette $`\mathrm{plaq}=5.841(47)\times 10^1`$.
All numerical results in this work were obtained from lattices of size $`L=8`$, $`N_t=4`$ with periodic spatial boundary conditions and anti-periodic temporal boundary conditions. The fifth direction was set to various values in the range $`[8,40]`$, the domain wall height $`m_0`$ was varied in the range $`[1.15,2.4]`$, the fermion mass was varied in the range $`[0.02,0.18]`$ and $`\beta `$ was varied in the range $`[4.65,5.95]`$. The molecular dynamics trajectory length was set to $`\tau =0.5`$ and the step size $`\delta \tau `$ was set to various values in the range $`[0.0078,0.02]`$ depending on the values of the other parameters. The CG stopping condition which is defined as the ratio of the norm of the residual vector over the norm of the source was set to $`10^6`$. This resulted in an average number of CG iterations ranging between 50 and 400 depending on the values of the other parameters.
The initial configuration was generally chosen to be in the phase opposite to that expected for the input parameters creating a very visible thermalization process in which the system should be seen to evolve into the correct phase. Typically $`100400`$ trajectories were needed to thermalize the lattice. The chiral condensate and Wilson line were measured in every sweep. The chiral condensate was measured using a standard “one-hit” stochastic estimator of the trace of $`D_F^1`$ with spin and $`s`$ coordinates restricted according to Eq. 13. Specifically we evaluated the quantities:
$`|W|`$ $`=`$ $`{\displaystyle \frac{1}{3L^3}}\left|{\displaystyle \underset{\stackrel{}{x}}{}}\mathrm{tr}[{\displaystyle \underset{lL(\stackrel{}{x})}{}}U_l]\right|`$ (14)
$`\overline{\psi }\psi `$ $`=`$ $`{\displaystyle \frac{1}{12L^3N_t}}\{\mathrm{tr}[s=0|1/D_F|s=L_s1P_\mathrm{R}]`$ (16)
$`+\mathrm{tr}[s=L_s1|1/D_F|s=0P_\mathrm{L}]\}.`$
Here $`U_l`$ identifies the $`SU(3)`$ gauge matrix corresponding to the link $`l`$ and the ordered product is taken for all links in the time-like line $`L(\stackrel{}{x})`$ with spatial coordinate $`\stackrel{}{x}`$. The somewhat unconventional normalization in Eq. 16 was used in our previous work and determines a spin and color average which for very large mass $`m_f`$ approaches $`1/m_f`$. (Note, here $`D_F`$ is the single-flavor Dirac operator defined in Eq. 5.)
## III Analytical considerations
In this section we summarize some of the analytically determined properties of domain wall fermions. These help guide our numerical investigations, which are done for finite and non-zero values for the three parameters of domain wall fermions, $`L_s`$, $`m_0`$, and $`m_f`$, as well as at finite bare coupling $`g_0`$.
### A $`L_s`$ dependence
For numerical simulations, the existence of the chiral limit for domain wall fermions and the rate of approach to it are of primary importance. The computational requirements for domain wall fermions grow as one power of $`L_s`$ from the simple increase in the number of operations. An additional slight increase in computational cost for larger $`L_s`$ comes from the decrease in the total quark mass due to smaller mixing between the chiral surface states, until the quark mass is dominated by the input $`m_f`$.
The axial Ward-Takahashi identities for domain wall fermions are the same as the continuum, except for an additional term which comes from the mixing of the left- and right-handed light surface states at the midpoint of the fifth dimension, $`L_s/2`$. At any lattice spacing this additional term vanishes as $`L_s\mathrm{}`$ for non-singlet axial symmetries . For the singlet axial symmetry, this extra term generates the axial anomaly. At strong coupling, the axial currents are conserved for $`L_s\mathrm{}`$ but, since the doubler fermions may enter the spectrum, these currents may not have the physical significance of axial currents.
For free domain wall fermions, the rate of approach to the chiral limit can be calculated. At finite $`L_s`$ the mixing of the chiral components is reflected in the fermion mass $`m_{\mathrm{eff}}`$. For the one flavor theory this effective mass is
$$m_{\mathrm{eff}}=m_0(2m_0)\left[m_f+(1m_0)^{L_s}\right]0<m_0<2,$$
(17)
$`m_{\mathrm{eff}}`$ has two pieces: one is proportional to the bare mass $`m_f`$ and the other expresses the residual mixing between the chiral modes bound to the domain walls. Since each bound chiral state decays exponentially with the distance from its wall, the residual mixing between them vanishes exponentially with $`L_s`$, with a decay constant of $`\mathrm{ln}|1m_0|`$. Notice that when $`L_s\mathrm{}`$, $`m_0`$ becomes an irrelevant parameter, provided it stays in the range (0,2).
In the free theory, one also finds that fermion states with non-zero four-momentum decay more slowly with the distance from the wall than do zero momentum states. The decay is controlled by the four-momentum and the value for $`m_0`$. Since the lattice momentum $`p_L^\mu =p^\mu a`$, where $`a`$ is the lattice spacing, the slower decay for modes with non-zero four-momentum is an $`O(a^2)`$ effect which should vanish in the continuum limit. In addition, for a given $`m_0`$, there is a critical four-momentum above which the fermions are no longer bound to the wall, but instead behave like massive, five-dimensional fermions. Of course, because these fermions are massive, they necessarily preserve the theory’s four-dimensional chiral symmetry since their propagation between the $`s=0`$ and $`s=L_s1`$ walls is exponentially suppressed.
For interacting theories, a simple expectation is for Eq. 17 to be replaced by
$$m_{\mathrm{eff}}=Z_m\left[m_f+ce^{\alpha L_s}\right].$$
(18)
The exponential dependence is seen perturbatively and proven to exist non-perturbatively, provided the gauge fields satisfy a smoothness condition. These analytic results support the expectation of exponential suppression of chiral symmetry breaking effects in the non-perturbative regime. However, this behavior may be best established by the sort of explicit numerical study reported here. Generally $`\alpha `$ should depend on $`m_0`$, allowing one to choose an optimal value for simulations at finite $`L_s`$. While in the free theory $`m_0=1`$ gives $`e^\alpha =0`$, for the interacting theory the variable character of fermion propagation in fluctuating background gauge fields makes decoupling the walls with a single value for $`m_0`$ unlikely, except at very weak coupling.
Close to the continuum limit, it can be argued that this form for the effective mass, an input quark mass plus a residual mass $`m_{\mathrm{res}}`$, should enter all long-distance observables. However, away from the continuum limit or for quantities that cannot be obtained from a low energy effective QCD Lagrangian this is not necessarily the case. Therefore, different observables may approach their $`L_s\mathrm{}`$ limit in different ways, depending on the momentum scales which enter the observable, and the corrections to the input quark mass, particularly at stronger couplings, may be more complicated. In a numerical investigation this has to be kept in mind. In this paper only the chiral condensate and pion susceptibility are considered. Work on larger lattices involving measurements of many fermionic operators is currently in progress.
Numerical simulations may well be the only way to determine the dependence of chiral symmetry breaking effects on $`L_s`$ for intermediate lattice spacings ($`1`$ to 3 GeV). While for full QCD, perturbative and non-perturbative arguments support exponential falloff with $`L_s`$, for quenched theories, where the lack of damping from a fermionic determinant can lead to configurations with unsuppressed small eigenvalues for the fermions, the large $`L_s`$ behavior is even more in need of determination through simulations. Some results from quenched QCD simulations have been discussed in Refs. .
### B $`m_0`$ dependence
For free domain wall fermions the number of light flavors is controlled by the value of $`m_0`$ . In particular $`m_0<0`$ corresponds to zero light flavors, $`0<m_0<2`$ to one, $`2<m_0<4`$ to four, and $`4<m_0<6`$ to six light flavors. The theory is symmetric under $`m_010m_0`$.
For the interacting theory the values of $`m_0`$ which distinguish between different numbers of flavors are changed. Light fermions first appear for $`m_0>0`$, the one to four flavor transition occurs for $`m_0>2`$, etc. and the theory is still symmetric about $`m_010m_0`$. This is expected perturbatively and seen numerically . There is also some numerical evidence that the transition between different numbers of flavors is smooth and spread out over a small region of $`m_0`$ . For the interacting theory, keeping $`m_0<2`$ guarantees that a theory with not more than one flavor is being studied.
While $`m_0`$ is an irrelevant parameter for $`L_s\mathrm{}`$, it is very important for simulations, not only in controlling the approach to the chiral limit and the flavor content of the theory, but also for insuring that light fermions with an average momentum given by the temperature are still bound to the walls. For the free theory, the range of four-momenta carried by states that are bound to the walls increases as $`m_0`$ increases from zero, as do the corresponding Dirac eigenvalues. As $`m_0`$ approaches one, the largest Dirac eigenvalues of these “bound” states become farther off-shell, with values $`1/a`$. As $`m_0`$ increases above 1, the number of these off-shell states continues to grow but rather than their eigenvalues increasing, instead their degeneracy increases beyond what would be seen for the large momentum states of a free theory. As $`m_0`$ increases further and approaches 2, some of these excess, degenerate states become more nearly on-shell until for $`m_0>2`$ one has the low-lying Dirac eigenvalues of a free, four-flavor theory. Thus, in the free case a choice of $`m_0`$ midway between 0 and 2 is best, giving the largest phase space for physical states bound to the walls, without adding additional flavors. Using this behavior as a guide for the interacting case, one expects that choosing $`m_0`$ midway between the value where a single light fermion is bound to the walls and four light fermions are bound allows the largest range of four-momentum for a single flavor of light quark bound to the walls.
### C Topology
An important property of the domain wall fermion Dirac operator is the presence of exact zero modes in the $`L_s=\mathrm{}`$ limit, as can be seen from the overlap formalism . These zero modes are related to the topological charge of the gauge field and as a result an approximate form of the index theorem is present on the lattice . Studies on semiclassical configurations show the presence of modes which are very close to zero modes even at finite $`L_s`$ and as a result make lattice studies of anomalous symmetry breaking possible
During simulations, field configurations of different winding number should show zero mode effects in fermionic observables. The efficiency with which the hybrid Monte Carlo can move the system between sectors of different winding is an important question, as are the long correlations along the fifth direction which develop for gauge field configurations where the topology is changing. These issues have been studied in numerical simulations of the dynamical Schwinger model where the hybrid Monte Carlo performed well and topology changing occurred. For this exploratory study of full QCD thermodynamics, the input quark masses are not small, so the effects of topology should not be particularly large.
## IV The finite temperature QCD phase transition
The previous sections have described the domain wall fermion formulation and important questions about it that need to be investigated numerically. Here we report on simulations of full QCD at finite temperature with domain wall fermions on $`8^3\times 4`$ lattices. Studying this system allows us to investigate domain wall fermions for full QCD and look for the presence of chirally broken and symmetric phases. The small volume makes scanning over many values for $`L_s`$, $`m_0`$, $`m_f`$ and $`g_0`$ possible, laying the foundation for more realistic simulations on larger volume.
Since the finite temperature transition of QCD is controlled by the chiral symmetries of the theory (for light quarks), using domain wall fermions to preserve the full global symmetries of the continuum should remove one systematic lattice error that is difficult to control. However, finite temperature simulations are generally only possible on relatively coarse lattices ($`a^1700`$ MeV for a lattice with $`N_t=4`$), where analytic results about domain wall fermions are lacking. The light chiral modes of domain wall fermions at weak coupling must exist at $`a^1700`$ MeV, in the full non-perturbative gauge field backgrounds, for thermodynamic simulations to be possible. If it is found that chiral modes exist on coarse lattices, the size of the $`m_{\mathrm{res}}(L_s,\beta )`$ and its dependence on $`L_s`$ and $`\beta `$ must be investigated. (As already mentioned, $`m_{\mathrm{res}}`$ is only a sensible quantity for low-energy observables and it must be demonstrated that various determinations of it are consistent. In this section we refer to $`m_{\mathrm{res}}`$, without specifying precisely how it may be determined, as a generic indicator of the mixing between the chiral modes.)
### A Locating the transition
Locating the phase transition in full QCD requires scanning values for four parameters ($`m_0`$, $`L_s`$, $`m_f`$, and $`\beta `$). Without any knowledge of the location of the transition, or if it even exists for domain wall fermions, choosing parameters for initial simulations is difficult. For staggered fermions, the critical coupling for the finite temperature phase transition for 2 flavors on an $`N_t=4`$ lattice is $`\beta _c=5.265`$ for $`m=0.01`$ and $`\beta _c=5.291`$ for $`m=0.025`$ . Since staggered and domain wall fermions both have their chiral limit at zero quark mass, the light quarks have the largest effect in the location of $`\beta _c`$ and both theories have the same number of light flavors, we used the staggered values as a rough guide.
Our first simulations with domain wall fermions were done at $`\beta =5.0`$ and $`\beta =5.4`$, with the hope that these would be above and below the transition region. $`L_s=8`$ and $`m_f=0.1`$ were chosen to keep the computational difficulty modest. We worked with $`m_0=1.65`$, since for quenched simulations this choice gave a reasonable falloff between the walls at $`\beta =5.7`$ and for quenched QCD, $`\beta =5.7`$ is close to $`\beta _c`$ for an $`N_t=4`$ lattice.
Although with this choice of $`m_0`$, the $`\beta `$ range being examined ($`5.0\beta 5.4`$) lies below the chiral transition, we describe this point first since it demonstrates our very first efforts in charting this parameter space and the difficulties we encountered. The evolution of $`\overline{\psi }\psi `$ for $`m_f=0.1`$ and $`\beta =5.0`$ is shown in the upper panel of Figure 1 and the lower panel is for $`\beta =5.4`$. The hybrid Monte Carlo was run with a step size of $`\delta \tau =0.025`$ and 20 steps per trajectory, giving an acceptance of 66% for $`\beta =5.0`$ and 70% for $`\beta =5.4`$. The evolution appears quite generic and the simulation presented no difficulty to the hybrid Monte Carlo. For $`\beta =5.0`$ the Wilson line expectation value was 0.0223(15) and for $`\beta =5.4`$ it was 0.0466(41). Both these values are small and indicate that both $`\beta `$ values correspond to the confined phase.
The chiral condensate $`\overline{\psi }\psi `$ was also measured for a variety of valence masses. In quenched QCD at zero temperature, extrapolations of $`\overline{\psi }\psi `$ to $`m_f=0`$ using quark masses from $`m_f=0.02`$ to $`m_f=0.1`$ were used to see that chiral modes existed for a particular $`m_0`$. The limit $`\overline{\psi }\psi (m_f0)`$ could only be non-zero if light chiral modes were present, provided $`L_s`$ is large enough that the residual mixing is unimportant. (For the current finite temperature case, $`\overline{\psi }\psi (m_f0)`$ can be zero either from the absence of chiral modes or because the system is in the symmetry-restored phase.)
Figure 2 shows that $`\overline{\psi }\psi `$ extrapolates to a non-zero value for both $`\beta =5.0`$ and $`\beta =5.4`$ and this value is not very sensitive to $`\beta `$. The values for the Wilson line indicate both $`\beta `$ values are in the confined phase, so the $`\overline{\psi }\psi `$ results show that light chiral modes are present with an unknown residual mixing. The insensitivity to $`\beta `$ is an interesting feature.
Next, instead of scanning larger values of $`\beta `$, we decided to change $`m_0`$ from 1.65 to 1.9, keeping all other parameters identical. (This reflects our initial search path in parameter space and does not imply the absence of a transition at $`m_0=1.65`$ and $`L_s=8`$.) The acceptance is 59% for $`\beta =5.0`$ and 71% for $`\beta =5.4`$ The Wilson line for $`\beta =5.0`$ is 0.030(2), while for $`\beta =5.4`$ it is 0.202(5), indicating that $`\beta =5.4`$ is likely deconfined. The evolutions show a very different behavior for the condensate evaluated at the dynamical quark mass. The value at $`\beta =5.0`$ has increased, part of which likely reflects the change with $`m_0`$ in the overlap between the five-dimensional light modes and the surfaces at $`s=0`$ and $`L_s1`$. The $`\beta =5.4`$ values are much smaller, consistent with the deconfined phase.
Figure 4 shows the valence quark extrapolation. The small value of $`\overline{\psi }\psi (m_f0)`$ suggests the restoration of chiral symmetry. Of course, there is a possibility that this small value might instead be caused by the loss of chiral modes. However, this is unlikely because we have seen that chiral modes do exist for $`\beta =5.0`$ and one expects that at the weaker $`\beta =5.4`$ coupling these chiral modes should be even more numerous. Therefore, we have preliminary evidence for two phases of full QCD with dynamical domain wall fermions.
To solidify the evidence for two different phases of QCD with domain wall fermions further simulations for $`L_s=8`$ and $`m_0=1.9`$ were done with dynamical quark masses of 0.14 and 0.18. These points are shown in Figure 5. The dashed line is the fit to the quenched extrapolation shown in Figure 4. There is not a large difference between the two extrapolations, although both full QCD extrapolations fall below the quenched extrapolations, indicating some suppression of small eigenvalues through the presence of the fermion determinant. In the next section, we study the dynamical mass extrapolation of $`\overline{\psi }\psi `$ for larger values of $`L_s`$ to see if the non-zero value for $`\overline{\psi }\psi (m_f0)`$ decreases with increasing $`L_s`$.
Additional simulations with $`m_f=0.1`$, $`L_s=8`$ and $`m_0=1.9`$ were done for $`\beta =`$ 5.2, 5.3 and 5.45, which produced the data for $`\overline{\psi }\psi `$ and the Wilson line shown in Figure 6. Crossover behavior is seen for both observables further supporting the identification of both a chirally broken and a chirally restored phase. These simulations are at a small value of $`L_s`$, so the contribution of $`m_{\mathrm{res}}`$ to the effective quark mass may be large. Since $`m_{\mathrm{res}}(\beta ,L_s)`$ is likely varying across the transition region, due to the change in $`\beta `$, the shape of the curves is expected to reflect this varying effective quark mass.
### B $`L_s`$ dependence in the two phases
With this evidence for two phases, we turned to exploring the $`L_s`$ dependence in each phase. For the confined phase, we chose $`\beta =5.2`$ to be at weaker coupling while still in this phase and in the deconfined phase we chose $`\beta =5.45`$, to be farther from the transition. Keeping $`m_0=1.90`$, simulations were done for many values of $`L_s`$ and the dynamical quark mass, $`m_f`$. Table I gives the parameters for $`\beta =5.2`$ and Table II gives them for $`\beta =5.45`$. A plot of the evolution of $`\overline{\psi }\psi `$ for $`\beta =5.20`$ and 5.45 is shown in Figure 7 for $`m_f=0.02`$ and $`L_s=16`$. With a step size of $`\delta \tau =1/64`$ the acceptance was 90%. Once again there is no evidence for difficulty in the hybrid Monte Carlo evolution of this system.
Figure 8 shows results for $`\overline{\psi }\psi `$ at $`\beta =5.2`$ plotted versus $`m_f`$ for $`L_s=8`$ and 16. The dashed lines are linear fits to the lowest three values for $`m_f`$ while the solid lines are quadratic fits to all values of $`m_f`$. The fits for $`L_s=8`$ are
$`\overline{\psi }\psi `$ $`=`$ $`0.0117(2)+0.095(2)m_f`$ (19)
$`\overline{\psi }\psi `$ $`=`$ $`0.0112(3)+0.114(5)m_f0.15(2)m_f^2`$ (20)
with $`N_{\mathrm{dof}}=1`$ and 2 and $`\chi ^2/N_{\mathrm{dof}}=3.7`$ and 0.4, respectively. The fits for $`L_s=16`$ are
$`\overline{\psi }\psi `$ $`=`$ $`0.0082(1)+0.089(2)m_f`$ (21)
$`\overline{\psi }\psi `$ $`=`$ $`0.0080(2)+0.099(3)m_f0.08(1)m_f^2`$ (22)
with $`N_{\mathrm{dof}}=1`$ and 2 and $`\chi ^2/N_{\mathrm{dof}}=0.03`$ and 0.5, respectively. The results shows a strong $`L_s`$ dependence to which we now turn.
Figure 9 shows $`\overline{\psi }\psi `$ for $`\beta =5.2`$ plotted versus $`L_s`$ for a variety of values of $`m_f`$. The curves are fits to the form $`c_0+c_1\mathrm{exp}(\alpha L_s)`$ for $`L_s=8`$ to 40. The fit parameters are
$`\overline{\psi }\psi `$ $`=`$ $`0.01527(4)+0.0188(8)\mathrm{exp}(0.149(5)L_s)m_f=0.1`$ (23)
$`\overline{\psi }\psi `$ $`=`$ $`0.00779(8)+0.014(1)\mathrm{exp}(0.116(8)L_s)m_f=0.02`$ (24)
$`\overline{\psi }\psi `$ $`=`$ $`0.0059(1)+0.014(1)\mathrm{exp}(0.11(1)L_s)m_f0.0`$ (25)
All fits have $`N_{\mathrm{dof}}=4`$ and give $`\chi ^2/N_{\mathrm{dof}}=5.1`$, 5.6 and 6.6, respectively. The $`m_f0`$ points are first found by extrapolating to $`m_f=0`$ at fixed $`L_s`$ and then fitting these values versus $`L_s`$. Although the values for $`\chi ^2`$ are somewhat large, the data is well fit by a function with exponential dependence on $`L_s`$. (Note these somewhat large $`\chi ^2`$ values can be caused by underestimates of the errors which may result if our Monte Carlo evolutions are not sufficiently long to allow proper control the long-time autocorrelations.)
Similar results have been obtained for $`\beta =5.45`$. Figure 10 shows the results for $`\overline{\psi }\psi `$ for $`\beta =5.45`$ for $`L_s=8`$ and 16. ($`L_s=24`$ and 32 results are tabulated below.) Again, the dashed lines are linear fits to the lowest three values for $`m_f`$ while the solid lines are quadratic fits to all values of $`m_f`$. The fits for $`L_s=8`$ are
$`\overline{\psi }\psi `$ $`=`$ $`0.00227(7)+0.095(1)m_f`$ (26)
$`\overline{\psi }\psi `$ $`=`$ $`0.00219(9)+0.099(2)m_f0.037(9)m_f^2`$ (27)
with $`N_{\mathrm{dof}}=1`$ and 2 and $`\chi ^2/N_{\mathrm{dof}}=0.6`$ and 0.1, respectively. The fits for $`L_s=16`$ are
$`\overline{\psi }\psi `$ $`=`$ $`0.00039(8)+0.100(2)m_f`$ (28)
$`\overline{\psi }\psi `$ $`=`$ $`0.00040(6)+0.100(3)m_f0.01(2)m_f^2`$ (29)
with $`N_{\mathrm{dof}}=1`$ and 2 and $`\chi ^2/N_{\mathrm{dof}}=0.09`$ and 0.02, respectively. Linear fits for the larger values of $`L_s`$ give
$`\overline{\psi }\psi `$ $`=`$ $`0.00016(8)+0.100(2)m_fL_s=24`$ (30)
$`\overline{\psi }\psi `$ $`=`$ $`0.00006(6)+0.099(1)m_fL_s=32`$ (31)
with $`N_{\mathrm{dof}}=1`$ for both $`L_s`$ and $`\chi ^2/N_{\mathrm{dof}}=0.01`$ and 7.1, respectively. We see that with increasing $`L_s`$, the extrapolated value for the condensate at $`m_f=0`$ decreases steadily.
Figure 11 shows $`\overline{\psi }\psi `$ for $`\beta =5.45`$ plotted versus $`L_s`$ for a variety of values of $`m_f`$. The curves are fits to the form $`c_0+c_1\mathrm{exp}(\alpha L_s)`$ for $`L_s=8`$ to 32. The fit parameters are
$`\overline{\psi }\psi `$ $`=`$ $`0.0102(1)+0.08(3)\mathrm{exp}(0.48(6)L_s)m_f=0.1`$ (32)
$`\overline{\psi }\psi `$ $`=`$ $`0.00599(4)+0.015(3)\mathrm{exp}(0.26(2)L_s)m_f=0.06`$ (33)
$`\overline{\psi }\psi `$ $`=`$ $`0.00213(4)+0.025(4)\mathrm{exp}(0.31(2)L_s)m_f=0.02`$ (34)
$`\overline{\psi }\psi `$ $`=`$ $`0.00010(5)+0.019(3)\mathrm{exp}(0.27(2)L_s)m_f0.0`$ (35)
All fits have $`N_{\mathrm{dof}}=3`$ and give $`\chi ^2/N_{\mathrm{dof}}=0.4`$, 4.8, 1.1 and 0.8, respectively. Here again the data strongly support exponential suppression of mixing between the walls for $`\overline{\psi }\psi `$.
For both the confined and deconfined cases, we see $`\overline{\psi }\psi `$ exponentially approaching a limiting value for large $`L_s`$ (which is zero in the deconfined case). At the stronger coupling of the confined phase, the decay constant is $`1/10`$, while in the deconfined phase it is $`1/4`$. One expects faster decay at weak coupling, but at present we do not know whether the different phases also play a role in the decay constant.
### C Studying the $`m_0`$ dependence of the transition
The parameter $`m_0`$ is relevant at finite lattice spacing, since it controls not only when there is a single light fermion bound to the domain walls but also the maximum momentum this fermion can have while still being bound. It is expected that this parameter will not have to be fine-tuned for domain wall fermions to work correctly, but care in choosing a value is necessary to get the correct number of light species and the maximum allowable phase space for light fermions in the thermal ensemble.
We have studied the characteristics of the transition region by choosing $`m_f=0.1`$, $`L_s=12`$ and simulating for values of $`\beta `$ near the phase transition for $`m_0=1.15`$, 1.4, 1.65, 1.8, 1.9, 2.0, 2.15 and 2.4. Tables III, IV, V, VI, VII, VIII, IX and X contain simulation parameters and results. For parameters where a deconfined thermal state was expected, the initial lattice was disordered, while an initial ordered lattice was used where a confined state was expected.
Figure 12 shows the expectation value of the magnitude of the Wilson line $`|W|`$ for these runs. A rapid crossover is seen for all values of $`m_0`$. The lines are the result of fitting the four points nearest the transition (five points where we have a point close to the transition) to the function
$$f(x)=c_0(c_1+\mathrm{tanh}[c_2(x\beta _c)]).$$
(36)
This is a phenomenologically useful form for determining the point of maximum slope for the Wilson line. The points far from the transition are not included in these fits, since this phenomenological function poorly represents the data there.
Figure 13 shows similar results for $`\overline{\psi }\psi `$ with the lines being a fit to Eq. 36. For $`m_0=1.15`$ and 1.4, the $`\overline{\psi }\psi `$ data do not allow even a rough determination of $`\beta _c`$. For small enough $`m_0`$, the light chiral modes should not exist and we have evidence for that at $`m_0=1.15`$. The value for $`\overline{\psi }\psi `$ is very small and shows little change even when the Wilson line shows evidence for the transition. In addition, the Wilson lines indicate the transition is very close to the value of 5.6925 for quenched QCD on a $`24^3\times 4`$ lattice supporting the conclusion that light fermion modes are not present in the simulations. The effects of the heavy modes are apparently quite well canceled by the Pauli-Villars fields.
Figure 14 gives estimates for $`\beta _c`$ determined from the Wilson line and $`\overline{\psi }\psi `$. These are in quite reasonable agreement, particularly given the phenomenological character of their determination. For $`m_01.2`$, $`\beta _c`$ is close to the quenched value and moves smoothly to smaller values as $`m_0`$ is increased. For these larger values for $`m_0`$, the light quark states appear and the maximum momentum for a state bound to the walls should increase. These light states make $`\overline{\psi }\psi `$ show crossover behavior and are required for our simulations to be proper studies of two-flavor QCD. At our largest value of $`m_0`$ (2.4), we may be approaching the transition from a two flavor theory to an eight flavor one (recall that the domain wall determinant is squared in our simulations, doubling the number of fermion flavors.)
## V Determining the residual mass
As mentioned in Section III, it can be expected that for long-distance physical quantities, the effects of mixing between the chiral wall states will result in a residual mass contribution to the total quark mass. This is just the statement that the dominant effect of the mixing, from the perspective of a low-energy effective Lagrangian, is to introduce another source for chiral symmetry breaking (beyond the input $`m_f`$), which takes the form of the operator $`m_{\mathrm{res}}\overline{\psi }\psi `$ at low energies. For a quantity like $`m_\pi ^2`$, whose dependence on chiral symmetry breaking can be expressed as a physical parameter times the total quark mass, the quark mass which enters should be $`m_f+m_{\mathrm{res}}`$.
However, for quantities whose sensitivity to chiral symmetry breaking effects extends up to the cutoff scale, such an argument does not go through. The chiral condensate, $`\overline{\psi }\psi `$ is such a quantity. For domain wall fermions with $`L_s\mathrm{}`$ (or staggered fermions), expanding in the input quark mass in the chirally broken phase gives
$$\overline{\psi }\psi =c_0+c_1m_f+O(m_f^2).$$
(37)
The coefficient $`c_1`$ is ultraviolet divergent in the continuum and therefore, on the lattice, gets large contributions from modes at the cutoff scale. For such an operator, the $`L_s`$ dependence is not reliably represented by just making the replacement $`m_fm_f+m_{\mathrm{res}}`$.
From this discussion, it is clear that although Figure 9 shows that the large $`L_s`$ limit for $`\overline{\psi }\psi `$ at $`m_f=0.02`$ has likely been reached by $`L_s40`$, one cannot conclude that the value for $`m_{\mathrm{res}}`$ has vanished. To measure $`m_{\mathrm{res}}`$, it is natural to look for effects in the pion mass, which is in turn governed by the axial Ward-Takahashi identity. This has been done in quenched simulations Refs. , at zero temperature, but here we are interested in determining $`m_{\mathrm{res}}`$ in the confined phase at finite temperature for small volumes for $`N_f=2`$ QCD.
Our small volumes preclude taking large separations in two-point functions to completely isolate the pion from other states. Thus a direct measurement of the pion mass or the overlap of the pion with any particular source is not possible here. Instead, we use the integrated form for the flavor non-singlet axial Ward-Takahashi identity and try to see the contributions of the pion. In the zero quark mass limit on infinite volumes, the pion contributions become poles. Thus we can look for the effects of these precursors of the pion poles, even when they do not completely dominate the Ward-Takahashi identity.
Starting from the flavor non-singlet axial Ward-Takahashi identity in and summing over all lattice points gives
$$\overline{\psi }\psi =m_f\chi _\pi +\mathrm{\Delta }J_5.$$
(38)
Here $`\psi `$ is the four-dimensional fermion field defined by Eq. 13 and the pseudoscalar susceptibility is (no sum on $`a`$)
$$\chi _\pi \frac{2}{4N_c}\underset{x}{}\overline{\psi }(x)\gamma _5\frac{\lambda ^a}{2}\psi (x)\overline{\psi }(0)\gamma _5\frac{\lambda ^a}{2}\psi (0),$$
(39)
(The factor of $`1/4N_c`$ is needed to match our normalization for $`\overline{\psi }\psi `$.) The additional contribution from chiral mixing due to finite $`L_s`$ is
$$\mathrm{\Delta }J_5\frac{2}{4N_c}\underset{x}{}j_5^a(x,L_s/2)\overline{\psi }(0)\gamma _5\frac{\lambda ^a}{2}\psi (0),$$
(40)
where
$`j_5^a(x,L_s/2)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\overline{\mathrm{\Psi }}(x,L_s/2)(1\gamma _5)\lambda ^a\mathrm{\Psi }(x,L_s/2+1)`$ (41)
$``$ $`{\displaystyle \frac{1}{4}}\overline{\mathrm{\Psi }}(x,L_s/2+1)(1+\gamma _5)\lambda ^a\mathrm{\Psi }(x,L_s/2)`$ (42)
is a pseudoscalar density at the midpoint of the fifth dimension which couples left- and right-handed degrees of freedom.
We have done extensive simulations for many values of $`L_s`$ with $`\beta =5.2`$, $`m_0=1.9`$ and $`m_f=0.02`$ to study the consequences of the Ward-Takahashi identity. At the time of these simulations, we were not measuring $`\mathrm{\Delta }J_5`$ explicitly. However, the other two terms in the Ward-Takahashi identity were measured, allowing a determination of the $`\mathrm{\Delta }J_5`$ term. Figure 15 shows $`\overline{\psi }\psi `$, $`\chi _\pi `$ and $`\mathrm{\Delta }J_5`$ for a variety of values of $`L_s`$. Fitting $`\mathrm{\Delta }J_5`$ to an exponential form for $`L_s=16`$ to 40 gives the solid line in the figure and the result
$$\mathrm{\Delta }J_5=0.0096(2)\mathrm{exp}(0.0191(9)L_s)\chi ^2/\mathrm{dof}=6.4/2$$
(43)
We see that our data is consistent with $`\mathrm{\Delta }J_5`$ vanishing as $`L_s\mathrm{}`$, although the decay constant is quite small, $`1/50`$.
Pion poles should dominate the Ward-Takahashi identity when the pions are light and the pions should become massless when $`m_f+m_{\mathrm{res}}=0`$. (This is only strictly true in the infinite volume limit.) Thus we look for the pseudoscalar susceptibility in large volumes for small total quark mass to behave as
$$\chi _\pi =a_1/(m_f+m_{\mathrm{res}})+a_0+𝒪(m_f+m_{\mathrm{res}}).$$
(44)
where the $`a_i`$ are independent of $`L_s`$ and $`m_f`$. This gives a pion pole (for large volumes) at $`m_f=m_{\mathrm{res}}`$, while $`a_0`$ gives the contribution to the susceptibility of modes whose mass is non-zero when the quark mass vanishes. Like $`\overline{\psi }\psi `$, $`a_0`$ receives contributions diverging as $`1/a^2`$ and hence may be sensitive to unphysical 5-dimensional modes. For this expression to be useful, we do not require the pole term to dominate the remaining terms, but it must make a large enough contribution to be visible.
The $`\mathrm{\Delta }J_5`$ term in Eq. 38 also has a pole contribution coming from the propagation of the conventional light pseudoscalar along the $`s=0`$ and $`L_s1`$ boundaries from $`0`$ to $`x`$. This light state has non-zero overlap with the midpoint pseudoscalar density for finite $`L_s`$, but this overlap should be exponentially suppressed. Therefore we expect $`\mathrm{\Delta }J_5`$ to also have a pole at $`m_f=m_{\mathrm{res}}`$, giving $`\mathrm{\Delta }J_5`$ the same form as $`\chi _\pi `$, namely
$$\mathrm{\Delta }J_5=b_1^{}/(m_f+m_{\mathrm{res}})+b_0^{}+𝒪(m_f+m_{\mathrm{res}}).$$
(45)
Considering the case where the pole terms dominate gives
$$\overline{\psi }\psi =\frac{a_1m+b_1^{}}{m+m_{\mathrm{res}}}$$
(46)
For $`\overline{\psi }\psi `$ to be finite in this case requires
$$a_1m+b_1^{}=a_1(m+m_{\mathrm{res}})$$
(47)
so the most general form for $`\mathrm{\Delta }J_5`$ is
$$\mathrm{\Delta }J_5=m_{\mathrm{res}}\chi _\pi +b_0+𝒪(m_f+m_{\mathrm{res}}),$$
(48)
Where $`b_0=b_0^{}m_{\mathrm{res}}a_0`$. Using this then gives
$$\overline{\psi }\psi =(m_f+m_{\mathrm{res}})\chi _\pi +b_0$$
(49)
up to terms linear in the quark mass.
Our procedure for extracting $`m_{\mathrm{res}}`$ from these small volumes involves measuring values for $`\chi _\pi `$ and $`\overline{\psi }\psi `$ for a variety of valence quark masses for a simulation with a fixed dynamical quark mass. Since the Ward-Takahashi identity is a consequence of the form of the domain wall fermion operator, independent of the weight used to generate the gauge field ensemble in which the fermionic observables are measured, it is satisfied by observables measured with valence masses. Of course, extrapolations in valence quark mass can lead to problems due to the gauge field ensemble including configurations with small fermion eigenvalues that are not present when a dynamical extrapolation is done. Here we have a small dynamical quark mass present in the generation of the gauge fields, so such effects are expected to be unimportant.
For a given $`L_s`$, we simultaneously fit $`\chi _\pi `$ and $`\overline{\psi }\psi `$ to the forms in Eqs. 44 and 49. These are four parameter fits for $`a_0,a_1,b_0`$ and $`m_{\mathrm{res}}`$ and the resulting value for $`m_{\mathrm{res}}`$ we refer to as $`m_{\mathrm{res}}^{(\mathrm{GMOR})}`$. (All measurements of the residual mass from low energy physics should agree. We use this notation to detail the explicit technique we have used for this determination.) We have used quark masses of 0.02, 0.06, 0.10 and 0.14 in our fits. These fits do not include possible correlations between the quantities computed for different values of $`m_f`$ because the correlation matrix itself is poorly determined.
The results are given in Table XI, where the errors are all from application of the jack knife method. Notice that $`b_0`$ is negative for all values of $`L_s`$, meaning that the non-pole contributions to $`\mathrm{\Delta }J_5`$ are smaller than $`m_{\mathrm{res}}a_0`$. We have then fit these values of $`m_{\mathrm{res}}^{(\mathrm{GMOR})}`$ and $`b_0`$ to the form $`c_0+c_1\mathrm{exp}(\alpha L_s)`$ and found
$`b_0`$ $`=`$ $`0.0104(4)\mathrm{exp}(0.016(2)L_s)\chi ^2/\mathrm{dof}=0.34(19)`$ (50)
$`m_{\mathrm{res}}^{(\mathrm{GMOR})}`$ $`=`$ $`0.185(6)\mathrm{exp}(0.0280(15)L_s)\chi ^2/\mathrm{dof}=0.28(25)`$ (51)
Figure 16 shows these values and the fits.
We can see that both $`m_{\mathrm{res}}^{(\mathrm{GMOR})}`$ and $`b_0`$ are falling exponentially, but with a very small decay constant $`1/50`$. This is in sharp contrast to the decay constant for $`\overline{\psi }\psi `$ which is $`1/10`$. This is further evidence for the distinction between the residual mass that enters in low-energy observables and the residual mixing which effects observables dependent on degrees of freedom at the cutoff scale.
Since our determination of the residual mass has been done for small volumes, one can worry about the finite volume effects. We have done a similar extraction of the residual mass and compared it with determinations of the residual mass from extrapolations of $`m_\pi ^2`$ for much larger volumes and find reasonable agreement . We are continuing to study various determinations of the residual mass.
## VI Conclusions
In this work the properties of domain wall fermions relevant to numerical simulations of full $`N_f=2`$ QCD at finite temperature were investigated on relatively small lattices of size $`8^3\times 4`$. Conventional numerical algorithms (the Hybrid Monte Carlo and the conjugate gradient) worked without any difficulty beyond the additional computational load of the fifth dimension. Evidence for both confined and deconfined phases was found and the $`L_s`$ and $`m_0`$ dependence of each phase was investigated.
The domain wall fermion action is expected to preserve the full chiral symmetries of QCD for large $`L_s`$. For the stronger couplings used for the confined phase simulations, the chiral condensate approached its asymptotic value for $`L_s3240`$. However, our determination of the residual mass effects present in low energy observables show a residual mass of $`0.06`$ for $`L_s=40`$. For the weaker couplings needed to study the deconfined, chirally restored phase, the residual mass effects are expected to be much smaller for the same $`L_s`$, although we have not yet measured the residual mass in this region.
In particular, it was found that for the two flavor theory there is a phase where the $`SU(2)\times SU(2)`$ chiral symmetry is broken spontaneously to a full $`SU(2)`$ flavor symmetry and a phase where the full $`SU(2)\times SU(2)`$ chiral symmetry is intact. For the values of $`L_s`$ we used, the dependence of observables on the coupling in the transition region is likely quite influenced by the change in the residual mass with the coupling. To suppress this effect will require larger values for $`L_s`$, thermodynamics studies at larger $`N_t`$ (and hence weaker coupling) or improved variants of domain wall fermions.
Our simulations show that domain wall fermions have passed one vital test for numerical work, light chiral modes exist at quite strong coupling. A second important result, which was expected from work with dynamical fermions in the Schwinger model , is that domain wall fermions do not present any problems to conventional dynamical fermion numerical algorithms. Given these results, we are pursuing simulations of the phase transition on larger lattices to achieve more physically meaningful results. The slow falloff of the residual mass with $`L_s`$ can be overcome with more computing power or, hopefully, improvements to the formulation. At present, this is all that stands in the way of simulating the $`N_f=2`$ QCD phase transition with three degenerate light pions at finite lattice spacing.
## Acknowledgments
The numerical calculations were done on the 400 Gflop QCDSP computer at Columbia University. This research was supported in part by the DOE under grant # DE-FG02-92ER40699 and for P. Vranas in part by NSF under grant # NSF-PHY96-05199.
## A Gamma matrices
The Dirac gamma matrices used in this work are:
$`\gamma _1=\left(\begin{array}{cccc}0& 0& 0& i\\ 0& 0& i& 0\\ 0& i& 0& 0\\ i& 0& 0& 0\end{array}\right),\gamma _2=\left(\begin{array}{cccc}0& 0& 0& 1\\ 0& 0& 1& 0\\ 0& 1& 0& 0\\ 1& 0& 0& 0\end{array}\right),`$ (A9)
$`\gamma _3=\left(\begin{array}{cccc}0& 0& i& 0\\ 0& 0& 0& i\\ i& 0& 0& 0\\ 0& i& 0& 0\end{array}\right),\gamma _4=\left(\begin{array}{cccc}0& 0& 1& 0\\ 0& 0& 0& 1\\ 1& 0& 0& 0\\ 0& 1& 0& 0\end{array}\right),`$ (A18)
$`\gamma _5=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right).`$ (A23)
## B Evolution Algorithm
As described in Section II, we use the Hybrid Monte Carlo ‘$`\mathrm{\Phi }`$’ algorithm of Gottlieb et al. extended to include the Pauli-Villars regulator fields. Further, we use a preconditioned variant of the Dirac operator specified in Eq. 5. In this Appendix we describe the resulting algorithm we use to evolve the gauge fields including the effects of the two flavors of domain wall quarks and the Pauli-Villars regulator fields.
Following this approach, we generate a Markov chain of gauge fields $`U_\mu (x)`$, pseudo-fermion fields $`\mathrm{\Phi }_F`$, Pauli-Villars fields $`\mathrm{\Phi }_{PV}`$ and conjugate momenta $`H_\mu (x)`$ according to the distribution:
$$Z=[dU][dH][d\mathrm{\Phi }_F^{}][d\mathrm{\Phi }_F][d\mathrm{\Phi }_{PV}^{}][d\mathrm{\Phi }_{PV}]e^{}$$
(B1)
where
$$=S_G+\frac{1}{2}\underset{x,\mu }{}H_\mu (x)^2+\mathrm{\Phi }_F^{}[\stackrel{~}{D}_F^{}\stackrel{~}{D}_F]^1\mathrm{\Phi }_F+\mathrm{\Phi }_{PV}^{}[\stackrel{~}{D}_F^{}\stackrel{~}{D}_F]_{m_f=1}\mathrm{\Phi }_{PV}.$$
(B2)
Here, the fields $`\mathrm{\Phi }_F`$ and $`\mathrm{\Phi }_{PV}`$ as well as the preconditioned operator $`\stackrel{~}{D}_F`$ are defined only on odd sites with
$$\stackrel{~}{D}_F=(5m_0)^2(D_F)_{oe}(D_F)_{eo}$$
(B3)
where $`(D_F)_{oe}`$ and $`(D_F)_{eo}`$ represent the DWF operator of Eq. 5 evaluated between odd and even or even and odd sites respectively. Note, even and odd are defined in a five-dimensional sense, e.g. for an even site the sum of all five coordinates is an even number. Eq. B3 employs the usual preconditioning scheme for Wilson fermions implemented in 5 dimensions. Similar considerations justify the form used for the Pauli-Villars action. Since $`det\stackrel{~}{D}_F=det\{(5m_0)D_F\}`$, we have rescaled both the fields $`\mathrm{\Phi }_F`$ and $`\mathrm{\Phi }_{PV}`$ to introduce the extra factor of $`(5m_0)`$ into Eq. B3 in order to simplify the subsequent algebra.
To begin a new HMC trajectory, we start with the values of the gauge fields $`U_\mu (x)`$ produced by the previous trajectory. We then choose Gaussian distributed fields $`\eta (x,s)_F`$, $`\eta (x,s)_{PV}`$ and $`H_\mu (x)`$ from which we construct the fields $`\mathrm{\Phi }_F=\stackrel{~}{D}_F\eta _F`$ and $`\mathrm{\Phi }_{PV}=(\stackrel{~}{D}_F^1|_{m_f=1})\eta _{PV}`$. Here we have introduced new field variables $`H_\mu (x)`$, conjugate to the link matrices, which are elements of the algebra of $`SU(3)`$, and hence traceless and hermitian.
Next, we carry out the molecular dynamics time evolution of the fields $`H_\mu (x)`$ and $`U_\mu (x)`$ according to equations of motion which are phase space volume preserving and conserve the fictitious 6-dimensional “energy” $``$ of Eq. B2. The first of these Hamilton-like equations determines the relation between $`U_\mu (x)`$ and the conjugate variable $`H_\mu (x)`$:
$$\frac{\mathrm{d}U_\mu (x)}{\mathrm{d}\tau }=iH_\mu (x)U_\mu (x).$$
(B4)
The second equation can be derived from the requirement that $``$ is $`\tau `$-independent. First, following Gottlieb et al. one writes:
$$\frac{\mathrm{d}}{\mathrm{d}\tau }=\underset{x,\mu }{}\mathrm{Tr}\left[iH_\mu (x)F_\mu (x)+\frac{\mathrm{d}H_\mu (x)}{\mathrm{d}\tau }H_\mu (x)\right].$$
(B5)
Then the constancy of $``$ is insured if for the second equation of motion we impose:
$$i\frac{\mathrm{d}H_\mu (x)}{\mathrm{d}\tau }=\left[F_\mu (x)\right]_{TA}.$$
(B6)
The subscript $`TA`$ indicates the traceless anti-hermitian part of the matrix, a restriction required by the traceless, hermitian character of the variables $`H_\mu (x)`$. (The definition of $`F_\mu (x)`$ implied by Eq. B5 makes $`F`$ anti-hermitian and it is only the traceless part of $`F`$ that enters that equation.)
Finally we will determine the specific form for the force term $`F_\mu (x)`$. This can be done by using the general formula
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}\tau }}\psi ^{}|D_F^{()}|\psi `$ $`=`$ $`{\displaystyle \frac{i}{2}}{\displaystyle \underset{x,s}{}}\{\psi ^{}(x,s)^{}H_\mu (x)U_\mu (x)(1\gamma ^\mu )\psi (x+\mu ,s)`$ (B8)
$`\psi ^{}(x+\mu ,s)^{}U_\mu (x)^{}H_\mu (x)(1\pm \gamma ^\mu )\psi (x,s)\}`$
which follows immediately from Eq’s. 7 and B5 where the lower choice of signs corresponds to the case of $`D_F^{}`$. Now we re-express the derivative:
$$\frac{\mathrm{d}}{\mathrm{d}\tau }\mathrm{\Phi }_F^{}[\stackrel{~}{D}_F^{}\stackrel{~}{D}_F]^1\mathrm{\Phi }_F=\chi _F^{}\left[\frac{\mathrm{d}}{\mathrm{d}\tau }\stackrel{~}{D}_F^{}\stackrel{~}{D}_F\right]\chi _F,$$
(B9)
where we construct $`\mathrm{\Phi }=\stackrel{~}{D}_F\eta _F`$ from the Gaussian source $`\eta _F`$ and then obtain $`\chi _F`$ by solving $`\stackrel{~}{D}_F^{}\stackrel{~}{D}_F\chi _F=\mathrm{\Phi }_F`$. Now we must evaluate
$`\chi _F^{}\left[{\displaystyle \frac{\mathrm{d}}{\mathrm{d}\tau }}\stackrel{~}{D}_F^{}\stackrel{~}{D}_F\right]\chi _F`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}\tau }}\chi _F|[(5m_0)^2(D_F^{})_{oe}(D_F^{})_{eo}]`$ (B11)
$`[(5m_0)^2(D_F)_{oe}(D_F)_{eo}]|\chi _F,`$
We will obtain eight terms by letting the derivative act on each of the four $`D_F`$ operators. Four of those terms will involve $`U_\mu (x)`$ and four $`U_\mu (x)^{}`$, with the final four terms being the hermitian conjugates of the first four. Combining Eq.’s B5, B8, B9 and B11, we find:
$`\mathrm{tr}\{H_\mu (x)F_\mu (x)\}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{s}{}}\{\chi _F(x,s)H_\mu (x)U_\mu (x)(1+\gamma ^\mu )(x+\mu ,s)|(D_F^{})_{eo}\stackrel{~}{D}_F|\chi _F`$ (B12)
$`+`$ $`\chi _F|(D_F^{})_{oe}|x,sH_\mu (x)U_\mu (x)(1+\gamma ^\mu )x+\mu ,s|\stackrel{~}{D}_F|\chi _F`$ (B13)
$`+`$ $`\chi _F|\stackrel{~}{D}_F^{}|x,sH_\mu (x)U_\mu (x)(1\gamma ^\mu )x+\mu ,s|(D_F)_{eo}|\chi _F`$ (B14)
$`+`$ $`\chi _F|\stackrel{~}{D}_F^{}(D_F)_{oe}|x,sH_\mu (x)U_\mu (x)(1\gamma ^\mu )\chi (x+\mu ,s)\mathrm{h}.\mathrm{c}.\}.`$ (B15)
This expression can be written in a very simple form if we define two new spinor quantities:
$`w(x,s)`$ $`=`$ $`\{\begin{array}{cc}x,s|(D_F^{})_{eo}\stackrel{~}{D}_F|\chi _F\hfill & (x,s)\text{ even}\hfill \\ x,s|\stackrel{~}{D}_F|\chi _F\hfill & (x,s)\text{ odd}\hfill \end{array}`$ (B18)
$`v(x,s)`$ $`=`$ $`\{\begin{array}{cc}x,s|(D_F)_{eo}|\chi _F\hfill & (x,s)\text{ even}\hfill \\ \chi _F(x,s)\hfill & (x,s)\text{ odd}\hfill \end{array}`$ (B21)
Using these quantities in Eq. B15 and factoring out the generator $`H_\mu (x)`$ gives:
$`F_{[F]\mu }(x)`$ $`=`$ $`{\displaystyle \frac{1}{2}}U_\mu (x){\displaystyle \underset{s}{}}\mathrm{tr}_{\mathrm{spin}}[(1\gamma _\mu )v(x+\widehat{\mu },s)w^{}(x,s)`$ (B23)
$`+(1+\gamma _\mu )w(x+\widehat{\mu },s)v^{}(x,s)]\mathrm{h}.\mathrm{c}.`$
where we have added now the subscript $`[F]`$ to distinguish this fermion force from that produced by the Pauli-Villars fields described below. Since there are no gauge fields in the extra direction, it is not surprising that this looks very similar to the Wilson fermion force with an additional sum over the s-direction.
The force term produced by the Pauli-Villars fields is closely related to that derived above. We need only replace the field $`\chi _F`$ with $`\mathrm{\Phi }_{PV}`$, set $`m_f=1`$ and change the sign of the resulting force:
$$F_{[PV]\mu }(x)=F_{[F]\mu }(x)|_{m_f=1,\chi _F=\mathrm{\Phi }_{\mathrm{PV}}}.$$
(B24)
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# References
A critical examination of the spin dynamics in high-$`T_C`$ cuprates
Ph. Bourges<sup>1</sup>, B. Keimer<sup>2</sup>, L.P. Regnault<sup>3</sup> and Y. Sidis<sup>1</sup>
| 1 | Laboratoire Léon Brillouin, CEA-CNRS, CE Saclay, 91191 Gif sur Yvette, France |
| --- | --- |
| 2 | Max-Planck-Institut für Festkörperforschung, 70569 Stuttgart, Germany |
| 3 | CEA Grenoble, Département de Recherche Fondamentale sur la matière Condensée, 38054 Grenoble cedex 9, France |
abstract
A critical examination of the spin dynamics in high-$`T_C`$ cuprates is made on the light of recent inelastic neutron scattering results obtained by different groups. The neutron data show that incommensurate magnetic peaks in YBCO belong to the same excitation as the resonance peak observed at $`(\pi /a,\pi /a)`$. Being only observed in the superconducting state, the incommensurability is then rather difficult to reconcile with a stripe picture. We also discuss the link between the resonance peak spectral weight and the superconducting condensation energy.
After more than ten years of intense investigations, the precise role of antiferromagnetic (AF) correlations for the mechanism of the high-temperature superconductivity remains a puzzling and open question. Since the early days, it has been obvious that both phenomena are clearly connected just by looking at the generic phase diagram of high-$`T_C`$ cuprates. Of course, a competitive role rather than cooperative between long-range antiferromagnetism and superconductivity was generally inferred as both phenomena are thought to occur in exclusion of each other. The next key question was then: are the dynamical AF correlations observed in the superconducting (SC) range of the phase diagram prejudicial or responsible for superconductivity ?
A necessary step to put some insight into this still unsolved question is the knowledge of the spectral weight of the spin susceptibility, $`\chi (Q,\mathrm{}\omega )`$. $`\chi (Q,\mathrm{}\omega )`$ would, for instance, enter the SC pairing interactions in any mechanism based on antiferromagnetism. As a matter of fact, Inelastic Neutron Scattering (INS) is the only technique which directly measures the full energy and momentum dependences of the imaginary part of the spin susceptibility. Further, the amplitude of $`Im\chi (Q,\mathrm{}\omega )`$ can be determined in absolute units by a calibration of the magnetic neutron intensity versus other scattering such as phonons. This has been done only recently for the high-$`T_C`$ cuprates and, brings essential insight for the relation between AF correlations and superconductivity as we shall see below.
This technique is limited by the need of large single crystals (of cm<sup>3</sup> size) usually difficult to grow in complex systems such as high-$`T_C`$ cuprates. This has reduced the number of systems which could be studied to a very few: $`\mathrm{La}_{2\mathrm{x}}\mathrm{Sr}_\mathrm{x}\mathrm{CuO}_4`$ (LSCO), $`\mathrm{YBa}_2\mathrm{Cu}_3\mathrm{O}_{6+\mathrm{x}}`$(YBCO) and only recently $`\mathrm{Bi}_2\mathrm{Sr}_2\mathrm{CaCu}_2\mathrm{O}_{8+\delta }`$ (BSCO). Further, INS spectra can be sometimes ambiguous to analyze as, for instance, neutron scattering also directly measures the phonon spectrum which is typically of the same order of magnitude. Of course, a lot of effort has been developed to overcome these difficulties. However, this situation has postponed the emergence of a full agreement between the different groups. However, clear unmistakable features have been established which have considerable impact on the role of AF fluctuations. Here, we shall emphasize some key aspects on the basis of published data by the different groups. It should be mentioned that as far as the raw data are concerned, a fairly good agreement can be noticed. Disagreements are rather related to the data analysis which sometimes leads to clearly different conclusions.
The resonance peak: a collective spin excitation of $`d`$-wave superconductors
Among the observed magnetic features , the ”AF resonance peak” observed below $`T_C`$ is certainly one of most important results which has been widely studied since its discovery in 1991 by Rossat-Mignod et al in YBCO<sub>6.92</sub>. When entering the SC state and only below $`T_C`$, a sharp (almost energy resolution limited) spin excitation appears in the neutron scattering data at an energy, $`E_r`$ and at the AF wavevector $`Q_{AF}=(\pi /a,\pi /a)`$ ($`Q_{AF}`$ is the propagation wavevector of the AF state of the insulating undoped parent compound, a=3.85 Å is the 2D square lattice parameter).
The striking characteristic of the resonance peak is actually its temperature dependence. Indeed, its energy, $`E_r`$, does not shift towards lower energy when approaching $`T_C`$ (a shift of at most $``$ 4 % can be inferred ) but its intensity is vanishing upon heating at the superconducting temperature $`T_C`$ for all doping levels, actually following an order parameter-like behavior. Recently, an attempt has been made to associate the vanishing of the resonance peak intensity in underdoped sample with the temperature T where the resistivity displays the so-called ”pseudo-gap” anomaly. This statement is not correct being based on an arbitrary analysis. Indeed, neither the data published by Dai et al nor our own data provide any justification for a separation in the normal state (NS) of the spectrum into resonant and nonresonant parts. No published temperature dependence of the neutron intensity at the resonance energy (or more correctly, at the energy transfer where the resonance peak appears in the SC state) suggests an anomaly at a temperature T larger than $`T_C`$. A clear upturn is systematically observed only at the SC transition temperature. In our opinion, this incorrect attribution of the ”onset of the resonance peak at T” has been made from the fact that the broad maximum of the spin susceptibility in the normal state occurs in some underdoped sample roughly at the same energy as the resonance peak . But, as a matter of fact, the apparent equivalence of the normal state energy and the resonance peak energy breaks down in underdoped samples closer to optimal doping .
Interestingly, the resonance energy scales with the SC temperature as: $`E_r5.2k_BT_C`$. This relation holds in the two systems where the resonance peak has been observed so far, YBCO and BSCO. This actually is not only valid at optimal doping but also remains correct on both sides of the high $`T_C`$ phase diagram: on the underdoped side, as experimentally realized for different oxygen contents in YBCO, as well as on the overdoped side as observed in a BSCO sample. This generic relationship of $`T_C`$ with the temperature-independent resonance energy calls for an explanation which is not obvious when one considers the different models usually invoked to interpret the resonance peak (See Refs for a discussion of these approaches).
Further, the resonance feature appears to be strongly sensitive to parameters which affect the superconducting properties. For instance, the substitution of Zn impurities within the CuO<sub>2</sub> plane in YBCO, known to strongly reduce the SC temperature ($`dT_C/dy12`$ K/%) without changing the doping level, induces a rapid vanishing of the resonance intensity: small amounts of zinc impurities ($`y`$ ranging from 0.5% to 2% in YBa<sub>2</sub>(Cu<sub>1-y</sub>Zn<sub>y</sub>)<sub>3</sub>O<sub>6+x</sub>) are sufficient to remove its spectral weight without strong renormalization of the resonance energy itself. In contrast, magnetic Ni impurities which are three times less efficient to remove superconductivity ($`dT_C/dy4`$ K/%) , have also less effect on the resonance peak intensity and keep the ratio $`E_r/k_BT_C`$ almost unchanged . This extreme sensitivity of the resonance feature to defects affecting the SC transition temperature then might explain why no resonance peak has been reported so far in the LSCO system whose maximum $`T_C`$ ($``$ 40 K) is anomalously low as compared to other single CuO<sub>2</sub> layer systems where $`T_C`$ can reach 90 K (Tl- or Hg- based system). The disorder which might be responsible for the reduction of $`T_C`$ in LSCO can also remove the resonance peak feature. Further, Zn and Ni impurities in YBCO also produce a systematic broadening in energy of the resonance peak, by $``$ 10 meV. Similar broadening found in BSCO can then be naturally accounted for by the presence of intrinsic defects in that system.
Until recently, the resonance peak has been widely described as a single commensurate excitation. Although this statement remains certainly correct in the slightly overdoped YBCO<sub>7</sub> system, we have recently demonstrated in YBCO<sub>6.85</sub> that the resonance peak actually exhibits a full dispersion curve away from $`(\pi /a,\pi /a)`$ momentum. This illustrates, on experimental grounds, that the resonance peak can be considered as a collective mode of the superconducting state of high-$`T_C`$ cuprates as theoretically proposed (see e.g. ). The observed downward dispersion actually relates the commensurate resonance peak with the incommensurate peaks observed at lower energy and recently reported in underdoped YBCO. By detailed temperature dependences of the neutron intensity at different wavevectors and energies, we have established a dispersion compatible with the following relationship,
$$E_r(q)=\sqrt{E_r^2(Q_{AF})(\alpha q)^2}$$
(1)
where $`q`$ is the wavevector measured from $`Q_{AF}=(\pi /a,\pi /a)`$. $`E_r(Q_{AF})=41`$ meV is the previous commensurate resonance energy, and $`\alpha `$ 125 meV.Å represents an isotropic dispersion relation. Certainly, the relation Eq. 1 is only a first approximation which needs to be refined. Indeed, the measured wavevector pattern at a fixed energy E= 35 meV located below $`E_r(Q_{AF})`$ exhibits an intensity modulation in the 2D $`(H,K)`$ momentum space shown in Fig. 1 with larger intensity in the directions (100) or (010) and lower intensity in the directions (110) or (1$`\overline{1}`$0). Such detailed momentum dependence (which reproduces the shape reported in YBCO<sub>6.6</sub> at 24.5 meV (below $`E_r(Q_{AF})=34`$ meV) as well as that discussed in ), implies a modification in the dispersion relation of Eq. 1. For instance, an anisotropy of $`\alpha `$ between the (100) and (110) directions should be added and would certainly account for the momentum pattern of the neutron intensity shown in Fig. 1. Although the resonance peak dispersion is, so far, only evidenced in one sample, YBCO<sub>6.85</sub>, we think it is a generic feature of the spin dynamics in the superconducting state over a wide part of the high-$`T_C`$ cuprate phase diagram. Data reported in Refs. are fully consistent with such an interpretation although this has not been discussed this way. For sure, more work is needed to generalize this conclusion, for instance, to give the actual doping dependence of the $`\alpha `$ parameter.
The observation of incommensurate peaks , in addition to the commensurate resonance peak, has stimulated several theoretical models in Fermi liquid-like theories. It has been discussed as a combined effect of both i) topology of the band structure and ii) anisotropic superconducting order parameter either at the level of the bare susceptibility or after taking into account of the interactions by a random phase approximation . Furthermore, a dispersive collective mode has been predicted to arise below the particle-hole spin-flip continuum in the $`d`$-wave superconducting state as a result of a momentum-dependent pole in the spin susceptibility pulled by antiferromagnetic interactions . Our recent observation of a downward dispersion supports the latter proposal. However, to fully establish the collective nature of the resonance peak, a necessary step will be to observe the particle-hole spin-flip continuum. In any case, our recent data demonstrate that superconductivity affects not only the energy lineshape of the spin susceptibility by inducing a resonance peak at $`(\pi /a,\pi /a)`$ but also that it drastically changes its momentum dependences.
”Incommensurate peaks” in YBCO: not an evidence for dynamical stripes
The observation of ”incommensurate peaks” at some energy transfers has often been interpreted as clearcut evidence of dynamical stripes in YBCO. Our recent detailed study basically rules out this conclusion (at least for near-optimally doped YBCO). Indeed, we established that the ”incommensurate peaks” are only observed in the superconducting state and are additionally closely related to the commensurate resonance peak by a continuous dispersion relation (Eq. 1) as discussed above. This puts the observation of the magnetic incommensurability in YBCO in a totally new perspective. Being energy-dependent, temperature-dependent and doping-independent, the ”discommensuration” is rather difficult to understand within a stripe picture where typically a characteristic distance (between charge stripes) needs to be observed. Without invoking any specific model, it becomes clear that their interpretation has to be necessarily related to the one made for the ”commensurate” resonance peak.
The vanishing at $`T_C`$ of the “incommensurate” excitations, we reported in YBCO<sub>6.85</sub> , can be actually anticipated over a wide part of the phase diagram. \[Notice that, even below $`T_C`$, it is still not established under which conditions and exactly in which doping range the “incommensurate” excitations are present in YBCO.\] Nevertheless, their disappearance in the normal state is actually consistent with the different data published so far. Indeed, the reports of normal state incommensurability in YBCO are rather scarce. At best, it is said that these incommensurate excitations remain in a small temperature window above $`T_C`$ (up to 70-75 K for $`T_C`$=63 K) and finally disappear upon heating. But, as this intensity is weak on top of a phononic background (always present in such unpolarized neutron scattering experiments) whose structure factor mimics an incommensurate-like intensity modulation, no clear conclusion can be made and, at least, requires further work. In any case, fluctuations of the SC state (in the conventional meaning) could also explain the persistence of “incommensurate” excitations in a small temperature range above $`T_C`$. Recently, it has been argued that these incommensurate magnetic fluctuations have a one-dimensional nature . This is based on measurements using a partially (half) detwinned YBCO<sub>6.6</sub> sample. Due to the above-mentioned phononic background and the scattering geometry used, this report is rather inconclusive: it is not proved that the observed effect is related to the magnetic scattering. Indeed, the detwinning of the sample can actually affect the background itself (for instance, if it is related to an $`a^{}`$-polarized phonon). To make their point clear, these authors have to demonstrate that the balance of intensity between $`a^{}`$ and $`b^{}`$ is not present at high temperature (where the magnetic intensity is weaker and commensurate) or present polarized neutron beam data.
Our results in YBCO<sub>6.85</sub> also contrast with those reported in the LSCO system where no change of the incommensurate peak position occurs across the superconducting temperature. In LSCO, the incommensurate peak structure begins to disappear only around room temperature . Further, the observed energy range where incommensurate peaks are observed is very different in the two systems (down to the lowest energies in LSCO but limited in a small energy range below $`E_r`$ in YBCO). Their similarity is then reduced to only the symmetry of the incommensurate pattern along the (100) or (010) directions seen in both systems. However, the actual “fortified castle”-like shape observed in YBCO looks rather different from the four well defined peaks observed in LSCO. This makes dubious the universality of the spin fluctuations claimed to occur in the two systems only based on the occurrence of “incommensurate” magnetic peaks. The origin of incommensurability in both systems likely requires a different scenario although common ingredients (such as Fermi surface topology) might be invoked.
As discussed above, our detailed study of the incommensurate magnetic peaks in YBCO shows that a standard interpretation within a ’stripe phase’ picture is inconsistent. However, it should be noticed that a situation of strongly disordered stripes, as recently theoretically discussed in , is still possible. This would correspond to the case where the AF correlation length is lower than the mean distance between stripes . And so, there is no $`\pi `$-phase shift from one AF cluster to the next one. These decorrelated AF clusters would give rise to the broad commensurate peaks observed in the normal state. However, the behavior of such objects in the superconducting state has not been addressed so far. This would be of great interest.
Resonance peak and Superconducting condensation energy
The knowledge of the spin-spin correlation function in absolute units is becoming a crucial topic for the description of the physical properties of high-$`T_C`$ cuprates. For instance, magnetic neutron scattering has been recently proposed to provide a direct measurement of the condensate fraction of a superconductor. A direct link with the high-$`T_C`$ mechanism has also addressed in the framework of the t-J model. The proposal is the following: if the SC pairing mechanism is due to AF exchange then the SC condensation energy, $`E_C`$, would be the energy gain between the normal state and the superconducting state of an exchange energy $`E_J`$ of the form:
$$E_J=\frac{3J}{2\pi (g\mu _B)^2}_{BZ}d^2q[\mathrm{cos}(q_xa)+\mathrm{cos}(q_ya)]𝑑\omega \frac{Im\chi (q,\omega )}{1\mathrm{exp}(\mathrm{}\omega /k_BT)}$$
(2)
where the sum over the wavevector is performed over the 2D Brillouin zone (BZ) and normalized by the BZ volume, $`(2\pi /a)^2`$. The condensation energy then reads,
$$E_C=E_J^{NS}E_J^{SC}$$
(3)
It is essential to realize that Eq. 3 is a subtle net difference of the magnetic fluctuations spectral weight between the normal state and the superconducting state additionally weighted by a momentum form factor $`[\mathrm{cos}(q_xa)+\mathrm{cos}(q_ya)]`$ corresponding to the Fourier transform of the AF exchange. It follows that the temperature dependent change in exchange energy crucially depends on a redistribution of the magnetic spectral weight in momentum. Indeed, according to Ref. the exchange energy differs from the total moment sum rule, $`W=_{BZ}d^2q𝑑\omega Im\chi (q,\omega )/(1\mathrm{exp}(\mathrm{}\omega /k_BT))`$, only by this momentum-dependent form factor. If one neglects this wavevector dependence in Eq. 2, Eq. 3 becomes meaningless as $`E_C`$ will necessarily be zero to satisfy the sum-rule. The wavevector form factor in Eq. 2 is then essential and cannot be neglected.
In a recent Report, Dai et al. have followed this idea and claim to have found a quantitative correspondence between the temperature derivative of the spectral weight of spin excitations in YBCO and the electronic specific heat $`\mathrm{C}_{\mathrm{el}}\mathrm{dE}_\mathrm{J}/\mathrm{dT}`$. We wish to point out that the analysis provided by Dai et al. fails at an elementary level as they fully neglected the wavevector form factor in Eq. 2 by rewriting $`E_J`$ as,
$$E_J\frac{3J}{\pi (g\mu _B)^2}_{BZ}d^2q𝑑\omega \frac{Im\chi ^{res}(q,\omega )}{1\mathrm{exp}(\mathrm{}\omega /k_BT)}$$
(4)
where $`Im\chi ^{res}(q,\omega )`$ is only the resonant part of the spin excitations. Eq. 4 is derived by assuming that the spins accounting for the resonance part are fully decorrelated in the normal state in contrast with the observation of AF dynamical correlations above $`T_C`$. They then conclude that a large part of the electronic specific heat is due to spin fluctuations. There is no doubt that the electronic specific heat and the spin fluctuations are related in some way: after all, they are ultimately attributable to the same strongly interacting electron system. However, the analysis of Dai et al. is much too crude to uncover this underlying relation. In an optimally doped sample, they finally obtain a contribution to the specific heat $``$ three times larger than the measured one (as found in ). In underdoped samples, the discrepancy is even bigger as the measured specific heat jump drastically falls down whereas the resonance peak spectral weight remains approximately constant for all doping as $`d^2q𝑑\omega Im\chi ^{res}(q,\omega )0.05\pm .02\mu _B^2`$, and so would be the calculated specific heat jump at $`T_C`$. \[It should be noticed that this absolute unit value has been independently obtained by the two different groups\]. Further, the attempt to relate the specific heat anomaly in the normal state with a speculated onset of the resonance peak at $`T^{}`$ (see above) is meaningless. Indeed, the most salient feature of the electronic specific heat is its pronounced increase with increasing doping in the normal state. By contrast, the magnetic spectral weight strongly decreases with increasing doping in the same temperature range. These discrepancies do not necessarily suggest that the proposal of Eq. 2 is not correct. It just means that the analysis performed in Ref., Eq. 4, relating the magnetic fluctuation spectrum and the electronic specific heat is invalid and inconclusive as it oversimplifies the physical content of Eqs. 2 and 3. As emphasized by Scalapino and White , the net difference in Eq. 3 will be very small and then difficult to estimate. To overcome this problem, Dai et al. have arbitrarily considered only the contribution of the resonance peak spectral weight around $`(\pi /a,\pi /a)`$ and at the energy $`E_r(Q_{AF})`$ (that they attempt to relate to the electronic specific heat). The actual change of the spin susceptibility across the SC temperature as discussed above (dispersion behavior such as Eq. 1) reveals that the estimate of Eq. 3 would be very subtle (especially in underdoped samples).
In conclusion, the resonance peak is certainly a key feature for the description of the physical properties of high-$`T_C`$ superconductors which has been widely reported at the commensurate AF wave vector. Further, the observation of its dispersion experimentally suggests its collective nature. It now emerges that the role of such a magnetic collective mode would be essential for the interpretation of physical properties of high-$`T_C`$ superconductors, for instance, to describe the complex spectral structure of the one-particle spectrum as reported by photoemission spectroscopy.
Acknowledgments:
We wish to thank A.H. Castro Neto, G. Deutscher, D. Pavuna for stimulating discussions at the Klosters conference.
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# Landscape statistics of the low autocorrelated binary string problem
## 1 Introduction
The Low Autocorrelated Binary String Problem (LABSP) consists of finding binary strings $`x`$ of length $`N`$ over the alphabet $`\{\pm 1\}`$ with low aperiodic off-peak autocorrelation $`R_k(x)=_{i=1}^{Nk}x_ix_{i+k}`$ for all lags $`k`$. These strings have technical applications such as the synchronization in digital communication systems and the modulation of radar pulses.
The quality of a string $`x`$ is measured by the fitness or energy function
$$(x)=\frac{1}{2N}\underset{k=1}{\overset{N1}{}}\left[\underset{i=1}{\overset{Nk}{}}x_ix_{i+k}\right]^2=\frac{1}{2N}\underset{k=1}{\overset{N1}{}}R_k(x)^2.$$
(1)
In most of the literature on the LABSP the merit factor $`F(x)=N^2/(4(x))`$ is used (see e.g. ): using $``$ instead is more convenient for explicit computations.
Recently there has been much interest in frustrated models without explicit disorder. The LABSP and related bit-string problems have served as model systems for this avenue of research . These investigations have lead to a claim that LABSP has a ‘golf-course’ type landscape structure, which would explain the fact that it has been identified as a particularly hard optimization problem for heuristic algorithms such as Simulated Annealing (see and the references therein).
The landscape of LABSP consists of a (dominant) 4-spin Hamiltonian plus an asymptotically negligible quadratic component. We note that the generic $`4`$-spin landscape is Derrida’s $`4`$-spin Hamiltonian which is a linear combination of all $`\left(\genfrac{}{}{0pt}{}{N}{4}\right)`$ distinct $`4`$-spin functions, while the LABSP Hamiltonian, on the other hand, only contains $`𝒪(N^3)`$ non-vanishing $`4`$-spin contributions. The landscape of the LABSP thus corresponds to a dilute 4-spin ferromagnet. Numerical simulations in show that the LABSP has by far more local optima than a generic 4-spin glass model, which corroborates the rather surprising finding that disordered ferromagnets have more metastable states than their spin-glass counterparts .
In this contribution we carry out a thorough investigation of the statistical properties of the energy landscape of LABSP aiming at to determine whether it has any peculiar features that would lead to a ‘golf-course’ structure, with vanishingly small correlations between the energies of neighboring states. To do so we carry out a comparison with four disordered models, namely, the random energy model (REM), the $`\pm 1`$ 4-spin glass model , a mean-field approximation to $``$ (MF) , which reproduces the results of Golay’s ergodicity assumption , and, finally, the $`\pm 1`$ 2-spin glass model . The replica analyses indicate that the first three models have a rather unusual spin-glass phase, where the overlap between any pair of different equilibrium states vanishes, while the last model has a normal spin-glass phase described by a continuous order parameter function.
The rest of this paper is organized in the following way. In section 2 we calculate analytically the average density of local minima of the disordered mean-field approximation to $``$ and show that it indeed describes very well the statistics of metastable states of the pure model. Rather surprisingly, we find that the value of the energy density at which the density of local minima vanishes coincides with the bound predicted by Golay , as well as with the ground-state energy predicted by the first step of replica-symmetry breaking calculations of the mean-field model . To properly compare the landscapes of the different models mentioned above, in section 3 we consider two global measures of landscape structure which have been introduced in the Simulated Annealing literature: depth and difficulty . We show that LABSP, the mean-field approximation, and the binary $`\pm 1`$ 2 and 4-spin glasses exhibit approximately the same qualitative behavior in these parameters, while the behavior pattern of the random energy model departs significantly from those. Finally, in section 4 we summarize our main results and present some concluding remarks.
## 2 Mean-field approximation
Bouchaud and Mézard and, independently, Marinari et al. have proposed the following disordered model, which is “as close as possible” to the pure model:
$$_d=\frac{1}{2N}\underset{k=1}{\overset{N1}{}}\left[\underset{i=1}{\overset{N}{}}\underset{ji}{\overset{N}{}}J_{ij}^kx_ix_j\right]^2$$
(2)
Here the coupling strengths $`J_{ij}^kJ_{ji}^k`$ are statistically independent random variables that can take on the value $`1`$ with probability $`\left(Nk\right)/N^2`$ and zero otherwise. Hence the average number of bonds in $``$ and $`_d`$ is the same, namely, $`Nk`$. Moreover, the pure model is recovered with the choice $`J_{ij}^k=\delta _{i+k,j}`$. Probably the most appealing feature of this model is that its high-temperature (replica-symmetric) free-energy is identical to that obtained by Bernasconi using Golay’s ergodicity assumption , in which the squared autocorrelations $`R_k^2`$ are treated as independent random variables. As the constraints of the one-dimensional geometry are lost in the disordered Hamiltonian $`_d`$, it can be viewed as the mean field version of $``$.
The thermodynamics of the disordered model (2) is interesting on its own since, similarly to the random energy model , it presents a first order transition at a certain temperature $`T_g`$, below which the overlap between any pair of different equilibrium states vanishes . In contrast to the random energy model, however, the degrees of freedom are not completely frozen for $`T<T_g`$, and the entropy vanishes linearly with $`T`$ as the temperature decreases towards zero. To better understand the low-temperature phase of the mean-field Hamiltonian $`_d`$, in the following we will calculate analytically the expected number of metastable states $`𝒩\left(ϵ\right)`$ with a given energy density $`ϵ`$.
The energy cost per site of flipping the spin $`x_i`$ is $`\delta _d^i=\mathrm{\Delta }_i`$ where
$$\mathrm{\Delta }_i=\underset{k}{}v_i^k\left(\underset{j}{}v_j^k2v_i^k\right)$$
(3)
with
$$v_i^k=\frac{1}{\sqrt{N}}\underset{ji}{}\left(J_{ij}^k+J_{ji}^k\right)x_ix_j.$$
(4)
We say that a state $`x=(x_1,\mathrm{},x_N)`$ is a strict local minimum if $`\mathrm{\Delta }_i<0`$ for all $`i`$; in the case that the equality $`\mathrm{\Delta }_i=0`$ holds for some $`i`$, we call $`x`$ a degenerate local minimum. In the forthcoming analysis, the choice of $``$ instead of $`<`$, which is customary in optimization theory, see e.g. , does not make any difference. In section 3, however, degeneracies will play a role.
The average number of local minima with energy density $`ϵ`$ can be written as
$$𝒩\left(ϵ\right)=\text{Tr}_x\delta \left[ϵ\frac{1}{N}_d\left(x\right)\right]\underset{i}{}\mathrm{\Theta }\left(\mathrm{\Delta }_i\right)$$
(5)
where $`\text{Tr}_x`$ denotes the summation over the $`2^N`$ spin configurations and $`\mathrm{}`$ stands for the average over the couplings $`J_{ij}^k`$. Here $`\mathrm{\Theta }\left(x\right)=1`$ if $`x>0`$ and $`0`$ otherwise, and $`\delta \left(x\right)`$ is the Dirac delta-function.
Using the integral representation of the delta-function we obtain
$`𝒩\left(ϵ\right)`$ $`=`$ $`N{\displaystyle \frac{d\widehat{ϵ}}{2\pi }\text{e}^{𝐢N\widehat{ϵ}ϵ}\underset{i}{}\frac{d\mathrm{\Delta }_id\widehat{\mathrm{\Delta }}_i}{2\pi }\mathrm{\Theta }\left(\mathrm{\Delta }_i\right)\text{e}^{𝐢\widehat{\mathrm{\Delta }}_i\mathrm{\Delta }_i}\underset{ik}{}\frac{dv_i^kd\widehat{v}_i^k}{2\pi }\text{e}^{𝐢v_i^k\widehat{v}_i^k}}`$ (6)
$`\times \mathrm{exp}\left\{𝐢{\displaystyle \frac{\widehat{ϵ}}{8}}{\displaystyle \underset{k}{}}\left({\displaystyle \underset{i}{}}v_i^k\right)^2𝐢{\displaystyle \underset{ik}{}}\widehat{\mathrm{\Delta }}_iv_i^k\left({\displaystyle \underset{j}{}}v_j^k2v_i^k\right)\right\}`$
$`\times \text{Tr}_x\mathrm{exp}\left[{\displaystyle \frac{𝐢}{\sqrt{N}}}{\displaystyle \underset{ik}{}}\widehat{v}_i^k{\displaystyle \underset{ji}{}}\left(J_{ij}^k+J_{ji}^k\right)x_ix_j\right].`$
The average over the couplings can be easily carried out and, in the thermodynamic limit $`N\mathrm{}`$, it yields
$`\mathrm{ln}\mathrm{}`$ $`=`$ $`{\displaystyle \underset{k}{}}(1{\displaystyle \frac{k}{N}})[{\displaystyle \frac{2𝐢}{\sqrt{N}}}\left({\displaystyle \frac{1}{\sqrt{N}}}{\displaystyle \underset{i}{}}x_i\right)\left({\displaystyle \frac{1}{\sqrt{N}}}{\displaystyle \underset{i}{}}\widehat{v}_i^kx_i\right)`$ (7)
$`+{\displaystyle \frac{1}{N}}{\displaystyle \underset{i}{}}\left(\widehat{v}_i^k\right)^2+\left({\displaystyle \frac{1}{N}}{\displaystyle \underset{i}{}}\widehat{v}_i^k\right)^2].`$
We note that this result could have been obtained by considering the couplings $`J_{ij}^k`$ as Gaussian independent random variables with means and variances equal to $`\left(1k/N\right)/N`$. To get a physical but nontrivial thermodynamic limit we must assume that the magnetization $`_ix_i`$ scales with $`N^{1/2}`$, which results then in the vanishing of the term that contains the dependence on the spin variables in eq.(7). Droping this term, the sum over the spin configurations yields simply $`2^N`$. As the remaining calculations are rather straightforward we will only sketch them in the sequel.
To carry out the integrals over $`v_i^k`$ and $`\widehat{v}_i^k`$ we introduce the auxiliary parameters $`Nq_k=_i\left(\widehat{v}_i^k\right)^2`$, $`Nm_k=_i\widehat{v}_i^k`$, and $`r_k=_iv_i^k`$. After performing the resulting Gaussian integrals we introduce the saddle-point parameters $`NM=_i\widehat{\mathrm{\Delta }}_i`$ and $`NQ=_i\widehat{\mathrm{\Delta }}_i^2`$ which allow the decoupling of the indices $`k`$ and $`i`$. The final result is
$`𝒩\left(ϵ\right)`$ $`=`$ $`2^NN^3{\displaystyle \frac{dMd\widehat{M}}{2\pi }\frac{dQd\widehat{Q}}{2\pi }\frac{d\widehat{ϵ}}{2\pi }\mathrm{exp}\left[𝐢N\left(M\widehat{M}+Q\widehat{Q}+ϵ\widehat{ϵ}\right)\right]}`$ (8)
$`\times \mathrm{exp}\left[N{\displaystyle _0^1}𝑑z\mathrm{ln}G_0(z,\widehat{ϵ},M,Q)+N\mathrm{ln}G_1(\widehat{M},\widehat{Q})\right]`$
where
$`G_0`$ $`=`$ $`{\displaystyle \frac{dqd\widehat{q}}{2\pi }\frac{dmd\widehat{m}}{2\pi }\frac{drd\widehat{r}}{2\pi }\mathrm{exp}\left[𝐢r\widehat{r}z\left(q+m^2\right)𝐢\frac{\widehat{ϵ}}{8}r^2\right]}`$ (9)
$`\mathrm{exp}\left[𝐢\widehat{m}\left(m\widehat{r}rM\right)+𝐢\widehat{q}\left(q\widehat{r}^2r^2Q2\widehat{r}rM+4𝐢M\right)\right]`$
and
$$G_1=\frac{d\mathrm{\Delta }d\widehat{\mathrm{\Delta }}}{2\pi }\mathrm{\Theta }\left(\mathrm{\Delta }\right)\mathrm{exp}\left[𝐢\widehat{Q}\widehat{\mathrm{\Delta }}^2+𝐢\widehat{\mathrm{\Delta }}\left(\mathrm{\Delta }\widehat{M}\right)\right].$$
(10)
The integrals in eq.(8) are then evaluated in the limit $`N\mathrm{}`$ by the standard saddle-point method, while the integrals in the equations for $`G_0`$ and $`G_1`$ are trivially performed. The final result for the exponent
$$\alpha \left(ϵ\right)=\frac{1}{N}\mathrm{ln}𝒩\left(ϵ\right)$$
(11)
is simply
$`\alpha \left(ϵ\right)`$ $`=`$ $`𝐢\left[\left(2+\widehat{M}\right)M+Q\widehat{Q}+ϵ\widehat{ϵ}\right]+\mathrm{ln}\text{erfc}\left[{\displaystyle \frac{\widehat{M}}{\left(4𝐢\widehat{Q}\right)^{1/2}}}\right]`$ (12)
$`{\displaystyle \frac{1}{2}}{\displaystyle _0^1}𝑑z\mathrm{ln}\left[1+8\left(QM^2\right)z^2+8𝐢Mz+𝐢\widehat{ϵ}z\right]`$
where the saddle-point parameters $`M`$, $`\widehat{M}`$, $`Q`$, $`\widehat{Q}`$, and $`\widehat{ϵ}`$ are determined so as to maximize $`\alpha `$. In particular, a brief analysis of the saddle-point equations indicates that $`M`$, $`\widehat{Q}`$ and $`\widehat{ϵ}`$ are imaginary so that $`\alpha `$ is real, as expected. Introducing the real parameters $`\mu =𝐢M`$, $`\beta =𝐢\widehat{ϵ}`$, $`\eta =\widehat{M}/\left(4𝐢\widehat{Q}\right)^{1/2}`$, and $`\xi =Q/M^2`$, we rewrite eq.(12) as
$`\alpha \left(ϵ\right)`$ $`=`$ $`2\mu {\displaystyle \frac{\eta ^2}{\xi }}+\beta ϵ+\mathrm{ln}\text{erfc}\left(\eta \right)`$ (13)
$`{\displaystyle \frac{1}{2}}{\displaystyle _0^1}𝑑z\mathrm{ln}\left[1+\left(\beta +8\mu \right)z+8\mu ^2\left(1+\xi \right)z^2\right]`$
where we have used the saddle-point equation $`\alpha /\widehat{Q}=0`$ to eliminate $`\widehat{Q}`$. We note that in eq. (13) the parameters $`\eta `$ and $`\mu `$ are decoupled which facilitates greatly the numerical problem of maximizing $`\alpha `$.
The number of local minima, regardless of their particular energy values, is obtained by maximizing $`\alpha `$ with respect to $`ϵ`$, which corresponds to setting $`\beta =0`$ in the saddle-point equations. In this case, the value of the energy density that maximizes $`\alpha `$, denoted by $`ϵ_t`$, can be interpreted as the typical (average) energy density of the local minima. We find $`\alpha =0.4394`$ and $`ϵ_t=0.0837`$. These results agree very well with the numerical data $`\alpha 0.4388\pm (7)`$ and $`ϵ_t0.0826\pm (6)`$, obtained through the exhaustive search for $`N20`$ and averaging over $`100`$ realizations of the couplings.
Moreover, an exhaustive search for $`N30`$ yields that the exponent governing the exponential growth of the number of local minima in the pure model $``$ is $`0.453\pm (7)`$ and the typical energy density of the minima is $`0.086\pm (2)`$. Hence, so far as the statistics of metastable states is concerned, the mean-field Hamiltonian $`_d`$ yields in fact a very close approximation to the pure Hamiltonian $``$. For the purpose of comparison we note that $`\alpha =0.1992`$ and $`\alpha =0.3552`$ for the binary $`\pm 1`$ 2-spin glass and 4-spin glass models , respectively, while $`\alpha =\mathrm{ln}20.6931`$ for the random energy model .
In Fig.1 we show the exponent $`\alpha `$ as a function of the energy density $`ϵ`$. For the sake of clarity we present only the region of positive values of $`\alpha `$. The lowest value of $`ϵ`$ at which the exponent $`\alpha `$ vanishes, denoted by $`ϵ_0`$, gives a lower bound to the ground-state energy density of the spin model defined by the Hamiltonian (2) . We find $`ϵ_0=0.0202845`$ which, within the numerical precision, is exactly the value predicted by the first step of replica-symmetry breaking as well as by Golay’s ergodicity hypothesis . This coincidence between the replica and the density of metastable states predictions for the ground-state energy occurs also in the random energy model .
A similar study of the symmetrized version of the mean-field Hamiltonian (2), in which $`J_{ij}^k=J_{ji}^k`$, yields exactly the same expression for the exponent $`\alpha `$, see eq.(13), provided that the energy density $`ϵ`$ is replaced by $`ϵ_s/2`$. Hence the symmetrization procedure results in a trivial rescaling the energy densities of the local minima, without affecting their number.
## 3 Energy Barriers and Basin Sizes
The picture that comes out of the replica approach to disordered spin models is that the phase space $`V`$ composed of the $`2^N`$ spin configurations is broken into several valleys connected by saddle points . The relative location and energetic properties of valleys and saddles are expected to determine e.g. the ease with which the ground state can be reached.
It will be convenient to introduce the notion of saddle-point energy $`E[s,w]`$ between two (not necessarily strict) minima $`s`$ and $`w`$. Denoting, for the sake of generality, the energy of state $`x`$ by $`f(x)`$, we can write
$$E[s,w]=\mathrm{min}\left\{\mathrm{max}\left[f(z)|z𝐩\right]|𝐩:\text{path from }s\text{ to }w\right\},$$
(14)
where a path $`𝐩`$ is a sequence of configurations connected by one-spin flips (or, more generally, by moves taken from any desired “move set”). The saddle-point energy $`E[s,w]`$ forms an ultrametric distance measure on the set of local minima, see e.g. . The barrier enclosing a local minimum is the height of the lowest saddle point that gives access to an energetically more favorable minimum. In symbols:
$$B(s)=\mathrm{min}\left\{E[s,w]f(s)|w:f(w)<f(s)\right\}$$
(15)
If $`B(s)=0`$ then the local minimum $`s`$ is marginally stable. It is easy to check that eq.(15) is equivalent to the definition of the depth of local minimum in . It agrees for metastable states with the more general definition of the depth of a “cycle” in the literature on inhomogeneous Markov chains .
The information contained in the energy barriers is conveniently summarized by two global parameters that e.g. determine the convergence behavior of Simulated Annealing and related algorithms. The depth of a landscape is defined as
$$𝖣=\mathrm{max}\left\{B(s)|s\text{ is not a global minimum }\right\}.$$
(16)
It can be shown that Simulated Annealing converges almost surely to a ground state if and only if the cooling schedule $`T_k`$ satisfies $`_{k0}\mathrm{exp}(𝖣/T_k)=\mathrm{}`$ . In order to make the depth comparable between different landscapes we shall consider below the dimensionless parameter $`𝖣/\sigma `$, where $`\sigma ^2`$ is the variance of the energy across the landscape. A related quantity is the (dimensionless) difficulty of the landscape, defined by
$$\psi =\mathrm{max}\left\{\frac{B(s)}{f(s)f(\mathrm{min})}|s\text{ is not a global minimum}\right\}$$
(17)
where $`f(\mathrm{min})`$ is the global energy minimum and the maximum is taken over non-global minima only. It is directly related to the optimal speed of convergence of Simulated Annealing.
Since a direct evaluation of eq.(14) would require the explicit constructions of all possible paths it does not provide a feasible algorithm for determining $`E[s,w]`$ even if $`N`$ is small enough to allow an exhaustive survey of the landscape. The values of $`E[s,w]`$ and $`B(s)`$ can, however, be retrieved from the barrier tree of the landscape. Barrier trees have been considered recently in the context of RNA folding and under the name “disconnectivity graphs” in the protein folding literature . In this contribution we use a modified version of the program barriers, which was developed for the analysis of RNA folding landscapes in . For the sake of completeness we briefly outline the definition of the barrier trees below.
For simplicity let us assume that the energies of any two spin configurations are distinct, i.e., there is a unique ordering of the spin configurations by their energies. The construction of the barrier tree starts from an energy-sorted list of all configurations in the landscape. We will need two lists of valleys throughout the calculation: The global minimum $`x[1]`$ belongs to the first active valley $`V_1`$, while the list of inactive valleys is empty initially. Going through this list of all configurations in the order of increasing energy we have three possibilities for the spin configuration $`x[k]`$ at step $`k`$.
* $`x[k]`$ has neighbors in exactly one of the active valleys $`V_i`$. Then $`x[k]`$ belong to $`V_i`$.
* $`x[k]`$ has no neighbor in any of the (active or inactive) valleys that we have found so far. Then $`x[k]`$ is a local minimum and determines a new active valley $`V_l`$. In the barrier tree $`x[k]`$ becomes a leaf.
* $`x[k]`$ has a neighbor in more than one active valley, say $`\{V_{i_1},V_{i_2},\mathrm{},V_{i_q}\}`$. Then it is a saddle point connecting these active valleys. In the barrier tree $`x[k]`$ becomes an internal node. In this case we add $`x[k]`$ to valley $`V_{i_1}`$ with the lowest energy. Then we copy the configurations of $`V_{i_2},\mathrm{},V_{i_q}`$ to $`V_{i_1}`$. Finally, the status of $`V_{i_2}`$ through $`V_{i_q}`$ is changed from active to inactive. This reflects the fact that from the point of view of a configuration with an energy higher than the saddle point $`x[k]`$, $`V_{i_1},\mathrm{},V_{i_q}`$ appear as a single valley that is subdivided only at lower energy. Consequently, after the highest saddle-point energy has been encountered, all valleys except for the globally optimal $`V_1`$ are in the inactive list.
The outcome of this procedure is a tree such as the one shown in Fig. 2. The leaves correspond to the valleys of the landscape, while the interior nodes denote the saddle points. The tree contains the information on all local minima and their connecting saddle points. Indeed, saddle-point energies, and energy barriers can be immediately read off the barrier trees.
A precise definition of valleys and saddle points in a landscape requires that we take into account the degeneracies in the energy function, i.e., the existence of distinct spin configurations with identical energies, in particular, the presence of neutrality, where neighboring configurations have identical energies . Degeneracies complicate the construction of the barrier tree, since the energy-sorting of the landscape is not unique any more. The simplest remedy is to use the same procedure as above starting from an arbitrary energy sorting. In this case the order of degenerate configurations in the list is arbitrary but fixed throughout the computation. Before proceeding to a configuration with strictly higher energy a simple clean-up step needs to be included in the tree-building algorithm: adjacent valleys with $`E[s,w]=f(s)=f(w)`$ are joined to a single valley. Note that the resulting barrier tree may still contain distinct valleys with the same energy, as the examples in Fig. 3 show. The leaves of the barrier trees are in general valleys which may contain more than one degenerate local minimum.
There is a clear visual difference between the barrier trees for LABSP and the mean-field approximation MF at the one hand, and the $`\pm 1`$ 4-spin Hamiltonian and the REM on the other hand. The main difference appears to be a much larger amount of degeneracy in LABSP/MF, in particular highly degenerate ground states. In fact, it can be shown that the pure Hamiltoninan (1) has many nontrivial symmetries, besides the trivial one where $`x`$ is replaced by $`x`$, which are then responsible for the high degeneracy observed in the tree barrier . Obviously, the disordered Hamiltonian (2) cannot have the same symmetries as the pure one, and so its high degeneracy stems simply from the extreme dilution of the couplings $`J_{ij}^k`$. All models, except REM, are symmetric under replacing $`x`$ by $`x`$, hence all states appear in pairs. We note that the barrier tree of the $`\pm 1`$ 4-spin model is reminiscent of the “funnels” discussed e.g. in protein folding, with a large energy difference between the two global optima and almost all local “traps”. In contrast, the REM shows, as expected, no relationship between energy and nearness of local minimum to the global one.
During the construction of the barrier tree it is easy only to compute the lowest barrier $`B^{}(s)`$ from $`s`$ to a local minimum that comes earlier in the list of configurations, instead of the lowest barrier $`B(s)`$ to a local minimum with strictly smaller energy. Clearly, $`B^{}(s)B(s)`$ since we take the minimum over a few more configurations than prescribed by eq.(15). In case of degenerate landscapes our version of the barriers program calculates $`B^{}(s)`$ which depends on the ordering of the degenerate configurations. We obtain, however, $`B(s)=B^{}(s)`$ for at least one of the valleys at each energy level. The fact that in eq.(16) we are required to maximize in particular over the barriers necessary to escape from any given energy level, however, implies that the values of depth and difficulty can be obtained directly from $`B^{}(s)`$ instead of $`B(s)`$. We note at this point that a modified version $`𝖣^{}𝖣`$ of the depth in which the maximum over all non-global minima is replaced by the maximum over all minima except one global minimum $`x^{}`$ can also be obtained by the simplified procedure above, since it can be shown that $`𝖣^{}`$ is independent of the choice of $`x^{}`$ . The parameter $`𝖣^{}`$ also appears in exact results on the convergence of Simulated Annealing.
Depth and Difficulty are shown in Figure 4 as a function of the number $`N`$ of spins. While there are (moderate) quantitative differences, there does not seem to be a qualitative difference between the LABSP, the mean-field approximation, and the discretized 2 and 4-spin models. Note that all landscapes with the exception of the Gaussian REM have constant scaled depth $`𝖣/\sigma `$, while $`𝖣/\sigma `$ increases linearly with the system size in the REM.
A linear regression of the difficulties yields the slopes $`0.595\pm 6`$, $`0.926\pm 14`$, and $`0.07\pm 2`$ for the mean-field Hamiltonian, the $`\pm 1`$-version of the $`4`$-spin model and the $`\pm 1`$ 2-spin model, respectively. As expected, the difficulty of the quadratic model is much smaller than the difficulty of the 4-spin model. For the sake of clarity, we have omitted the data about the mean difficulty of the REM since it is too large as compared to those shown in the figure. Moreover the width of its difficulty distribution is also so large that the mean value is not physically meaningful.
Additional information on the local minima can be traced during the construction of the barrier tree. We say that a configuration $`x`$ belongs to the basin of the local minimum $`s`$ if $`s`$ is the endpoint of the gradient walk (steepest descent) starting in $`x`$. (Recall that each step of a gradient walk goes to the lowest energy neighbor.) By determining the valley to which the lowest energy neighbor of $`x[k]`$ belongs we may for instance record the basin size of each local minimum. In a landscape without neutrality the gradient walk is uniquely determined by the initial condition, hence the basins form a partition of the set of configurations. We neglect the effects of neutrality in our numerical data by directing the gradient walk to the first possibility in the energy-sorted list of configurations.
Computationally we find, for all models but the REM, that there is an approximate linear relationship between the energy of a local minimum and the logarithm of the size of its gradient-walk basin of attraction, see Figure 5. The fact that the deepest valleys have the largest basins of attraction can be understood as a consequence of the correlation between neighboring spin configurations in all landscapes with the exception of the REM, for which all low-energy minima have essentially the same size of basin of attraction.
## 4 Discussion
The performance evaluation of local search heuristics, in particular Simulated Annealing, in typical instances of optimization problems is a relatively new subject, where the existing criteria for measuring the hardness or difficulty of a problem are still not widely known or accepted, as compared to e.g. the more traditional worst-case analysis. In fact, on the one hand, one expects that the average number of local minima may serve as a measure of the problem hardness, while, on the other hand, one must concede that only local minima separated by high energy barriers are potential traps for the search heuristic. In this paper we combine the concepts of depth and difficulty from the Simulated Annealing literature to the average density of states calculations from the statistical mechanics of disordered systems to obtain a reasonably complete statistical description of the energy landscapes of several classic disordered models. The motivation is to compare the statistical features of these landscapes with the properties of a rather puzzling deterministic problem — the low autocorrelated binary string problem (LABSP) — which has been identified as a particularly hard optimization problem for search heuristics such as Simulated Annealing.
Our results indicate that there is only a quantitative difference between the depths and difficulties, as defined in the Simulated Annealing literature, of all models investigated, with the exception of the random energy model (REM) for which the complete lack of correlations between the energies of neighboring configurations results in a genuine golf-course type landscape. Hence, we have found no evidences of a golf-course like structure in the LABSP landscape, which resembles much more a correlated spin-glass model than the REM. It must be emphasized that although the pure LABSP model (1) may have a glass phase characterized by uncorrelated equilibrium states (at least its mean-field version has such a phase ), the mere existence of this phase is no evidence of a golf-course like structure which, as mentioned above, requires vanishing correlations between the energy values of neighboring spin configurations.
Perhaps the “golf-course” conjecture stems simply from the fact that for large $`N`$ the LABSP is a much more difficult problem for Simulated Annealing than the familiar quadratic spin glass, as shown in Fig. 4. Interestingly, the pairwise comparison between the problems indicates that those problems with the larger number of local minima have also the larger difficulty, the only exception being LABSP and the 4-spin glass model. It should therefore be interesting to use these two problems as a test-bed for validating the hardness criteria proposed in the Simulated Annealing literature.
Finally, our analysis has shown that the disordered, mean-field Hamiltonian $`_d`$, eq.(2), describes surprisingly well the qualitative (e.g. the barrier trees) as well as the quantitative (e.g. number and typical energy of local minima) features of the pure model $``$, eq.(1).
### Acknowledgements
Special thanks to Christoph Flamm for the source code of his program barriers. The work of JFF was supported in part by Conselho Nacional de Desenvolvimento Científico e Tecnológico (CNPq). The work of PFS was supported in part by the Austrian Fonds zur Förderung der Wissenschaftlichen Forschung, Proj. No. 13093-GEN. FFF is supported by Fundação de Amparo à Pesquisa do Estado de São Paulo (FAPESP). We thank FAPESP for supporting PFS’s visit to São Carlos, where part of his work was done.
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# Dephasing times in quantum dots due to elastic LO phonon-carrier collisions
Semiconductor Quantum Dots (QDs) may pave the way to optoelectronic devices with vastly superior performance compared to present devices \[\]. A detailed understanding of QD’s optical properties is therefore of utmost importance, both from a basic science and from a technological point of view. The homogeneous linewidth $`\mathrm{\Gamma }`$ of an optical transition, and the dephasing time, or, equivalently, the polarization relaxation time, $`T_2=2/\mathrm{\Gamma }`$, are basic characteristics for interaction of light with quantum systems. In lasers, the homogeneous linewidth can strongly effect the laser gain. The dephasing time defines the time scale on which coherent interaction of light with medium takes place \[\], and, in particular, it gives the ultimate time scale for realization of coherent control in a quantum system (see, for instance, Ref. \[\] and references therein).
In bulk and Quantum Well (QW) semiconductors the dephasing time is usually taken to be the intraband relaxation time \[\]. At low carrier densities, this time is set by LO phonon-carrier scattering while at higher carrier densities carrier-carrier collisions are also important. The energy spectrum of carriers in QDs is discrete, and in this sense the interaction processes in QDs are qualitatively different from those in QWs or bulk materials. For example, if the LO phonon energy does not coincide with the separation of energy levels in a QD, LO phonon-carrier interaction can not lead to carrier relaxation between these levels (this is the so called “phonon bottleneck” \[\]). At first sight, then, in the absence of this efficient mechanism much longer relaxation and dephasing times could be expected.
Carrier relaxation is a result of inelastic carrier–phonon and carrier–carrier collisions. In contrast, the dephasing time $`T_2`$ is influenced by both inelastic and elastic collisions. Uskov et al. \[\] demonstrated recently that elastic collisions of two-dimensional (2D) carriers with carriers confined in self-assembled QDs can lead to substantial collisional broadening ($`T_2`$ 0.1 – 1 ps at moderate densities of 2D carriers) of QD spectral lines. In this work we study elastic collisions between QD-carriers and LO-phonons: these collisions can disturb the phase of the carrier wave functions in QDs without changing the populations of the carrier energy levels, and accordingly QD spectral lines will be broadened. We will show that these processes can lead to dephasing times $`T_2`$ of some hundreds of femtoseconds for typical QDs at room temperature. These values are of the order of the dephasing times in bulk materials and QWs, and are in accordance with recent room temperature measurements \[\].
Dephasing and broadening of spectral lines in QDs at low temperatures ($`<`$ 50 K) are usually attributed to interactions of QD carriers with acoustic phonons \[\]. In this case, the continuum of acoustic phonons leads to a continuous band of satellites around the Zero Phonon Line (ZPL). This band is considered as homogeneous broadening of the formerly discrete QD spectral line, and the dephasing time is calculated from the width of the band. The present paper considers broadening and dephasing via a polar coupling to LO phonons with fixed frequency. The broadening of the ZPL by LO phonons is interpreted as follows. QD carriers can virtually absorb and then emit (or emit and then absorb) LO phonons. This second order process does not change the final energy of the carrier, but it does change the phase of the carrier wave function. This change of the phase, in turn, implies dephasing of the dipole for the considered optical transition, and accordingly a broadening of the spectral line \[\].
The scattering against LO phonons gives rise to fluctuations $`\mathrm{\Delta }E(t)`$ in the transition energy between the electron and hole QD levels: $`E(t)=E^{(0)}+E_s+\mathrm{\Delta }E(t)`$, where $`E^{(0)}`$ is the unperturbed transition energy, and $`E_s`$ is the energy-shift due to the LO phonons. For simplicity, we ignore the Coulombic electron-hole interaction. In a classical treatment the line structure function of the optical transition is given by the Fourier transform of the correlation function $`\psi _f(\tau )=\overline{f(t+\tau )f^{}(t)}`$ for the oscillator $`f(t)`$, (here the overline indicates an average over fluctuations)
$$f(t)\mathrm{exp}\left\{(i/\mathrm{})\left[(E^{(0)}+E_s)t+_0^t𝑑t_1\mathrm{\Delta }E(t_1)\right]\right\}.$$
(1)
The correlation function $`\psi _f(\tau )`$ can be expressed as \[\]
$`\psi _f(\tau )`$ $``$ $`\mathrm{exp}[i(E^{(0)}+E_s)(\tau /\mathrm{})]`$ (2)
$`\times `$ $`\mathrm{exp}\left[{\displaystyle \frac{1}{2\mathrm{}^2}}{\displaystyle _0^\tau }𝑑t_1{\displaystyle _0^\tau }𝑑t_2\overline{\mathrm{\Delta }E(t_1)\mathrm{\Delta }E(t_2)}\right].`$ (3)
In this work we must use the quantum mechanical analog of Eq. (1): the fluctuations in the transition energy are due to interactions with optical phonons, which obey Bose statistics. Thus, instead of $`\psi _f(\tau )`$ we consider \[\]
$$A(\tau )=\widehat{T}\mathrm{exp}\left\{(i/\mathrm{})\left[E^{(0)}\tau +_t^{t+\tau }𝑑t^{}\widehat{H}(t^{})\right]\right\},$$
(4)
where $`\widehat{T}`$ is the time-ordering operator and $`\widehat{H}(t)`$ is the effective interaction Hamiltonian describing the fluctuations in the transition energy. $`\mathrm{}`$ indicates an average over the phonon ensemble. The result of our calculations, to be described below, is that the lineshape function, which is just the Fourier transform of (4), can be expressed in a compact form:
$`A(\omega )`$ $`=`$ $`{\displaystyle 𝑑\tau e^{i(\mathrm{}\omega E^{(0)}E_s)(\tau /\mathrm{})}}`$ (5)
$`\times `$ $`\mathrm{exp}\left\{{\displaystyle \frac{\sigma ^2}{2\gamma _{LO}}}[|\tau |+{\displaystyle \frac{1}{2\gamma _{LO}}}(e^{2\gamma _{LO}|\tau |}1)]\right\}.`$ (6)
The shift $`E_s`$ and the parameter $`\mathrm{}\sigma `$, which has the meaning of average quadratic fluctuation of the spectral line, depend on temperature, phonon interaction mechanism, and quantum dot geometry, and explicit expressions for them are given below. $`\tau _{LO}=1/2\gamma _{LO}`$ is the life-time of optical phonons, which is finite because of various interaction mechanisms (such as phonon-phonon scattering or boundary scattering), and it enters our theory as a parameter which must be calculated separately, or extracted from experiment. If $`\sigma 2\gamma _{LO}`$, $`A(\omega )`$ has a Lorentzian line form with FWHM $`\mathrm{\Gamma }=\sigma ^2/\gamma _{LO}`$, while if $`\sigma 2\gamma _{LO}`$ Eq. (6) gives rise to a Gaussian line with width $`\mathrm{\Gamma }=2\sqrt{2\mathrm{ln}2}\sigma `$.
In order to arrive at Eq. (6) the following steps were needed. First, we derive an effective Hamiltonian describing the interaction between quantum dot charge carriers and optical phonons. Next, the time-ordered expectation value (4) must be evaluated. Finally, the developed formalism will be applied to a concrete model of a quantum dot, and the temperature dependence of $`\mathrm{\Gamma }`$ is computed.
The carrier-phonon interaction is ($`x=e`$ for electrons, $`x=h`$ for holes)
$$\widehat{U}^x(𝐫,t)=\underset{𝐪}{}C_q^x[a_𝐪(t)e^{i𝐪𝐫}+\mathrm{h}.\mathrm{c}.],$$
(7)
where $`a_𝐪`$ is the annihilation operator for the LO phonon with the wave vector $`𝐪`$. For polar carrier–phonon coupling $`C_q^e=C_q^h`$. The finite life-time of the optical phonons gives rise to the following expression for the phonon Green function \[\]:
$`D(𝐪,\tau )`$ $``$ $`\widehat{T}\{[a_𝐪(\tau )+a_𝐪^{}(\tau )][a_𝐪(0)+a_𝐪^{}(0)]\}`$ (8)
$`=`$ $`e^{\gamma _{LO}|\tau |}\left[(\overline{n}+1)e^{i\omega _{LO}|\tau |}+\overline{n}e^{i\omega _{LO}|\tau |}\right]`$ (9)
where $`\overline{n}=1/[\mathrm{exp}(\mathrm{}\omega _{LO}/k_BT)1]`$, and $`\omega _{LO}`$ is the LO phonon frequency. Assuming that the phonon energy $`\mathrm{}\omega _{LO}`$ is much less than the energy separation between the QD levels, an effective Hamiltonian $`\widehat{H}(t)`$ can be derived by neglecting real transitions beween the levels \[\], and applying perturbation theory in the carrier–phonon interaction. Up to second order one finds
$`\widehat{H}^{(1)}(t)`$ $`=`$ $`{\displaystyle \underset{𝐪}{}}\left[f(𝐪)a_𝐪(t)+f^{}(𝐪)a_𝐪^{}(t)\right],`$ (10)
$`\widehat{H}^{(2)}(t)`$ $`=`$ $`{\displaystyle \underset{\mathrm{𝐪𝐪}^{}}{}}[a_𝐪(t)a_𝐪^{}(t)F(𝐪,𝐪^{})`$ (11)
$`+a_𝐪(t)a_𝐪^{}^{}(t)F(𝐪,𝐪^{})+a_𝐪^{}(t)a_𝐪^{}(t)F(𝐪,𝐪^{})`$ (12)
$`+a_𝐪^{}(t)a_𝐪^{}^{}(t)F(𝐪,𝐪^{})].`$ (13)
Here
$`f(𝐪)`$ $`=`$ $`C_q^eM_𝐪^{e1}+C_q^hM_𝐪^{h1},`$ (14)
$`F(𝐪,𝐪^{})`$ $`=`$ $`|C_qC_q^{}|{\displaystyle \underset{p>1}{}}\left[{\displaystyle \frac{M_𝐪^{ep}M_{𝐪^{}}^{ep}{}_{}{}^{}}{E_{ep}^0E_{e1}^0}}+{\displaystyle \frac{M_𝐪^{hp}M_{𝐪^{}}^{hp}{}_{}{}^{}}{E_{hp}^0E_{h1}^0}}\right].`$ (15)
$`E_{ep}^0(E_{hp}^0)`$ is the unperturbed electron (hole) energy with the wave function $`\psi _{ep}`$ ($`\psi _{hp}`$), which enters via the matrix element
$$M_𝐪^{xp}=𝑑𝐫\psi _{xp}^{}(𝐫)\mathrm{exp}(i𝐪𝐫)\psi _{x1}(𝐫).$$
(16)
We consider broadening of the transition between electron and hole ground states ($`p=1`$), so that $`E^{(0)}=E_g+E_{e1}^0+E_{p1}^0`$, where $`E_g`$ is the energy gap of the QD material.
The first order term Eq. (10) is just the standard Huang–Rhys Hamiltonian \[\], while the second order term (13) is an extension of this theory. A similar second order term has earlier been shown to lead to broadening of the ZPL in case of impurities in doped crystals \[\]. Below we show that the Hamiltonian (10) – (13), which only involves virtual transitions between the QD levels, leads to a broadening of the ZPL in QD’s in accordance with experiments \[\]. The existence of the nearby levels is essential for obtaining the broadening. The effects of mixing of the QD levels due to transitions between the levels on intensities of phonon satellites have recently been considered in \[\].
The evaluation of (4) with the effective Hamiltonian (10) – (13). can be carried out exactly using the cumulant technique \[\]. The result is (note that cross-terms vanish because they involve unequal number of phonons)
$`A(\tau )`$ $`=`$ $`\mathrm{exp}\{i[E^{(0)}+{\displaystyle \underset{j=1,2}{}}\widehat{H}^{(j)}](\tau /\mathrm{})`$ (17)
$`{\displaystyle \frac{1}{2\mathrm{}^2}}{\displaystyle _0^\tau }dt_1{\displaystyle _0^\tau }dt_2{\displaystyle \underset{j=1,2}{}}\widehat{T}\mathrm{\Delta }\widehat{H}^{(j)}(t_1)\mathrm{\Delta }\widehat{H}^{(j)}(t_2)\},`$ (18)
where $`\mathrm{\Delta }\widehat{H}^{(j)}(t)=\widehat{H}^{(j)}(t)\widehat{H}^{(j)}`$. It is interesting to note the formal similarity with the classical result Eq. (3). The phonon averages can be carried out with the help of Wick’s theorem, which stipulates that the various operator averages implicit in (17) are given as the sum of all pairwise averages. The time-ordered expectation value due to $`\widehat{H}^{(1)}`$ is just $`_0^t𝑑t_1_0^t𝑑t_2_𝐪|f(𝐪)|^2D(𝐪,t_1t_2)`$; this expression is well-known, say, from studies of phonon-assisted tunneling in resonant-tunneling diodes \[\]. In our case this term gives the Huang–Rhys shift $`E_s^{(1)}=S\mathrm{}\omega _{LO}`$ of the spectral line, where $`S=_𝐪|f(𝐪)/(\mathrm{}\omega _{LO})|^2`$ is the Huang–Rhys parameter \[\], as well as the satellites at energies $`\pm n\mathrm{}\omega _{LO}`$ from the main line. For the present purposes it is important to realize that the first-order term $`\widehat{H}^{(1)}`$ cannot lead to broadening of the ZPL if $`\gamma _{LO}/\omega _{LO}1`$, see also Ref.\[\]. A straightforward calculation shows that the contribution due to $`\widehat{H}^{(2)}`$ is obtained by a replacement $`|f(𝐪)|^2D(t_1t_2)F(𝐪,𝐪^{})[F(𝐪^{},𝐪)+F(𝐪,𝐪^{})]D(𝐪,t_1t_2)D(𝐪^{},t_1t_2)`$. This expression generates two kinds of terms. (i) Terms proportional to $`\mathrm{exp}[\pm 2i\omega _{LO}(t_1t_2)]`$; together with $`\widehat{H}^{(2)}`$ in (17) they give the spectral shift in the second-order theory. They are also relevant for the two–phonon satellites at $`\pm n2\mathrm{}\omega _{LO}`$ of the main line, which however is not our focus here. (ii) Terms without these fast time-oscillations which determine the zero-phonon linewidth. Collecting all the terms describing the spectral shift and broadening, we find
$`E_s^{(2)}`$ $`=`$ $`(2\overline{n}+1)(S_1+S_2)\mathrm{}\omega _{LO},`$ (19)
$`\sigma `$ $`=`$ $`\sqrt{4\overline{n}(\overline{n}+1)S_2}\omega _{LO},`$ (20)
where
$`S_1`$ $`=`$ $`{\displaystyle \underset{𝐪}{}}F(𝐪,𝐪)/(\mathrm{}\omega _{LO}),`$ (21)
$`S_2`$ $`=`$ $`{\displaystyle \underset{\mathrm{𝐪𝐪}^{}}{}}F(𝐪,𝐪^{})[F(𝐪^{},𝐪)+F(𝐪,𝐪^{})]/2(\mathrm{}\omega _{LO})^2.`$ (22)
The second-order shift (19) together with the first-order shift $`E_s^{(1)}`$ gives the total shift $`E_s`$ of the ZPL in (6). The parameters $`S_1`$ and $`S_2`$ are extensions of the parameter $`S`$ of the Huang–Rhys theory. Our numerical calculations show that $`S_1,S_21`$, and also $`S_1S_2`$. Eq. (20) together with (6) allows an analysis of the spectral lineshape, and, in particular, its width as a function of temperature.
As a practical application, we consider a cylindrical quantum dot of radius $`R`$ and height $`H`$, with infinite confining potential. In the calculation of the average fluctuation $`\sigma `$, we include in Eq. (15) only the leading terms in the sum over the states $`p`$ . If $`2R>H`$, two states have the same energy, and are the closest to the ground state. Taking into account this two-fold degeneracy, we get
$$\mathrm{}\sigma =\sqrt{\overline{n}(\overline{n}+1)}\frac{\mathrm{}\omega _{LO}}{\mathrm{\Delta }\stackrel{~}{E}}g(H/R)\frac{K}{R}$$
(23)
where $`\mathrm{\Delta }\stackrel{~}{E}=\mathrm{}^2/(2(m_e+m_h)R^2)[\alpha _{11}^2\alpha _{10}^2]`$. $`\alpha _{lm}`$ is the $`l`$-th root of the Bessel function of the order $`m`$. The material parameter $`K=e^2(1/ϵ_{\mathrm{}}1/ϵ)/(\pi ^2ϵ_0)`$ has the value 26.0 meV$``$nm for GaAs. The dimensionless function $`g(H/R)`$ is shown in the inset of Fig. 1.
In the numerical calculation of the width $`\mathrm{\Gamma }`$ (see Fig. 1) we use $`\tau _{LO}=1.5`$ ps \[\] so that $`\mathrm{}\gamma _{LO}=0.22`$ meV and $`(\omega _{LO}\tau _{LO})^10.01`$. From (23) one finds that at $`T=300`$ K the average quadratic fluctuation is $`\mathrm{}\sigma 2.2`$ meV (corresponding to $`S_2210^3`$ in Eq. (20)), so that the condition $`\sigma 2\gamma _{LO}`$ is satisfied, and consequently the spectral line has a Gaussian shape. At lower temperatures $`\sigma `$ decreases and can reach the regime $`\sigma 2\gamma _{LO}`$, where the Lorentzian shape obtains. At $`T=`$ 300 K the linewidth $`\mathrm{}\mathrm{\Gamma }`$ is 5.1 meV, corresponding to $`T_2250`$ fs. Taking into account higher excited states in the sum (15) in evaluating $`\sigma `$ typically reduces $`T_2`$ by a factor of $`1.52`$. Thus, one can expect room temperature dephasing times of the order of $`100200`$ fs. This is in good agreement with experimental values ($`200300`$ fs) \[\]. The LO-phonon contribution to the linewidth is proportional to $`\overline{n}(\overline{n}+1)`$ if $`\mathrm{\Gamma }1/\tau _{LO}`$, and thus vanishes as the temperature tends to zero. For example, for $`T<50`$ K the calculated values of $`\mathrm{}\mathrm{\Gamma }`$ are below 20 $`\mu `$eV (see Figure 1), which is below the linewidths found in experiment ($`50150`$ $`\mu `$eV) \[\]. On the other hand, the acoustic phonon related linewidth does not vanish in this limit \[\], and we conclude that the LO-mechanism is dominant only at elevated temperatures, say $`T>`$ 100 K. The present broadening mechanism also provides a possible way of understanding fast carrier relaxation in QDs at elevated temperatures \[\]: due to the substantial broading of the levels their differences no longer have to match strictly the LO-phonon energy, thereby weakening the phonon bottleneck effect.
It is also instructive to consider the classical limit of our calculations. Then $`\overline{n}`$ 1, and thus for $`\mathrm{\Gamma }1/\tau _{LO}`$ we find $`\mathrm{\Gamma }\overline{n}^2`$. This result also follows from Eq. (3) if we treat the phonons as classical random fields perturbing the QD carriers. This classical result allows one to establish an analogy with broadening in atomic gases \[\]. There, elastic atomic collisions lead to dephasing of optical transitions, and the homogeneous linewidth in the Lorentzian limit is proportional to the density of the colliding particles. In our case the broadening of ZPL is the result of elastic second order interactions between phonons and the carriers, and hence $`\mathrm{\Gamma }\overline{n}^2`$. Note also that the LO phonon lifetime $`\tau _{LO}`$ plays the same role as the collisional time in atomic gases \[\].
We finally consider the effect of the shape and size of the QD on the dephasing time. Eq. (23) implies that $`\sigma L`$, where $`L`$ characterizes the size of the QD. In the Lorentz case $`\mathrm{\Gamma }\sigma ^2L^2`$, while in the Gaussian case $`\mathrm{\Gamma }\sigma L`$. From Fig. 1 we can see that $`\mathrm{\Gamma }`$ can be modified by a factor 2–3 by changing the ratio $`H/R`$ of the cylinder. Note also that $`\sigma `$ and $`\mathrm{\Gamma }`$ depend on the carrier masses, $`\sigma (m_e+m_h)`$. These characteristic dependencies point towards a possibility to engineer the spectral linewidths.
In conclusion, we have developed a theory for dephasing of an optical transition in QDs due to second-order elastic interaction with LO-phonons. The theory results in expressions which are easily evaluated numerically, and can also be used to draw qualitative conclusions about the dependence of the dephasing time on physical parameters defining the QD. At room temperature, despite of an apparent phonon bottleneck, the considered mechanism leads to dephasing times of the order of 200 fs, i.e. of the same order as in bulk, and also as observed in QDs.
One of us (A. V. U.) thanks Ministry of Education, Science, Culture and Sport of Japan for support during his work at TUAT. We acknowledge D. Birkedal, P. Borri, and J. Hvam for useful comments.
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# The Theory of Gravitation in the Space - Time with Fractal Dimensions and Modified Lorents Transformations
## I Introduction
The general relativity theory (GR) is one of the most known and used the theories of gravitation field. It is elegant, beautiful from physical point of view and explains all experimental facts (truly, it must be point out that others gravitational theories explain them too, as example see ). The one of the known lacks of the GR consists in impossibility to include other physical fields (except gravitation field) in the frame of GR, i.e. in impossibility to look on other fields as on the characteristics of metric of Riemann geometry. The equations of GR based on assumption that all systems of frame are equivalent, the absolute systems of references are absent thou the space and the time in the presence of gravitational fields are inhomogeneous. In the fractal theory of time and space ( see -) all the physical fields are included in the fractional dimensions of the time and the space, the time and the space fields are real fields and any system of frames are an absolute system of frame, the Lorents transformations and known physical laws fulfilled as not rigorously laws and are only very good approach for describing characteristics of the time and the space. For the case when fractional corrections to integer dimensions are small all the equations of the fractal theory practically coincide with equations of known physical theories or may give only non-sufficient corrections. Here we present the theory for the fields with spin equal two in the fractal time and space on the base of the equations given in ,, ( in the paper used the modified Lorents transformations with corrections given by the fractal theory of time for moving with speed of light(-) for constructing the theories of scalar, vector and spinor field). The equation of this theory differs from equations of the theories , because the modified Lorents transformations give four systems of gravitational equations for four different gravitational fields (two with real energies and two with imaginary energies) instead of only one system in theories , .
## II Generalized fractional derivatives
Following -, we will consider both time and space as the initial real material fields existing in the world and generating all other physical fields by means of their fractional dimensions. For the gravitational fields equations construction may be used the principle of minimum fractal dimensions functional as it was made in -. The aim of the paper is to include (in addition to Lagrangians used in , ) in Lagrangians of fields the additional members that give corrections to Lorents transformations ,, in the domain of velocities $`vc`$. For describing the functions defined on multifractal sets it is necessary to introduce the generalized fractional derivatives (see , , ). Therefore, we introduce as in the cited papers the integral functionals (both left-sided and right-sided) which are suitable to describe the dynamics of functions defined on multifractal sets of time and space (generalized fractional derivatives (GFD), see -, ) and replace by GFD the usual derivatives and integral respect to time and space coordinates in the fractional dimensions functional. These functionals GFD are simple and natural generalization of the Riemann-Liouville fractional derivatives and integrals:
$$D_{+,t}^df(t)=\left(\frac{d}{dt}\right)^n_a^t\frac{f(t^{})dt^{}}{\mathrm{\Gamma }(nd(t^{}))(tt^{})^{d(t^{})n+1}}$$
(1)
$$D_{,t}^df(t)=(1)^n\left(\frac{d}{dt}\right)^n_t^b\frac{f(t^{})dt^{}}{\mathrm{\Gamma }(nd(t^{}))(t^{}t)^{d(t^{})n+1}}$$
(2)
where $`\mathrm{\Gamma }(x)`$ is Euler’s gamma function, and $`a`$ and $`b`$ are some constants from $`[0,\mathrm{})`$. In these definitions, as usually, $`n=\{d\}+1`$ ,where $`\{d\}`$ is the integer part of $`d`$ if $`d0`$ (i.e. $`n1d<n`$) and $`n=0`$ for $`d<0`$. If $`d=const`$, the generalized fractional derivatives (GFD) (1)-(2) coincide with the Riemann - Liouville fractional derivatives ($`d0`$) or fractional integrals ($`d<0`$). When $`d=n+\epsilon (t),\epsilon (t)0`$, GFD can be represented by means of integer derivatives and integrals. There are relations between GFD and ordinary derivatives for $`d_\alpha `$ near integer values. If $`d_\alpha `$ $`n`$ where $`n`$ is an integer , ( for example $`d_\alpha `$=$`1+\epsilon (𝐫,(t),t)`$, $`\alpha =𝐫,t)`$, in that case it is possible represent GFD by approximate relations (see ,)
$$D_{+,x_\alpha }^{1+ϵ}f(𝐫(t),t)=\frac{}{x_\alpha }f(𝐫(t),t)+\frac{}{x_\alpha }[\epsilon (𝐫(t),t)f(𝐫(t),t)]$$
(3)
For $`n=1`$, i.e. $`d=1+\epsilon `$, $`\left|\epsilon \right|<<1`$ it is possible to obtain:
$$D_{+,t}^{1+\epsilon }f(t)\frac{}{t}f(t)+a\frac{}{t}\left[\epsilon (𝐫(t),t)f(t)\right]$$
(4)
where $`a`$ is constant and defined by the choice of the rules of regularization of integrals (1)-(2) (for more detail see ,, ). The selection of the rule of regularization that gives a real additives for usual derivative in (3) yield $`a=0.5`$ for $`d<1`$ and $`a=1.077`$ for $`d>1`$ . The functions under integral sign in (1)-(2) we consider as the generalized functions defined on the set of the finite functions . The notions of GFD, similar to (1)-(2), also defined and for the space variables $`𝐫`$.
In the definitions of GFD (1)-(2) the connections between fractal dimensions of time $`d_t(𝐫(t),t)`$ and characteristics of physical fields (say, potentials $`\mathrm{\Phi }_i(𝐫(t),t),i=1,2,..)`$ or densities of Lagrangians $`L_i`$) are determined, following , by the relation
$$d_t(𝐫(t),t)=1+\underset{i}{}\beta _iL_i(\mathrm{\Phi }_i(𝐫(t),t))$$
(5)
where $`L_i`$ are densities of energy of physical fields, $`\beta _i`$ are dimensional constants with physical dimension of $`[L_i]^1`$ (it is worth to choose $`\beta _i^{}`$ in the form $`\beta _i^{}=a^1\beta _i`$ for the sake of independence from regularization constant). The definition of time as the system of subsets and definition the FD $`d`$ (see ( 5)) connects the value of fractional (fractal) dimension $`d_t(r(t),t)`$ with each time instant $`t`$. Thus $`d_t`$ depends both on time $`t`$ and coordinates $`𝐫`$. If $`d_t=1`$ (the absence of physical fields) the set of time has topological dimension equal to unity. The multifractal model of time allows, as will be shown early (, ) , to consider the divergence of energy of masses moving with speed of light in the special relativity theory as the result of the requirement of rigorous validity of the conservation laws in the presence of physical fields that is valid only for closed systems. In our theory there are an approximate fulfillment of conservation laws as in the fractal theory of time and space the Universe is treated as an open system defined on the measure carrier (the closed system is the Universe together with the measure carrier).
The gravitational equations may be received by using the principle of minimum to functional of fractal dimensions with dependencies of GFD (, ) or by replacing the ordinary derivatives in proper physical equations by GFD () (the results will be the same). In this paper we generalize the theory of gravitational fields , , by including in the equations received in these papers the results of paper that took into account the modified Lorents transformations.
## III The based equations for physical fields with spin equal two
For generalization of the gravitational field theory presented in , by means of construction the equations in the fractal time and space with modified Lorents transformations we write at first the field equations for scalar function $`\mathrm{\Phi }`$ of paper
$`(\mathrm{}^24a_0^2{\displaystyle \frac{^4}{t^4}})\mathrm{\Phi }(𝐫,t)=E_0^4\mathrm{\Phi }(𝐫,t)`$ (6)
where $`\mathrm{}`$ is D’Alamber operator ($`\mathrm{}=\mathrm{\Delta }\frac{^2}{t^2}`$, $`\mathrm{\Delta }`$ is Laplasian), $`\mathrm{\Phi }`$ are functions describing particles or fields. For scalar $`\mathrm{\Phi }`$ equation (6) describes the scalar field in the space with fractal dimensions of time that originate the all physical fields ($`a_00)`$ . The corrections in (6) to the usual D’Alamber equation are the result of modifying the Lorentz transformations. The last is consequences of fractal nature of time. Now we may write the equations with taken into account both phenomenon: the influences of multifractal structure of time ( use in the equations the generalized Riemann-Liouville fractional derivatives (GFD) instead of ordinary derivatives) and corrections to equations from modified Lorents transformations received in . In that case the equation (6) take the form
$`(D_{,t}^{d_t}D_{+,t}^{d_t}`$ $``$ $`\mathrm{\Delta })^2\mathrm{\Phi }(𝐫,t)=[E_0^4+`$ (7)
$`+`$ $`4a_0^2(D_{,t}^{d_t}D_{+,t}^{d_t})^2]\mathrm{\Phi }(𝐫,t)`$ (8)
where functionals $`D_{+,t}^d`$ and $`D_{,t}^d`$ defined by (1), (2) and (5). It is useful to receive from these equations of the fourth order (7) the four equations of the second order. It is possible if use the Dirac type four-component matrices $`\alpha _i`$ where $`(i=1,2,3,4):\alpha _i\alpha _j+\alpha _j\alpha _i=\delta _{ij}`$. Than we have four equations of second order for the fields both with real energies (two equations) and with imaginary energies (two equations):
$`(D_{,t}^{d_t}D_{+,t}^{d_t}`$ $``$ $`\mathrm{\Delta })I\mathrm{\Phi }_i(𝐫,t)=[\alpha _1E_0^2+`$ (9)
$`+`$ $`2a_0\alpha _2(D_{,t}^{d_t}D_{+,t}^{d_t})]\mathrm{\Phi }_i(𝐫,t)`$ (10)
where $`\alpha _1`$ and $`\alpha _2`$ may be chosen as in
$`\alpha _1=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right),\alpha _2=\left(\begin{array}{cccc}0& 0& 0& 1\\ 0& 0& 1& 0\\ 0& 1& 0& 0\\ 1& 0& 0& 0\end{array}\right)`$ (19)
and $`I`$ is the 4-component unit matrix. This situation is the same as for vector, spinor or Proca fields described in the .
## IV The selection of the models of gravity
The generalization of gravitational theory in fractal space on the base of equations (9) is possible by using the two models of describing the gravity fields :
a) The first model based on Einstein representations the space-time as the continuous set of points described by Riemannian geometry. In that case in the fractal theory of time and space the measure carrier must be defined as the Riemannian sets with integer dimensions. On this sets we construct the fractal sets of time and space with dimensions defined by Lagrangians densities of energy for all physical fields (see (5) ). On the fractal sets the laws of physics will broken because of the sets of space and time become the open systems (see statistical theory of open system in \[klim\] ) connected with the measure carrier. An arbitrary theories for any physical fields will include (in the fractal space) the influences of Riemannian geometry because of Riemannian carrier of measure. So the Riemannian geometry of fields will be not consequences of characteristics of fields (for example, gravitational field with spin equal two), but characteristics of measure carrier. In that case we may describe the gravitational field (if take into account the influences of all physical fields on behaviour of gravitational field) by using the principle of minimum of fractal dimensions functional and Euler equations with generalized fractional derivatives (GFD) as we introduced in the equations of paper , , . In this model it is necessary to use the covariant derivatives in the fractal Riemannian time-space as it was made in , and improve them by introducing the four equations for the gravitational field tensor $`\mathrm{\Phi }^{\mu \nu }`$ and introducing the corrections from modified Lorents transformations (the introducing of the last corrections was demonstrated above). In that case the new gravitational equations will describe four fields: the two gravitational fields with real energies (if use the analogy with Dirac theory these fields may have different sign of the gravitational charges) and the two gravitational fields with imaginary energies . All these fields exist in the Riemann time-space with fractal dimensions. In the case of integer dimensions of time and space the received equations (and the theory) coincide with equations of ordinary GR. In this paper we consider also the another gravity model with more analogy to Logunov-Mestvirichvili model of gravity .
b) The second model for describing the gravitation fields in the fractal time and space ( by GFD using ) consists in the selection of other the measure carrier. The selection the measure carrier is: the measure carrier selected as the flat four-dimensions pseudo Euclidean Minkowski time-space. Our fractal Universe in that case defined as the multifractal sets on the pseudo Euclidean Minkowski time-space ( the model of a measure carrier selection in the case a) was the model of Riemannian time-space). Let us select (as a base) the system of reference which coincide for FD equal to unit with Cartesian system of reference (we remind that in the fractal theory of time and space there are only an absolute systems of reference but if FD of time and space near integer the principle of equivalence of all reference systems is valid with grate exactness ). The equations of the gravitation fields in that case will be similar to the equations of the theory in which all derivatives replaced on GFD and metric tensor $`\gamma ^{\mu \nu }`$ are contained the functions (functionals) originated by fractional dimensions (i.e. it must be the function of $`L`$ where $`L`$ is Lagrangians energy densities of gravitation fields). Beside these corrections must be taken into account the corrections (the main corrections) from modifying of Lorents transformations and the presence as result of it of four sorts of the gravitational fields (two with real and two with imaginary energies). The differences the theory of gravitation based on above statements from the theory in that case are: 1) the real gravitational fields in the theory originated by fractional dimensions of time, but not postulated as in ; 2) the ordinary derivatives replaced by GFD for taking into account the FD of time; 3) we took into account also the modification of Lorents transformations for time with fractional dimensions; 4) the time-space is fractal and only the measure carrier is Minkowski time-space with integer dimensions..
## V The gravitational equations defined on Minkowski time-space measure carrier
It is convenient to use the designations of the theory for the equations construction of gravitation fractal theory in the multifractal Universe defined on the pseudo Euclidean Minkowski time-space (this space may have any dimensions but integer). The equations for gravitation field tensor $`\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }=\sqrt{\gamma }\mathrm{\Phi }^{\mu \nu }`$ ($`\gamma =det(\gamma _{\mu \nu })`$, $`\stackrel{~}{t}^{\mu \nu }`$= $`\sqrt{\gamma }t^{\mu \nu }`$, $`L`$ \- is a Lagrangians density of physical fields (see in details , )) have form
$`\gamma ^{\alpha \beta }D_{,\alpha }^{d_i}D_{+,\beta }^{d_i}I\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }=\alpha _1b^2\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }+\lambda \stackrel{~}{t}^{\mu \nu }(\gamma ^{\mu \nu },\mathrm{\Phi }_A)+`$ (20)
$`+2a_0\gamma ^{44}D_{,t}^{d_i}D_{+,t}^{d_i}\alpha _2\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }`$ (21)
$`\stackrel{~}{t}^{\mu \nu }=2{\displaystyle \frac{\delta L}{\delta \gamma _{\mu \nu }}}`$ (22)
$$D_{\pm ,\mu }^{d_i}\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }=0$$
(23)
In equations (20) the tensors of gravitational fields $`\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }`$ included in the form of column
$`\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }=\left(\begin{array}{c}\stackrel{~}{\mathrm{\Phi }}_1^{\mu \nu }\\ \stackrel{~}{\mathrm{\Phi }_2}^{\mu \nu }\\ \stackrel{~}{\mathrm{\Phi }_3}^{\mu \nu }\\ \stackrel{~}{\mathrm{\Phi }_4}^{\mu \nu }\end{array}\right)`$ (28)
The metric tensor $`\gamma ^{\alpha \beta }`$ is a function (or functional) of the tensors of gravitational fields $`\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }`$ and defined on the fractal pseudo Euclidean Minkowski time-space. These dependencies the $`\gamma ^{\alpha \beta }`$ from$`\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }`$ originated by dependencies the interval $`dS^2`$ from $`\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }`$ because the last in the fractal Minkowski space has the complicated form (see ) and may be expand in powers of $`\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }`$. For expansion $`\gamma ^{\mu \nu }`$ in powers of $`\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }`$ obtain
$`\gamma ^{\mu \nu }(\stackrel{~}{\mathrm{\Phi }}^{\mu \nu },D_{+,t}^{d_t}\stackrel{~}{\mathrm{\Phi }}^{\mu \nu },..)=\gamma ^{\mu \nu }+`$ (29)
$`+{\displaystyle A_{\alpha \beta }^{\mu \nu }\stackrel{~}{\mathrm{\Phi }}^{\alpha \beta }}+\mathrm{}`$ (30)
were $`A_{\alpha \beta }^{\mu \nu }`$ are coefficients of expansion and depend at coordinate and time. So if it is possible when for large distances from center of gravity $`r_0`$ ($`r_0<<r`$) to limit oneself by two first members of (29) we may write
$`\gamma ^{\mu \nu }(\stackrel{~}{\mathrm{\Phi }}^{\alpha \beta })\gamma ^{\mu \nu }+\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }`$ (31)
The equation (23) describes the boundary conditions for $`\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }`$ on the Universe surface and $`b`$ is a mass of graviton and play role of parameter expanding the domain of existence of GFD.
## VI Gravitational fields defined on the Riemann space measure carrier
As the carrier of a measure is the Riemann space with an integer dimensions we obtain the determination for covariant derivatives in Riemann space with fractional dimensions
$$D_{\pm ,\alpha }^{d_i}t^{\mu \nu }=D_{\pm ,\alpha }^{d_i}t^{\mu \nu }+\gamma _{\alpha \beta }^\nu t^{\mu \beta }i=t,r$$
(32)
where $`t^{\mu \nu }`$ is the energy-momentum tensor and $`\gamma ^{\mu \nu }`$ is the metric tensor of the Riemann ”four-dimension space with fractional dimensions”, $`D_{\pm ,\alpha }^{d_i}`$ are GFD, $`\gamma _{\alpha \beta }^\nu `$ are Christoffel symbols
$$\gamma _{\alpha \beta }^\nu =\frac{1}{2}\gamma ^{\nu \sigma }(D_{\pm ,\alpha }^{d_i}\gamma _{\beta \sigma }+D_{\pm ,\beta }^{d_i}\gamma _{\alpha \sigma }+D_{\pm ,\sigma }^{d_i}\gamma _{\alpha \beta })$$
(33)
The equations for gravitation field tensor $`\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }`$than read
$`\gamma ^{\alpha \beta }D_{,\alpha }^{\nu ,d_i}D_{+,\beta }^{\nu ,d_i}I\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }=b^2\alpha _1\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }+\lambda \stackrel{~}{t}^{\mu \nu }(\gamma ^{\mu \nu },\mathrm{\Phi }_A)+`$ (34)
$`+2a_0D_{,+}^{d_t}D_{+,t}^{d_t}\alpha 2\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }`$ (35)
where $`b`$ is a constant value that necessary to introduce for using more broad sets of functions with GFD and it after calculations must be put zero. The $`\mathrm{\Phi }^{\mu \nu }`$ is a four column matrix. So we have again four equations for gravitational fields with real and imaginary energies. The equation for curvature tensor (with GFD ) have an usual form, but it will be four equations for different the curvature tensors $`R(i),i=1,2,3,4`$ and necessary take into account corrections in the covariant derivatives from fractal nature of space and of modifying Lorents transformations
$$R_i^{\mu \nu }\frac{1}{2}\gamma ^{\mu \nu }R_i=\frac{8\pi }{\sqrt{g}}T_i^{\mu \nu }$$
(36)
$$D_{\pm ,\mu }^{d_i}\stackrel{~}{g}_i^{\mu \nu }=0$$
(37)
The equation (37) describes the boundary conditions for $`g^{\mu \nu }`$ on the Universe surface. For the case of weak fields the generalized covariant derivatives may be represented as (see )
$$D_{\pm ,\alpha }^{d_i}t^{\mu \nu }^{}D_{\pm ,\alpha }^{d_i}t^{\mu \nu }+^{\prime \prime }D_{\pm ,\alpha }^{d_i}t^{\mu \nu }$$
(38)
The $`{}_{}{}^{}D_{\pm ,\alpha }^{d_i}`$ in (38) describes the contribution from FD of time and space, the member $`{}_{}{}^{\prime \prime }D_{\pm ,\alpha }^{d_i}`$ describes the contribution from Riemann space with integer dimensions.
Let us see what differences between very similar equations (20) and (34). The equation (20) differs from equation (34) based on the Riemannian measure by three aspects:
a) the metric tensor $`\gamma ^{\mu \nu }`$ in (20) determined on the Minkowski space with fractional dimensions;
b) equations differs by dependencies of metrics tensor $`\gamma ^{\mu \nu }`$ from $`L`$ ( because in the (20) there are no dependencies in the $`\gamma ^{\mu \nu }`$ from $`L`$ originated by the Riemann metric tensor), there are only dependencies originating by FD;
c) the reason of appearance in equation (20) of the dependencies the $`\gamma ^{\mu \nu }`$ at $`L`$ lay in the originate it by the only fractal dimensions of time and space. If FD are integer the (20) coincide with equation of the theory . If FD integer in (34) these equations coincide with equations of GR. For weak fields GFD may be represented only by FD covariant derivatives (only two members in the right part of (29)) and in that case (20) may be represented by metric tensor $`g^{\mu \nu }`$ of an ”effective” Riemann space with integer dimensions as in (see also ). So (20) gives the equations GR too. We pay attention that the corresponding results of the theory for connections between metric tensor $`\gamma ^{\mu \nu }`$ of Minkowski space with ”effective” metric tensor $`g^{\mu \nu }`$ of Riemann space and gravitation tensor (thou they are valid) are the special case of our theory. In general case the metric tensor of Minkowski space are complicated function of gravitation field tensor. We took into account the alterations in equations originated by modified Lorents transformations in models with both measures (Riemann and Minkowski).
## VII Conclusion
In this paper we considered the two models of gravity theories defined on the multifractal sets of time and space. From our point of view there are two main approach to theories of gravity: the first is the approach of Einstein’s theory in which the gravitation fields and forces are no exist and its role play the curvature of the Riemannian time-space originated by Riemannian geometry. The second approach is the approach of postulating the real gravitational fields and forces made in the Logunov-Mestvirichvili theory . The last theory treats the gravity as an usual real field in the flat pseudo Euclidean time-space and it is the very attractive feature of this theory. The results of both theories coincide on the distances far from the gravitational radius of centre of gravity. On the distances of order the gravitational radius (if our Universe is multifractal set of space and time points defined on the measure carrier) the both theories it seems are not correct. In the Universe with multifractal dimensions of time and space on these distances the main role will play the integral characteristics of GFD and all equations become not differential but integral equations without containing of any infinity. We presented in these paper the new theory of gravity: gravity theory in the time and the space with fractional dimensions. This theory use the idea and results of works \- , \- and take into account the corrections to SR given by the theory of almost inertial systems in the time with fractional dimensions . In other words this gravitation theory expand the main results of on the gravitational fields.
Let us enumerate now the main results the theory presented in this paper:
1) The theory gives four sorts of different gravitation fields: two fields with real energy ( these fields differs by the sign of their energies (the field for gravitons and the field for anti-gravitons) and two fields with imaginary energies. The situation is the same as for vector (electromagnetic) and spinor (electron-positron) fields considered in ;
2) The interactions for each of both fields with real energies with imaginary energies fields are different. This gives the possibility to introduce the assumption about existence of new characteristics of gravitational fields ( ”quasi-spin” ) for explaining these facts;
3) The consideration of two models of a measure carriers ( the measure on Riemannian space and the measure on Minkowski space) are made.
4) The presented theory coincide with GR or the theory Logunov-Mestvirichvili for case when FD of time and space become integer.
5) In the fractal time and space the ordinary derivatives and integral must be replaced by GFD. So the main idea of this work may be used for generalization of all gravitational theories not considered here , including quantum theories of gravitation.
6) In this paper adopted the point of view: our Universe is multifractal sets of time and space ”points” (see details in the \- ). As any multifractal set it defined on a measure carrier. Thus the Universe is an open system (statistical theory of open systems see in the , ) and the all physical conservation laws (energy, mass and so on) fulfill as the very good approach. The exact the conservation laws fulfill only for closed system: the Universe plus the measure carrier. So the correct selection of measure is a very serious task and it is the task of near future. In the domains of Universe where the correction to integer time and space dimensions are very small (in such domain of Universe we live and such domains are in distances far away from gravitational radius) the exchange by energy, mass, momentum and so on between the Universe and the measure carrier is very small too and it may be neglected. It is necessary nevertheless remember about continuous exchange by energy between Universe and the measure carrier ( absorption and emission of energy) in every place of our Universe in the frame of presented fractal theory. The Universe never lost its energy it seems in that case and the far energy future of Universe is not so sad.
(7) In this paper we considered the characteristics of gravitational fields (characteristics electromagnetic and electron-positron fields were considered in . Naturally the algorithm used in the paper may be applied to any fields (electro-weak, Lee-Yang, quarks and so on) in domains where the fractional correction to the dimensions of time are small. In that case will be true the main results of this paper: every physical fields must be replaced by four fields with the real and imaginary energies. So it seems very likely that all physical fields must have their imaginary twins (if our Universe is fractal).
8) Nobody knows what the time and the space dimensions has our Universe. If the dimensions of time and space are fractional the presented in this paper the theories of gravitational fields will be true ( if at least one of the selections the measures carrier are valid) and will describe the reality of our Universe. As was stressed in the one of methods of verification the fractal theory of time and space is to accelerate the charge particle to speed of light that in the time with fractional dimensions is possible (because for spaces with FD of time the SR was modified in the narrow domain of velocities near velocity of light (see - ) .
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# INFNFE-03-00 A new approach to multi-jet calculations in hadron collisions
## 1 Introduction
The ability to evaluate production rates for multi-jet final states will be fundamental at the LHC to study a large class of processes, within and beyond the SM. A necessary feature of any multi-jet calculation is the possibility to properly evolve the purely partonic final state, for which exact fixed-order perturbative calculations can be performed, into the observable hadronic final state. This evolution is best performed using shower Monte Carlo calculations. The accurate description of color-coherence effects, furthermore, requires a careful bookkeeping of the contribution to the matrix elements of all possible color configurations. The goal of the algorithm described in this talk is to allow the effective calculation of multi-parton matrix elements, allowing the separation, to the leading order in $`1/N_c^2`$ ($`N_c=3`$ being the number of colors), of the independent color configurations. This technique allows an unweighting of the color configurations, and allows the merging of the parton level calculation with the HERWIG Monte Carlo.
The key element of the strategy is the use of the algorithm ALPHA, introduced in Ref. for the evaluation of arbitrary multi-parton matrix elements. The algorithm is built up from an iteration of matrix multiplication and is intended for the authomatic calculation of tree-level amplitudes. It has a complexity growing like a power in the number of particles, compared to the factorial-like growth that one expects from naive diagram counting.
The ALPHA algorithm will be reviewed in section (2), and its application to multi-jets physics in section (3).
## 2 The ALPHA algorithm
In reference , a new approach to the computation of tree level scattering amplitudes was introduced. This approach, based on the numerical Legendre transform of the effective action, is particularly useful for the authomatic calculation of multi-particle processes. This technique was implemented in a FORTRAN code which has been succesfully used to study several intricated electroweak processes of interest both at LEP and at the NLC . Among the key features of the algorithm are:
* Suitable (easy to implement) for the authomatic calculation of the scattering amplitudes.
* Inclusion of mass effects is straightforward
* The CPU cost grows like $`K^n`$ ($`K`$ constant, $`n`$ number of external particles) as opposed to the factorial growth of the number of Feynman graphs.
We review in this Section the ALPHA algorithm; the interested reader can find a more detailed discussion in the original paper , which includes an explicit analytic example for the $`\lambda \varphi ^3`$ theory.
Let $`\mathrm{\Gamma }`$ be the one-particle-irreducible generator of the Green functions for a given theory. Then the computation of the S-matrix requires the evaluation of the Legendre transform, Z, of $`\mathrm{\Gamma }`$:
$$Z(J^\alpha )=\mathrm{\Gamma }(\varphi ^\alpha )+J^\alpha (x)\varphi ^\alpha (x)$$
where $`\varphi ^\alpha `$ are the classical fields defined as the solutions of
$$J^\alpha =\frac{\delta \mathrm{\Gamma }}{\delta \varphi ^\alpha },$$
and the $`J^\alpha `$ play the role of classical sources.
For concretness we will develop in some detail the case of the scattering amplitue for $`n`$ external gluons. At tree-level $`\mathrm{\Gamma }`$ coincide with the Lagrangian. In momentum space
$`_{YM}(A)`$ $`=`$ $`1/2(p_\mu A_\nu ^ap_\nu A_\mu ^a)^2`$
$`+gf_{abc}(p_\mu A_\nu ^ap_\nu A_\mu ^a)A_\mu ^bA_\nu ^c`$
$`(B_{\mu \nu }^a)^22gf_{abc}B_{\mu \nu }^aA_\mu ^bA_\nu ^c+J_\mu ^aA_\mu ^a`$
where $`B_{\mu \nu }^a`$ is an auxiliary field and it is introduced in order to deal with trilinear interactions only. The sources $`J_\mu ^a`$ is a standard source term: $`J(p)=_{j=1}^nϵ_\mu ^a\delta (pp_j)`$, i.e. it contains the relevant excitation for the external particles. Notice that because of this choice only a finite number of momenta (linear combinations of $`p_j`$) enters into the problem which is reduced to a problem with a finite number of degrees of freedom.
The fields $`A(p)`$ are then found as solutions of the equation of motion. Let us stick to Feynman gauge for definitness.
$`A_\mu ^a(p)`$ $`=`$ $`{\displaystyle \frac{g}{p^2}}f_{abc}[2(kp)A^b(q)A_\mu ^c(k)`$
$`2q_\mu A^b(q)A_\mu ^c(k)`$
$`B_{\mu \nu }^b(q)A_\nu ^c(k)]+{\displaystyle \frac{1}{p^2}}J_\mu (p)`$
$`B_{\mu \nu }^a(p)`$ $`=`$ $`gf_{abc}A_\mu ^b(q)A_\nu ^c(k)`$
where momentum conservation is understood.
The equations of motion (LABEL:motion) are solved iteratively (expansion in $`g`$) and this implies that the problem is solved with a loop of matrices multiplication, more suitable for numerical implementation than the standard approach. The initialization steps are
$`A_\mu ^a(p_j)`$ $`=`$ $`ϵ_\mu ^a(p_j)`$
$`B_{\mu \nu }^a(p_j)`$ $`=`$ $`0`$
and the subsequent steps
$`A_\mu ^a(p_j+p_k)`$ $`=`$
$`{\displaystyle \frac{gf_{abc}}{(p_j+p_k)^2}}[2(p_j+2p_k)A^b(p_j)A_\mu ^c(p_k)`$
$`+2(p_k+2p_j)A^b(p_k)A_\mu ^c(p_j)`$
$`2(p_j+p_k)_\mu A^b(p_k)A^c(p_j)`$
$`B_{\mu \nu }^b(p_j)A_\mu ^c(p_k)]`$
$`B_{\mu \nu }^a(p_j+p_k)`$ $`=`$ $`gf_{abc}A^b(p_k)A_\mu ^c(p_j)`$
The step (LABEL:iter) is then iterated until we have constructed $`A`$ and $`B`$ fields with up to $`n`$ momenta. An important remark is in order here. Performing the iteration (LABEL:iter) we drop terms containing $`ϵ^a(p_j)ϵ^a(p_j)`$ i.e. twice the same external momenta. In fact, although legitimate, they don’t contribute to the amplitude.
Finally the amplitude is
$$𝒜(p_{j_1},\mathrm{},p_{j_n})=_{YM}(A)$$
(4)
where the $`A`$ fields are obtained in eqns. (LABEL:init,LABEL:iter)<sup>1</sup><sup>1</sup>1Using the equation of motion it is possible to halve the required number of iteration step as well as to reduce the number of contributions entering the final expression of the amplitude.. Notice that the prescription to drop terms containing at least twice the same external momenta ($`ϵ^a(p_j)ϵ^a(p_j)`$) is still kept and it is because of this prescription toghether with the choice (LABEL:init) for the initial step that neither truncation nor functional derivation is actually required.
A remark is in order here. The algorithm sketched above has an important feature: it provides a very compact way of storing the relevant information. Indeed the number of contraction $`A^a(p_{j_1}+\mathrm{}+p_{j_m})`$ (see eq. LABEL:iter) which is needed is of the order $`2^n`$, each of them requiring roughly the same CPU time to be computed. Therefore the CPU cost of the algorithm grows like a constant to the power n, n being the number of external particles, as opposed to the factorial growth of the number of Feynman graphs. To hint how does this work let us develop in more detail the case of the amplitude for five external gluons. For simplicity we neglect the four gluons coupling. The $`A`$ fields are given by
$`A_\mu ^a(p_j)`$ $`=`$ $`ϵ_\mu ^a(p_j)j=1,\mathrm{},5`$
$`A_\mu ^a(p_j+p_k)`$ $`=`$
$`{\displaystyle \frac{gf_{abc}}{(p_j+p_k)^2}}[2(p_j+2p_k)A^b(p_j)A_\mu ^c(p_k)`$
$`+2(p_k+2p_j)A^b(p_k)A_\mu ^c(p_j)`$
$`2(p_j+p_k)_\mu A^b(p_k)A^c(p_j)]`$
$`j=1,\mathrm{},5`$
and the ALPHA amplitude by
$`𝒜`$ $`=`$ $`g{\displaystyle \underset{j=1}{\overset{5}{}}}A^a(p_j)_\alpha `$
$`{\displaystyle \underset{klmnj}{\overset{5}{}}}f_{abc}\tau _{\alpha \beta \gamma }A^b(p_k+p_l)_\beta A^c(p_m+p_n)_\gamma `$
$`\tau _{\alpha \beta \gamma }`$ $`=`$ $`g_{\alpha \beta }(p_jp_kp_l)_\gamma `$
$`+g_{\gamma \alpha }(p_m+p_np_j)_\beta `$
$`+g_{\beta \gamma }(p_k+p_lp_mp_n)_\alpha `$
notice that the summations over $`k,l,m,n`$, as well as over lorentz and color indices, are carried over before the multiplication by $`A^a(p_j)`$. It is this feature extensively used in the construction of both the $`A`$ fields and the amplitude which turns the factorial growth of the number of Feynman graphs to a power law growth of the CPU cost of the ALPHA algorithm. Notice that in eqns. (LABEL:suba,LABEL:ampl5) $`A(p_j+p_k+p_m)`$ and $`A(p_j+p_k+p_m+p_l)`$ are not computed (and used). Morover neither the kinetic term nor the source one appears in the expression (LABEL:ampl5) of the amplitude. These simplifications occur because of the equation of motion as discussed in more detail in .
## 3 Multi-jets processes
Multi-jet final states play an important role in the study of high-energy collisions. They provide in fact interesting signatures for several phenomena, both within the Standard Model (e.g. top-pair production), and beyond it (e.g. multi-jet decays of supersymmetric particles such as gluinos and squarks). The accurate determination of the properties of these phenomena requires a good understanding of the properties of the usually large multi-jet QCD backgrounds, which can distort the shapes of signal distributions and affect the measurement of quantities such as resonances’ masses.
There are several reasons for wanting to improve the tools currently available to perform these calculations.
1. First of all, interesting final states with larger jet multiplicities will become available with the next generation of colliders (LHC and NLC).
2. Secondly, one would like to be able to complement the calculation of parton-level matrix elements with the evaluation of the full hadronic structure of the final state.
The goal of the algorithm described in this Section is to allow the effective calculation of multi-parton matrix elements, allowing the separation, to the leading order in $`1/N_c^2`$ ($`N_c=3`$ being the number of colors), of the independent color configurations. This technique allows an unweighting of the color configurations, and allows the merging of the parton level calculation with the HERWIG Monte Carlo. The key ingredient of our strategy is the ALPHA algorithm outlined in the previous section. It’s power-law growing in complexity is a necessary feature of any attempt to evaluate matrix elements for processes with large numbers of external particles, since the number of Feynman diagrams grows very quickly (see table 1) beyond any reasonable value.
Once the hard scattering matrix element is known the prescription to correctly generate the parton-shower associated to a given event in the large-$`N_c`$ limit is the following:
1. Calculate the $`(n1)!`$ dual amplitudes corresponding to all possible planar color configurations. This can be done with ALPHA using $`N_c`$ large enough .
2. Extract the most likely color configuration for this event on a statistical basis, according to the relative contribution of the single configurations to the total event weight <sup>2</sup><sup>2</sup>2Defining $`w_i=|A_i|^2`$ for each color flow $`i`$, and $`W_i=_{k=1,\mathrm{},i}w_k/_{k=1,\mathrm{},n}w_k`$, the $`j`$-th color structure will be selected if $`W_{j1}\xi <W_j`$, for a random number $`\xi `$ uniformly distributed over the interval $`[0,1]`$.. Since each dual amplitude is gauge invariant, the choice of color-configurations is also a gauge-invariant operation.
3. Develop the PS out of each initial and final-state parton, starting from the selected color configuration. This step can be carried out by feeding the generated event to a Monte-Carlo program such as HERWIG, which is precisely designed to turn partons into jets starting from an assigned color-ordered configuration.
Notice that, if the dual amplitudes are evaluated for a specific helicity configuration, HERWIG will also include spin-correlation effects in the evolution of the parton shower .
As a result, use of the dual-amplitude representation of a multi-gluon amplitude allows to accurately describe not only the large-angle correlations in multi-jet final states, but also the full shower evolution of the initial- and final-state partons with the same accuracy available in HERWIG for the description of 2-jet final states.
In alternative to the above prescription, one can use ALPHA to calculate the matrix elements for external states with assigned colors. Since these states are all orthogonal, such an approach is particularly efficient if one wants to use a Monte Carlo method to sum over all possible color states. The program will then extract through a standard unweighting (at the leading order in $`1/N_c^2`$) a specific color flow from all possible color flows contributing to a given orthogonal color state. This color flow is then suitable as an initial condition for the shower evolution. The advantage of this approach is twofold: first the number of dual amplitudes contributing to the amplitude for external states with assigned colors is substantially smaller than the total one and second dual amplitudes are required only for accepted events which, in general, are a small fraction of the generated ones. Further details can be found in .
### 3.1 Results
As an example of our technique, we present here results for the following two parton-level processes:
$`gg`$ $``$ $`8g`$
$`q\overline{q}`$ $``$ $`8g.`$ (7)
For comparisons, we also computed the above reactions in the simple approximation first suggested by Kunszt and Stirling . This approximation (hereafter referred to as SPHEL) consists in assuming that the average value of maximally helicity violating amplitudes is equal to the average value of all other non-zero amplitudes.
The kinematic configuration and the cut values used in our numerical examples are as follows:
$`\sqrt{\widehat{s}}=1500\mathrm{GeV},`$ $`p_{T_i}>60\mathrm{GeV},|\eta _i|<2,`$
$`\mathrm{\Delta }R_{ij}>0.7.`$
These values, and the choice of a fixed strong coupling $`\alpha _s=0.12`$, only serve for illustrative purposes.
In Fig. 1, we show the distribution of the minimum gluon transverse momentum for both processes. In Fig 2 we plot the distributions for the maximum gluon transverse momentum.
## 4 Conclusions
We have reviewed the ALPHA algorithm to evaluate the exact, tree-level matrix elements in the context of multi-parton processes in QCD. The algorithm is suitable (easy to implement) for the authomatic calculation of the scattering amplitudes and it allows the inclusion of mass effects in a straightforward manner. The CPU cost of the algorithm grows like $`K^n`$ ($`K`$ constant, $`n`$ number of external particles) as opposed to the factorial growth of the number of Feynman graphs.
This technique has been tested for processes such as $`ggn`$ gluons and $`q\overline{q}n`$ gluons, with $`n`$ up to 9. We discussed how the summation over colour configurations allows the construction of parton-level event generators suitable to interfacing with a parton-shower evolution including the effects of colour-coherence. This will eventually lead to a fully exclusive, hadron-level description of multi-jet final states, accurately incorporating the dynamics of large jet-jet separation angles.
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# On the geometric simple connectivity of open manifolds IMRN 2004, to appear
## 1 Statements of the results
The problem we address in this paper is whether 1-handles are necessary in a handle decomposition of a simply connected manifold. Moreover we investigate when it is possible to kill 1-handles within the proper homotopy type of a given open manifold.
The relation between algebraic connectivity and geometric connectivity (in various forms) was explored first by E.C. Zeeman (see ) in connection with the Poincaré conjecture. Zeeman’s definition of the geometric $`k`$-connectivity of a manifold amounts to asking that any $`k`$-dimensional compact can be engulfed in a ball. His main result was the equivalence of algebraic $`k`$-connectivity and geometric $`k`$-connectivity for $`n`$-manifolds, under the condition $`kn3`$. Notice that it makes no difference whether one considers open and compact manifolds.
Later C.T.C. Wall () introduced another concept of geometric connectivity using handle theory which was further developed by V. Poénaru in his work around the Poincaré conjecture. A similar equivalence between the geometric and algebraic connectivities holds *in the compact case* but this time one has to replace the previous codimension condition by $`kn4`$. In this respect all results in low codimension are hard results. There is also a non-compact version of this definition which we can state precisely as follows:
###### Definition 1.1.
A (possibly non-compact) manifold, which might have nonempty boundary, is geometrically $`k`$-connected (abbreviated. g.$`k`$.c.) if there exists a proper handlebody decomposition without $`j`$-handles, for $`1jk`$.
One should emphasize that now the compact and non-compact situations are no longer the same. The geometric connectivity is a consequence of the algebraic connectivity only under additional hypotheses concerning the ends. The purpose of this paper is to partially characterize these additional conditions.
###### Remark 1.1.
Handle decomposition are known to exist for all manifolds in the topological, PL and smooth settings, except in the case of topological 4-manifolds. In the latter case the existence of a handlebody decomposition is equivalent to that of a PL (or smooth) structure. However in the open case such a smooth structure always exist (in dimension 4). Although most results below can be restated and proved for other categories, we will restrict ourselves to considering PL manifolds and handle decompositions in the sequel.
We will be mainly concerned with geometric simple connectivity (abbreviated. g.s.c.) in the sequel. A related concept, relevant only in the non-compact case, is:
###### Definition 1.2.
A (possibly non-compact) polyhedron $`P`$ is weakly geometrically simply connected (abbreviated w.g.s.c.) if $`P=_{j=1}^{\mathrm{}}K_j`$, where $`K_1K_2\mathrm{}K_j\mathrm{}`$ is an exhaustion by compact connected sub-polyhedra with $`\pi _1(K_j)=0`$. Alternatively, any compact subspace is contained in a simply connected sub-polyhedron.
Notice that a w.g.s.c. polyhedron is simply connected.
###### Remark 1.2.
* The w.g.s.c. spaces with which we will be concerned in the sequel are usually manifolds. Similar definitions can be given in the case of topological (respectively smooth) manifolds where we require the exhaustions to be by topological (respectively smooth) submanifolds. All results below hold true for this setting too (provided handlebodies exists) except those concerning Dehn exhaustibility, since the later is essentially a PL concept.
* The w.g.s.c. is much more flexible than the g.s.c., the latter making sense only for manifolds, and enables us to work within the realm of polyhedra. However one can easily show (see below) that w.g.s.c. and g.s.c. are equivalent for non-compact manifolds of dimension different from $`4`$ (under the additional irreducibility assumption for dimension $`3`$). Its invariance under proper homotopy equivalences expresses the persistence of a geometric property (not being g.s.c.) with respect to some higher dimensional manipulations (as taking the product with a ball) of open manifolds.
The first result of this paper is (see Theorem 3.15 and Proposition 3.24):
###### Theorem 1.3.
If an open $`n`$-manifold is w.g.s.c., then it is end compressible. Conversely, in dimension $`n4`$, if an open, simply connected manifold is end compressible, then it is g.s.c. (one has to assume irreducibility for $`n=3`$).
###### Remark 1.4.
A similar result holds more generally for non-compact manifolds with boundary, with the appropriate definition of end compressibility.
End compressibility (see Definition 3.7) is an algebraic condition which is defined in terms of the fundamental groups of the submanifolds which form an exhaustion. Notice that end compressibility is weaker than simple connectivity at infinity for $`n>3`$.
###### Remark 1.5.
The result above should be compared with Siebenmann’s obstructions to finding a boundary for an open manifold of dimension greater than 5 (see and for a thorough discussion of this and related topics). The intermediary result permitting to kill 1-handles in this framework is theorem 3.10 p.16 from : Let $`W^n`$ be an open smooth $`n`$-manifold with $`n5`$ and $``$ an isolated end. Assume that the end $``$ has stable $`\pi _1`$ and its $`\pi _1()`$ is finitely presented. Then there exists arbitrarily small 1-neighborhoods of $``$ i.e. connected submanifolds $`V^nW^n`$ whose complements have compact closure, having compact connected boundary $`V^n`$ such that $`\pi _1()\pi _1(V)`$ and $`\pi _1(V)\pi _1(V)`$ are isomorphisms. It is easy to see that this implies that the 1-neighborhood $`V^n`$ is g.s.c. This is the principal step towards canceling the handles of $`W^n`$ hence obtaining a collar. One notices that the hypotheses in Siebenmann’s theorem are stronger than the end compressibility but the conclusion is stronger too. In particular an arbitrary w.g.s.c. manifold need not have a well-defined fundamental group at infinity, as is the case for $`\pi _1`$-stable ends. However we think that the relationship between the $`\pi _1`$-(semi)stability of ends and end compressibility would deserve further investigation.
The full power of the $`\pi _1`$-stability is used to cancel more than 1-handles. Actually L. Siebenmann considered tame ends, which means that $``$ is $`\pi _1`$-stable and it has arbitrarily small neighborhoods which are finitely dominated. The tameness condition is strong enough to insure (see theorem 4.5 of ) that all $`k`$-handles can be canceled for $`kn3`$. One more obstruction (the end obstruction) is actually needed in order to be able to cancel the $`(n2)`$-handles (which it turns out to imply the existence of a collar). There exist tame ends which are not collared (i.e. with non-vanishing end obstruction), as well as $`\pi _1`$-stable ends with finitely presented $`\pi _1()`$ which are not tame. Thus the obstructions for killing properly the handles of index $`1\lambda k`$ should be weaker than the tameness of the end for $`kn3`$ and must coincide with Siebenmann’s for $`k=n2`$.
###### Remark 1.6.
It might be worthy to compare our approach with the results from . First the w.g.s.c. is the analogue of the reverse collaring. According to (, Prop.8.5, p.93) a space $`W`$ is reverse collared if it has an exhaustion by compacts $`K_j`$ for which the inclusions $`K_jW`$ are homotopy equivalences (while the w.g.s.c. asks only that these inclusions be 1-connected), and hence a simply connected reverse collared space is w.g.s.c.
The right extension of the w.g.s.c. to non-simply connected spaces, which is suitable for applications to 3-manifolds, is the Tucker property (see ), which is also a proper homotopy invariant and can be formulated as a group theoretical property for coverings. This is again weaker than the reverse tameness/collaring (see , Prop.8.9 and Prop.11.13). Moreover there is a big difference between the extended theory and the present one: while the w.g.s.c. is the property of having no extra 1-handles, the Tucker property expresses the fact that some handlebody decomposition needs only finitely many 1-handles, without any control on their number.
In this respect our first result is a sharpening of the theory of reverse collared manifolds specific to the realm of simply connected spaces.
The g.s.c. is mostly interesting in low dimensions, for instance in dimension 3 where it implies the simple connectivity at infinity. However its importance relies on its proper homotopy invariance, which has been discovered in a particular form by V. Poénaru (see also ), enabling us to transform the low dimensional problem ”is $`W^3`$ g.s.c.?” into a high dimensional one, for example ”is $`W^3\times D^n`$ w.g.s.c.?”. We provide in this paper a criterion permitting to check the answer to the high dimensional question, in terms of an arbitrary exhaustion by compact submanifolds. This criterion is expressed algebraically as the end compressibility of the manifold and is closer to the forward tameness from , rather than the reverse tameness.
###### Remark 1.7.
If $`W^k`$ is compact and simply-connected then the product $`W^k\times D^n`$ with a closed $`n`$-disk is g.s.c. if $`n+k5`$. However there exist non-compact $`n`$-manifolds with boundary which are simply connected but not end compressible (hence not w.g.s.c.) in any dimension $`n`$, for instance $`W^3\times M^n`$ where $`\pi _1^{\mathrm{}}W^30`$. Notice that $`W^k\times \mathrm{int}(D^n)`$ is g.s.c. for $`n1`$ since $`\pi _1^{\mathrm{}}(W^k\times \mathrm{int}(D^n))=0`$.
We will also prove (see Theorem 4.8):
###### Theorem 1.8.
There exist uncountably many open contractible $`n`$-manifolds for any $`n4`$ which are not w.g.s.c.
The original motivation for this paper was to try to kill 1-handles of open 3-manifolds at least stably (i.e. after stabilizing the $`3`$-manifold). The meaning of the word stably in , where such results first arose, is to do so at the expense of taking products with some high dimensional compact ball. This was extended in by allowing the 3-manifold to be replaced by any other polyhedron having the same proper homotopy type. The analogous result is true for $`n5`$ (for $`n=4`$ only a weaker statement holds true (see Theorem 5.2):
###### Theorem 1.9.
If a non-compact manifold of dimension $`n4`$ is proper homotopically dominated by a w.g.s.c. polyhedron, then it is w.g.s.c. A non-compact $`4`$-manifold proper homotopically dominated by a w.g.s.c. polyhedron is end compressible.
Since proper homotopy equivalence implies proper homotopy domination we obtain:
###### Corollary 1.10.
If a non-compact manifold of dimension $`n4`$ is proper homotopy equivalent to a w.g.s.c. polyhedron, then it is w.g.s.c.
###### Remark 1.11.
This criterion is an essential ingredient in Poénaru’s proof (see ) of the covering space conjecture: If $`M^3`$ is a closed, irreducible, aspherical 3-manifold, then the universal covering space of $`M^3`$ is $`𝐑^3`$. Further developments suggest a similar result in higher dimensions, by replacing the simple connectivity at infinity conclusion with the weaker w.g.s.c. We will state below a group-theoretical conjecture abstracting this purely 3-dimensional result.
It is very probable that there exist examples of open 4-manifolds which are not w.g.s.c., but their products with a closed ball are w.g.s.c. Thus in some sense the previous result is sharp.
The dimension 4 deserves special attention also because one expects that the w.g.s.c. and the g.s.c. are not equivalent. Specifically V. Poénaru conjectured that:
###### Conjecture 1.12 (Poénaru Conjecture).
If the interior of a compact contractible 4-manifold with boundary a homology sphere is g.s.c. then the compact 4-manifold is also g.s.c.
###### Remark 1.13.
* A consequence of this conjecture, for the particular case of the product of a homotopy 3-disk $`\mathrm{\Delta }^3`$ with an interval, is the Poincaré conjecture in dimension 3. This follows from the two results announced by V. Poénaru:
Theorem: If $`\mathrm{\Sigma }^3`$ is a homotopy 3-sphere such that $`\mathrm{\Sigma }^3\times [0,1]`$ is g.s.c. then $`\mathrm{\Sigma }^3`$ is g.s.c. (hence standard).
Theorem: If $`\mathrm{\Delta }^3`$ is a homotopy 3-disk then $`\mathrm{int}(\mathrm{\Delta }^3\times [0,1]\mathrm{}_{\mathrm{}}S^2\times D^2)`$ is g.s.c., where $`\mathrm{}`$ denotes the boundary connected sum.
* The differentiable Poincaré conjecture in dimension 4 is widely believed to be false. One reasonable reformulation of it would be the following: A smooth homotopy 4-sphere (equivalently, homeomorphic to $`S^4`$) that is g.s.c. should be diffeomorphic to $`S^4`$.
* The two conjectures above (Poénaru’s and the reformulated Poincaré) conjectures imply also the smooth Schoenflies conjecture in dimension 4, which states that a $`3`$-sphere smoothly embedded in $`S^4`$ bounds a smoothly embedded $`4`$-ball. In fact by a celebrated result of B. Mazur any such Schoenflies ball has its interior diffeomorphic to $`𝐑^4`$, hence g.s.c.
An immediate corollary would be that the interior of a Poénaru-Mazur 4-manifold may be w.g.s.c. but not g.s.c., because some (compact) Poénaru-Mazur 4-manifolds are known to be not g.s.c. (the geometrization conjecture implies this statement for all 4-manifolds whose boundary is not a homotopy sphere). The proof is due to A. Casson and it was based on partial positive solutions to the following algebraic conjecture \[13, p.117\] and \[16, p.403\].
###### Conjecture 1.14 (Kervaire Conjecture).
Suppose one adds an equal number of generators $`\alpha _1,\mathrm{},\alpha _n`$ and relations $`r_1,\mathrm{},r_n`$ to a non-trivial group $`G`$, then the group $`\frac{G\alpha _1,\mathrm{},\alpha _n}{r_1,\mathrm{},r_n}`$ that one obtains is also non-trivial.
A. Casson showed that certain $`4`$-manifolds $`(W^4,W^4)`$ have no handle decompositions without $`1`$-handles by showing that if they did, then $`\pi _1(W^4)`$ violates the Kervaire conjecture. Our aim would be to show that most contractible 4-manifolds are not g.s.c., and the method of the proof is to reduce this statement to the compact case. However our methods permit us to obtain only a weaker result, in which one shows that the interior of such a manifold cannot have handlebody decompositions without 1-handles, if the decomposition has also some additional properties (see Theorems 7.11 and 7.12 for precise statements):
###### Theorem 1.15.
Assume that we have a proper handlebody decomposition without 1-handles for the interior of a Poénaru-Mazur 4-manifold. If there exists a far away intermediary level 3-manifold $`M^3`$ whose homology is represented by disjoint embedded surfaces and whose fundamental group projects to the trivial group on the boundary, then the compact 4-manifold is also g.s.c. There always exists a collection of immersed surfaces, which might have non-trivial intersections and self-intersections along homologically trivial curves, that fulfills the previous requirements.
###### Remark 1.16.
Almost all of this paper deals with geometric 1-connectivity. However the results can be reformulated for higher geometric connectivities within the same range of codimensions.
We wish to emphasize that there is a strong group theoretical flavour in the w.g.s.c. for universal covering spaces. In this respect the universal covering conjecture in dimension three (see ) would be a first step in a more general program. Let us define a finitely presented, infinite group $`\mathrm{\Gamma }`$ to be w.g.s.c. if there exists a compact polyhedron with fundamental group $`\mathrm{\Gamma }`$ whose universal covering space is w.g.s.c. It is not hard to show that this definition does not depend on the particular polyhedron one chooses but only on the group. This is part of a more general philosophy, due to M. Gromov, in which infinite groups are considered as geometric objects. This agrees with the idea that killing 1-handles of manifolds is a group theoretical problem in topological disguise. The authors think that the following might well be true:
###### Conjecture 1.17.
Fundamental groups of closed aspherical manifolds are w.g.s.c.
This will be a far reaching generalization of the three dimensional result announced by V. Poénaru in . It is worthy to note that all reasonable examples of groups (e.g. word hyperbolic, semi-hyperbolic, $`CAT(0)`$, group extensions, one relator groups) are w.g.s.c. It would be interesting to find an example of a finitely presented group which fails to be w.g.s.c. Notice that the well-known examples of M. Davis of Coxeter groups which are fundamental groups of aspherical manifolds whose universal covering spaces are not simply connected at infinity are actually $`CAT(0)`$ hence w.g.s.c. However one might expect a direct connection between the semi-stability of finitely presented groups, the quasi-simple filtrated groups (see ) and the w.g.s.c. We will address these questions in a future paper.
### Outline of the paper
In section 2, we compare w.g.s.c., g.s.c and the simple connectivity at infinity (s.c.i.), showing that w.g.s.c and g.s.c. are equivalent in high dimensions and presenting some motivating examples.
Section 3 contains the core of the paper, where we introduce the algebraic conditions and prove their relation to w.g.s.c. We then construct uncountably many Whitehead-type manifolds in section 4, and show that there are uncountably many manifolds that are not geometrically simply-connected.
In section 5, we show that end-compressibility is a proper-homotopy invariant. Finally, in sections 6 and 7, we turn to the $`4`$-dimensional case.
Acknowledgements. We are indebted to David Gabai, Ross Geoghegan, Valentin Poénaru, Vlad Sergiescu and Larry Siebenmann for valuable discussions and suggestions. We are grateful to the referee for his comments which permitted us to improve the clarity and the readability of this paper. Part of this work has been done when the first author visited Tokyo Institute of Technology, which he wishes to thank for their support and hospitality, and especially to Teruaki Kitano, Tomoyoshi Yoshida and Akio Kawauchi. Part of this work was done while the second author was supported by a Sloan Dissertation fellowship.
## 2 On the g.s.c.
### 2.1 Killing 1-handles of 3-manifolds after stabilization
We start with some motivating remarks about the compact 3-dimensional situation, for the sake of comparison.
###### Definition 2.1.
The geometric 1-defect $`ϵ(M^n)`$ of the compact manifold $`M^n`$ is $`ϵ(M^n)=\mu _1(M^n)\mathrm{rank}\pi _1(M)`$, where $`\mu _1(M^n)`$ is the minimal number of 1-handles in a handlebody decomposition and $`\mathrm{rank}\pi _1(M)`$ is the minimal number of generators of $`\pi _1(M)`$.
###### Remark 2.1.
The defect (i.e., the geometric 1-defect) is always non-negative. There exist examples (see ) of 3-manifolds $`M^3`$ with $`\mathrm{rank}\pi _1(M^3)=2`$ and defect $`ϵ(M^3)=1`$. No examples of 3-manifolds with larger defect, nor of closed 4-manifolds with positive defect are presently known. However it’s probably true that $`ϵ(M^3\times [0,1])=0`$, for all closed 3-manifolds though this might be difficult to settle even for the explicit examples of M. Boileau and H. Zieschang.
The defect is meaningless in high dimensions because of:
###### Proposition 2.2.
For a compact manifold $`ϵ(M^n)=0`$ holds, if $`n5`$.
###### Proof.
The proof is standard. Consider a 2-complex $`K^2`$ associated to a presentation of $`\pi _1(M^n)`$ with the minimal number $`r`$ of generators. By general position there exists an embedding $`K^2M^n`$ inducing an isomorphism of fundamental groups. Then a regular neighborhood of $`K^2`$ in $`M^n`$ has a handlebody decomposition with $`r`$ 1-handles. Since the complement is 1-connected then by () it is g.s.c. for $`n5`$ and this yields the claim. ∎
###### Corollary 2.3.
For a closed 3-manifold $`M^3`$ one has $`ϵ(M^3\times D^2)=0`$.
###### Remark 2.4.
As a consequence if $`\mathrm{\Sigma }^3`$ is a homotopy 3-sphere then $`\mathrm{\Sigma }^3\times D^2`$ is g.s.c. Results of Mazur (, improved by Milnor in dimension 3) show that $`\mathrm{\Sigma }^3\times D^3=S^3\times D^3`$, but it is still unknown whether $`\mathrm{\Sigma }^3\times D^2=S^3\times D^2`$ holds. An earlier result of Poénaru states that $`(\mathrm{\Sigma }^3nD^3)\times D^2=(S^3nD^3)\times D^2`$ for some $`n1`$. More recently, Poénaru’s program reduced the Poincaré Conjecture to the g.s.c. of $`\mathrm{\Sigma }^3\times [0,1]`$.
### 2.2 $`\pi _1^{\mathrm{}}`$ and g.s.c.
The remarks which follow are intended to (partially) clarify the relation between g.s.c. and the simple connectivity at infinity (which will be abbreviated as s.c.i. in the sequel), in general.
Recall that a space $`X`$ is s.c.i. (and one also writes $`\pi _1^{\mathrm{}}(X)=0`$) if for any compact $`KX`$ there exists a larger compact $`KLX`$ having the property that, any loop $`lXL`$ is null homotopic within $`XK`$. This is an important tameness condition for the ends of the space. The following result was proved in (, Thm. 1):
###### Proposition 2.5.
Let $`W^n`$ be an open simply connected $`n`$-manifold of dimension $`n5`$. If $`\pi _1^{\mathrm{}}(W^n)=0`$ then $`W^n`$ is g.s.c.
###### Remark 2.6.
The converse fails as the following examples show. Namely, for any $`n5`$ there exist open $`n`$-manifolds $`W^n`$ which are geometrically $`(n4)`$-connected but $`\pi _1^{\mathrm{}}(W^n)0`$.
There exist compact contractible $`n`$-manifolds $`M^n`$ with $`\pi _1(M^n)0`$, for any $`n4`$ (see ). Since $`k`$-connected compact $`n`$-manifolds are geometrically $`k`$-connected if $`kn4`$ (see ), these manifolds are geometrically $`(n4)`$-connected. Let us consider now $`W^n=\mathrm{int}(M^n)`$, which is diffeomorphic to $`M^n_{M^nM^n\times \{0\}}M^n\times [0,1)`$. Any Morse function on $`M^n`$ extends over $`\mathrm{int}(M^n)`$ to a proper one which has no critical points in the open collar $`M^n\times [0,1)`$, hence $`\mathrm{int}(M^n)`$ is also geometrically $`(n4)`$-connected. On the other hand $`\pi _1(M^n)0`$ implies $`\pi _1^{\mathrm{}}(W^n)0`$.
However the following partial converse holds:
###### Proposition 2.7.
Let $`W^n`$ be a non-compact simply connected $`n`$-manifold which has a proper handlebody decomposition
1. without $`1`$ or $`(n2)`$handles, or
2. without $`(n1)`$ or $`(n2)`$handles.
Then $`\pi _1^{\mathrm{}}(W^n)=0`$.
###### Remark 2.8.
When $`n=3`$ this simply says that 1-handles are necessary unless $`\pi _1^{\mathrm{}}(W^3)=0`$.
Proof of Proposition 2.7. Consider the handlebody decomposition $`W^n=B^n_{j=1}^{\mathrm{}}h_j^{i_j}`$, where $`h_j^{i_j}`$ is an $`i_j`$-handle ($`B^n=h_0^0`$). Set $`X_m=B^n_{j=1}^mh_j^{i_j}`$, for $`m0`$. Assume that this decomposition has no $`1`$ nor $`(n2)`$handles. Since there are no 1-handles it follows that $`\pi _1(X_j)=0`$ for any $`j`$ (it is only here one uses the g.s.c.).
###### Lemma 2.9.
If $`X^n`$ is a compact simply connected $`n`$-manifold having a handlebody decomposition without $`(n2)`$handles then $`\pi _1(X^n)=0`$.
###### Proof.
Reversing the handlebody decomposition of $`X^n`$ one finds a decomposition from $`X^n`$ without 2-handles. One slides the handles to be attached in increasing order of their indices. Using Van Kampen Theorem it follows that $`\pi _1(X^n)=\pi _1(X^n)𝐅(r)`$, where $`r`$ is the number of 1-handles, and thus $`\pi _1(X^n)=0`$. ∎
###### Lemma 2.10.
If the compact submanifolds $`\mathrm{}X_mX_{m+1}\mathrm{}`$ exhausting the simply connected manifold $`W^n`$ satisfy $`\pi _1(X_m)=0`$, for all $`m`$, then $`\pi _1^{\mathrm{}}(W^n)=0`$.
###### Proof.
For $`n=3`$ this is clear. Thus we suppose $`n4`$. For any compact $`KW^n`$ choose some $`X_mK`$ such that $`X_mK=\mathrm{}`$. Consider a loop $`lW^nX_m`$. Then $`l`$ bounds an immersed (for $`n5`$ embedded) 2-disk $`\delta ^2`$ in $`W^n`$. We can assume that $`\delta ^2`$ is transversal to $`X_m`$. Thus it intersects $`X_m`$ along a collection of circles $`l_1,\mathrm{},l_pX_m`$. Since $`\pi _1(X_m)=0`$ one is able to cap off the loops $`l_j`$ by some immersed 2-disks $`\delta _jX_m`$. Excising the subsurface $`\delta ^2X_m`$ and replacing it by the disks $`\delta _j`$ one obtains an immersed 2-disk bounding $`l`$ in $`W^nK`$. ∎
This proves the first case of the Proposition 2.7. In order to prove the second case choose some connected compact subset $`KW^n`$. By compactness there exists $`k`$ such that $`KX_k`$. Let $`r`$ be large enough (this exists by the properness) such that any handle $`h_p^{i_p}`$ whose attaching zone touches the lateral surface of one of the handles $`h_1^{i_1},h_2^{i_2},\mathrm{},h_k^{i_k}`$ satisfies $`pr`$. The following claim will prove the Proposition 2.7:
###### Proposition 2.11.
Any loop $`l`$ in $`W^nX_r`$ is null-homotopic in $`W^nK`$.
###### Proof.
Actually the following more general engulfing result holds:
###### Proposition 2.12.
If $`C^2`$ is a 2-dimensional polyhedron whose boundary $`C^2`$ is contained in $`W^nX_r`$ then there exists an isotopy of $`W^n`$ (with compact support), fixing $`C^2`$ and moving $`C^2`$ into $`W^nK`$.
###### Proof.
Suppose that $`C^2X_m`$. One reverses the handlebody decomposition of $`X_m`$ and obtains a decomposition from $`X_m`$ without 1- or 2-handles. Assume that we can move $`C^2`$ such that it misses the last $`jr1`$ handles. By general position there exists an isotopy (fixing the last $`j`$ handles) making $`C^2`$ disjoint of the co-core ball of the $`(j1)`$-th handle, since the co-core disk has dimension at most $`n3`$. The uniqueness of the regular neighborhood implies that we can move $`C^2`$ out of the $`(j1)`$-th handle (see e.g. ), by an isotopy which is identity on the last $`j`$ handles. This proves the Proposition 2.12. ∎
This yields the result of Proposition 2.11 by taking for $`C^2`$ any 2-disk parameterizing a null homotopy of $`l`$. ∎
### 2.3 G.s.c. and w.g.s.c.
###### Proposition 2.13.
The non-compact manifold $`W^n`$ ($`n4`$), which one supposes to be irreducible if $`n=3`$, is w.g.s.c. if and only if it is g.s.c.
###### Proof.
For $`n=3`$ it is well-known that g.s.c. is equivalent to w.g.s.c. which is also equivalent to s.c.i. if the manifold is irreducible. For $`n5`$ this is a consequence of Wall’s result stating the equivalence of g.s.c. and simple connectivity in the compact case (see ). If $`W^n`$ is w.g.s.c. then it has an exhaustion by compact simply connected sub-manifolds $`M_j`$ (by taking suitable regular neighborhoods of the polyhedra). One can also refine the exhaustion such that the boundaries are disjoint. Then the pairs $`(\mathrm{cl}(M_{j+1}M_j),M_j)`$ are 1-connected, hence () they have a handlebody decomposition without 1-handles. Gluing together these intermediary decompositions we obtain a proper handlebody decomposition as claimed. ∎
## 3 W.g.s.c. and end compressibility
In this section, we show that w.g.s.c. is equivalent to an algebraic condition which we call *end compressibility*. This in turn implies infinitely many conditions, *end $`k`$-compressibility* for ordinals $`k`$, and is equivalent to all these plus a finiteness condition.
Using the above, we give explicit examples of open manifolds that are not w.g.s.c.
### 3.1 Algebraic preliminaries
In this section we introduce various algebraic notions of compressibility and study the relations between these. This will be applied in a topological context in subsequent sections, where compressibility corresponds to being able to attach enough two handles, and stable-compressibility refers to the same after possibly attaching some $`1`$-handles.
###### Definition 3.1.
A pair $`(\phi :AB,\psi :AC)`$ of group morphisms is strongly compressible if $`\phi (\mathrm{ker}\psi )=\phi (A)`$.
###### Remark 3.1.
Strong compressibility is symmetric in the arguments $`(\phi ,\psi )`$ i.e. $`\phi (\mathrm{ker}\psi )=\phi (A)`$ is equivalent to $`\psi (\mathrm{ker}\phi )=\psi (A)`$. The proof is an elementary diagram chase.
###### Definition 3.2.
A pair $`(\phi :AB,\psi :AC)`$ of group morphisms, is stably compressible if there exists some free group $`𝐅(r)`$ on finitely many generators, and a morphism $`\beta :𝐅(r)C`$, such that the pair $`(\phi 1_{𝐅(r)}:A𝐅(r)B𝐅(r),\psi \beta :A𝐅(r)C)`$ is strongly compressible.
We shall see that stable-incompressibility implies infinitely many conditions, indexed by the ordinals, on the pair of morphisms. We first define a series of groups (analogous to the lower central series).
###### Definition 3.3.
Consider a fixed pair $`(\phi :AB,\psi :AC)`$ of group morphisms. We define inductively a subgroup $`G_\alpha C`$ for any ordinal $`\alpha `$. Set $`G_0=C`$. If $`G_\alpha `$ is defined for every $`\alpha <\beta `$ (i.e. $`\beta `$ is a limit ordinal) then set $`G_\beta =_{\alpha <\beta }G_\alpha `$. Further set $`G_{\alpha +1}=𝒩(\psi (\mathrm{ker}\phi ),G_\alpha )G_\alpha `$ for any other ordinal, where $`𝒩(K,G)`$ is the smallest normal group containing $`K`$ in $`G`$.
The groups $`G_\alpha `$ form a decreasing sequence of subgroups of $`C`$. Using Zorn’s lemma there exists an infimum of the lattice of groups $`G_\alpha `$, ordered by the inclusion, which we denote by $`G_{\mathrm{}}=_\alpha G_\alpha `$ (over all ordinals $`\alpha `$).
###### Definition 3.4.
The pair $`(\phi :AB,\psi :AC)`$ is said to be $`\alpha `$-compressible if $`\psi (A)G_\alpha `$ (where $`\alpha `$ is an ordinal or $`\mathrm{}`$).
###### Lemma 3.2.
Given a subgroup $`LC`$ there exists a maximal subgroup $`\mathrm{\Gamma }=\mathrm{\Gamma }(L,C)`$ of $`C`$ so that $`L𝒩(L,\mathrm{\Gamma })=\mathrm{\Gamma }`$.
###### Proof.
There exists at least one group $`\mathrm{\Gamma }`$, for instance $`\mathrm{\Gamma }=L`$. Further if $`\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }^{}`$ verify the condition $`𝒩(L,\mathrm{\Gamma })=\mathrm{\Gamma }`$, then their product $`\mathrm{\Gamma }\mathrm{\Gamma }^{}`$ does. Thus, Zorn’s Lemma says that a maximal element for the lattice of subgroups verifying this property (the order relation is the inclusion) exists. ∎
###### Lemma 3.3.
We have $`\mathrm{\Gamma }(L,C)=G_{\mathrm{}}`$, where $`L=\psi (\mathrm{ker}\phi )`$.
###### Proof.
First, $`G_{\mathrm{}}`$ satisfies the condition $`𝒩(L,\mathrm{\Gamma })=\mathrm{\Gamma }`$ otherwise the minimality will be contradicted. Pick an arbitrary $`\mathrm{\Gamma }`$ satisfying this condition. If $`\mathrm{\Gamma }G_\alpha `$ it follows that $`\mathrm{\Gamma }=𝒩(L,\mathrm{\Gamma })𝒩(L,G_\alpha )=G_{\alpha +1}`$, hence by a transfinite induction we derive our claim. ∎
###### Definition 3.5.
One says that $`K`$ is full in $`\mathrm{\Gamma }`$ if $`𝒩(K,\mathrm{\Gamma })=\mathrm{\Gamma }`$. If we have a pair and $`\psi (\mathrm{ker}\phi )`$ is full in $`\mathrm{\Gamma }`$, then we call $`\mathrm{\Gamma }`$ admissible.
###### Remark 3.4.
If $`\mathrm{\Gamma }`$ is admissible then $`\psi (\mathrm{ker}\phi )G_{\mathrm{}}`$, since $`G_{\mathrm{}}`$ is the largest group with this property.
###### Proposition 3.5.
If the pair $`(\phi :AB,\psi :AC)`$, where $`A,B`$ and $`C`$ are finitely generated and $`\phi (A)`$ is finitely presented is stably compressible then it is $`\mathrm{}`$-compressible and there exists a subgroup $`\mathrm{\Gamma }G_{\mathrm{}}C`$ which is normally finitely generated within $`C`$ and such that $`\psi (\mathrm{ker}\phi )`$ is full in $`\mathrm{\Gamma }`$.
Conversely, if the pair $`(\phi :AB,\psi :AC)`$ is $`\mathrm{}`$-compressible and there exists a finitely generated subgroup $`\mathrm{\Gamma }G_{\mathrm{}}`$ such that $`\psi (\mathrm{ker}\phi )`$ is full in $`\mathrm{\Gamma }`$, then the pair is stably compressible.
###### Proof.
We set $`K=\mathrm{ker}\phi `$ in the sequel. We establish first:
###### Lemma 3.6.
If the pair $`(\phi :AB,\psi :AC)`$ is stably compressible then it is $`\mathrm{}`$-compressible.
###### Proof.
We will use a transfinite recurrence with the inductive steps provided by the next two lemmas. Set $`\beta :𝐅(r)C`$ for the morphism making the pair $`(\phi 1,\psi \beta )`$ strongly compressible.
###### Lemma 3.7.
If $`\beta (𝐅(r))G_i`$ and $`\psi (A)G_i`$ then $`\psi (A)G_{i+1}`$.
###### Proof.
By hypothesis $`\phi 1(\mathrm{ker}\psi \beta )\phi 1(A𝐅(r))`$. Alternatively, for any $`b\phi (A)BB𝐅(r)`$ there exists some $`xA𝐅(r)`$ such that $`\phi 1(x)=b`$ and $`\psi \beta (x)=1`$. One can write uniquely $`x`$ in normal form (see , Thm.1.2., p.175) as $`x=a_1f_1a_2f_2\mathrm{}a_mf_m`$ where $`a_jA,f_j𝐅(r)`$ are non-trivial (except maybe $`f_m)`$. Then $`\phi 1(x)=\phi (a_1)f_1\phi (a_2)f_2\mathrm{}\phi (a_m)f_m`$.
Since the normal form is unique in $`B𝐅(r)`$ one derives that $`x`$ has the following property. There exists a sequence $`p_0=1<p_1<\mathrm{}<p_lm`$ of integers for which
$$\phi (a_{p_j})=b_j1B,\text{ where }b_1b_2\mathrm{}b_j=b,$$
$$\phi (a_j)=1,\text{ for all }j\{p_0,p_1,\mathrm{},p_l\},$$
and
$$f_{p_j}f_{p_j+1}\mathrm{}f_{p_{j+1}1}=1,\text{ (for all }j\text{, with the convention }p_{l+1}=m\text{).}$$
Furthermore $`1=\psi \beta (x)`$ implies that (recall that $`K=\mathrm{ker}\phi `$)
$$1\psi (a_1K)\beta (f_1)\psi (K)\beta (f_2)\psi (K)\mathrm{}$$
$$\mathrm{}\beta (f_{p_11})\psi (a_{p_1}K)\beta (f_{p_1})\psi (K)\beta (f_{p_1+1})\psi (K)\mathrm{}\beta (f_m).$$
However each partial product starting at the $`p_j`$-th term and ending at the $`(p_{j+1}1)`$-th term is a product of conjugates of $`\psi (K)`$ by elements from the image of $`\beta `$:
$$\beta (f_{p_j})\psi (K)\beta (f_{p_j+1})\psi (K)\mathrm{}\psi (K)\beta (f_{p_{j+1}1})=$$
$$\underset{i=0}{\overset{p_{j+1}p_j1}{}}\left(\beta \left(\underset{k=0}{\overset{i}{}}f_{p_j+k}\right)\psi (K)\beta \left(\underset{k=0}{\overset{i}{}}f_{p_j+k}\right)^1\right)𝒩(\psi (K),G_i)=G_{i+1}.$$
We used above the inclusions $`\psi (K)\psi (A)G_i`$ and $`\beta (𝐅(r))G_i`$. Therefore
$$1\psi (a_1K)G_{i+1}\psi (a_{p_1}K)G_{i+1}\mathrm{}\psi (a_{p_l}K)G_{i+1}=$$
$$=\psi (aK)G_{i+1},$$
for any $`aA`$ such that $`\phi (a)=b`$. This implies that $`\psi (aK)G_{i+1}`$ and hence $`\psi (A)G_{i+1}`$. ∎
###### Lemma 3.8.
If $`\beta (𝐅(r))G_{i1}`$ and $`\psi (A)G_i`$ then $`\beta (𝐅(r))G_i`$.
###### Proof.
One can use the symmetry of the algebraic compressibility and then the argument from the previous lemma. Alternatively, choose $`f𝐅(r)B𝐅(r)`$ and some $`xA𝐅(r)`$ such that $`\phi 1(x)=f`$ and $`\psi \beta (x)=1`$. Using the normal form as above we find this time $`1\beta (f)𝒩(\psi (K),G_{i1})`$ hence $`\beta (𝐅(r))G_i`$. ∎
Using in an alternate way the two previous lemmas one gets lemma 3.6. ∎
###### Lemma 3.9.
Let $`\beta :𝐅(r)C`$ be a homomorphism such that $`(\phi 1,\psi \beta )`$ is strongly compressible. Set $`\beta (𝐅(r))=H`$. Then $`\psi (K)`$ is full in $`\psi (K)H`$. In particular if $`A,B,C`$ are finitely generated and $`\phi (A)`$ is finitely presented then the subgroup $`\mathrm{\Gamma }=\psi (K)H`$ is finitely generated.
###### Proof.
We already saw that $`HG_{\mathrm{}}`$. Set $`W(L;X)=\{xx=_ig_ix_ig_i^1,g_iX,,x_iL\}`$ for two subgroups $`L,XC`$. The proof we used to show that $`\psi (A)G_{\mathrm{}}`$ and $`HG_{\mathrm{}}`$ actually yields $`\psi (A)W(\psi (K),H)`$ and respectively $`HW(\psi (K),H)`$. We remark now that $`W(\psi (K),H)=𝒩(\psi (K),\psi (K)H)`$. The left inclusion is obvious. The other inclusion consists in writing any element $`gxg^1`$ with $`g\psi (K)H`$, $`x\psi (K)`$ as a product of conjugates by elements of $`H`$. This might be done by recurrence on the length of $`g`$, by using the following trick. If $`g=y_1a_1y_2a_2`$, $`a_i\psi (K),y_iH`$ then $`gxg^1=y_1(a_1y_2a_1^1)(a_1a_2xa_2^1a_1^1)a_1y_2a_1^1y_1^1`$.
Consequently the fact that $`HW(\psi (K),H)`$ implies
$$W(\psi (K),H)W(\psi (K),W(\psi (K),H))=𝒩(\psi (K),𝒩(\psi (K),\psi (K)H))$$
$$𝒩(\psi (K),\psi (K)H)=W(\psi (K),H),$$
hence all inclusions are equalities. Also $`\psi (K)H𝒩(\psi (K),\psi (K)H)`$ since both components $`\psi (K)`$ and $`H`$ are contained in $`𝒩(\psi (K),\psi (K)H)`$. This shows that $`𝒩(\psi (K),\psi (K)H)=\psi (K)H`$ hence $`\psi (K)`$ is full in $`\psi (K)H`$.
We take therefore $`\mathrm{\Gamma }=\psi (K)H`$. It suffices to show now that each of the groups $`K`$ and $`H`$ are finitely generated. $`H`$ is finitely generated since it is the image of $`𝐅(r)`$. Furthermore $`K`$ is finitely generated since $`A/K=\phi (A)`$ is finitely presented and $`A`$ is finitely generated. The theorem of Neumann (, p.52) shows that $`K`$ must be normally finitely generated. This proves the claim. ∎
###### Lemma 3.10.
Assume that $`\psi (K)`$ is full in $`\mathrm{\Gamma }G_{\mathrm{}}`$, where $`\mathrm{\Gamma }`$ is finitely generated. If the pair $`(\phi ,\psi )`$ is $`\mathrm{}`$-compressible then it is stably compressible.
###### Proof.
Consider $`r`$ big enough and a surjective homomorphism $`\beta :𝐅(r)\mathrm{\Gamma }`$. This implies that $`\psi \beta (A𝐅(r))=\psi (K)\mathrm{\Gamma }=\mathrm{\Gamma }`$. We have to show that any $`x\mathrm{\Gamma }`$ is in $`\psi \beta (\mathrm{ker}\phi 1)`$.
Recall that $`𝒩(\psi (K),\mathrm{\Gamma })=\mathrm{\Gamma }`$. Set $`{}_{}{}^{g}x=gxg^1`$. Then $`x=_i^{g_i}x_i`$ can be written as a product of conjugates of elements $`x_i\psi (K)`$ by elements $`g_i\mathrm{\Gamma }`$. Choose $`f_i𝐅(r)`$ so that $`\beta (f_i)=g_i`$ and $`y_iK`$ so that $`\psi (y_i)=x_i`$. Then $`\psi \beta (_if_i^1y_if_i)=x`$ and $`\phi 1(_if_i^1y_if_i)=_if_i^1\phi (y_i)f_i=1`$, since $`y_iK`$. ∎
Then the proposition 3.5 follows. ∎
### 3.2 End compressible manifolds
###### Definition 3.6.
The pair of spaces $`(T^{},T)`$ is strongly compressible (respectively stably compressible) if for each component $`S_j`$ of $`T`$ one chooses a component $`V_j`$ of $`T^{}\mathrm{int}(T)`$ such that $`S_jV_j`$ so that the pair $`(_j\pi _1(S_j)\pi _1(T),_j\pi _1(S_j)_j\pi _1(V_j))`$ is strongly compressible (respectively stably compressible). The morphisms are induced by the obvious inclusions. Similarly, the pair of spaces $`(T^{},T)`$ is said to be $`\alpha `$-compressible if the pair of morphisms from above is $`\alpha `$\- compressible. Set also $`G_{\mathrm{}}(T,T^{})`$ for the $`G_{\mathrm{}}`$ group associated to this pair of morphisms.
###### Remark 3.11.
These morphisms are not uniquely defined and depend on the various choices of base points in each component. However the compressibility does not depend on the particular choice of the representative.
###### Definition 3.7.
The open manifold $`W^n`$ is end compressible (respectively end $`k`$-compressible) if every exhaustion of $`W^n`$ by compact submanifolds $`T_i^n`$, such that $`\pi _1(T_i^n)\pi _1(T_i^n)`$ is a surjection, has a refinement
$$W^n=\underset{i=1}{\overset{\mathrm{}}{}}T_i^n,T_i^n\mathrm{int}(T_{i+1}^n),$$
such that:
1. all pairs $`(T_{i+1}^n,T_i^n)`$ are stably-compressible (respectively $`k`$-compressible).
2. if $`S_{i,j}^{n1}`$ denote the components of $`T_i^n`$ then the homomorphism $`_j:\pi _1(S_{i,j}^{n1})\pi _1(T_i^n)`$ induced by the inclusion is surjective.
3. any component of $`T_{i+1}^n\mathrm{int}(T_i^n)`$ intersect $`T_i^n`$ along precisely one component.
###### Remark 3.12.
As in the case of the compressibility the condition 2 above is independent of the homomorphism we chose, which might depend on the base points in each component.
###### Remark 3.13.
* One can ask that each connected component of $`T_{i+1}^n\mathrm{int}(T_i^n)`$ has exactly one boundary component from $`T_i^n`$. By adding to an arbitrary given $`T_i^n`$ the regular neighborhoods of arcs in $`T_{i+1}^n\mathrm{int}(T_i^n)`$ joining different connected component this condition will be fulfilled.
* Any simply-connected manifold $`W`$ of dimension at least $`5`$ has an exhaustion by $`T_i^n`$ that have the property that the natural maps $`\pi _1(T_i^n)\pi _1(T_i^n)`$ are surjective for all $`i`$. A proof will be given in the next section (see lemmas 3.26 and 3.28). Thus the above condition is never vacuous.
We shall henceforth assume that exhaustions have the property that the natural maps $`\pi _1(T_i^n)\pi _1(T_i^n)`$ are surjective for all $`i`$.
###### Remark 3.14.
The condition that the pair $`(T_{i+1}^n,T_i^n)`$ is stably-compressible is implied by (and later it will be proved that it is equivalent to) the pair of conditions
1. $`(T_{i+1}^n,T_i^n)`$ is $`\mathrm{}`$-compressible
2. There exists an admissible subgroup $`\mathrm{\Gamma }_i`$ of $`G_{\mathrm{}}(T_i^n,T_{i+1}^n)`$ which is finitely presented.
###### Theorem 3.15.
Any w.g.s.c. open $`n`$-manifold $`W^n`$ is end compressible. Conversely, for $`n4`$, $`W^n`$ is end compressible if and only if it is w.g.s.c.
###### Remark 3.16.
Notice that the end compressibility of $`W^3`$ implies that of $`W^3\times D^2`$. As a consequence of this result for $`n5`$ we will derive that $`W^3\times D^2`$ is w.g.s.c. and the invariance of the w.g.s.c. under proper homotopies (see theorem 5.2) will imply the result of the theorem for $`n=3`$. We will restrict then for the proof to $`n5`$.
###### Remark 3.17.
It is an important issue to know whether the stable-compressibility of one particular exhaustion implies the stable-compressibility of some refinement of any exhaustion. This is a corollary of our theorem 3.15 and proposition 3.18. In fact if $`W^n`$ is as above then $`W^n\times D^k`$ has one stably-compressible exhaustion. Take $`n+k5`$ to insure that $`W^k\times D^n`$ is w.g.s.c. and use the proposition 3.18. In particular any product exhaustion has a stably-compressible refinement, and the claim follows.
### 3.3 Proof of Theorem 3.15
Let us consider an exhaustion $`\{T_i^n\}_{i=1,\mathrm{}}`$ of $`W^n`$ by compact submanifolds, and fix some index $`i`$. The following result is the main tool in checking that specific manifolds are not w.g.s.c.
###### Proposition 3.18.
Any exhaustion as above of the w.g.s.c. manifold $`W^n`$ has a refinement for which consecutive terms fulfill the conditions:
1. all pairs $`(T_{i+1}^n,T_i^n)`$ are stably-compressible.
2. if $`S_{i,j}^{n1}`$ denote the components of $`T_i^n`$ the map $`\pi _1(S_{i,j}^{n1})\pi _1(T_i^n)`$ induced by the inclusion is surjective.
###### Proof.
Since $`W^n`$ is w.g.s.c. there exists a compact 1-connected submanifold $`M^n`$ of $`W^n`$ such that $`T_i^nM^n`$. We can suppose $`M^nT_{i+1}^n`$, without loss of generality. From now on we will focus on the pair $`(T_{i+1}^n,T_i^n)`$ and suppress the index $`i`$, and denote it $`(T^{},T)`$, for the sake of notational simplicity.
###### Lemma 3.19.
The pair $`(T^{},T)`$ is stably compressible.
###### Proof.
Let $`\phi :\pi _1(T)\pi _1(T)`$ and $`\psi :\pi _1(T)\pi _1(T^{}\mathrm{int}(T))`$ be the homomorphisms induced by the inclusions $`TT`$, and $`TT^{}\mathrm{int}(T)`$. If $`T`$ has several components then we choose base points in each component and set $`\pi _1(T)=_j\pi _1(S_j)`$ for notational simplicity.
Let us consider a handlebody decomposition of $`M^n\mathrm{int}(T)`$ (respectively a connected component) from $`T`$,
$$M^n\mathrm{int}(T)=\left(T\times [0,1]\right)\underset{\lambda =1}{\overset{n1}{}}\left(_jh_j^\lambda \right),$$
where $`h_j^\lambda `$ is a handle of index $`\lambda `$. We suppose the handles are attached in increasing order of their index. Since the distinct components of $`T`$ are not connected outside $`T`$ (by remark 3.13) the 1-handles which are added have the extremities in the same connected component of $`T`$. Set $`M_2^nM^n`$ (respectively $`M_1^n`$) for the submanifold obtained by attaching to $`T`$ only the handles $`h_j^\lambda `$ of index $`\lambda 2`$ (respectively those of index $`\lambda 1`$). Then $`\pi _1(M_2^n)=0`$, because adding higher index handles does not affect the fundamental group and we know that $`\pi _1(M^n)=0`$.
###### Lemma 3.20.
The pair $`(T^{},M_1^n)`$ is strongly compressible.
###### Proof.
Let $`\{\gamma _j\}_{j=1,p}M_1^n`$ be the set of attaching circles for the 2-handles of $`M_2^n`$ and $`\{\delta _j^2\}_{j=1,p}`$ be the corresponding core of the 2-handles $`h_j^2`$ ($`j=1,p`$). Since $`\delta _j^2`$ is a 2-disk embedded in $`M^n\mathrm{int}(T)`$ it follows that the homotopy class $`[\gamma _j]`$ vanishes in $`\pi _1(M^n\mathrm{int}(T))`$.
Let $`\mathrm{\Gamma }\pi _1(M_1^n)`$ be the normal subgroup generated by the homotopy classes of the curves $`\{\gamma _j\}_{j=1,p}`$ which are contained in $`M_1^n`$. Notice that this amounts to picking base points which are joined to the loops. Therefore the image of $`\mathrm{\Gamma }`$ under the map $`\pi _1(M_1^n)\pi _1(M_2^n\mathrm{int}(M_1^n))`$, induced by the inclusion, is zero. In particular its image in $`\pi _1(T^{}\mathrm{int}(M_1^n))`$ is zero.
On the other hand the images of the classes $`[\gamma _j]`$ in $`\pi _1(M_1^n)`$ normally generate all of the group $`\pi _1(M_1^n)`$ because $`\pi _1(M_2^n)=1`$.
These two properties are equivalent to the strong compressibility of the pair $`(M_2^n,M_1^n)`$ which in turn implies that of $`(T^{},M_1^n)`$. ∎
Rest of the proof of Lemma 3.19: Assume now that the number of 1-handles $`h_j^1`$ is $`r`$. Notice that $`\pi _1(M_1^n\mathrm{int}(T))\pi _1(T)𝐅(r)`$, because it can be obtained from $`T\times [0,1]`$ by adding 1-handles and the 1-handles we added do not join distinct boundary components, so that each one contributes with a free factor. In particular the inclusion $`TM_1^n\mathrm{int}(T)`$ induces a monomorphism $`\pi _1(T)\pi _1(T)𝐅(r)`$. Observe that $`M_1^n\mathrm{int}(T)`$ can also be obtained from $`M_1^n\times [0,1]`$ by adding $`(n1)`$-handles hence the inclusion $`M_1^nM_1^n\mathrm{int}(T)`$ induces the isomorphism $`\pi _1(M_1^n)\pi _1(T)𝐅(r)`$. The same reasoning gives the isomorphism $`\pi _1(M_1^n)\pi _1(T)𝐅(r)`$.
In particular we can view the subgroup $`\mathrm{\Gamma }`$ as a subgroup of $`\pi _1(T)𝐅(r)`$. The previous lemma tells us that $`\mathrm{\Gamma }`$ lies in the kernel of $`\pi _1(T)𝐅(r)\pi _1(T^{}\mathrm{int}(T))`$ and also projects epimorphically onto $`\pi _1(T)𝐅(r)`$. The identification of the respective maps with the morphisms induced by inclusions yields our claim. ∎
This finishes the proof of Proposition 3.18. ∎
Conversely, assume that $`W^n`$ has an exhaustion in which consecutive pairs are stably compressible. Then it is sufficient to show the following:
###### Proposition 3.21.
If $`(T^{},T)`$ is a stably compressible pair of $`n`$-manifolds and $`n5`$ then $`T\mathrm{int}(M^n)T^{}`$ where $`M^n`$ is a compact submanifold with $`\pi _1(M^n)=0`$.
###### Proof.
One can realize the homomorphism $`\beta :𝐅(r)\pi _1(T^{}\mathrm{int}(T))`$ by a disjoint union of bouquets of circles $`^rS^1T^{}\mathrm{int}(T)`$. There is one bouquet in each connected component of $`T^{}\mathrm{int}(T)`$. One joins each wedge point to the unique connected component of $`T`$ for which that is possible by an arc, and set $`M_1^n`$ for the manifold obtained from $`T`$ by adding a regular neighborhood of the bouquets $`^rS^1`$ in $`T^{}`$ (plus the extra arcs). This is equivalent to adding 1-handles with the induced framing.
###### Lemma 3.22.
The kernel $`\mathrm{ker}\psi \beta \pi _1(M_1^n)`$ is normally generated by a finite number of elements $`\gamma _1,\gamma _2,\mathrm{},\gamma _p`$.
###### Proof.
Consider a finite presentation $`𝐅(k)/H\pi _1(T)`$. We know that $`\pi _1(M_1^n)=\pi _1(T)𝐅(r)`$. Furthermore the composition map
$$\lambda :𝐅(k)𝐅(r)\pi _1(T)𝐅(r)\stackrel{\psi \beta }{}\mathrm{\Gamma },$$
is surjective (since $`\beta `$ is). The first map is the free product of the natural projection with the identity. Therefore $`𝐅(k+r)/\mathrm{ker}\lambda \mathrm{\Gamma }`$ is a presentation of the group $`\mathrm{\Gamma }`$. The Theorem of Neumann (see , p.52) states that any presentation on finitely many generators of a finitely presented group has a presentation on these generators with only finitely many of the given relations. Applying this to $`\mathrm{\Gamma }`$ one derives that there exist finitely many elements which normally generate $`\mathrm{ker}\lambda `$ in $`𝐅(k+r)`$. Then the images of these elements in $`\pi _1(T)𝐅(r)`$ normally generate $`\mathrm{ker}\psi \beta `$ (the projection $`𝐅(k)𝐅(r)\pi _1(T)𝐅(r)`$ is surjective). This yields the claim. ∎
###### Lemma 3.23.
The elements $`\gamma _i`$ are also in the kernel of $`\pi _1(M_1^n)\pi _1(T^{}\mathrm{int}(M_1^n))`$.
###### Proof.
The map $`\pi _1(T^{}\mathrm{int}(M_1^n))\pi _1(T^{}\mathrm{int}(T))`$ induced by the inclusion is injective because $`T^{}\mathrm{int}(T)`$ is obtained from $`T^{}\mathrm{int}(M_1^n)`$ by adding $`(n1)`$-handles (dual to the 1-handles from which one gets $`M_1^n`$ starting from $`T`$), and $`n5`$. Thus the map $`\pi _1(M_1^n)\pi _1(T^{}\mathrm{int}(T))`$ factors through
$$\pi _1(M_1^n)\pi _1(T^{}\mathrm{int}(M_1^n))\pi _1(T^{}\mathrm{int}(T)),$$
and any element in the kernel must be in the kernel of the first map, as stated. ∎
The dimension restriction $`n5`$ implies that we can assume $`\gamma _j`$ are represented by embedded loops having only the base point in common. Then $`\gamma _j`$ bound singular 2-disks $`D_j^2T^{}\mathrm{int}(M_1^n)`$. By a general position argument, one can arrange such that the 2-disks $`D_j^2`$ are embedded in $`T^{}\mathrm{int}(M_1^n)`$ and have disjoint interiors.
As a consequence the manifold $`M^n`$ obtained from $`M_1^n`$ by attaching 2-handles along the $`\gamma _j`$’s (with the induced framing) can be embedded in $`T^{}\mathrm{int}(T)`$. Moreover $`M^n`$ is a compact manifold whose fundamental group is the quotient of $`\pi _1(M_1^n)`$ by the subgroup normally generated by the elements $`\phi 1(\gamma _j)`$’s. The group $`\phi 1(\mathrm{ker}\psi \beta )`$ is normally generated by the elements $`\phi 1(\gamma _j)`$. By hypothesis the pair $`(\phi 1,\psi \beta )`$ is compressible hence $`\phi 1(\mathrm{ker}\psi \beta )`$ contains $`\phi 1(\pi _1(M_1^n))=\phi 1(\pi _1(T)𝐅(r))=\phi (\pi _1(T))𝐅(r)`$. Next
$$\frac{\pi _1(M_1^n)}{\phi 1(\mathrm{ker}\psi \beta )}=\frac{\pi _1(T)𝐅(r)}{\phi (\pi _1(T))𝐅(r)}\frac{\pi _1(T)}{\phi (\pi _1(T))}=1,$$
since $`\phi `$ has been supposed surjective. Therefore the quotient of $`\pi _1(M_1^n)`$ by the subgroup normally generated by the elements $`\phi 1(\gamma _j)`$ is trivial. ∎
### 3.4 End 1-compressibility is trivial for $`n5`$
We defined an infinite sequence of obstructions (namely $`k`$-compressibility for each $`k`$) to the w.g.s.c. However the first obstruction is trivial i.e. equivalent to the simple connectivity, in dimensions $`n4`$. In fact the main result of this section establishes the following:
###### Proposition 3.24.
End 1-compressibility and simple connectivity (s.c.) are equivalent for open $`n`$-manifolds of dimension $`n5`$.
###### Proof.
We first consider a simpler case:
###### Proposition 3.25.
The result holds in the case of a manifold $`W^n`$ of dimension at least $`5`$ with one end.
###### Proof.
In this case $`W^n`$ has an exhaustion $`T_i`$ with $`T_i`$ connected for all $`i`$.
###### Lemma 3.26.
$`W^n`$ has an exhaustion such that the map $`\phi :\pi _1(T_i)\pi _1(T_i)`$ induced by inclusion is a surjection for each $`i`$.
###### Proof.
As $`W^n`$ is simply connected, by taking a refinement we can assume that each inclusion map $`\pi _1(T_i)\pi _1(T_{i+1})`$ is the zero map. As usual we denote $`T_i`$ and $`T_{i+1}`$ by $`T`$ and $`T^{}`$ respectively.
Now, take a handle-decomposition of $`T`$ starting with the boundary $`T`$. Suppose the core of each $`1`$-handle of this decomposition is homotopically trivial in $`T`$, then it is immediate that $`\phi `$ is a surjection. We will enlarge $`T`$ by adding some $`1`$-handles and $`2`$-handles (that are embedded in $`T^{}`$) at $`T`$ in order to achieve this.
Namely, let $`\gamma `$ be the core of a $`1`$-handle. By hypothesis, there is a disc $`D^2`$ in $`T^{}`$ bounding the core of each of the $`1`$-handles, which we take to be transversal to $`T`$. As the dimension of $`W^n`$ is at least $`5`$, $`D^2`$ can be taken to be embedded. Notice that the 2-disks corresponding to all 1-handles can also be made disjoint, by general position. Thus $`D^2`$ intersects $`T^{}\mathrm{int}(T)`$ in a collection of embedded disjoint planar surfaces. The neighborhood of each disc component of this intersection can be regarded as a $`2`$-handle (embedded in $`T^{}`$) which we add to $`T`$ at $`T`$. For components of $`D^2(T^{}\mathrm{int}(T))=D^2\mathrm{int}(T)`$ with more than one boundary component, we take embedded arcs joining distinct boundary components. We add to $`T`$ a neighborhood of each arc, which can be regarded as a $`1`$-handle. After doing this for a finite collection of arcs, $`D^2D^2T`$ becomes a union of discs. Now we add $`2`$-handles as before. The disc $`D^2`$ that $`\gamma `$ bounds is now in $`T`$.
Further, the dual handles to the handles added are of dimension at least $`3`$. In particular we can extend the previous handle-decomposition to a new one for (the new) $`T`$ starting at (the new) $`T`$ with no new $`1`$-handles. Thus, after performing the above operation for the core of each $`1`$-handle of the original handle decomposition, the core of each $`1`$-handle of the resulting handle-decomposition of $`T`$ starting at $`T`$ bounds a disc in $`T`$. Thus $`\varphi `$ is a surjection. ∎
The above exhaustion is in fact $`1`$-compressible by the following algebraic lemma.
###### Lemma 3.27.
Suppose that we have a square of maps verifying the Van Kampen theorem:
$$\begin{array}{ccc}A& \stackrel{\psi }{}& C\\ \phi & & \gamma & & \\ B& \stackrel{\beta }{}& D\end{array}$$
Let $`\xi :AD`$. Suppose also that $`\phi `$ is surjective and $`\beta `$ is the zero map. Then $`\psi (A)𝒩(\psi (\mathrm{ker}\phi ))=𝒩(\psi (\mathrm{ker}\phi ),C)`$.
###### Proof.
Observe that $`\xi (\mathrm{ker}\phi )=0`$, hence $`\psi (\mathrm{ker}\phi )\mathrm{ker}\gamma `$. Hence we can define another diagram with $`A^{}=A/\mathrm{ker}\phi `$, $`B^{}=B`$, $`C^{}=C/𝒩(\psi (\mathrm{ker}\phi ))`$, $`D^{}=D`$:
$$\begin{array}{ccc}A^{}=A/\mathrm{ker}\phi & \stackrel{\psi }{}& C^{}=C/𝒩(\psi (\mathrm{ker}\phi ))\\ \phi & & \gamma & & \\ B^{}=B& \stackrel{\beta }{}& D^{}=D\end{array}$$
Notice that the map $`A/\mathrm{ker}\phi C/𝒩(\psi (\mathrm{ker}\phi ))`$ is well-defined. Again this diagram verifies the Van Kampen theorem. For this diagram, the induced $`\phi `$ is an isomorphism. It is immediate then that the universal (freest) $`D^{}`$ must be $`C^{}`$. In fact $`D^{}=C^{}A^{}/𝒩(\{\psi (a)a^1,aA^{}\})`$. Consider the map $`C^{}C^{}A^{}/𝒩(\{\psi (a)a^1,aA^{}\})C^{}`$, where the second arrow consists in replacing any occurrence of $`a`$ by the element $`\psi (a)`$ and taking the product in $`C^{}`$. This composition is the identity and the first map is a surjection, hence the map $`C^{}D^{}`$ is an isomorphism.
The map induced by $`\xi `$ is zero since $`\beta `$ is the zero. But the map $`\xi :A^{}D^{}`$ is the map $`\psi :A^{}C^{}`$ followed by an isomorphism, hence $`\psi (A^{})=0`$. This is equivalent to $`\psi (A)𝒩(\psi (\mathrm{ker}\phi ))`$. ∎
Note that the above lemma is purely algebraic, and in particular independent of dimension. The two lemmas immediately give us the proposition for one-ended manifolds $`W^n`$ with $`n5`$. ∎
The general case. We now consider the general case of a simply-connected open manifold $`W^n`$ of dimension at least $`5`$, with possibly more than one end. We shall choose the exhaustion $`T_i`$ with more care in this case.
We will make use of the following construction several times. Start with a compact submanifold $`A^n`$ of codimension $`0`$, with possibly more than one boundary component. Assume for simplicity (by enlarging $`A^n`$ if necessary) that no complementary component of $`A^n`$ is pre-compact. As $`W^n`$ is simply connected, we can find a compact submanifold $`B^n`$ containing $`A^n`$ in its interior such that the inclusion map on fundamental groups is the zero map. Further, we can do this by thickening and then adding the neighborhood of a $`2`$-complex, i.e., a collection of $`1`$-handles and $`2`$-handles. Namely, for each generator $`\gamma `$ of $`\pi _1(A^n)`$, we can find a disc $`D^2`$ that $`\gamma `$ bounds, and then add $`1`$-handles and $`2`$-handles as in lemma 3.26. Thus, as $`n5`$, the boundary components of $`B^n`$ correspond to those of $`A^n`$. We repeat this with $`B^n`$ in place of $`A^n`$ to get another submanifold $`C^n`$.
Observe that as a consequence of this and the simple-connectivity of $`W^n`$, the inclusion map $`\pi _1(A^nV^n)\pi _1(B^nV^n)`$ is the zero map for any component $`V^n`$ of $`W^nA^n`$. More generally if $`Z^nV^n`$, then $`\mathrm{ker}(\pi _1(A^nZ^n)\pi _1(B^nZ^n))=\mathrm{ker}(\pi _1(A^nZ^n)\pi _1(B^nZ^n(W^nV^n)))`$. Similar results hold with $`B^n`$ and $`C^n`$ in place of $`A^n`$ and $`B^n`$.
Now start with some $`A_1^n`$ as above and construct $`B_1^n`$ and $`C_1^n`$. Thicken $`C_1^n`$ slightly to get $`T=T_1`$. We will eventually choose a $`T_2=T^{}`$, but for now we merely note that it can (and so it will) be chosen in such a manner that the inclusion map $`\pi _1(T)\pi _1(T^{})`$ is the zero map. Let $`S_j,j=1,\mathrm{},n`$ be the boundary components of $`T`$. Let $`X_j`$ be the union of the component of $`T\mathrm{int}(B^n)`$ containing $`S_j`$ and $`B^n`$, and define $`Y_j`$ analogously with $`C^n`$ in place of $`B^n`$. Denote the image of $`\pi _1(X_j)`$ in $`\pi _1(Y_j)`$ by $`\overline{\pi }(X_j)`$. We then have a natural map $`\phi :\pi _1(S_j)\overline{\pi }(X_j)`$. Let $`V_j`$ be the component of $`T^{}\mathrm{int}(T)`$ containing $`S_j`$.
###### Lemma 3.28.
By adding $`1`$-handles and $`2`$-handles to $`S_j`$, we can ensure that $`\pi _1(S_j)`$ surjects onto $`\overline{\pi }(X_j)`$.
###### Proof.
The proof is essentially the same as that of lemma 3.26. We start with a handle-decomposition for $`X_j`$ starting from $`S_j`$. We shall ensure that the image in $`\overline{\pi }(X_j)`$ of the core of each $`1`$-handles is trivial. Namely, for each core, we take a disc $`D`$ that it bounds. By the above remarks, we can, and do, ensure that the disc lies in $`Y_j`$, and in particular does not intersect any boundary component of $`T`$ except $`S_j`$. As in lemma 3.26 we may now add $`1`$-handles and $`2`$-handles to $`S_j`$ to achieve the desired result. ∎
Notice that the changes made to $`T`$ in the above lemma do not affect $`S_k`$, $`X_k`$ and $`Y_k`$ for $`kj`$. Hence, by repeated application of the above lemma, we can ensure that all the maps $`\pi _1(S_j)\overline{\pi }(X_j)`$ are surjections. Also notice that the preceding remarks show that $`\mathrm{ker}\phi =\mathrm{ker}(\pi _1(S_j)\pi _1(T))`$. Now take $`A=\pi _1(S_j)`$, $`B=\overline{\pi }(X_j)`$ and $`C=\pi _1(V_j)`$, and let $`D`$ be the image of $`\pi _1(V_jX_j)`$ in $`\pi _1(V_jY_j)`$. Then, by the preceding remarks and lemma 3.28, the diagram
$$\begin{array}{ccc}A& \stackrel{\psi }{}& C\\ \phi & & \gamma & & \\ B& \stackrel{\beta }{}& D\end{array}$$
satisfies the hypothesis of the lemma 3.27. The 1-compressibility for the pair $`(T^{},T)`$ follows.
Now we continue the process inductively. Suppose $`T_k`$ has been defined, choose $`A_{k+1}`$ so that it contains $`T_k`$ and also in such a manner as to ensure that $`A_i`$’s exhaust $`M`$. Then find $`B_{k+1}`$, $`C_{k+1}`$ and $`T_{k+1}`$ as above. The rest follows as above. ∎
###### Remark 3.29.
For the case of simply-connected, one-ended (hence contractible) $`3`$-manifolds, a theorem of Luft says that $`M`$ can be exhausted by a union of homotopy handlebodies. These satisfy the conclusion of lemma 3.26, hence the proposition still holds. More generally, we can apply the sphere theorem to deduce that we have an exhaustion by connected sums of homotopy handlebodies. It follows that each pair $`(T,T^{})`$ of this exhaustion is $`1`$-compressible as we can decompose $`T`$ and consider each component separately without affecting $`\phi `$ or $`\psi `$.
## 4 Examples of contractible manifolds
### 4.1 Uncountably many Whitehead-type manifolds
Recall the following definition from :
###### Definition 4.1.
A Whitehead link $`T_0^nT_1^n`$ is a null-homotopic embedding of the solid torus $`T_0^n`$ in the (interior of the) unknotted solid torus $`T_1^n`$ lying in $`S^n`$ such that the pair $`(T_1^n,T_0^n)`$ is (boundary) incompressible.
The solid $`n`$-torus is $`T^n=D^2\times S^1\times S^1\mathrm{}\times S^1`$. By iterating the ambient homeomorphism which sends $`T_0^n`$ onto $`T_1^n`$ one obtains an ascending sequence $`T_0^nT_1^nT_2^n\mathrm{}`$ whose union is called a Whitehead-type $`n`$-manifold. A Whitehead-type manifold is open contractible and not s.c.i. D.G. Wright () gave a recurrent procedure to construct many Whitehead links in dimensions $`n3`$. One shows below that this construction provides uncountably many distinct contractible manifolds.
We introduce an invariant for pairs of solid tori which generalizes the wrapping number in dimension 3. Moreover this provides invariants for open manifolds of Whitehead-type answering a question raised in .
###### Definition 4.2.
A spine of the solid torus $`T^n`$ is an embedded $`t^{n2}=\{\}\times S^1\times S^1\times \mathrm{}\times S^1T^n`$ having a trivial normal bundle in $`T^n`$. This gives $`T^n`$ the structure of a trivial 2-disk bundle over $`t^{n2}`$.
###### Remark 4.1.
Although the spine is not uniquely defined, its isotopy class within the solid torus is.
Consider a pair of solid tori $`T_0^nT^n`$. We fix some spine $`t^{n2}`$ for $`T^n`$. To specify the embedding of $`T_0^n`$ is the same as giving the embedding of a spine $`t_0^{n2}`$ of $`T_0^n`$ in $`T^n`$. The isotopy class of the embedding $`t_0^{n2}T^n`$ is therefore uniquely defined by the pair. Let us pick-up a Riemannian metric $`g`$ on the torus $`T^n`$ such that $`T^n`$ is identified with the regular neighborhood of radius $`r`$ around $`t^{n2}`$. We denote this by $`T^n=t^{n2}[r]`$, and suppose for simplicity that $`r=1`$. Then $`t^{n2}[\lambda ]`$ for $`\lambda 1`$ will denote the radius $`\lambda `$ tube around $`t^{n2}`$ in this metric.
###### Definition 4.3.
The wrapping number of the Whitehead link $`T_0^nT^n`$ is defined as follows:
$$w(T^n,T_0^n)=\underset{\epsilon 0}{lim}\underset{t_0^n(t_0^nt^{n2}[\epsilon ])}{inf}\frac{\mathrm{vol}(t_0^{n2})}{\mathrm{vol}(t^{n2})},$$
where $`(t_0^nT^n)`$ is the set of all embeddings of the spine $`t_0^{n2}`$ of $`T_0^n`$ in the given isotopy class, and $`\mathrm{vol}`$ is the $`(n2)`$-dimensional volume.
###### Remark 4.2.
Notice that a priori this definition might depend on the particular choice of the spine $`t^{n2}`$ and on the metric $`g`$.
###### Proposition 4.3.
The wrapping number is a topological invariant of the pair $`(T^n,T_0^n)`$.
###### Proof.
There is a natural projection map on the spine $`\pi :T^nt^{n2}`$, which is the fiber bundle projection of $`T^n`$ (with fiber a 2-disk). When both $`T^n`$ and $`t^{n2}`$ are fixed then such a projection map is also defined only up to isotopy. Set therefore
$$l(T^n,T_0^n)=\underset{t_0^n(t_0^nT^n)}{inf}\underset{xt^{n2}}{inf}\mathrm{}\left\{\pi ^1(x)t_0^{n2}\right\}.$$
Since $`inf_{xt^{n2}}\mathrm{}\left\{\pi ^1(x)t_0^{n2}\right\}`$ does not depend on the particular projection map (in the fixed isotopy class) this number represent a topological invariant of the pair $`(T^n,T_0^n)`$. Hence the claim follows from the following result:
###### Proposition 4.4.
$`w(T^n,T_0^n)=l(T^n,T_0^n).`$
###### Proof.
Consider a position of $`t_0^{n2}`$ for which the minimum value $`l(T^n,T_0^n)`$ is attained. A small isotopy make $`t_0^{n2}`$ transversal to $`\pi `$. Then, for this precise position of $`t_0^{n2}`$ there exists some number $`M`$ such that
$$\mathrm{}\left\{\pi ^1(x)t_0^{n2}\right\}M,\text{ for any }xt^{n2}.$$
Denote by $`\mu `$ the Lebesgue measure on $`t^{n2}`$.
###### Lemma 4.5.
For any $`\epsilon >0`$ one can move $`t_0^{n2}`$ in $`T^n`$ by an ambient isotopy such that the following conditions are fulfilled:
$$\mathrm{}\left\{\pi ^1(x)t_0^{n2}\right\}M,\text{ for any }xt^{n2}.$$
$$\mu \left(\left\{xt^{n2}\right|\mathrm{}\left\{\pi ^1(x)t_0^{n2}\right\}>l(T^n,T_0^n)\}\right)<\epsilon .$$
###### Proof.
The set $`U=\left\{xt^{n2}\right|\mathrm{}\left\{\pi ^1(x)t_0^{n2}\right\}=l(T^n,T_0^n)\}`$ is an open subset of positive measure. Consider then a flow $`\phi _t`$ on the torus $`t^{n2}`$ which expands a small ball contained in $`U`$ into the complement of a measure $`\epsilon `$ set (e.g. a small tubular neighborhood of a spine of the 1-holed torus). Extend this flow as $`1_{D^2}\times \phi _t`$ all over $`T^n`$ and consider its action on $`t_0^{n2}`$. ∎
###### Lemma 4.6.
$`w(T^n,T_0^n)l(T^n,T_0^n).`$
###### Proof.
The map $`\pi `$ is the projection of the metric tube around $`t^{n2}`$ on its spine, hence the Jacobian $`Jac(\pi |_{t_0^{n2}})`$ has bounded norm $`|Jac(\pi |_{t_0^{n2}})|1`$. It follows that
$$\frac{\mathrm{vol}(t_0^{n2})}{\mathrm{vol}(t^{n2})}=\frac{\pi ^{}𝑑\mu }{𝑑\mu }\frac{_{\pi ^1(U)}|Jac(\pi |_{t_0^{n2}})|d\mu }{_U𝑑\mu }+M\epsilon $$
$$l(T^n,T_0^n)(1\epsilon )+M\epsilon ,$$
for any $`\epsilon >0`$, hence the claim. ∎
###### Lemma 4.7.
$`w(T^n,T_0^n)l(T^n,T_0^n).`$
###### Proof.
Set $`\lambda _t:t^{n2}[\delta ]t^{n2}[t\delta ]`$ for the map given in coordinates by $`\lambda _t(p,x)=(tp,x)`$, $`pD^2,xt^{n2}`$. Here the projection $`\pi `$ provides a global trivialisation of $`t^{n2}[\delta ]T^n`$. Then
$$\underset{t0}{lim}|Jac(\pi |_{t_0^{n2}}\lambda _t)|=1.$$
Therefore for $`t`$ close enough to 0 one derives
$$\underset{t0}{lim}\frac{\mathrm{vol}(\lambda _t(t_0^{n2}))}{\mathrm{vol}(t^{n2})}=\underset{t0}{lim}\frac{_{\lambda _t(t_0^{n2})}|Jac(\pi |_{t_0^{n2}})|d\mu }{_{t^{n2}}𝑑\mu }l(T^n,T_0^n).$$
Since the position of $`t_0^{n2}`$ was chosen arbitrary, this inequality survives after passing to the infimum and the claim follows. ∎
###### Theorem 4.8.
There exist uncountably many Whitehead-type manifolds for $`n5`$.
###### Proof.
The proof here follows the same pattern as that given by McMillan () for the 3-dimensional case. Let us establish first the following useful property of the wrapping number:
###### Proposition 4.9.
If $`T_0^nT_1^nT_2^n`$ then $`w(T_2^n,T_0^n)=w(T_2^n,T_1^n)w(T_1^n,T_0^n)`$.
###### Proof.
This is a consequence of the two lemmas below:
###### Lemma 4.10.
$`l(T_2^n,T_0^n)l(T_2^n,T_1^n)l(T_1^n,T_0^n)`$.
###### Proof.
Consider $`t_0^{n2}T_1^nT_2^n`$, where $`T_1^n`$ is a very thin tube around $`t_1^{n2}`$, and the two projections to $`\pi _2:t_1^{n2}t_2^{n2}`$ and $`\pi _1:t_0^{n2}t_1^{n2}`$ respectively. Using Lemma 4.5 one can assume that the conditions
$$\mu \left(\left\{xt_1^{n2}\right|\mathrm{}\left\{\pi _1^1(x)\right\}=l(T_1^n,T_0^n)\}\right)>1\epsilon ,$$
$$\mu \left(\left\{xt_2^{n2}\right|\mathrm{}\left\{\pi _2^1(x)\right\}=l(T_2^n,T_1^n)\}\right)>1\epsilon ,$$
hold. For small enough $`ϵ`$ one derives that
$$\mu \left(\left\{xt_2^{n2}\right|\mathrm{}\left\{(\pi _2\pi _1)^1(x)\right\}=l(T_2^n,T_1^n)l(T_1^n,T_0^n)\}\right)>0.$$
This proves that the minimal cardinal of the $`(\pi _2\pi _1)^1(x)`$ is not greater than $`l(T_2^n,T_1^n)l(T_1^n,T_0^n)`$, hence the claim. ∎
###### Lemma 4.11.
$`w(T_2^n,T_0^n)w(T_2^n,T_1^n)w(T_1^n,T_0^n)`$.
###### Proof.
We can assume that $`w(T_2^n,T_1^n)0`$. Consider an embedding of the $`(n2)`$-torus $`s_1^{n2}T_2=t_2^{n2}[\epsilon ]`$ for which the value of $`\frac{\mathrm{vol}(t_1^{n2})}{\mathrm{vol}(t_2^{n2})}`$ (as function of $`t_1^{n2}`$) is closed to the infimum in the isotopy class. We will assume that in all formulas below the tori lay in their respective isotopy classes. Then
$$\left(\underset{t_0^{n2}s_1^{n2}[2\epsilon ]}{inf}\frac{\mathrm{vol}(t_0^{n2})}{\mathrm{vol}(s_1^{n2})}\right)\left(\underset{t_1^{n2}t_2^{n2}[\epsilon ]}{inf}\frac{\mathrm{vol}(t_1^{n2})}{\mathrm{vol}(t_2^{n2})}\right)$$
$$\left(\underset{t_0^{n2}s_1^{n2}[2\epsilon ]}{inf}\frac{\mathrm{vol}(t_0^{n2})}{\mathrm{vol}(t_1^{n2})}\right)\frac{\mathrm{vol}(s_1^{n2})}{\mathrm{vol}(t_2^{n2})}=$$
$$=\underset{t_0^{n2}s_1^{n2}[2\epsilon ]}{inf}\frac{\mathrm{vol}(t_0^{n2})}{\mathrm{vol}(t_2^{n2})}\underset{t_0^{n2}t_2^{n2}[\epsilon ]}{inf}\frac{\mathrm{vol}(t_0^{n2})}{\mathrm{vol}(t_2^{n2})}.$$
The last inequality follows from the fact that $`s_1^{n2}[2\epsilon ]t_2^{n2}[\epsilon ]`$. In fact $`w(T_2^n,T_1^n)0`$ implies that $`s_1^{n2}`$ intersects any 2-disk $`D^2\times \{\}`$ (i.e. any fiber of the projection $`\pi _2:T_2^nt_2^{n2}`$) of $`T_2^n=t_2^{n2}[\epsilon ]`$ in at least one point. Then the transversal disk $`D^2\times \{\}`$ of radius $`\epsilon `$ is therefore contained in the tube $`s_1^{n2}[2\epsilon ]`$ of radius $`2\epsilon `$ around $`s_1^{n2}`$, establishing the claimed inclusion.
On the other hand the following holds
$$\underset{\epsilon 0}{lim}\underset{t_0^{n2}s_1^{n2}[2\epsilon ]}{inf}\frac{\mathrm{vol}(t_0^{n2})}{\mathrm{vol}(s_1^{n2})}=w(T_1^n,T_0^n),$$
due to the topological invariance of the wrapping number. Letting $`\epsilon `$ go to 0 in the previous inequality yields the claim. ∎
###### Proposition 4.12.
There exist Whitehead links whose wrapping number has the form $`2^{n2}p`$ for any natural number $`p`$.
###### Proof.
The claim is well-known for $`n=3`$. One uses Wright’s construction () of Whitehead links by induction on the dimension. If $`T_0^nT_1^n`$ is a Whitehead link then set $`T^{n+1}=T_1^n\times S^1`$. Consider the projection $`q`$ of the solid torus $`T_0^n\times S^1D^2\times S^1\times \mathrm{}\times S^1`$ onto $`D^2\times S^1`$ (the first and the last factors). Choose some Whitehead link $`L^3D^2\times S^1`$, and set then $`Q^{n+1}=q^1(L^3)`$. The pair $`Q^{n+1}T^{n+1}`$ is a Whitehead link of dimension $`n+1`$. The Proposition then is an immediate consequence of:
###### Lemma 4.13.
$`w(T^{n+1},Q^{n+1})=w(T_1^n,T_0^n)w(D^2\times S^1,L^3)`$.
###### Proof.
From the multiplicativity of $`w`$ and the triviality of the projection $`q`$ it is sufficient to prove that $`w(T_1^n\times S^1,T_0^n\times S^1)=w(T_1^n,T_0^n)`$. This formula can be checked directly using $`l`$ instead of $`w`$. ∎
###### Proposition 4.14.
For any sequence $`𝐩=p_0,p_1,\mathrm{}`$ of positive integers consider a Whitehead-type manifold $`W^n(𝐩)=_{k=1}^{\mathrm{}}T_k^n`$, where $`w(T_{k+1}^n,T_k^n)=2^{n2}p_k`$. If the sequences $`𝐩`$ and $`𝐪`$ have infinitely many non-overlapping prime factors then the manifolds $`W^n(𝐩)`$ and $`W^n(𝐪)`$ are not PL homeomorphic.
###### Proof.
The proof is similar to that of (, p.375). Set $`W^n(𝐩)=_{k=1}^{\mathrm{}}T_k^n`$, $`W^n(𝐪)=_{k=1}^{\mathrm{}}s_k^n`$, where $`T_k^n`$,$`\stackrel{~}{T}_k^n`$, are tori, as above. If $`h:W^n(𝐪)W^n(𝐩)`$ is a PL homeomorphism, there exist integers $`j,k`$ such that $`T_0^nint(h(\stackrel{~}{T}_j^n))`$, $`q_k`$ has a prime factor which occurs in $`𝐪`$ but not in $`𝐩`$, $`k>j+1`$ and $`h(\stackrel{~}{T}_k^n)\mathrm{int}(T_m^n)`$. We have therefore
$$w(T_m^n,T_0^n)=w((T_m^n,h(\stackrel{~}{T}_k^n))w(h(\stackrel{~}{T}_k^n),h(\stackrel{~}{T}_j^n))w(h(\stackrel{~}{T}_j^n,T_0^n)).$$
We have obtained a contradiction because $`q_k`$ divides $`w(h(\stackrel{~}{T}_k^n),h(\stackrel{~}{T}_j^n))`$ but not the left hand side (which is non-zero also). ∎
### 4.2 Open manifolds which are not w.g.s.c.
In general the tower of obstructions we defined in the previous sections is not trivial as is shown below:
###### Theorem 4.15.
For uncountably many Whitehead-type manifolds $`W^n`$ of dimension $`n3`$ the manifolds $`W^n\times N^k`$ are not $`\mathrm{}`$-compressible for any closed simply connected $`k`$-manifold $`N^k`$.
###### Proof.
It is sufficient to consider the case of the Whitehead-type manifolds since the pair of groups appearing in the product exhaustions are the same as this case.
We start with the 3-dimensional case, and take for $`W^3`$ the classical Whitehead manifold. Recall that $`W^3`$ is an increasing union of solid tori $`T_i`$, with $`T_i`$ embedded in $`T_{i+1}`$ as a neighborhood of a Whitehead link.
We shall first show that the pair $`(T_{i+1},T_i)`$ is not $`2`$-weakly compressible, and hence not stably compressible. We then extend this argument to show that any pair of the form $`(T_{i+n},T_i)`$ is not $`n+1`$-weakly compressible, and hence not stably compressible. By proposition 3.18, it follows that $`W^3\times N^k`$ is not $`\mathrm{}`$-compressible, and hence not w.g.s.c.
Let $`T`$ and $`T^{}`$ be as usual and let $`M^3=T^{}\mathrm{int}(T)`$, and fix a base point $`pT`$. Then $`C=G_0=\pi _1(M^3)`$ in our usual notation. Note that $`\mathrm{ker}(\phi )`$ is normally generated by the meridian of $`T`$ and hence $`\pi _1(M^3)/𝒩(\mathrm{ker}(\phi ),C)=\pi _1(T^{})=`$. Thus $`G_1=𝒩(\mathrm{ker}(\phi ),C)`$ consists of the homologically trivial elements in $`\pi _1(M^3)`$.
Consider now the cover $`\stackrel{~}{M^3}`$ of $`M^3`$ with fundamental group $`G_1`$. This is $`^3`$ with the neighborhood of an infinite component link, say indexed by the integers, deleted. Further each pair of adjacent components has linking number $`1`$. Pick a lift $`p^{}`$ of the base point $`p`$, which we use for all the fundamental groups we consider.
In this cover, $`\psi (A)`$ is the image of the bounding torus $`T`$ of the component of this link containing $`p^{}`$, and $`\psi (\mathrm{ker}(\phi ))`$ is generated by the meridian of this component. Thus, $`G_1/𝒩(\psi (\mathrm{ker}(\phi ),G_1)`$ is the fundamental group of $`\stackrel{~}{M^3}_TD^2\times S^1`$, i.e., of $`\stackrel{~}{M^3}`$ with a solid torus glued along $`T`$ to kill the meridian. But, because of the linking, the longitude $`\lambda T`$ is not trivial in this group, i.e. $`\lambda G_2=𝒩(\psi (\mathrm{ker}(\phi ),G_1)`$. Since $`\lambda \psi (A)`$, we see that the Whitehead link is not $`2`$-compressible.
We shall now consider a pair $`(T^{},T)`$ in some refinement of the given exhaustion. This is homeomorphic to a pair of the form $`(T_n,T_1)`$ for some $`n`$. As before pick a base point $`pT`$.
Let $`M_i=T_iT_1`$ and let $`M=M_n`$. Note that $`M_iM_{i+1}`$. In terms of earlier notation, $`M_n=T^{}T`$ and $`\pi _1(M)=C=G_0`$. Further $`\mathrm{ker}(\phi )`$ is normally generated by the meridian of $`T_1`$.
We have a sequence of subgroups $`G_kG_0=\pi _1(M)`$ and hence covers $`M^j`$ of $`M`$ corresponding to these subgroups. Pick lifts $`p^k`$ of the base point $`p`$ to these covers. Then $`𝒩(\psi (\mathrm{ker}(\phi ),G_k)`$ is generated by the meridian of the component of the inverse image of $`T_1`$ containing $`p^k`$. As the meridian is in $`\mathrm{ker}(\phi )`$, and each $`G_i`$ is the normal subgroup generated by $`\mathrm{ker}(\phi )`$ in $`G_{i1}`$, we see inductively that the lift of the meridian is a closed curve in $`M^k`$ so that the previous sentence makes sense.
Let $`N^i`$ be the result of gluing a solid torus or cylinder to $`M^i`$ along the component containing $`p^k`$ so that the meridian is killed. Then by the above $`G_k/𝒩(\psi (\mathrm{ker}(\phi ),G_k)=\pi _1(N^i)`$.
We shall prove by induction the following lemma.
###### Lemma 4.16.
$`T^{nk}`$ lifts to $`N^k`$, or equivalently, $`M_{nk}`$ lifts to $`M^k`$. Furthermore the longitude of the lift of $`T_{nk}`$ is a non-trivial element in $`\pi _1(N^k)`$.
###### Proof.
The case when $`k=1`$ is the above special case. Suppose now that the statement is true for $`k`$.
As the longitude of the lift of $`T_{nk}`$ is a non-trivial element in $`\pi _1(N^k)`$, in $`M^{k+1}`$ the inverse image of $`M_{nk}`$ is a cylinder with a sequence of linked lifts of $`M_{n(k+1)}`$ deleted. Thus, $`M_{n(k+1)}`$ lifts to $`M^{k+1}`$, and its longitude is linked with other lifts. It follows that the longitude of the lift of $`T_{n(k+1)}`$ is non-trivial in $`\pi _1(N^{k+1})`$. ∎
As a subgroup of $`G_{n1}`$, $`\psi (A)`$ is the image of the lift of $`T_1`$ containing the base point. As in the special case, as the longitude of this torus is a non-trivial element of $`\pi _1(N^{n1})`$, it follows that $`\psi (A)G^n`$. Thus $`(T_n,T_1)`$ is not $`n`$-compressible.
This ends the proof of the claim for the Whitehead manifold. Observe however that the same proof works for uncountably many similar manifolds – namely we may embed $`T_i`$ in $`T_{i+1}`$ as a link similar to the Whitehead link that winds around the solid torus several times.
We will use now a recurrence on the dimension and the results of the previous section in order to settle the higher dimensional situation. Consider for simplicity $`n=4`$ and a Whitehead-type manifold $`W^4`$ which is the ascending union of solid tori as in Wright’s construction. We use the notations from lemma 4.12 below. Then the pair of tori $`(T^4,Q^4)`$ is constructed out of the two Whitehead links in one dimension less $`(T_1^3,T_0^3)`$ and $`(D^2\times S^1,L^3)`$.
As above $`\mathrm{ker}(\phi )`$ is normally generated by the meridian and $`G_1`$ consists of homologically trivial elements of $`\pi _1(T^4\mathrm{int}(Q^4))`$, by using Van Kampen and the fact that $`\pi _1(T^4)`$ is abelian. The cover $`\stackrel{~}{M^4}`$ of $`M^4=T^4\mathrm{int}(Q^4)`$ with fundamental group $`G_1`$ is $`𝐑^4`$ with a thick infinite link deleted. There is an obvious $`𝐙^2`$ action on the components of this link, and so we can label the boundary tori as $`T_{i,j}`$, for integer $`i,j`$.
Let $`\lambda `$ be the longitude curve having the parameters $`(n,k)`$ on the torus $`T_{0,0}`$. Then one can compute the linking numbers $`\mathrm{lk}(\lambda ,T_{0,1})=k`$ and $`\mathrm{lk}(\lambda ,T_{1,0})=n`$. This follows because both links used in the construction were the standard Whitehead link. Variations which yield non-zero linking numbers are also convenient for our purposes. Consequently for non-zero $`n,k`$ we obtained an element which is non-trivial in $`G_2`$ hence the pair of solid tori is not $`2`$-compressible. A similar argument goes through the higher compressibility as well. Using suitable variations in choosing the links and mixing the pairs of solid tori as in previous section yields uncountably many examples as in the theorem. ∎
###### Remark 4.17.
It follows that $`W^3\times D^k`$ is not w.g.s.c. using the criterion from . However the previous theorem is more precise regarding the failure of g.s.c. for these product manifolds.
## 5 The proper homotopy invariance of the w.g.s.c.
### 5.1 Dehn exhaustibility
We study in this section to what extent the w.g.s.c. is a proper homotopy invariant.
###### Definition 5.1.
A polyhedron $`M`$ is (proper) homotopically dominated by the polyhedron $`X`$ if there exists a map $`f:MX`$ such that the mapping cylinder $`Z_f`$ (properly) retracts on $`M`$.
###### Remark 5.1.
A proper homotopy equivalence is the simplest example of a proper homotopically domination.
The main result of this section is:
###### Theorem 5.2.
For $`n4`$ a non-compact $`n`$-manifold is w.g.s.c. if and only if it is proper homotopically dominated by a w.g.s.c. polyhedron.
###### Remark 5.3.
It seems that the result does not hold, as stated, for $`n=4`$ (see also the next section).
Proof of Theorem 5.2. The main ingredient of the proof is the following notion, weaker than the w.g.s.c., introduced by Poénaru:
###### Definition 5.2.
The simply-connected $`n`$-manifold $`W^n`$ is Dehn exhaustible if, for any compact $`KW^n`$ there exists some simply connected compact polyhedron $`L`$ and a commutative diagram
$$\begin{array}{ccc}K& \stackrel{f}{}& L\\ & i& g\\ & & W^n\end{array}$$
where $`i`$ is the inclusion, $`f`$ is an embedding, $`g`$ is an immersion and $`f(K)M_2(g)=\mathrm{}`$. Here $`M_2(g)`$ is the set of double points, namely $`M_2(g)=\{xL;\mathrm{}g^1(g(x))2\}L`$. If $`n=3`$ then one asks the map $`g`$ to be a generic immersion, which means here that it has no triple points.
The first step is to establish:
###### Proposition 5.4.
An open simply-connected manifold which is proper homotopically dominated by a w.g.s.c. polyhedron is Dehn exhaustible.
###### Proof.
The proof given in for the 3-dimensional statement extends without any essential modification, and we skip the details. ∎
###### Remark 5.5.
Poénaru proved a Dehn-type lemma (see , p.333-339) which states that a Dehn exhaustible 3-manifold is w.g.s.c. This settles the dimension 3 case.
###### Lemma 5.6.
If the open simply-connected $`n`$-manifold $`W^n`$ is Dehn exhaustible and $`n5`$ then it is w.g.s.c.
###### Proof.
Consider a connected compact submanifold $`KW^n`$. Assume that there exists a compact polyhedron $`M^n`$ with $`\pi _1(M^n)=0`$ and an immersion $`F`$
$$\begin{array}{ccc}K& \stackrel{f}{}& M\\ & i& F\\ & & W^n\end{array}$$
such that $`M_2(F)K=\mathrm{}`$.
###### Lemma 5.7.
One can suppose that $`M^n`$ is a manifold.
###### Proof.
The polyhedron $`M^n`$ is endowed with an immersion $`F`$ into the manifold $`W^n`$. Among all abstract regular neighborhoods (i.e. thickenings) of $`M^n`$ there is a $`n`$-dimensional one $`U(M^n,F)`$, which is called the regular neighborhood determined by the immersion, such that the following conditions are fulfilled:
1. $`F:M^nW^n`$ extends to an immersion $`\stackrel{~}{F}:U(M^n,F)W^n`$.
2. The image of $`\stackrel{~}{F}(U(M^n,F))W^n`$ is the regular neighborhood of the polyhedron $`F(M^n)`$ in $`W^n`$.
The construction of the PL regular neighborhood determined by an immersion of polyhedra is given in . The authors were building on the case of an immersion of manifolds, considered previously in . Moreover, if one replaces $`M^n`$ by the manifold $`U(M^n,F)`$ and $`F`$ by $`\stackrel{~}{F}`$ we are in the conditions required by the Dehn-type lemma. ∎
Consider now a handlebody decomposition of $`M^nf(K)`$ and let $`N_2^n`$ be the union of $`f(K)`$ with the handles of index 1 and 2. Then $`\pi _1(N_2^n)=0`$. Let $`\delta _j^2,\delta _j^1`$ be the cores of these extra 1- and 2-handles. By using a small homotopy of $`F`$ one can replace $`F(\delta _j^2)W^n`$ by some embedded 2-disks $`d_j^2W^n`$ with the same boundary. Also by general position these 2-disks can be chosen to have disjoint interiors. Both assertions follow from the assumptions $`n5`$, and $`M_2(F)f(K)=\mathrm{}`$. This implies that the restriction of the new map $`F^{}`$, obtained by perturbing $`F`$, to $`\delta _j^2`$ (and $`\delta _j^1`$) is an embedding into $`W^nK`$. Using the uniqueness of the regular neighborhood it follows that $`F^{}`$ can be chosen to be an embedding on $`N_2^n`$. In particular $`K`$ is engulfed in the 1-connected compact $`F^{}(N_2^n)`$. ∎
### 5.2 Dehn exhaustibility and end compressibility in dimension 4
###### Proposition 5.8.
An open 4-manifold is end compressible if and only if it is Dehn exhaustible.
###### Proof.
We have to reconsider the proof of Proposition 3.21. Everything works as above except that the disks $`\delta _j^2`$ cannot be anymore embedded, but only (generically) immersed. They may have finitely many double points in their interior. Then the manifold $`M^4`$ obtained by adding 2-handles along the $`\gamma _j`$ has a generic immersion $`F:MT^{}`$, whose double points $`M_2(F)`$ are outside of $`T`$. This implies that $`W^4`$ is Dehn exhaustible.
Conversely assume that $`W^4`$ is Dehn exhaustible. Let $`K^4`$ be a compact submanifold of $`W^4`$ and $`M^4`$ be the immersible simply connected polyhedron provided by the Dehn exhaustibility property. Lemma 5.7 allow us to assume that $`M^4`$ and $`F(M^4)`$ are 4-manifolds. Consider now to the proof of the first claim from Theorem 5.2. It is sufficient to consider the case when $`M^4=M_2^4`$ i.e. $`M^4`$ is obtained from $`K^4`$ by adding 1- and 2-handles. If $`\mathrm{\Gamma }\pi _1(M_1^4)`$ is the normal subgroup generated by the attaching curves of the 2-handles of $`M^4`$ then the same argument yields:
$$\mathrm{\Gamma }\mathrm{ker}(\pi _1(M_1^4)\pi _1(M^4\mathrm{int}(K^4))).$$
Since $`F`$ is a generic immersion we can suppose that $`F`$ is an embedding of the cores of the 1-handles and so $`F|_{M_1^4}`$ is an embedding.
We have $`F(M^4\mathrm{int}(K^4))F(M^4)\mathrm{int}(K^4)`$ because the double points of $`F`$ are outside $`K^4`$. Now the homomorphism induced by $`F`$ on the left side of the diagram
$$\begin{array}{cccc}\pi _1(M_1^4)& & \pi _1(M^4\mathrm{int}(K^4))& \\ F& & F& \\ \pi _1(F(M_1^4))& & \pi _1(F(M^4)\mathrm{int}(K^4))& \end{array}$$
is an isomorphism and we derive that
$$F(\mathrm{\Gamma })\mathrm{ker}(\pi _1(F(M_1^4))\pi _1(F(M^4)\mathrm{int}(K^4))).$$
Meantime $`F(\mathrm{\Gamma })`$ surjects onto $`\pi _1(F(M_1^4))`$ under the map $`\pi _1(F(M_1^4))\pi _1(F(M_1^4))`$. But $`F(M_1^4)`$ is homeomorphic to $`M_1^4`$, hence it is obtained from $`K^4`$ by adding 1-handles. This shows that the pair $`(F(M^4),K^4)`$ is stably compressible, from which one obtains the end compressibility of $`W^4`$ as in the proof of Theorem 3.15. ∎
###### Remark 5.9.
If the open 4-manifold $`W^4`$ is Dehn exhaustible then $`W^4\times [0,1]`$ is also Dehn exhaustible hence w.g.s.c. Therefore an example of an open 4-manifold $`W^4`$ which is end compressible but which is not w.g.s.c. will show that the result of Theorem 5.2 cannot be extended to dimension 4, as stated. Such examples are very likely to exist, as the Dehn lemma is known to fail in dimension 4 (by S. Akbulut’s examples).
### 5.3 Proper-homotopy invariance of the end compressibility
It would be interesting to have a soft version of the theorem 5.2 for the end compressibility situation. Notice that the definition of the end compressibility extends word by word to non-compact polyhedra. One uses instead of the boundary of manifolds the frontier of a polyhedron.
###### Remark 5.10.
One expects that the following be true. If a polyhedron $`M`$ is proper homotopically dominated by an end compressible polyhedron $`X`$ then $`M`$ is also end compressible. The only ingredient lacking for the complete proof is the analogue of the remark 3.17 for polyhedra: if one exhaustion is stably-compressible then all exhaustions have stably-compressible exhaustions.
Along the same lines we have:
###### Proposition 5.11.
If there is a degree one map $`f:X^nM^n`$ between one ended manifolds of the same dimension, then if $`X^n`$ is end-compressible so is $`M^n`$.
###### Proof.
We use the fact that degree-one maps are surjective on fundamental group. Given an exhaustion $`\{L_j\}`$ of $`M^n`$, pull it back to $`\{f^1(L_j)\}`$ of $`X^n`$. Notice that $`f^1(L_j)=f^1(L_j)`$ where $``$ stands for the frontier.
One needs then the following approximation by manifolds result. Given two $`n`$-complexes $`K_1\mathrm{int}(K_2)𝐑^n`$ there exist regular neighborhoods $`K_j^\epsilon 𝐑^n`$ such that $`K_j\mathrm{int}(K_j^\epsilon )`$, $`K_j`$ is homotopy equivalent to $`K_j^\epsilon `$, $`K_2\mathrm{int}(K_1)`$ is homotopy equivalent to $`K_2^\epsilon K_1^\epsilon `$, and moreover $`K_1^\epsilon `$ is homotopy equivalent to $`K_1`$. This uses essentially the fact that $`K_j`$ are of codimension zero in $`𝐑^n`$.
Now, the hypothesis applied to the approximating exhaustion consisting of submanifolds implies the existence of a stably-compressible refinement of $`\{f^1(L_j)\}`$. Since degree-one maps are surjective on fundamental groups the lemma below permits to descend to $`M^n`$.
###### Lemma 5.12.
Suppose that the triple $`(A,B,C)`$ of groups surjects onto $`(A^{},B^{},C^{})`$ i.e. we have three surjections with diagrams commuting. Then if $`(A,B,C)`$ is strongly (or stably) compressible then so is $`(A^{},B^{},C^{})`$.
The proof is straightforward.
The only subtlety above is to make sure the inverse image of boundary components is connected (else we can connect them up in the one-ended case). ∎
###### Remark 5.13.
In the many-ended case, we need to say that we have a degree-one map between each pair of ends (not 2 ends mapping to one, with one of them having degree 2 and the other -1). This holds in particular for a proper map that has degree one and is injective on ends.
## 6 G.s.c. for 4-manifolds
### 6.1 W.g.s.c. versus g.s.c.
###### Definition 6.1.
A geometric Poénaru-Mazur-type manifold $`M^4`$ is a compact simply connected $`4`$-manifold satisfying the following conditions:
1. $`H_2(M^4)=0`$.
2. the boundary $`M^4`$ is connected and $`\pi _1`$-dominates a virtually geometric 3-manifold group, i.e. there exists a surjective homomorphism
$$\pi _1(M^4)\pi _1(N^3),$$
onto the (non-trivial) fundamental group of a virtually geometric 3-manifold $`N^3`$.
###### Proposition 6.1.
The interior $`\mathrm{int}(M^4)`$ of a geometric Poénaru-Mazur-type manifold $`M^4`$ does not have a proper handlebody decomposition without 1-handles with the boundary of a cofinal subset of the intermediate manifolds obtained on a finite number of handle additions being homology spheres.
### 6.2 Casson’s proof of Proposition 6.1
The main ingredient is the following proposition extending an unpublished result of A. Casson:
###### Proposition 6.2.
Consider the 4-dimensional (compact) cobordism $`(W^4,M^3,N^3)`$ such that $`(W^4,M^3)`$ is 1-connected. Assume moreover that the following conditions are satisfied:
1. $`H_2(W^4,M^3;𝐐)=0`$, both $`M^3`$ and $`N^3`$ are connected.
2. $`\pi _1(N^3)`$ is a group which $`\pi _1`$-dominates a virtually geometric non-trivial 3-manifold group. Let $`K`$ be the kernel of this epimorphism.
3. $`b_1(W^4)b_1(N^3)`$, where $`b_1`$ denotes the first Betti number.
4. The map $`\pi _1(N^3)\pi _1(W^4)`$ induced by the inclusion $`N^3W^4`$ has kernel strictly bigger than the subgroup $`K`$. In particular this is true if this map is trivial.
Then any handlebody decomposition of $`W^4`$ from $`M^3`$ has 1-handles i.e. the pair $`(W^4,M^3)`$ is not g.s.c.
###### Remark 6.3.
A necessary condition for the g.s.c. of $`(W^4,M^3)`$ is that the map $`\pi _1(M^3)\pi _1(W^4)`$, induced by the inclusion $`M^3W^4`$, be onto. In fact adding 2-handles amounts to introducing new relations to the fundamental group of the boundary, whereas the latter is not affected by higher dimensional handle additions.
Casson’s result was based on partial positive solutions to the Kervaire Conjecture 1.14. One proves that certain $`4`$-manifolds $`(N,N)`$ have no handle decompositions without $`1`$-handles by showing that if they did, then $`\pi _1(N)`$ violates the Kervaire conjecture. Casson’s argument works to the extent that the Kervaire conjecture is known to be true. Casson originally applied it using a theorem of M. Gerstenhaber and O.S. Rothaus (), which said that the Kervaire conjecture holds for subgroups of a compact Lie group. Subsequently, O.S. Rothaus () showed that the conjecture in fact holds for residually finite groups. Since residual finiteness for all $`3`$-manifold groups is implied by the geometrization conjecture, Casson’s argument works in particular for all manifolds satisfying the geometrization conjecture. A simple argument (Remark 6.4 below) extends the class of groups for which the Kervaire conjecture is known further.
###### Proposition 6.4.
If some non-trivial quotient $`Q`$ of a group $`G`$ satisfies the Kervaire conjecture, then so does $`G`$. In particular if a finitely generated group $`G`$ has a proper finite-index subgroup, then G satisfies the Kervaire conjecture (since finite groups satisfy the Kervaire conjecture by ).
###### Proof.
Let $`\varphi :GQ`$ be the quotient map. Assume that $`Q`$ satisfies the Kervaire conjecture. Suppose that G violates the Kervaire conjecture. Then we have generators $`\alpha _1,\mathrm{},\alpha _n`$ and relations such that $`\frac{G\alpha _1,\mathrm{},\alpha _n}{r_1,\mathrm{},r_n}`$ is the trivial group. Let $`\varphi ^{}:G\alpha _1,\mathrm{},\alpha _nQ\overline{\alpha _1},\mathrm{},\overline{\alpha _n}`$ be the map extending $`\varphi `$ by mapping $`\alpha _i`$ to $`\overline{\alpha _i}`$. This is clearly a surjection, and induces a surjective map $`\overline{\varphi }:\frac{G\alpha _1,\mathrm{},\alpha _n}{r_1,\mathrm{},r_n}\frac{Q\overline{\alpha _1},\mathrm{},\overline{\alpha _n}}{\varphi ^{}(r_1),\mathrm{},\varphi ^{}(r_n)}`$. But since the domain of the surjection $`\overline{\varphi }`$ is trivial, so is the codomain. But this means that $`\frac{Q\overline{\alpha _1},\mathrm{},\overline{\alpha _n}}{\varphi ^{}(r_1),\mathrm{},\varphi ^{}(r_n)}`$ is trivial, and so $`Q`$ violates the Kervaire conjecture, a contradiction. ∎
###### Proof of Proposition 6.2.
Suppose that
$$W^4=M^3\times [0,1]\mathrm{}_k2\text{-handles}\mathrm{}_r3\text{-handles},$$
(with some 0-handle or 4-handle added if one boundary component is empty). It is well-known (see ) that the homology groups $`H_{}(W^4,M^3)`$ are the same as those of a differential complex $`C_{}`$, whose component $`C_j`$ is the free module generated by the $`j`$-handles. Therefore this complex has the form:
$$0𝐙^r𝐙^k0.$$
Thus $`H_2(W^4,M^3;𝐐)=0`$ implies that $`kr`$ holds.
Consider now the handlebody decomposition is turned up-side down:
$$W^4=N^3\times [0,1]\mathrm{}_r1\text{-handles }\mathrm{}_k2\text{-handles},$$
(plus possibly one 0-handle or 4-handle if the respective boundary component is empty).
By the van Kampen theorem it follows that $`\pi _1(W^4)`$ is obtained from $`\pi _1(N^3)=\pi _1(N^3\times [0,1])`$ by adding one generator for each 1-handle and one relation for each 2-handle. Therefore
$$\pi _1(W^4)=\pi _1(N^3)𝐅(r)/W(k),$$
where $`𝐅(r)`$ is the free group on $`r`$ generators $`x_1,\mathrm{},x_r`$ and $`W(k)`$ is a normal subgroup of the free product generated also by $`k`$ words $`Y_1,\mathrm{},Y_k`$.
Consider a virtually geometric 3-manifold $`L^3`$ such that $`\pi _1(M^3)\pi _1(L^3)`$ is surjective. If $`L^3`$ is a geometric 3-manifold then its fundamental group is residually finite (see e.g. , Thm.3.3, p.364). Let $`d_{ij}`$ be the degree of the letter $`x_j`$ in the word representing $`Y_i`$. The result of Rothaus (, Thm. 18, p.611) states that for any locally residually finite group $`G`$ and choice of words $`Y_i`$ such that $`𝐝=(d_{ij})_{i,j}`$ is of (maximal) rank $`k`$, the natural morphism $`GG𝐅(r)/W(k)`$ is an injection. We have therefore a commutative diagram
$$\begin{array}{ccc}\pi _1(M^3)& & \pi _1(W^4)\\ & & \\ \pi _1(L^3)& & \pi _1(L^3)𝐅(r)/W(k)\end{array}$$
whose vertical arrows are surjections. The kernel of the map induced by inclusion, $`\pi _1(M^3)\pi _1(W^4)`$ is contained in $`K`$. This contradicts our hypothesis.
On the other hand if the rank of $`𝐝`$ is not maximal then by considering the abelianisations one derives $`H_1(G𝐅(r)/W(k))H_1(G)𝐙`$, hence $`b_1(W^4)b_1(N^3)+1`$, which is also false. ∎
###### Corollary 6.5.
Consider a 4-manifold $`W^4`$ which is compact connected simply-connected with non-simply connected boundary $`M`$. If the boundary is (virtually) geometric and $`H_2(W^4)=0`$ then $`W^4`$ is not g.s.c.
###### Proof of Proposition 6.1.
Assume now that $`\mathrm{int}(M^4)`$ admitted a proper handlebody decomposition without 1-handles. One identifies $`\mathrm{int}(M^4)`$ with $`M^4_{MM\times \{0\}}M\times [0,1)`$. We can truncate the handle decomposition at a finite stage in order to obtain a manifold $`Q^4`$ such that $`Q^4M^4\times (0,1)`$, because the decomposition is proper. We can suppose that $`Q^4`$ is connected since $`\mathrm{int}(M^4)`$ has one end. Then $`Q^4`$ is g.s.c. hence $`\pi _1(Q^4)=0`$.
By hypothesis, we can choose $`Q^4`$ to be a homology sphere. Then $`Q^4`$ separates the cylinder $`M^4\times [0,1]`$ into two manifolds with boundary which, by Mayer-Vietoris, have the homology of $`S^3`$. This implies that $`H_2(Q^4)=0`$ (again by Mayer-Vietoris).
Let us consider now the map $`f:Q^4M^4\times [0,1]M^4`$, the composition of the inclusion with the obvious projection.
###### Lemma 6.6.
The map $`f`$ has degree one hence induces a surjection on the fundamental groups.
###### Proof.
The 3-manifold $`Q^4`$ separates the two components of the boundary. In particular the generic arc joining $`M^4\times \{0\}`$ to $`M^4\times \{1\}`$ intersects transversally $`Q^4`$ in a number of points, which counted with the sign sum up to 1 (or -1). If properly interpreted this is the same as claiming the degree of $`f`$ is one.
It is well-known that a degree one map between orientable 3-manifolds induce a surjective map on the fundamental group (more generally, the image of the homomorphism induced by a degree $`d`$ is a subgroup whose index is bounded by $`d`$). ∎
This shows that $`\pi _1(Q^4)\pi _1(M^4)`$ is surjective. On the other hand $`\pi _1(M^4)`$ surjects onto a non-trivial residually finite group. Since $`\pi _1(Q^4)\pi _1(Q^4)=1`$ is the trivial map, the argument we used previously (from Rothaus’ theorem) gives us a contradiction. This settles our claim. ∎
## 7 Handle decompositions without 1-handles in dimension 4
### 7.1 Open tame 4-manifolds
###### Definition 7.1.
An exhaustion of a 4-manifold is g.s.c. if it corresponds to a proper sequence of handle additions with no 1-handles. Alternatively one has a proper Morse function, which we will refer to as time, with words like past and future having obvious meanings, with no critical points of index one. The inverse images of regular points are $`3`$-manifolds, which we refer to as the manifold at that time.
We assume henceforth that we have a g.s.c. handle decomposition of the interior $`\mathrm{int}(W^4)`$ of $`(W^4,W^4)`$, a compact four manifold with boundary a homology $`3`$-sphere and $`H_2(W^4)=0`$.
Now let $`(K_i^4,K_i^4),i`$ denote the $`4`$-manifolds obtained by successively attaching handles to the zero handle $`(B^4,S^3)`$, that is if $`t:(W^4,W^4)`$ is the Morse function time, then $`(K_i^4,K_i^4)=t^1((\mathrm{},a_i])`$, with $`a_i`$ being points lying between pairs of critical values of the Morse function.
###### Lemma 7.1.
$`K_{i+1}^4`$ is obtained from $`K_i^4`$ by one of the following:
* A 0-frame surgery about a homologically trivial knot in $`K_i^4`$.
* Cutting along a non-separating $`2`$-sphere in $`K_i^4`$ and capping off the result by attaching a $`3`$-ball.
These correspond respectively to attaching $`2`$-handles and $`3`$-handles to $`(K_i^4,K_i^4)`$.
###### Proof.
Since attaching $`2`$-handles and $`3`$-handles correspond to surgery and cutting along $`2`$-spheres respectively, we merely have to show that the surgery is 0-frame about a homologically trivial curve and the spheres along which one cuts are non-separating.
First note that the absence of $`1`$-handles implies $`H_1(K_i^4)=0=\pi _1(K_i^4)`$, for all $`i`$. Further, each $`K_i^4`$ is connected because $`\mathrm{int}(W^4)`$ has one end. Thus the $`2`$-spheres along which any $`K_i^4`$ is split have to be non-separating.
Using Mayer-Vietoris, the fact that $`H_2(W^4)=0`$, and the long exact sequence in homology we derive that $`H_2(K_i^4)=H_2(K_i^4)`$. Also adding a $`3`$-handle decreases the rank of $`H_2(K_i^4)`$ by one hence every surgery increases the rank of $`H_2(K_i^4)`$ by one unit. But this means that the surgery must be a zero-frame surgery about a homologically trivial curve. ∎
For $`i`$ large enough, $`K_i^4`$ lies in a collar $`W^4\times [0,\mathrm{})`$ hence we have a map $`f_i:K_i^4W^4`$ which is the composition of the inclusion with the projection. By the Lemma 5.1 the maps $`f_i`$ are of degree one and induce surjections $`\varphi _i:\pi _1(K_i^4)\pi _1(W^4)`$. Here and henceforth we always assume that the index $`i`$ is large enough so that $`f_i`$ is defined.
###### Lemma 7.2.
The homotopy class of a curve along which surgery is performed is in the kernel of $`\varphi _i:\pi _1(K_i^4)\pi _1(W^4)`$.
###### Proof.
If a surgery is performed along a curve $`\gamma `$, this means that a 2-handle is attached along the curve in the $`4`$-manifold $`W^4`$. Hence $`\gamma `$ bounds a disk in $`W^4\times [0,\mathrm{})`$, which projects to a disk bounded by $`f_i(\gamma )`$ in $`W^4`$. ∎
###### Remark 7.3.
The maps $`\varphi _i`$ and $`\varphi _{i+1}`$ are related in a natural way. To define the map $`\varphi _{i+1}`$, take a generic curve $`\gamma `$ representing any given element of $`\pi _1(K_{i+1}^4)`$. If $`K_{i+1}^4`$ is obtained from $`K_i^4`$ by splitting along a sphere, then $`\gamma `$ is a curve in $`K_i^4`$, and so we can simply take its image. On the other hand, if a surgery was performed, then we may assume that $`\gamma `$ lies off the solid torus that has been attached, and hence lies in $`K_i^4`$, so we can take its image as before. This map is well-defined by lemma 7.2.
###### Definition 7.2.
A curve $`\gamma ^{}K_i^4`$ is a descendant of the surgery curve $`\gamma K_i^4`$ if it is homotopic to it in $`K_i^4`$ (though not in general homotopic to $`\gamma `$ after the surgery). A curve $`\gamma K_i^4`$ is said to persist till $`K_{i+n}^4`$ if some descendant of $`\gamma `$ persists, i.e., we can homotope $`\gamma `$ in $`K_i^4`$ so that it is disjoint from all the future $`2`$-spheres on which 3-handles are attached while passing from $`K_i^4=M_i^3`$ to $`K_{i+n}^4=M_{i+n}^3`$.
###### Definition 7.3.
A curve $`\gamma K_i^4`$ is said to die by $`K_{i+n}^4`$ if it is homotopically trivial in the $`4`$-manifold obtained by attaching $`2`$-handles to $`K_i^4`$ along the curves in $`K_i^4`$ where surgeries are performed in the process of passing to $`K_{i+n}^4`$, or equivalently, $`\gamma `$ is trivial in the group obtained by adding relations to $`\pi _1(M_i^3)`$ corresponding to curves along which the surgery is performed.
We prove now a key property of the sequence $`K_i^4`$.
###### Lemma 7.4.
For each $`i`$, there is a uniform $`n=n(i)`$ such that any curve $`\gamma K_i^4`$, $`\gamma \mathrm{ker}\varphi _i`$ that persists till $`K_{i+n}^4`$ dies by $`K_{i+n}^4`$.
###### Proof.
We can find $`x[0,\mathrm{})`$ so that $`W^4\times \{x\}`$ is entirely after $`K_i^4`$, and $`n_1`$ so that $`W^4\times \{x\}K_{i+n_1^4}`$, because the handlebody decomposition is proper. We then define $`n`$ by repeating this process once, i.e. $`W^4\times \{x_1+\epsilon \}K_{i+n}^4`$, for some $`x_1+\epsilon >x_1>x`$ for which $`W^4\times \{x_1\}`$ is entirely after $`K_{i+n_1}^4`$. Consider $`\gamma \mathrm{ker}\varphi _i`$ which persists till $`K_{i+n_1}^4`$. This means that there is an annulus properly embedded in $`K_{i+n}^4\mathrm{int}(K_i^4)`$, whose boundary curves are $`\gamma `$ and $`\stackrel{~}{\gamma }K_{i+n}^4W^4\times [x_1,x_1+\epsilon )`$. Since $`\gamma \mathrm{ker}\varphi _i`$ it bounds a disc in $`W^4\times [x_1,\mathrm{})`$. This disc together with the above annulus ensure that $`\gamma `$ dies by $`K_{i+n}^4`$, as they bound together a disc entirely in $`K_{i+n}^4\mathrm{int}(K_i^4)`$, and $`3`$-handles do not affect the fundamental group. ∎
### 7.2 The structure theorem
Suppose henceforth that we have a sequence of connected $`3`$-manifolds $`M_i^3W^4\times [0,\mathrm{})`$ and associated maps onto $`f_i:M_i^3W^4`$ that satisfies the properties of $`K_i`$ stated above. Specifically one asks that:
* The maps $`f_i:M_i^3W^4`$ are of degree one, hence inducing surjection $`\varphi _i:\pi _1(M_i^3)\pi _1(W^4)`$.
* $`M_{i+1}`$ is obtained from $`M_i^3`$ either by a 0-frame surgery along a homologically trivial knot in $`M_i^3`$, or else by cutting along a non-separating 2-sphere in $`M_i^3`$.
* The surgery curves in $`M_i^3`$ belong to $`\mathrm{ker}\varphi _i`$.
* The maps $`\varphi _i`$ and $`\varphi _{i+1}`$ are related as in Remark 7.3.
* For any $`i`$ there exists some $`n=n(i)`$ such that any curve in $`M_i^3`$ which persists till $`M_{i+n}^3`$ dies by $`M_{i+n}^3`$.
We show in this section that, after possibly changing the order of attaching handles, any handle decomposition without $`1`$-handles is of a particular form.
We first describe a procedure for attempting to construct a handle decomposition for $`\mathrm{int}(W^4)`$ starting with a partial handle decomposition, with boundary $`M_i^3`$. In general, $`M_i^3`$ has non-trivial homology. It follows readily from the proof of lemma 7.1 that $`H_1(M_i^3)`$ is a torsion free abelian group. The only way we can remove homology is by splitting along spheres. To this end, we take a collection of surfaces representing the homology, perform surgeries along curves in these surfaces so that they compress down to spheres, and then split along these spheres. By doing the surgeries, we have created new homology, and hence have to take new surfaces representing this homology and continue this procedure. In addition to this, we may need to perform other surgeries to get rid of the homologically trivial portion of the kernel of $`\varphi _i:\pi _1(M_i^3)\pi _1(W^4)`$.
The above construction may meet obstructions, since the surgeries have to be performed about curves that are homologically trivial as well as lie in the kernel of $`\varphi _i`$, hence it may not be always possible to perform enough of them to compress the surfaces to spheres. The construction terminates at some finite stage if at that stage all the homology is represented by spheres and no surgery off these surfaces is necessary.
###### Theorem 7.5.
After possibly changing the order of attaching handles, any handle decomposition without $`1`$-handles may be described as follows. We have a collection of surfaces $`F_j^2(i)`$, with disjoint simple closed curves $`l_{j,k}F_j^2(i)`$ and a generic immersion $`\psi _i:_jF_j^2(i)M_i^3`$ such that:
* The surfaces represent the homology of $`M_i^3`$, i.e. $`\psi _i`$ induces a surjection $`\psi _i:H_2(_jF_j^2(i))H_2(M_i^3)`$.
* The immersion $`\psi _i`$ has only ordinary double points and the restriction to each individual surface $`F_j(i)`$ is an embedding. The double curves of $`\psi _i`$ are among the curves $`l_{j,k}`$. Their images $`\psi _i(l_{j,k})`$ are called seams.
* When compressed along the seams (i.e. by adding 2-handles along them) the surfaces $`\psi _i(F_j^2(i))`$ become unions of spheres.
* The seams are homologically trivial curves in $`M_i^3`$ and lie in the kernel of $`\varphi _i`$.
* The pull backs (see the definition below) of the surfaces $`\psi _m(F_k^2(m))M_m^3`$ for $`m>i`$, which are surfaces with boundary in $`M_i^3`$, can only intersect the $`F_j(i)`$’s either transversely at the seams or by having some boundary components along the seams.
We attach $`2`$-handles along all the seams of $`M_i^3`$, and possibly also along some curves that are completely off the surfaces $`F_j^2(i)`$ in $`M_i^3`$ and have no intersection with any future surface $`F_k(m)`$, $`m>i`$. We then attach $`3`$-handles along the 2-spheres obtained by compressing the surfaces $`\psi _i(F_j^2(i))`$. Iterating this procedure gives us the handle decomposition.
We will see that once we construct the surfaces, all of the properties follow automatically.
Let $`F^2M_{i+n}^3`$ be an embedded surface. We let $`M_i^3=K_i^4`$, where $`K_i^4`$ is the bounded component in $`\mathrm{int}(W^4)`$. Then $`K_{i+n}^4\mathrm{int}(K_i^4)=M_i^3\times [0,\epsilon ]h_j^2h_k^3`$, where $`h_j^m`$ are the attached $`m`$-handles.
###### Lemma 7.6.
There exists an isotopy of $`K_{i+n}^4\mathrm{int}(K_i^4)`$ such that $`F^2M_i^3\times \{\epsilon \}h_j^2`$ and $`F^2h_j^2=_k\delta _{j,k}^2`$, where $`\delta _{j,k}^2`$ are disjoint 2-disks properly embedded in the pair $`(h_j^2,_ah_j^2)`$ (here $`_ah_j^2`$ denotes the attachment zone of the handle, which is a solid torus), which are parallel to the core of the handle. Moreover $`\delta _{j,k}^2(_ah_j^2))`$ are concentric circles on the torus, parallel to the 0-framing of the attaching circle.
###### Proof.
It follows from a transversality argument that the image of $`F`$ intersects only the 2-handles, along 2-disks. Further it is sufficient to see that the circles $`\delta _{j,k}^2`$ are homotopic to the 0-framing since in $`K_{i+n}^4\mathrm{int}(K_i^4)`$ homotopy implies isotopy for circles. If one circle is null-homotopic then it can be removed by means of an ambient isotopy. If a circle turns $`p`$-times around the longitude, then it cannot bound a disk in $`h_j^2`$ unless $`p=1`$. ∎
###### Definition 7.4.
Consider a parallel copy in $`M_i^3=M_i^3\times \{0\}`$ of the surface with boundary $`F^2=F^2_{j,k}\delta _{j,k}^2M_i^3\times \{\epsilon \}_j_ah_j^2`$, and use standardly embedded annuli in the torus $`_ah_j^2`$, which join the parallel circles to the central knot in order to get a surface with boundary on the surgery loci. We calls this a pull back of the surface $`F^2M_{i+n}^3`$.
###### Lemma 7.7.
Let $`\alpha :\pi _1(M_i^3)H_1(M_i^3)`$ be the Hurewicz map. Then $`\varphi _i(\mathrm{ker}(\alpha ))=\pi _1(W^4)`$, i.e. the pair $`(\alpha ,\varphi _i)`$ is strongly compressible.
###### Proof.
Consider the diagram $`\begin{array}{ccc}\mathrm{\Gamma }& \stackrel{\varphi }{}& G\\ \pi & & \\ \mathrm{\Gamma }_{ab}& \stackrel{\varphi _{ab}}{}& G_{ab}\end{array}`$ where $`\varphi `$ is surjective, and the subscript $`ab`$ means abelianisation. Then it is automatically that $`\pi (\mathrm{ker}\varphi )=\mathrm{ker}\varphi _{ab}`$. Since $`H_1(W^4)=0`$, and the strong compressibility is symmetric, the result follows. ∎
Now, let $`n=n(i)`$ be as in the conclusion of lemma 7.4. We consider a maximal set of disjoint non-parallel essential 2-spheres (which is uniquely defined up to isotopy) and pull back these spheres up to time $`i`$ to get a collection of planar surfaces, whose union is a 2-dimensional polyhedron $`\mathrm{\Sigma }_iM_i^3`$.
###### Lemma 7.8.
If $`\iota `$ denotes the map induced by the inclusion $`\pi _1(M_i^3\mathrm{\Sigma }_i)\pi _1(M_i^3)`$ then the restriction
$$\varphi _i:\iota (\pi _1(M_i^3\mathrm{\Sigma }_i))\mathrm{ker}(\alpha )\pi _1(W^4)$$
is surjective.
###### Proof.
The pull-backs in $`M_i^3`$ of spheres $`S_m^2M_{i+j}^3`$ are planar surfaces with boundary components being the loci of future surgeries. Further, after compressing the spheres $`S_m^2`$ of $`M_{i+j}`$ (hence arriving into $`M_{i+j+k}^3`$) we have a surjection $`\varphi _{i+j+k}`$, thus the map $`\pi _1(M_{i+j}^3S_m^2)\pi _1(W^4)`$ is also surjective. This means that there exist curves in the complement of the planar surfaces in $`M_i^3`$ mapping to every element of $`\pi _1(W^4)`$. Moreover, by the above lemma, we have such curves that are homologically trivial in $`M_{i+j}^3`$, and hence in $`M_i^3`$ as all surgery curves are null-homologous. ∎
###### Lemma 7.9.
$`i_{}:H_1(M_i^3\mathrm{\Sigma }_i)H_1(M_i^3)`$ is the zero map.
###### Proof.
If not then there exists a curve $`\gamma M_i^3\mathrm{\Sigma }_i`$ that represents a non-trivial element of $`H_1(M_i^3)`$. Modifying by a homologically trivial element if necessary, we may assume that $`\gamma ker(\varphi _i)`$. By the previous lemma $`\gamma `$ persists. The group $`\pi _1(K_{i+n}^4\mathrm{int}(K_i^4))`$ is the quotient of $`\pi _1(M_i^3)`$ by the relations generated by the surgery curves, which are homologically trivial. In particular $`H_1(K_{i+n}^4\mathrm{int}(K_i^4))=H_1(M_i^3)`$. Then the class of $`\gamma H_1(K_{i+n}^4\mathrm{int}(K_i^4))`$ is non-zero since its image in $`H_1(M_i^3)`$ is non-zero by hypothesis. This gives the required contradiction. ∎
We are now in a position to prove the structure theorem. The images of the immersion $`\psi _i`$ is obtained from the polyhedron $`\mathrm{\Sigma }_i`$ by stitching together several planar surfaces along the boundary knots. These knots will be the seams of the surfaces. It is clear by construction that we have all the desired properties as soon as we show that there are enough planar surfaces to be stitched together to represent all the homology.
To see this, we consider the reduced homology exact-sequence of the pair $`(M_i^3,M_i^3\mathrm{\Sigma }_i)`$, and use the fact that $`M_i^3\mathrm{\Sigma }`$ is connected, since $`M_{i+n}^3`$ is, as well as lemma 7.9. Thus, we have the exact sequence
$$\mathrm{}H_1(M_i^3\mathrm{\Sigma }_i)H_1(M_i^3)H_1(M_i^3,M_i^3\mathrm{\Sigma }_i)\stackrel{~}{H}_0(M_i^3\mathrm{\Sigma }_i)$$
which gives the exact sequence
$$0H_1(M_i^3)H_1(M_i^3,M_i^3\mathrm{\Sigma }_i)0$$
which together with an application of Alexander duality gives $`H_1(M_i^3)H_1(M_i^3,M_i^3\mathrm{\Sigma }_i)H^2(\mathrm{\Sigma }_i)`$. Further, as the isomorphisms $`H_1(M_i^3)H^2(M_i^3)`$ and $`H_1(M_i^3,M_i^3\mathrm{\Sigma }_i)H^2(\mathrm{\Sigma }_i)`$, given respectively by Poincaré and Alexander duality, are obtained by taking cup products with the fundamental class, the diagram
$$\begin{array}{ccc}H^2(M_i^3)& & H^2(\mathrm{\Sigma }_i)\\ & & \\ H_1(M_i^3)& & H_1(M_i^3,M_i^3\mathrm{\Sigma }_i)\end{array}$$
commutes.
Thus the inclusion of $`\mathrm{\Sigma }_i`$ in $`M_i^3`$ gives an isomorphism $`H^2(M_i^3)H^2(\mathrm{\Sigma }_i)`$. Since $`H_2(M_i^3)`$ and $`H_2(\mathrm{\Sigma }_i)`$ have no torsion, the cap product induces perfect pairings $`H_2(M_i^3)\times H^2(M_i^3)`$ and $`H_2(\mathrm{\Sigma }_i)\times H^2(\mathrm{\Sigma }_i)`$. Therefore, by duality, the map $`H_2(\mathrm{\Sigma }_i)H_2(M_i^3)`$ induced by inclusion is also an isomorphism.
Now take a basis for $`H^2(\mathrm{\Sigma }_i)`$. Each element of this basis can be looked at as an integral linear combination of the planar surfaces (as in cellular homology), with trivial boundary. We obtain a surface corresponding to each such homology class by taking copies of the planar surfaces, with the number and orientation determined by the coefficient. Since the homology classes are cycles, these planar surfaces can be glued together at the boundaries to form closed, oriented, immersed surfaces. Without loss of generality, we can assume these to be connected.
###### Remark 7.10.
By doing surgeries on the seams of $`\mathrm{\Sigma }_iM_i^3`$ some new homology is created (the homology of $`M_{i+1}^3`$) One constructs naturally surfaces representing the homology of $`M_{i+1}^3`$, as follows. One considers generalized Seifert surfaces in $`M_i^3`$ of the loci of the surgeries, which are surfaces which might have boundary components along other seams. Then one caps-off the boundaries by using the cores of the 2-handles which are added and push the closed surfaces into $`M_{i+1}^3`$. Notice that we can consider also some Seifert surface whose boundary components are seams in some $`M_{i+n}^3`$ for $`n>1`$.
### 7.3 On Casson finiteness
Suppose we do have a $`4`$-manifold $`(W^4,W^4)`$ with a g.s.c. handle decomposition of its interior. Since there may be infinitely many handles, we cannot use Casson’s argument. However, we note that we can use Casson’s argument if we can show that
* $`(W^4,W^4)`$ has a (finite) handle decomposition without $`1`$-handles.
* Some $`(Z^4,Z^4)`$ has a handle decomposition without $`1`$-handles, where $`Z^4`$ is compact, contractible with $`\pi _1(Z^4)=\pi _1(W^4)`$.
* Some $`(Z^4,Z^4)`$ has a handle decomposition without $`1`$-handles, where $`Z^4`$ is compact, contractible and there is a surjection $`\pi _1(Z^4)\pi _1(W^4)`$ (by Proposition 6.4).
Thus, we can apply Casson’s argument if we show finiteness, or some weak form of finiteness such as the latter statements above.
We now assume that the handle decomposition is as in the conclusion of Theorem 7.5. We will change our measures of time so that passing from $`M_i^3`$ to $`M_{i+1}^3`$ consists of performing all the surgeries required to compress the surfaces, splitting along the 2-spheres, and also performing the necessary surgeries off the surface.
In $`M_i^3`$, we have a collection of embedded surfaces representing all the homology of $`M_i^3`$. We see that we have Casson finiteness in a special case.
###### Theorem 7.11.
If there exists $`i`$ so that the immersion $`\psi _i:_jF_j^2(i)M_i^3`$ is actually an embedding and $`\varphi _i(\psi _i(F_j^2(i))=\{1\}\pi _1(W^4)`$ then $`\pi _1(W^4)`$ violates the Kervaire conjecture.
###### Proof.
Let $`k`$ be the rank of $`H_1(M_i^3)`$ and $`P_j,1jk`$ be the fundamental groups of the surfaces. Since the surfaces are disjoint, $`\pi _1(M_i^3)`$ is obtained by HNN extensions from the fundamental group $`G`$ of the complement of the surfaces. Thus, if $`\psi _j`$ are the gluing maps, we have
$$\pi _1(M_i^3)=G,t_1,\mathrm{},t_k;t_jxt_j^1=\psi _j(x)xP_j$$
Now, since $`\varphi _i(P_j)=1`$ and $`\varphi _i(G)=\pi _1(W^4)`$, $`\pi _1(M_i^3)`$ surjects onto $`\pi _1(W^4),t_1,\mathrm{},t_n`$, the group obtained by adding $`k`$ generators to $`\pi _1(W^4)`$. But, $`M_i^3`$ is obtained by using $`n`$ 2-handles and $`nk`$ 3-handles. Thus, as in Casson’s theorem, $`\pi _1(M_i^3)`$ is killed by adding $`nk`$ generators and $`n`$ relations. This implies that $`\pi _1(W^4)`$ is killed by adding $`n`$ generators and $`n`$ relations. ∎
###### Theorem 7.12.
There exists always an immersion as in the structure theorem with the additional property that $`\varphi _i(\psi _i(F_j^2(i))=\{1\}\pi _1(W^4)`$ holds true.
###### Proof.
By construction the images of the seams (which are roughly speaking half the generators of the fundamental group) are null-homotopic. If the fundamental groups of the generalized Seifert surfaces from the previous remark map to the trivial group, then after doing surgery on the seams we obtain surfaces $`F_j^2(i+1)`$ representing homology with trivial $`\pi _1`$ images by $`\varphi _{i+1}`$. Thus, it suffices to show that we obtain this condition for a choice of Seifert surfaces for all seams, at some time in the future.
Fix $`i`$ large enough so that $`M_i^3`$ is in the collar $`W^4\times [0,\mathrm{})`$.
###### Lemma 7.13.
There exists some $`n^{}=n^{}(i)`$ such that, whenever a 2-sphere immersed in $`\mathrm{int}(W^4)K_{i+n^{}}^4`$ bound a 3-ball immersed in the collar, it actually bounds a 3-ball immersed in $`\mathrm{int}(W^4)K_i^4`$.
###### Proof.
Choose $`n^{}`$ large enough so that a small collar $`W^4\times [x,y]K_{i+n^{}}^4\mathrm{int}(K_i^4)`$. Then use the horizontal flow to send $`W^4\times [0,y]`$ into $`W^4\times [x,y]`$ by preserving the right side boundary. This yields a ball in the complement of $`K_i^4`$. ∎
We need the following analogue, for $`\pi _2`$ instead of $`\pi _1`$, of Lemma 7.4.
###### Lemma 7.14.
There exists $`k=k(i)`$ such that any immersed 2-sphere in $`K_i^4`$ (respectively in the intersection of $`K_i^4`$ with the collar) which bounds a 3-ball in $`\mathrm{int}(W^4)`$ (respectively in the collar) does so in $`K_{i+k}^4`$ (respectively in the collar).
###### Proof.
The same trick we used in the proof of the previous lemma applies. ∎
Choose now $`n`$ large enough such that $`K_{i+n}^4K_{i+n^{}}^4`$ contains a non-trivial collar $`W^4\times [x,y]`$, and $`n>k(i+n^{})`$ provided by the Lemma 7.4. Consider a surgery curve $`\gamma M_j^3`$ for some $`i+n^{}j<i+n`$.
###### Lemma 7.15.
There exists a generalized Seifert surface for $`\gamma `$ in $`M_i^3`$, so that the other boundary components are surgery curves from $`M_m^3`$, with $`mi+n`$, and whose fundamental group maps to the trivial group under $`\varphi _i`$.
###### Proof.
The curve $`\gamma `$ bounds a disc in the $`2`$-handle attached to it. Further, as it can be pulled back, say along an annulus, to time $`M_{i+n^{}}^3`$, and then dies by $`M_{i+n}^3`$, it bounds another disc consisting of the annulus and the disc by which it dies. These discs together form an immersed $`2`$-sphere. Consider the class $`\nu \pi _2(W^4)`$ of this 2-sphere by using the projection of the collar on $`W^4`$. We can realize the element $`\nu `$ by an immersed 2-sphere in a small collar $`W^4\times [x,y]K_{i+n}^4K_i^4`$. Therefore by modifying the initial 2-sphere by this sphere (which is far from $`\gamma `$) in the small collar one finds an immersed 2-sphere whose image in $`\pi _2(W^4)`$ is trivial. Since $`i`$ was large enough $`K_{i+n}^4K_i^4`$ is a subset of a larger collar $`W^4\times [0,z]`$. Then the 2-sphere we constructed bounds a 3-ball in $`W^4\times [0,z]`$, and so by Lemma 7.13 it also does so in $`X^4=K_{i+n}^4K_i^4`$. Let $`\mu :S^2X^4`$ denote this immersion realizing a trivial element of $`\pi _2(X)`$. Then $`\mu `$ lifts to a map $`\stackrel{~}{\mu }:S^2\stackrel{~}{X}^4`$, where $`\stackrel{~}{X}^4`$ is the universal covering space of $`X^4`$. Since $`\mu `$ is null-homotopic the homology class of $`[\stackrel{~}{\mu }]=0H_2(\stackrel{~}{X}^4)`$ is trivial, when interpreting $`\stackrel{~}{\mu }`$ as a 2-cycle in $`\stackrel{~}{X}^4`$.
The homology of $`\stackrel{~}{X}^4`$ is computed from the $`\pi _1(X^4)`$-equivariant complex associated to the handle decomposition, whose generators in degree $`d`$ are the $`d`$-handles attached to $`K_i^4`$ in order to get $`K_{i+n}^4`$. Therefore one has then the following relation in this differential complex:
$$[\stackrel{~}{\mu }]=\underset{j}{}c_j[h_j^3],c_j.$$
The action of the algebraic boundary operator $``$ on the element $`[h_j^3]`$ can be described geometrically as the class of the 2-cycle which represents the attachment 2-sphere $`^+h_j^3`$ of the 3-handle $`h_j^3`$. Consequently the previous formula can be rewritten as
$$\stackrel{~}{\mu }=\underset{j}{}c_j^+h_j^3+\underset{k}{}d_kL_k,c_j,d_k$$
where $`L_k`$ are closed surfaces (actually these are closed 2-cycles, but they can be represented by surfaces by the well-known results of R. Thom) with the property that
$$[L_k][\delta _m^2]=0,k,m$$
($`\delta _m^2`$ denotes the core of the 2-handle $`h_m^2`$). Let us compute explicitly the boundary operator on the 3-handles, in terms of the surfaces we have in the 2-complex $`\mathrm{\Sigma }_i`$. Set
$$[h_j^3]=\underset{k}{}m_{jk}[h_k^2].$$
Then the coefficient $`m_{jk}`$ is the number of times the boundary $`^+h_j^3`$ runs over the core of $`h_k^2`$. But the 2-sphere $`^+h_j^3`$, when pulled back in $`M_i^3`$, is a planar surface in $`M_i^3`$ whose boundary circles (i.e. at seams) are capped-off by the core disks $`\delta _k^2`$ of the 2-handles $`h_k^2`$. Therefore the number $`m_{jk}`$ is the number of times the seam $`^+h_k^2`$ appears in the planar surface which is a pull-back of $`^+h_j^3`$. In particular the coefficient of a $`2`$-handle vanishes in a $`3`$-cycle only if the boundaries of the planar surface glue together to close up at the corresponding surgery locus.
Thus the pull-backs of the surfaces $`_jc_j^+h_j^3`$ give a surface $`F^2`$ in $`M_i^3`$ with boundary the curve $`\gamma `$ with which we started, plus some other curves along which surgery is performed by time $`i+n`$. As this is in fact a closed cycle in the universal cover, the surface $`F^2`$ lifts to a surface in $`\stackrel{~}{X}^4`$, with a single boundary component, corresponding to a curve which is not surgered by the time $`i+n`$. Therefore the map $`\pi _1(F^2)\pi _1(X^4)\pi _1(W^4)`$ factors through $`\pi _1(\stackrel{~}{X}^4)=1`$, hence the image of $`\pi _1(F^2)`$ in $`\pi _1(W^4)`$ is trivial. ∎
Thus, after surgering along the curves up to the $`M_{i+n}^3`$ we do have the required Seifert surfaces to compress to get embedded surfaces with trivial $`\pi _1(W^4)`$ image. ∎
###### Proposition 7.16.
If $`\gamma _{i,k}M_i^3`$ are the surgery curves (i.e. the seams) then $`\gamma _{i,k}LCS_{\mathrm{}}(\pi _1(M_i^3))`$, where $`LCS_s(G)`$ is the lower central series of the group $`G`$, $`LCS_1(G)=G,LCS_{s+1}(G)=[G,LCS_s(G)]`$, and $`LCS_{\mathrm{}}(G)=_{s=1}^{\mathrm{}}LCS_s(G)`$.
###### Proof.
We will express each surgery locus $`\gamma `$ as a product of commutators of the form $`[\alpha _i,\beta _i]`$, with each $`\alpha _i`$ being conjugate to a surgery locus (possibly $`\gamma `$ itself). It then follows readily that $`\gamma LCS_{\mathrm{}}(\pi _1(M_i^3))`$, as now if each $`\gamma _{i,k}LCS_s(\pi _1(M_i^3))`$, then each $`\gamma _{i,k}LCS_{s+1}(\pi _1(M_i^3))`$.
Suppose now $`\gamma =\gamma _{i,k}`$ is a surgery locus. Then the $`0`$-frame surgery along $`\gamma `$ creates homology in $`M_{i+1}^3`$ which by our structure theorem is represented by a surface $`S^2=\psi _{i+1}(F_j^2(i+1))`$. The pullback of $`S^2`$ to $`i`$ gives a surface with boundary along seams, and being compressed to a sphere by the seams, so that the algebraic multiplicity of $`\gamma `$ is $`1`$ while that of all other seams is $`0`$. In terms of the fundamental group, this translates to the relation that was claimed. ∎
Since, we have immersed surfaces of the required form, the obstruction we encounter is in making these surfaces disjoint at some finite stage. Note that for a finite decomposition, we do indeed have disjoint surfaces representing the homology after finitely many surgeries, since we in fact have a family of such spheres.
###### Proposition 7.17.
There exists a $`2`$-complex $`\mathrm{\Sigma }=_{i=1}^{\mathrm{}}\mathrm{\Sigma }_i`$ with intersections along double curves, coming from a handle-decomposition as above, where all the seams are trivial in homology, but which does not carry disjoint, embedded surfaces representing all of the homology.
###### Proof.
For the first stage, take two surfaces of genus $`2`$, and let them intersect transversely along two curves (which we call seams) that are disjoint and homologically independent in each surface. Next, take as Seifert surfaces for these curves once punctured surfaces of genus $`2`$ intersecting in a similar manner, and glue their boundary to the above-mentioned curves of intersection. Repeat this process to obtain the complex.
At the first stage, we cannot have embedded, disjoint surfaces representing the homology as the cup product of the surfaces is non-trivial. As the surfaces are compact, we must terminate at some finite stage. We will prove that if we can have disjoint surfaces at the stage $`k+1`$, then we do at stage $`k`$. This will suffice to give the contradiction.
Now, we know the complex cannot be embedded in the first stage. Suppose we did have disjoint embedded surface $`F_1`$ and $`F_2`$ at stage $`k+1`$. Since these form a basis for the homology, they contain curves on them that are the seams at the first stage with algebraically non-zero multiplicity, i.e., the collection of curves representing the seam is not homologically trivial in the intersection of the first stage with the surface. Further, some copy of the first seam must bound a subsurface $`F_i^{}`$ in each of the surfaces, for otherwise the surface contains a curve dual to the seam. For, the cup product of such a dual curve with the homology class of the other surface is non-trivial, hence it must intersect the other surface, contradicting the hypothesis that the surfaces are disjoint. Similarly, at the other seam we get surfaces $`F_i^{\prime \prime }`$.
By deleting the first stage surfaces and capping off the first stage seams by attaching discs, we get a complex exactly as before with the $`(j+1)`$th stage having become the $`j`$th stage. Further, the $`F_1^{}`$ and $`F_2^{\prime \prime }`$ now give disjoint, embedded surfaces representing the homology that are supported by stages up to $`k`$. This suffices as above to complete the induction argument.
It is easy to construct a handle-decomposition corresponding to this complex. Figure 1 shows a construction of tori with one curve of intersection. Here we have used the notation of Kirby calculus, with the thickened curves being an unlink along each component of which $`0`$-frame surgery has been performed. It is easy to see that the same construction can give surfaces of genus $`2`$ intersecting in $`2`$ curves. On attaching the first two $`2`$-handles, the boundary is $`(S^2\times S^1)\mathrm{\#}(S^2\times S^1)`$. Since the curves of intersection are unknots, after surgery they bound spheres. Further, it is easy to see by cutting along these that the boundary is $`(S^2\times S^1)\mathrm{\#}(S^2\times S^1)`$ after attaching the $`2`$-handles and $`3`$-handles as well. Repeating this process, we obtain our embedding.
Thus we have an infinite handle-decomposition satisfying our hypothesis for which this $`2`$-complex is $`\mathrm{\Sigma }`$. ∎
### 7.4 A wild example
We will construct an example of an open, contractible $`4`$-manifold that is not tame, and that has a handle-decomposition without $`1`$-handles.
###### Theorem 7.18.
There is a proper handle-decomposition of an open, contractible $`4`$-manifold $`W^4`$ such that $`W^4`$ is not the interior of a compact $`4`$-manifold. In particular $`W^4`$ does not have a finite handle-decomposition.
###### Proof.
We will take a variant of the example in the last section. Namely, we construct an explicit handle-decomposition according to a canonical form.
Start with a $`0`$-handle and attach to its boundary three $`2`$-handles along an unlink. the resulting manifold has boundary $`(S^2\times S^1)\mathrm{\#}(S^2\times S^1)\mathrm{\#}(S^2\times S^1)`$ obtained by $`0`$-frame-surgery about each component of an unlink with $`3`$ components. We now take as Seifert surfaces for these components surfaces of genus $`2`$, so that each pair intersects in a single curve, so that the curves of intersection form an unlink and are unlinked with the original curves.
Now, attach $`2`$-handles along the curves of intersections, and then $`3`$-handles along the Seifert surfaces compressed to spheres by adding discs in the $`2`$-handles just attached. It is easy to see that the resulting manifold once more has boundary $`(S^2\times S^1)\mathrm{\#}(S^2\times S^1)\mathrm{\#}(S^2\times S^1)`$. Thus, we may iterate this process. Further, the generators of the fundamental group at any stage are the commutators of the generators at the previous stage.
Suppose $`W^4`$ is in fact tame. Then, we may use the results of the previous sections. Now, by construction no curve dies as only trivial relations have been added. Thus every element in kernel($`\varphi _i`$) must fail to persist by some uniform time. In particular, the image of the group after that time in the present (curves that persist beyond that time) must inject under $`\varphi _i`$. But we know that it also surjects. Thus, we must have an isomorphism.
Thus, there is a unique element mapping onto each element of $`\pi _1(W^4)`$. Hence this element must persist till infinity as we have a surjection at all times. On the other hand, since the limit of the lower central series of the free group is trivial, no non-trivial element persists. This gives a contradiction unless $`\pi _1(W^4)`$ is trivial.
But there are non-trivial elements that do persist beyond any give time. As no element dies, we again get a contradiction. ∎
### 7.5 Further obstructions from Gauge theory
To further explore some of the subtleties that one might encounter in trying to construct a handle decomposition without $`1`$-handles for a contractible manifold, given one for its interior, we consider a more general situation. We will consider sequences of $`3`$-manifolds $`M_i^3`$ that begin with $`S^3`$. As before, we require that each manifold comes from the previous one by $`0`$-frame surgery about a homologically trivial curve or by splitting along a non-separating $`S^2`$ and capping off. Also, we require degree-one maps $`f_i`$ to a common manifold $`N^3`$, related as before. We will say that the sequence limits to $`N^3`$ if any curve that persists dies as in lemma 7.4.
In this situation, our main question generalizes to a relative version, namely, given any such sequence $`\{M_i^3\}`$, with $`M_k^3`$ an element in the sequence, is there a finite sequence that agrees up to $`M_k^3`$ with the old sequence and whose final term is $`N^3`$?
We will show that there is an obstruction to completing certain sequences to finite sequences when $`N^3=S^3`$. We do not know whether there are infinite sequences limiting to $`N^3`$ in this case.
Let $`𝔓^3`$ denote the Poincaré homology sphere. Observe that we cannot pass from this to $`S^3`$ by $`0`$-frame surgery about homologically trivial curves and capping-off non-separating spheres. For, if we could, $`𝔓^3`$ would bound a manifold with $`H^2=_k\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right]`$, which is impossible as $`𝔓^3`$ has Rochlin invariant $`1`$. On the other hand, for the same reason, $`𝔓^3`$ cannot be part of any sequence of the above form.
Using Donaldson’s theorem , we have a similar result for the connected sum $`𝔓^3\mathrm{\#}𝔓^3`$ of $`𝔓^3`$ with itself. The main part of the proof of this lemma was communicated to us by R. Gompf.
###### Lemma 7.19.
One cannot pass from $`𝔓^3\mathrm{\#}𝔓^3`$ to $`S^3`$ by $`0`$-frame surgery along homologically trivial curves and capping off non-separating $`S^2`$’s.
###### Proof.
If we did have such a sequence of surgeries, then $`𝔓^3\mathrm{\#}𝔓^3`$ bounds a $`4`$-manifold $`M_1^4`$ with $`H^2=_k\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right]`$ , with a half-basis formed by embedded spheres. Now glue this to a manifold with form $`E_8E_8`$ which is bounded by $`𝔓^3\mathrm{\#}𝔓^3`$ to get $`M^4`$.
We can surger out the disjoint family of $`S^2`$’s from $`M^4`$ to get a $`4`$-manifold with form $`E_8E_8`$ and trivial $`H_1`$. This contradicts Donaldson’s theorem. ∎
We still do not have a sequence as claimed. For, Cassson’s argument shows that $`𝔓^3\mathrm{\#}𝔓^3`$ cannot be part of a sequence. To obtain such a sequence, we will construct a manifold $`N^3`$ that can be obtained by $`0`$-frame surgery on algebraically unlinked $`2`$-handles from each of $`S^3`$ and $`𝔓^3\mathrm{\#}𝔓^3`$. Thus, $`N^3`$ is part of a sequence. On the other hand, if we had a sequence starting at $`N^3`$ that terminated at $`S^3`$, then we would have one starting at $`𝔓^3\mathrm{\#}𝔓^3`$, which contradicts the above lemma.
To construct $`N^3`$, take a contractible $`4`$-manifold $`K^4`$ that bounds $`𝔓^3\mathrm{\#}𝔓^3`$. By Freedman’s theorem (), this exists, and can moreover be smoothed after taking connected sums with sufficiently many copies of $`S^2\times S^2`$. Take a handle decomposition of $`K^4`$. This may include $`1`$-handles, but these must be boundaries of $`2`$-handles. Hence, by handle-slides, we can ensure that each $`1`$-handle is, at the homological level, a boundary of a $`2`$-handle and is not part of the boundary of any other $`2`$-handle. Replacing the $`1`$-handle by a $`2`$-handle does not change the boundary, and changes $`H^2(K^4)`$ to $`H^2(K^4)\left(_k\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right]\right)`$. We do this dually with $`3`$-handles too. Sliding $`2`$-handles over the new ones, we can ensure that the attaching maps of the $`2`$-handles are having the same algebraic linking (and framing) structure as a disjoint union of Hopf links.
Now take $`N^3`$ obtained from $`S^3`$ by attaching half the links, so that these are pairwise algebraically unlinked. The manifold $`N^3`$ has the required properties.
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# 1 Introduction
## 1 Introduction
The AdS/CFT correspondence gives a remarkable relation between string theory on a spacetime and a certain conformal field theory (CFT) on the boundary of this spacetime . In particular the near horizon geometry of D3 branes gives the space $`AdS_5\times S^5`$, and string theory on this space is conjectured to be dual to N=4 supersymmetric Yang-Mills on the boundary of the $`AdS_5`$. When the string theory is weakly coupled, tree level supergravity is a valid low energy approximation. The dual CFT at this point is a strongly coupled Yang-Mills theory, which cannot therefore be studied perturbatively. On the other hand weakly coupled Yang-Mills theory is dual to string theory in a domain of parameters where the latter cannot be approximated by supergravity on a gently curved spacetime. In spite of this fact, it turns out that certain quantities computed in the supergravity limit of string theory agree with their corresponding dual quantities in the Yang-Mills theory, where the latter computation is done at weak coupling. One believes that such an agreement is due to the supersymmetry which is present in the theory; this supersymmetry would for example protect the dimensions of chiral operators from changing when the coupling is varied. Interestingly, the values of 3-point correlation functions of chiral operators are also found to agree, when we compare the tree level supergravity calculation on AdS space with the computation in the free Yang-Mills theory ; the latter is just the result obtained by Wick contractions among the fields in the chiral operators. It is not clear if the 3-point function of chiral operators is protected against change of coupling at all values of $`N`$; the above result just tells us that the large $`N`$ results agree between the weak and strong coupling limits.<sup>1</sup><sup>1</sup>1See however for an analysis of the finite $`N`$ case.
One is thus led to ask: are the 3-point functions of chiral operators protected in the other cases of the AdS/CFT correspondence? In particular we will be interested in the case of the D1-D5 system , which gives a near-horizon geometry $`AdS_3\times S^3\times M`$, where $`M`$ is a torus $`T^4`$ or a $`K3`$ space. This system is of great interest for the issues related to black holes, since it yields, upon addition of momentum excitations, a supersymmetric configuration which has a classical (i.e. not Planck size) horizon. In particular, the Bekenstein entropy computed from the classical horizon area agrees with the count of microstates for the extremal and near extremal black holes . Further, the low energy Hawking radiation from the hole can be understood in terms of a unitary microscopic process, not only qualitatively but also quantitatively, since one finds an agreement of spin dependence and radiation rates between the semiclassically computed radiation and the microscopic calculation . While it is possible to use simple models for the low energy dynamics of the D1-D5 system when one is computing the coupling to massless modes of the supergravity theory, it is believed that the exact description of this CFT must be in terms of a sigma model with target space being a deformation of the orbifold $`M^N/S_N`$, which is the symmetric orbifold of $`N`$ copies of $`M`$. (Here $`N=n_1n_5`$, with $`n_1`$ being the number of D1 branes and $`n_5`$ being the number of D5 branes, and we must take the low energy limit of the sigma model to obtain the desired CFT.) In particular we may consider the ‘orbifold point’ where the target space is exactly the orbifold $`M^N/S_N`$ with no deformation. It was argued in that this CFT does correspond to a certain point in the moduli space of string theories on $`AdS_3\times S^3\times M`$, but at this point the string theory is in a strongly coupled domain where it cannot be approximated by tree level supergravity on a smooth background. The orbifold point is the closest we can get to a ‘free’ theory on the CFT side, and thus this point is the analogue of free N=4 supersymmetric Yang-Mills in the D3 brane example. Thus one would like to compare the three point functions of chiral operators in the supergravity limit with the 3-point functions at the orbifold point, to see if we have an analogue of the surprising agreement that was found in the case of the $`AdS_5`$ \- 4-d Yang-Mills duality.
The orbifold group in our case is $`S_N`$, the permutation group of $`N`$ elements. This group is nonabelian, in contrast to the cyclic group $`Z_N`$ which has been studied more extensively in the past for computation of correlation functions in orbifold theories . Though there are some results in the literature for general orbifolds , the study of nonabelian orbifolds is much less developed than for abelian orbifolds.
It turns out however that the case of the $`S_N`$ orbifolds has its own set of simplifications which make it possible to develop a technique for computation of correlation functions for these theories. The essential quantities that we wish to compute are the correlation functions of ‘twist operators’, in the CFT that arises from the infra-red limit of the 2-d sigma model with target space $`M^N/S_N`$. If we circle the insertion of a twist operator of the permutation group $`S_N`$, different copies of the target space $`M`$ permute into each other. We pass to the covering space of the 2-d base space, such that on this covering space the fields of the CFT are single-valued. For the special case where the orbifold group is $`S_N`$, the path integral on the base space with twist operators inserted becomes a path integral over the covering space for the CFT with only one copy of $`M`$, with no operator insertions. Thus the correlation functions of twist operators can be rewritten as partition functions on Riemann surfaces of different genus for the CFT arising from one copy of $`M`$.
In the simplest case, which is also the case giving the leading contribution at large $`N`$, the genus of the covering surface is zero, and we just get the partition function on the sphere. But the metric on this sphere is determined by the orders and locations of the twist operators. We can write this metric as $`g=e^\varphi \widehat{g}`$, for some fiducial metric $`\widehat{g}`$, provided we take into account the conformal anomaly given by the Liouville action for $`\varphi `$. It turns out that $`\varphi `$ is harmonic outside a finite number of isolated points, so the Liouville action can be computed by observing the local behavior of the covering surface at these points. In this manner we can compute any correlation function of twist operators. For given operators we get contributions from only a finite range of genera for the covering surface.
We compute the 2-point function of twist operators, from which we recover their well known scaling dimensions. We then compute the contribution to 3-point functions which comes from covering surfaces of genus zero. This gives the complete result for the fusion coefficients of twist operators for a subset of cases (the ‘single overlap’ cases) and the leading result at large $`N`$ for all cases. We then compute the 4-point function for twist operators of order $`2`$; this correlator has contributions from genus 0 and genus 1, and we compute both contributions.
We find a certain ‘universality’ in the correlation functions, since the Liouville action depends only on the central charge $`c`$ of the CFT with target space $`M`$. The contribution from higher genus covering surfaces will involve the partition functions at those genera, while the leading order result coming from the covering surface of genus zero will depend only on $`c`$. These observations generalize the fact that the dimensions of the twist operators depend only on $`c`$.
We address only bosonic operators in this paper, though we expect the extension to the supersymmetric case to be relatively straightforward. Thus we also do not address the comparison to supergravity; this comparison should be carried out only for the supersymmetric case. But we do observe some features of our correlation functions that accord with some of the patterns that emerge in the supergravity computation. In particular we find a similarity between the condition for the 3-point correlator to have a contribution from genus zero covering surfaces, and the condition for primaries to fuse together in the SU(2) Wess-Zumino-Witten model.
There are several earlier works that relate to the problem we are studying, in particular . We will mention these in more detail in the context where they appear.
The plan of this paper is the following. Section 2 describes our method of computing correlation functions. In section 3 we compute the 2-point functions, and thus recovers the scaling dimensions of the twist operators. In section 4 we construct the map to the covering space for 3-point functions, for the case where the covering space is a sphere. In section 5 we find the Liouville action associated to this map. In section 6 we use the above results to obtain the contribution to 3-point functions from covering surfaces of genus zero. Section 7 discusses an example of a 4-point function. Section 8 is a discussion.
## 2 Computing correlation functions through the Liouville action
### 2.1 The twist operators
Let us consider for simplicity the case $`M=R`$, i.e. the noncompact real line. Then the sigma model target space $`M^N/S_N`$ can be described through a collection of $`N`$ free bosons $`X^1,X^2,\mathrm{}X^N`$, living on a plane parameterized by the complex coordinate $`z`$. We will see later that we can extend the analysis directly to other CFTs. Let the $`z`$ plane have the flat metric
$$ds^2=dzd\overline{z}.$$
(2.1)
The CFT is defined through a path integral over the values of the $`X^i`$, with action
$$S=d^2z2_zX^i_{\overline{z}}X^i.$$
(2.2)
We now make the definition of this partition function more precise. We cut off the $`z`$ plane at a large radius
$$|z|=\frac{1}{\delta },\delta \mathrm{small}.$$
(2.3)
We want boundary conditions at this large circle to represent the fact that the identity operator has been inserted at infinity. We will explain the norm of this boundary state later on.
Thus we imagine that our CFT is defined on a ‘sample’ of size $`1/\delta `$ – correlation functions are to be computed by putting the operators at $`|z_i|<<1/\delta `$ and we take $`\delta `$ to zero at the end of the calculation. (An exception will be the insertion of an operator at infinity, which we will have to define separately.) We write the path integral for a single boson on the $`z`$ plane as
$$Z_\delta =DXe^S.$$
(2.4)
The path integral for $`N`$ bosons, with no twist operator insertions, is $`(Z_\delta )^N`$.
The twist operator $`\sigma _{12}(z_1)`$ can be described through the following. Cut a circular hole of radius $`ϵ`$ in the $`z`$ plane about the point $`z_1`$. While the path integral over $`X^i,i=3,\mathrm{}N`$ is left unchanged, we modify the boundary conditions on $`X^1,X^2`$ such that as we go around the hole at $`z_1`$ we get
$$X^1X^2,X^2X^1.$$
(2.5)
We call this operator $`\sigma _{12}^ϵ(z_1)`$, where the $`ϵ`$ in the superscript reminds us of the regulation used to define the twist. Note that we still have to define more precisely the state that is inserted at the edge of the hole of radius $`ϵ`$ \- we will do this below. If we want to maintain the boundary condition at the circle $`|z|=1/\delta `$ (and introduce no twist there) we must insert at some point $`z_2`$ another such twist operator $`\sigma _{12}`$; the size of the circular hole around $`z_2`$ is also $`ϵ`$. Let us compute the path integral with these boundary conditions, and call it $`Z_{ϵ,\delta }[\sigma _2(z_1),\sigma _2(z_2)]`$. Then we define the correlation function
$$\sigma ^ϵ(z_1)\sigma ^ϵ(z_2)_\delta \frac{Z_{ϵ,\delta }[\sigma _2(z_1),\sigma _2(z_2)]}{(Z_\delta )^N}.$$
(2.6)
We will later define rescaled twist operators, and the cutoffs $`ϵ,\delta `$ will disappear from all final answers.
### 2.2 Path integral on the covering space
The functions $`X^1,X^2`$ above were not single valued in the $`z`$ plane due to the insertion of the twist operators. We wish to pass to a covering space where these functions would become one single valued function. Since the fields $`X^3,\mathrm{}X^N`$ are not involved in the twist, they cancel out in the RHS of (2.6). We thus consider only $`X^1,X^2`$ in the following.
Consider a configuration of the fields $`X^1,X^2`$ which contributes to the path integral. Consider a simply connected patch of the $`z`$ plane, which excludes the holes around $`z_1,z_2`$. Over this patch there are two functions defined: one for each field, though due to the twists there is no global way to label the functions uniquely as being $`X^1`$ or $`X^2`$. Take any one of the functions over this patch - call it $`X(z)`$. To construct the covering surface $`\mathrm{\Sigma }`$, let this open set in $`z`$ be a patch on $`\mathrm{\Sigma }`$, with the complex structure given by $`z`$, and the metric also equal to the metric (2.1) of the $`z`$ plane. There is one field $`X`$ that will be defined over $`\mathrm{\Sigma }`$, and over this patch let it be the above mentioned function $`X(z)`$. Now consider another such simply connected open patch, partly overlapping with the first, and use it to define another patch on $`\mathrm{\Sigma }`$. Clearly, as we go around the point $`z_1`$, following these overlapping patches, the surface $`\mathrm{\Sigma }`$ will look locally like the Riemann surface of a function
$$t=(zz_1)^{1/2}.$$
(2.7)
Further, the functions $`X^1,X^2`$ will be both encoded in the function $`X`$ which will be single valued on $`\mathrm{\Sigma }`$, and the action of the configuration of $`X^1,X^2`$ will be reproduced if in each patch we use for $`X`$ the action
$$S=_{\mathrm{patch}}d^2z2_zX_{\overline{z}}X.$$
(2.8)
The coordinate $`z`$ cannot be a globally well defined coordinate for $`\mathrm{\Sigma }`$, since a generic value of $`z`$ corresponds to two points on $`\mathrm{\Sigma }`$. We will call our choice of coordinate on $`\mathrm{\Sigma }`$ as $`t`$, which will be locally holomorphically related to $`z`$.
In (2.8) we must evaluate the path integral on the patch using the metric induced from the $`z`$ plane on the patch; the path integral depends in general on the metric and the physical problem is defined through the metric chosen on the $`z`$ plane. Thus in terms of the coordinate $`t`$ used to describe $`\mathrm{\Sigma }`$ we will have
$$ds^2=dzd\overline{z}=|\frac{dz}{dt}|^2dtd\overline{t}.$$
(2.9)
For the example above we can take
$$z=a\frac{t^2}{2t1}.$$
(2.10)
Near the location $`z=0,t=0`$, we evidently have $`zt^2`$, so that $`t`$ parameterizes the covering surface near $`z=0`$. But the map between $`z`$ and $`t`$ is singular also at $`z=a,t=1`$, since near $`t=1`$
$$\frac{dz}{dt}2a(t1),zaa(t1)^2.$$
(2.11)
Thus the twist operators are located at $`z_1=0`$ and $`z_2=a`$.
The covering space $`\mathrm{\Sigma }`$ (parameterized by $`t`$) is a double cover of the original sphere parameterized by $`z`$. We had cut the $`z`$ plane at $`|z|=1/\delta `$, and inserted the identity there. This boundary in the $`z`$ space corresponds to two regions in the $`t`$ space: $`z\mathrm{}`$ maps to $`t\mathrm{}`$ as well as $`t1/2`$. Thus in the $`t`$ plane we will have a boundary near infinity, as well as in a small disc cut out around $`t=1/2`$. We will have the identity inserted at each of these boundaries; the precise state with norm will be defined in the subsection below.
We now note that we have simply a path integral over a single free boson $`X`$ on $`\mathrm{\Sigma }`$ \- there is no twist operator left in the problem, and any boundaries present on $`\mathrm{\Sigma }`$ carry only the identity state. We will show below how to close these holes in $`\mathrm{\Sigma }`$, and then we will have just a path integral over a closed surface to compute.
The method of passing to a covering space to analyze orbifold correlation functions has been studied by many authors, for example . The observation that for symmetric orbifolds one gets a single copy of the target space with no nontrivial operator insertions on the covering space is implicit in . The notion of passing to the covering space to take into account the twist operators for symmetric orbifolds is also used in the computation of partition functions in . The $`Z_2`$ orbifold is the same as the $`S_2`$ orbifold, and the map to the covering space was used in to find the 4-point correlation of twist operators for the $`Z_2`$ orbifold of a complex boson.
We will depart from the usual way of computing the correlation functions of twist operators, and use a different way which we describe below. The usual computation for $`\sigma \sigma `$ (and the one adopted in ) proceeds by first finding
$$f\frac{X^i(z)X^i(w)\sigma (z_1)\sigma (z_2)}{\sigma (z_1)\sigma (z_2)}$$
(2.12)
by looking at the singularities of $`f`$ as a function of $`z,w`$, and constructing a function with these singularities. One then takes the limit $`zw`$, subtracts the singularity and constructs the stress tensor $`T=\frac{1}{2}X^iX^i(w)`$. Next, one uses the conformal Ward identity to relate $`T\sigma \sigma `$ to $`\sigma \sigma `$, thus obtaining an expression for $`_z\mathrm{log}\sigma (z)\sigma (0)`$. Solving this equation gives the functional form of the 2-point function, and the dimension of $`\sigma `$ can be read off from the solution.
A similar analysis can be done for the 3-point function, but the functional form of the 3-point function of primary fields is determined by their dimensions, and would tell us nothing new. One cannot find the fusion coefficients $`C_{ijk}`$ between the twist operators from the 3-point analysis because the method does not determine the overall normalization of the correlator. Thus to find the $`C_{ijk}`$ one applies the method to the 4-point function $`\sigma \sigma \sigma \sigma `$, finds the functional form of this correlator, and then uses factorization to extract the $`C_{ijk}`$.
To be able to use such a method one must have a simple stress tensor which can be written as a product of fields, each of which has a simple known behavior near the twist operators. One must also use inspection to construct the correlators like (2.12). Further, to find the $`C_{ijk}`$ we need to go up to the 4-point function.
The method we use will apply to $`S_N`$ orbifolds, but not for example to $`Z_N`$ orbifolds with $`N>2`$. On the other hand we will not need that the stress tensor have a simple form (in fact we will not use the stress tensor at all). Further, we can compute the $`C_{ijk}`$ using only the 2 and 3-point functions. Thus the method is suited to the computation of correlation functions for CFTs arising from sigma models with target space $`M^N/S_N`$, which arise in D-brae physics. The method also brings out the fact that many quantities for symmetric orbifolds are ‘universal’ in the sense that they do not depend on the details of the the manifold $`M`$.
### 2.3 Closing the punctures
For the discussion below the order and number of twist operators can be arbitrary, but for explicitness we assume that the covering space $`\mathrm{\Sigma }`$ has the topology of a sphere, and we use the correlator $`\sigma _2\sigma _2`$ as an illustrative example. It will be evident that no new issues arise for other correlators or when $`\mathrm{\Sigma }`$ has higher genus; we will mention the changes for higher genus where relevant.
As it stands the covering surface $`\mathrm{\Sigma }`$ that we have constructed has several ‘holes’ in it. We will now give a prescription for closing these holes, thus making a closed surface which we also call $`\mathrm{\Sigma }`$. The prescription for closing the holes will amount to defining precisely the states to be inserted at various boundaries. If the surface is closed then we can use the Liouville action to find the path integral after change of metric; on an open surface the boundary states can change as well.
The holes are of the following kinds:
(i) The holes in the finite $`z`$ plane at the insertion of the twist operators. These holes are circles with radius $`ϵ`$ in the $`z`$ plane, and lift to holes in $`\mathrm{\Sigma }`$ under the map $`t(z)`$.
(ii) The holes in $`\mathrm{\Sigma }`$ at finite values of $`t`$, arising from the fact that we have cut the $`z`$ space at $`|z|=1/\delta `$. These holes are located at points $`t_0`$ where the map behaves as $`z\frac{1}{tt_0}`$.
(iii) The hole in $`\mathrm{\Sigma }`$ at $`t=\mathrm{}`$, which also arises from the fact that the $`z`$ space is cut at $`|z|=1/\delta `$. If there is no twist operator at $`z=\mathrm{}`$, then we have $`zt`$ for large $`t`$, and the hole in $`\mathrm{\Sigma }`$ is the image of $`|z|=1/\delta `$ under this map.
We first complete our definition of the twist operators by defining the state inserted at the edge of the cut out hole; this addresses holes of type (i) above. We had said above that the twist operator $`\sigma _{12}`$ imposes the boundary condition (2.5), but this does not specify the operator completely. In fact there are an infinite number of operators, with increasing dimensions, which all create the same twist and in general create some further excitation of the fields. We define $`\sigma _{12}`$ to be the operator from this family with lowest dimension. After mapping the problem to the $`t`$ space, we need to ask what operator insertions at the points $`t=0`$, $`t=1`$ give the slowest power law fall off for the partition function when the distance $`a`$ between the twist operators is increased. The answer is of course that we must insert a multiple of the identity operator at the punctures in the $`t`$ space.
But we will also need to know the norm of this state, and would thus like to construct it through a path integral. Thus let the covering surface be locally defined through (2.7). For $`|t|>ϵ^{1/2}`$ the metric on the covering surface is the one induced from the $`z`$ plane
$$ds^2=dzd\overline{z}=dtd\overline{t}|\frac{dz}{dt}|^2=4|t|^2dtd\overline{t},|t|>ϵ^{1/2}.$$
(2.13)
We ‘close the hole’ in the $`t`$ space by choosing, for $`|t|<ϵ^{1/2}`$, the metric
$$ds^2=4ϵdtd\stackrel{~}{t},|t|<ϵ^{1/2}.$$
(2.14)
Thus we have glued in a ‘flat patch’ in the $`t`$ space to close the hole created in the definition of the twist operator. The metric is continuous across the boundary $`|t|=ϵ^{1/2}`$, but there is curvature concentrated along this boundary. The path integral of $`X`$ over the disc $`|t|<ϵ^{1/2}`$ creates the required state along the edge of the hole. The map (2.7) is only the leading order approximation to the actual map in general, but our prescription is to ‘close with a flat patch’ the hole in the $`t`$ space, where the hole is the image on $`\mathrm{\Sigma }`$ of a circular hole in the $`z`$ plane. As $`ϵ0`$, the small departure of the map from the form (2.7) will cease to matter.
Note that we could have chosen a different metric to replace the choice (2.14) inside the hole, but this would just correspond to a different overall normalization of the twist operator. (Thus it would be like taking a different choice of $`ϵ`$.) Once we make the choice (2.14) then we must use the same construction of the twist operator in all correlators, and then the non–universal choices in the definitions will cancel out.
The other holes, of types (ii) and (iii), arise from the hole at infinity in the $`z`$ plane, and we proceed by first replacing the $`z`$ plane by a closed surface. We take another disc with radius $`1/\delta `$ (parameterized by a coordinate $`\stackrel{~}{z}`$) and glue it to the boundary of the $`z`$ plane. Thus we get a sphere with metric given by
$`ds^2`$ $`=`$ $`dzd\overline{z},|z|<{\displaystyle \frac{1}{\delta }},`$
$`=`$ $`d\stackrel{~}{z}d\overline{\stackrel{~}{z}},|\stackrel{~}{z}|<{\displaystyle \frac{1}{\delta }},`$
$`\stackrel{~}{z}`$ $`=`$ $`{\displaystyle \frac{1}{\delta ^2}}{\displaystyle \frac{1}{z}}.`$ (2.15)
The path integral over the second disc defines a state at the boundary of the first disc. This state is proportional to the identity. But further, our explicit construction gives the state a known norm, which is something we needed to completely define the path integrals like (2.4) and (2.6).
Since the $`z`$ space is closed at infinity, we find that the holes of type (ii) and (iii) are now automatically closed in $`\mathrm{\Sigma }`$ – since we make $`\mathrm{\Sigma }`$ as a cover of the $`z`$ sphere with metric on every patch induced from the metric on this $`z`$ sphere.
The space $`\mathrm{\Sigma }`$ is now a closed surface with a certain metric, and the path integral giving $`Z_{ϵ,\delta }[\sigma _2(z_1),\sigma _2(z_2)]`$ in (2.6) is to be carried out on this closed surface.
### 2.4 The method of calculation
We had found above that if we have twist operators $`\sigma _{12}`$ at $`z=0`$, $`z=a`$, then the partition function with the twist operators inserted equals that for a single field $`X`$ on the double cover $`\mathrm{\Sigma }`$ of the $`z`$ plane given by the map (2.10). $`\mathrm{\Sigma }`$ is also a sphere with infinitesimal holes cut out, but since only the identity operator is inserted at these punctures, we can close the holes and just get the partition function of $`X`$ on a closed surface $`\mathrm{\Sigma }`$.
At this point one might wonder that since this partition function is some given number, how does it depend on the parameter $`a`$, which gave the separation between the twist operators in the $`z`$ plane? The point is that even though $`\mathrm{\Sigma }`$ is a sphere for all $`a`$, the metric on this sphere depends on $`a`$ – this is evident from (2.9).
We can compute the partition function of $`X`$ on $`\mathrm{\Sigma }`$ using some fixed fiducial metric $`\widehat{g}`$ on the $`t`$ space, but we must then take into account the conformal anomaly, which says that if $`ds^2=e^\varphi d\widehat{s}^2`$, then the partition function $`Z^{(s)}`$ computed with the metric $`ds^2`$ is related to the partition function $`Z^{(\widehat{s})}`$ computed with $`d\widehat{s}^2`$ through
$$Z^{(s)}=e^{S_L}Z^{(\widehat{s})},$$
(2.16)
where
$$S_L=\frac{c}{96\pi }d^2t\sqrt{g^{(\widehat{s})}}[_\mu \varphi _\nu \varphi g^{(\widehat{s})\mu \nu }+2R^{(\widehat{s})}\varphi ]$$
(2.17)
is the Liouville action . Here $`c`$ is the central charge of the CFT. Since we are considering the theory of a single free field $`X`$ on $`\mathrm{\Sigma }`$, we have $`c=1`$.
Let us choose the fiducial metric $`\widehat{g}`$ on $`\mathrm{\Sigma }`$ to be (in the case where $`\mathrm{\Sigma }`$ is a sphere)
$`d\widehat{s}^2`$ $`=`$ $`dtd\overline{t},|t|<{\displaystyle \frac{1}{\delta ^{}}}`$
$`=`$ $`d\stackrel{~}{t}d\overline{\stackrel{~}{t}},|\stackrel{~}{t}|<{\displaystyle \frac{1}{\delta ^{}}},`$
$`\stackrel{~}{t}`$ $`=`$ $`{\displaystyle \frac{1}{\delta ^2}}{\displaystyle \frac{1}{t}}.`$ (2.18)
We will let $`\delta ^{}0`$ at the end. Thus we have chosen the fiducial metric on the $`t`$ space to be the flat metric of a plane up to a large radius $`1/\delta ^{}`$, after which we glue an identical disc at the boundary to obtain the topology of a sphere, just as we did for the $`z`$ space.
From (2.16), (2.17) we see that if we increase $`\varphi `$ by a constant, then $`Z`$ changes by a known factor. Using that fact that for a sphere $`\sqrt{g}R=8\pi `$, we find that the partition function of $`X`$ on the sphere with metric (2.18) is
$$Z_\delta ^{}=Q(\delta ^{})^{\frac{c}{3}}=Q(\delta ^{})^{\frac{1}{3}}.$$
(2.19)
Thus we will have $`Z^{(\widehat{s})}=Z_\delta ^{}`$. Here $`Q`$ is a constant that is regularization dependent and cannot determined by anything that we have chosen so far. ($`Q`$ determines the size of the sphere for which the partition function will attain the value unity; since the CFT has no inbuilt scale we cannot find the value of this size in any absolute way.) $`Q`$ will cancel out in all final calculations.
The partition function of one boson on the $`z`$ sphere with the metric (2.15) is $`Z_\delta `$ (cf. eq. (2.4)), and we have
$$Z_\delta =Q\delta ^{\frac{c}{3}}=Q\delta ^{\frac{1}{3}}.$$
(2.20)
### 2.5 Contributions to $`S_L`$
The partition function with twist operators inserted can be written as
$$Z_{ϵ,\delta }[\sigma _{n_1}(z_1),\mathrm{},\sigma _{n_k}(z_k)]=e^{S_L}Z^{(\widehat{s})}.$$
(2.21)
Thus the computation of the correlation function boils down to computing $`S_L`$. There are three types of contributions to $`S_L`$, which we will analyze separately
$$S_L=S_L^{(1)}+S_L^{(2)}+S_L^{(3)},$$
(2.22)
$`S_L^{(1)}`$ will give the essential numerical contributions to the correlation functions (as well as regulation dependent quantities), while $`S_L^{(2)}`$ and $`S_L^{(3)}`$ give only regulation dependent quantities; regulation parameters cancels out at the end.
(a) We have cut out various discs from the $`z`$ plane where the physical theory is defined: we have removed infinity by taking $`|z|<1/\delta `$ and have also cut out circles of radius $`ϵ`$ around the twist operator insertions. Let us call this region of $`z`$ the ‘regular region’. This ‘regular’ region of the $`z`$ space has an image in the $`t`$ space, which we call the ‘regular region’ on $`\mathrm{\Sigma }`$. On $`\mathrm{\Sigma }`$ we will find, apart from the obvious cuts around the images of the twist operators and a cut near $`|t|=\mathrm{}`$, further possible cuts around images of $`z=\mathrm{}`$ as discussed in subsection 2.3. Let the contribution to $`S_L`$ from this ‘regular region’ of $`\mathrm{\Sigma }`$ be called $`S_L^{(1)}`$.
To evaluate (2.21) we need to choose a fiducial metric on the $`t`$ space. Suppose that the map $`z(t)`$ has the form $`zbt`$ as $`t\mathrm{}`$ . (When there is no twist operator at infinity the map can be taken to have this form.) Let this fiducial metric $`d\widehat{s}^2`$ be of the form (2.18) with
$$\frac{1}{\delta }<\frac{b}{\delta ^{}}.$$
(2.23)
With this choice the boundary $`|z|=1/\delta `$ gets mapped to a curve inside the disc $`|t|<1/\delta ^{}`$ (i.e. into the ‘first half’ of the $`t`$ sphere).
In this ‘regular region’ of $`\mathrm{\Sigma }`$, the fiducial metric (2.18) is flat, and so there is no contribution from the $`R\varphi `$ term in (2.17). Thus we have
$$S_L^{(1)}=\frac{1}{96\pi }d^2t[_\mu \varphi ^\mu \varphi ],$$
(2.24)
where the integral extends over the region described above. We rewrite (2.24) as
$$S_L=\frac{1}{96\pi }d^2t[\varphi _\mu ^\mu \varphi ]+\frac{1}{96\pi }_{}\varphi _n\varphi .$$
(2.25)
Here $``$ is the boundary of the ‘regular region’ of $`\mathrm{\Sigma }`$, and $`_n`$ is the normal derivative at the boundary. From (2.9) we find that
$$\varphi =\mathrm{log}[\frac{dz}{dt}]+\mathrm{log}[\frac{d\overline{z}}{d\overline{t}}],$$
(2.26)
so that
$$_\mu ^\mu \varphi =4_t_{\overline{t}}\varphi =0,$$
(2.27)
and we get
$$S_L=\frac{1}{96\pi }_{}\varphi _n\varphi .$$
(2.28)
The boundaries of the ‘regular region’ are of two kinds: those arising from the holes of size $`|zz_i|=ϵ`$ cut around the twist operators, and those arising from the cutoff at infinity ($`|z|=1/\delta `$). Consider the boundary of the hole arising from some twist operator $`\sigma _n(z_i)`$. We regulated the twist operator by choosing this hole to be a circle in the $`z`$ plane, so we start by looking at a segment of the boundary using the coordinate $`z`$. We have
$$_n=\frac{1}{|z|}(z_z+\overline{z}_{\overline{z}}).$$
(2.29)
Writing $`z=|z|e^{i\theta }`$, one finds
$$𝑑s=|z|𝑑\theta =\frac{|z|}{i}\frac{dz}{z}=\frac{|z|}{i}\frac{d\overline{z}}{\overline{z}}.$$
(2.30)
Thus we get
$$_{}𝑑s\varphi _n\varphi =i𝑑z\varphi _z\varphi +c.c.$$
(2.31)
Since $`z`$ is holomorphically dependent on $`t`$, we can write
$$dz_z\varphi =dt_t\varphi .$$
(2.32)
We can thus write for the contribution to $`S_L`$ from any hole
$$\frac{1}{96\pi }_{}ds\varphi _n\varphi =\frac{1}{96\pi }[idt\varphi _t\varphi +c.c.],$$
(2.33)
where $`\varphi `$ is given through (2.26). A similar analysis applies to all the other boundaries of the ‘regular region’ on $`\mathrm{\Sigma }`$, and we compute (2.33) for each such boundary. Since the ‘holes’ on $`\mathrm{\Sigma }`$ are infinitesimal size punctures, computing (2.33) needs only the leading order behavior of $`\varphi `$ at the punctures.
(b) We had cut out holes of radius $`ϵ`$ in the $`z`$ plane around the insertions of the twist operators, and these gave corresponding holes in the ‘regular region’ of $`\mathrm{\Sigma }`$. We now compute the contribution to $`S_L`$ from the part $`H`$ of $`\mathrm{\Sigma }`$ that is used to close such a hole. Since we had closed these holes with the flat metric (2.14), and since the fiducial metric we use on $`\mathrm{\Sigma }`$ is also flat in $`H`$ ($`d\widehat{s}^2=dtd\overline{t}`$), we get $`\varphi =\mathrm{constant}`$, and so there is no contribution from the kinetic term in (2.17). Note that at the boundary of $`H`$ we have $`_t\varphi `$ nonzero but bounded, then since the area of this boundary is zero (the boundary is one-dimensional) we get no contribution to the kinetic term from the boundary either. The curvature term in (2.17) is zero, since the curvature of the fiducial metric is zero throughout the region where the twist operators are inserted. Thus we get no contribution to $`S_L`$ from these regions $`H`$ of $`\mathrm{\Sigma }`$.
(c) Now consider the contributions from the points that have finite $`t`$, but $`z\mathrm{}`$. The ‘regular region’ on $`\mathrm{\Sigma }`$ had excluded the image of $`|z|>1/\delta `$. This image will have a small disc $`D`$ around some finite $`t_0`$, if we have
$$z\frac{\alpha }{tt_0}+\beta +\mathrm{}$$
(2.34)
The fiducial metric we are using on $`\mathrm{\Sigma }`$ is flat here, so there is no contribution from the curvature term in (2.17). The region inside the disc $`D`$ has a metric induced from the ‘second half’ of the $`z`$ sphere (i.e. the part parameterized by $`\stackrel{~}{z}`$ in (2.15)) so that the metric is $`ds^2=d\stackrel{~}{z}d\overline{\stackrel{~}{z}}`$. Thus
$`\stackrel{~}{z}`$ $`=`$ $`{\displaystyle \frac{1}{\delta ^2}}{\displaystyle \frac{1}{z}}{\displaystyle \frac{1}{\delta ^2\alpha }}(tt_0){\displaystyle \frac{\beta }{\delta ^2\alpha ^2}}(tt_0)^2,`$
$`\varphi `$ $`=`$ $`\mathrm{log}{\displaystyle \frac{d\stackrel{~}{z}}{dt}}+c.c.\mathrm{log}{\displaystyle \frac{1}{\delta ^2\alpha }}{\displaystyle \frac{2\beta }{\alpha }}(tt_0)+c.c.,`$
$`_t\varphi `$ $``$ $`{\displaystyle \frac{2\beta }{\alpha }}.`$ (2.35)
The area of the disc $`D`$ in the fiducial metric is $`\pi |tt_0|^2\pi (|\alpha |\delta )^2`$. As $`\delta 0`$, we find that $`d^2t_t\varphi _{\overline{t}}\varphi 0`$. Thus we get no contribution to $`S_L`$ from these images of the cut at infinity.
(d) Now we look at the region of $`\mathrm{\Sigma }`$ near $`t=\mathrm{}`$. Let $`zbt`$ for large $`t`$. Let the image of $`|z|=1/\delta `$ be the contour $`C`$ on $`\mathrm{\Sigma }`$. By the choice (2.23) and the fact that $`C`$ satisfies $`|t|\frac{1}{b\delta }`$, we find that $`C`$ is inside the curve $`|t|=1/\delta ^{}`$. Let the contribution to $`S_L`$ from the region between $`C`$ and $`|t|=1/\delta ^{}`$ be called $`S_L^{(2)}`$.
Since the fiducial metric (2.18) is flat in this region, there is no contribution from the curvature term in (2.17). For the kinetic term we have
$`\stackrel{~}{z}`$ $`=`$ $`{\displaystyle \frac{1}{\delta ^2z}}{\displaystyle \frac{1}{\delta ^2bt}},{\displaystyle \frac{d\stackrel{~}{z}}{dt}}{\displaystyle \frac{1}{\delta ^2b}}{\displaystyle \frac{1}{t^2}},`$
$`_t\varphi `$ $``$ $`{\displaystyle \frac{2}{t}},{\displaystyle d^2t\varphi \varphi }32\pi \mathrm{log}{\displaystyle \frac{\delta ^{}}{|b|\delta }}.`$ (2.36)
and we get
$$S_L^{(2)}=\frac{1}{3}\mathrm{log}\frac{\delta ^{}}{|b|\delta }.$$
(2.37)
(e) Moving further outwards in the $`t`$ plane, we find a ‘ring of curvature’ at $`|t|=1/\delta ^{}`$ (cf. eq. (2.18)). At this ring we have
$`\varphi \mathrm{log}{\displaystyle \frac{(\delta ^{})^2}{\delta ^2b}}+c.c.,{\displaystyle }d^2tR\varphi =8\pi \varphi ,`$ (2.38)
which gives a contribution to $`S_L`$ equal to
$$S_L^{(3)}=\frac{1}{3}\mathrm{log}\frac{(\delta ^{})^2}{\delta ^2|b|}.$$
(2.39)
The kinetic term in $`\varphi `$ has no contribution at this ring. Further, the region $`|t|>1/\delta `$ gives no contribution to $`S_L`$, since the curvature of the fiducial metric is zero, and the map gives $`\varphi =\mathrm{constant}`$.
### 2.6 The correlator in terms of the Liouville field
Let us collect all the above contributions together. Note that
$$S_L^{(2)}+S_L^{(3)}=\frac{1}{3}\mathrm{log}\frac{\delta ^{}}{\delta }.$$
(2.40)
so that the variable $`b`$ drops out of this combination.
Now let us go back to the expression (2.6) that we want to evaluate :
$$\sigma _n^ϵ(0)\sigma _n^ϵ(a)_\delta =\frac{Z_{ϵ,\delta }[\sigma _n(z_1),\sigma _n(z_2)]}{(Z_\delta )^N}=e^{S_L}\frac{Z^{(\widehat{s})}}{(Z_\delta )^n},$$
(2.41)
here we used equation (2.16). Taking into account the relations (2.19) and (2.20) we finally get:
$$\sigma _n^ϵ(0)\sigma _n^ϵ(a)_\delta =e^{S_L}\left(\frac{\delta ^n}{\delta ^{}}\right)^{1/3}Q^{1n}.$$
(2.42)
Substituting the expression for the Liouville action, $`S_L`$, we conclude that
$$\sigma _n^ϵ(0)\sigma _n^ϵ(a)_\delta =e^{S_L^{(1)}}e^{S_L^{(2)}+S_L^{(3)}}\left(\frac{\delta ^n}{\delta ^{}}\right)^{1/3}Q^{1n}=e^{S_L^{(1)}}\delta ^{\frac{n1}{3}}Q^{1n}.$$
(2.43)
Thus we observe a cancellation of $`\delta ^{}`$, which served only to choose a fiducial metric on $`\mathrm{\Sigma }`$ and thus should not appear in any final result. The only quantity that needs computation is $`S_L^{(1)}`$ using (2.33).
Let us mention that formula (2.41) has a simple extension to the case of a general correlation function:
$$\sigma _{n_1}^ϵ(z_1)\mathrm{}\sigma _{n_k}^ϵ(z_k)_\delta =e^{S_L}\frac{Z^{(\widehat{s})}}{(Z_\delta )^s},$$
(2.44)
where $`Z^{(\widehat{s})}`$ is a partition function of the covering Riemann surface $`\mathrm{\Sigma }`$ with the fiducial metric $`d\widehat{s}^2`$ ($`\mathrm{\Sigma }`$ may have any genus), and $`s`$ is a number of fields involved in nontrivial permutation ($`s=n`$ in the case of the two point function (2.41)). The partition function $`Z^{(\widehat{s})}`$ may depend on the moduli of the surface $`\mathrm{\Sigma }`$ and its size (there are no moduli in the case of the sphere and the size is parameterized by $`\delta ^{}`$).
## 3 The 2-point function
### 3.1 The calculation
Let us apply the above scheme to evaluate the 2-point function of twist operators. If one of the twist operators corresponds to the permutation
$$(1\mathrm{}n),$$
(3.1)
then the other one should correspond to the permutation
$$(n\mathrm{}1),$$
(3.2)
since otherwise the correlation function vanishes. Thus we can write
$$\sigma _n(0)\sigma _n(a)$$
(3.3)
instead of $`\sigma _{(1\mathrm{}n)}(0)\sigma _{(n\mathrm{}1)}(a)`$ without causing confusion.
The generalization of the map (2.10) to the case of $`\sigma _n`$ is
$$z=a\frac{t^n}{t^n(t1)^n}.$$
(3.4)
For this map we have
$`\varphi `$ $`=`$ $`\mathrm{log}|{\displaystyle \frac{dz}{dt}}|^2=\mathrm{log}[an{\displaystyle \frac{t^{n1}(t1)^{n1}}{(t^n(t1)^n)^2}}]+c.c.,`$ (3.5)
$`{\displaystyle \frac{d\varphi }{dt}}`$ $`=`$ $`{\displaystyle \frac{(2t+n1)(t1)^n(2tn1)t^n}{t(t1)((t1)^nt^n)}}`$
This map has the branch points located at
$$t=0z=0\text{and}t=1z=a.$$
(3.6)
There are $`n`$ images of the point $`z=\mathrm{}`$ in $`t`$ plane:
$$t_k=\frac{1}{1\alpha _k},\alpha _k=e^{\frac{2\pi ik}{n}},k=0,1,\mathrm{},n1.$$
(3.7)
We note that $`\alpha _0=1`$ gives $`t=\mathrm{}`$.
Let us compute the contribution (2.33) for the point $`z=0`$. Near this point we have:
$`z(1)^{n+1}at^n,|t|{\displaystyle \frac{|z|^{1/n}}{a^{1/n}}},`$ (3.8)
$`\varphi \mathrm{log}[ant^{n1}],_t\varphi {\displaystyle \frac{n1}{t}}.`$
Then we get the contribution to the Liouville action (2.25):
$$S_L(t=0)=\frac{n1}{12}\left[\frac{1}{n}\mathrm{log}|a|+\mathrm{log}(nϵ^{\frac{n1}{n}})\right].$$
(3.9)
By a reflection symmetry $`t1t`$, $`zaz`$, we get the same contribution from the other branch point:
$$S_L(t=1)=\frac{n1}{12}\left[\frac{1}{n}\mathrm{log}|a|+\mathrm{log}(nϵ^{\frac{n1}{n}})\right].$$
(3.10)
Now we look at the images of infinity. First we note that the integral over the boundary located near $`t=\mathrm{}`$ will give zero, since $`d\varphi /dt`$ goes like $`1/t^2`$, the length of the circle goes like $`t`$ and the value of $`\varphi `$ is at best logarithmic in $`t`$. But we do get a contribution from the images of $`z=\mathrm{}`$ located at finite points in $`t`$ plane. Note that
$$\text{if}t=\frac{1}{1\alpha _k}+x,\text{then}(t1)^nt^n\frac{xn}{\alpha _k(1\alpha _k)^{n2}}.$$
(3.11)
This leads to
$`z{\displaystyle \frac{a\alpha _k}{n(1\alpha _k)^2}}{\displaystyle \frac{1}{x}},x{\displaystyle \frac{a\alpha _k}{n(1\alpha _k)^2}}{\displaystyle \frac{1}{z}},`$ (3.12)
$`\varphi =\mathrm{log}\left[a^1n(1\alpha _k)^2\alpha _k^{n1}z^2\right],_t\varphi {\displaystyle \frac{2}{t(1\alpha _k)^1}}.`$ (3.13)
The point $`t=t_k`$ we are considering gives following contribution to the Liouville action:
$$S_L(t=t_k)=\frac{1}{6}\mathrm{log}\left[a^1n(1\alpha _k)^2\alpha _k^{n1}\delta ^2\right].$$
(3.14)
Thus the total contribution from the images of infinity is:
$$S_L(z=\mathrm{})=\frac{1}{6}\underset{k=1}{\overset{n1}{}}\mathrm{log}\left[\frac{n(1\alpha _k)^2\alpha _k^{n1}}{a\delta ^2}\right]=\frac{n1}{6}\mathrm{log}[|a|\delta ^2]+\frac{n+1}{6}\mathrm{log}[n].$$
(3.15)
We have used the following properties of $`\alpha _k`$:
$$\underset{k=1}{\overset{n1}{}}\alpha _k=1;\underset{k=1}{\overset{n1}{}}(q\alpha _k)=\frac{q^n1}{q1}n,\text{if}q1,$$
(3.16)
which follow from the fact that $`\{\alpha _k\}`$ is the set of different solutions of the equation $`\alpha ^n1=0`$ and $`\alpha _0=1`$.
Adding all the contributions together, we get an expression for the interesting part of the Liouville action:
$$S_L^{(1)}=\frac{1}{6}\left[(n\frac{1}{n})\mathrm{log}|a|+\frac{(n1)^2}{n}\mathrm{log}ϵ+2(n1)\mathrm{log}\delta 2\mathrm{log}n\right].$$
(3.17)
This leads to the final expression for the correlation function (see (2.43)):
$$\sigma _n^ϵ(0)\sigma _n^ϵ(a)_\delta =e^{S_L^{(1)}}\delta ^{\frac{n1}{3}}Q^{1n}=a^{\frac{1}{6}(n\frac{1}{n})}C_nϵ^{A_n}Q^{B_n},$$
(3.18)
$$A_n=\frac{(n1)^2}{6n},B_n=1n,C_n=n^{1/3}.$$
(3.19)
Thus we read off the dimension $`\mathrm{\Delta }_n`$ of $`\sigma _n`$
$$\mathrm{\Delta }_n=\frac{1}{24}(n\frac{1}{n})$$
(3.20)
The other constants in (3.18) are to be absorbed into the normalization of $`\sigma _n`$. We will discuss this renormalization after computing the 3-point functions.
### 3.2 ‘Universality’ of the 2-point function
The theory we have considered above is that of the orbifold $`M^N/S_N`$ where the manifold $`M`$ is just $`R`$, the real line. If $`M`$ was $`R^d`$ instead, we could treat the $`d`$ different species of fields independently, and obtain
$$\mathrm{\Delta }_n=\frac{c}{24}(n\frac{1}{n})$$
(3.21)
where $`c=d`$ is the central charge of the CFT for one copy of $`M=R^d`$. But we see that we would obtain the result (3.21) for the symmetric orbifold with any choice of $`M`$; we just use the value of $`c`$ for the CFT on $`M`$. Around the insertion of the twist operator we permute the copies of $`M`$, but the definition of the twist operator does not involve directly the structure of $`M`$ itself. The Liouville action (2.17) determines the correlation function using only the value $`c`$ of the CFT. Thus we recover the result (3.21) for any $`M`$.
This ‘universality’ of $`\mathrm{\Delta }_n`$ is well known, and the value of $`\mathrm{\Delta }_n`$ can be deduced from the following standard argument. Consider the CFT on a cylinder parameterized by $`w=x+iy,0<y<2\pi `$. At $`x\mathrm{}`$ let the state be the vacuum of the orbifold CFT $`M^N/S_N`$. Since there is no twist, each copy of $`M`$ gives its own contribution to the vacuum energy, which thus equals $`\frac{c}{24}N`$. Now insert the twist operator $`\sigma _n`$ at $`w=0`$, and look at the state for $`x\mathrm{}`$. The copies of $`M`$ not involved in the twist contribute $`\frac{c}{24}`$ each as before, but those that are twisted by $`\sigma _n`$ turn into effectively one copy of $`M`$ defined on a circle of length $`2\pi n`$. Thus the latter set contribute $`\frac{c}{24n}`$ to the vacuum energy. The change in the energy between $`x=\mathrm{}`$ and $`x=\mathrm{}`$ gives the dimension of $`\sigma _n`$ (since the state at $`x\mathrm{}`$ is the vacuum)
$$\frac{c}{24n}[\frac{cn}{24}]=\frac{c}{24}(n\frac{1}{n})=\mathrm{\Delta }_n$$
(3.22)
Thus while our calculation of the 2-point function has not taught us anything new, we have obtained a scheme that will yield the higher point functions for symmetric orbifolds using an extension of the same universal features that gave the value of $`\mathrm{\Delta }_n`$ in the above argument.
## 4 The map for the 3-point function
### 4.1 Genus of the covering surface
Let us first discuss the nature of the covering surface $`\mathrm{\Sigma }`$ for the case where we have an arbitrary number of twist operators in the correlation function $`\sigma _{n_1}\sigma _{n_2}\mathrm{}\sigma _{n_k}`$. The CFT is still defined on the plane $`z`$, which we will for the moment regard as a sphere by including the point at infinity.
At the insertion of the operator $`\sigma _{n_j}(z_j)`$ the covering surface $`\mathrm{\Sigma }`$ has a branch point of order $`n_j`$, which means that $`n_j`$ sheets of $`\mathrm{\Sigma }`$ meet at $`z_j`$. One says that the ramification order at $`z_j`$ is $`r_j=n_j1`$. Suppose further that over a generic point $`z`$ here are $`s`$ sheets of the covering surface $`\mathrm{\Sigma }`$. Then the genus $`g`$ of $`\mathrm{\Sigma }`$ is given by the Riemann–Hurwitz formula:
$$g=\frac{1}{2}\underset{j}{}r_js+1$$
(4.1)
Let us now consider the 3-point function. We require each twist operator to correspond to a single cycle of the permutation group, and regard the product of two cycles to represent the product of two different twist operators. Let the cycles have lengths $`n,m,q`$ respectively. It is easy to see that we can obtain covering surfaces $`\mathrm{\Sigma }`$ of various genera. For example, if we have
$$\sigma _{12}\sigma _{13}\sigma _{123}$$
(4.2)
as the three permutations, then we have $`r_1=1,r_2=1,r_3=2,s=3`$, and we get $`g=0`$. On the other hand with
$$\sigma _{123}\sigma _{123}\sigma _{123}$$
(4.3)
we get $`r_1=r_2=r_3=2,s=3`$ and we get $`g=1`$. (This genus 1 surface is a singular limit of the torus, however.)
Let us concentrate on the case where we get $`g=0`$. Without loss of generality we can take the first permutation $`\sigma _n`$ to be the cycle
$$(1,2,\mathrm{}k,k+1,\mathrm{}n)$$
(4.4)
The second permutation is restricted by the requirement that when composed with (4.4) it yields a single cycle (which would be the conjugate permutation of the third twist operator). In addition we must have a sufficiently small number of indices in the result of the first two permutations so that we do get $`g=0`$. A little inspection shows that $`\sigma _m`$ must have the form
$$(k,k1,\mathrm{}1,n+1,n+2,\mathrm{}n+mk)$$
(4.5)
Thus the elements $`1,2,\mathrm{}k`$ of the first permutation occur in the second permutation in the reverse order, and then we have a new set of elements $`n+1,\mathrm{}n+mk`$. These two permutations compose to give the cycle $`\sigma _m\sigma _n`$ equal to
$$(k+1,k+2,\mathrm{}n,1,n+1,n+2,\mathrm{}n+mk)$$
(4.6)
Thus $`\sigma _q`$ must be the inverse of the cycle (4.6), and we have
$$q=n+m2k+1$$
(4.7)
Note that the number of ‘overlaps’ (i.e., common indices) between $`\sigma _n`$ and $`\sigma _m`$ is $`k`$. Note that we must have $`k1`$ in order that the product $`\sigma _m\sigma _n`$ be a single cycle rather than just a product of two cycles. Also note that if we have $`q=n+m1`$, then since $`sq`$, (4.1) gives that $`\mathrm{\Sigma }`$ must have genus zero (this will be a ‘single overlap’ correlator).
Let $`\mathrm{\Sigma }`$ be the covering surface that corresponds to the insertions
$`\sigma _n(z_1)\sigma _m(z_2)\sigma _q(z_3)`$. Then the number of sheets of $`\mathrm{\Sigma }`$ over a generic point $`z`$ is just the total number of indices used in the permutations
$$s=n+mk=\frac{1}{2}(n+m+q1)$$
(4.8)
Thus the genus of $`\mathrm{\Sigma }`$ is
$$\frac{n1}{2}+\frac{m1}{2}+\frac{q1}{2}s+1=0$$
(4.9)
### 4.2 The map for the case $`g=0`$
We are looking for a covering surface of the sphere that is ramified at three points on the sphere, with a finite order of ramification at each point. We look for the map from $`z`$ to $`\mathrm{\Sigma }`$ as a ratio of two polynomials
$$z=\frac{f_1(t)}{f_2(t)};$$
(4.10)
the existence of such a map will be evident from its explicit construction. By using the $`SL(2,C)`$ symmetry group of the $`z`$ sphere, we will place the twist operators $`\sigma _n,\sigma _m,\sigma _q`$ at $`z=0,z=a,z=\mathrm{}`$ respectively. We can assume without loss of generality that
$$nq,mq.$$
(4.11)
Note that we had placed a cutoff in the $`z`$ plane to remove the region at infinity, and will not be immediately clear how to normalize a twist that occurs around the circle at infinity. We will discuss this issue of normalization later.
By making an $`SL(2,C)`$ transformation $`t^{}=\frac{at+b}{ct+d}`$ of the surface $`\mathrm{\Sigma }`$, which we assume is parameterized by the coordinate $`t`$, we can take
$$z(t=0)=0,z(t=\mathrm{})=\mathrm{},z(t=1)=a,$$
(4.12)
Note that this $`SL(2,C)`$ transformation maintains the form (4.10) of $`z`$ to be a ratio of two polynomials, and we will use the symbols $`f_1,f_2`$ to denote the polynomials after the choice (4.12) has been made.
Since we need $`s`$ values of $`t`$ for a generic value of $`z`$, with $`s`$ given by (4.8), the relation (4.10) should give a polynomial equation of order $`s`$ for $`t`$. Thus the degrees $`d_1,d_2`$ of the polynomials $`f_1,f_2`$ should satisfy:
$$\mathrm{max}(d_1,d_2)=s=\frac{1}{2}(n+m+q1).$$
Since we have chosen $`t=\mathrm{}`$ for $`z=\mathrm{}`$, we get $`d_1>d_2`$, and we have
$$d_1=\frac{1}{2}(n+m+q1).$$
(4.13)
The requirement of the proper behavior at infinity ($`zt^q`$) then gives:
$$d_2=d_1q=\frac{1}{2}(n+mq1).$$
(4.14)
Finally, the the number of indices common between the permutations $`\sigma _n(0)`$ and $`\sigma _m(a)`$ (the overlap) is
$$\frac{1}{2}(n+m(q1))=d_2+1.$$
(4.15)
Let us now look at the structure required of the map (4.10). For $`z0`$ we need
$$z=t^n(C_0+O(t))$$
(4.16)
Similarly for $`t1`$ we need
$$z=a+(t1)^m(C_1+O(t1))$$
(4.17)
Then we find
$$f_2^2\frac{dz}{dt}=f_1^{}f_2f_2^{}f_1=Ct^{n1}(1t)^{m1}.$$
(4.18)
($`C`$ is a constant). The last step follows on noting that the expression $`f_1^{}f_2f_2^{}f_1`$ is a polynomial of degree $`d_1+d_21=n+m2`$, and the behavior of $`z`$ near $`z=0,z=a`$ already provides all the possible zeros of this polynomial $`f_2^2\frac{dz}{dt}`$. The expression in (4.18) is just the Wronskian of $`f_1,f_2`$, and our knowledge of this Wronskian given an easy way to find these polynomials. We seek a second order linear differential equation whose solutions are the linear span $`f=\alpha f_1+\beta f_2`$. Such an equation is found by observing that
$$\left|\begin{array}{ccc}f& f^{}& f^{\prime \prime }\\ f_1& f_1^{}& f_1^{\prime \prime }\\ f_2& f_2^{}& f_2^{\prime \prime }\end{array}\right|=0$$
(4.19)
so that we get the equation
$$Wf^{\prime \prime }W^{}f^{}+c(t)f=0$$
(4.20)
where
$$W=f_2f_1^{}f_1f_2^{},c(t)=f_2^{}f_1^{\prime \prime }f_1^{}f_2^{\prime \prime }$$
(4.21)
Here $`W`$ is given by (4.18). The coefficient $`W^{}`$ of $`f^{}`$ is
$$W^{}=Ct^{n2}(t1)^{m2}[(n1)(n+m2)t]$$
(4.22)
The coefficient $`c(t)`$ must be a polynomial of degree $`n+m4`$ but in fact we can argue further that it must have the form
$$\gamma t^{n2}(1t)^{m2},\gamma =\mathrm{constant}$$
(4.23)
To see this look at the equation near $`t=0`$. Let $`c(t)\alpha t^k`$ with $`k<n2`$. Then the equation reads
$$t^{n1k}f^{\prime \prime }(n1)t^{n2k}f^{}+\frac{\alpha }{C}f=0$$
(4.24)
Note that the two polynomials $`f_1,f_2`$ which solve the equation must not have a common root $`t=0`$, since we assume that (4.10) is already expressed in reduced form. Thus at least one of the solutions must go like $`fconstant`$ at $`t=0`$, which is in contradiction with (4.24) since the first two terms on the LHS vanish while the last does not ($`a0`$ by definition). Thus $`c(t)`$ has a zero of order at least $`n2`$ at $`t=0`$, and by a similar argument, a zero of order at least $`m2`$ at $`t=1`$. Thus the result (4.23) follows.
Dividing through by $`Ct^{n2}(1t)^{m2}`$ we can write the equation (4.20) as
$$t(1t)f^{\prime \prime }+[(n1)+(n+m2)t]f^{}+\stackrel{~}{\gamma }=0$$
(4.25)
Let us now look at $`t\mathrm{}`$, and let the solutions to the above equation go like $`t^p`$. Then we get
$$p(p1)+p(m+n2)+\stackrel{~}{\gamma }=0.$$
(4.26)
which has the solutions
$$p_\pm =\frac{1}{2}\left(m+n1\pm \sqrt{(m+n1)^2+4\stackrel{~}{\gamma }}\right)$$
(4.27)
But since we have a twist operator of order $`q`$ at infinity, we must have
$$p_+p_{}=q,$$
(4.28)
This gives
$$\stackrel{~}{\gamma }=\frac{1}{4}(qmn+1)(q+m+n1)=d_1d_2.$$
(4.29)
Thus we have found the equation which is satisfied by both $`f_1`$ and $`f_2`$:
$$t(1t)y^{\prime \prime }+(n+1(d_1d_2+1)t)y^{}d_1d_2y=0,$$
(4.30)
which is the hypergeometric equation. Its general solution is given by
$$y=AF(d_1,d_2;n+1;t)+Bt^nF(d_1+n,d_2+n;n+1;t).$$
(4.31)
The map we are looking for can be written as
$$z=a\frac{d_2!d_1!}{n!(d_1n)!}\frac{\mathrm{\Gamma }(1n)}{\mathrm{\Gamma }(1n+d_2)}t^n\frac{F(d_1+n,d_2+n;n+1;t)}{F(d_1,d_2;n+1;t)},$$
(4.32)
where we have chosen the normalizations of $`f_1,f_2`$ such that the $`t=1`$ maps to $`z=a`$. In our case $`d_1`$, $`d_2`$ and $`n`$ are integers. Some of the individual terms in the above expression are undefined for integer $`d_1`$, $`d_2`$, $`n`$ and a limit should be taken from non–integer values of $`n`$ (while keeping $`d_1,d_2`$ fixed at their integer values). We can write the result in a well defined way by using Jacobi polynomials, which are a set of orthogonal polynomials defined through the hypergeometric function
$`P_n^{(\alpha ,\beta )}(x)`$ $``$ $`\left(\begin{array}{c}n+\alpha \\ n\end{array}\right)F(n,n+\alpha +\beta +1;\alpha +1;{\displaystyle \frac{1x}{2}})`$ (4.35)
$`=`$ $`{\displaystyle \frac{1}{n!}}{\displaystyle \underset{\nu =0}{\overset{n}{}}}\left(\begin{array}{c}n\\ \nu \end{array}\right)(n+\alpha +\beta +1)\mathrm{}(n+\alpha +\beta +\nu )`$ (4.39)
$`(\alpha +\nu +1)\mathrm{}(\alpha +n)({\displaystyle \frac{x1}{2}})^\nu `$
Then (4.32) becomes
$$z=at^nP_{d_1n}^{(n,d_1d_2+n1)}(12t)\left[P_{d_2}^{(n,d_1d_2+n1)}(12t)\right]^1.$$
(4.40)
We will have occasion to use the Wronskian of the polynomials later, and we define $`\stackrel{~}{W}`$ to be normalized as follows
$`\stackrel{~}{W}(t)`$ $`=`$ $`{\displaystyle \frac{d}{dt}}\left[t^nP_{d_1n}^{(n,d_1d_2+n1)}(12t)\right]P_{d_2}^{(n,d_1d_2+n1)}(12t)`$ (4.41)
$``$ $`t^nP_{d_1n}^{(n,d_1d_2+n1)}(12t){\displaystyle \frac{d}{dt}}P_{d_2}^{(n,d_1d_2+n1)}(12t)`$
$`=`$ $`{\displaystyle \frac{nd_1!}{n!d_2!(d_1n)!}}{\displaystyle \frac{\mathrm{\Gamma }(d_2n+1)}{\mathrm{\Gamma }(1n)}}t^{n1}(1t)^{d_1+d_2n},`$
We will also have occasion to use the relation (4.32) containing hypergeometric functions, and we define
$`W(t)`$ $`=`$ $`\left[t^nF(d_1+n,d_2+n;n+1;t)\right]^{}F(d_1,d_2;n+1;t)`$ (4.42)
$``$ $`t^nF(d_1+n,d_2+n;n+1;t)F^{}(d_1,d_2;n+1;t)`$
$`=`$ $`nt^{n1}(1t)^{d_1+d_2n}.`$
We will calculate the three point function using the map (4.32), (4.40) in the next section.
## 5 The Liouville action for the 3-point function
Let us evaluate the three point function
$$\sigma _n(0)\sigma _m(a)\sigma _q(\mathrm{})$$
(5.1)
using the map (4.32), (4.40).
Recall that we cut circles of radius $`ϵ`$ in the $`z`$ plane around the twist operators at $`z=0`$ and $`z=a`$ to regularize these twist operators. But unlike the case of the 2-point function discussed in section 3, now we have the twist operator $`\sigma _q`$ inserted at infinity. This means that the fields $`X^I`$ have boundary conditions around $`z=\mathrm{}`$ such that $`q`$ of the $`X^I`$ form a cycle under rotation around the circle $`X^{i_1}X^{i_2}\mathrm{}X^{i_q}X^{i_1}`$, while the remaining fields $`X^I`$ are single valued around this circle. Note that if the covering surface $`\mathrm{\Sigma }`$ has $`s`$ sheets over a generic $`z`$ then there will be $`sq`$ such single valued fields $`X^I`$.
The covering surface $`\mathrm{\Sigma }`$ will have punctures at $`t=0`$ and $`t=1`$ corresponding to $`z=0`$ and $`z=a`$ respectively. In addition it will have punctures corresponding to the ‘puncture at infinity’ in the $`z`$ plane. These latter punctures are of two kinds. The first kind of puncture in the $`t`$ plane will correspond to the place where $`q`$ sheets meet in the $`z`$ plane - i.e., the lift of the point where the twist operator was inserted. But we will also have $`sq`$ other punctures in the $`t`$ plane that correspond to the cut at $`|z|=1/\delta `$ for the $`X^I`$ that are single valued around $`z=\mathrm{}`$. We will choose (when defining the ‘regular region’) a cutoff at value $`|z|=1/\stackrel{~}{\delta }`$ for the first kind of puncture (i.e. the puncture arising from fields $`X^I`$ that are twisted at $`z=\mathrm{}`$) and a value $`|z|=1/\delta `$ for the second kind of puncture (i.e. punctures for fields $`X^I`$ which are not twisted at infinity). We will see that both $`\delta `$ and $`\stackrel{~}{\delta }`$ cancel from all final results.
### 5.1 The contribution from $`z=0,t=0`$
Let us first consider the point $`z=0`$ which gives $`t=0`$. Near this point the map (4.32) gives:
$$za\frac{d_2!d_1!}{n!(d_1n)!}\frac{\mathrm{\Gamma }(1n)}{\mathrm{\Gamma }(1n+d_2)}t^n,t\left(\frac{zn!(d_1n)!\mathrm{\Gamma }(1n+d_2)}{ad_2!d_1!\mathrm{\Gamma }(1n)}\right)^{1/n}.$$
(5.2)
Note that by using the relation
$$\mathrm{\Gamma }(x)\mathrm{\Gamma }(1x)=\frac{\pi }{\mathrm{sin}(\pi x)}$$
(5.3)
we can write
$$\frac{\mathrm{\Gamma }(1n)}{\mathrm{\Gamma }(1n+d_2)}=\frac{\mathrm{\Gamma }(nd_2)}{\mathrm{\Gamma }(n)}\frac{\mathrm{sin}(\pi (nd_2))}{\mathrm{sin}(\pi n)}=\frac{(n1)!}{(nd_21)!}(1)^{d_2},$$
(5.4)
so that the $`\mathrm{\Gamma }`$ functions in the above expressions are in reality well defined.
The Liouville field and its derivative are given by:
$$\varphi \mathrm{log}\left(\frac{nad_2!d_1!}{n!(d_1n)!}\frac{(nd_21)!}{(n1)!}t^{n1}\right)+c.c.,_t\varphi \frac{n1}{t},$$
(5.5)
where we have dropped the factor $`(1)^{d_2}`$ in (5.4) since $`\varphi `$ is the real part of the logarithm. Substituting these values into the expression for the Liouville action:
$$S_L=\frac{i}{96\pi }𝑑t\varphi _t\varphi ,$$
(5.6)
we get a contribution from the point $`t=0`$:
$$S_L(t=0)=\frac{n1}{12}\mathrm{log}\left(nϵ^{\frac{n1}{n}}\right)\frac{n1}{12n}\mathrm{log}\left(a\frac{d_2!d_1!}{n!(d_1n)!}\frac{(nd_21)!}{(n1)!}\right),$$
(5.7)
where we note that the integration in (5.6) is performed along the circle
$$|t|=\left(\frac{ϵn!(d_1n)!(n1)!}{ad_2!d_1!(nd_21)!}\right)^{1/n}.$$
(5.8)
A simplification analogous to (5.4) will occur in many relations below, but for simplicity we leave the $`\mathrm{\Gamma }`$ functions in the form where they have negative arguments; we replace them with factorials of positive numbers only in the final expressions.
### 5.2 The contribution from $`z=a,t=1`$
Let us look at the point $`t=1`$. Using the expression for Wronskian (4.42), we find the derivative of the map (4.32):
$$\frac{dz}{dt}=a\frac{d_2!d_1!}{n!(d_1n)!}\frac{\mathrm{\Gamma }(1n)}{\mathrm{\Gamma }(1n+d_2)}\frac{nt^{n1}(1t)^{d_1+d_2n}}{\left[F(d_1,d_2;n+1;t)\right]^2},$$
(5.9)
which can be combined with known property of hypergeometric function:
$$F(a,b;c;1)=\frac{\mathrm{\Gamma }(c)\mathrm{\Gamma }(cab)}{\mathrm{\Gamma }(ca)\mathrm{\Gamma }(cb)}$$
(5.10)
to give the result:
$`z`$ $``$ $`a\beta a(1t)^{d_1+d_2n+1},`$ (5.11)
$`\beta `$ $`=`$ $`{\displaystyle \frac{n}{d_1+d_2n+1}}{\displaystyle \frac{d_1!d_2!(d_1n)!}{n!\left[(d_1+d_2n)!\right]^2}}{\displaystyle \frac{\mathrm{\Gamma }(1n+d_2)}{\mathrm{\Gamma }(1n)}}.`$ (5.12)
Our usual analysis gives:
$`z`$ $``$ $`a\beta a(1t)^{d_1+d_2n+1},1t\left({\displaystyle \frac{za}{a\beta }}\right)^{\frac{1}{d_1+d_2n+1}},`$
$`\varphi `$ $``$ $`\mathrm{log}\left(a\beta (d_1+d_2n+1)(1t)^{d_1+d_2n}\right)+c.c.,_t\varphi {\displaystyle \frac{d_1+d_2n}{1t}},`$
$`S_L(t=1)`$ $`=`$ $`{\displaystyle \frac{d_1+d_2n}{12}}\mathrm{log}(d_1+d_2n+1)`$
$``$ $`{\displaystyle \frac{(d_1+d_2n)^2}{12(d_1+d_2n+1)}}\mathrm{log}ϵ{\displaystyle \frac{(d_1+d_2n)}{12(d_1+d_2n+1)}}\mathrm{log}(a|\beta |).`$
Note that $`d_1+d_2n+1=m`$, so that we can rewrite the contribution from $`t=1`$ in a way which makes it look more symmetrical with the contribution from $`t=0`$. But we will defer all such simplifications to the final expressions for the fusion coefficients.
### 5.3 The contribution from $`z=\mathrm{}`$
To analyze the contribution from the point $`t=\mathrm{}`$ it is convenient to look at the map written in terms of Jacobi polynomials (4.40). Then one can use the Rodrigues’ formula to represent the Jacobi polynomials in the form:
$$P_k^{(\alpha \beta )}(x)=2^k\underset{j=0}{\overset{k}{}}\left(\begin{array}{c}k+\alpha \\ j\end{array}\right)\left(\begin{array}{c}k+\beta \\ kj\end{array}\right)(x1)^{kj}(x+1)^j.$$
(5.14)
The limit $`x\mathrm{}`$ gives:
$$P_k^{(\alpha \beta )}(x)x^k2^k\underset{j=0}{\overset{k}{}}\left(\begin{array}{c}k+\alpha \\ j\end{array}\right)\left(\begin{array}{c}k+\beta \\ kj\end{array}\right)=x^k2^k\left(\begin{array}{c}2k+\alpha +\beta \\ k\end{array}\right).$$
(5.15)
Substitution of this limit into the expression (4.40) gives the behavior near $`t=\mathrm{}`$ :
$`z`$ $``$ $`a\gamma (1)^{d_1d_2n}t^{d_1d_2},t\left((1)^{d_1d_2n}{\displaystyle \frac{z}{a\gamma }}\right)^{\frac{1}{d_1d_2}}`$
$`\gamma `$ $`=`$ $`{\displaystyle \frac{d_2!(d_1d_21)!}{(d_1n)!(nd_21)!}}{\displaystyle \frac{\mathrm{\Gamma }(d_1)}{\mathrm{\Gamma }(d_2d_1)}},`$ (5.16)
$`\varphi `$ $``$ $`\mathrm{log}\left(a\gamma (d_1d_2)t^{d_1d_21}\right)+c.c.,_t\varphi {\displaystyle \frac{d_1d_21}{t}}.`$
Consider first the point $`z=\mathrm{},t=\mathrm{}`$. Recall that we have taken the ‘regular region’ on $`\mathrm{\Sigma }`$ to be bounded by the image of $`1/\stackrel{~}{\delta }`$ (rather than $`1/\delta `$) when a twist operator is inserted. The contour around the puncture at infinity in the $`t`$ plane should be taken to go clockwise rather than anti-clockwise, so that it looks like a normal anti-clockwise contour in the local coordinate $`t^{}=1/t`$ around the puncture. Thus to compute the contribution from this puncture we should follow our usual procedure but reverse the overall sign. The result reads:
$`S_L(t=\mathrm{})`$ $`=`$ $`(1)[{\displaystyle \frac{d_1d_21}{12}}\mathrm{log}(d_1d_2)`$
$`+`$ $`{\displaystyle \frac{(d_1d_21)^2}{12(d_1d_2)}}\mathrm{log}\stackrel{~}{\delta }{\displaystyle \frac{d_1d_21}{12(d_1d_2)}}\mathrm{log}(a|\gamma |)].`$
Finally let us analyze the images of $`z=\mathrm{}`$ that give finite values $`t_i`$ of $`t`$. At each of these points the map $`tz`$ is one–to–one, in contrast to the above case $`z=\mathrm{},t=\mathrm{}`$ where $`q`$ values of $`t`$ correspond to each value of $`z`$ in a neighborhood of the puncture. Further, there is no sign reversal for the contour of integration around these punctures when we use the coordinate $`t`$ to describe the contour.
Looking at the structure of the map (4.40) one can easily identify the locations of the $`t_i`$: they coincide with zeroes of the denominator. So to evaluate the contribution to the Liouville action from the $`t_i`$ we will need some information about zeroes of Jacobi polynomials. Using the fact that Jacobi polynomials have only simple zeroes we can expand the map (4.40) around any of the $`t_i`$:
$`z`$ $``$ $`at_i^n{\displaystyle \frac{P_{d_1n}^{(n,d_1d_2+n1)}(12t_i)}{P_{}^{}{}_{d_2}{}^{(n,d_1d_2+n1)}(12t_i)}}{\displaystyle \frac{1}{tt_i}}{\displaystyle \frac{a\xi _i}{tt_i}},`$ (5.18)
$`\xi _i`$ $`=`$ $`t_i^n{\displaystyle \frac{P_{d_1n}^{(n,d_1d_2+n1)}(12t_i)}{\frac{d}{dt}P_{d_2}^{(n,d_1d_2+n1)}(12t_i)}}.`$ (5.19)
Then everything can be evaluated in terms of $`\xi _i`$:
$`tt_i`$ $``$ $`{\displaystyle \frac{a\xi _i}{z}},\varphi \mathrm{log}\left({\displaystyle \frac{a\xi _i}{(tt_i)^2}}\right)+c.c.,_t\varphi {\displaystyle \frac{2}{tt_i}}`$
$`S_L(t=t_i)`$ $`=`$ $`{\displaystyle \frac{1}{6}}\mathrm{log}(\delta ^2a\xi _i).`$ (5.20)
Collecting the contributions from all the $`t_i`$ we get:
$$S_L(\text{all }t_i)=\frac{d_2}{6}\mathrm{log}(\delta ^2a)\frac{1}{6}\underset{i=1}{\overset{d_2}{}}\mathrm{log}(\xi _i),$$
(5.21)
and we only need to evaluate the product of $`\xi _i`$. Note that the regularization parameter $`\delta `$ we use here has the same meaning as one considered in section 3.
This product can be written in terms of the Wronskian (4.41) and the discriminant of Jacobi polynomials. To see this we first rewrite (5.19) in terms of zeroes of Jacobi polynomials. If $`z=P/Q`$, then the Wronskian (4.41) is $`\stackrel{~}{W}=P^{}QPQ^{}=PQ^{}`$ at a zero of $`Q`$. Writing any of the $`\xi _i`$ as $`\xi =P/Q^{}=PQ^{}/Q^2`$ we find
$$\xi _i=\stackrel{~}{W}(t_i)\left[2a_0\underset{ji}{}(x_ix_j)\right]^2,$$
(5.22)
where $`a_0`$ is the coefficient in front of the highest power in the polynomial $`P_{d_2}^{(n,d_1d_2+n1)}`$; it can be evaluated using (5.15). The $`x_i`$ are the zeros of the polynomial $`Q(x)`$ in the denominator.
Applying the general definition of the discriminant to Jacobi polynomials
$$D_{d_2}^{(n,d_1d_2+n1)}a_0^{2d_22}\underset{i<j}{}(x_ix_j)^2,$$
(5.23)
we get
$$\underset{i=1}{\overset{d_2}{}}\xi _i=(1)^{d_2}2^{2d_2}a_0^{2(d_22)}\left[D_{d_2}^{(n,d_1d_2+n1)}\right]^2\underset{i=1}{\overset{d_2}{}}\stackrel{~}{W}(t_i).$$
(5.24)
The discriminant of Jacobi polynomials can be evaluated :
$`𝒟`$ $``$ $`D_{d_2}^{(n,d_1d_2+n1)}=2^{d_2(d_21)}`$
$`\times `$ $`{\displaystyle \underset{j=1}{\overset{d_2}{}}}j^{j+22d_2}(jn)^{j1}(jd_1d_2+n1)^{j1}(jd_11)^{d_2j}.`$
To evaluate the right hand side of (5.24) we only need the expressions for
$$\underset{i=1}{\overset{d_2}{}}t_i\text{and}\underset{i=1}{\overset{d_2}{}}(1t_i).$$
(5.26)
Let us consider the general Jacobi polynomial:
$$P_k^{(\alpha \beta )}(12t)=(2)^ka_0t^k+\mathrm{}+a_{k+1}=(2)^kb_0(t1)^k+\mathrm{}+b_{k+1},$$
(5.27)
Obviously $`b_0=a_0`$. By taking the limits $`t\mathrm{}`$, $`t0`$ and $`t1`$ in the above expression we find:
$`{\displaystyle \underset{i=1}{\overset{d_2}{}}}t_i`$ $`=`$ $`{\displaystyle \frac{a_{k+1}}{2^ka_0}}={\displaystyle \frac{\mathrm{\Gamma }(k+\alpha +1)\mathrm{\Gamma }(k+\alpha +\beta +1)}{\mathrm{\Gamma }(\alpha +1)\mathrm{\Gamma }(2k+\alpha +\beta +1)}},`$ (5.28)
$`{\displaystyle \underset{i=1}{\overset{d_2}{}}}(1t_i)`$ $`=`$ $`{\displaystyle \frac{b_{k+1}}{2^kb_0}}={\displaystyle \frac{\mathrm{\Gamma }(k+\beta +1)\mathrm{\Gamma }(k+\alpha +\beta +1)}{\mathrm{\Gamma }(\beta +1)\mathrm{\Gamma }(2k+\alpha +\beta +1)}}.`$ (5.29)
Collecting all contributions together, we get
$`\mathrm{log}{\displaystyle \underset{i=1}{\overset{d_2}{}}}\xi _i`$ $`=`$ $`2d_2(d_21)\mathrm{log}2+d_2\mathrm{log}n2\mathrm{log}𝒟(3d_24)\mathrm{log}d_2!`$
$`+`$ $`d_2\mathrm{log}\left[{\displaystyle \frac{d_1!}{n!(d_1n)!}}\right]+(n+d_21)\mathrm{log}{\displaystyle \frac{(n1)!}{(nd_21)!}}`$
$`+`$ $`(d_1d_2+3)\mathrm{log}{\displaystyle \frac{(d_1d_2)!}{d_1!}}+(d_1+d_2n)\mathrm{log}{\displaystyle \frac{(d_1+d_2n)!}{(d_1n)!}}.`$
### 5.4 The total Liouville action.
Collecting the contributions from the different branching points we obtain the final expression for the Liouville action
$`S_L^{(1)}`$ $`=`$ $`\left({\displaystyle \frac{(n1)^2}{12n}}+{\displaystyle \frac{(d_1+d_2n)^2}{12(d_1+d_2n+1)}}\right)\mathrm{log}ϵ{\displaystyle \frac{(d_1d_21)^2}{12(d_1d_2)}}\mathrm{log}\stackrel{~}{\delta }`$
$``$ $`{\displaystyle \frac{d_2}{3}}\mathrm{log}\delta \left({\displaystyle \frac{n1}{12n}}+{\displaystyle \frac{d_1+d_2n}{12(d_1+d_2n+1)}}{\displaystyle \frac{d_1d_21}{12(d_1d_2)}}+{\displaystyle \frac{d_2}{6}}\right)\mathrm{log}a`$
$``$ $`{\displaystyle \frac{n1}{12}}\mathrm{log}n{\displaystyle \frac{d_1+d_2n}{12}}\mathrm{log}(d_1+d_2n+1)`$
$`+`$ $`{\displaystyle \frac{d_1d_21}{12}}\mathrm{log}(d_1d_2){\displaystyle \frac{n1}{12n}}\mathrm{log}\left({\displaystyle \frac{d_1!d_2!}{n!(d_1n)!}}{\displaystyle \frac{\mathrm{\Gamma }(1n)}{\mathrm{\Gamma }(1n+d_2)}}\right)`$
$``$ $`{\displaystyle \frac{d_1+d_2n}{12(d_1+d_2n+1)}}\mathrm{log}|\beta |+{\displaystyle \frac{d_1d_21}{12(d_1d_2)}}\mathrm{log}|\gamma |{\displaystyle \frac{1}{6}}\mathrm{log}\left({\displaystyle \underset{i=1}{\overset{d_2}{}}}\xi _i\right).`$
The values of $`\beta `$ and $`\gamma `$ are given by (5.12) and (5.16), and the last term is given through (LABEL:ProdXi). According to (2.44), the three point function is given by:
$$\sigma _n^ϵ(0)\sigma _m^ϵ(a)\sigma _q^{\stackrel{~}{\delta }}(\mathrm{})_\delta =e^{S_L^{(1)}}e^{S_L^{(2)}+S_L^{(3)}}\frac{Z^{(\widehat{s})}}{(Z_\delta )^s},$$
(5.32)
where $`s`$ is number of fields involved in permutation; it is defined by (4.8). Note that we have not determined the values of $`S_L^{(2)}`$ and $`S_L^{(3)}`$ for the case under consideration, as we will see these quantities will cancel in the final answer.
## 6 Normalizing the twist operators
This Liouville action (5.4) yields the correlation function for twist operators with the regularization parameters $`ϵ,\delta ,\stackrel{~}{\delta }`$. We immediately see that the power of $`a`$ in the correlator is
$`a:`$ $`\left({\displaystyle \frac{n1}{12n}}+{\displaystyle \frac{d_1+d_2n}{12(d_1+d_2n+1)}}{\displaystyle \frac{d_1d_21}{12(d_1d_2)}}+{\displaystyle \frac{d_2}{6}}\right)`$ (6.1)
$`+{\displaystyle \frac{1}{2}}\left(M_a^n+M_a^mM_a^q\right)={\displaystyle \frac{1}{6}}\left(n{\displaystyle \frac{1}{n}}+m{\displaystyle \frac{1}{m}}q+{\displaystyle \frac{1}{q}}\right),`$
which agrees with the expected $`a`$ dependence of the 3-point function
$$\sigma _n^ϵ(0)\sigma _m^ϵ(a)\sigma _q^{\stackrel{~}{\delta }}(\mathrm{})|a|^{2(\mathrm{\Delta }_m+\mathrm{\Delta }_n\mathrm{\Delta }_q)}.$$
(6.2)
To obtain the final correlation functions and fusion coefficients we have two sources of renormalization coefficients that need to be considered:
(a) We have to normalize the operators $`\sigma _n^ϵ`$ such that their 2-point functions are set to unity at unit $`z`$ separation; at this point we should find that the parameters $`ϵ,\delta ,\stackrel{~}{\delta }`$ disappear from the 3-point (and higher point) functions as well. After we normalize the twist operators $`\sigma _n^ϵ`$ in this way we will call them $`\sigma _n`$.
(b) The CFT had $`N`$ fields $`X^I`$, though only $`n`$ of them are affected by the twist operator $`\sigma _n^ϵ`$. However at the end of the calculation of any correlation function of the operators $`\sigma _{n_i}^ϵ`$ we must sum over all the possible ways that the $`n_i`$ fields that are twisted can be chosen from the total set of $`N`$ fields. Thus we will have to define operators $`O_n`$ that are sums over conjugacy classes of the permutation group, and these operators $`O_n`$ are the only ones that will finally be well defined operators in the CFT . The correctly normalized $`O_n`$ will thus have combinatoric factors multiplying the normalized operators $`\sigma _n`$.
We choose to arrive at the final normalized operators $`O_n`$ in these two steps since the calculations involved in steps (a) and (b) are quite different; further when $`n_i<<N`$ the factor coming from (b) is just a power of $`N`$ which is easily found.
### 6.1 Normalizing the $`\sigma _n^ϵ`$
Let us define the normalized twist operators $`\sigma _n`$ by requiring
$$\sigma _n(0)\sigma _n(a)=\frac{1}{|a|^{4\mathrm{\Delta }_n}}.$$
(6.3)
From (3.18) we see that
$$\sigma _n=D_n\sigma _n^ϵ,D_n=[C_nϵ^{A_n}Q^{B_n}]^{1/2}.$$
(6.4)
Let the Operator Product Expansion (OPE) have the form
$$\sigma _m(a)\sigma _n(0)\frac{|C_{nmq}^\sigma |^2}{|a|^{2(\mathrm{\Delta }_n+\mathrm{\Delta }_m\mathrm{\Delta }_q)}}\sigma _q(0)+\mathrm{},$$
(6.5)
where we have written the OPE for holomorphic and anti-holomorphic blocks combined, and we have put a superscript $`\sigma `$ on the fusion coefficients $`C_{nmq}^\sigma `$ to remind ourselves that these are not the final fusion coefficients of the physical operators $`O_n`$. With the normalization (6.3) we will get
$`\sigma _n(z_1)\sigma _m(z_2)\sigma _q(z_3)=`$
$`{\displaystyle \frac{|C_{nmq}^\sigma |^2}{|z_1z_2|^{2(\mathrm{\Delta }_n+\mathrm{\Delta }_m\mathrm{\Delta }_q)}|z_2z_3|^{2(\mathrm{\Delta }_m+\mathrm{\Delta }_q\mathrm{\Delta }_n)}|z_3z_1|^{2(\mathrm{\Delta }_q+\mathrm{\Delta }_n\mathrm{\Delta }_m)}}}`$ (6.6)
We have computed the 3-point functions and should thus be able to get the fusion coefficients $`C_{nmq}^\sigma `$ from (6.6). However while two of our twist operators were inserted at finite points in the $`z`$ plane, the last one was inserted at infinity. Putting one of the points at infinity simplified the calculation, but it also creates the following problem: unlike the twist operators at $`z=0,z=a`$ which are normalized through (6.4) it is not clear what is the normalization of the twist operator that is inserted at infinity. (We could think of this operator as inserted at a puncture on the sphere at infinity, and therefore no different from the other insertions, but we have chosen the flat metric on the $`z`$ plane and thus made infinity a special region carrying curvature.) To get around this problem we adopt the following scheme. If we have the OPE (6.5) then we will get
$$\frac{\sigma _n(0)\sigma _m(a)\sigma _q^{\stackrel{~}{\delta }}(\mathrm{})}{\sigma _q(0)\sigma _q^{\stackrel{~}{\delta }}(\mathrm{})}=\frac{|C_{nmq}^\sigma |^2}{|a|^{2(\mathrm{\Delta }_n+\mathrm{\Delta }_m\mathrm{\Delta }_q)}}$$
(6.7)
and we will thus not need to know the normalization of the operator at infinity. $`\sigma _n(0)\sigma _m(a)\sigma _q^{\stackrel{~}{\delta }}(\mathrm{})`$ can be found from our 3-point function calculation together with the normalization factors for $`\sigma _n,\sigma _m`$ from (6.3). To compute the denominator we must find the 2-point function with one operator at infinity. (We had earlier computed the 2-point function with both operators in the finite $`z`$ plane since if we put one operator at infinity then we loose any position dependence in the correlator and cannot extract the scaling dimensions.)
To evaluate the two point function
$$\sigma _n(0)\sigma _n(\mathrm{})$$
(6.8)
we consider the map:
$$z=bt^n.$$
(6.9)
This map has order $`n`$ ramification points at $`t=0`$ and $`t=\mathrm{}`$ and the usual calculations give:
$`\varphi =\mathrm{log}[nbt^{n1}]+c.c.,_t\varphi ={\displaystyle \frac{n1}{t}},`$
$`S_L(t=0)={\displaystyle \frac{n1}{12}}\mathrm{log}[nb^{1/n}ϵ^{(n1)/n}],`$ (6.10)
$`S_L(t=\mathrm{})=(1){\displaystyle \frac{n1}{12}}\mathrm{log}[nb^{1/n}\stackrel{~}{\delta }^{(1n)/n}].`$ (6.11)
As before we cut a hole of size $`ϵ`$ around the origin and put the twist at infinity on a boundary at $`|z|=1/\stackrel{~}{\delta }`$. We also have an extra negative sign for the cut at infinity since the contour that goes anti-clockwise in the local coordinate at $`1/t`$ near $`t=\mathrm{}`$ goes clockwise in the coordinate $`t`$. Collecting both contributions we get:
$`\sigma _n^ϵ(0)\sigma _n^{\stackrel{~}{\delta }}(\mathrm{})`$ $`=`$ $`ϵ^{F_n}\stackrel{~}{\delta }^{F_n}e^{S_L^{(2)}+S_L^{(3)}}{\displaystyle \frac{Z^{(\widehat{s})}}{(Z_\delta )^n}},`$ (6.12)
$`F_n`$ $`=`$ $`{\displaystyle \frac{(n1)^2}{12n}}.`$ (6.13)
The correlator is $`b`$–independent as expected. The values of $`S_L^{(2)}`$, $`S_L^{(3)}`$ and partition function $`Z^{(\widehat{s})}`$ depend on $`\delta ,\delta ^{},\stackrel{~}{\delta }`$, but these expressions are are the same as in (5.32) and so they cancel in the final answer.
The $`C_{n,m,q}^\sigma `$ are then given by
$`|C_{n,m,q}^\sigma |^2`$ $`=`$ $`{\displaystyle \frac{\sigma _n(0)\sigma _m(a)\sigma _q^{\stackrel{~}{\delta }}(\mathrm{})}{\sigma _q(0)\sigma _q^{\stackrel{~}{\delta }}(\mathrm{})}}`$
$`=`$ $`{\displaystyle \frac{\sigma _n^ϵ(0)\sigma _m^ϵ(a)\sigma _q^{\stackrel{~}{\delta }}(\mathrm{})}{\sigma _q^ϵ(0)\sigma _q^{\stackrel{~}{\delta }}(\mathrm{})}}\sqrt{{\displaystyle \frac{\sigma _q^ϵ(0)\sigma _q^ϵ(1)}{\sigma _n^ϵ(0)\sigma _n^ϵ(1)\sigma _m^ϵ(0)\sigma _m^ϵ(1)}}}.`$
As a consistency check of our procedure we look at powers of various regularization parameters: they should cancel in any physical quantity. For $`|C_{n,m,q}^\sigma |^2`$ we have the following powers for $`ϵ,\delta ,\stackrel{~}{\delta }`$, $`Q`$:
$`ϵ:`$ $`\left({\displaystyle \frac{(n1)^2}{12n}}+{\displaystyle \frac{(m1)^2}{12m}}\right)F_q{\displaystyle \frac{1}{2}}\left(A_n+A_mA_q\right)=0,`$ (6.15)
$`\stackrel{~}{\delta }:`$ $`{\displaystyle \frac{(q1)^2}{12q}}F_q=0,`$ (6.16)
$`\delta :`$ $`{\displaystyle \frac{d_2}{3}}s+q=0,`$ (6.17)
$`Q:`$ $`1s{\displaystyle \frac{1}{2}}\left(B_n+B_m+B_q\right)=0.`$ (6.18)
We used the expressions (4.13) and (4.14) for $`d_1`$ and $`d_2`$, the values of $`A_n,B_n,F_n`$ from (3.19) and (6.13) and the genus relation (4.9).
We finally get (for the contribution from $`\mathrm{\Sigma }`$ of genus zero) the logarithm of the fusion coefficient
$`\mathrm{log}|C_{n,m,q}^\sigma |^2`$ $`=`$ $`{\displaystyle \frac{1}{6}}\mathrm{log}\left({\displaystyle \frac{q}{mn}}\right){\displaystyle \frac{n1}{12}}\mathrm{log}n{\displaystyle \frac{m1}{12}}\mathrm{log}m+{\displaystyle \frac{q1}{12}}\mathrm{log}(q)`$
$``$ $`{\displaystyle \frac{n1}{12n}}\mathrm{log}\left({\displaystyle \frac{d_1!d_2!}{n!(n1)!}}{\displaystyle \frac{(d_1m)!}{(d_1n)!}}\right)`$
$``$ $`{\displaystyle \frac{m1}{12m}}\mathrm{log}\left({\displaystyle \frac{d_1!d_2!}{m!(m1)!}}{\displaystyle \frac{(d_1n)!}{(d_1m)!}}\right)`$
$`+`$ $`{\displaystyle \frac{q1}{12q}}\mathrm{log}\left({\displaystyle \frac{(q1)!d_2!}{(d_1n)!(d_1m)!}}{\displaystyle \frac{(d_1d_2)!)}{d_1!}}\right){\displaystyle \frac{1}{6}}\mathrm{log}\left({\displaystyle \underset{i=1}{\overset{d_2}{}}}\xi _i\right).`$
The expression for the product of $`\xi _i`$ is given by (LABEL:ProdXi).
The coefficients $`C_{n,m,q}^\sigma `$ must be symmetric in the indices $`m,n,q`$. We have written (6.1) in such a way that all the terms except the last one show a manifest symmetry between $`m`$ and $`n`$. It can be shown without much difficulty that the last term (given through (LABEL:ProdXi) and (5.3)) is also symmetric in $`m`$ and $`n`$. In particular $`d_1`$ and $`d_2`$ are symmetric in $`m`$ and $`n`$, and $`d_1d_2+n1=m`$, so that the Jacobi polynomial whose discriminant is calculated in (5.3) is $`P_{d_2}^{n,m}`$.
On the other hand it is not at all obvious that the expression (LABEL:ProdXi) is symmetric under the interchange of $`q`$ with either $`n`$ or $`m`$. Note that $`d_2+1`$ is the number of elements that overlap between the permutations $`\sigma _n`$ and $`\sigma _m`$, and the product in (5.3) runs over the range $`j=1\mathrm{}d_2`$. This number $`d_2+1`$ is in general different from the number of overlapping elements between the permutations $`\sigma _q`$ and $`\sigma _n`$ or between $`\sigma _q`$ and $`\sigma _m`$, and thus there is no simple way to write (LABEL:ProdXi) in form that makes its total symmetry manifest. Nevertheless, this expression is indeed symmetric in all three arguments $`n,m,q`$, as can be checked by evaluating the expression through a symbolic manipulation program. Verifying this symmetry provides a useful check of all our calculations for the 3-point function.
### 6.2 Two special cases
Due to the structure of the discriminant (5.3), the general expression for the fusion coefficient looks complicated for an arbitrary value of $`d_2`$. However there are two important cases where significant simplifications occur. These are the cases of one and two overlaps (we recall that the number of common indices in $`\sigma _n`$ and $`\sigma _m`$ is $`d_2+1`$). One can see that for $`d_2=0`$ and $`d_2=1`$ the discriminant $`𝒟=1`$. Let us analyze both these cases.
For one overlap we have:
$$d_2=0,d_1=q=m+n1,$$
(6.20)
and the logarithm of fusion coefficient is given by:
$`\mathrm{log}|C_{n,m,m+n1}^\sigma |^2={\displaystyle \frac{1}{12}}\left(n+{\displaystyle \frac{1}{n}}\right)\mathrm{log}n{\displaystyle \frac{1}{12}}\left(m+{\displaystyle \frac{1}{m}}\right)\mathrm{log}m`$ (6.21)
$`+{\displaystyle \frac{1}{12}}\left(q+{\displaystyle \frac{1}{n}}+{\displaystyle \frac{1}{m}}1\right)\mathrm{log}q{\displaystyle \frac{1}{12}}\left(1+{\displaystyle \frac{1}{q}}{\displaystyle \frac{1}{n}}{\displaystyle \frac{1}{m}}\right)\mathrm{log}\left({\displaystyle \frac{(q1)!}{(m1)!(n1)!}}\right).`$
In particular we get:
$$|C_{223}^\sigma |^2=2^{\frac{4}{9}}3^{\frac{1}{4}}.$$
(6.22)
The case of two overlaps corresponds to
$$d_2=1,q=m+n3,d_1=q+1,$$
(6.23)
and the result reads:
$`\mathrm{log}|C_{n,m,m+n3}^\sigma |^2={\displaystyle \frac{1}{12}}\left({\displaystyle \frac{1}{n}}+{\displaystyle \frac{1}{m}}{\displaystyle \frac{1}{m+n3}}3\right)\mathrm{log}\left({\displaystyle \frac{(m+n3)!}{(m1)!(n1)!}}\right)`$
$`{\displaystyle \frac{n^2+1}{12n}}\mathrm{log}n{\displaystyle \frac{m^2+1}{12m}}\mathrm{log}m+{\displaystyle \frac{1}{12}}\left(2+{\displaystyle \frac{(m+n4)^2}{m+n3}}\right)\mathrm{log}(m+n3)+`$
$$\frac{1}{12}\left(\frac{nm}{mn}2n+\frac{m+n4}{m+n3}\right)\mathrm{log}(n1)+\frac{1}{12}\left(\frac{mn}{mn}2m+\frac{m+n4}{m+n3}\right)\mathrm{log}(m1)$$
$$+\frac{1}{12}\left(2(m+n)5+\frac{1}{n}+\frac{1}{m}+\frac{1}{m+n3}\right)\mathrm{log}(m+n2).$$
(6.24)
In particular for $`m=3`$ we get:
$`\mathrm{log}|C_{n3n}^\sigma |^2`$ $`=`$ $`{\displaystyle \frac{2}{9}}\mathrm{log}n{\displaystyle \frac{1}{6}}\left(n+{\displaystyle \frac{1}{n}}{\displaystyle \frac{2}{3}}\right)\mathrm{log}(n1)`$ (6.25)
$`+`$ $`{\displaystyle \frac{1}{6}}\left(n+{\displaystyle \frac{1}{n}}+{\displaystyle \frac{2}{3}}\right)\mathrm{log}(n+1){\displaystyle \frac{2}{9}}\mathrm{log}2{\displaystyle \frac{5}{18}}\mathrm{log}3.`$
From this expression we can extract the value of $`C_{232}^\sigma `$ and check that it equals the value of $`C_{223}^\sigma `$ given by (6.22).
### 6.3 Combinatoric factors and large N limit.
The twist operators we have considered so far do not represent proper fields in the conformal field theory. In the orbifold CFT there is one twist field for each conjugacy class of the permutation group, not for each element of the group . The true CFT operators that represent the twist fields can be constructed by summing over the group orbit:
$$O_n=\frac{\lambda _n}{N!}\underset{hG}{}\sigma _{h(1\mathrm{}n)h^1}.$$
(6.26)
Here $`G`$ is the permutation group $`S_N`$ and the normalization constant $`\lambda _n`$ will be determined below.
Using normalization condition for the $`\sigma `$ operators:
$$\sigma _n(0)\sigma _n(1)=1$$
(6.27)
we find:
$$O_n(0)O_n(1)=\frac{\lambda _n^2}{N!}\sigma _{(1\mathrm{}n)}\sigma _{h(1\mathrm{}n)h^1}=\lambda _n^2n\frac{(Nn)!}{N!}\sigma _n(0)\sigma _n(1).$$
(6.28)
Requiring the normalisation $`O_n(0)O_n(1)=1`$ we find the value of $`\lambda _n`$:
$$\lambda _n=\left[\frac{n(Nn)!}{N!}\right]^{1/2}.$$
(6.29)
Let us now look at the three point function. First we consider the combinatorics for the $`g=0`$ cases that we worked with above; the permutation structure was described in (4.4), (4.5), (4.6). Simple combinatorics yields
$`O_n(0)O_m(1)O_q(z)=({\displaystyle \frac{1}{N!}})^3\lambda _n\lambda _m\lambda _q`$
$`\times nmq{\displaystyle \frac{N!}{(Ns)!}}(Nn)!(Nm)!(Nq)!\sigma _n(0)\sigma _m(1)\sigma _q(z).`$ (6.30)
One way of getting this expression is to note that $`s`$ different indices are involved in the permutation, and we can select these indices, in the order in which they appear when the permutations are written out, in $`N!/(Ns)!`$ ways. Having obtained the indices for any given permutation, we ask how many elements out of the sum over group elements yields this set of indices in the permutation; the answer for $`\sigma _n`$ for example is $`(Nn)!`$, since only the permutations of the remaining $`Nn`$ elements leave the indices in $`\sigma _n`$ untouched. Finally, we note that any permutation $`\sigma _k`$ can be written in $`k`$ equivalent ways since we can begin the set of indices with any index that we choose from the set; this leads to the factors $`nmq`$.
Substituting the values of $`\lambda _i`$ we get the final result:
$$O_n(0)O_m(1)O_q(z)=\frac{\sqrt{mnq(Nn)!(Nm)!(Nq)!}}{(Ns)!\sqrt{N!}}\sigma _n(0)\sigma _m(1)\sigma _q(z).$$
with $`s=\frac{1}{2}(n+m+q1)`$.
Now we analyse the behavior of the combinatoric factors for arbitrary genus $`g`$ but in the limit where $`N`$ is taken to be large while the orders of twist operators ($`m`$, $`n`$ and $`q`$) as well as the parameter $`g`$ are kept fixed. There are $`s`$ different fields $`X^i`$ involved in the 3-point function, and these fields can be selected in $`N^s`$ ways. Similarily the 2-point function of $`\sigma _n`$ will go as $`N^n`$ since $`n`$ different fields are to be selected. Thus the 3-point function of normalised twist operators will behave as
$$N^{s\frac{n+m+q}{2}}=N^{(g+\frac{1}{2})}$$
(6.31)
(which can also be obtained from (6.3)).
Thus in the large $`N`$ limit the contributions from surfaces with high genus will be suppressed, and the leading order the answer can be obtained by considering only contributions from the sphere ($`g=0`$). This is presicely the case that we have analysed in detail, and knowing the amplitude $`\sigma _n(0)\sigma _m(1)\sigma _q(z)`$ one can easily extract the leading order of the CFT correlation function:
$$O_n(0)O_m(1)O_q(z)=\sqrt{\frac{1}{N}}\sqrt{mnq}\sigma _n(0)\sigma _m(1)\sigma _q(z)_{sphere}+O\left(\frac{1}{N^{3/2}}\right).$$
(6.32)
## 7 Four Point Function.
In this section we compute specific examples of 4-point functions, without attempting to analyze the most general case. The computations illustrate interesting features which arise in our approach for four and higher point functions. In particular we will also need to compute a genus one correlation function. We will also be able to verify specific examples of the fusion coefficients computed in the last section as they will be recovered through factorization of the 4-point functions.
### 7.1 An example of a 4-point function on a sphere.
Let us start with a map that has branch points appropriate for a 4-point correlation function of the form
$$\sigma _n(0)\sigma _2(1)\sigma _2(w)\sigma _n(\mathrm{})$$
(7.1)
Consider the map
$$z=Ct^n\frac{ta}{t1},$$
(7.2)
where the parameter $`a`$ will be related with coordinate $`w`$ and the value of coefficient $`C`$ will be determined below. The map (7.2) has two obvious ramification points: $`z=0`$ and $`z=\mathrm{}`$, both of them give an $`n`$–th order branch point for nonzero values of $`a`$. For a general value of $`a`$ the map (7.2) has two more ramification points; to find them we should look at the equation
$$\frac{dz}{dt}=0.$$
(7.3)
For general value of $`a`$ this equation reads:
$$t^{n1}(nt^2t((n1)a+(n+1))+an)=0.$$
(7.4)
The first factor corresponds to the obvious fact that at the point $`t=0`$ we have a ramification point of $`n`$–th order, while the positions of the two “implicit” points of second order are given by:
$$t_\pm =\frac{1}{2n}\left((n1)a+n+1\pm \sqrt{(a1)((n1)^2a(n+1)^2)}\right).$$
(7.5)
One of these points should correspond to $`z=1`$; we let this be the point $`t_+`$. The other must correspond to $`z=w`$. By requiring $`z(t_+)=1`$ we determine the value of coefficient $`C`$:
$$C=t_+^n\frac{t_+1}{t_+a},$$
(7.6)
and we note that in what follows we will have
$$w=t_{}.$$
(7.7)
Now we will analyze contributions to the Liouville action coming from the different ramification points. Let us start from the point $`t=0`$. If $`a0`$ the map (7.2) near this point has the form:
$$zCat^n$$
(7.8)
and the inverse map is
$$t\left(\frac{z}{aC}\right)^{1/n}.$$
(7.9)
The Liouville field and its derivative are given by:
$$\varphi =\mathrm{log}\left(\frac{dz}{dt}\right)+c.c.\mathrm{log}(nCat^{n1})+c.c.,_t\varphi \frac{n1}{t}.$$
(7.10)
As usual we will cut a hole of radius $`ϵ`$ around the point $`z=0`$. Then the contribution to the Liouville action coming from the integration over the boundary of this hole is
$$S_L(t=0)=\frac{n1}{12}\mathrm{log}\left[n(aC)^{1/n}ϵ^{\frac{n1}{n}}\right].$$
(7.11)
The same analysis near $`t=\mathrm{}`$ gives:
$`zCt^n,t\left({\displaystyle \frac{z}{C}}\right)^{1/n},`$
$`\varphi \mathrm{log}(nCt^{n1})+c.c.,_t\varphi {\displaystyle \frac{n1}{t}},`$
$`S_L(t=\mathrm{})=(1){\displaystyle \frac{n1}{12}}\mathrm{log}\left[nC^{1/n}\stackrel{~}{\delta }^{\frac{1n}{n}}\right].`$ (7.12)
Here we have cut a large circle of radius $`1/\stackrel{~}{\delta }`$ in $`z`$ plane and the factor of $`(1)`$ in the last equation comes from the fact that we go around the point $`t=\mathrm{}`$ clockwise.
Near the point $`t=t_{}`$ we get:
$`zz_{}+\xi _{}(tt_{})^2,tt_{}\left({\displaystyle \frac{z}{\xi _{}}}\right)^{1/2},`$
$`\xi _{}={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{d^2z}{dt^2}}\right)_{t=t_{}}{\displaystyle \frac{nz_{}(t_{}t_+)}{2t_{}(t_{}a)(t_{}1)}},`$
$`\varphi {\displaystyle \frac{1}{2}}\mathrm{log}(4(zz_{})\xi _{})+c.c.,_t\varphi {\displaystyle \frac{1}{t}},`$
$`S_L(t=t_{})={\displaystyle \frac{1}{24}}\mathrm{log}\left[4ϵ\xi _{}\right].`$ (7.13)
To get a contribution for $`t=t_+`$ one should make a replacement $`+`$ in the last expression; we also note that $`z_+=1`$. Thus we get:
$`S_L(t=t_+)={\displaystyle \frac{1}{24}}\mathrm{log}\left[4ϵ\xi _+\right],`$ (7.14)
$`\xi _+{\displaystyle \frac{n(t_+t_{})}{2t_+(t_+a)(t_+1)}}.`$
Finally we should consider the images of $`z=\mathrm{}`$ that give finite values for $`t`$ \- there will be a puncture here corresponding to the boundary of the $`|z|`$ plane. As before we let this circle in the $`z`$ plane have a radius $`1/\delta `$. The map (7.2) has only one such image: $`t=1`$. The result is
$`zC{\displaystyle \frac{1a}{t1}},t1{\displaystyle \frac{C(1a)}{z}},`$
$`\varphi \mathrm{log}({\displaystyle \frac{z^2}{C(1a)}})+c.c.,_t\varphi {\displaystyle \frac{2}{t1}}`$
$`S_L(t=1)={\displaystyle \frac{1}{6}}\mathrm{log}\left[C(1a)\delta ^2\right].`$ (7.15)
Collecting all this information together and making obvious simplifications we finally get the expression for the Liouville action corresponding to our four point function:
$`S_L^{(1)}`$ $`=`$ $`\left({\displaystyle \frac{(n1)^2}{12n}}{\displaystyle \frac{1}{12}}\right)\mathrm{log}ϵ{\displaystyle \frac{(n1)^2}{12n}}\mathrm{log}\stackrel{~}{\delta }{\displaystyle \frac{1}{3}}\mathrm{log}\delta `$ (7.16)
$``$ $`{\displaystyle \frac{1}{4}}\mathrm{log}C{\displaystyle \frac{1}{12}}\mathrm{log}2{\displaystyle \frac{1}{12}}\left({\displaystyle \frac{n+1}{2}}{\displaystyle \frac{1}{n}}\right)\mathrm{log}a{\displaystyle \frac{1}{8}}\mathrm{log}(1a)`$
$``$ $`{\displaystyle \frac{1}{24}}\mathrm{log}\left[(n1)^2a(n+1)^2\right]{\displaystyle \frac{1}{12}}\mathrm{log}n.`$
The above expression gives the correlator $`\sigma _n^ϵ(0)\sigma _2^ϵ(1)\sigma _2^ϵ(w)\sigma _n^{\stackrel{~}{\delta }}(\mathrm{})`$. To compute the correlation function of normalised twist operators, with one point at infinity, defined as
$`\sigma _n(0)\sigma _2(1)\sigma _2(w)\sigma _n(\mathrm{})`$ $``$ $`\underset{|z|\mathrm{}}{lim}|z|^{4\mathrm{\Delta }_n}\sigma _n(0)\sigma _2(1)\sigma _2(w)\sigma _n(z)`$ (7.17)
$`=`$ $`{\displaystyle \frac{\sigma _n(0)\sigma _2(1)\sigma _2(w)\sigma _n^{\stackrel{~}{\delta }}(\mathrm{})}{\sigma _n(0)\sigma _n^{\stackrel{~}{\delta }}(\mathrm{})}}`$
we use arguments similar to those in subsection 6.1. Then we find
$`_4`$ $``$ $`\sigma _n(0)\sigma _2(1)\sigma _2(w)\sigma _n(\mathrm{})`$ (7.18)
$`=`$ $`{\displaystyle \frac{\sigma _n^ϵ(0)\sigma _2^ϵ(1)\sigma _2^ϵ(w)\sigma _n^{\stackrel{~}{\delta }}(\mathrm{})}{\sigma _n^ϵ(0)\sigma _n^{\stackrel{~}{\delta }}(\mathrm{})}}[\sigma _2^ϵ(0)\sigma _2^ϵ(1)]^1`$
This leads to
$`\mathrm{log}_4`$ $`=`$ $`S_L^{(1)}\mathrm{log}\left(ϵ^{F_n}\stackrel{~}{\delta }^{F_n}\right)\mathrm{log}\left(C_2ϵ^{A_2}Q^{B_2}\right)+{\displaystyle \frac{1}{3}}(n+1n)\mathrm{log}\delta `$ (7.19)
$`+`$ $`((n1)n)\mathrm{log}Q.`$
where $`A_n,C_n,F_n`$ are given in (3.19) and (6.13). The sourse of the last two terms is the fact that the numerator has $`n+1`$ fields transforming nontrivially, while the denominator has only $`n`$:
$`\sigma _n^ϵ(0)\sigma _n^{\stackrel{~}{\delta }}(\mathrm{})`$ $`=`$ $`ϵ^{F_n}\stackrel{~}{\delta }^{F_n}e^{S_L^{(2)}+S_L^{(3)}}{\displaystyle \frac{Z^{(\widehat{s})}}{(Z_\delta )^n}},`$ (7.20)
$`\sigma _n^ϵ(0)\sigma _2^ϵ(1)\sigma _2^ϵ(w)\sigma _n^\delta (\mathrm{})`$ $`=`$ $`e^{S_L^{(1)}}e^{S_L^{(2)}+S_L^{(3)}}{\displaystyle \frac{Z^{(\widehat{s})}}{(Z_\delta )^{n+1}}}`$ (7.21)
We observe that the powers of regularization parameters cancel:
$`\mathrm{log}ϵ:`$ $`{\displaystyle \frac{(n1)^2}{12n}}{\displaystyle \frac{1}{12}}F_nA_2=0,`$
$`\mathrm{log}\stackrel{~}{\delta }:`$ $`{\displaystyle \frac{(n1)^2}{12n}}F_n=0,\mathrm{log}\delta :{\displaystyle \frac{1}{3}}+{\displaystyle \frac{1}{3}}=0,`$ (7.22)
$`\mathrm{log}Q:`$ $`n+(n1)+(21)=0.`$
This cancellation gives a consistency check on the calculations. The expression for the logarithm of the normalized four point function is given by:
$`\mathrm{log}_4`$ $`=`$ $`{\displaystyle \frac{1}{4}}\mathrm{log}C{\displaystyle \frac{1}{12}}\left({\displaystyle \frac{n+1}{2}}{\displaystyle \frac{1}{n}}\right)\mathrm{log}a{\displaystyle \frac{1}{8}}\mathrm{log}(1a)`$
$``$ $`{\displaystyle \frac{1}{24}}\mathrm{log}\left[(n1)^2a(n+1)^2\right]{\displaystyle \frac{1}{12}}\mathrm{log}n{\displaystyle \frac{5}{12}}\mathrm{log}2.`$ (7.23)
The value of $`C`$ is given by (7.6).
### 7.2 Analysis of the 4-point function
Let us step back from the above calculation and think about the structure of a 4-point function $`\sigma _n(0)\sigma _2(1)\sigma _2(w)\sigma _n(\mathrm{})`$. Consider the limit $`w0`$, and ask what operators are produced in the OPE of $`\sigma _n`$ and $`\sigma _2`$. There are three possibilities, which must all be considered when we make the CFT operators $`O_j`$ out of the sum over indices in the $`\sigma _j`$:
(a) The indices of $`\sigma _2`$ and the indices of $`\sigma _n`$ have no overlap - i.e., the operators are of the form $`\sigma _{12}\sigma _{34\mathrm{}n+2}`$. In this case the other two operators must have no overlapping indices either, and the entire 4-point function factors into two different parts $`\sigma _2\sigma _2\sigma _n\sigma _n`$. The covering surfaces are separate for the two parts, and we just multiply together the correlation functions obtained from the covering surfaces for the 2-point functions.
(b) The indices of $`\sigma _2`$ and the indices of $`\sigma _n`$ have one overlap - i.e., the operators are of the form $`\sigma _{12}\sigma _{23\mathrm{}n+1}`$. The OPE then produces the operator $`\sigma _{123\mathrm{}n+1}=\sigma _{n+1}`$. The other two operators in the correlator must also have a singe overlap so that they can produce $`\sigma _{n+1}`$. The genus of the surface thus produced is seen to be
$$g=\frac{1+1+(n1)+(n1)}{2}(n+1)+1=0$$
(7.24)
This case in fact corresponds to the surface that was constructed in the subsection above. Note that if we take $`\sigma _{12}`$ around $`\sigma _{23\mathrm{}n+1}`$ then it becomes $`\sigma _{13}`$. The OPE of $`\sigma _{13}`$ with $`\sigma _{23\mathrm{}n+1}`$ is still an operator of the form $`\sigma _{n+1}`$. On the other hand if we take the two $`\sigma _2`$ operators near each other, then we get the identity if we have $`\sigma _{12}\sigma _{12}`$, but we get $`\sigma _{123}`$ if move the operators through a path such that they become $`\sigma _{12}`$ and $`\sigma _{13}`$. In fact by moving the various operators around each other on the $`z`$ plane, we can also get from the same correlator OPEs of the form $`\sigma _{12}\sigma _{34}`$ (which is nonsingular) and $`\sigma _{12}\sigma _{n,n1,\mathrm{}21}`$ which produces an operator $`\sigma _{n1}`$. Thus we should find singularities in the 4-point function arising from this surface to correspond to all these possibilities.
(c) The indices of $`\sigma _2`$ and the indices of $`\sigma _n`$ have two overlaps, and the total number of indices involved in the correlator is $`s=n`$. (Note that case (b) above also could be brought to a form where $`\sigma _2`$ and $`\sigma _n`$ have two overlaps, but the number of indices involved there overall was $`n+1`$.) The other two operators in the correlator must have a similar overlap of indices, since otherwise they cannot produce an operator that has only $`n`$ distinct indices. In this case the genus of the covering surface is
$$g=\frac{1+1+(n1)+(n1)}{2}n+1=1.$$
(7.25)
We see that the correlator $`O_2O_2O_2O_2`$ will have contributions from correlators $`\sigma _2\sigma _2\sigma _2\sigma _2`$ that give genus 0 and genus 1 surfaces, but no other surfaces. The genus 0 case is contained in the analysis in subsection 7.1, and we will study the genus 1 case in subsection 7.4 below. Note that for the genus 0 case we have many combinations of indices for the $`\sigma _2`$ operators as discussed in (b) above, but these all arise from different branches of the same function (7.23). We must thus add the results from these branches (as well as the disconnected part (case (a) above) and the genus 1 contribution) to obtain the complete 4-point function of the $`O_2`$ operators. We will not carry out the explicit addition since we expect the result to be simpler in the supersymmetric case, which we hope to present elsewhere.
### 7.3 Analysis of the $`g=0`$ contribution
In this subsection, we analyse the correlator computed in (7.23), which corresponds to case (b) above, to check if it reproduces the expected short distance limits.
#### 7.3.1 The limit $`w0`$.
First let us consider the limit $`w0`$, which corresponds to $`t_{}0`$ and $`a0`$. For small values of $`a`$ we have:
$`t_+`$ $``$ $`{\displaystyle \frac{n+1}{n}},t_{}{\displaystyle \frac{an}{n+1}},C{\displaystyle \frac{n^n}{(n+1)^{n+1}}},`$
$`z_{}`$ $``$ $`n^{2n}(n+1)^{22n}a^{n+1},`$
$`\mathrm{log}_4`$ $``$ $`{\displaystyle \frac{1}{12}}\left({\displaystyle \frac{1}{n(n+1)}}{\displaystyle \frac{1}{2}}\right)\mathrm{log}z_{}{\displaystyle \frac{5}{12}}\mathrm{log}2`$
$``$ $`\left({\displaystyle \frac{n}{6}}+{\displaystyle \frac{1}{6(n+1)}}+{\displaystyle \frac{1}{12}}\right)\mathrm{log}n+\left({\displaystyle \frac{n}{6}}+{\displaystyle \frac{1}{6n}}+{\displaystyle \frac{1}{12}}\right)\mathrm{log}(n+1).`$
One can see that the correct singularity $`(z_{})^{2(\mathrm{\Delta }_n+\mathrm{\Delta }_2\mathrm{\Delta }_{n+1})}`$ is reproduced. Using the expressions for three point functions we derived before one can check that in the limit $`w0`$:
$$\mathrm{log}_42(\mathrm{\Delta }_n+\mathrm{\Delta }_2\mathrm{\Delta }_{n+1})\mathrm{log}w+2\mathrm{log}|C_{n,2,n+1}^\sigma |^2,$$
(7.27)
which agrees with anticipated factorization.
#### 7.3.2 The limit $`w1`$
Let us now consider a limit $`a1`$, which corresponds to one of two possible ways for point $`w`$ to approach $`1`$. After introducing $`b=a1`$ we get:
$`t_\pm `$ $``$ $`1+{\displaystyle \frac{b(n1)}{2n}}\pm {\displaystyle \frac{1}{2n}}\sqrt{4nb},C1,`$
$`t_\pm a`$ $``$ $`\pm \sqrt{{\displaystyle \frac{b}{n}}}{\displaystyle \frac{b}{n}}{\displaystyle \frac{n+1}{2}},`$
$`z_{}z_+`$ $`=`$ $`{\displaystyle \frac{z_{}}{z_+}}1\left(12\sqrt{{\displaystyle \frac{b}{n}}}\right)^n\left(1(n1)\sqrt{{\displaystyle \frac{b}{n}}}\right)`$
$`\times `$ $`\left(1(n+1)\sqrt{{\displaystyle \frac{b}{n}}}\right)14i\sqrt{nb},`$
$`\mathrm{log}_4`$ $``$ $`{\displaystyle \frac{1}{4}}\mathrm{log}(w1)=4\mathrm{\Delta }_2\mathrm{log}(w1).`$ (7.28)
This singularity corresponds to $`\sigma _2`$ and $`\sigma _2`$ fusing to the identity.
There is another limit ($`a\frac{(n+1)^2}{(n1)^2}`$) which also corresponds to $`w1`$. Introducing $`b=a(n+1)^2/(n1)^2`$, we get:
$`t_\pm `$ $``$ $`{\displaystyle \frac{n+1}{n1}}\pm \sqrt{{\displaystyle \frac{b}{n}}},C\left({\displaystyle \frac{n1}{n+1}}\right)^{n+1},`$
$`{\displaystyle \frac{dz}{dt}}`$ $`=`$ $`{\displaystyle \frac{Cnt^n}{(t1)^2}}(tt_+)(tt_{}){\displaystyle \frac{n(n1)^4}{4(n+1)^2}}(tt_+)(tt_{})`$
$`z_{}z_+`$ $``$ $`{\displaystyle \frac{n(n1)^4}{4(n+1)^2}}{\displaystyle _{t_+}^t_{}}𝑑t(tt_+)(tt_{})={\displaystyle \frac{(n1)^4b^{3/2}}{3\sqrt{n}(n+1)^2}},`$
$`\mathrm{log}_4`$ $``$ $`{\displaystyle \frac{1}{36}}\mathrm{log}(w1)+\mathrm{log}(n1)\left[{\displaystyle \frac{1}{9}}{\displaystyle \frac{1}{6}}\left(n+{\displaystyle \frac{1}{n}}\right)\right]{\displaystyle \frac{2}{9}}\mathrm{log}n`$
$`+`$ $`\mathrm{log}(n+1)\left[{\displaystyle \frac{1}{6}}\left(n+1+{\displaystyle \frac{1}{n}}\right){\displaystyle \frac{1}{18}}\right]{\displaystyle \frac{2}{3}}\mathrm{log}2{\displaystyle \frac{1}{36}}\mathrm{log}3.`$
Using the equations (6.25) and (6.22) one can see that
$$\mathrm{log}_42(\mathrm{\Delta }_2+\mathrm{\Delta }_2\mathrm{\Delta }_3)\mathrm{log}(w1)+\mathrm{log}|C_{223}^\sigma |^2+\mathrm{log}|C_{n3n}^\sigma |^2.$$
(7.30)
This corresponds to merging $`\sigma _2`$ and $`\sigma _2`$ to $`\sigma _3`$.
#### 7.3.3 The limit $`w\mathrm{}`$.
We now look at the remaining limits of the expression (7.23) at which the four point function becomes singular. They emerge at the points where the coefficient $`C`$ goes either to $`0`$ or infinity, i.e. if the value of $`t_+`$ approaches one of the points: $`0,1,a,\mathrm{}`$. Substituting this to the quadratic equation (7.4), we get the candidates for the critical values of $`a`$: $`0,1,\mathrm{}`$. Two of these limits we already considered, now we analyze the last possibility: $`a\mathrm{}`$. In this limit we have:
$$t_+a\frac{n1}{n}\frac{1}{n(n1)},t_{}\frac{n}{n1}.$$
(7.31)
So it is convenient to keep a point $`z_{}`$ fixed and vary the value of $`z_+`$ instead<sup>2</sup><sup>2</sup>2Note that the definition of $`t_+`$ and $`t_{}`$ depend on the choice of branch for a multivalued function; in particular $`t_+`$ and $`t_{}`$ interchange if one goes along a small circle around the point $`a=1`$.. Thus we get:
$`C`$ $`=`$ $`t_{}^n{\displaystyle \frac{t_{}1}{t_{}a}}\left({\displaystyle \frac{n1}{n}}\right)^n{\displaystyle \frac{1}{a(n1)}},`$
$`z_+`$ $``$ $`a^{n1}\left({\displaystyle \frac{n1}{n}}\right)^{2n}{\displaystyle \frac{1}{(n1)^2}},`$
$`\mathrm{log}_4`$ $``$ $`\left({\displaystyle \frac{1}{24}}+{\displaystyle \frac{1}{12n(n1)}}\right)\mathrm{log}z_++\mathrm{log}(n1)\left[{\displaystyle \frac{1}{12}}{\displaystyle \frac{1}{6}}\left(n+{\displaystyle \frac{1}{n}}\right)\right]`$ (7.32)
$`+`$ $`\left[{\displaystyle \frac{1}{6}}\left(n1+{\displaystyle \frac{1}{n1}}\right)+{\displaystyle \frac{1}{12}}\right]\mathrm{log}n{\displaystyle \frac{5}{12}}\mathrm{log}2.`$
This expression can be rewritten in terms of three point functions:
$$\mathrm{log}_42(\mathrm{\Delta }_n\mathrm{\Delta }_2\mathrm{\Delta }_{n1})+2\mathrm{log}|C_{n1,2,n}^\sigma |^2,$$
(7.33)
thus it corresponds to the factorization of the following type:
$$\left(\sigma _{(1\mathrm{}n)}\sigma _{(12)}\right)\left(\sigma _{(12)}\sigma _{(1\mathrm{}n)}\right).$$
(7.34)
Thus the four point function reproduces the anticipated factorizations.
### 7.4 The $`g=1`$ correlator $`\sigma _{12}\sigma _{12}\sigma _{12}\sigma _{12}`$
Let us consider the case $`n=2`$ in (7.1), so that we have the correlator
$$\sigma _2(0)\sigma _2(1)\sigma _2(w)\sigma _2(\mathrm{})$$
(7.35)
We wish to have the number of sheets over a generic point in the $`z`$ plane to be $`2`$; this gives $`g=1`$ for the covering surface $`\mathrm{\Sigma }`$. Each branch point is of order $`2`$, so we seek a map of the form
$$\frac{dz}{dt}=\alpha [z(z1)(zw)(zz_{\mathrm{}})]^{1/2}$$
(7.36)
We choose not to put any branch point at infinity explicitly, the limit $`z_{\mathrm{}}\mathrm{}`$ will be taken in the end of the calculation. This equation may be solved using the Weierstrass function $`𝒫`$ and the solution in the $`z_{\mathrm{}}\mathrm{}`$ limit is given by:
$$z(t)=\frac{𝒫(t)e_1}{e_2e_1}$$
(7.37)
where
$`e_1`$ $`=`$ $`𝒫({\displaystyle \frac{1}{2}}),e_2=𝒫({\displaystyle \frac{\tau }{2}}),e_3=𝒫({\displaystyle \frac{1}{2}}+{\displaystyle \frac{\tau }{2}})`$
$`w`$ $`=`$ $`\left({\displaystyle \frac{\theta _3(\tau )}{\theta _4(\tau )}}\right)^4={\displaystyle \frac{e_3e_1}{e_2e_1}}`$ (7.38)
The coordinate $`t`$ describes a torus given by modding out the complex plane with translations by $`1`$ and $`\tau `$.
We choose the fiducial metric on the torus to be that flat metric $`d\widehat{s}^2=dtd\overline{t}`$. Then we calculate the contribution to the Liouville action from the point $`z=0`$. Near this point
$`{\displaystyle \frac{dz}{dt}}`$ $``$ $`\alpha z^{1/2}\sqrt{wz_{\mathrm{}}}`$
$`\varphi `$ $`=`$ $`\mathrm{log}{\displaystyle \frac{dz}{dt}}+c.c.\mathrm{log}\left(\alpha z^{1/2}\sqrt{wz_{\mathrm{}}}\right)+c.c.`$
$`_t\varphi `$ $`=`$ $`{\displaystyle \frac{dz}{dt}}_z\varphi ,_z\varphi {\displaystyle \frac{1}{2z}}`$
We can write $`dt_t\varphi `$=$`dz_z\varphi `$ for any infinitesimal segment of the contour of integration around the puncture, but we must circle the $`z`$ plane puncture twice to circle the $`t`$ space puncture once. We find it easier to work in the $`z`$ plane instead of the $`t`$ space to evaluate the integral around the puncture. We recall that the puncture has a radius $`ϵ`$ in the $`z`$ plane, and we put in a factor of $`2`$ at the end to account for the relation between $`z`$ and $`t`$ contours. Then we get for the contribution to $`S_L`$:
$$S_L(z=0)=\frac{1}{24}\mathrm{log}\left[\alpha ^2ϵ|wz_{\mathrm{}}|\right]$$
(7.40)
Similarly from the points $`z=1`$, $`z=w`$ and $`z=z_{\mathrm{}}`$ we get the contributions
$`S_L(z=1)={\displaystyle \frac{1}{24}}\mathrm{log}\left[\alpha ^2ϵ|1w||1z_{\mathrm{}}|\right],`$ (7.41)
$`S_L(z=w)={\displaystyle \frac{1}{24}}\mathrm{log}\left[\alpha ^2ϵ|w||1w||wz_{\mathrm{}}|\right],`$ (7.42)
$`S_L(z=z_{\mathrm{}})={\displaystyle \frac{1}{24}}\mathrm{log}\left[\alpha ^2ϵ|z_{\mathrm{}}||z_{\mathrm{}}||wz_{\mathrm{}}|\right].`$ (7.43)
Now we look at the image of infinity, where we have cut a circle $`|z|=1/\delta `$. We have
$$\frac{dz}{dt}\alpha z^2,\varphi \mathrm{log}[\alpha z^2]+c.c.,_z\varphi \frac{2}{z}$$
(7.44)
Noting that we must take the $`z`$ plane contour clockwise, and putting in the factor of $`2`$ to relate the contour to the $`t`$ space contour, we get the contribution
$$S_L(z=\mathrm{})=\frac{1}{3}\mathrm{log}\left[\alpha \delta ^2\right]$$
(7.45)
Adding up all contributions we get
$`S_L^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{6}}\mathrm{log}ϵ{\displaystyle \frac{2}{3}}\mathrm{log}\delta {\displaystyle \frac{1}{12}}\mathrm{log}\left|w(1w)z_{\mathrm{}}(z_{\mathrm{}}1)(z_{\mathrm{}}w)\right|`$ (7.46)
$``$ $`{\displaystyle \frac{1}{6}}\mathrm{log}ϵ{\displaystyle \frac{2}{3}}\mathrm{log}\delta {\displaystyle \frac{1}{12}}\mathrm{log}\left|w(1w)\right|{\displaystyle \frac{1}{4}}\mathrm{log}|z_{\mathrm{}}|`$
Note that $`\alpha `$ does not appear in the final result, as one could anticipate from the fact that this constant can be absorbed in the rescaling of the $`t`$ plane. In the case under consideration there is no contribution to the Liouville action coming from the $`|z|>1/\delta `$ region: there is no curvature on the torus $`t`$ and
$$\frac{d\stackrel{~}{z}}{dt}=\frac{1}{\delta ^2z^2}\frac{dz}{dt}\frac{\alpha }{\delta ^2},$$
(7.47)
giving a constant $`\varphi `$ to leading order and thus a vanishing kinetic term for $`\varphi `$. Thus $`S_L=S_L^{(1)}`$ and the general expression (2.44) gives:
$$\sigma _2(0)\sigma _2(1)\sigma _2(w)\sigma _2(z_{\mathrm{}})_\delta =e^{S_L^{(1)}}\frac{Z^{(\widehat{s})}}{(Z_\delta )^2}.$$
(7.48)
To obtain the normalized 4-point function we write
$`\sigma _2(0)\sigma _2(1)\sigma _2(w)\sigma _2(\mathrm{})`$ $``$ $`\underset{|z_{\mathrm{}}|\mathrm{}}{lim}|z_{\mathrm{}}|^{4\mathrm{\Delta }_2}\sigma _2(0)\sigma _2(1)\sigma _2(w)\sigma _2(z_{\mathrm{}})`$ (7.49)
$`=`$ $`\underset{|z_{\mathrm{}}|\mathrm{}}{lim}|z_{\mathrm{}}|^{4\mathrm{\Delta }_2}{\displaystyle \frac{\sigma _2^ϵ(0)\sigma _2^ϵ(1)\sigma _2^ϵ(w)\sigma _2^{\stackrel{~}{\delta }}(z_{\mathrm{}})}{\sigma _2^ϵ(0)\sigma _2^ϵ(1)^2}}`$
$`=`$ $`2^{2/3}|w(1w)|^{\frac{1}{12}}Z_\tau `$
Here we have used the fact that in this case the partition function $`Z^{(\widehat{s})}`$ in (2.16) is that on the flat torus with modular parameter $`\tau `$ given through (7.38).
Since the group $`S^2`$ equals the group $`Z_2`$, we can compare (7.49) with the 4-point function obtained for $`\sigma _2`$ operators for the $`Z_2`$ orbifold in . (7.49) agrees with (4.13) of for the case of a noncompact boson field and with (4.16) for the compact boson field.
One observes that if the fields $`X^i`$ are noncompact bosons, then as $`w1`$ we find a factor $`\mathrm{log}(w1)`$ in the OPE in addition to the expected power $`(w1)^{1/4}`$. We suggest the following interpretation of this logarithm. There is a continuous family of momentum modes for the noncompact boson, with energy going to zero. If we do not orbifold the target space, then momentum conservation allows only a definite momentum mode to appear in the OPE of two fields. But the orbifolding destroys the translation invariance in $`X^1X^2`$, and nonzero momentum modes can be exchanged between sets of operators where each set does not carry any net momentum charge. The exchange of such modes (with dimensions accumulating to zero) between the pair $`\sigma _2(w)\sigma _2(1)`$ and the pair $`\sigma _2(0)\sigma _2(\mathrm{})`$ gives rise to the logarithm. Of course when the boson is compact, this logarithm disappears, as can be verified from (7.49) or the equivalent results in .
## 8 Discussion
The motivation for our study of correlation functions of symmetric orbifolds was the fact that the dual of the $`AdS_3\times S^3\times M`$ spacetime (which arises in black hole studies) is the CFT arising from the low energy limit of the D1-D5 system, and the D1-D5 system is believed to be a deformation of an orbifold CFT (with the undeformed orbifold as a special point in moduli space). To study this duality we must really study the supersymmetric orbifold theory, while in this paper we have just studied the bosonic theory. It turns out however that the supersymmetric orbifold can be studied with only a small extension of what we have done here. Following we can bosonize the fermions. Then if we go to the covering space $`\mathrm{\Sigma }`$ near the insertion of twist operator then we find only the following difference from the bosonic case – at the location of the twist operator we do not have the identity operator, but instead a ‘charge operator’ of the form $`P(X_a,^2X_a,\mathrm{})e^{i{\scriptscriptstyle k_aX_a}}`$. Here $`X_a`$ are the bosons that arise from bosonizing the complex fermions, and $`P`$ is a polynomial expression in its arguments. It is easy to compute the correlation function of these charge operators on the covering space $`\mathrm{\Sigma }`$, and then we have the twist correlation functions for the supersymmetric theory. We will present this calculation elsewhere, but here we note that many properties of interest for the supersymmetric correlation functions can be already seen from the bosonic analysis that we have done here. In this section we recall the features of the $`AdS/CFT`$ duality map and analyse some properties of the 3-point functions in the CFT.
### 8.1 ‘Universality’ of the correlation functions
We have mentioned before that while we have discsussed the orbifold theory $`R^N/S_N`$ (where the coordinate of $`R`$ gave $`X`$, a real scalar field), we could replace the CFT of $`X`$ by any other CFT of our choice, and the calculations performed here would remain essentially the same. When the covering surface $`\mathrm{\Sigma }`$ had genus zero, the results depend only on the value of $`c`$, and thus if we had $`(T^4)^N/S_N`$ theory or a $`(K_3)^N/S_N`$ theory, then we would simply choose $`c=4`$ in the Liouville action (2.17) (instead of $`c=1`$). If $`\mathrm{\Sigma }`$ had $`g=1`$, then we would need to put in the partition function (of a single copy) of $`T^4`$ or $`K_3`$ for the value of $`Z^{(\widehat{s})}`$ in (2.16). But apart from the value of $`c`$ and the value of partition functions on $`\mathrm{\Sigma }`$ there is no change in the calculation. Thus in particular the 3-point functions that we have computed at genus zero are universal in the sense that if we take them to the power $`c`$ then we get the 3-point functions for any CFT of the form $`M^N/S_N`$ with the CFT on $`M`$ having cenral charge $`c`$.
There is a small change in the calculation when we consider the supersymmetric case. The fermions from different copies of $`M`$ anticommute, and the twist operators carry a representation of the $`R`$ symmetry. As a consequence the dimensions of the twist operators are not given by (3.21), but for the supersymmetric theory based on $`M=T^4`$ are given by $`\frac{1}{2}(n1)`$. However as mentioned above, our analysis can be extended with small modifications to such theories as well.
Note that our method does not work if we have an orbifold group other than $`S_N`$. Thus for example if we had a $`Z_N`$ orbifold of a complex boson , then we could go to the covering space over a twist operator $`\sigma _n`$, but not write the CFT in terms of an unconstrained field on this covering space. The reason is that we have $`n`$ sheets or more of the cover over any point in the base space, but the central charge of the theory is just $`2`$, and so we cannot attribute one scalar field to each sheet of the cover. Thus our method, and its associated universalities, are special to $`S_N`$ orbifolds, where a twist operator just permutes copies of a given CFT but does not exploit any special symmetry of the CFT itself.
### 8.2 The genus expansion and the fusion rules of WZW models
We have studied the orbifold CFT on the plane, but found that the correlation functions can be organized in a genus expansion, arising from the genus of the covering surface $`\mathrm{\Sigma }`$. In the large $`N`$ limit the contribution of a higher genus surface goes like $`1/N^{g+\frac{1}{2}}`$. This situation is similar to that in the Yang-Mills theory that is dual to $`AdS_5\times S^5`$. The Yang-Mills theory has correlation functions that can be expanded in a genus expansion, with higher genus surfaces supressed by $`1/N^g`$. In the Yang-Mills theory the genus expansion has its origins in the structure of Feynman diagrams for fields carring two indices (the ‘double line representation’ of gauge bosons). In our case we have quite a different origin for the genus expansion. In the case of $`AdS_5\times S^5`$ it is believed that the genus expansion of the dual Yang-Mills theory is related to the genus expansion of string theory on this spacetime, though the precise relationship is not clear. It would be interesting if the genus expansion we have for the $`D1D5`$ CFT would be related to the genus expansion of the string theory on $`AdS_3\times S^3\times M`$.
In this context we observe the following relation. It was argued in that the orbifold CFT $`M^N/S_N`$ indeed corresponds to a point in the D1-D5 system moduli space. Further, at this point we have the number of 1-branes ($`n_1`$) and of 5-branes ($`n_5`$) given by $`n_5=1,N=n_1n_5=n_1`$. The dual string theory is in general an $`SU(2)`$ Wess-Zumino-Witten (WZW) model , though at the orbifold point of the CFT this string theory is complicated to analyze. The twist operators $`\sigma _n`$, $`n=1\mathrm{}N`$ of the CFT ($`\sigma _1=`$Identity) correspond to WZW primaries with $`j=(n1)/2,0j\frac{N1}{2}`$. Since in a usual WZW model we have $`0jk/2`$, we set $`k=N1`$.
The fusion rules for the WZW model, which give the 3-point functions of the string theory on the sphere (tree level) are as follows. The spins $`j`$ follow the rules for spin addition in $`SU(2)`$, except that there is also a ‘truncation from above’
$`(j_1,j_2)`$ $``$ $`j_3`$
$`|j_1j_2|`$ $``$ $`j_3|j_1+j_2|,j_1+j_2+j_3k`$ (8.1)
Now consider the 3-point function in the orbifold CFT, for the case where the genus of the covering surface $`\mathrm{\Sigma }`$ is $`g=0`$ . The ramification order of $`\mathrm{\Sigma }`$ at the insertion of $`\sigma _{n_i}`$ is $`r_i=(n_i1)=2j_i`$. The rules in (4.4), (4.5), (4.6) translate to $`|j_1j_2|j_3|j_1+j_2|`$. Further, the number of sheets $`s`$ is bounded as $`sN`$. Then the relation (4.1) gives
$$\underset{i}{}\frac{r_i}{2}=g1+s1+Nj_1+j_2+j_3k$$
(8.2)
While (8.2) is a relation for the bosonic orbifold theory, we expect an essentially similar relation for the supersymmetric case. Thus we observe a similarity between the $`g=0`$ 3-point functions of the WZW model (8.1) and of the CFT (8.2).
At genus $`g=1`$ however, we find that any three spins $`j_1,j_2,j_3`$ can give a nonzero 3-point function in the string theory. In the orbifold CFT, however, we get only a slight relaxation of the rule (8.2): we get $`j_1+j_2+j_3k+1`$. Roughly speaking we can reproduce this rule in the string theory if we require that in the string theory one loop diagram there be a way to draw the lines such that only spins $`j1/2`$ be allowed to circulate in the loop. Of course we are outside the domain of any good perturbation expansion at this point, since if the spins are of order $`k`$ then there is no small parameter in the theory to expand in, and thus there is no requirement that there be an exact relation between the rules in a WZW string theory and the rules in the orbifold CFT.
We note that in the 3-point functions of chiral primaries that were studied had ‘one overlap’ in their indices. This corresponds to $`j_1+j_2=j_3`$ in the above fusion rules, and since for the supersymmetric case the dimension is linear in the charge, we also have $`\mathrm{\Delta }_1+\mathrm{\Delta }_2=\mathrm{\Delta }_3`$. This corresponds to the case of ‘extremal’ correlation functions in the language of . In the 3-point correlators for this special case were found by an elegant recursion relation, which arises from the fact that there is no singularity in the OPE, and thus the duality relation of conformal blocks becomes a ‘chiral ring’ type of associativity law among the fusion coefficients. It is not clear however how to extend this method to the non-extremal case, and one motivation for the present work was to develop a scheme to compute the correlators for $`j_1+j_2<j_3`$, which corresponds to more than one overlap. In the case of one overlap we have extended our calculations to the supersymmetric case, and found results in agreement with .
### 8.3 3-point couplings and the stringy exclusion principle.
In the $`AdS_5\times S^5`$ case the 3-point couplings of supergravity agree with the large $`N`$ limit of the 3-point functions in the free Yang-Mills theory; thus there is a nonrenormalization of this correlator as the coupling $`g`$ is varied. It is not clear if a similar result holds for the $`AdS_3\times S^3\times M`$ case, and even less clear what nonrenormalization theorems hold at finite $`N`$. But it is nevertheless interesting to ask how the correlators in the orbifold CFT behave as we go from infinite $`N`$ to finite $`N`$, and in particular what happens as we approach the limits of the stringy exclusion principle.
Thus we examine the ratio
$$R(m,n,q;\overline{N})\frac{\sqrt{\overline{N}}O_nO_mO_q_{\overline{N}}}{lim_N\mathrm{}\sqrt{N}O_nO_mO_q_N}$$
(8.3)
where the subscripts on the correlator give the value of $`N`$. We have rescaled the correlators by $`\sqrt{N}`$ to obtain the effective coupling of the 3-point function; the correlator itself goes as $`1/\sqrt{N}`$. For $`n,m,q<<\overline{N}`$ we expect $`R1`$, while as $`n,m,q`$ become order $`\overline{N}`$ we expect that $`R`$ will fall to zero. We take the case of the 3-point function with single overlap, and further set $`m=n`$. Then we have $`q=2n1`$, and we write
$$R(n,n,2n1;\overline{N})R(n;\overline{N})$$
(8.4)
It is easy to see that for the case of single overlap the correlators $`\sigma _n\sigma _m\sigma _{m+n1}`$ can get a contribution only from surfaces $`\mathrm{\Sigma }`$ with $`g=0`$, for which case we have done a complete calculation of the correlator and its combinatorics. Note further that in the ratio (8.3) the actual value of $`\sigma _n\sigma _m\sigma _{m+n1}`$ will cancel, and the value of $`R`$ will be determined by combinatorial factors. These factors are expected to be the same for the bosonic and supersymmetric cases.
In the figure we plot $`R(n;\overline{N})`$ versus $`n`$ (for $`\overline{N}=1000`$). We see that $`R`$ drops significantly after $`n`$ exceeds $`\sqrt{\overline{N}}`$. This effect can be traced to the fact that the number of ways to select $`s`$ ordered indices from $`\overline{N}`$ indices is
$`\overline{N}(\overline{N}1)`$ $`\mathrm{}`$ $`(\overline{N}s+1)=\overline{N}^s(1{\displaystyle \frac{1}{\overline{N}}})(1{\displaystyle \frac{2}{\overline{N}}})\mathrm{}(1{\displaystyle \frac{s1}{\overline{N}}})`$ (8.5)
$``$ $`\overline{N}^s(1{\displaystyle \frac{1}{\overline{N}}}{\displaystyle \underset{j=1}{\overset{s1}{}}}j)=\overline{N}^s(1{\displaystyle \frac{1}{\overline{N}}}{\displaystyle \frac{s(s1)}{2}})`$
If the CFT 3-point function is indeed not renormalized for finite $`N`$, then the above result has interesting implications. The coupling between three gravitons would then be a constant for low energies ($`n<<\overline{N}`$) but would drop rapidly for high energies. Thus the behavior of high frequency modes would not follow a naive ‘equivalence principle’.
This issue may be relevant to Hawking’s derivation of black hole radiation, where we need to make a change of coordinates to study the high frequency modes near the horizon. For these modes to evolve as used in the derivation, we use implicitly the naive value of the following cubic coupling: that of a low energy graviton (representing the attraction of the hole) and two high energy quanta (representing the high energy mode emerging from the horizon, getting redshifted by the attraction of the hole). If this coupling differs from the one expected from naive gravitational physics, then the semiclassical derivation of Hawking radiation may require modification, with corresponding implications for the information paradox.
### 8.4 Conclusion
It would be important to pursue further the study of the supersymmetric case, and to compare with the dual superstring theory. The subset of correlators computed in was compared to supergravity in , but it was a little unclear how closely the two calculations agreed. A better picture may emerge when we look at the complete set of correlators of the supersymmetric side, which is possible to do by extending our computation here to include the R-charges carried by the twist operators in the supersymmetric case.
It was argued recently in that the CFT of the D1-D5 system exhibits a duality to a set of spacetimes, of which the AdS space is only one member. If the 3-point functions are protected against coupling changes then we should see a reflection of this fact in correlators at the orbifold point.
## Acknowledgements
We are grateful to A. Jevicki, M. Mihailescu, S. Ramgoolam, and S. Frolov for patiently explaining their results to us, and to L. Rastelli for extensive discussions in the early phase of this work. We also benefited greatly from discussions with S. Das, E. D’Hoker, C. Imbimbo, F. Larsen, E. Martinec, S. Mukhi, S. Sethi, and S.T. Yau.
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# How an Anomalous Cusp Bifurcates
## Abstract
We study the pattern of activated trajectories in a double well system without detailed balance, in the weak noise limit. The pattern may contain cusps and other singular features, which are similar to the caustics of geometrical optics. Their presence is reflected in the quasipotential of the system, much as phase transitions are reflected in the free energy of a thermodynamic system. By tuning system parameters, a cusp may be made to coincide with the saddle point. Such an anomalous cusp is analogous to a nonclassical critical point. We derive a scaling law, and nonpolynomial ‘equations of state’, that govern its bifurcation into conventional cusps.
The optimal trajectory concept has been widely used in the theory of noise-activated transitions . In the weak noise limit, when transitions between stable states become exponentially rare, one or at most a few trajectories in the system state space are singled out as escape paths of least resistance. Also, between the energetically lowest stable state and any other state there are at most a few dominant activated trajectories. Such optimal trajectories, which are determined by a ‘least energy expended’ or ‘least action’ variational principle, are experimentally observable . In systems that have the property of detailed balance, they are time-reversed relaxational trajectories. But in nonequilibrium systems, which lack detailed balance, the optimal trajectory pattern extending from a stable state may be more complicated.
Optimal trajectories are similar in many ways to the rays of geometrical optics, which characterize in the short wavelength limit the waves emanating from a point source. That is because optical rays may be computed variationally too, from a ‘least optical depth’ principle. In a medium with inhomogeneous index of refraction, it is common for a ray family to bounce off a caustic surface, leaving the region behind in shadow: not illuminated, or illuminated only indirectly. Other singular features with a catastrophe-theoretic interpretation may be produced . In a noise-driven system in which detailed balance is violated, the pattern of optimal trajectories may contain similar features . See Fig. 1.
To understand the crossing of optimal trajectories, it is useful to look at the quasipotential of the noise-driven system. If $`ϵ`$ is the noise strength (e.g., $`ϵ=kT`$ in thermal systems), and $`\rho (𝐱)`$ denotes the stationary probability density of the system at state $`𝐱`$, then a quasipotential $`W=W(𝐱)`$ may be defined phenomenologically by
$$\rho (𝐱)\text{const}\times \mathrm{exp}[W(𝐱)/ϵ],ϵ0.$$
(1)
This definition makes sense whether or not the system dynamics are conservative, and whether or not the noise acts so as to preserve detailed balance. $`W`$ equals zero at the energetically lowest stable state(s), and $`W(𝐱)`$ is essentially the minimum energy needed to excite the system to state $`𝐱`$. It may be computed as a line integral along the optimal trajectory extending to $`𝐱`$.
Formally, $`W`$ is multivalued at any state, such as the states near a caustic, that is reached by more than one optimal trajectory emanating from the energetically lowest state(s). But by (1), the least value is dominant, and the trajectories giving rise to others are unphysical. The state space of a noise-driven system is typically partitioned by ‘switching surfaces’, on which dominance switches between branches of $`W`$.
The switching of dominance resembles a first-order phase transition in a condensed matter system. The similarity is unsurprising, since phase transitions with classical critical exponents also have a catastrophe-theoretic interpretation . Consider, for example, a ferromagnetic system with extensive order parameter $`m`$ (magnetization), in a magnetic field $`h`$. Its thermodynamics are determined by a free energy function $`\mathrm{\Psi }=\mathrm{\Psi }(T,h)`$. Below the critical temperature $`T_c`$, $`\mathrm{\Psi }`$ and $`m=\mathrm{\Psi }/h`$ are multivalued. If the phase transition is classical, i.e., of mean-field form, $`\mathrm{\Psi }`$ is three-valued in a sharp-tipped region of the $`(T,h)`$ plane bounded by ‘spinodals’ of the form $`h\pm \text{const}\times (T_cT)^{3/2}`$. That is because the leading terms in the Legendre transform $`\stackrel{~}{\mathrm{\Psi }}^{(h)}(T,m)hm\mathrm{\Psi }(T,h)`$ are of Ginzburg–Landau type:
$$\stackrel{~}{\mathrm{\Psi }}^{(h)}(T,m)C_2(TT_c)m^2/2+C_4m^4/4.$$
(2)
In the catastrophe-theoretic sense , the spinodals are fold caustics. Each is the projection of a fold in the graph of $`\mathrm{\Psi }`$, which is a two-dimensional surface, onto the $`(T,h)`$ plane. The critical point $`(T_c,0)`$ from which the spinodals extend is a cusp catastrophe: the projection of the point on the graph of $`\mathrm{\Psi }`$ at which the two folds join. A first-order phase transition line, on which dominance switches between branches of $`\mathrm{\Psi }`$, extends from $`(0,0)`$ to $`(T_c,0)`$.
Switching lines in two-dimensional noise-driven systems are clearly analogous to first-order phase transition lines, and caustics to spinodals. Caustics typically terminate at cusps, and switching lines also frequently terminate at cusps. So cusps, which are very common, are analogous to second-order critical points . They are physically important because at any cusp, the prefactor ‘const’ in (1), which in general is $`𝐱`$-dependent, diverges.
In previous work , we pointed out that in many noise-driven two-dimensional double well systems without detailed balance, a cusp may be moved to coincide with the saddle point between the wells, by tuning parameters. If they coincide, the Kramers ($`ϵ0`$) limit of noise-induced interwell transitions is greatly affected. The prefactor in the Kramers transition rate formula becomes anomalous: it acquires a negative power of $`ϵ`$.
Precisely at criticality, we were able to approximate the quasipotential $`W`$ near any such ‘anomalous cusp’. Our expression differed from the polynomial ‘normal forms’ of conventional catastrophe theory. In thermodynamics, it would correspond to a nonclassical phase transition. In catastrophe theory, it would be interpreted as a nongeneric catastrophe: one of the few such of physical relevance to have been discovered since the work of Berry and Mount on the short wavelength limit of scattering .
In this Letter, we extend Ref. by analysing the ‘unfolding’ of an anomalous cusp in a typical two-dimensional noise-activated system, as a parameter is moved toward or away from criticality. We explain how it may bifurcate into conventional cusps. Our scaling law for the bifurcation yields a corresponding law for the divergence of the Kramers prefactor .
Consider the following double well model, which is similar to models of blocking dynamics in glassy systems, where a particle is coupled to a randomly fluctuating barrier whose position is coupled to the particle motion . Let $`x`$ (a particle position variable) and $`y`$ (a barrier state variable) be nontrivially coupled in such a way that the values $`\pm 1`$ for $`x`$ and $`0`$ for $`y`$ are stable. If $`x`$ and $`y`$ are overdamped and are driven by white noise of strength $`ϵ`$, their joint dynamics could be modeled by Langevin equations
$`\dot{x}`$ $`=`$ $`\lambda _x\left[x\left(1x^2\right)\alpha xy^2\right]+ϵ^{1/2}\eta _x(t)`$ (3)
$`\dot{y}`$ $`=`$ $`|\lambda _y|(1+x^2)y+ϵ^{1/2}\eta _y(t).`$ (4)
The parameters $`\lambda _x>0`$ and $`\lambda _y<0`$ determine the time scales on which $`x`$ and $`y`$ evolve, and govern the all-important relaxational behavior near the saddle point $`(0,0)`$, where $`(\dot{x},\dot{y})(\lambda _xx,|\lambda _y|y)`$. The forcing terms $`(\eta _x,\eta _y)`$ are a pair of independent Gaussian white noises, so that $`\eta _i(s)\eta _j(t)`$ equals $`\delta _{ij}\delta (st)`$. The parameter $`\alpha `$ controls the absence of detailed balance: only when $`\alpha `$ equals $`\mu |\lambda _y|/\lambda _x`$ is there detailed balance, since only in that case is the drift field derived from a potential.
Our results are insensitive to the details of the coupling between $`x`$ and $`y`$, so long as the model is symmetric through $`x=0`$ and $`y=0`$. To compute the pattern of optimal trajectories emanating from the bottom of the $`x<0`$ well or the $`x>0`$ well, we use the fact that in any multidimensional noise-driven system with vector Langevin equation $`\dot{𝐱}=𝐮(𝐱)+ϵ^{1/2}𝜼(t)`$, the optimal trajectories are really zero-energy Hamiltonian trajectories, generated by the Wentzell–Freidlin Hamiltonian
$$H(𝐱,𝐩)=𝐩^2/2+𝐮(𝐱)𝐩.$$
(5)
That is because the associated Hamilton’s principle is
$$\delta L(𝐱,\dot{𝐱})𝑑t=\delta \left|\dot{𝐱}𝐮(𝐱)\right|^2𝑑t=0,$$
which is clearly a ‘least energy expended’ principle. The conjugate momentum $`𝐩L/\dot{𝐱}`$ equals $`\dot{𝐱}𝐮`$, which measures the system’s motion against the drift.
Figure 2 was obtained from (4) and (5) by integrating Hamilton’s equations outward, at zero energy, from $`(1,0)`$, i.e., from the bottom of the right-hand well. A small portion of the $`x<0`$ half-plane is reached by optimal trajectories, but the rest is in shadow. In phase space, which is four-dimensional, the optimal trajectories trace out a two-dimensional manifold, called a Lagrangian manifold. This manifold lies ‘above’ only a small part of the $`x<0`$ half-plane. It folds over, covering the shaded portion of the $`(x,y)`$ plane more than once. In the shaded region, the momentum $`𝐩`$ and the quasipotential $`W`$, which equals $`𝐩𝑑𝐱`$, are two-valued ($`x<0`$) or three-valued ($`x>0`$). At any point $`𝐱`$, $`𝐩`$ equals $`\mathbf{}W(𝐱)`$.
In Fig. 2, the parameter $`\alpha `$ is chosen to be slightly greater than a certain critical value, $`\alpha _c`$. If there is detailed balance, the optimal trajectory pattern contains no singular features, but if $`\alpha `$ is increased through $`\alpha _c`$, a cusp $`(x_c,0)`$ emerges from the saddle point at $`(0,0)`$ and moves toward $`(1,0)`$. This phenomenon is not peculiar to the model defined by (4). In any symmetric two-dimensional double well system that violates detailed balance and has a tunable parameter, a similar cusp may be born. The focusing of optimal trajectories at the cusp resembles the focusing of rays in a radially symmetric optical system.
The value $`\alpha _c`$ can be computed from the second-order variational equation $`\delta ^2L𝑑t=0`$, which is a criterion for bifurcation. On physical grounds, when $`\alpha <\alpha _c`$, $`\delta ^2L𝑑t`$ computed along the on-axis trajectory to the saddle is positive, but when $`\alpha >\alpha _c`$, it is negative. In the model (4), $`\alpha _c`$ turns out (cf. Ref. ) to equal $`2\mu (\mu +1)`$.
The cusp $`(x_c,0)`$ that is present when $`\alpha >\alpha _c`$ resembles a second-order critical point. $`W`$ is three-valued in the sharp-tipped region extending from it, which is bounded by ‘spinodals’ of the form $`y\pm \text{const}\times (x_cx)^{3/2}`$. (See Fig. 2.) Moreover, there is a switching line extending from the saddle at $`(0,0)`$ to $`(x_c,0)`$. As noted, this line is analogous to a first-order phase transition line.
What remains to be understood is how the cusp is born at $`\alpha =\alpha _c`$. In a three-dimensional space with coordinates $`(x,y;\alpha )`$, there is a line of second-order critical points in the $`y=0`$ plane that extends from $`(0,0;\alpha _c)`$ to $`(1,0;+\mathrm{})`$. By analogy with thermodynamics, one might expect $`(0,0;\alpha _c)`$ to be a third-order critical point. At any fixed $`\alpha >\alpha _c`$, the leading terms in the Legendre transform $`\stackrel{~}{W}^{(y)}(x,p_y)yp_yW`$, close to the cusp, are known to be of Ginzburg–Landau type :
$$\stackrel{~}{W}^{(y)}(x,p_y)C_2(\alpha )\left[xx_c(\alpha )\right]p_y^2/2+C_4(\alpha )p_y^4/4.$$
(6)
One might expect that the correct three-dimensional generalization would be a higher-degree polynomial in $`x`$, $`p_y`$, and $`\alpha \alpha _c`$. That would allow the birth of the cusp to be viewed as a classical phase transition, or one of the generic (polynomial) elementary catastrophes .
In Ref. we presented initial evidence against this. At criticality ($`\alpha =\alpha _c`$), we were able to construct a scaling solution for $`W`$, valid near the $`x`$-axis close to the saddle. The equation satisfied by the scaling function contained non-integer powers: in fact, powers that depended continuously on the model parameter $`\mu `$.
By linearizing Hamiltonian dynamics near the saddle, we have now characterized fully the behavior of the quasipotential $`W`$ near a singular point like $`(0,0;\alpha _c)`$. Our chief new result is a cubic equation satisfied by the double Legendre transform of $`W`$. It defines a higher-order, but nonclassical, critical point. We have also extended our $`\alpha =\alpha _c`$ scaling law to the case when $`|\alpha \alpha _c|`$ is nonzero but small. These results should extend to any symmetric double well system with a tunable parameter.
Figure 3, which was obtained in the ‘unbroken phase’ ($`\alpha <\alpha _c`$), sheds light on behavior near criticality. The crucial feature is the two caustics, relics of which appeared in Fig. 2. They form part of the boundary of the ‘illuminated’ region. Each caustic extends from a cusp, which is located very close to the $`y`$-axis separatrix between the two wells. As $`\alpha \alpha _{c}^{}{}_{}{}^{}`$, the cusps neck down to the saddle. Any further increase in $`\alpha `$ causes the $`x`$-axis cusp to be born, and to move toward positive $`x`$.
The merged cusp at $`(0,0)`$, when $`\alpha =\alpha _c`$, is truly anomalous. It is the projection of a point on the boundary of the Lagrangian manifold, rather than of a point in its interior. So it is a boundary catastrophe: a singular point of a sensitive kind. To explain its bifurcation into cusps on the $`x`$-axis or $`y`$-axis, we must approximate the optimal trajectory pattern in a neighborhood of the saddle, for $`|\alpha \alpha _c|`$ nonzero but small.
We accordingly linearize Hamilton’s equations, which any optimal trajectory must satisfy, near the point $`(x,y;p_x,p_y)=(0,0;0,0)`$ in phase space. A simple analysis (cf. Ref. ) shows that this fixed point has two stable directions, $`𝐞_s=(0,1;0,0)`$ and $`\stackrel{~}{𝐞}_s=(1,0;2\lambda _x,0)`$, and two unstable directions, $`𝐞_u=(1,0;0,0)`$ and $`\stackrel{~}{𝐞}_u=(0,1;0,2|\lambda _y|)`$. The zero-momentum directions (no tilde) are eigendirections for relaxational trajectories, which follow the drift. In the linear approximation, any optimal trajectory near $`(x,y)=(0,0)`$ must satisfy
$`(x,y;p_x,p_y)`$ $``$ $`C_se^{|\lambda _y|t}𝐞_s+\stackrel{~}{C}_se^{\lambda _xt}\stackrel{~}{𝐞}_s`$ (8)
$`+C_ue^{\lambda _xt}𝐞_u+\stackrel{~}{C}_ue^{|\lambda _y|t}\stackrel{~}{𝐞}_u,`$
where the $`C`$’s are trajectory-specific constants.
We now index the ‘fan’ of optimal trajectories that approach the saddle point, as in Figs. 2 and 3, by $`s`$. The normalization of this index variable is somewhat arbitrary. A reasonable choice would be for it to denote distance from the $`x`$-axis (at a fixed $`x>0`$, near the saddle). With this choice, $`s=0`$ will correspond to the uphill optimal trajectory that climbs toward $`(0,0)`$ along the positive $`x`$-axis. If each coefficient in (8) can be expanded in $`s`$ about $`s=0`$, then by symmetry considerations
$`C_s`$ $`=`$ $`a_1s+a_3s^3+\mathrm{},`$ (10)
$`\stackrel{~}{C}_s`$ $`=`$ $`b_0+b_2s^2+\mathrm{},`$ (11)
$`C_u`$ $`=`$ $`c_2s^2+c_4s^4+\mathrm{},`$ (12)
$`\stackrel{~}{C}_u`$ $`=`$ $`d_1s+d_3s^3+\mathrm{}.`$ (13)
We identify the passage through criticality, as $`\alpha `$ is increased through $`\alpha _c`$, with the passing through zero of the coefficients $`c_2`$ and $`d_1`$. So, setting $`\delta \alpha _c\alpha `$, we take $`c_2`$ and $`d_1`$ to be linearly proportional to $`\delta `$, to leading order.
Eq. (8) comprises four scalar equations. Eliminating $`s`$ and $`t`$ among them, we can derive ‘equations of state’ relating $`\delta `$ and any three of the phase space coordinates $`x`$, $`y`$, $`p_x`$, and $`p_y`$. When $`|\delta |1`$, the equation relating $`x`$, $`y`$, $`p_y`$, and $`\delta `$ (near the $`x`$-axis), and the equation relating $`x`$, $`y`$, $`p_x`$, and $`\delta `$ (near the $`y`$-axis), turn out to be, respectively,
$`0`$ $`=`$ $`(p_y2|\lambda _y|y)^3`$ (16)
$`+k_1\delta x^{2\mu }(p_y2|\lambda _y|y)+k_0x^{4\mu }p_y`$
$`0`$ $`=`$ $`(p_x+2\lambda _xx)^3`$ (18)
$`+\mathrm{}_1\delta p_x^{2\mu 2}y^2(p_x+2\lambda _xx)+\mathrm{}_0y^4p_x^{4\mu 3}`$
Here $`k_1`$, $`k_0`$, $`\mathrm{}_1`$, and $`\mathrm{}_0`$ are positive constants. At criticality ($`\delta =0`$), (16)–(18) reduce to the equations we previously obtained by an altogether different technique .
By definition, the cusp $`(x_c,0)`$ is the point on the positive $`x`$-axis where $`W`$ or $`(p_x,p_y)=\mathbf{}W`$ stops being multivalued as a function of $`(x,y)`$, as $`x`$ increases from $`0`$. It is easy to verify from (16) that when $`\delta `$ is small and negative (i.e., $`\alpha \alpha _c`$ is small and positive), the cusp location $`x_c`$ satisfies $`x_c(\delta )^{1/2\mu }`$. The $`\delta `$-dependence of the parent cusps, which have $`y=\pm y_c`$, can be computed from (18). They are the points close to the $`y`$-axis where $`(p_x,p_y)`$ first becomes multivalued, as $`|y|`$ increases from zero. When $`\delta 0^+`$ (i.e., $`\alpha \alpha _{c}^{}{}_{}{}^{}`$), we find that the parent cusps neck down at the rate $`y_c\delta ^{3/2\mu }`$. That is so when $`1\mu <3/2`$, at least; for other $`\mu `$, the prediction is that there are no parent cusps, and no necking down. All these predictions have been numerically confirmed .
Despite the continuously varying exponents, the emergence of the $`x`$-axis cusp is surprisingly similar to a second-order phase transition. Recall that close to the ferromagnetic critical point defined by (2), scaled magnetization $`M`$ and scaled magnetic field $`H`$ are related by
$$M^3\pm MH=0,$$
(19)
or equivalently by the scaling law $`M=\varphi _\pm (H)`$. Here $`MAm/|TT_c|^{1/2}`$ and $`HBh/|TT_c|^{3/2}`$, with $`A\sqrt{C_4/C_2}`$ and $`BA^3/C_4`$. The plus (minus) applies when $`TT_c`$ is positive (negative). Eq. (16) may be rewritten in the form (19), provided that one defines
$`M`$ $``$ $`(p_y2|\lambda _y|y)/|k_1\delta x^{2\mu }+k_0x^{4\mu }|^{1/2}`$ (20)
$`H`$ $``$ $`2|\lambda _y|k_0x^{4\mu }y/|k_1\delta x^{2\mu }+k_0x^{4\mu }|^{3/2}.`$ (21)
The plus (minus) applies when $`k_1\delta x^{2\mu }+k_0x^{4\mu }`$ is positive (negative). The law $`M=\varphi _\pm (H)`$, in which $`\delta `$ appears implicitly, provides a unified description of the $`x`$-axis behavior both at criticality ($`\delta =0`$) and away ($`\delta 0`$).
The most striking consequence of this approach is a general scaling law, showing how the quasipotential varies as $`(x,y;\alpha )`$ moves away from $`(0,0;\alpha _c)`$, in any direction. It can be written using the double Legendre transform $`\stackrel{~}{W}^{(x,y)}(p_x,p_y)𝐱𝐩W`$, which equals $`𝐱𝑑𝐩`$. Near the saddle,
$`W(x,y)`$ $``$ $`W(0,0)\lambda _xx^2+|\lambda _y|y^2`$ (22)
$`\stackrel{~}{W}^{(x,y)}(p_x,p_y)`$ $``$ $`W(0,0)p_x^2/4\lambda _x+p_y^2/4|\lambda _y|.`$ (23)
Let $`R=R(p_x,p_y)`$ denote the difference between $`\stackrel{~}{W}^{(x,y)}`$ and the right-hand side of (23). Since $`\stackrel{~}{W}^{(x,y)}=𝐱\dot{𝐩}𝑑t`$, $`R`$ may be expressed in terms of $`s`$ and $`t`$ by employing (8) and (3). By eliminating $`s`$ and $`t`$ from the formulæ for $`R(s,t)`$, $`p_x(s,t)`$, and $`p_y(s,t)`$, we find to leading order in $`\delta `$
$$R^3+m_1\delta p_x^{2\mu }p_y^2Rm_0p_x^{4\mu }p_y^4=0,$$
(24)
where $`m_1`$ and $`m_0`$ are positive constants. This is the extension to $`\delta 0`$ of a formula derived in Ref. .
In phase transition language, the cubic equation (24) fully characterizes the nonclassical structure of the critical point $`(x,y;\alpha )=(0,0;\alpha _c)`$, i.e., $`(p_x,p_y;\alpha )=(0,0;\alpha _c)`$. Equation (24) is also interesting from a catastrophe theory point of view. Nongeneric catastrophes, which are difficult to classify, may in general be perturbed in an infinite number of ways so as to yield generic catastrophes . But (24) describes the bifurcation of a nongeneric catastrophe (the anomalous cusp) into conventional cusps, in a unique, physically determined way.
This research was supported in part by NSF grant PHY-9800979. A portion was completed while the authors were in residence at the Aspen Center for Physics.
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# Magnetic field measurements in white dwarfs. Magnetic field, rotation and spectrum of 40 Eri B
## Introduction
The magnetic observations of dwarfs (WDs) have been carried out in the Laboratory of Stellar Physics of SAO since 1989. Originally the tasks were to detect magnetic fields in normal DA WDs observing them with an accuracy of about a few kG and to study the distribution of WDs over their surface magnetic fields (magnetic field function), to observe known magnetic WDs in order to find their rotational periods and to study the magnetic fields – rotation relation. Until 1995 the observations were fulfilled with the hydrogen–line magnetometer of the 6–m telescope (Stol’, 1991; Bychkov et al., 1991; Shtol’ et al., 1997; Fabrika et al., 1997). These observations and such observations of other authors (mainly by Schmidt & Smith, 1995) have shown that WDs with magnetic fields $`B\begin{array}{c}>\\ \end{array}10`$ kG are not numerous.
It was found in the later statistical studies (Fabrika & Valyavin, 1999; Valyavin & Fabrika, 1999) that the frequency of magnetic $`B\begin{array}{c}>\\ \end{array}1`$ MG WDs is about 2 $`\%`$ among hot (young stars) and it is about 20 $`\%`$ among cool (old) stars, i. e. the magnetic field does certainly evolve in WDs. It was also found that the magnetic field function is a power one, $`P_BB^\alpha `$, with a power index $`\alpha 1.3`$. Normalization properties of the magnetic field function allowed to estimate the frequency of magnetic WDs in the weak–field limit. We may expect the minimum surface magnetic field strength to be 1–10 kG in hot WDs ($`T>10000`$ K) and in cool WDs ($`T<10000`$ K) it is 10–50 kG. To confirm these conclusions by direct observations and to study the weak–field part of the magnetic field function, we decided to observe the brightest DA WDs with a higher accuracy.
## Zeeman spectroscopy with the Main Stellar Spectrograph
In 1995–1997 Zeeman spectral observations of bright DA WDs were carried out on the Main Stellar Spectrograph (MSS) of the 6–m telescope with a circular polarization analyzer (Valyavin et al., 1997; Fabrika & Valyavin, 1999). The time–resolved spectroscopy with 3 – 5 min exposure times was aimed at both detection of possible rotational variability in individual spectra and deriving high accuracy estimates of average magnetic field in sum spectra. Each image in this mode of observations consists of two simultaneous spectra of opposite circular polarization which are splitted in the analyzer. Two different orientations of the first quarter–wave plate are used to record opposite circular polarization spectra at the same place of the CCD detector in two adjacent exposures. We used this configuration of the higher spectral resolution with a hope to detect magnetic shifts in narrow non–LTE H$`\alpha `$ profiles of DA ($`8000÷22000`$ K) WDs. It was expected that observations in this mode could provide the desired accuracy of about 1 kG. In some brightest WDs such an accuracy has actually been reached, however without positive detections. For example in time–resolved Zeeman spectroscopy of WD 0713+584 we did not detect any periodical signals during a time of observations of 113 min. The total average effective magnetic field estimate has been obtained in this star particular $`B_e=0.1\pm 1.0`$ kG.
The brightest degenerate star 40 Eri B was observed in the same mode on September 14, 1995. The average magnetic field, $`B_e=0.5\pm 0.37`$ kG, during a time of observations of 81 min has been found. However in another observing run, on December 5, 1995, we found about 4 -hours’ sinusoidal variations of the effective field of 40 Eri B during 215 min of observations. These observations are displayed in Fig. 1. They consist of 54 (3 min) individual measurements averaged in 17 bins. The best fit is $`B_e=A+B\mathrm{cos}(2\pi t/P+\varphi )`$ with $`A=510\pm 520`$ G, $`B=2300\pm 700`$ G and a period $`P4`$ hours. It was concluded that for reliable detection of magnetic field and rotational period in 40 Eri B such mode of observations must be continued.
The MSS observations of 40 Eri B might be interpreting not only as the 4 - hours’ periodical variations. The period suspected cannot be considered as reliable, because the time of observations is comparable with the period, the scatter of the individual points is high, and one can say only that a variability of the magnetic field was detected that time. In addition we found other periodical signals in those data, one of them was $`37`$ min. There was also a period of about 10 min detected in the previous observing run on September 14, 1995. Indeed, we could say that some probable magnetic field variations were registered that time.
In order to demonstrate the magnetic shift of the H$`\alpha `$ line in opposite circular polarization spectra, we present here the results of analysis of two sum double–polarized spectra obtained from the individual spectra for the December 5, 1995 observing run. The first sum spectrum consists of individual spectra showing the maximal positive polarization (the top points in Fig. 1, situated between the 24th and 25th hours of the time axis), and the other spectrum consists of the zero polarization spectra (the “cross–over points” in Fig. 1, situated between the 23th and 24th hours). In Fig. 2 two spectra of different polarization from the first summed image in the region of the H$`\alpha `$ line are shown (bottom). The magnetic shift in H$`\alpha `$ is clearly visible. There are some $`\mathrm{H}_2\mathrm{O}`$ atmospheric absorption lines in the region not showing a shift. The result of subtraction of these spectra , i. e. the polarization (V–parameter) is shown in the middle of Fig. 2. The typical S–wave polarization in the line profile is detected. At the top of Fig. 2 the polarization obtained from the second (“cross–over”) image is shown. No detectable polarization is observed there.
## Observations with the SP-124 spectrograph
In 1997 we changed the observational mode to the medium resolution spectrograph SP–124. This spectrograph was equiped with a PM–CCD (1024x1024) providing spectra of a high quality (Neizvestnyi et al., 1998) and with a new polarimetric analyzer having a better transparency (Bychkov et al., 2000). The analyzer is equiped with a rotatable quarter–wave plate. The spectral resolution in this mode is 5–6 $`\AA `$ (2.3 $`\AA `$/pix). With such a resolution we have lost the possibility of making measurements of the narrow non–LTE H$`\alpha `$ peak, instead we can study two hydrogen lines H$`\alpha `$ and H$`\beta `$ simultaneously, measuring central cores of these lines. Otherwise the method is the same, the left– and right–circular polarized Zeeman spectra are obtained simultaneously on the detector. The main advantage of this mode is that one can obtain Zeeman spectra of WDs with high signal–to–noise because of the better CCD cosmetic and stability (and the lower spectral resolution), besides it provides for recording of two hydrogen lines in the spectrum.
In Table 1 we present new results of magnetic field measurements in 5 DA degenerates in the moderate resolution mode. N is the number of spectra obtained. We found no magnetic field in four stars, but a reliable magnetic field has been detected in WD 1953—011. In a massive Zeeman spectroscopy of DA WDs Schmidt & Smith (1995) found a magnetic field, $`B_e=15.1\pm 6.6`$ kG, of WD 1953–011. This result suggested WD 1953–011 to be a magnetic WD candidate. This star was also suspected as magnetic by Koester et al. (1998) in their high–resolution spectroscopy observations of the NLTE H$`\alpha `$ cores in DA WDs. In this program of searching for rotation in WDs the narrow H$`\alpha `$ cores are fitted with broadened NLTE models, and thus projections of rotation velocities are derived. In normal DA stars the rotation velocities are extremely small with typical upper limits for v sin i of about 15 km/s. In WD 1953–011 a formal v sin i has been found to be 173 $`\pm `$ 10km/s, however a clear Zeeman splitting has been detected in two independent spectra. The splitting corresponds to a surface magnetic field $`B_s93`$ kG (Koester et al., 1998). In our observations the magnetic field of WD 1953–011 was firmly detected, $`B_e=28\pm 6`$ kG. In Fig. 3 the spectrum (the sum of two Zeeman spectra) of WD 1953–011 is shown as well as the circular polarization derived from the Zeeman spectra. The S–wave polarization pattern is clearly seen both in H$`\alpha `$ and H$`\beta `$.
The magnetic nature of WD 1953–011 was thus confirmed in direct Zeeman spectroscopy. Koester et al. (1998) discussed the fact of considerable systematic difference in magnetic fields estimated in high–resolution spectroscopy from H$`\alpha `$ NLTE core splitting and those found in Zeeman spectroscopy of moderate resolution. This could be an effect connected here with different parts of the line profile which are measured by these different methods. However it is hardly probable for a magnetic field to be notably different depending on the star atmosphere height. The longitudinal magnetic field varies with rotational phase over wide limits, from 0 to almost the polar field $`B_p`$, depending on orientation; at the same time the surface (integral) field varies with rotation over much narrower limits. In a model of a dipolar magnetic field with occasional orientation of the dipole axis the surface and longitudinal (effective) magnetic fields are related by a statistical relationship (Angel et al., 1981) $`B_s3B_e`$. In this particular star WD 1953–011 our result, $`B_e28`$ kG, and that by Koester et al. (1998), $`B_s93`$ kG, agree well with the statistical relation.
The known magnetic degenerate star WD 0009+501 was observed on September 1, 1999 in the same mode, the time–resolved Zeeman spectroscopy. WD 0009+501 has been discovered as magnetic by Schmidt & Smith (1994). In 12 observations for 4 observing nights they found a variable magnetic field in the range $`B_e=+9÷100`$ kG with $`\sigma (B_e)=\pm 8÷15`$ kG. The field obviously varies because of rotation of the star. A power spectrum analysis has shown a variety of peaks between the 2-nd and 20th hours, with some preference for shorter periods (Schmidt & Smith 1994). There were about 12 minutes between the spectra in an observational pair, and no variability of magnetic field was detected inside the pairs, but obvious variations were detected in adjacent observations separated by about 5 hours. They observed a nearly full amplitude of variability (about 90 kG) during an interval of 5 hours.
Our observing run of WD 0009+501 consists of 21 continuous Zeeman spectra, the first 10 spectra were taken with a 10 -minut of exposure, and the next 11 spectra with 5 -minut of exposure. Results of measurements of the H$`\alpha `$ line are shown in Fig. 4. The total profile of the line in opposite polarization spectra was measured by a Gauss–analysis. The measurements of H$`\beta `$ are not shown here, they are of about the same behaviour, but less accurate. The observations cover a 2.6 -hours’ interval. The variability of magnetic field is clearly seen in Fig. 4 (filled circles). It was not possible to derive the period accurately from this observational series alone. One can conclude from our data only that the rotational period of WD 0009+501 is about 2 hours. However with addition of the data by Schmidt & Smith (1994) and analyzing these total data together we have found the rotational period, it is 1.83 hours. The same unique period has been identified in their and our power spectra. In Fig. 4 the 10 close separated in time observations by Schmidt & Smith (1994) are shown as open squares. They have been shifted by arbitrary phase of the 1.83–hours period. The curve in the figure is sinusoidal fitting of our data. The average magnetic field is $`<B_e>=42.3\pm 5.4`$ kG, the semi–amplitude is $`32.0\pm 6.8`$ kG.
In Fig. 5 we present both the summed unpolarized spectrum of WD 0009+501 (bottom) consisting of 8 spectra showing maximal magnetic field and the corresponding circularly polarized spectrum (top). The polarization is clearly seen both in H$`\alpha `$ and H$`\beta `$.
### Magnetic field of 40 Eridani B
The brightest DA star 40 Eri B was observed in the SP-124 mode on January 25, 27, and 28 1999 (hereafter — the first, second and third night). The aim was to check the magnetic field variability suspected in our observations in the MSS–mode earlier (see above). A total of 66 Zeeman images were taken (correspondingly 10, 26 and 30 images on these particular nights) with an exposure time of 5 minutes. On the whole the observations lasted for 70, 175 and 192 minutes, respectively on these nights, but they included the time for rotation of the quarter–wave plate and other observational operations. We rotated the plate every 5 (sometimes 10) exposures. This operation is very important for both to take correctly into account many systematic effects connected with nonideal orientation of the slit and the polarimetric analyzer and to project the spectra of opposite polarization onto the same pixels of the detector. Nevertheless, in the case of magnetic variables like our target 40 Eri B is, this procedure alone can not solve the problem of systematic shifts. For this reason we took spectra of standard stars in the same mode just before and after the observations of 40 Eri B. As standards, we used different bright stars which were chosen depending on a current observational program and situation. These are the cool stars $`ϵ`$ Tau, $`\alpha `$ Cep or neighbouring stars. The instrumental line shifts in opposite polarization spectra have been measured in the standards to control any possible systematic shifts with time. In the first star all lines have been measured by the cross–correlation method in restricted ($`100\AA `$) spectral regions around H$`\alpha `$ and H$`\beta `$. Other standard stars are of about the same spectral class as the target or have clearly visible hydrogen lines, so we have measured these two lines in their spectra. One cannot expect effective magnetic fields $`>1`$ kG in bright stars selected by chance (Monin et al., 2000), and all these stars could serve as standards to control our “magnetic zero” variability with an accuracy of $`1`$ kG or better. Instrumental shifts of about 1 –3 kG were really detected in these observations for 2 – 3 hours of observing, and the data were corrected for these shifts by a linear interpolation.
Two hydrogen lines H$`\alpha `$ and H$`\beta `$ have been analysed in 40 Eri B. All the spectra have been normalized with the same window, 180 $`\AA `$. Relative shifts in the circular polarized spectra have been measured by the Gauss–analysis method. Only central line cores of about 30 $`\AA `$ in width were thus measured in the normalized spectra. Both the normalization procedure and the spectral filter parameters were optimized and they were used in a strictly the same way in all the spectra. Not all the spectra we obtained were finally used in magnetic field measurements. The target is low enough at the 6–m telescope latitude (the zenith distance is $`51^{}`$ at the culmination). We observed 40 Eri B around the culmination and tried to obtain the observations as long as possible, nevertheless, depending on weather conditions not all the spectra were obtained with a desired quality ($`\mathrm{S}/\mathrm{N}>80÷90`$ in an individual spectrum of one polarization). We selected the best spectra before the measurements of magnetic field, and only the ones where the hydrogen line cores were not distorted by cosmic ray particles. Particle traces in the steep line cores, even being corrected, may result in distortion of the line shifts.
Results of magnetic field measuremets in H$`\alpha `$ line from individual Zeeman spectra of 40 Eri B are shown in Fig. 6 for the first (top) and the third (bottom) nights of the January 1999 observing run. The crosses indicate the magnetic field of the target, the stars indicate that of the standards. On the first night the stars $`\alpha `$ Cep (before) and $`ϵ`$ Tau (after) were observed as standards. On the third night there were neighbouring stars, a star situated about $`19^{}`$ to north (before) and 40 Eri A (after). The first star was unfortunately not bright enough on the last night to provide a desired accuracy of the zero–field point ($`<1`$ kG). This may result in some uncertainty as a total shift (a linear trend) in calibration of the first points in this series of observations. One can see the variable magnetic field in 40 Eri B detected in H$`\alpha `$.
Fig. 6: The H$`\alpha `$ line magnetic field variability of 40 Eridani B obtained on the first and third nights of observation
Power spectrum analysis of the data from all three nights together, both Fourier and a least–squars search for sinusoidal signals display several peaks from about 2 hours to 5 – 6 hours. These peaks are not independent, they are on a common origin consisting in variable, periodical and probably, not sinusoidal signal. Magnetic phase curves in H$`\alpha `$ and H$`\beta `$ of two best periods, $`2^\mathrm{h}\mathrm{\hspace{0.17em}25}^\mathrm{m}`$ and $`5^\mathrm{h}\mathrm{\hspace{0.17em}17}^\mathrm{m}`$, are displayed in Fig. 7 and 8 respectively. Both the H$`\alpha `$ and the H$`\beta `$ lines confirm the periodical variations, they are about the same both in phase and in amplitude, though the scatter in the H$`\beta `$ data is greater. The individual data are presented by open circles, the mean values averaged in the phase bins $`0.1\mathrm{\Delta }\varphi `$ with their r.m.s. errors are shown by filled circles. We present in these figures also the sinusoidal fitting of the data as $`B_e=A+B\mathrm{sin}(2\pi t/P+\varphi _0)`$. The phase $`\varphi _0`$ was fixed the same in both lines. The parameters of the curves have been found as follows.
t
Fig. 7: The first period magnetic field phase curve for H$`\alpha `$ (top) and H$`\beta `$ (bottom)
Fig. 8: The same as in Fig. 7 for the second suspected
period
$`P=2^\mathrm{h}\mathrm{\hspace{0.17em}25}^\mathrm{m}`$
H$`\alpha `$: $`A=300\pm 550`$ G, $`B=2500\pm 900`$ G,
H$`\beta `$: $`A=400\pm 1000`$ G, $`B=4000\pm 1400`$ G,
$`P=5^\mathrm{h}\mathrm{\hspace{0.17em}17}^\mathrm{m}`$
H$`\alpha `$: $`A=300\pm 550`$ G, $`B=4400\pm 700`$ G,
H$`\beta `$: $`A=400\pm 900`$ G, $`B=4000\pm 1400`$ G.
A number of possible systematic shifts (for instance, the atmospheric dispersion) may influence the results changing the positiones of these two lines in antiphase, because the lines are situated at different edges of the spectral range covered in our spectroscopy. The correlated variability means that systematic effects do not distort the results notably. However 3 points do not agree with the H$`\alpha `$ magnetic curves, but agree well with those in H$`\beta `$. They are marked by the crosses in Figs. 7, 8. They belong to the second night. We can not exclude here the influence of systematic effects.
The setting the equalizing of the zero phases to be equal in H$`\alpha `$ and H$`\beta `$ resulted in the sinusoidal fits parameters are not best. For instance, the H$`\alpha `$ magnetic field amplitude in the first period has been found to be too small in the formal fit. The zero phases are really slightly different in H$`\alpha `$ and H$`\beta `$, and this is due to the individual data scatter. We conclude that the H$`\alpha `$ and H$`\beta `$ data both show the same periodical variations of the effective magnetic field. The semi–amplitude of the variations, $`B_{max}4000÷5000`$ G, and average field is about zero $`\pm 500`$ G. The direct average magnetic field derived from the H$`\alpha `$ data only (without the 3 measurements indicated by crosses in the figures) is $`<B_e>=200\pm 500`$ G. The ratio $`B_{max}/<B_e>`$ is found to be very high. From the data in Fig. 7, 8One can make a conclusion on the star and its magnetic field orientation. If the magnetic field of 40 Eridani B is a central dipole, then the rotational axis inclination to the line of sight is high, $`i90^{}`$, and the magnetic axis inclination to the rotational axis is about the same, $`\beta 90^{}`$.
We may conclude that these new data confirm the previous result obtained in the 1995 observations with MSS, where we found a variable magnetic field in 40 Eri B with a semi–amplitude of $`2300\pm 700`$ G changing on a time–scale of about 4 hours. However we cannot determine firmly a real sole period of the variability. We could select two possible periods in the variability. The first one ($`2^\mathrm{h}\mathrm{\hspace{0.17em}25}^\mathrm{m}`$) was determined with a much better accuracy than the second one ($`5^\mathrm{h}\mathrm{\hspace{0.17em}17}^\mathrm{m}`$). The latter period is longer than the longest observational series on the second and third nights (none of the observations covers the whole 5 hours’ period, they appear as fragments of the magnetic phase curve in Fig. 8). Nevertheless, the 5 hours’ phase curve demonstrates that if the rotational period of 40 Eri B is $`{}_{}{}^{>}5`$ hours, the magnetic curve must be non–sinusoidal (a non–dipolar magnetic field).
The rotation and the magnetic field could be tested together through their impact on the line broadening. The two effects broaden the central core of the H$`\alpha `$ line independently. The surface magnetic field splitting is $`\mathrm{\Delta }\mathrm{V}_{\mathrm{mf}}=1.84\mathrm{B}_\mathrm{s}`$ km/s kG. The rotational broadening is $`\mathrm{\Delta }\mathrm{V}_{\mathrm{rot}}=\mathrm{V}_{\mathrm{rot}}\mathrm{sin}\mathrm{i}=4\pi \mathrm{R}/(\mathrm{P}/\mathrm{sin}\mathrm{i})`$, where R is the radius and P is the period. The total broadening $`\mathrm{\Delta }\mathrm{V}=\sqrt{\mathrm{\Delta }\mathrm{V}_{\mathrm{mf}}^2+\mathrm{\Delta }\mathrm{V}_{\mathrm{rot}}^2}`$ is the observed quantity. In high–resolution spectroscopy the central H$`\alpha `$ core profile being fitted to broadened NLTE models yields an estimate of the velocity $`\mathrm{\Delta }\mathrm{V}`$. Heber et al. (1997) presented such a study of 40 Eri B; they found the broadening of less than $`\mathrm{\Delta }\mathrm{V}<8`$ km/s to be the upper limit at a $`3\sigma `$ level. We show in Fig. 9 the relation between $`\mathrm{B}_\mathrm{s}`$ and $`\mathrm{P}/\mathrm{sin}\mathrm{i}`$ for this star , where the radius $`\mathrm{R}=8.9810^8`$ cm (Reid, 1996) is accepted. Two curves in the figure indicate the total broadening $`\mathrm{\Delta }\mathrm{V}=10`$ and 6 km/s.
The permitted region for the magnetic field and rotation of 40 Eridani B is one below the curves. We present these two curves keeping in mind that the magnetic broadening can be variable with rotational phase. It varies depending on the dipole parameters and on the unknown phase of the suspected (2 –5 hours’) rotational period, when the spectra for the $`v\mathrm{sin}i`$ analysis were taken (Reid, 1996; Heber et al., 1997). The surface magnetic field strength may change by a factor of $`1÷2`$ in the course of rotation. This estimate follows from the study of magnetic Ap stars by Mathys et al. (1997). They collected data on 42 well–studied magnetic Ap stars, their magnetic curves and the mean magnetic field modulus. The ratio $`q=<B>_{max}/<B>_{min}`$ of the observed maximum and minimum values of the mean surface magnetic field modulus varies from star to star over the limits indicated above. About 60 $`\%`$ of the stars have $`q=1.01.1`$, and the other stars have this ratio $`q=1.11.9`$. Taking into account the magnetic broadening as a probable variable contributor to the total line broadening and $`\mathrm{\Delta }\mathrm{V}<8`$ km/s as the observed $`3\sigma `$ upper limit (Heber et al., 1997), we can adopt $`\mathrm{\Delta }\mathrm{V}<10`$ km/s as a very probable upper limit of the H$`\alpha `$ core width in this star. Two vertical lines in Fig. 9 correspond to the rotational periods $`2^\mathrm{h}\mathrm{\hspace{0.17em}25}^\mathrm{m}`$ and $`5^\mathrm{h}\mathrm{\hspace{0.17em}17}^\mathrm{m}`$ (at $`\mathrm{sin}i=1`$), which we discussed above. The permitted region for 40 Eridani B is that both below the upper curve and to the right of the $`2^\mathrm{h}\mathrm{\hspace{0.17em}25}^\mathrm{m}`$ vertical line. We can conclude that the high–accuracy Zeeman observations must be continued to clarify the rotation of the star and its magnetic field curve. Fig. 9 demonstrates the possibilities. In the near future we will be able to know both the rotational period and the magnetic field structure and orientation in 40 Eridani B.
### Spectrum of 40 Eridani B with ultrahigh signal–to–noise
The presence of heavy elements in the atmospheres of very hot DA (H–rich) white dwarfs is well established. These are a group of stars with effective temperatures in excess of 55000 K (Feige 24, G 191–B2B). Being originally discovered with IUE (Bruhweiler & Kondo, 1981), the heavy elements are extensively studied in UV and far UV spectra of the hottest DAs, which show the presence of absorption lines of C, N, O, Si, S, P, Fe and Ni (Sion et al., 1992; Barstow et al, 1993; Vennes et al., 1996). The heavy elements are separated to the photosphere by the radiative levitation. There is evidence for the stratification of Fe in the atmosphere of G 191–B2B (Barstow et. al., 1999), Fe is stratified with increasing abundance at greater depth. Stratification of elements is obtained self–consistently in atmospheric model atmospheres with account for gravitational settling and radiative levitation (Dreizler & Wolf, 1999).
Cool DAZ degenerate stars are believed to accrete interstellar gas, which enriches their photospheres. A few DA stars display Ca, Mg, Fe lines both in the UV and in visible regions. They are G 74–7 (7300 K) (Lacombe et al., 1983; Billeres et al., 1997), G 29–38 (11000 K) (Koester et al., 1997), G 238–44 (20000 K) (Holberg et al., 1997). Some lines (Mg II $`\lambda `$ 4481) originating from the excited lower level demonstrated that they cannot arise in the interstellar gas. Recently Zuckerman & Reid (1998) have observed 38 cool DA stars with HIRES on the Keck I telescope. They have searched for the Ca II K line which could indicate a gas accretion. They have found that the CaII K line was registered in spectra of about 20 $`\%`$ of the stars. In this situation the border between DA and DAZ stars looks rather illusory. It is very important to get more information on the signs of the elements in the atmospheres of white dwarfs both from UV spectra and in visible ones using very high signal–to–noise spectroscopy. In the controversy between line–free DC and DA/DB stars such a spectroscopy has drastically reduced the number of DC stars (Greenstein, 1986).
In the medium–temperature atmospheres $`{}_{}{}^{<}20000`$ K, as it is in 40 Eridani B, the radiative pressure in not so strong as to separate heavy elements and to enrich the photosphere. As it follows from the simulations by Chayers et al. (1995) of expected equilibrium abundances of heavy elements levitating up to the photosphere, in DA stars with a temperatures $`{}_{}{}^{<}20000`$ K one may expect Al/H and Si/H $`{}_{}{}^{<}10_{}^{9}`$, and other elements are less levitating at such low temperatures. The theory predicts all heavy elements in the photospheres of DA white dwarfs have to settle down by gravitational sedimentation, which is a very rapid and effective process. For instance, the settling time for metals at a temperature of 15000 K is about 3–4 days, and it is only 0.6 day for He (Paquette et al., 1986).
Another possibility of supplying the photospheres of DAs with heavy elements is a gas accretion (Alcock & Illarionov, 1980a; 1980b) from interstellar medium (ISM) or stellar winds from a companion in binaries. Accreting mass rate depends strongly on differential velocity $`V`$ between a star and the local ISM, $`\dot{M_a}V^3`$. The distribution of the differencial velocities between DA stars and the ISM is asymmetrical (Holberg et al., 1999), the distribution is nearly uniform between \+ 80 and – 20 km/s. This assumes that the lines measured are related to the stars or the local ISM disturbed by stellar winds. Basing on the central part of the distribution, it can be suggested that the relative velocities may be $`{}_{}{}^{>}\pm 20`$ km/s. The interstellar gas accretion rate is $`\dot{M_a}=4\pi m_H\widehat{n}(GM)^2V^3\mathrm{3.2\hspace{0.17em}10}^{16}\widehat{n}_1m_{0.6}^2v_{20}^3M_{}`$/y, where the star’s mass is $`m_{0.6}=M/0.6M_{}`$, the relative velocity is $`v_{20}=V/20`$ km/s, the number density of the ISM (out of dependence on its ionization state) is $`\widehat{n}_1=\widehat{n}/1\mathrm{c}\mathrm{m}^3`$.
With regard to the interaction of stars with the ISM any star can be in one of the two states: i) the mass loss and ii) the ISM gas accretion. The mass loss of hot stars is much more effective than the radiative pressure in prevention of accretion. However in the case ii) where the star does accrete a gas, under some conditions the accretion may be eventually stopped at the border of magnetosphere and the gas can be expelled by the rotating magnetic field — the propeller (Shvartsman, 1970; 1971; Davidson & Ostriker, 1973; Illarionov & Sunyaev, 1975). Even if a magnetosphere is in the propeller regime, some portions of the gas may penetrate inside the magnetosphere and reach the stellar surface due to various plasma instabilities. A study of tracks of the elements in white dwarfs may help to understand their status of interaction with the ISM.
The photosphere temperature of 40 Eridani B is $`T_e=16500`$ K. The radiative levitation cannot supply the photosphere with heavy elements, but accretion can. If the magnetic field discussed above is dipolar, it can prevent the accretion depending on the relative velocity $`V`$ between the star and the ISM. Greenstein (1980) reported a discovery with IUE an of absorption line near $`\lambda 1391`$ with an equivalent width of 3 Å, which could be Si IV or, wich is most probable, the (0, 5) transition of the H<sub>2</sub> Lyman band. The last interpretation suggests presence of cool gas in the upper atmosphere of the star with a column density $`\mathrm{N}(\mathrm{H}_2)10^{15}\mathrm{cm}^2`$ (Greenstein, 1980). In spite of this the line was not confirmed in the later IUE study (Bruhweiler & Kondo, 1983), the same feature has been reported to exist in other white dwarfs (Wegner, 1982). No stellar photospheric features has been reported by Holberg et al. (1998) in the IUE total co–added spectrum of 40 Eridani B. Only weak interstellar features due to Si II $`\lambda \mathrm{\hspace{0.17em}1260}`$, C II $`\lambda \mathrm{\hspace{0.17em}1334}`$, and O I $`\lambda \mathrm{\hspace{0.17em}1302}`$ have been detected there.
We have a good chance to check the theory and to search for the tracks of elenents in 40 Eridani B in the visible range. The total exposure time over the 3 nights of the Zeeman observations was 5 hours. We shifted the grating by 30 – 50 Å on each night in order to avoid the nonuniform pixel–to–pixel sensitivity of the CCD. The superflat, superbias and superdark images were prepared and standard procedures of spectrum reduction were applied. We obtained a summed unpolarized spectrum of the star (Fig. 10). In the range 5500 – 6000 Å about $`410^6`$ counts obtained in the total spectrum, and about $`1.510^6÷210^6`$ counts in the blue $`<4500`$ Å and red $`>6500`$ Å ranges. The spectrum in the region $`56006200`$ Å is shown in Fig. 11. There are 4 spectra — for the 3-rd, 2-nd, 1-st nights and the total spectrum in the figure from bottom to upwards accordingly. The spectra were smoothed and normalized with a window of about 50 Å, so any information on broader spectral details was lost in the spectra. THe spectral resolution in the summed spectra is comparatively low, 6 Å. Unfortunately the Ca II K line was beyond the spectral range registered (the main goal of the spectroscopy was to study the magnetic field, and hydrogen H$`\alpha `$ line was the most important).
One can see a number of weak absorption lines in the spectrum, and the strongest of them repeatedly appeared in the spectrum of each night. In Fig. 12 we present the final background–subtracted spectrum (bold line) and that with no background subtraction (thin line) in the same spectral range, as well as the corresponding total background spectrum (top) which has been divided into 20 and shifted up for the best visual inspection. A special study of the background subtraction was carried out. We found that any possible errors in the background normalization and subtraction could not change the resulting absorption lines found in the 40 Eridani B. The Moon phases were 0.7 – 0.85 during these observations, and the background contribution in the spectrum was $`1\%`$.
In Fig. 12 are also shown positions of telluric lines (Curcio et al., 1964), they are either positions of the strongest line in the group of lines or positions of groups of lines convolved with our spectral resolution. In spite of the rather low resolution we certainly observe these lines in the spectrum. The expected position of the strongest He I line 5876 Å is also shown. This line intensity is less than 0.001. A preliminary inspection of this range and other ranges of the spectrum has shown that we observe numerous weak lines of Ca I, Na I, Fe I, Fe II, Mg I, Mg II, Si II at a level of intensities of $`0.1÷0.5\%`$. The ultraweak absorption spectrum of the star corresponds to a temperature of 5000 – 8000 , so they are not photospheric features. These lines could be interstellar by origin, but not all of them, because many lines have excited the low level, and Na I $`\mathrm{D}_1,\mathrm{D}_2`$ lines, for example, are not the strongest among other lines observed. An analysis of these absorption features will be made elsewhere. The signal–to–noise in the total spectrum is in the range 1000 – 2000, but the correct value can be found only after detailed identification of the spectrum.
We conclude preliminarily that a rich absorption spectrum of heavy traced elements is present in 40 Eridani B. Assuming that the photospheric He I line are broadened to FWHM $`20`$ Å we find that its equivalent width is $`W_\lambda <10`$ mÅ. In a weak absorption line approximation, $`W_\lambda =\pi e^2\lambda ^2Nf/mc^2`$, the upper limit of the He I line intensity gives the upper limit of the column number density of this element $`N(HeI)<10^{11}\mathrm{cm}^2`$ or He abundance $`N(He)/N(H)<10^7`$. This confirms the theory that He has indeed diffused under the photosphere in the hydrogen–rich degenerate.
The traced absorption spectrum does not agree with the photospheric temperature; this suggests that we have observed either i) the uppermost atmosphere gas in this star, or ii) circumstellar gas, or iii) interstellar gas. We believe that the traced heavy elements are supplied through the gas accretion process, this gas may come both from the interstellar medium and from the comparatively close companion 40 Eri C. The local ISM consists of a few fluxes (Bochkarev, 1990) which move with some dispersion in about similar directions. In the heliocentric frame the average velocity of the fluxes is $`20÷30`$ km/s, and the common direction is $`\alpha 90^{},\delta 0^{}÷10^{}`$. This direction is not far from that to 40 Eridani B. Comparing with the real radial velocity of the star (Reid, 1996), $`V_r=44`$ km/s, we find that the relative velocity is high enough. Direct measurements of the ISM line of sight velocity have been obtained by Holberg et al. (1998), \+ 7 km/s; this gives the differential radial velocity between the star and the ISM as 51 km/s. In a tangential direction the proper motion corresponds to a velocity of 94 km/s. We can find that the total relative velocity of 40 Eridani B and its local ISM is $`V=100120`$ km/s. Accepting a density of the ISM close to the Sun $`\widehat{n}_{0.1}=\widehat{n}/0.1\mathrm{cm}^3`$, find the mass accretion rate $`\dot{M_a}2.210^{19}\widehat{n}_{0.1}m_{0.6}^2v_{100}^3M_{}`$/y.
The companion of 40 Erinadi C is an M4e dwarf separated from it by the $`a=40`$ a.e. orbit. The companion’s wind is captured by the white dwarf, and the capture radius is $`R_c2GM/(V_{orb}^2+V_w^2)`$, where the relative orbital velocity $`V_{orb}5`$ km/s is much less than the probable stellar wind velocity from the companion, $`V_w100`$ km/s. One may expect the white dwarf to accrete the wind gas at a rate $`\dot{M_a}=\dot{M_w}(\pi R_c^2/4\pi a^2)1.610^4m_{0.6}v_{100}^2`$. Assuming the mass loss rate in the wind of 40 Eridani C to be $`\dot{M_w}10^{15}M_{}`$/y, we obtain about the same value of the accretion rate as that found from ISM, $`\dot{M_a}10^{19}M_{}`$/y. Both the circumstellar gas density and its velocity field around 40 Eridani B may change in a rather complex way depending on the companions’ winds interaction and orbital phase.
Using the above discussed magnetic field strength and rotational period we can find whether the accreting gas reaches the stellar surface or not. If the magnetic field is dipolar with the pole strength $`B_p`$, the Alfven radius (radius of stopping of accretion) is $`R_A^{7/2}=B_p^2R^6/\dot{M_a}(8GM)^{1/2}`$, where the stellar radius $`R910^8`$ cm. The co–rotation radius is $`R_c^{1/3}=GMP^2/4\pi ^2`$, where P is the rotational period. The accretion is permitted when $`R_A<R_c`$ or $`\dot{M_a}>\dot{M_{cr}}=8^{1/2}(2\pi )^{7/3}(GM)^{5/3}P^{7/3}B_p^2R^6910^{20}m_{0.6}^{5/3}p_3^{7/3}b_5^2M_{}`$/y, where $`p_3=P/3`$ hours, $`b_5=B_p/5`$ kG.
We find that 40 Eridani B can accrete the gas both from the ISM and from the companions onto the surface. At the same time, because we have found that $`\dot{M_a}_{}^>\dot{M_{cr}}`$, the accreting gas will form a circumstellar rotating envelope in the magnetosphere at a distance $`rR_A410^{11}`$ cm. The circumstellar gas can produce the ultraweak absorption spectrum observed.
* The work was supported by the RFBR grant N 98–02–16554.
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# Optimal Manipulations with Qubits: Universal Quantum Entanglers
## I Introduction
A pure quantum state of two systems $`A`$ and $`B`$ is said to be entangled if it is not a product of a state for $`A`$ and a state for $`B`$. Two systems in an entangled state are correlated, and these correlations are intrinsically quantum mechanical . For example, one must use entangled states in order to produce violations of Bell inequalities or in the test of local realism proposed by Hardy . Entangled states also play a key role in quantum information, in particular they are essential in quantum teleportation and in superdense coding . In quantum computers entanglement is one of the features of quantum mechanics which give these machines their power .
Here we would like to consider the problem of how to produce entanglement. In particular, if we are given particles, or systems, $`A`$ and $`B`$ in the pure states $`|\mathrm{\Psi }__A`$ and $`|\mathrm{\Phi }__B`$, respectively, we would like to produce the state $`(|\mathrm{\Psi }__A|\mathrm{\Phi }__B+|\mathrm{\Phi }__A|\mathrm{\Psi }__B)`$ (up to normalization). Formally, we are looking for the symmetrization map
$`𝒮:|\mathrm{\Psi }|\mathrm{\Phi }(|\mathrm{\Psi }|\mathrm{\Phi }+|\mathrm{\Phi }|\mathrm{\Psi }).`$ (1)
In what follows, where possible we omit explicit subscripts $`A`$ and $`B`$. The order in which the vectors are written in the tensor products implicitly denotes to which system they belong (i.e. the left vector corresponds to the system $`A`$, while the right vector corresponds to the system $`B`$). We assume that the two quantum systems (e.g., qubits) are physically distinguishable. For instance they could be located in different regions of space. The task is to entangle their internal degrees of freedom.
That the symmetrization cannot be done perfectly via a unitary transformation can be shown by the following argument. We consider the case in which $`|\mathrm{\Psi }`$ and $`|\mathrm{\Phi }`$ are both qubits. A perfect transformation would have to transform the basis vectors as
$`|00|v_0`$ $``$ $`|00|v_1`$ (2)
$`|01|v_0`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}(|01+|10)|v_2`$ (3)
$`|10|v_0`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}(|01+|10)|v_3`$ (4)
$`|11|v_0`$ $``$ $`|11|v_4,`$ (5)
where the $`|v_j`$, for $`j=0,4`$, are normalized “machine” vectors, i. e. we assume that the entangler itself has its own degrees of freedom. In addition, it is assumed that the entangler is always initially in the same state, $`|v_0`$. Unitarity requires that $`v_2|v_3=0`$. Now let us consider the case where the input vectors are $`|\mathrm{\Psi }=\alpha |0+\beta |1`$ and $`|\mathrm{\Phi }=|0`$ (i.e. the state of the qubit $`A`$ is unknown, while the qubit $`B`$ is in a known state). The transformation (5) gives us
$$|\mathrm{\Psi }|0\alpha |00|v_1+\frac{\beta }{\sqrt{2}}(|01+|10)|v_3,$$
(6)
whereas what it should produce is a vector proportional to $`|\mathrm{\Psi }|0+|0|\mathrm{\Psi }`$, which in the basis $`|0`$, $`|1`$ reads
$$|\mathrm{\Psi }|0|\mathrm{\Psi }|0+|0|\mathrm{\Psi }=2\alpha |00+\beta (|01+|10).$$
(7)
The vectors in the right-hand sides of Eqs.(6) and (7) are clearly not the same, no matter what choice is made for $`|v_1`$ and $`|v_3`$. Therefore, we need to search for devices which will produce approximate versions of the desired state or will produce this state but with a probability which is less than one.
One way of creating a symmetrized state out of two independent systems is by means of a measurement - that is the two systems are optimally measured and their states are estimated. Based on this estimation a two-particle entangled state is prepared. If we begin with two qubits prepared so that one of the states is known ($`|0`$) while the other is unknown ($`|\mathrm{\Psi }`$), we need only estimate the state of one of the particles and this can be performed with a fidelity equal to 2/3 . The information gained from the optimal measurement is then used in the preparation procedure. This is discussed in Section II A.
We shall present quantum mechanical entangling transformations which generate entangled states with much higher fidelity than can be achieved by measuring the input particles. In Section II B we briefly discuss a probabilistic symmetrization (entanglement) which can be realized via a controlled-SWAP gate. The probability of success in this procedure is input-state dependent. In Section III we present the optimal input-state independent quantum entangler and we also study the inseparabity of the outputs of this entangler. In Section IV we show that the universal-NOT gate can also serve as a very interesting entangling device.
## II State-dependent symmetrization
We shall first look at two examples of processes which produce entangled states, for which the quality of the output depends on the input state. That is, these procedures work better for some states than for others. The first is perhaps the most obvious method, we simply measure the input state. We shall consider a more limited problem in this case, entangling an unknown with a known state. The output state resulting from this procedure is only an approximation to the desired one. The second is a probabilistic method; the output when it is produced is ideal, but the probability of successfully producing it is less than one. In this case we shall consider the full problem of entangling two unknown states.
### A Entanglement via measurement
Our task is to entangle an input qubit in an unknown state with a reference qubit in a known state $`|0`$. That is, we want to realize the symmetrization map $`|0__A|\mathrm{\Psi }__B|\mathrm{\Psi }^{(id)}_{_{AB}}`$ with the output parameterized as
$`|\mathrm{\Psi }^{(id)}_{_{AB}}={\displaystyle \frac{2\mathrm{cos}\frac{\vartheta }{2}|00+\sqrt{2}\mathrm{e}^{i\phi }\mathrm{sin}\frac{\vartheta }{2}|+}{\sqrt{2\left(1+\mathrm{cos}^2\frac{\vartheta }{2}\right)}}},`$ (8)
The approach we will discuss here is as follows: firstly, the unknown single-qubit state $`|\mathrm{\Psi }`$ is measured and then using the information gained thereby an approximate version of the desired output is constructed. In order to specify this procedure in more detail, we must describe what measurement is to be made and how its results will be used to construct the output state. The quality of the output will be determined by calculating the fidelity between the actual output and the desired output. We shall first examine a specific strategy and then find an upper bound on the fidelity for a wide class of measurement-based procedures.
Our first measurement-based scenario can then be realized in the following way. In the case of a single input qubit the optimal way to estimate the state, is to measure it along a randomly chosen direction in the two-dimensional Hilbert space. Therefore, the first step in implementing the measurement-based procedure is choosing a random vector $`|\eta `$, where
$`|\eta =\mathrm{cos}{\displaystyle \frac{\vartheta ^{}}{2}}|0+\mathrm{e}^{i\phi ^{}}\mathrm{sin}{\displaystyle \frac{\vartheta ^{}}{2}}|1,`$ (9)
and measuring $`|\mathrm{\Psi }`$ along it. If the result is positive, then the output is taken to be $`|\mathrm{\Phi }_{_{AB}}`$, and if negative, the output is $`|\stackrel{~}{\mathrm{\Phi }}_{_{AB}}`$, where
$`|\mathrm{\Phi }_{_{AB}}`$ $`=`$ $`{\displaystyle \frac{|0|\eta +|\eta |0}{\sqrt{2\left(1+\mathrm{cos}^2\frac{\vartheta ^{}}{2}\right)}}}`$ (10)
$`=`$ $`{\displaystyle \frac{2\mathrm{cos}\frac{\vartheta ^{}}{2}|00+\sqrt{2}\mathrm{e}^{i\phi ^{}}\mathrm{sin}\frac{\vartheta ^{}}{2}|+}{\sqrt{2\left(1+\mathrm{cos}^2\frac{\vartheta ^{}}{2}\right)}}}`$ (11)
and
$`|\stackrel{~}{\mathrm{\Phi }}_{_{AB}}`$ $`=`$ $`{\displaystyle \frac{|0|\eta ^{}+|\eta ^{}|0}{\sqrt{2\left(1+\mathrm{sin}^2\frac{\vartheta ^{}}{2}\right)}}}`$ (12)
$`=`$ $`{\displaystyle \frac{2\mathrm{e}^{i\phi ^{}}\mathrm{sin}\frac{\vartheta ^{}}{2}|00\sqrt{2}\mathrm{cos}\frac{\vartheta ^{}}{2}|+}{\sqrt{2\left(1+\mathrm{sin}^2\frac{\vartheta ^{}}{2}\right)}}}`$ (13)
where the state $`|\eta ^{}`$ is the state orthogonal to $`|\eta `$,
$`|\eta ^{}=\mathrm{e}^{i\phi ^{}}\mathrm{sin}{\displaystyle \frac{\vartheta ^{}}{2}}|0\mathrm{cos}{\displaystyle \frac{\vartheta ^{}}{2}}|1.`$ (14)
For a particular orientation of the measurement apparatus, i.e. for the particular choice of the state $`|\eta `$ this measurement-based scenario gives the two-qubit output density matrix
$$\rho ^{(out)}(\vartheta ,\phi |\vartheta ^{},\phi ^{})=|\mathrm{\Psi }|\eta |^2|\mathrm{\Phi }\mathrm{\Phi }|+|\mathrm{\Psi }|\eta ^{}|^2|\stackrel{~}{\mathrm{\Phi }}\stackrel{~}{\mathrm{\Phi }}|.$$
(15)
To get the final output density matrix one averages this over all possible choices of the measurement (i.e. over all vectors $`|\eta `$)
$$\rho ^{(out)}(\vartheta ,\phi )=\frac{1}{4\pi }_0^{2\pi }𝑑\phi ^{}_0^\pi 𝑑\vartheta ^{}\mathrm{sin}\vartheta ^{}\rho ^{(out)}(\vartheta ,\phi |\vartheta ^{},\phi ^{}).$$
(16)
Finally, the fidelity can be found by computing the matrix element of this density matrix in the ideal output state, $`|\mathrm{\Psi }^{(id)}_{_{AB}}`$,
$$(\vartheta ,\phi )=\mathrm{\Psi }^{(id)}|\rho ^{(out)}(\vartheta ,\phi )|\mathrm{\Psi }^{(id)}.$$
(17)
This fidelity depends on the input state, and this dependence can be eliminated if we average over all input states
$$\overline{}=𝑑\mathrm{\Omega }(\vartheta ,\phi ).$$
(18)
This is the proper fidelity to use to judge how well our proposed strategy performs if we assume that all input states are equally probable. A more explicit expression for it is
$`\overline{}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}{\displaystyle _0^{2\pi }}𝑑\phi {\displaystyle _0^{2\pi }}𝑑\phi ^{}{\displaystyle _0^\pi }\mathrm{sin}\vartheta d\vartheta {\displaystyle _0^\pi }\mathrm{sin}\vartheta ^{}d\vartheta ^{}`$ (19)
$`\times `$ $`\left[|\eta |\mathrm{\Psi }|^2|\mathrm{\Psi }|\mathrm{\Phi }|^2+|\eta ^{}|\mathrm{\Psi }|^2|\mathrm{\Psi }|\stackrel{~}{\mathrm{\Phi }}|^2\right],`$ (20)
Explicitly evaluating this integral we find
$`\overline{}=54+112(\mathrm{ln}2)^2154.5\mathrm{ln}20.719`$ (21)
which is a bit larger than $`2/3`$, the fidelity of the estimation of a state of a single qubit.
Let us now generalize this procedure. We shall again begin by choosing a random vector $`|\eta `$, but now according to a distribution $`q(\vartheta ^{},\phi ^{})`$, which we shall leave unspecified for now. The output density matrix is taken to be either $`\rho _1(\eta )`$ if the measurement result is positive or $`\rho _0(\eta )`$ if it is negative, where
$$\rho _j(\eta )=d\mathrm{\Omega }^{\prime \prime }p_j(\vartheta ^{\prime \prime },\phi |\vartheta ^{},\phi ^{})|\mathrm{\Gamma }(\vartheta ^{\prime \prime },\phi ^{\prime \prime })\mathrm{\Gamma }(\vartheta ^{\prime \prime },\phi ^{\prime \prime })|,$$
(22)
with $`j=0,1`$, and
$`|\mathrm{\Gamma }(\vartheta ^{\prime \prime },\phi ^{\prime \prime })_{_{AB}}=\mathrm{cos}{\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}|00_{AB}+e^{i\varphi ^{\prime \prime }}\mathrm{sin}{\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}|+_{AB}.`$ (23)
The conditional probabilities $`p_j`$ will also be left unspecified; this allows us to consider a wide class of measurement-based strategies. The output density matrix, for a particular $`|\eta `$ is then
$$\rho (\eta )=|\eta |\mathrm{\Psi }|^2\rho _1(\eta )+|\eta ^{}|\mathrm{\Psi }|^2\rho _0(\eta ).$$
(24)
Averaging over $`|\eta `$ gives us the final output density matrix
$$\rho ^{(out)}(\vartheta ,\phi )=𝑑\mathrm{\Omega }^{}\rho (\eta )q(\vartheta ^{},\phi ^{}),$$
(25)
and the fidelities for a specific input state and averaged over all input states are given by Eqs. (17) and (18), respectively, but with $`\rho ^{(out)}`$ computed from Eq. (25) instead of Eq. (16). In particular we have that
$`\overline{}={\displaystyle 𝑑\mathrm{\Omega }^{}𝑑\mathrm{\Omega }^{\prime \prime }\underset{j=0}{\overset{1}{}}f_j(\vartheta ^{\prime \prime },\phi ^{\prime \prime };\vartheta ^{},\phi ^{})P_j(\vartheta ^{\prime \prime },\phi ^{\prime \prime };\vartheta ^{},\phi ^{})},`$ (26)
where
$$P_j(\vartheta ^{\prime \prime },\phi ^{\prime \prime };\vartheta ^{},\phi ^{})=p_j(\vartheta ^{\prime \prime },\phi ^{\prime \prime }|\vartheta ^{},\phi ^{})q(\vartheta ^{},\phi ^{}),$$
(27)
for $`j=0,1`$, and
$`f_0`$ $`=`$ $`{\displaystyle 𝑑\mathrm{\Omega }\frac{1}{2(1+\mathrm{cos}^2(\vartheta /2))}}|2\mathrm{cos}{\displaystyle \frac{\vartheta }{2}}\mathrm{cos}{\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}`$ (29)
$`+\sqrt{2}e^{i(\phi ^{\prime \prime }\phi )}\mathrm{sin}{\displaystyle \frac{\vartheta }{2}}\mathrm{sin}{\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}|^2|\mathrm{\Psi }|\eta ^{}|^2`$
$`f_1`$ $`=`$ $`{\displaystyle 𝑑\mathrm{\Omega }\frac{1}{2(1+\mathrm{cos}^2(\vartheta /2))}}|2\mathrm{cos}{\displaystyle \frac{\vartheta }{2}}\mathrm{cos}{\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}`$ (31)
$`+\sqrt{2}e^{i(\phi ^{\prime \prime }\phi )}\mathrm{sin}{\displaystyle \frac{\vartheta }{2}}\mathrm{sin}{\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}|^2|\mathrm{\Psi }|\eta |^2.`$
What we can now do is to find an upper bound for the fidelity, $`\overline{}`$, for any distribution of the vector $`|\eta `$ and any prescription for using the result of the measurement along $`|\eta `$ to manufacture the entangled state. We note that for $`j=0,1`$
$$1=𝑑\mathrm{\Omega }^{}𝑑\mathrm{\Omega }^{\prime \prime }P_j(\vartheta ^{\prime \prime },\phi ^{\prime \prime };\vartheta ^{},\phi ^{}),$$
(32)
which implies that
$$\overline{}sup|f_0|+sup|f_1|,$$
(33)
where the supremums are taken over the range $`0\vartheta ^{},\vartheta ^{\prime \prime }\pi `$ and $`0\phi ^{},\phi ^{\prime \prime }<2\pi `$.
Our first task is to find explicit expressions for the functions $`f_0`$ and $`f_1`$. We have that
$`f_0`$ $`=`$ $`d_1\mathrm{cos}^2{\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}\mathrm{sin}^2{\displaystyle \frac{\vartheta ^{}}{2}}+d_2\mathrm{cos}^2{\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}\mathrm{cos}^2{\displaystyle \frac{\vartheta ^{}}{2}}+{\displaystyle \frac{1}{2}}d_2\mathrm{sin}^2{\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}\mathrm{sin}^2{\displaystyle \frac{\vartheta ^{}}{2}}`$ (35)
$`+{\displaystyle \frac{1}{2}}d_3\mathrm{sin}^2{\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}\mathrm{cos}^2{\displaystyle \frac{\vartheta ^{}}{2}}\sqrt{2}d_2\mathrm{cos}(\phi ^{\prime \prime }\phi ^{})\mathrm{cos}{\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}\mathrm{cos}{\displaystyle \frac{\vartheta ^{}}{2}}\mathrm{sin}{\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}\mathrm{sin}{\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}`$
$`f_1`$ $`=`$ $`d_1\mathrm{cos}^2{\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}\mathrm{cos}^2{\displaystyle \frac{\vartheta ^{}}{2}}+d_2\mathrm{cos}^2{\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}\mathrm{sin}^2{\displaystyle \frac{\vartheta ^{}}{2}}+{\displaystyle \frac{1}{2}}d_2\mathrm{sin}^2{\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}\mathrm{cos}^2{\displaystyle \frac{\vartheta ^{}}{2}}`$ (37)
$`+{\displaystyle \frac{1}{2}}d_3\mathrm{sin}^2{\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}\mathrm{sin}^2{\displaystyle \frac{\vartheta ^{}}{2}}+\sqrt{2}d_2\mathrm{cos}(\phi ^{\prime \prime }\phi ^{})\mathrm{cos}{\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}\mathrm{cos}{\displaystyle \frac{\vartheta ^{}}{2}}\mathrm{sin}{\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}\mathrm{sin}{\displaystyle \frac{\vartheta ^{\prime \prime }}{2}}`$
where
$`d_1`$ $`=`$ $`2\mathrm{ln}21`$ (38)
$`d_2`$ $`=`$ $`34\mathrm{ln}2`$ (39)
$`d_3`$ $`=`$ $`8\mathrm{ln}25.`$ (40)
From the above equations it is clear that in order to maximize $`f_0`$ we need to choose $`\phi ^{\prime \prime }\phi ^{}=\pi `$ and to maximize $`f_1`$ we need to choose $`\phi ^{\prime \prime }\phi ^{}=0`$. Making these choices and simplifying the resulting expressions we find that
$`f_0(\vartheta ^{\prime \prime },\pi ;\vartheta ^{},0)`$ $`=`$ $`{\displaystyle \frac{1}{4}}[1+c_1\mathrm{cos}\vartheta ^{\prime \prime }c_2\mathrm{cos}\vartheta ^{\prime \prime }\mathrm{cos}\vartheta ^{}`$ (42)
$`+c_3\mathrm{sin}\vartheta ^{\prime \prime }\mathrm{sin}\vartheta ^{}]`$
$`f_1(\vartheta ^{\prime \prime },0;\vartheta ^{},0)`$ $`=`$ $`{\displaystyle \frac{1}{4}}[1+c_1\mathrm{cos}\vartheta ^{\prime \prime }+c_2\mathrm{cos}\vartheta ^{\prime \prime }\mathrm{cos}\vartheta ^{}`$ (44)
$`+c_3\mathrm{sin}\vartheta ^{\prime \prime }\mathrm{sin}\vartheta ^{}],`$
where
$`c_1`$ $`=`$ $`34\mathrm{ln}2`$ (45)
$`c_2`$ $`=`$ $`12\mathrm{ln}28`$ (46)
$`c_3`$ $`=`$ $`\sqrt{2}(34\mathrm{ln}2).`$ (47)
These functions can now be maximized. The maximum of $`f_0`$ occurs at $`\vartheta ^{}=\pi `$ and $`\vartheta ^{\prime \prime }=0`$, and the maximum of $`f_1`$ occurs when $`\vartheta ^{}=0`$ and $`\vartheta ^{\prime \prime }=0`$. The maximum values of both functions are the same and are approximately equal to $`0.386`$. This implies that the fidelity for this kind of a measurement-based strategy must satisfy
$$\overline{}4\mathrm{ln}220.773.$$
(48)
As we shall see, a method which maintains quantum coherences at all stages of the process can do better than this.
### B Controlled-SWAP gate
We now begin with systems $`A`$ and $`B`$ of the same physical origin. Their pure states are described by vectors in the $`D`$-dimensional Hilbert space $``$, so that both together are described by $``$. Let $`\{|u_j|j=1,\mathrm{}D\}`$ be an orthonormal basis for $``$. System $`A`$ is in the state
$$|\mathrm{\Psi }__A=\underset{j=1}{\overset{D}{}}c_j|u_j__A,$$
(49)
and system $`B`$ is in the state
$$|\mathrm{\Phi }__B=\underset{j=1}{\overset{D}{}}d_j|u_j__B.$$
(50)
Our objective is to produce the (entangled) symmetrized state \[see Eq. (1)\]
$$|\mathrm{\Psi }|\mathrm{\Phi }+|\mathrm{\Phi }|\mathrm{\Psi }=\underset{j=1}{\overset{D}{}}\underset{k=1}{\overset{D}{}}(c_jd_k+c_kd_j)|u_j|u_k,$$
(51)
(here we omit the normalization factor).
Recently Barenco et al. have shown that the entanglement (symmetrization) of the form (1) can be performed when the two input qubits interact via a controlled-SWAP (Fredkin) gate with an ancilla initially prepared in a specific state. The entanglement is achieved when a conditional measurement is performed on the ancilla. Exactly the same scenario can be used not only for qubits but for arbitrary quantum systems. To show this we briefly review the operation of the controlled-SWAP gate.
This gate has three inputs. The first, the control bit, is a qubit. The second and third are for $`D`$-dimensional systems. The control bit is unaffected by the action of the gate. If the control bit is $`|0`$, then the gate does nothing, i.e. the output state is the same as the input state. If the control bit is $`|1`$, then the two $`D`$-dimensional states are swapped. This can be accomplished by the following explicit unitary transformation:
$`|0|u_j|u_k`$ $``$ $`|0|u_j|u_k;`$ (52)
$`|1|u_j|u_k`$ $``$ $`|1|u_k|u_j.`$ (53)
Summarizing, the action of our controlled-SWAP gate is,
$`|0|\mathrm{\Psi }|\mathrm{\Phi }`$ $``$ $`|0|\mathrm{\Psi }|\mathrm{\Phi };`$ (54)
$`|1|\mathrm{\Psi }|\mathrm{\Phi }`$ $``$ $`|1|\mathrm{\Phi }|\mathrm{\Psi }.`$ (55)
We now define the qubit states
$`|v_+={\displaystyle \frac{1}{\sqrt{2}}}(|0+|1);|v_{}={\displaystyle \frac{1}{\sqrt{2}}}(|0|1),`$ (56)
and take the input state of the controlled-SWAP gate to be $`|v_+|\mathrm{\Psi }__A|\mathrm{\Phi }__B`$. Using the SWAP transformation (55) we find that the output state is
$`|\mathrm{\Psi }^{(out)}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|0|\mathrm{\Psi }|\mathrm{\Phi }+|1|\mathrm{\Phi }|\mathrm{\Psi }\right)`$ (57)
$`=`$ $`{\displaystyle \frac{1}{2}}|v_+\left(|\mathrm{\Psi }|\mathrm{\Phi }+|\mathrm{\Phi }|\mathrm{\Psi }\right)`$ (58)
$`+`$ $`{\displaystyle \frac{1}{2}}|v_{}\left(|\mathrm{\Psi }|\mathrm{\Phi }|\mathrm{\Phi }|\mathrm{\Psi }\right).`$ (59)
If we now measure the qubit in the $`|v_\pm `$ basis we obtain the states $`(|\mathrm{\Psi }|\mathrm{\Phi }\pm |\mathrm{\Phi }|\mathrm{\Psi })`$ with probabilities $`(1\pm |\mathrm{\Psi }|\mathrm{\Phi }|^2)/2`$, respectively. As we see the probability of generation of a particular entangled state explicitly depends on the (unknown) states of the two systems. In particular, let us assume we begin with two orthogonal qubits, $`|\mathrm{\Psi }`$ and $`|\mathrm{\Psi }^{}`$. Then either of the maximally entangled state, $`(|\mathrm{\Psi }|\mathrm{\Psi }^{}\pm |\mathrm{\Psi }^{}|\mathrm{\Psi })/\sqrt{2}`$ can prepared with probability $`1/2`$.
We stress that the probability of the success in this entanglement (symmetrization) procedure is input-state dependent. In what follows our task will be to find a “machine” which entangles the input with a constant (i.e. input-state independent) fidelity. This covariance property of the entangler with respect to unitary transformations performed on the input qubits makes the entangler universal.
## III Universal entanglers
Suppose we again consider the problem of constructing a device which will entangle a qubit in an arbitrary unknown state $`|\mathrm{\Psi }=\alpha |0+\beta |1`$ with a qubit in a known, reference state, which we shall take to be the basis state $`|0`$. Before we proceed further we have to specify properties of the entangling map. In fact, we can consider two maps. The symmetrization map
$`𝒮:|0__A|\mathrm{\Psi }__B|\mathrm{\Psi }^{(id)}_{_{AB}}=N_s\left(|\mathrm{\Psi }|0+|0|\mathrm{\Psi }\right),`$ (60)
and the anti-symmetrization map
$`𝒜:|0__A|\mathrm{\Psi }__B|\overline{\mathrm{\Psi }}^{(id)}_{_{AB}}=N_a\left(|\mathrm{\Psi }|0|0|\mathrm{\Psi }\right),`$ (61)
where $`N_{a,s}`$ are corresponding normalization factors. As we have shown in the introduction perfect entanglers for arbitrary unknown states cannot be constructed. So the task of the physically realizable symmetric (anti-symmetric) entangler is to produce outputs as close as possible to the ideally entangled states $`|\mathrm{\Psi }^{(id)}_{_{AB}}`$ ($`|\overline{\mathrm{\Psi }}^{(id)}_{_{AB}}`$). In what follows we will quantify the quality of the performance of the universal entangler with the help of the fidelity
$$:=\mathrm{\Psi }^{(id)}|\rho ^{(out)}|\mathrm{\Psi }^{(id)}.$$
(62)
We shall impose the condition that the value of this fidelity does not depend on the input. The fidelity (62) is a good measure of the accuracy with which the entangler produces the desired output state, but we would also like to evaluate the degree of entanglement of the actual output state. Here, however, we have a problem which is due to the fact that it is still not clear how to quantify the entanglement of a quantum system which is in a mixed state. When a bipartite system is in a pure state, then the von Neumann entropy of subsystems can serve as a measure of entanglement. In the case of impure states more sophisticated measures are required (see for instance ).
In terms of the basis vectors, the input state is $`\alpha |00+\beta |01`$, and the ideal output state in the case of symmetrization is
$$|\mathrm{\Psi }^{(id)}=\frac{\left(2\alpha |00+\sqrt{2}\beta |+\right)}{(4|\alpha |^2+2|\beta |^2)^{1/2}},$$
(63)
while in the case of the anti-symmetrization we have
$$|\overline{\mathrm{\Psi }}^{(id)}=|,$$
(64)
where $`|\pm `$ are symmetric and anti-symmetric Bell states in the given basis
$$|\pm =\frac{1}{\sqrt{2}}(|01\pm |10).$$
(65)
In what follows we will briefly discuss the anti-symmetric entangler and then we will concentrate on the symmetric entangler.
### A Entanglement via anti-symmetrization
Recently Alber studied a quantum entangler which takes as an input a quantum-mechanical system prepared in an unknown pure state $`|\mathrm{\Psi }__A`$ and a reference (known) state (let us say $`|0__A`$ ) and at the output generates a two particle entangled state $`\rho _{_{AB}}^{(out)}`$ which is optimally entangled. Alber imposed two constraints on the output of the universal quantum entangler
$`\mathrm{Tr}__A\left[\rho _{_{AB}}^{(out)}\right]=\mathrm{Tr}__B\left[\rho _{_{AB}}^{(out)}\right]={\displaystyle \frac{𝟙}{D}}`$ (66)
and
$`S\left[\rho _{_{AB}}^{(out)}\right]\mathrm{minimum}.`$ (67)
Where $`D`$ is the dimensionality of the Hilbert space of the system $`A`$ ($`B`$) and $`S`$ is the von Neumann entropy $`S=\mathrm{Tr}\rho \mathrm{ln}\rho `$ associated with a given density operator $`\rho `$. The first condition corresponds to the requirement that the subsystems at the output are in the maximally mixed state while the second conditions guarantees that the whole system is as close as possible to a pure two-particle state. Alber has found the solution for this problem. It turns out that the two-particle state which is produced by the optimal (with respect to the above conditions), universal entangler is independent of the input state $`|\mathrm{\Psi }`$ and is equal to a maximally disordered mixture of all possible anti-symmetric Bell states. In the case of qubits ($`D=2`$) there is only one possible anti-symmetric Bell state $`|`$. That is, Alber’s machine realizes the anti-symmetric entangler. We see that the universality of Alber’s entangler means that all inputs are mapped to a single output (the anti-symmetric Bell state $`|`$), so the ideal output state is a priori known, and one could instead build a device which just prepares the known output state. In the antisymmetric entangler the information initially encoded in the qubit $`A`$ is completely lost. But our task is different, we want to redistribute the initial unknown information encoded in the state of the qubit $`A`$, into the entangled state of two qubits. Therefore we will analyze universal entanglement via symmetrization, because the ideal state (63) directly contains information about the initial state of the qubit $`A`$. In other words, we consider the entangling procedure not only as the way to generate the state with highest possible entanglement but also we require that this state contains as much information about the input(s) as possible.
### B Entanglement via symmetrization
Let us now construct a machine which entangles an unknown state with the known state $`|0`$. Taking into account the basic features of the symmetrization transformation (60) we can assume that the basis vectors transform as
$`|00|v_0`$ $``$ $`|00|w_0+|+|x_0`$ (68)
$`|01|v_0`$ $``$ $`|00|w_1+|+|x_1,`$ (69)
where $`|w_0`$, $`|w_1`$, $`|x_0`$, and $`|x_1`$ are states of the entangler itself. The entangler is initially always prepared in the state $`|v_0`$.
We want to impose the condition that the fidelity between the actual output state and the ideal output state be independent of the state $`|\mathrm{\Psi }`$, but before doing so let us state the restrictions which unitarity places on the machine vectors. These are
$`w_0^2+x_0^2`$ $`=`$ $`1;`$ (70)
$`w_1^2+x_1^2`$ $`=`$ $`1;`$ (71)
$`w_0|w_1+x_0|x_1`$ $`=`$ $`0,`$ (72)
where $`x^2x|x`$. We now calculate the output two-qubit density matrix $`\rho ^{(out)}`$ by using the transformation in Eq. (68) to find the full output density matrix and then tracing out the machine degrees of freedom. We then find the fidelity (62) by taking the matrix element of this density matrix in the ideal output state. Our task is to find the machine vectors $`|x_j`$ and $`|w_j`$ ($`j=0,1`$) such that the fidelity $``$ does not depend on the input state $`|\mathrm{\Psi }`$ and simultaneously is as close as possible to unity.
We find that if we choose $`|x_0`$ to be orthogonal to each of the other machine vectors and $`|w_1`$ to be orthogonal to $`|x_0`$ and $`|w_0`$, then the output fidelity will be independent of the phases of $`\alpha `$ and $`\beta `$. Making these choices we find that
$``$ $`=`$ $`N^1\{2|\alpha |^4w_0^2+|\beta |^4x_1^2+|\alpha |^2|\beta |^2`$ (73)
$`\times `$ $`[\sqrt{2}(w_0|x_1+x_1|w_0)+2w_1^2+x_0^2]\},`$ (74)
where $`N=2|\alpha |^2+|\beta |^2`$.
In order for this expression to be independent of $`|\alpha |`$ and $`|\beta |`$ it is necessary that the expression in the curly brackets be proportional to
$$(2|\alpha |^2+|\beta |^2)(|\alpha |^2+|\beta |^2)=2|\alpha |^4+3|\alpha |^2|\beta |^2+|\beta |^4.$$
(75)
Comparing this expression to Eq. (73) we see that
$`w_0`$ $`=`$ $`x_1`$ (76)
$`3w_0^2`$ $`=`$ $`\sqrt{2}(x_1|w_0+w_0|x_1)+2w_1^2+x_0^2.`$ (77)
If these conditions are satisfied, then the fidelity is simply equal to $`w_0^2`$, so that we want to make this quantity as large as possible. If we now make use of the unitarity conditions and the two equations above, we find that
$$1\frac{2}{3}\sqrt{2}\mathrm{cos}\mu =\frac{1w_0^2}{w_0^2},$$
(78)
where
$$\mathrm{cos}\mu =\frac{x_1|w_0+w_0|x_1}{2w_0^2}.$$
(79)
From Eq. (78) we see that $`w_0^2`$ will be a maximum when $`\mathrm{cos}\mu =1`$, which implies that $`|w_0`$ and $`|x_1`$ are parallel. When this condition is satisfied, we find that
$$=w_0^2=\frac{9+3\sqrt{2}}{14},$$
(80)
which gives $`0.946`$ as the approximate value of the fidelity. This means that the output state $`\rho ^{(out)}`$ is indeed very close to the ideal state, and it should be remembered that this fidelity is the same for all input states.
We can summarize our results for the machine vectors as follows. From the above analysis we see that we can take the machine state space to be three dimensional. Define
$$\mathrm{cos}\theta =\left[\frac{9+3\sqrt{2}}{14}\right]^{1/2};\mathrm{sin}\theta =\left[\frac{53\sqrt{2}}{14}\right]^{1/2},$$
(81)
and let $`\{|v_j|j=1,\mathrm{}3\}`$ be an orthonormal basis for the machine vector space. We then have
$`|w_0`$ $`=`$ $`\mathrm{cos}\theta |v_1`$ (82)
$`|w_1`$ $`=`$ $`\mathrm{sin}\theta |v_2`$ (83)
$`|x_0`$ $`=`$ $`\mathrm{sin}\theta |v_3`$ (84)
$`|x_1`$ $`=`$ $`\mathrm{cos}\theta |v_1,`$ (85)
and our transformation in terms of basis vectors becomes
$`|00|v_0`$ $``$ $`\mathrm{cos}\theta |00|v_1+\mathrm{sin}\theta |+|v_3;`$ (86)
$`|01|v_0`$ $``$ $`\mathrm{sin}\theta |00|v_2+\mathrm{cos}\theta |+|v_1.`$ (87)
By construction this is the optimal entangling transformation which entangles an unknown pure state with a known reference state.
Alternatively, for $`|\mathrm{\Psi }=\alpha |0+\beta |1`$ we can rewrite this transformation in the form
$`|0|\mathrm{\Psi }|v_0`$ $``$ $`\mathrm{cos}\theta (\alpha |00+\beta |+)|v_1`$ (89)
$`+\mathrm{sin}\theta \left(\alpha |+|v_3+\beta |00|v_2\right).`$
When the trace over the entangler is performed we obtain the density operator $`\rho _{_{AB}}^{(out)}`$ describing the two qubits $`A`$ and $`B`$ at the output of the quantum entangler
$`\rho _{AB}^{(out)}`$ $`=`$ $`(|\alpha |^2\mathrm{cos}^2\theta +|\beta |^2\mathrm{sin}^2\theta )|0000|`$ (90)
$`+`$ $`(|\alpha |^2\mathrm{sin}^2\theta +|\beta |^2\mathrm{cos}^2\theta )|++|`$ (91)
$`+`$ $`\mathrm{cos}^2\theta (\alpha \beta ^{}|00+|+\alpha ^{}\beta |+00|)`$ (92)
It is important to stress that the fidelity (62) associated with the output state (89) is input state independent.
### C Remarks
Throughout this paper we have utilized the fidelity (62) as the measure of the performance of the quantum entangler. The universality (covariance) of the entangler is expressed in the fact that the value of the fidelity $``$ is equal for all input states. We note that this covariance constraint is equivalent to the requirement that the Bures distance defined as
$`d_B(\rho _1,\rho _2)=\sqrt{2}\left(1\mathrm{Tr}\sqrt{\widehat{\rho }_1^{1/2}\widehat{\rho }_2\widehat{\rho }_1^{1/2}}\right)^{1/2},`$ (93)
between the ideal state $`|\mathrm{\Psi }^{(id)}`$ and the output of the entangler $`\rho _{_{AB}}^{(out)}`$ is constant. In our particular case we find the Bures distance to be
$`d_B=2\mathrm{sin}(\theta /2)0.0541,`$ (94)
for all inputs. This distance is very small indeed. It is important to note that the Hilbert-Schmidt norm
$`d_{HS}(\rho _1,\rho _2)=\left[\mathrm{Tr}(\rho _1\rho _2)^2\right]^{1/2},`$ (95)
which in our case can be expressed as
$`d_{HS}=\left[12+\mathrm{Tr}\left(\rho _{_{AB}}^{(out)}\right)^2\right]^{1/2},`$ (96)
is not input-state independent because $`\mathrm{Tr}\left(\rho _{_{AB}}^{(out)}\right)^2`$ depends on the initial state. This is closely related to the fact that the von Neumann entropy of the state $`\rho _{_{AB}}^{(out)}`$ is state dependent (see below).
### D Inseparability of the output qubits
We note that the entanglement between the two qubits prepared in the state $`|\mathrm{\Psi }^{(id)}`$ depends on the particular form of the state $`|\mathrm{\Psi }=\alpha |0+\beta |1`$. Because $`|\mathrm{\Psi }^{(id)}`$ is a pure state we can quantify the degree of entanglement via the von Neumann entropy $`S`$ of one of the two qubits under consideration, i.e. $`S_A=\mathrm{Tr}[\rho _A\mathrm{ln}\rho _A]`$ (obviously $`S_A=S_B`$). For $`\alpha =1`$ the entropy is equal to zero, which corresponds to a completely disentangled state (we note that in this case $`|\mathrm{\Psi }^{(id)}=|0|0`$). The entropy takes the maximal value $`S=\mathrm{ln}2`$ for $`\alpha =0`$ when $`|\mathrm{\Psi }^{(id)}=(|0|1+|1|0)/\sqrt{2}`$. We plot this entropy in Fig. 1 (see line 1). The entropy of the individual particle (qubit) at the output of the entangler, i.e. $`\rho __A^{(out)}=\mathrm{Tr}\rho _{_{AB}}^{(out)}`$ is always larger than in the ideal case (see line 2 in Fig. 1). Nevertheless, for the case $`\alpha =0`$ we have in this case $`S(\alpha =0)=0.998\mathrm{ln}2`$, i.e. this entropy is very close to the entropy of a qubit in the ideal case. Unfortunately, this entropy in the case of an impure two-particle state cannot be used as a measure of entanglement.
It is interesting to find the entropy of the two-particle state $`\rho _{_{AB}}^{(out)}`$ at the output of the entangler as a function of the initial state (in the ideal case the two-particle system is always considered to be in a pure state with $`S=0`$). We plot this entropy in Fig. 2. We see that the total entropy of the output is state-dependent and it takes the minimal value for $`\alpha ^2=1/2`$. Therefore the entropy of the subsystems does not indicate whether they are entangled.
We need to check whether the two qubits $`A`$ and $`B`$ at the output are indeed quantum-mechanically entangled. Quantum-mechanical entanglement of two qubits formally means that the density operator of these two qubits is represented by an inseparable matrix (see ). It follows from the Peres-Horodecki theorem that the necessary and sufficient condition of inseparability of the two-qubit density matrix $`\rho _{_{AB}}`$ is that the corresponding partially transposed matrix $`\rho _{_{AB}}^{T_2}`$ has at least one negative eigenvalue.
For instance, let us consider the state $`|\mathrm{\Psi }=\alpha |0+\beta |1`$ with real amplitudes $`\alpha `$ and $`\beta `$. The partially transposed matrix corresponding to the state $`|\mathrm{\Psi }^{(id)}`$ given by Eq.(63) has one negative eigenvalue
$`E(\alpha )={\displaystyle \frac{\alpha ^21}{2(\alpha ^2+1)}}.`$ (97)
We plot this eigenvalue in Fig. 3 (see line 1). We see that the eigenvalue is negative for all values of $`\alpha `$ except $`\alpha =1`$ when $`|\mathrm{\Psi }^{(id)}=|0|0`$. The minimal value of the eigenvalue is achieved for $`\alpha =0`$ when the two qubits are in the maximally entangled state $`(|01+|10)/\sqrt{2}`$.
Now we utilize the Peres-Horodecki theorem to check whether the state $`\rho _{_{AB}}^{(out)}`$ given by Eq. (90) describes an entangled state of two qubits. Firstly, we find that the partially transposed matrix corresponding to the density operator (90) has one eigenvalue which is negative for all values of $`\alpha `$ (here we assume $`\alpha `$ and $`\beta `$ to be real). In particular, this eigenvalue for $`\alpha =0`$ is
$`E(\alpha =0)={\displaystyle \frac{1}{2}}\left[\mathrm{cos}^2\theta \left(\mathrm{cos}^4\theta +\mathrm{sin}^4\theta \right)^{1/2}\right],`$ (98)
which is the minimal value ($`0.447`$) of the negative eigenvalue. On the other hand the maximal value ($`0.001`$) is attained for $`\alpha =1`$
$`E(\alpha =1)={\displaystyle \frac{1}{2}}\left[\mathrm{sin}^2\theta \left(\mathrm{cos}^4\theta +\mathrm{sin}^4\theta \right)^{1/2}\right].`$ (99)
The complete dependence of $`E(\alpha )`$ is shown in Fig. 3. From this figure we clearly see that the output density operator is inseparable for an arbitrary input considered in this Section. We note, that if the entanglement is measured in terms of the tangle as introduced by Wootters then the negative eigenvalues $`E`$ of the partially transposed density operators perfectly reflect the degree of entanglement between the two qubits in our cases.
By construction the fidelity of the entangler in this case is constant but the actual degree of entanglement is state-dependent. This suggests that it would be interesting to find an entangler, whose output states have the same degree of entanglement irrespective of the input, yet still carry information about the input.
## IV Entanglement via universal NOT gate
Even though the negative eigenvalue of the partially transposed density matrix cannot be directly used as the measure of entanglement, we see that the degree of entanglement between two qubits generated in the entangler (87) depends on the input state. In what follows we describe a different type of the entangler, which out of a single qubit $`|\mathrm{\Psi }`$ generates a two-qubit state as close as possible to the state
$`|\mathrm{\Psi }|\{\mathrm{\Psi },\mathrm{\Psi }^{}\}(|\mathrm{\Psi }|\mathrm{\Psi }^{}+|\mathrm{\Psi }^{}|\mathrm{\Psi })/\sqrt{2}.`$ (100)
We will present an entangler which not only produces the state which is as close as possible to the ideal state $`|\{\mathrm{\Psi },\mathrm{\Psi }^{}\}`$ but also has the property that the fidelity does not depend on the input state. In addition, the degree of entanglement also does not depend on the input. This type of the entangler implicitly assumes creation of the state $`|\mathrm{\Psi }^{}`$ from the input $`|\mathrm{\Psi }`$. That is, we face the problem of creating an orthogonal state from unknown input.
It is not a problem to complement a classical bit, i.e. to change the value of a bit, a $`0`$ to a $`1`$ and vice versa. This is accomplished by a NOT gate. Complementing a qubit, however, is another matter. The complement of a qubit $`|\mathrm{\Psi }`$ is the qubit $`|\mathrm{\Psi }^{}`$ which is orthogonal to it. But it is not possible to build a device which will take an arbitrary (unknown) qubit and transform it into the qubit orthogonal to it. As shown in Ref. the ideal universal-NOT (U-NOT) operation corresponds to the inversion of the Bloch (Poincaré) sphere. This inversion preserves angles (related in a simple way to the scalar product $`|\mathrm{\Phi }|\mathrm{\Psi }|`$ of rays), so by Wigner’s Theorem the ideal U-NOT must be implemented either by a unitary or by an anti-unitary operation. Unitary operations correspond to proper rotations of the Poincaré sphere, whereas anti-unitary operations correspond to orthogonal transformations with determinant $`1`$. Clearly, the U-NOT operation is of the latter kind, and an anti-unitary operator $`\mathrm{\Theta }`$ (unique up to a phase) implementing it is
$$\mathrm{\Theta }\left(\alpha |0+\beta |1\right)=\beta ^{}|0\alpha ^{}|1.$$
(101)
The difficulty with anti-unitarily implemented symmetries is that they are not completely positive, i.e., they cannot be applied to a small system, leaving the rest of the world alone.
Because we cannot design a perfect Universal-NOT gate, we have introduced in Ref. an approximate optimal U-NOT gate (an analogous spin-flip operation has recently been introduced by Gisin and Popescu ). This device takes as an input the qubit $`A`$ in the state $`|\mathrm{\Psi }`$ and generates at the output a qubit in a mixed state as close as possible to the orthogonal state $`|\mathrm{\Psi }^{}`$. The role of the U-NOT gate is played by two additional (ancilla) qubits $`B`$ and $`C`$. So, all together the transformation involves three qubits and it can be explicitly written as
$`|\mathrm{\Psi }__A|X_{_{BC}}`$ $``$ $`\gamma _0|\mathrm{\Psi },\mathrm{\Psi }_{_{AB}}|\mathrm{\Psi }^{}__C`$ (103)
$`+\gamma _1|\{\mathrm{\Psi },\mathrm{\Psi }^{}\}_{_{AB}}|\mathrm{\Psi }__C,`$
where $`|X_{_{BC}}`$ is the initial state of the U-NOT gate; $`\gamma _0=\sqrt{2/3}`$ and $`\gamma _1=\sqrt{1/3}`$. In this particular transformation the qubit $`C`$ at the output is in the state which is as orthogonal as possible to the input state. The fidelity of this transformation is input-state independent and is equal to $`=2/3`$.
### A U-NOT as the entangler
It is interesting to note that the two-qubit state $`\rho _{_{AB}}^{(out)}`$ at the output of the U-NOT gate (103) has the form
$`\rho _{_{AB}}^{(out)}=\gamma _1^2|\{\mathrm{\Psi },\mathrm{\Psi }^{}\}\{\mathrm{\Psi },\mathrm{\Psi }^{}\}|+\gamma _0^2|\mathrm{\Psi }\mathrm{\Psi }\mathrm{\Psi }\mathrm{\Psi }|.`$ (104)
The mean fidelity between the state $`\rho _{_{AB}}^{(out)}`$ and the ideal output (100) is input-state independent and takes the value $`=1/3`$. This again corresponds to the fact that the Bures distance between the actual output of the entangler and the ideal output is input state independent and equal to $`d_B=(22/\sqrt{3})^{1/2}`$. We can easily check that the partially transposed matrix corresponding to the density operator (104) has one negative eigenvalue $`E=(2\sqrt{5})/6`$ which is constant and does not depend on the initial input state $`|\mathrm{\Psi }`$.
We note that the Universal NOT gate (103) acts also a quantum cloner, i.e. the two qubits $`A`$ and $`B`$ are the optimal clones of the input (for details see Refs.). It is the optimality of the transformation (103) with respect to cloning and the generation of the optimally orthogonal state (i.e. the universal NOT gate) which indicates that the transformation (103) also serves as the optimal universal entangler.
### B Proof of optimality
Our proof of the optimality of the entangler (100) via the U-NOT gate is based on the recent idea of Gisin that the impossibility of instantaneous signaling generates upper bounds on the fidelity of particular quantum-mechanical processes. To be more specific, we have shown earlier that the impossibility of the ideal (perfect) entangler is due to the linearity of quantum mechanics. On the other hand, another consequence of the linearity of quantum mechanics is the fact that the entangled quantum-mechanical states cannot be used for super-luminal communication. Gisin has shown that this no-signaling constraint implies bounds on the fidelity of universal cloning and the universal U-NOT gate. In the case of cloning the bound on fidelity is $`=5/6`$, while in the case of the U-NOT gate the bound is $`=2/3`$. We note that the transformation (103) achieves both these bounds when used as the cloner or the U-NOT gate, respectively. Recently Alber used this idea of Gisin to prove that the upper bound in the fidelity of the anti-symmetric entangling is equal to unity. The no-signaling constraint can also be used to derive an upper bound on the fidelity of the entangling operation given in Eq. (100) . We will present a proof, which is based on the methods developed in reference , that this upper bound is $`=1/3`$, which means that the U-NOT gate (103) serves as the optimal universal entangler in the sense of Eq. (100).
We consider a process in which a single particle input state is mapped into a two particle output state. The input state can be represented as
$$\rho ^{(in)}(\stackrel{}{m})=\frac{1}{2}(𝟙+\stackrel{}{𝕞}\stackrel{}{\sigma }),$$
(105)
where $`\stackrel{}{m}`$ is a real vector whose length is less than or equal to unity. The most general two-particle output state, which is hermitian and has a trace equal to one, can be expressed as
$`\rho ^{(out)}(\stackrel{}{m})`$ $`=`$ $`{\displaystyle \frac{1}{4}}[𝟙+\stackrel{}{𝕒}\stackrel{}{\sigma }𝟙+𝟙\stackrel{}{𝕓}\stackrel{}{\sigma }`$ (106)
$`+`$ $`{\displaystyle \underset{j,k=x,y,z}{}}t_{jk}\sigma _j\sigma _k],`$ (107)
where $`\stackrel{}{a}`$, $`\stackrel{}{b}`$, and $`t_{jk}`$ are functions of $`\stackrel{}{m}`$. The requirement that the reduced density matrixes of the two output particles be the same, which we shall impose, implies that $`\stackrel{}{a}=\stackrel{}{b}`$.
We now want to impose the requirement of covariance. This means that if $`\rho ^{(in)}(\stackrel{}{m})`$ is mapped onto $`\rho ^{(out)}(\stackrel{}{m})`$, and if $`u`$ is a matrix in $`SU(2)`$, then the input state $`u\rho ^{(in)}(\stackrel{}{m})u^1`$ will be mapped onto the output state $`uu\rho ^{(out)}(\stackrel{}{m})u^1u^1`$. Another way of stating this condition is obtained by noting that if we express $`u`$ as
$$u=\mathrm{exp}(i\theta \widehat{e}\stackrel{}{\sigma }/2),$$
(108)
where $`\widehat{e}`$ is a unit vector corresponding to the rotation axis and $`\theta `$ is the rotation angle, then
$$u(\stackrel{}{m}\stackrel{}{\sigma })u^1=\stackrel{}{m}^{}\stackrel{}{\sigma },$$
(109)
where $`\stackrel{}{m}^{}=R(\widehat{e},\theta )\stackrel{}{m}`$. The rotation matrix, $`R(\widehat{e},\theta )`$, is the $`3\times 3`$ matrix which rotates a vector about the axis $`\widehat{e}`$ by an angle $`\theta `$, and it is given explicitly by
$$R(\widehat{e},\theta )=\mathrm{exp}(\theta \widehat{e}\stackrel{}{K}),$$
(110)
where
$`K_x=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 1\\ 0& 1& 0\end{array}\right),`$ (114)
$`K_y=\left(\begin{array}{ccc}0& 0& 1\\ 0& 0& 0\\ 1& 0& 0\end{array}\right),`$ (118)
$`K_z=\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& 0\end{array}\right).`$ (122)
We have that
$$u\rho ^{(in)}(\stackrel{}{m})u^1=\rho ^{(in)}(R\stackrel{}{m}),$$
(123)
which will be mapped to $`\rho ^{(out)}(R\stackrel{}{m})`$, so that the covariance condition can now be expressed as
$$\rho ^{(out)}(R\stackrel{}{m})=uu\rho ^{(out)}(\stackrel{}{m})u^1u^1.$$
(124)
Now let us examine the consequences of this relation. We shall first consider the terms linear in $`\stackrel{}{\sigma }`$ and let $`R`$ be a rotation about $`\stackrel{}{m}`$ by a very small angle $`\theta `$. We have that
$$\stackrel{}{a}(R\stackrel{}{m})=R\stackrel{}{a}(\stackrel{}{m}),$$
(125)
which for our choice of rotation becomes
$$\stackrel{}{a}(\stackrel{}{m})=(𝟙+\theta \widehat{𝕞}\stackrel{}{𝕂})\stackrel{}{𝕒}(\stackrel{}{𝕞}),$$
(126)
or
$$\widehat{m}\stackrel{}{K}\stackrel{}{a}(\stackrel{}{m})=0,$$
(127)
where $`\widehat{m}`$ is a unit vector in the direction of $`\stackrel{}{m}`$. This implies that $`\widehat{m}\times \stackrel{}{a}(\stackrel{}{m})=0`$, so that $`\stackrel{}{a}(\stackrel{}{m})`$ is parallel to $`\stackrel{}{m}`$, and we can write $`\stackrel{}{a}(\stackrel{}{m})=a(\stackrel{}{m})\stackrel{}{m}`$. If we now substitute this result back into Eq. (125) and consider a general rotation $`R`$, we have that
$$a(R\stackrel{}{m})=a(\stackrel{}{m}).$$
(128)
This implies that $`a(\stackrel{}{m})`$ is a constant, which, following , we shall denote by $`\eta `$.
Now let us see what covariance implies about the terms quadratic in $`\stackrel{}{\sigma }`$. Application of the covariance condition, Eq. (124), to these terms gives
$$t_{jk}(R\stackrel{}{m})=\underset{j^{},k^{}}{}R_{jj^{}}R_{kk^{}}t_{j^{}k^{}}(\stackrel{}{m}).$$
(129)
If we again choose $`R`$ to be a rotation about $`\stackrel{}{m}`$ by a small angle $`\theta `$, we find the condition
$$0=\underset{j^{}}{}(\widehat{m}\stackrel{}{K})_{jj^{}}t_{j^{}k}(\stackrel{}{m})+\underset{k^{}}{}(\widehat{m}\stackrel{}{K})_{kk^{}}t_{jk^{}}(\stackrel{}{m}).$$
(130)
If we choose $`\stackrel{}{m}`$ to be in the $`z`$ direction, in particular $`\stackrel{}{m}=\widehat{z}`$, we find, as did Gisin, that $`t_{xx}=t_{yy}`$, $`t_{xy}=t_{yx}`$, and $`t_{xz}=t_{zx}=t_{yz}=t_{zy}=0`$, where all of these are evaluated at $`\stackrel{}{m}=\widehat{z}`$. We now want to impose the no signaling condition
$$\rho ^{(out)}(\widehat{z})+\rho ^{(out)}(\widehat{z})=\rho ^{(out)}(\widehat{x})+\rho ^{(out)}(\widehat{x}),$$
(131)
and to do so we need to find all of the density matrixes in the above equation in terms of $`t_{jk}(\widehat{z})`$. This can be done by applying the covariance condition, Eq. (124), to $`\rho ^{(out)}(\widehat{z})`$ and making the proper choice of $`R`$. When these results are substituted into Eq. (131) we find that $`t_{xx}(\widehat{z})=t_{yy}(\widehat{z})=t_{zz}(\widehat{z})`$, and we shall designate this common value by $`t(\widehat{z})`$. We then have that
$$\rho ^{(out)}(\widehat{z})=\frac{1}{4}\left(\begin{array}{cccc}1+2\eta +t& 0& 0& 0\\ 0& 1t& 2(t+it_{xy})& 0\\ 0& 2(tit_{xy})& 1t& 0\\ 0& 0& 0& 12\eta +t\end{array}\right).$$
(132)
The basis in which the matrix is expressed is $`\{|+\widehat{z},+\widehat{z},|+\widehat{z},\widehat{z},|\widehat{z},+\widehat{z},|\widehat{z},\widehat{z}\}`$, where $`\sigma _z|\pm \widehat{z}=\pm |\widehat{z}`$. This matrix must be positive, which implies that the eigenvalues
$$\frac{1}{4}(1\pm 2\eta +t);\frac{1}{4}(1t\pm 2\sqrt{t^2+t_{xy}^2})$$
(133)
must be nonnegative.
For an input state $`\rho ^{(in)}(\widehat{z})`$ our desired output state is $`(|+\widehat{z},\widehat{z}+|\widehat{z},+\widehat{z})/\sqrt{2}`$, and this implies that the fidelity of $`\rho ^{(out)}`$ is
$$=\frac{1+t}{4}.$$
(134)
This is clearly maximized when $`t`$ is as large as possible, and examining the eigenvalues of $`\rho ^{(out)}`$, this happens when $`t_{xy}=0`$ and $`t=1/3`$. Substituting this into the expression for the fidelity, we see that the maximum fidelity is $`1/3`$. This means that the no-signaling constraint specifies the upper bound on the fidelity of the symmetric entangling which is exactly the same one as achieved by the U-NOT gate. This proves that the entangling via the U-NOT gate is optimal.
### C Remark
We note that using the universal NOT gate one can also produce an entangled state of the form (60). Specifically, the U-NOT gate allows Charlie (C) to produce an entangled state, consisting of $`|\mathrm{\Psi }`$ and one of two known states, which is shared by Alice (A) and Bob (B). In order to see how this can be accomplished it is useful to express the state on the right-hand side of Eq. (103) as
$`\sqrt{{\displaystyle \frac{1}{3}}}(|\mathrm{\Psi }__A|\mathrm{\Phi }_{}_{_{BC}}+|\mathrm{\Psi }__B|\mathrm{\Phi }_{}_{_{AC}}),`$ (135)
where
$`|\mathrm{\Phi }_{}={\displaystyle \frac{(|\mathrm{\Psi }|\mathrm{\Psi }^{}|\mathrm{\Psi }^{}|\mathrm{\Psi })}{\sqrt{2}}}={\displaystyle \frac{(|0|1|1|0)}{\sqrt{2}}}`$ (136)
is the singlet state. Charlie now measures his particle along the axis corresponding to the states $`|0`$ and $`|1`$. Whatever result he obtains for his particle, the other two particles will be in an entangled state shared by Alice and Bob. For example, if Charlie finds his particle in the state $`|1`$, Alice and Bob share the state in Eq. (60). Note that Charlie can choose the states with which the state $`|\mathrm{\Psi }`$ will be entangled by choosing the axis along which to measure his particle.
This implies that if we want to produce either the entangled state of $`|\mathrm{\Psi }`$ with $`|0`$ or the entangled state of $`|\mathrm{\Psi }`$ with $`|1`$, and we don’t care which one we get, this can be done with perfect fidelity. Perhaps a better way of stating this is that if we want to entangle $`|\mathrm{\Psi }`$ with one of two orthogonal states, this can be done perfectly, and we will know with which state it is entangled.
## V Conclusions
In this paper we have studied various possibilities for entangling two qubits so the initial information about their preparation is preserved. We have studied a specific situation when the state of one of the qubits is known while the second state is arbitrary. We have shown that entanglement via symmetrization in this case can be performed with a very high fidelity (much higher than the fidelity of estimation). This type of entanglement can be very useful for stabilization of the storage of an (unknown) quantum state of one qubit against environmental interaction and a random imprecision . We have shown that the U-NOT gate optimally implements the entanglement transformation $`|\mathrm{\Psi }|\mathrm{\Psi }|\mathrm{\Psi }^{}+|\mathrm{\Psi }^{}|\mathrm{\Psi }`$. This means that the transformation (103) is very special indeed - it describes the optimal cloning, the optimal U-NOT transformation as well as the optimal entangler.
###### Acknowledgements.
We thank Nicolas Gisin and Christoph Simon for helpful discussions. This work was supported by the National Science Foundation under grant PHY-9970507 and by the IST project EQUIP under the contract IST-1999-11053.
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# Type Ia Supernova Explosion Models
## 1 INTRODUCTION
Changes in the appearance of the night sky, visible with the naked eye, have always called for explanations (and speculations). But, although “new stars”, i.e. novae and supernovae, are observed by humans for thousands of years, the modern era of supernova research began only about one century ago on August 31, 1885, when Hartwig discovered a “nova” near the center of the Andromeda galaxy, which became invisible about 18 months later. In 1919 Lundmark estimated the distance of M31 to be about $`7\times 10^5`$ lyr, and by that time it became obvious that Hartwig’s nova had been several 1000 times brighter than a normal nova (Lundmark 1920). It was also Lundmark (1921) who first suggested an association between the supernova observed by Chinese astronomers in 1054 and the Crab nebula.
A similar event as S Andromeda was observed in 1895 in NGC 5253 (“nova” Z Centauri), and this time the “new star” appeared to be 5 times brighter than the entire galaxy, but it was not before 1934 that a clear distinction between classical novae and supernovae was made (Baade & Zwicky 1934). Systematic searches, performed predominantly by Zwicky, lead to the discovery of 54 supernovae in the years up to 1956 and, due to improved observational techniques, 82 further supernovae were discovered in the years from 1958 to 1963, all of course in external galaxies (e.g., Zwicky 1965).
Until 1937 spectrograms of supernovae were very rare, and what was known seemed to be not too different from common novae. This changed with the very bright (m$`{}_{V}{}^{}8.4`$) supernova SN1937c in IC 4182 which had spectral features very different from any object that had been observed before (Popper 1937). All of the other supernovae discovered in the following years showed very little dispersion in their maximum luminosity and their post-maximum spectra looked very similar at a given time. Based on this finding Wilson (1939) and Zwicky (1938a) suggested to use supernovae as distance indicators.
In 1940 it became clear, however, that there exist at least two distinctly different classes of supernovae. SN1940c in NGC 4725 had a spectrum very different from all other previously observed supernovae for which good data were available at that time, leading Minkowski (1940) to introduce the names “Type I” for those with spectra like SN1937c and “Type II” for SN 1940c-like events, representing supernovae without and with Balmer-lines of hydrogen near maximum light.
Whether or no the spectral differences also reflect a different explosion mechanism was not known. In contrast, the scenario originally suggested by Zwicky (1938b), that a supernova occurs as the transition from an ordinary star to a neutron star and gains its energy from the gravitational binding of the newly born compact object, was for many years the only explanation. Hoyle & Fowler (1960) were the first to discover that thermonuclear burning in an electron-degenerate stellar core might trigger an explosion and (possibly) the disruption of the star. Together with the idea that the light curves could be powered by the decay-energy of freshly produced radioactive <sup>56</sup>Ni (Truran, Arnett, & Cameron 1967; Colgate & McKee 1969) this scenario is now the generally accepted one for a sub-class of all Type I supernovae called Type Ia today. It is a bit amusing to note that all supernovae (besides the Crab nebula) on which Zwicky had based his core-collapse hypothesis were in fact of Type Ia and most likely belonged to the other group, whereas the first core-collapse supernova, SN1940c, was observed only about one year after he published his paper.
To be more precise, supernovae which do not show hydrogen lines in their spectra but a strong silicon P Cygni – feature near maximum light are named Type Ia (Wheeler & Harkness 1990). They are believed to be the result of thermonuclear disruptions of white dwarfs, either consisting of carbon and oxygen with a mass close to the Chandrasekhar-mass, or of a low-mass C+O core mantled by a layer of helium, the so-called sub-Chandrasekhar-mass models (see the recent reviews by Woosley (1997b); Woosley & Weaver (1994a, 1994b) and Nomoto et al. (1994b, 1997)). The main arguments in favor of this interpretation include: (i) the apparent lack of neutron stars in some of the historical galactic supernovae (e.g. SN1006, SN1572, SN1604); (ii) the rather homogeneous appearance of this sub-class; (iii) the excellent fits to the light curves, which can be obtained from the simple assumption that a few tenths of a solar mass of <sup>56</sup>Ni is produced during the explosion; (iv) the good agreement with the observed spectra of typical Type Ia supernovae. Several of these observational aspects are discussed in some detail in Section 2, together with their cosmological implications. Questions concerning the nature of the progenitor stars are addressed in Section 3, and models of light curves and spectra are reviewed in Section 5.
But having good arguments in favor of a particular explosion scenario does not mean that this scenario is indeed the right one. Besides that one would like to understand the physics of the explosion, the fact that the increasing amount of data also indicates that there is a certain diversity among the Type Ia supernovae seems to contradict a single class of progenitor stars or a single explosion mechanism. Moreover, the desire for using them as distance indicator makes it necessary to search for possible systematic deviations from uniformity. Here, again, theory can make important contributions. In Section 4, therefore, we discuss the physics of thermonuclear combustion, its implementation into numerical models of exploding white dwarfs, and the results of recent computer simulations. A summary and conclusions follow in Section 6.
## 2 OBSERVATIONS
The efforts to systematically obtain observational data of SNe Ia near and far have gained tremendous momentum in recent years. This is primarily a result of the unequaled potential of SNe Ia to act as “standardizable” candles (Branch & Tammann 1992; Riess, Press, & Kirshner 1996; Hamuy et al. 1995; Tripp 1998) for the measurement of the cosmological expansion rate (Hamuy et al. 1996b; Branch 1998) and its variation with lookback time (Perlmutter et al. 1999; Schmidt et al. 1998; Riess et al. 1998). For theorists, this development presents both a challenge to help understand the correlations among the observables and an opportunity to use the wealth of new data to constrain the zoo of existing explosion models. There exist a number of excellent reviews about SNe Ia observations in general (Filippenko 1997b), their spectral properties (Filippenko 1997a), photometry in the IR and optical bands (Meikle et al. 1996; Meikle et al. 1997), and their use for measuring the Hubble constant (Branch 1998). Recent books that cover a variety of observational and theoretical aspects of type Ia supernovae are Ruiz-Lapuente, Canal, & Isern (1997) and Niemeyer & Truran (2000). Below, we highlight those aspects of SN Ia observations that most directly influence theoretical model building at the current time.
### 2.1 General Properties
The classification of SNe Ia is based on spectroscopic features: the absence of hydrogen absorption lines, distinguishing them from Type II supernovae, and the presence of strong silicon lines in the early and maximum spectrum, classifying them as Types Ia (Wheeler & Harkness 1990).
The spectral properties, absolute magnitudes, and light curve shapes of the majority of SN Ia are remarkably homogeneous, exhibiting only subtle spectroscopic and photometric differences (Branch & Tammann 1992; Hamuy et al. 1996c; Branch 1998). It was believed until recently that approximately 85% of all observed events belong to this class of normal (“Branch-normal”, Branch, Fisher, & Nugent 1993) SNe Ia, represented for example by SNe 1972E, 1981B, 1989B, and 1994D. However, the peculiarity rate can be as high as 30 % as suggested by Li et al. (2000).
The optical spectra of normal SN Ia’s contain neutral and singly ionized lines of Si, Ca, Mg, S, and O at maximum light, indicating that the outer layers of the ejecta are mainly composed of intermediate mass elements (Filippenko 1997b). Permitted Fe II lines dominate the spectra roughly two weeks after maximum when the photosphere begins to penetrate Fe-rich ejecta (Harkness 1991). In the nebular phase of the light curve tail, beginning approximately one month after peak brightness, forbidden Fe II, Fe III, and Co III emission lines become the dominant spectral features (Axelrod 1980). Some Ca II remains observable in absorption even at late times (Filippenko 1997a). The decrease of Co lines (Axelrod 1980) and the relative intensity of Co III and Fe III (Kuchner et al. 1994) give evidence that the light curve tail is powered by radioactive decay of <sup>56</sup>Co (Truran et al. 1967; Colgate & McKee 1969).
The early spectra can be explained by resonant scattering of a thermal continuum with P Cygni-profiles whose absorption component is blue-shifted according to ejecta velocities of up to a few times $`10^4`$ km/s, rapidly decreasing with time in the early phase (Filippenko 1997a). Different lines have different expansion velocities (Patat et al. 1996), suggesting a layered structure of the explosion products.
Photometrically, SN Ia rise to maximum light in the period of approximately 20 days (Riess et al. 1999b) reaching
$$M_\mathrm{B}M_\mathrm{V}19.30\pm 0.03+5\mathrm{log}(H_0/60)$$
(1)
with a dispersion of $`\sigma _M0.3`$ (Hamuy et al. 1996b). It is followed by a first rapid decline of about three magnitudes in a matter of one month. Later, the light curve tail falls off in an exponential manner at a rate of approximately one magnitude per month. In the I-band, normal SNe Ia rise to a second maximum approximately two days after the first maximum (Meikle et al. 1997).
It is especially interesting that the two most abundant elements in the universe, hydrogen and helium, so far have not been unambiguously detected in SN Ia spectra (Filippenko (1997a), but see Meikle et al. (1996) for a possible identification of He) and there are no indications for radio emission of SNe Ia. Cumming et al. (1996) failed to find any signatures of H in the early-time spectrum of SN 1994D and used this fact to constrain the mass accretion rate of of the progenitor wind (Lundqvist & Cumming 1997). The later spectrum of SN 1994D also did not exhibit narrow H$`\alpha `$ features (Filippenko 1997b). Another direct constraint for the progenitor system accretion rate comes from the non-detection of radio emission from SN 1986G (Eck et al. 1995), used by Boffi & Branch (1995) to rule out symbiotic systems as a possible progenitor of this event.
### 2.2 Diversity and Correlations
Early suggestions (Pskovskii 1977; Branch 1981) that the existing inhomogeneities among SN Ia observables are strongly intercorrelated are now established beyond doubt (Hamuy et al. 1996a; Filippenko 1997a). Branch (1998) offers a recent summary of correlations between spectroscopic line strengths, ejecta velocities, colors, peak absolute magnitudes, and light curve shapes. Roughly speaking, SNe Ia appear to be arrangeable in a one-parameter sequence according to explosion strength, wherein the weaker explosions are less luminous, redder, and have a faster declining light curve and slower ejecta velocities than the more energetic events (Branch 1998). The relation between the width of the light curve around maximum and the peak brightness is the most prominent of all correlations (Pskovskii 1977; Phillips 1993). Parameterized either by the decline rate $`\mathrm{\Delta }m_{15}`$ (Phillips 1993; Hamuy et al. 1996a), a “stretch parameter” (Perlmutter et al. 1997), or a multi-parameter nonlinear fit in multiple colors (Riess et al. 1996), it was used to renormalize the peak magnitudes of a variety of observed events, substantially reducing the dispersion of absolute brightnesses (Riess et al. 1996; Tripp 1998). This correction procedure is a central ingredient of all current cosmological surveys that use SNe Ia as distance indicators (Perlmutter et al. 1999; Schmidt et al. 1998).
SN 1991bg and SN 1992K are well-studied examples for red, fast, and subluminous supernovae (Filippenko et al. 1992a; Leibundgut et al. 1993; Hamuy et al. 1994; Turatto et al. 1996). Their V, I, and R-band light curve declined unusually quickly, skipping the second maximum in I, and their spectrum showed a high abundance of intermediate mass elements (including Ti II) with low expansion velocities but only little iron (Filippenko et al. 1992a). Models for the nebular spectra and light curve of SN 1991bg consistently imply that the total mass of <sup>56</sup>Ni in the ejecta was very low ($`0.07`$ M) (Mazzali et al. 1997a). On the other side of the luminosity function, SN 1991T is typically mentioned as the most striking representative of bright, energetic events with broad light curves (Phillips et al. 1992; Jeffery et al. 1992; Filippenko et al. 1992b; Ruiz-Lapuente et al. 1992; Spyromilio et al. 1992). Rather than the expected Si II and Ca II, its early spectrum displayed high-excitation lines of Fe III but returned to normal a few months after maximum (Filippenko et al. 1992b).
Peculiar events like SN 1991T and SN 1991bg were suggested to belong to different subgroups of SNe Ia than the normal majority, created by different explosion mechanisms (Mazzali et al. 1997a; Filippenko et al. 1992b; Fisher et al. 1999). The overall SN Ia luminosity function seems to be very steep on the bright end (Vaughan et al. 1995), indicating that “normal” events are essentially the brightest while the full class may contain a large number of undetected subluminous SNe Ia (Livio 2000). New results (Li et al. 2000) indicate, however, that the luminosity function may be shallower than anticipated.
There is also mounting evidence that SN Ia observables are correlated with the host stellar population (Branch 1998). SNe Ia in red or early-type galaxies show, on average, slower ejecta velocities, faster light curves, and are dimmer by $`0.2`$ to 0.3 mag than those in blue or late-type galaxies (Hamuy et al. 1995, 1996a; Branch, Romanishin, & Baron 1996). The SN Ia rate per unit luminosity is nearly a factor of two higher in late-type galaxies than in early-type ones (Cappellaro et al. 1997). In addition, the outer regions of spirals appear to give rise to similarly dim SNe Ia as ellipticals whereas the inner regions harbor a wider variety of explosion strengths (Wang, Höflich, & Wheeler 1997). When corrected for the difference in light curve shape, the variation of absolute magnitudes with galaxy type vanishes along with the dispersion of the former. This fact is crucial for cosmological SN Ia surveys, making the variations with stellar population consistent with the assumption of a single explosion strength parameter (Perlmutter et al. 1999; Riess et al. 1998).
### 2.3 Nearby and Distant SNe Ia
Following a long and successful tradition of using relatively nearby ($`z0.1`$, comprised mostly of the sample discovered by the Calán/Tololo survey (Hamuy et al. 1996a)) SNe Ia for determining the Hubble constant (Branch 1998), the field of SN Ia cosmology has recently seen a lot of activity expanding the range of observed events out to larger redshift, $`z1`$. Systematic searches involving a series of wide-field images taken at epochs separated by 3-4 weeks, in addition to pre-scheduled follow-up observations to obtain detailed spectroscopy and photometry of selected events, have allowed two independent groups of observers – the Supernova Cosmology Project (SCP) (Perlmutter et al. 1999) and the High-$`z`$ Supernova Search Team (Schmidt et al. 1998) – to collect data of more than 50 high-redshift SNe. Extending the Hubble diagram out to $`z1`$ one can, given a sufficient number of data points over a wide range of $`z`$, determine the density parameters for matter and cosmological constant, $`\mathrm{\Omega }_\mathrm{M}`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$, independently (Goobar & Perlmutter 1995) or, in other words, constrain the equation of state of the universe (Garnavich et al. 1998). Both groups come to a spectacular conclusion (Riess et al. 1998; Perlmutter et al. 1999): The distant SNe are too dim by $`0.25`$ mag to be consistent with a purely matter dominated, flat or open FRW universe. Interpreted as being a consequence of a larger than expected distance, this discrepancy can be resolved only if $`\mathrm{\Omega }_\mathrm{\Lambda }`$ is non-zero, implying the existence of an energy component with negative pressure. In fact, the SN Ia data is consistent with a spatially flat universe made up of two parts vacuum energy and one part matter.
Both groups discuss in detail the precautions that were taken to avoid systematic contaminations of the detection of cosmological acceleration, including SN Ia evolution, extinction, and demagnification by gravitational lensing. All of these effects would, in all but the most contrived scenarios, give rise to an increasing deviation from the $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$-case for higher redshift, while the effect of a non-zero cosmological constant should become less significant as $`z`$ grows. Thus, the degeneracy between a systematic overestimation of the intrinsic SN Ia luminosity and cosmological acceleration can be broken when sufficiently many events at $`z0.85`$ are observed (Filippenko & Riess 2000). Meanwhile, the only way to support the cosmological interpretation is by “…adding to the list of ways in which they are similar while failing to discern any way in which they are different” (Riess et al. 1999a). This program has been successful until recently: The list of similarities between nearby and distant SNe Ia includes spectra near maximum brightness (Riess et al. 1998) and the distributions of brightness differences, light curve correction factors, and $`BV`$ color excesses of both samples (Perlmutter et al. 1999). Moreover, while the nearby sample covers a range of stellar populations similar to the one expected out to $`z1`$, a separation of the low-$`z`$ data into sub-samples arising from different progenitor populations shows no systematic shift of the distance estimates (Filippenko & Riess 2000). However, a recent comparison of the rise times of more than 20 nearby SNe (Riess et al. 1999b) with those determined for the SCP high-redshift events gives preliminary evidence for a difference of roughly 2.5 days. This result was disputed by Aldering, Knop, & Nugent (2000) who conclude that the rise times of local and distant supernovae are statistically consistent.
### 2.4 Summary: Observational Requirements for Explosion Models
To summarize the main observational constraints, any viable scenario for the SN Ia explosion mechanism has to satisfy the following (necessary but probably not sufficient) requirements:
1. Agreement of the ejecta composition and velocity with observed spectra and light curves. In general, the explosion must be sufficiently powerful (i.e., produce enough <sup>56</sup>Ni) and produce a substantial amount of high-velocity intermediate mass elements in the outer layers. Furthermore, the isotopic abundances of “normal” SNe Ia must not deviate significantly from those found in the solar system.
2. Robustness of the explosion mechanism. In order to account for the homogeneity of normal SNe Ia, the standard model should not give rise to widely different outcomes depending on the fine-tuning of model parameters or initial conditions.
3. Intrinsic variability. While the basic model should be robust with respect to small fluctuations, it must contain at least one parameter that can plausibly account for the observed sequence of explosion strengths.
4. Correlation with progenitor system. The explosion strength parameter must be causally connected with the state of the progenitor white dwarf in order to explain the observed variations as a function of the host stellar population.
## 3 LIGHT CURVE AND SPECTRA MODELING
Next we have to discuss the problem of coupling the interior physics of an exploding white dwarf to what is finally observed, namely light curves and spectra, by means of radiative transfer calculations. For many astrophysical applications this problem is not solved, and SN Ia are no exceptions. In fact, radiation transport is even more complex in Type Ia’s than for most other cases.
A rough sketch of the processes involved can illustrate some of the difficulties (see, e.g., Mazzali & Lucy (1993); Eastman & Pinto (1993)). Unlike most other objects we know in astrophysics SN Ia do not contain any hydrogen. Therefore the opacities are always dominated either by electron scattering (in the optical) or by a huge number of atomic lines (in the UV). In the beginning, the supernova is an opaque expanding sphere of matter into which energy is injected from radioactive decay. This could happen in a very inhomogeneous manner, as will discussed later. As the matter expands diffusion times eventually get shorter than the expansion time and the supernova becomes visual. However, because the star is rapidly expanding the Doppler-shift of atomic lines causes important effects. For example, a photon emitted somewhere in the supernova may find the surrounding matter more or less transparent until it finds a line Doppler-shifted such that it is trapped in that line and scatters many times. As a consequence, the spectrum might look thermal although the photon “temperature” has nothing in common with the matter temperature.
It is also obvious that radiation transport in SN Ia is very non-local and that the methods used commonly in models of stellar atmospheres need refinements. As a consequence, there is no agreement yet among the groups modeling light-curves and spectra as to what the best approach is. Therefore it can happen that even if the same model for the interior physics of the supernova is inserted into one of the existing codes for modeling light-curves and spectra the predictions for what should be “observed” could be different, again a very unpleasant situation. Things get even worse because all such models treat the exploding star as being spherically symmetric, an assumption that is at least questionable, given the complex combustion physics discussed below.
In the following subsections we outline some of the commonly used numerical techniques and also discuss their predictions for SN Ia spectra and light curves. For more details on the techniques used by the various groups, we refer readers to the articles by Eastman (1997); Blinnikov (1997); Pinto (1997); Baron, Hauschildt, & Mezzacappa (1997); Mazzali et al. (1997b); Höflich et al. (1997), and Ruiz-Lapuente (1997).
### 3.1 Radiative transfer in Type Ia supernovae
In principle, the equations which have to be solved are well-known, either in form of the Boltzmann transport equation for photons or as a transport equation for the monochromatic intensities. However, to solve this time dependent, frequency dependent radiation transport problem, including the need of treating the atoms in non-LTE, is very expensive, even in spherically symmetric situations. Therefore, approximations of various kinds are usually made which give rise to controversial discussions.
Conceptually, it is best to formulate and solve the transport equation in the co-moving (Lagrangian) frame (cf. Mihalas & Weibel Mihalas (1984)). This makes the transport equation appear simpler, but causes problems in calculating the “co-moving” opacity, in particular if the effect of spectral lines on the opacity of an expanding shell of matter is important, as in the case of SN Ia (Karp et al. 1977).
There are different ways to construct approximate solutions of the transport equation. One can integrate over frequency and replace the opacity terms by appropriate means, leaving a single (averaged) transport equation. Unfortunately, in order to compute the flux-mean opacity one has to know the solution of the transport equation. Frequently the flux-mean is replaced, for example, by the Rosseland mean, allowing for solutions, but at the expense of consistency (see, e.g., Eastman (1997)).
Another way out is to replace the transport equation by its moment expansion introducing, however, the problem of closure. In its simplest form, the diffusion approximation, the radiation field is assumed to be isotropic, the time rate of change of the flux is ignored, and the flux is expressed in terms of the gradient of the mean intensity of the radiation field. Replacing the mean intensity by the Planck function and closing the moment expansion by relating the radiation energy density and pressure via an Eddington factor (equal to 1/3 for isotropic radiation) finally leads to a set of equations that can be solved (Mihalas & Weibel Mihalas 1984).
Again, this simple approach has several short-comings that are obvious. First of all, the transition from an optically thick to thin medium at the photosphere requires a special treatment mainly because the radiation field is no longer isotropic. One can compensate for this effect by either putting in a flux-limiter or a variable Eddington factor to describe the transition from diffusion to free streaming, but both approaches are not fully satisfactory since it is difficult to calibrate the newly invented parameters (e.g., Kunasz (1984); Fu (1987); Blinnikov & Nadyoshin (1991); Mair et al. (1992); Stone, Mihalas, & Norman (1992); Yin & Miller (1995)).
Alternatively, one can bin frequency space into groups and solve the set of fully time dependent coupled monochromatic transport equations for each bin. In this approach the problem remains to compute average opacities for each frequency bin. Moreover, because of computer limitations, in all practical applications the number of bins cannot be large which introduces considerable errors, given the strong frequency dependence of the line-opacities (Blinnikov & Nadyoshin 1991; Eastman 1997) (see also Fig. 1).
Finally, in order to get synthetic spectra one might apply Monte-Carlo techniques, as was done by Mazzali et al. (1997b) and Lucy (1999). Here the assumption is that the supernova envelope is in homologous spherical expansion and that the luminosity and the photospheric radius are given. The formation of spectral lines is then computed by considering the propagation of a wave packet emitted from the photosphere subject to electron scattering and interaction with lines. Line formation is assumed to occur by coherent scattering, and the line profiles and escape probabilities are calculated in the Sobolev approximation. While this approach appears to be a powerful tool to get synthetic spectra it lacks consistency since the properties of the photosphere have to be calculated be other means.
But having a numerical scheme at hand to solve the transport equation is not sufficient. It is even more important to have accurate opacities. The basic problem, namely that at short wavelengths the opacity is dominated by a huge number of weak lines, was mentioned before. In practice this means that because the list included in anyone’s code is certainly incomplete and the available information may not always be very accurate it is difficult to estimate possible errors. Moreover, there is no general agreement among the different groups calculating SN Ia lightcurves and spectra on how to correct the opacities for Doppler-shifts of the lines, caused by the expansion of the supernova. The so-called “expansion opacity” (see Fig. 1) that should be used in approaches based on the diffusion equation as well as on moment expansions of the transport equation is still discussed in a controversial manner (Pauldrach et al. 1996; Blinnikov 1997; Baron et al. 1997; Eastman 1997; Höflich et al. 1997; Mazzali et al. 1997b; Pinto 1997).
Other open questions include the relative importance of absorption and scattering of photons in lines, and whether or not one can calculate the occupation numbers of atomic levels in equilibrium (LTE) or has to do it by means of the Saha-equation (NLTE) (Pauldrach et al. 1996; Nugent et al. 1997; Höflich, Wheeler, & Thielemann 1998b).
### 3.2 Results of numerical studies
Despite of the problems discussed in the previous Subsection radiation hydrodynamic models have been used widely as a diagnostic tool for SN Ia. These studies include computations of $`\gamma `$-ray (Burrows & The 1990; Müller, Höflich, & Khokhlov 1991; Burrows, Shankar, & van Riper 1991; Shigeyama et al. 1993; Ruiz-Lapuente et al. 1993b; Timmes & Woosley 1997; Höflich, Wheeler, & Khokhlov 1998a; Watanabe et al. 1999), UV and optical (Branch & Venkatakrishna 1986; Ruiz-Lapuente et al. 1992; Nugent et al. 1995; Pauldrach et al. 1996; Höflich et al. 1997; Nugent et al. 1997; Hatano, Branch, & Baron 1998; Höflich et al. 1998b; Lentz et al. 2000; Fisher et al. 1999; Lucy 1999; Lentz et al. 1999), and of infrared lightcurves and spectra (Spyromilio, Pinto, & Eastman 1994; Höflich 1995; Wheeler et al. 1998). All studies are based on the assumption, that the explosion remains on average spherically symmetric, an assumption which is questionable, as will be discussed in Sect. 5. Although spherical symmetry might be a good approximation for temperatures, densities, and velocities, the spatial distribution of the products of explosive nuclear burning is expected to be very non-spherical, and it is the distribution of the heavier elements, both in real and in velocity space, which determines to a large extent lightcurves and spectra.
With the possible exception of SN 1991T, where a 2 – 3$`\sigma `$ detection of the <sup>56</sup>Co decay-lines at 847keV and 1238keV has been reported (Morris et al. 1997) (see, however, Leising et al. (1995)), only upper limits on $`\gamma `$-ray line-emission from SN Ia are known. On the basis of the models this is not surprising since the flux limits of detectors such as COMPTEL on GRO ( 10<sup>-5</sup> photons per cm<sup>2</sup> and second (Schönfelder et al. 1996)) allows detections out to distances of about 15 Mpc in the most favorable cases, i.e. delayed-detonation models producing lots of <sup>56</sup>Ni in the outer parts of the supernova (Timmes & Woosley 1997). In fact, the tentative detection of decay-lines from SN 1991T at a distance of about 13 Mpc can be explained by certain delayed-detonation models and was even predicted by some of them (Müller et al. (1991); see also Sect. 5).
Synthetic (optical and UV) spectra of hydrodynamic models of SN Ia have been computed by several groups (Höflich & Khokhlov 1996; Nugent et al. 1997) and have been compared with the observations. The bottom line of these investigations is that Chandrasekhar-mass deflagration models are in good agreement with observations of “Branch-normals” such as SN 1992A and SN 1994D (Höflich & Khokhlov 1996; Nugent et al. 1997), and delayed-detonation are equally good. The reason is that in both classes of models the burning front starts by propagating out slowly, giving the star some time to expand. The front then speeds up to higher velocities, e.g. to a fair fraction of the sound velocity for deflagration models and to supersonic velocity for detonations, which is necessary to match the obseved high velocities of the ejecta. But as far as the amount of radioactive Ni is concerned, the predictions of both classes of models are not too different (Nugent et al. 1997). It also appears that sub-Chandrasekhar models cannot explain the observed UV-flux and the colors of normal SNe Ia (Khokhlov, Müller, & Höflich 1993). Moreover, although sub-Chandrasekhar models eject considerable amounts of He, according to the synthetic spectra He-lines should not be seen, eliminating them as a tool to distinguish between the models (Nugent et al. 1997).
In the infrared SN Ia do show non-monotonic behavior (Elias et al. 1985) and, as for the bolometric lightcurves, a correlation between peak-brightness and lightcurve-shape seems to exist (Contardo & Leibundgut 1998; Contardo 1999). Therefore calculations of IR lightcurves and spectra are of importance and they might prove to be a good diagnostic tool. Broad-band IR lightcurves have been computed by Höflich, Khokhlov, & Wheeler (1995) with the result that the second IR-peak can be explained as an opacity effect. Although the fits were not perfect the general behavior, again, was consistent with both, the deflagration and the delayed-detonation models. Detailed early IR-spectra have been calculated only recently (Wheeler et al. 1998) and the models provide a good physical understanding of the spectra. Again, a comparison between several of the delayed-detonation models and SN 1994D gave good agreement, but one might suspect that certain deflagrations would do equally well. However, in principle, synthetic IR spectra are sensitive to the boundary between explosive C and O and between complete and incomplete Si burning (Wheeler et al. 1998) and should provide some information on the progenitors and the explosion mechanism.
In conclusion, models of SN Ia lightcurves and spectra can fit the observations well but, so far, their predictive power is limited. The fact that multi-dimensional effects are ignored and that the opacities as well as the radiative-transfer codes have obvious shortcomings make it difficult to derive strong constraints on the explosion mechanism. It appears, however, that while it seems to be difficult to distinguish between pure deflagrations and delayed-detonations on the basis of synthetic lightcurves and spectra, sub-Chandrasekhar models cannot fit normal SN Ia equally well.
## 4 PROGENITOR SYSTEMS
In contrast to supernovae from collapsing massive stars for which in two cases the progenitor star was identified and some of its properties could be inferred directly from observations (SN1987A in the LMC (Blanco 1987; Gilmozzi 1987; Gilmozzi et al. 1987; Hillebrandt et al. 1987) and SN1993J in M81 (Benson et al. 1993; Schmidt et al. 1993; Nomoto et al. 1993; Podsiadlowski 1993)), there is not a single case known where we have this kind of information for the progenitor of a SN Ia. This is not too surprising, given the fact that their progenitors are most likely faint compact dwarf stars and not red or blue supergiants. Therefore we have to rely on indirect means if we want to determine their nature.
The standard procedure is then to eliminate all potential candidates if some of their properties disagree with either observations or physical principles, and to hope that one is left with a single and unique solution. Unfortunately, for the progenitors of Type Ia supernovae this cannot be done unambiguously, the problem being the lack of strong candidates that pass all possible tests beyond doubt.
In this Section we will first repeat the major constraints which have to be imposed on the progenitor systems and then discuss the presently favored candidates, Chandrasekhar-mass C+O white dwarfs and low-mass C+O white dwarf cores embedded in a shell of helium, in some detail. It will be shown, however, that even if we could single out a particular progenitor system this would narrow the parameter space for the initial conditions at the onset of the explosion, but might not determine them sufficiently well, in particular if we are aiming at a quantitative understanding. Some of the discussion given below follows recent reviews of Renzini (1996) and Livio (2000).
### 4.1 Observational constraints on Type Ia progenitors
As was already discussed in Sect. 2, SNe Ia are (spectroscopically) defined by the absence of emission lines of hydrogen and the presence of a (blue-shifted) Si II absorption line with a rest-wavelength of 6355Å near maximum light. The first finding requires that the atmosphere of the exploding star contains no or at most 0.1 M of hydrogen, and the second one indicates that some nuclear processing takes place and that products of nuclear burning are ejected in the explosion. Mean velocities of the ejecta, as inferred from spectral fits, are around 5,000 km/s and peak velocities exceeding 20,000 km/s are observed, which is consistent with fusing about 1 M of carbon and oxygen into Fe-group elements or intermediate-mass elements such as Si or Ca. The presence of some UV-flux, the width of the peak of the early light curve, and the fact that radioactive-decay models (<sup>56</sup>Ni $``$ <sup>56</sup>Co $``$ <sup>56</sup>Fe) can fit the emission very well, all point towards compact progenitor stars with radii of less than about 10,000 km.
After about two weeks the typical SN Ia spectrum changes from being dominated by lines of intermediate-mass nuclei to being dominated by Fe II. Since also a Co III feature is identified at later stages this adds evidence to the interpretation that they are indeed thermonuclear explosions of rather compact stars, leaving the cores of stars with main sequence masses near 6 to 8 M or white dwarfs as potential candidates. Moreover, the energetics of the explosion and the spectra seem to exclude He white dwarfs (Nomoto & Sugimoto 1977; Woosley, Taam, & Weaver 1986), mainly because such white dwarfs would undergo very violent detonations.
Next one notes that most SNe Ia, of order 85 %, have very similar peak luminosities, light curves, and spectra. The dispersion in peak blue and visual brightness is only of order 0.2 - 0.3 magnitudes calling for a very homogeneous class of progenitors. It is mainly this observational fact that seems to single out Chandrasekhar-mass white dwarfs as their progenitors. Since the ratio of energy to mass determines the velocity profile of the exploding star the homogeneity would be explained in a very natural way. However, as has been discussed in Sect. (2), there exist also significant differences among the various SNe Ia which may indicate that this simple interpretation is not fully correct. The difference in peak-brightness, ranging from sub-luminous events like SN 1991bg in NGC 4374 ($`B_{max}`$ = -16.54 (Turatto et al. 1996), as compared to the mean of the “Branch-normals” of $`B_{max}`$ -19 (Hamuy et al. 1996c)) to bright ones like SN 1991T, which was about 0.5 magnitudes brighter in $`B`$ than a typical Type Ia in the Virgo cluster (Mazzali, Danziger, & Turatto 1995), is commonly attributed to different <sup>56</sup>Ni-masses produced in the explosion. They range from about 0.07 M for SN 1991bg ( see, e.g., Mazzali et al. (1997a)) to at least 0.92 M for SN 1991T ((Khokhlov et al. 1993); see however Fisher et al. (1999)), with typically 0.6 M for normal SNe Ia (Höflich & Khokhlov 1996; Nugent et al. 1997). It is hard to see how this rather large range can be accommodated in a single class of models.
The stellar populations in which SNe Ia show up include spiral arms as well as elliptical galaxies, with some weak indication that they might be more efficiently produced in young populations (Bartunov, Tsvetkov, & Filimonova 1994). Again, if we insist on a single class of progenitors, the very fact that they do occur in ellipticals would rule out massive stars as potential candidates. On the other hand side, the observations may tell us that there is not a unique class of progenitors. In particular, the fact that the bright and slowly declining ones (like SN 1991T) are absent in elliptical and S0 galaxies may point towards different progenitor classes (Hamuy et al. 1996c).
All in all, the observational findings summarized so far are consistent with the assumption that Type Ia supernovae are the result of thermonuclear disruptions of white dwarfs, C+O white dwarfs being the favored model. The diversity among them must then be attributed to the history and nature of the white dwarf prior to the explosion and/or to the physics of thermonuclear burning during the event. It cannot be excluded, however, that at least some SNe Ia have a different origin, such as accretion-induced collapse of massive O-Ne-Mg (or O-Ne) white dwarfs for SN 1991bg–like objects (Nomoto et al. 1994a, 1995, 1996; Fryer et al. 1999). Also it is not clear whether or not there is a clear-cut distinction between Type Ib/c supernovae, defined by the absence of the Si II feature, and the (faint) SNe Ia. The former are believed to reflect the core-collapse of a massive star, its hydrogen-rich envelope being pealed off due to mass-loss in a binary system. For example, SN1987K started out as a SN II with H lines in its spectrum, but changed into a SN Ib/c-like spectrum after 6 months (Filippenko 1988), supporting this interpretation. It should be noted that SN 1991bg-like objects are not often observed, but that this may well be a selection effect. Suntzeff (1996), for example, argues that up to 40% of all Type Ia’s could perhaps belong to that sub-group.
### 4.2 Pre-supernova evolution of binary stars
Despite of all these uncertainties it is the current understanding and believe that the progenitors of SNe Ia are C+O white dwarfs in binary systems evolving to the stage of explosion by mass-overflow from the companion (single-degenerate scenario) or by the merger of two white dwarfs (double-degenerate scenario). Binary evolution of some sort is necessary because C+O white dwarfs a typically born with a mass around 0.6 M (Homeier et al. 1998) but need to be near the Chandrasekhar mass or to accumulate a shell of helium in order to explode. In this Subsection we will summarize the arguments in favor and against both scenarios.
Double-degenerates as potential Type Ia progenitors had many ups and downs in the past, beginning with the classic papers of Iben & Tutukov (1984) and Webbink (1984). The arguments in favor are that such binaries should exist as a consequence of stellar evolution, they would explain very naturally the absence of hydrogen, and they could, in principle, be an easy way to approach a critical mass. In fact, several candidate-systems of binary white dwarfs have recently been identified but most of the short-period ones (at present 8 systems are known with orbital periods of less than half a day), which could merge in a Hubble-time due to the emission of gravitational radiation, have a mass less than $`M_{\mathrm{chan}}`$ (Saffer, Livio, & Yungelson (1998); see also Livio (2000) for a recent review). There is only one system known (KPD 0422+5421; Koen, Orosz, & Wade (1998)) with a mass which, within the errors, could exceed $`M_{\mathrm{chan}}`$ , a surprisingly small number. None-the-less it is argued that from population synthesis one could arrive at about the right frequency of sufficiently massive mergers (Livio 2000).
Besides the lack of convincing direct observational evidence for sufficiently many appropriate binary systems, the homogeneity of “typical” SNe Ia may be an argument against this class of progenitors. It is not easy to see how the merging of two white dwarfs of (likely) different mass, composition, and angular momentum with different impact parameters, etc., will always lead to the same burning conditions and, therefore, the production of a nearly equal amount of <sup>56</sup>Ni. Moreover, some investigations of white dwarf mergers seem to indicate that an off-center ignition will convert carbon and oxygen into oxygen, neon, and magnesium, leading to gravitational collapse rather than a thermonuclear disruption (Woosley & Weaver 1986a; Saio & Nomoto 1985, 1998; Mochkovitch & Livio 1990). Finally, based on their galactic chemical evolution model, Kobayashi et al. (1998) claim that double-degenerate mergers lead to inconsistencies with the observed O/Fe as a function of metallicity, but this statement is certainly model-dependent. In any case, mergers might, if they are not responsible for the bulk of the SNe Ia, still account for some peculiar ones, such as the super-luminous SN 1991T -like explosions.
Single-degenerate models are in general favored today. They consist of a low-mass white dwarf accreting matter from the companion-star until either it reaches $`M_{\mathrm{chan}}`$ or a layer of helium has formed on-top of its C+O core that can ignite and possibly drive a burning front into the carbon and oxygen fuel. This track to thermonuclear explosions of white dwarfs was first discussed by Whelan & Iben (1973); Nomoto (1982b); Iben & Tutukov (1984) and Paczynski (1985). The major problem of these models has always been that nearly all possible accretion rates can be ruled out by rather strong arguments (Munari & Renzini 1992; Cassisi, Castellani, & Tornambe 1996; Tutukov & Yungelson 1996; Livio et al. 1996; King & Van Teeseling 1998; Kato & Hachisu 1999; Cassisi, Iben, & Tornambe 1998). In short, it is believed that white dwarfs accreting hydrogen at a low rate undergo nova eruptions and lose more mass in the outburst than they have accreted prior to it (e.g. Beer (1974); Gehrz, Truran, & Williams (1993)). At moderate accretion rates, a degenerate layer of helium is thought to form which might flash and could give rise to sub-Chandrasekhar explosions (which have other problems, as will be discussed later). Next, still higher accretion rates can lead to quiet hydrostatic burning of H and He, but these systems should be so bright that they could easily be detected, but it is not clear beyond doubt that they coincide with any of the known symbiotic or cataclysmic binaries. Very high accretion rates, finally would form an extended H-rich red giant envelope around the white dwarf the debris of which are not seen in the explosions (Nomoto, Nariai, & Sugimoto 1979) (see also Fig. 2). Therefore, it is very uncertain if white dwarfs accreting hydrogen from a companion star can ever reach the $`M_{\mathrm{chan}}`$ (Cassisi et al. 1998).
Some of these arguments may be questioned, however. Firstly, a class of binary systems has recently been discovered, the so-called “Supersoft X-ray Sources”, which are best interpreted as white dwarfs accreting hydrogen-rich matter at such a high rate that H burns steadily (Truemper et al. 1991; Greiner, Hasinger, & Kahabka 1991; Van Den Heuvel et al. 1992; Southwell et al. 1996; Kahabka & Van Den Heuvel 1997). It appears that if these white dwarfs could retain the accreted gas they might be good candidates for SN Ia progenitors. In principle, they could accrete a few tenths of a solar mass with a typical accretion rate of a few 10<sup>-7</sup>M/yr over the estimated lifetime of such systems of several 10<sup>9</sup> years. Since most of them are heavily extinct their total number might be sufficiently high (Di Stefano & Rappaport (1994); see also Livio (1996) and Yungelson et al. (1996)), although this statement is certainly model-dependent. However, some of the supersoft X-ray sources are known to be variable in X-rays (but not in the optical wave-bands) on time-scales of weeks (Pakull et al. 1993), too short to be related with the H-burning shell, possibly indicating substantial changes in the accretion rates. It may therefore not be justified to assume that the accretion rates we see now are sustained over several 10<sup>9</sup> years. But their very existence provides a first and strong case for the single-degenerate scenario.
Secondly, also the minimum accretion rate at which hydrogen burns quietly without a nova outburst is rather uncertain. All models that compute this rate ignore important pieces of physics and, therefore, their predictions could be off by orders of magnitude. For example, classical nova outbursts require that the accreted hydrogen-rich envelope of the white dwarf is also heavily enriched in C and O from the white dwarf’s core (see, e.g., Starrfield et al. (1972); Sparks, Starrfield, & Truran (1976); Starrfield, Truran, & Sparks (1978); Truran (1982)). One possible explanation has been that convective mixing and dredge-up might happen during the thermonuclear runaway, but recent numerical simulations indicate that this mechanism is insufficient (Kercek, Hillebrandt, & Truran 1999). In contrast to spherically symmetric models their 3-D simulations lead to a phase of quiet H-burning for accretion rates as low as 5 $`\times `$ 10<sup>-9</sup>M/yr for a white dwarf of 1 M rather than a nova outburst with mass-loss from the core. Other short-comings include the assumption of spherical accretion with zero entropy, the neglect of magnetic fields, etc. So the dividing line between steady hydrogen burning and nova eruptions might leave some room for SN Ia progenitors.
Finally, it has been argued that the interaction of a wind from the white dwarf with the accretion flow from lobe-filling low mass red giant may open a wider path to Type Ia supernovae. In a series of papersHachisu, Kato, & Nomoto (1996); Hachisu et al. (1999b); Hachisu, Kato, & Nomoto (1999a) discuss the effect that when the mass accretion rate exceeds a certain critical value the envelope solution on the white dwarf is no longer static but corresponds to a strong wind. The strong wind stabilizes the mass transfer and limits the accretion rate and the white dwarf can burn hydrogen steadily. However, their model assumes spherical accretion onto and a spherical wind from the white dwarf which seem to be contradicting assumptions. But the idea should certainly be followed up.
### 4.3 Evolution to ignition
In what follows we will assume most of the time that SN Ia progenitors are Chandrasekhar-mass C+O white dwarfs because, as was discussed in the previous sections, this class of models seems to fit best the “typical” or “average” Type Ia. In this subsection we also will not discuss models in which two degenerate white dwarfs merge and form a critical mass for the ignition of carbon, mainly because the merging process will, in reality, be very complex and it is difficult to construct realistic explosion models (although with increasing computational resources it may be possible in the future).
But even if we consider only Chandrasekhar-mass white dwarfs as progenitor candidates the information that is needed in order to model the explosion cannot be obtained easily. In particular, the thermal structure and the chemical composition are very uncertain. The C/O-ratio, for example, has to be know throughout the white dwarf, but this ratio depends on the main sequence mass of its progenitor and the metallicity of the gas from which it formed (Umeda et al. 1999; Wellstein & Langer 1999). It was found that, depending on the main sequence mass, the central C/O can vary from 0.4 to 0.6, considerably less than assumed in most supernova models.
Next, the thermal structure of a white dwarf on its way to an explosion depends on the (convective) URCA-process (Paczynski 1973; Iben 1978, 1982; Barkat & Wheeler 1990; Mochkovitch 1996). The URCA-pairs A = 21, 23, and 25 (such as, i.e., <sup>21</sup>Ne/<sup>21</sup>F, …) can lead to either heating or cooling, and possibly even to a temperature inversion near the center of the white dwarf. The abundances of the URCA-pairs depends again on the initial metallicity which could, thus, affect the thermal structure of the white dwarf. Unfortunately, the convection in the degenerate star is likely to be non-local, time-dependent, 3-dimensional, and very sub-sonic, but needs to be modeled over very long (secular) time-scales. It is not likely that in the near future we will be able to model these processes in a realistic manner, even on super-computers.
Due to these difficulties, numerical studies of the explosion rely on ad-hoc assumptions fixing the initial conditions, which are usually chosen to be as simple as possible. Realistic simulations have to be multi-dimensional, as will be explained in the next Section, and therefore numerical studies can only investigate a small fraction of the available parameter space. The failure or success of a particular model to explain certain observational results may, therefore, not be conclusive.
## 5 EXPLOSION MODELING
Numerical models are needed to provide the density, temperature, composition, and velocity fields of the supernova ejecta that result from the thermonuclear explosion of a white dwarf, accepted by most researchers as the “standard model” for SNe Ia (Sec. (2), (4)). This information can then be used to compute the resulting light curve and spectra with the help of radiation transport codes (Sec. (3)) or to compare the relative distribution of isotopes with the observed solar abundances.
To a very good approximation, the exploding white dwarf material can be described as a fully ionized plasma with varying degrees of electron degeneracy, satisfying the fluid approximation. The governing equations are the hydrodynamical equations for mass, species, momentum, and energy transport including gravitational acceleration, viscosity, heat and mass diffusion (Landau & Lifshitz 1963), and nuclear energy generation (Arnett 1996). They must be supplemented by an equation of state for an ideal gas of nuclei, an arbitrarily relativistic and degenerate electron gas, radiation, and electron-positron pair production and annihilation (Cox & Giuli 1968). The gravitational potential is calculated with the help of the Poisson equation. In numerical simulations that fully resolve the relevant length scales for dissipation, diffusion, and nuclear burning it is possible to obtain the energy generation rate from a nuclear reaction network (Timmes 1999, for a recent overview, see) and the diffusion coefficients from an evaluation of the kinetic transport mechanisms (Nandkumar & Pethick 1984). If, on the other hand, these scales are unresolved – as is usually the case in simulations on scales of the stellar radius – subgrid-scale models are required to compute (or parameterize) the effective large-scale transport coefficients and burning rates, which are more or less unrelated to the respective microphysical quantities (Khokhlov 1995; Niemeyer & Hillebrandt 1995b).
Initial conditions can be obtained from hydrostatic spherically symmetric models of the accreting white dwarf or – for Chandrasekhar mass progenitors – from the Chandrasekhar equation for a fully degenerate, zero temperature white dwarf (Kippenhahn & Weigert 1989). Given the initial conditions and symmetries specifying the boundary conditions, the dynamics of the explosion can in principle be determined by numerically integrating the equations of motion. Müller (1998) gives a detailed account of some current numerical techniques used for modeling supernovae.
Until the mid-nineties, most work on SN Ia explosions was done studying one-dimensional (1D), spherically symmetric models. This approach inherently lacks some important aspects of multidimensional thermonuclear burning relevant for $`M_{\mathrm{chan}}`$ -explosion models , e.g. off-center flame ignition, flame instabilities, and turbulence, which have to be mimicked by means of a spherical flame front with an undetermined turbulent flame speed e.g., Nomoto, Sugimoto, & Neo (1976); Nomoto, Thielemann, & Yokoi (1984); Woosley & Weaver (1986a); Woosley (1990). In spite of these caveats, 1D models still represent the only reasonable approach to combine the hydrodynamics with detailed nucleosynthesis calculations and to carry out parameter studies of explosion scenarios. In fact, most of the phenomenology of SN Ia explosions and virtually all of the model predictions for spectra and light curves are based on spherically symmetric models. Several recent articles (Woosley 1990; Nomoto et al. 1996; Höflich & Khokhlov 1996; Iwamoto et al. 1999) describe the methodology and trends observed in these studies, as well as their implications regarding the cosmological supernova surveys (Höflich et al. 1998b; Ruiz-Lapuente & Canal 1998; Umeda et al. 1999; Sorokina, Blinnikov, & Bartunov 1999).
Following the pioneering work of Müller & Arnett (1982, 1986), some groups have explored the dynamics of two-dimensional (2D) (Livne 1993; Arnett & Livne 1994a, 1994b; Niemeyer & Hillebrandt 1995b; Niemeyer, Hillebrandt, & Woosley 1996; Arnett 1997; Reinecke, Hillebrandt, & Niemeyer 1999a) and three-dimensional (3D) (Khokhlov 1994, 1995; Bravo & Garcia-Senz 1997; Benz 1997) explosion models, triggering the development of numerical algorithms for representing thin propagating surfaces in large-scale simulations (Khokhlov 1993a; Niemeyer & Hillebrandt 1995b; Bravo & Garcia-Senz 1995; Arnett 1997; Garcia-Senz, Bravo, & Serichol 1998; Reinecke et al. 1999b). It has also become possible to perform 2D and 3D direct numerical simulations (DNS), i.e. fully resolving the relevant burning and diffusion scales, of microscopic flame instabilities and flame-turbulence interactions (Niemeyer & Hillebrandt 1995a; Khokhlov 1995; Niemeyer & Hillebrandt 1997; Niemeyer, Bushe, & Ruetsch 1999).
### 5.1 Chandrasekhar Mass Explosion Models
Given the overall homogeneity of SNe Ia (Sec. (2.1)), the good agreement of parameterized 1D $`M_{\mathrm{chan}}`$ -models with observed spectra and light curves, and their reasonable nucleosynthetic yields, the bulk of normal SNe Ia is generally assumed to consist of exploding $`M_{\mathrm{chan}}`$ C+O white dwarfs (Hoyle & Fowler 1960; Arnett 1969; Hansen & Wheeler 1969). In spite of three decades of work on the hydrodynamics of this explosion mechanism (beginning with Arnett (1969)), no clear consensus has been reached whether the star explodes as a result of a subsonic nuclear deflagration that becomes strongly turbulent (Ivanova, Imshennik, & Chechetkin 1974; Buchler & Mazurek 1975; Nomoto et al. 1976, 1984; Woosley, Axelrod, & Weaver 1984), or whether this turbulent flame phase is followed by a delayed detonation during the expansion (Khokhlov 1991a, 1991b; Woosley & Weaver 1994a) or after one or many pulses (Khokhlov 1991b; Arnett & Livne 1994a, 1994b). Only the prompt detonation mechanism is agreed to be inconsistent with SN Ia spectra as it fails to produce sufficient amounts of intermediate mass elements (Arnett 1969; Arnett, Truran, & Woosley 1971).
This apparently slow progress is essentially a consequence of the overwhelming complexity of turbulent flame physics and deflagration-detonation transitions (DDTs) (Williams 1985; Zeldovich et al. 1985) that makes first-principle predictions based on $`M_{\mathrm{chan}}`$ -explosion models nearly impossible. The existence of an initial subsonic flame phase is, it seems, an unavoidable ingredient of all $`M_{\mathrm{chan}}`$ -models (and only those) where it is required to pre-expand the stellar material prior to its nuclear consumption in order to avoid the almost exclusive production of iron-peaked nuclei (Nomoto et al. 1976, 1984; Woosley & Weaver 1986a).
Guided by parameterized 1D models that yield estimates for the values for the turbulent flame speed $`S_\mathrm{t}`$ and the DDT transition density $`\rho _{\mathrm{DDT}}`$ (e.g., Höflich & Khokhlov 1996), a lot of work has been done recently on the physics of buoyancy-driven, turbulent thermonuclear flames in exploding $`M_{\mathrm{chan}}`$ -white dwarfs. The close analogy with thin chemical premixed flames has been exploited to develop a conceptual framework that covers all scales from the white dwarf radius to the microscopic flame thickness and dissipation scales (Khokhlov 1995; Niemeyer & Woosley 1997). In the following discussion of nuclear combustion (Sec. (5.1.1)), flame ignition (Sec. (5.1.2)), and the various scenarios for $`M_{\mathrm{chan}}`$ explosions characterized by the sequence of combustion modes (Sec. (5.1.3) – Sec. (5.1.6)) we will emphasize the current understanding of physical processes rather than empirical fits of light curves and spectra.
#### 5.1.1 FLAMES, TURBULENCE, AND DETONATIONS
Owing to the strong temperature dependence of the nuclear reaction rates, $`\dot{S}T^{12}`$ at $`T10^{10}`$ K (Hansen & Kawaler 1994, p. 247), nuclear burning during the explosion is confined to microscopically thin layers that propagate either conductively as subsonic deflagrations (“flames”) or by shock compression as supersonic detonations (Courant & Friedrichs 1948; Landau & Lifshitz 1963, chap. XIV). Both modes are hydrodynamically unstable to spatial perturbations as can be shown by linear perturbation analysis. In the nonlinear regime, the burning fronts are either stabilized by forming a cellular structure or become fully turbulent – either way, the total burning rate increases as a result of flame surface growth (Lewis & von Elbe 1961; Williams 1985; Zeldovich et al. 1985). Neither flames nor detonations can be resolved in explosion simulations on stellar scales and therefore have to be represented by numerical models.
When the fuel exceeds a critical temperature $`T_\mathrm{c}`$ where burning proceeds nearly instantaneously compared with the fluid motions (see Timmes & Woosley (1992) for a suitable definition of $`T_\mathrm{c}`$), a thin reaction zone forms at the interface between burned and unburned material. It propagates into the surrounding fuel by one of two mechanisms allowed by the Rankine-Hugoniot jump conditions: a deflagration (“flame”) or a detonation (cf. fig. 2.5 in Williams (1985)).
If the overpressure created by the heat of the burning products is sufficiently high, a hydrodynamical shock wave forms that ignites the fuel by compressional heating. A self-sustaining combustion front that propagates by shock-heating is called a detonation. Detonations generally move supersonically and therefore do not allow the unburned medium to expand before it is burned. Their speed depends mainly on the total amount of energy released per unit mass, $`ϵ`$, and is therefore more robustly computable than deflagration velocities. A good estimate for the velocity of planar strong detonations is the Chapman-Jouget velocity (Lewis & von Elbe 1961; Zeldovich et al. 1985; Williams 1985, and references therein). The nucleosynthesis, speed, structure, and stability of planar detonations in degenerate C+O material was analyzed by Imshennik & Khokhlov (1984); Khokhlov (1988, 1989, 1993b), and recently by Kriminski, Bychkov, & Liberman (1998) and Imshennik et al. (1999) who claim that C+O detonations are one-dimensionally unstable and therefore cannot occur in exploding white dwarfs above a critical density of $`2\times 10^7`$ g cm<sup>-3</sup> (Kriminski et al. 1998) (cf. Sec. (5.1.3)).
If, on the other hand, the initial overpressure is too weak, the temperature gradient at the fuel-ashes interface steepens until an equilibrium between heat diffusion (carried out predominantly by electron-ion collisions) and energy generation is reached. The resulting combustion front consists of a diffusion zone that heats up the fuel to $`T_\mathrm{c}`$, followed by a thin reaction layer where the fuel is consumed and energy is generated. It is called a deflagration or simply a flame and moves subsonically with respect to the unburned material (Landau & Lifshitz 1963). Flames, unlike detonations, may therefore be strongly affected by turbulent velocity fluctuations of the fuel. Only if the unburned material is at rest, a unique laminar flame speed $`S_\mathrm{l}`$ can be found which depends on the detailed interaction of burning and diffusion within the flame region (e.g., Zeldovich et al. 1985). Following Landau & Lifshitz (1963), it can be estimated by assuming that in order for burning and diffusion to be in equilibrium, the respective time scales, $`\tau _\mathrm{b}ϵ/\dot{w}`$ and $`\tau _\mathrm{d}\delta ^2/\kappa `$, where $`\delta `$ is the flame thickness and $`\kappa `$ is the thermal diffusivity, must be similar: $`\tau _\mathrm{b}\tau _\mathrm{d}`$. Defining $`S_\mathrm{l}=\delta /\tau _\mathrm{b}`$, one finds $`S_\mathrm{l}(\kappa \dot{w}/ϵ)^{1/2}`$, where $`\dot{w}`$ should be evaluated at $`TT_\mathrm{c}`$ (Timmes & Woosley 1992). This is only a crude estimate due to the strong $`T`$-dependence of $`\dot{w}`$. Numerical solutions of the full equations of hydrodynamics including nuclear energy generation and heat diffusion are needed to obtain more accurate values for $`S_\mathrm{l}`$ as a function of $`\rho `$ and fuel composition. Laminar thermonuclear carbon and oxygen flames at high to intermediate densities were investigated by Buchler, Colgate, & Mazurek (1980); Ivanova, Imshennik, & Chechetkin (1982); Woosley & Weaver (1986b), and, using a variety of different techniques and nuclear networks, by Timmes & Woosley (1992). For the purpose of SN Ia explosion modeling, one needs to know the laminar flame speed $`S_\mathrm{l}10^7\mathrm{}10^4`$ cm s<sup>-1</sup> for $`\rho 10^9\mathrm{}10^7`$ g cm<sup>-3</sup>, the flame thickness $`\delta =10^4\mathrm{}1`$ cm (defined here as the width of the thermal pre-heating layer ahead of the much thinner reaction front), and the density contrast between burned and unburned material $`\mu =\mathrm{\Delta }\rho /\rho =0.2\mathrm{}0.5`$ (all values quoted here assume a composition of $`X_\mathrm{C}=X_\mathrm{O}=0.5`$, Timmes & Woosley (1992)). The thermal expansion parameter $`\mu `$ reflects the partial lifting of electron degeneracy in the burning products, and is much lower than the typical value found in chemical, ideal gas systems (Williams 1985).
Observed on scales much larger than $`\delta `$, the internal reaction-diffusion structure can be neglected and the flame can be approximated as a density jump that propagates locally with the normal speed $`S_\mathrm{l}`$. This “thin flame” approximation allows a linear stability analysis of the front with respect to spatial perturbations. The result shows that thin flames are linearly unstable on all wavelengths. It was discovered first by Landau (1944) and Darrieus (1944) and is hence called the “Landau-Darrieus” (LD) instability. Subject to the LD instability, perturbations grow until a web of cellular structures forms and stabilizes the front at finite perturbation amplitudes (Zeldovich 1966). The LD instability therefore does not, in general, lead to the production of turbulence. In the context of SN Ia models, the nonlinear LD instability was studied by Blinnikov & Sasorov (1996), using a statistical approach based on the Frankel equation, and by Niemeyer & Hillebrandt (1995a) employing 2D hydrodynamics and a one-step burning rate. Both groups concluded that the cellular stabilization mechanism precludes a strong acceleration of the burning front as a result of the LD instability. However, Blinnikov & Sasorov (1996) mention the possible breakdown of stabilization at low stellar densities (i.e., high $`\mu `$) which is also indicated by the lowest density run of Niemeyer & Hillebrandt (1995a) – this may be important in the framework of active turbulent combustion (see below). The linear growth rate of LD unstable thermonuclear flames with arbitrary equation of state was derived by Bychkov & Liberman (1995a). The same authors also found a one-dimensional, pulsational instability of degenerate C+O flames (Bychkov & Liberman 1995b) which was later disputed by Blinnikov (1996).
The best studied and probably most important hydrodynamical effect for modeling SN Ia explosions is the Rayleigh-Taylor (RT) instability (Rayleigh 1883; Chandrasekhar 1961) resulting from the buoyancy of hot, burned fluid with respect to the dense, unburned material. Several groups have investigated the RT instability of nuclear flames in SNe Ia by means of numerical hydrodynamical simulations (Müller & Arnett 1982, 1986; Livne 1993; Khokhlov 1994, 1995; Niemeyer & Hillebrandt 1995b). After more than five decades of experimental and numerical work, the basic phenomenology of nonlinear RT mixing is fairly well understood (Fermi 1951; Layzer 1955; Sharp 1984; Read 1984; Youngs 1984): Subject to the RT instability, small surface perturbations grow until they form bubbles (or “mushrooms”) that begin to float upward while spikes of dense fluid fall down. In the nonlinear regime, bubbles of various sizes interact and create a foamy RT mixing layer whose vertical extent $`h_{\mathrm{RT}}`$ grows with time $`t`$ according to a self-similar growth law, $`h_{\mathrm{RT}}=\alpha g(\mu /2)t^2`$, where $`\alpha `$ is a dimensionless constant ($`\alpha 0.05`$) and $`g`$ is the background gravitational acceleration (Sharp 1984; Youngs 1984; Read 1984).
Secondary instabilities related to the velocity shear along the bubble surfaces (Niemeyer & Hillebrandt 1997) quickly lead to the production of turbulent velocity fluctuations that cascade from the size of the largest bubbles ($`10^7`$ cm) down to the microscopic Kolmogorov scale, $`l_\mathrm{k}10^4`$ cm where they are dissipated (Niemeyer & Hillebrandt 1995b; Khokhlov 1995). Since no computer is capable of resolving this range of scales, one has to resort to statistical or scaling approximations of those length scales that are not properly resolved. The most prominent scaling relation in turbulence research is Kolmogorov’s law for the cascade of velocity fluctuations, stating that in the case of isotropy and statistical stationarity, the mean velocity $`v`$ of turbulent eddies with size $`l`$ scales as $`vl^{1/3}`$ (Kolmogorov 1941). Knowledge of the eddy velocity as a function of length scale is important to classify the burning regime of the turbulent combustion front (Niemeyer & Woosley 1997; Niemeyer & Kerstein 1997; Khokhlov, Oran, & Wheeler 1997). The ratio of the laminar flame speed and the turbulent velocity on the scale of the flame thickness, $`K=S_\mathrm{l}/v(\delta )`$, plays an important role: if $`K1`$, the laminar flame structure is nearly unaffected by turbulent fluctuations. Turbulence does, however, wrinkle and deform the flame on scales $`l`$ where $`S_\mathrm{l}v(l)`$, i.e. above the Gibson scale $`l_\mathrm{g}`$ defined by $`S_\mathrm{l}=v(l_\mathrm{g})`$ (Peters 1988). These wrinkles increase the flame surface area and therefore the total energy generation rate of the turbulent front (Damköhler 1940). In other words, the turbulent flame speed, $`S_\mathrm{t}`$, defined as the mean overall propagation velocity of the turbulent flame front, becomes larger than the laminar speed $`S_\mathrm{l}`$. If the turbulence is sufficiently strong, $`v(L)S_\mathrm{l}`$, the turbulent flame speed becomes independent of the laminar speed, and therefore of the microphysics of burning and diffusion, and scales only with the velocity of the largest turbulent eddy (Damköhler 1940; Clavin 1994):
$$S_\mathrm{t}v(L).$$
(2)
Because of the unperturbed laminar flame properties on very small scales, and the wrinkling of the flame on large scales, the burning regime where $`K1`$ is called the corrugated flamelet regime (Pope 1987; Clavin 1994).
As the density of the white dwarf material declines and the laminar flamelets become slower and thicker, it is plausible that at some point turbulence significantly alters the thermal flame structure (Khokhlov et al. 1997; Niemeyer & Woosley 1997). This marks the end of the flamelet regime and the beginning of the distributed burning, or distributed reaction zone, regime (e.g., Pope 1987). So far, modeling the distributed burning regime in exploding white dwarfs has not been attempted explicitely since neither nuclear burning and diffusion nor turbulent mixing can be properly described by simplified prescriptions. Phenomenologically, the laminar flame structure is believed to be disrupted by turbulence and to form a distribution of reaction zones with various lengths and thicknesses. In order to find the critical density for the transition between both regimes, we need to formulate a specific criterion for flamelet breakdown. A criterion for the transition between both regimes is discussed in Niemeyer & Woosley (1997); Niemeyer & Kerstein (1997) and Khokhlov et al. (1997):
$$l_{\mathrm{cutoff}}\delta .$$
(3)
Inserting the results of Timmes & Woosley (1992) for $`S_\mathrm{l}`$ and $`\delta `$ as functions of density, and using a typical turbulence velocity $`v(10^6\text{cm})10^7`$ cm s<sup>-1</sup>, the transition from flamelet to distributed burning can be shown to occur at a density of $`\rho _{\mathrm{dis}}10^7`$ g cm<sup>-3</sup> (Niemeyer & Kerstein 1997).
The close coincidence of $`\rho _{\mathrm{dis}}`$ and the preferred value for $`\rho _{\mathrm{DDT}}`$ (Höflich & Khokhlov 1996; Nomoto et al. 1996) inspired some authors (Niemeyer & Woosley 1997; Khokhlov et al. 1997) to suggest that both are related by local flame quenching and re-ignition via the Zeldovich induction time gradient mechanism (Zeldovich et al. 1970), whereby a macroscopic region with a uniform temperature gradient can give birth to a supersonic spontaneous combustion wave that steepens into a detonation (Woosley 1990, and references therein). In the context of the SN Ia explosion mechanism, this effect was first analyzed by Blinnikov & Khokhlov (1986, 1987). Whether or not the gradient mechanism can account for DDTs in the delayed detonation scenario for SNe Ia is still controversial; while Khokhlov et al. (1997) conclude that it can, Niemeyer (1999) – using arguments based on incompressible computations of microscopic flame-turbulence interactions by Niemeyer et al. (1999) – states that thermonuclear flames may be too robust with respect to turbulent quenching to allow the formation of a sufficiently uniform temperature gradient.
Assuming that the nonlinear RT instability dominates the turbulent flow that advects the flame, the passive-surface description of the flame neglects the additional stirring caused by thermal expansion within the flame brush itself, accelerating the burnt material in random directions. Both the spectrum and cutoff scale may be affected by “active” turbulent combustion (Kerstein 1996; Niemeyer & Woosley 1997). Although the small expansion coefficient $`\mu `$ indicates that the effect is weak compared to chemical flames, a quantitative answer is still missing.
Finally, we note that some authors also studied the multidimensional instability of detonations in degenerate C+O matter (Boisseau et al. 1996; Gamezo et al. 1999), finding unsteady front propagation, the formation of a cellular front structure and locally incomplete burning in multidimensional C+O detonations. These effects may have interesting implications for SN Ia scenarios involving a detonation phase.
#### 5.1.2 FLAME IGNITION
As the white dwarf grows close to the Chandrasekhar mass $`M_{\mathrm{chan}}`$ $`1.4M_{}`$, the energy budget near the core is governed by plasmon neutrino losses and compressional heating. The neutrino losses increase with growing central density until the latter reaches approximately $`2\times 10^9`$ g cm<sup>-3</sup> (Woosley & Weaver 1986a). At this point, plasmon creation becomes strongly suppressed while electron screening of nuclear reactions enhances the energy generation rate until it begins to exceed the neutrino losses. This “smoldering” of the core region marks the beginning of the thermonuclear runaway (Arnett 1969; Arnett 1971; Woosley & Weaver 1986a). During the following $``$ 1000 years, the core experiences internally heated convection with progressively smaller turnover time scales $`\tau _\mathrm{c}`$. Simultaneously, the typical time scale for thermonuclear burning, $`\tau _\mathrm{b}`$, drops even faster as a result of the rising core temperature and the steep temperature dependence of the nuclear reaction rates.
During this period, the entropy and temperature evolution of the core is affected by the convective URCA process, a convectively driven electron capture-beta decay cycle leading to neutrino-antineutrino losses. It was first described in this context by Paczynski (1972) who argued it would cause net cooling and therefore delay the runaway. Since then, the convective URCA process was revisited by several authors (e.g, Bruenn 1973; Iben 1982; Barkat & Wheeler 1990; Mochkovitch 1996) who alternately claimed that it results in overall heating or cooling. The most recent analysis (Stein, Barkat, & Wheeler 1999) concludes that while the URCA neutrinos carry away energy, they cannot cool the core globally but instead slow down the convective motions.
At $`T7\times 10^8`$ K, $`\tau _\mathrm{c}`$ and $`\tau _\mathrm{b}`$ become comparable, indicating that convective plumes burn at the same rate as they circulate (Nomoto et al. 1984; Woosley & Weaver 1986a). Experimental or numerical data describing this regime of strong reactive convection is not available, but several groups are planning to conduct numerical experiments at the time this article is written. At $`T1.5\times 10^9`$ K, $`\tau _\mathrm{b}`$ becomes extremely small compared with $`\tau _\mathrm{c}`$, and carbon and oxygen virtually burn in place. A new equilibrium between energy generation and transport is found on much smaller length scales, $`l10^4`$ cm, where thermal conduction by degenerate electrons balances nuclear energy input (Timmes & Woosley 1992). The flame is born.
The evolution of the runaway immediately prior to ignition of the flame is crucial for determining its initial location and shape. Using a simple toy model, Garcia-Senz & Woosley (1995) found that under certain conditions, burning bubbles subject to buoyancy and drag forces can rise a few hundred km before flame formation, suggesting a high probability for off-center ignition at multiple, unconnected points. As a consequence, more material burns at lower densities, thus producing higher amounts of intermediate mass elements than a centrally ignited explosion. In a parameter study, Niemeyer et al. (1996) and Reinecke et al. (1999a) demonstrated the significant influence of the location and number of initially ignited spots on the final explosion energetics and nucleosynthesis.
#### 5.1.3 PROMPT DETONATION
The first hydrodynamical simulation of an exploding $`M_{\mathrm{chan}}`$ -white dwarf (Arnett 1969) assumed that the thermonuclear combustion commences as a detonation wave, consuming the entire star at the speed of sound. Given no time to expand prior to being burned, the C+O material in this scenario is transformed almost completely into iron-peaked nuclei and thus fails to produce significant amounts of intermediate mass elements, in contradiction to observations (Filippenko 1997a, 1997b). It is for this reason that prompt detonations are generally considered ruled out as viable candidates for the SN Ia explosion mechanism.
In addition to the empirical evidence, the ignition of a detonation in the high density medium of the white dwarf core was argued to be an unlikely event. In spite of the smallness of the critical mass for detonation at $`\rho 2\times 10^9`$ g cm<sup>-3</sup> (Niemeyer & Woosley 1997; Khokhlov et al. 1997) and the correspondingly large number of critical volumes in the core ($`10^{18}`$), the stringent uniformity condition for the temperature gradient of the runaway region (Blinnikov & Khokhlov 1986, 1987) was shown to be violated even by the minute amounts of heat dissipated by convective motions (Niemeyer & Woosley 1997). A different argument against the occurrence of a prompt detonation in C+O white dwarf cores was given by Kriminski et al. (1998), who found that C+O detonations may be subject to self-quenching at high material densities ($`\rho >2\times 10^7`$ g cm<sup>-3</sup>) (see also Imshennik et al. 1999).
#### 5.1.4 PURE TURBULENT DEFLAGRATION
Once ignited (Sec. (5.1.2)), the subsonic thermonuclear flame becomes highly convoluted as a result of turbulence produced by the various flame instabilities (Sec. (5.1.1)). It continues to burn through the star until it either transitions into a detonation or is quenched by expansion. The key questions with regard to explosion modeling are: a) What is the effective turbulent flame speed $`S_\mathrm{t}`$ as a function of time, b) Is the total amount of energy released during the deflagration phase enough to unbind the star and produce a healthy explosion, and c) Does the resulting ejecta composition and velocity agree with observations?
By far the most work has been done on 1D models, ignoring the multidimensionality of the flame physics and instead parameterizing $`S_\mathrm{t}`$ in order to answer b) and c) above (see Woosley & Weaver 1986a; Nomoto et al. 1996, for reviews). One of the most successful examples, model W7 of Nomoto et al. (1984), clearly demonstrates the excellent agreement of “fast” deflagration models with SN Ia spectra and light curves. $`S_\mathrm{t}`$ has been parameterized differently by different authors, for instance as a constant fraction of the local sound speed (Höflich & Khokhlov 1996; Iwamoto et al. 1999), using time-dependent convection theory (Nomoto et al. 1976; Buchler & Mazurek 1975; Nomoto et al. 1984; Woosley et al. 1984), or with a phenomenological fractal model describing the multiscale character of the wrinkled flame surface (Woosley 1990; Woosley 1997b). All of these studies essentially agree that very good agreement with the observations is obtained if $`S_\mathrm{t}`$ accelerates up to roughly 30 % of the sound speed. There remains a problem with the overproduction of neutron rich iron-group isotopes in fast deflagration models (Woosley et al. 1984; Thielemann, Nomoto, & Yokoi 1986; Iwamoto et al. 1999), but this may be alleviated in multiple dimensions (see below). Turning this argument around, Woosley (1997a) argues that <sup>48</sup>Ca can only be produced by carbon burning in the very high density regime of a $`M_{\mathrm{chan}}`$ white dwarf core, providing a clue that a few SNe Ia need to be $`M_{\mathrm{chan}}`$ explosions igniting at $`\rho 2\times 10^9`$ g cm<sup>-3</sup>. A slightly different approach to 1D SN Ia modeling was taken by Niemeyer & Woosley (1997), employing the self-similar growth rate of RT mixing regions (Sec. (5.1.1)) to prescribe the turbulent flame speed. Here, all the free parameters are fixed by independent simulations or experiments. The result shows a successful explosion, albeit short on intermediate mass elements, suggesting that the employed flame model is still too simplistic.
A number of authors have studied multidimensional deflagrations in exploding $`M_{\mathrm{chan}}`$ -white dwarfs using a variety of hydrodynamical methods (Livne 1993; Arnett & Livne 1994a; Khokhlov 1995; Niemeyer & Hillebrandt 1995b; Niemeyer et al. 1996; Reinecke et al. 1999a). The problem of simulating subsonic flames in large-scale simulations has two aspects: the representation of the thin, propagating surface separating hot and cold material with different densities, and the prescription of the local propagation velocity $`S_\mathrm{t}(\mathrm{\Delta })`$ of this surface as a function of the hydrodynamical state of the large-scale calculation with numerical resolution $`\mathrm{\Delta }`$. The former problem has been addressed with artificial reaction-diffusion fronts in PPM (Khokhlov 1995; Niemeyer & Hillebrandt 1995b; Niemeyer et al. 1996) and SPH (Garcia-Senz et al. 1998), a PPM-specific flame tracking technique (Arnett 1997), and a hybrid flame capuring/tracking method based on level sets (Reinecke et al. 1999b) (see Fig. 3). Regarding the flame speed prescription, some authors assigned the local front propagation velocity assuming that the flame is laminar on unresolved scales $`l<\mathrm{\Delta }`$ (Arnett & Livne 1994a), by postulating that $`S_\mathrm{t}(\mathrm{\Delta })`$ is dominated by the terminal rise velocity of $`\mathrm{\Delta }`$-sized bubbles (Khokhlov 1995), or by using Eq. (2) together with a subgrid-scale model for the unresolved turbulent kinetic energy providing $`v(\mathrm{\Delta })`$ (Niemeyer & Hillebrandt 1995b; Niemeyer et al. 1996; Reinecke et al. 1999a).
In most multidimensional calculations on stellar scales to date, the effective turbulent flame speed stayed below the required 30 % of the sound speed. The detailed outcome of the explosion is controversial; while some calculations show that the star remains gravitationally bound after the deflagration phase has ceased (Khokhlov 1995), others indicate that $`S_\mathrm{t}`$ may be large enough to produce a weak but definitely unbound explosion (Niemeyer et al. 1996). These discrepancies can probably be attributed to differences in the description of the turbulent flame and to numerical resolution effects that plague all multidimensional calculations.
Niemeyer & Woosley (1997) and Niemeyer (1999) speculate about additional physics that can increase the burning rate in turbulent deflagration models, in particular multipoint ignition and active turbulent combustion (ATC), i.e. the generation of additional turbulence by thermal expansion within the turbulent flame brush. ATC can, in principle, explain the acceleration of $`S_\mathrm{t}`$ up to some fraction of the sound speed (Kerstein 1996), but its effectiveness is so far unknown. Multipoint ‘ignition, on the other hand, has already been shown to significantly increase the total energy release compared to single-point ignition models (Niemeyer et al. 1996; Reinecke et al. 1999a). Furthermore, it allows more material to burn at lower densities, thus alleviating the nucleosynthesis problem of 1D fast deflagration models (Niemeyer et al. 1996).
We conclude the discussion of the pure turbulent deflagration scenario with a checklist of the model requirements summarized in Sec. (2.4). Assuming that some combination of buoyancy, ATC, and multipoint ignition can drive the effective turbulent flame speed to $`30\%`$ of the sound speed – which is not evident from multidimensional simulations – one can conclude from 1D simulations that pure deflagration models readily comply with all observational constraints. Most authors agree that $`S_\mathrm{t}`$ decouples from microphysics on large enough scales and becomes dominated by essentially universal hydrodynamical effects, making the scenario intrinsically robust. A noteworthy exception is the location and number of ignition points that can strongly influence the explosion outcome and may be a possible candidate for the mechanism giving rise to the explosion strength variability. Other possible sources of variations include the ignition density and the accretion rate of the progenitor system (Umeda et al. 1999; Iwamoto et al. 1999). All of these effects may potentially vary with composition and metallicity and can therefore account for the dependence on the progenitor stellar population.
#### 5.1.5 DELAYED DETONATION
Turbulent deflagrations can sometimes be observed to undergo spontaneous transitions to detonations (deflagration-detonation transitions, DDTs) in terrestrial combustion experiments (e.g., Williams 1985, pp. 217–219). Thus inspired, it was suggested that DDTs may occur in the late phase of a $`M_{\mathrm{chan}}`$ -explosion, providing an elegant explanation for the initial slow burning required to pre-expand the star, followed by a fast combustion mode that produces large amounts of high-velocity intermediate mass elements (Khokhlov 1991a; Woosley & Weaver 1994a). Many 1D simulations have meanwhile demonstrated the capability of the delayed detonation scenario to provide excellent fits to SN Ia spectra and light curves (Woosley 1990; Höflich & Khokhlov 1996), as well as reasonable nucleosynthesis products with regard to solar abundances (Khokhlov 1991b; Iwamoto et al. 1999). In the best fit models, the initial flame phase has a rather slow velocity of roughly one percent of the sound speed and transitions to detonation at a density of $`\rho _{\mathrm{DDT}}10^7`$ g cm<sup>-3</sup> (Höflich & Khokhlov 1996; Iwamoto et al. 1999). The transition density was also found to be a convenient parameter to explain the observed sequence of explosion strengths (Höflich & Khokhlov 1996).
Various mechanisms for DDT were discussed in the early literature on delayed detonations (see Niemeyer & Woosley (1997) and references therein). Recent investigations have focussed on the induction time gradient mechanism (Zeldovich et al. 1970; Lee, Knystautas, & Yoshikawa 1978), analyzed in the context of SNe Ia by Blinnikov & Khokhlov (1986) and Blinnikov & Khokhlov (1987). It was realized by Khokhlov et al. (1997) and Niemeyer & Woosley (1997) that a necessary criterion for this mechanism is the local disruption of the flame sheet by turbulent eddies, or, in other words, the transition of the burning regime from “flamelet” to “distributed” burning (Sec. (5.1.1)). Simple estimates (Niemeyer & Kerstein 1997) show that this transition should occur at roughly $`10^7`$ g cm<sup>-3</sup>, providing a plausible explanation for the delay of the detonation.
The critical length (or mass) scale over which the temperature gradient must be held fixed in order to allow the spontaneous combustion wave to turn into a detonation was computed by Khokhlov et al. (1997) and Niemeyer & Woosley (1997); it is a few orders of magnitude thicker than the final detonation front and depends very sensitively on composition and density.
The virtues of the delayed detonation scenario can again be summarized by completing the checklist of Sec. (2.4). It is undisputed that suitably tuned delayed detonations satisfy all the constraints given by SN Ia spectra, light curves, and nucleosynthesis. If $`\rho _{\mathrm{DDT}}`$ is indeed determined by the transition of burning regimes – which in turn might be composition dependent (Umeda et al. 1999) – the scenario is also fairly robust and $`\rho _{\mathrm{DDT}}`$ may represent the explosion strength parameter. Note that in this case, the variability induced by multipoint ignition needs to be explained away. If, on the other hand, thermonuclear flames are confirmed to be almost unquenchable, the favorite mechanism for DDTs becomes questionable (Niemeyer 1999). Moreover, should the mechanism DDT rely on rare, strong turbulent fluctuations one must ask about those events that fail to ignite a detonation following the slow deflagration phase which, on its own, cannot give rise to a viable SN Ia explosion. They might end up as pulsational delayed detonations or as unobservably dim, as yet unclassified explosions. Multidimensional simulations of the turbulent flame phase may soon answer whether the turbulent flame speed is closer to 1 % or 30 % of the speed of sound and hence decide whether DDTs are a necessary ingredient of SN Ia explosion models.
#### 5.1.6 PULSATIONAL DELAYED DETONATION
In this variety of the delayed detonation scenario, the first turbulent deflagration phase fails to release enough energy to unbind the star which subsequently pulses and triggers a detonation upon recollapse (Nomoto et al. 1976; Khokhlov 1991b). This model was studied in 1D by Höflich & Khokhlov (1996) and Woosley (1997b) (who calls it “pulsed detonation of the first type”) and in 2D by Arnett & Livne (1994b). Höflich & Khokhlov (1996) report that it produces little <sup>56</sup>Ni but a substantial amount of Si and Ca and may therefore explain very subluminous events like SN 1991bg. Woosley (1997b), using a fractal flame parameterization, also considered “pulsed deflagrations”, i.e. re-ignition occurs as a deflagration rather than a detonation, and “pulsed detonations of the second type” in which the burning also re-ignites as a flame but later accelerates and touches off a detonation . This latter model closely resembles the standard delayed detonation, whereas the former may or may not produce a healthy explosion, depending on the prescribed speed of the rekindled flame (Woosley 1997b).
Obtaining a DDT by means of the gradient mechanism is considerably more plausible after one or several pulses than during the first expansion phase (Khokhlov et al. 1997) as the laminar flame thickness becomes macroscopically large during the expansion, allowing the fuel to be preheated, and turbulence is significantly enhanced during the collapse.
The “checklist” for pulsational delayed detonations looks similar to that of simple delayed detonations (see above), with somewhat less emphasis on the improbability for DDT. Some fine-tuning of the initial flame speed is needed to obtain a large enough pulse in order to achieve a sufficient degree of mixing, while avoiding to unbind the star in a very weak explosion (Niemeyer & Woosley 1997). Again, these “fizzles” may be very subluminous and may have escaped discovery. We finally note that all pulsational models are in conflict with multidimensional simulations that predict an unbound star after the first deflagration phase.
### 5.2 Sub-Chandrasekhar Mass Models
C+O white dwarfs below the Chandrasekhar mass do not reach the critical density and temperature for explosive carbon burning by accretion, and therefore need to be ignited by an external trigger. Detonations in the accreted He layer were suggested to drive a strong enough shock into the C+O core to initiate a secondary carbon detonation (Weaver & Woosley 1980; Nomoto 1980, 1982a; Woosley, Weaver, & Taam 1980; Sutherland & Wheeler 1984; Iben & Tutukov 1984). The nucleosynthesis and light curves of Sub-$`M_{\mathrm{chan}}`$ models, also known as “helium ignitors” or “edge-lit detonations”, were investigated in 1D (Woosley & Weaver 1994b; Höflich & Khokhlov 1996) and 2D (Livne & Arnett 1995) and found to be superficially consistent with SNe Ia, especially subluminous ones (Ruiz-Lapuente et al. 1993a). Their ejecta structure is characterized almost inevitably by an outer layer of high-velocity Ni and He above the intermediate mass elements and the inner Fe/Ni core.
These models are favored mostly by the statistics of possible SN Ia progenitor systems (Yungelson & Livio 1998; Livio 2000) and by the straightforward explanation of the one-parameter strength sequence in terms of the WD mass (Ruiz-Lapuente, Burkert, & Canal 1995). However, they appear to be severely challenged both photometrically and spectroscopically: Owing to the heating by radioactive <sup>56</sup>Ni in the outer layer they are somewhat too blue at maximum brightness and their light curve rises and declines too steeply (Höflich & Khokhlov 1996; Nugent et al. 1997; Höflich et al. 1997). Perhaps even more stringent is the generic prediction of He-ignitors to exhibit signatures of high-velocity Ni and He, rather than Si and Ca, in the early and maximum spectra which is in strong disagreement with observations (Nugent et al. 1997; Höflich et al. 1997).
With respect to the explosion mechanism itself, the most crucial question is whether and where the He detonation manages to shock the C+O core sufficiently to create a carbon detonation. 1D models, by virtue of their built-in spherical symmetry, robustly (and unphysically) predict a perfect convergence of the inward propagating pressure wave and subsequent carbon ignition near the core (Woosley & Weaver 1994b). Some 2D simulations indicate that the C+O detonation is born off-center but still due to the convergence of the He-driven shock near the symmetry axis of the calculation (Livne 1990; Livne & Glasner 1991) while others find a direct initiation of the carbon detonation along the circle where the He detonation intersects the C+O core (Livne 1997; Arnett 1997; Wiggins & Falle 1997; Wiggins, Sharpe, & Falle 1998). Using 3D SPH simulations, Benz (1997) failed to see carbon ignition in all but the highest resolution calculations, where carbon was ignited directly at the interface rather than by shock convergence. Further, C ignition is facilitated if the He detonation starts at some distance above the interface, allowing the build-up of a fully developed pressure spike before it hits the carbon (Benz 1997). This result was confirmed by recent 3D SPH simulations (Garcia-Senz, Bravo, & Woosley 1999) that also examined the effect of multiple He ignition points, finding enhanced production of intermediate mass elements in this case. Hence, multidimensional SPH and PPM simulations presently confirm the validity of He-driven carbon detonations, in particular by direct ignition, but they also demonstrate the need for very high numerical resolution in order to obtain mutually consistent results (Arnett 1997; Benz 1997).
To summarize, sub-$`M_{\mathrm{chan}}`$ models are most severly constrained by their prediction of an outer layer of high-velocity Ni and He. Should further research conclude that spectra, colors, and light curves are less contaminated by this layer than presently thought, they represent an attractive class of candidates for SNe Ia, especially subluminous ones, from the point of view of progenitor statistics and the one-parameter explosion strength family. Note, however, that the SN Ia luminosity function in this scenario is directly linked to the distribution of white dwarf masses, predicting a more gradual decline on the bright side of the luminosity function than indicated by observations (Vaughan et al. 1995; Livio 2000). The explosion mechanism itself appears realistic, at least in the direct carbon ignition mode, but more work is needed to firmly establish the conditions for ignition of the secondary carbon detonation.
### 5.3 Merging White Dwarfs
The most obvious strength of the merging white dwarfs, or double-degenerate, scenario for SNe Ia (Webbink 1984; Iben & Tutukov 1984; Paczynski 1985) is the natural explanation for the lack of hydrogen in SN Ia spectra (Livio 2000) (cf. Sec. (2.1)). Furthermore, in contrast to the elusive progenitor systems for single degenerate scenarios, there is meanwhile some evidence for the existence of double degenerate binary systems (Saffer et al. 1998) despite earlier suspicions to the contrary (e.g., Bragaglia 1997). These systems are bound to merge as a consequence of gravitational wave emission with about the right statistics (Livio 2000) and give rise to some extreme astrophysical event, albeit not necessarily a SN Ia.
Spherically symmetric models of detonating merged systems, parameterized as C+O white dwarfs with thick envelopes, were analyzed by Höflich, Khokhlov, & Müller (1992); Khokhlov et al. (1993) and Höflich & Khokhlov (1996), giving reasonable agreement with SN Ia light curves. 3D SPH simulations of white dwarfs mergers (Benz et al. 1990; Rasio & Shapiro 1995; Mochkovitch, Guerrero, & Segretain 1997) show the disruption of the less massive star in a matter of a few orbital times, followed by the formation of a thick hot accretion disk around the more massive companion. The further evolution hinges crucially on the effective accretion rate of the disk: In case $`\dot{M}`$ is larger than a few times $`10^6`$ M yr<sup>-1</sup>, the most likely outcome is off-center carbon ignition leading to an inward propagating flame that converts the star into O+Ne+Mg (Nomoto & Iben 1985; Saio & Nomoto 1985; Kawai, Saio, & Nomoto 1987; Timmes, Woosley, & Taam 1994; Saio & Nomoto 1998). This configuration, in turn, is gravitationally unstable owing to electron capture onto <sup>24</sup>Mg and will undergo accretion-induced collapse (AIC) to form a neutron star (Saio & Nomoto 1985; Mochkovitch & Livio 1990; Nomoto & Kondo 1991). A recent re-examination of Coulomb corrections to the equation of state of material in nuclear statistical equilibrium indicates that AIC in merged white dwarf systems is even more likely than previously anticipated (Bravo & Garcia-Senz 1999).
Dimensional analysis of the expected turbulent viscosity due to MHD instabilities (Balbus, Hawley, & Stone 1996) suggests that it is very difficult to avoid such high accretion rates (Mochkovitch & Livio 1990; Livio 2000). Even under the unphysical assumption that angular momentum transport is dominated entirely by microscopic electron-gas viscosity, the expected life time of $`10^9`$ yrs (Mochkovitch & Livio 1990; Mochkovitch et al. 1997) and high UV luminosity of these accretion systems would predict the existence of $`10^7`$ such objects in the Galaxy, none of which have been observed (Livio 2000).
A possible solution to the collapse problem is to ignite carbon burning as a detonation rather than a flame immediately during the merger event, either in the core of the more massive star (Shigeyama et al. 1992) or at the contact surface (D Arnett & PA Pinto, private communication). This alternative clearly warrants further study.
To summarize, the merging white dwarf scenario has to overcome the crucial problem of avoiding accretion-induced collapse before it can be seriously considered as a SN Ia candidate. Its key strengths are a plausible explanation for the progenitor history yielding reasonable predictions for SN Ia rates, the straightforward explanation of the absence of H and He in SN Ia spectra, and the existence of a simple parameter for the explosion strength family (i.e., the mass of the merged system).
## 6 SUMMARY
In this review we have outlined our present understanding of Type Ia supernovae, summarizing briefly the observational constraints, but putting more weight on models of the explosion. From the tremendous amount of work carried out over the last couple of years it has become obvious that the physics of SNe Ia is very complex, ranging from the possibility of very different progenitors to the complexity of the physics leading to the explosion and the complicated processes which couple the interior physics to observable quantities. None of these problems is fully understood yet, but what one is tempted to state is that, from a theorist’s point of view, it appears to be a miracle that all the complexity seems to average out in a mysterious way to make the class so homogeneous. In contrast, as it stands, a save prediction from theory seems to be that SNe Ia should get more divers with increasing observed sample sizes. If, however, homogeneity would continue to hold this would certainly add support to the Chandrasekhar-mass single-degenerate scenario. On the other hand, even an increasing diversity would not rule out Chandrasekhar-mass single-degenerate progenitors for most of them. In contrast, there are ways to explain how the diversity is absorbed in a one parameter family of transformations, such as the Phillips-relation or modifications of it. For example, we have argued that the size of the convective core of the white dwarf prior to the explosion might provide a physical reason for such a relation.
As far as the explosion/combustion physics and the numerical simulations are concerned significant recent progress has made the models more realistic (and reliable). Thanks to ever increasing computer resources 3-dimensional simulations have become feasible which treat the full star with good spatial resolution and realistic input physics. Already the results of 2-dimensional simulations indicate that pure deflagrations waves in Chandrasekhar-mass C+O white dwarfs can lead to explosions, and one can expect that going to three dimensions, because of the increasing surface area of the nuclear flames, should add to the explosion energy. If confirmed, this would eliminate pulsational detonations from the list of potential models. On the side of the combustion physics, the burning in the distributed regime at low densities needs to be explored further, but it is not clear anymore whether a transition from a deflagration to a detonation in that regime is needed for successful models. In fact, according to recent studies such a transition appears to be rather unlikely.
Finally, sub-Chandrasekhar-mass models seem to face problems, both from the observations and from theory, leaving us with the conclusion that we seem to be lucky and Nature was kind to us and singled out from all possibilities the simplest solution, namely a Chandrasekhar-mass C+O white dwarf and a nuclear deflagration wave, to make a Type Ia supernova explosion.
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# LNF-00/017(P)UCRHEP-T278 Leptogenesis from Neutralino Decay with Nonholomorphic R-Parity Violation
## 1 Introduction
The creation of a lepton asymmetry, i.e. leptogenesis , which gets converted into the present observed baryon asymmetry of the Universe, is closely related to the mechanism by which neutrinos obtain mass. In general, all models of Majorana neutrino masses with the same low-energy particle content as that of the standard model are equivalent in the sense that they are all characterized by the same nonrenormalizable dimension-five operator $`\mathrm{\Lambda }^1\nu _i\nu _j\varphi ^0\varphi ^0`$. Different models of neutrino mass are merely different realizations of this operator. They become distinguishable only at high energies, and since their interactions must violate lepton number, leptogenesis is a very natural possibility. For the canonical seesaw mechanism and the Higgs triplet model , leptogenesis does indeed occur naturally . On the other hand, if neutrino masses are obtained radiatively , not only is leptogenesis difficult to achieve, the mechanism by which the former is accomplished leads naturally to the erasure of any primordial baryon asymmetry of the Universe . This is especially true in supersymmetric models of neutrino mass with R-parity violation. In a recent article , it was pointed out that leptogenesis is still possible in this case, provided that certain conditions regarding the R-parity violating terms are satisfied. Here we study this model in detail.
In Section 2 we write down the superpotential of the lepton-number violating (but baryon-number conserving) extension of the supersymmetric standard model, together with all possible soft supersymmetric breaking terms, including the nonholomorphic terms . In Section 3 we consider bilinear R-parity violation and how leptogenesis is related to neutrino mass in this limited scenario. We find it to be negligible for realistic values of $`m_\nu `$. In Section 4 we discuss how leptogenesis may occur without being constrained by neutrino mass in an expanded scenario. We assume negligible (enhanced) mixing between doublet (singlet) sleptons and charged Higgs bosons by allowing nonholomorphic soft supersymmetry breaking terms. In Section 5 we present the details of our calculations using the Boltzmann equations for obtaining the eventual lepton asymmetry. In Section 6, the complete charged-scalar mass matrix is displayed and analyzed. In Section 7, a new two-loop mechanism for neutrino mass is proposed. Finally in Section 8, there are some concluding remarks.
## 2 Superpotential and Soft Supersymmetry Breaking
In an unrestricted supersymmetric extension of the standard model of particle interactions, the chiral scalar superfields allow baryon-number violating terms which are not necessarily suppressed. These dangerous terms are usually avoided by assuming a conserved discrete quantum number for each particle called R-parity, which is defined as
$$\mathrm{R}(1)^{3B+L+2J},$$
(1)
where $`B`$ is its baryon number, $`L`$ its lepton number, and $`J`$ its spin angular momentum. With this definition, the standard-model particles have R = +1 and their supersymmetric partners have R = $`1`$. We can list the three families of leptons and quarks of the standard model using the notation where all superfields are considered left-handed:
$`L_i=(\nu _i,e_{i_L})(1,2,1/2),e_i^c(1,1,1),`$ (2)
$`Q_i=(u_i,d_i)(3,2,1/6),`$ (3)
$`u_i^c(3^{},1,2/3),d_i^c(3^{},1,1/3),`$ (4)
where $`i`$ is the family index, and the two Higgs doublets are given by
$`H_1=(h_1^0,h_1^{})(1,2,1/2),`$ (5)
$`H_2=(h_2^+,h_2^0)(1,2,1/2),`$ (6)
where the $`SU(3)_C\times SU(2)_L\times U(1)_Y`$ content of each superfield is also indicated. If R-parity is conserved, the superpotential is restricted to have only the terms
$$W=\mu H_1H_2+f_{ij}^eH_1L_ie_j^c+f_{ij}^dH_1Q_id_j^c+f_{ij}^uH_2Q_iu_j^c.$$
(7)
In this case, both baryon and lepton numbers are conserved. However, to forbid proton decay, it is sufficient to conserve either baryon number or lepton number (because the final state of the proton decay must contain a lepton or antilepton). If only baryon number or only lepton number is violated (thus R-parity is also violated), the conservation of the other quantum number is enough to satisfy all present experimental constraints. This has motivated numerous studies of R-parity violating models.
If R-parity is violated, the superpotential contains the additional terms
$$W^{}=\mu _iL_iH_2+\lambda _{ijk}L_iL_je_k^c+\lambda _{ijk}^{}L_iQ_jd_k^c+\lambda _{ijk}^{\prime \prime }u_i^cd_j^cd_k^c.$$
(8)
We cannot have all of these terms because then the proton will decay very quickly. We may choose only the lepton-number violating terms or only the baryon-number violating terms. Following the overwhelming choice of many others, we consider only the former case and set $`\lambda _{ijk}^{\prime \prime }=0`$. The remaining terms may now induce nonzero neutrino masses, either from mixing with the neutralino mass matrix, or in one-loop order. Although these terms are allowed, we do not know if they originate from any fundamental theory, so the couplings are considered free parameters, constrained only by experiment.
Other free parameters exist in the minimal supersymmetric standard model (MSSM), i.e. the soft supersymmetry breaking terms, which do not introduce quadratic divergences to the renormalized theory. Usually only the holomorphic terms are considered which come from the chiral superpotential interacting with gravity, together with the gaugino masses. The most general such Lagrangian conserving R-parity is:
$`_{soft}`$ $`=`$ $`\stackrel{~}{L}_i^a(M_L^2)_{ij}\stackrel{~}{L}_j^a\stackrel{~}{e}_i^c(M_e^2)_{ij}\stackrel{~}{e}_j^c\stackrel{~}{Q}_i^a(M_Q^2)_{ij}\stackrel{~}{Q}_j^a\stackrel{~}{u}_i^c(M_u^2)_{ij}\stackrel{~}{u}_j^c`$ (9)
$`\stackrel{~}{d}_i^c(M_d^2)_{ij}\stackrel{~}{d}_j^cM_{H_1}^2H_1^aH_1^aM_{H_2}^2H_2^aH_2^a\epsilon _{ab}(BH_1^aH_2^b+h.c.)`$
$`\epsilon _{ab}((A_ef_e)_{ij}H_1^a\stackrel{~}{L}_i^b\stackrel{~}{e}_j^c+(A_uf_u)_{ij}H_2^b\stackrel{~}{Q}_i^a\stackrel{~}{u}_j^c+(A_df_d)_{ij}H_1^a\stackrel{~}{Q}_i^b\stackrel{~}{d}_j^c+h.c.)`$
$`{\displaystyle \frac{1}{2}}(M_3\stackrel{~}{g}\stackrel{~}{g}+M_2\stackrel{~}{W}\stackrel{~}{W}+M_1\stackrel{~}{B}\stackrel{~}{B}+h.c.).`$
If R-parity is violated, more soft terms may be present, i.e.
$$_{soft}^{R/}=\epsilon _{ab}\left(B_i^{}\stackrel{~}{L}_i^aH_2^b+A_{ijk}^e\stackrel{~}{L}_i^a\stackrel{~}{L}_j^b\stackrel{~}{e}_k^c+A_{ijk}^d\stackrel{~}{Q}_i^a\stackrel{~}{L}_j^b\stackrel{~}{d}_k^c\right)A_{ijk}^S\stackrel{~}{u}_i^c\stackrel{~}{d}_j^c\stackrel{~}{d}_k^c+h.c..$$
(10)
We follow the convention that the coupling constants of all the R-parity conserving soft terms are denoted without a prime, while the R-parity violating terms are denoted with a prime.
Since the soft terms may originate from gravity couplings, there is no clear reason for them to come only from the renormalizable chiral superpotential. They may also originate from nonrenormalizable terms, which can be functions of both left and right chiral superfields. Such nonholomorphic terms allow new supersymmetry-breaking soft terms in the Lagrangian. The most general set of nonholomorphic soft terms conserving R-parity is:
$$_{soft}^{NH}=N_{ij}^eH_2^a\stackrel{~}{L}_i^a\stackrel{~}{e}_j^cN_{ij}^dH_2^a\stackrel{~}{Q}_i^a\stackrel{~}{d}_j^cN_{ij}^uH_1^a\stackrel{~}{Q}_i^a\stackrel{~}{u}_j^c+h.c..$$
(11)
Similarly, nonholomorphic soft terms breaking R-parity are:
$`_{soft}^{NHR/}`$ $`=`$ $`N_i^BH_1^a\stackrel{~}{L}_i^aN_i^eH_2^aH_1^a\stackrel{~}{e}_i^cN_{ijk}^u\stackrel{~}{L}_i^a\stackrel{~}{Q}_j^a\stackrel{~}{u}_k^c`$ (12)
$`N_{ijk}^S\stackrel{~}{u}_i^c\stackrel{~}{e}_j^c\stackrel{~}{d}_k^cN_{ijk}^d\epsilon _{ab}\stackrel{~}{Q}_i^a\stackrel{~}{Q}_j^b\stackrel{~}{d}_k^c+h.c.`$
In our convention, all “N” constants are for nonholomorphic terms. Since we assume lepton-number violation but baryon-number conservation, it implies $`\lambda _{ijk}^{\prime \prime }=A_{ijk}^S=N_{ijk}^d=0`$.
## 3 Bilinear R-Parity Violation
Before coming to the explicit model in the next section, we first look at the possibility of generating a lepton asymmetry from bilinear R-parity violating terms. This case is a good illustration of the general problems involved. One immediate consequence of the violation of lepton number through the bilinear terms is the mixing of the neutrinos with the neutralinos. In the MSSM there are four neutralinos, the U(1) gaugino ($`\stackrel{~}{B}`$), the SU(2) gaugino ($`\stackrel{~}{W}_3`$), and the two Higgsinos ($`\stackrel{~}{h}_1^0`$, $`\stackrel{~}{h}_2^0`$). When lepton number is violated through the R-parity violating terms, it is possible to assign zero lepton number to what we usually regard as lepton superfields . Suppose only the $`\tau `$ neutrino mixes with the neutralinos, then both $`\tau `$ and $`\nu _\tau `$ may be assigned an effective vanishing lepton number while the other leptons remain leptons. However, in the general three-family case, all three neutrinos may mix with the neutralinos, and the scalar partners of the neutrinos ($`\stackrel{~}{\nu }_i`$) may all acquire nonzero vacuum expectation values (VEVs).
In the most general case the neutralino mass matrix with all the seven fields in the basis $`[\stackrel{~}{B},\stackrel{~}{W}_3,\stackrel{~}{h}_1^0,\stackrel{~}{h}_2^0,\nu _1,\nu _2,\nu _3]`$ is given by
$$=\left[\begin{array}{ccccccc}M_1& 0& sr_Zv_1& sr_Zv_2& sr_Zv_{\nu _1}& sr_Zv_{\nu _2}& sr_Zv_{\nu _3}\\ 0& M_2& cr_Zv_1& cr_Zv_2& cr_Zv_{\nu _1}& cr_Zv_{\nu _2}& cr_Zv_{\nu _3}\\ sr_Zv_1& cr_Zv_1& 0& \mu & 0& 0& 0\\ sr_Zv_2& cr_Zv_2& \mu & 0& \mu _1& \mu _2& \mu _3\\ sr_Zv_{\nu _1}& cr_Zv_{\nu _1}& 0& \mu _1& 0& 0& 0\\ sr_Zv_{\nu _2}& cr_Zv_{\nu _2}& 0& \mu _2& 0& 0& 0\\ sr_Zv_{\nu _3}& cr_Zv_{\nu _3}& 0& \mu _3& 0& 0& 0\end{array}\right],$$
(13)
where $`s=\mathrm{sin}\theta _W`$, $`c=\mathrm{cos}\theta _W`$, $`r_Z=M_Z/v`$, and $`v_1`$, $`v_2`$, $`v_{\nu _i}`$ are the VEVs of $`h_1^0`$, $`h_2^0`$, and $`\stackrel{~}{\nu }_i`$ respectively, with $`v_1^2+v_2^2+v_\nu ^2=v^2`$ (246 GeV)<sup>2</sup> and $`v_\nu ^2=v_{\nu _1}^2+v_{\nu _2}^2+v_{\nu _3}^2`$. We also define $`\mathrm{tan}\beta =v_2/(v_1^2+v_\nu ^2)^{1/2}`$.
To understand the structure of the above $`7\times 7`$ mass matrix, let us assume that $`\mu `$ is the dominant term, then $`\stackrel{~}{h}_{1,2}^0`$ form a heavy Dirac particle of mass $`\mu `$ which mixes very little with the other physical fields. Removing these heavy fields will then give us the reduced $`5\times 5`$ matrix in the basis ($`\stackrel{~}{B},\stackrel{~}{W}_3,\nu _1,\nu _2,\nu _3`$):
$$=\left[\begin{array}{ccccc}M_1s^2\delta & sc\delta & sϵ_1& sϵ_2& sϵ_3\\ sc\delta & M_2c^2\delta & cϵ_1& cϵ_2& cϵ_3\\ sϵ_1& cϵ_1& 0& 0& 0\\ sϵ_2& cϵ_2& 0& 0& 0\\ sϵ_3& cϵ_3& 0& 0& 0\end{array}\right],$$
(14)
where
$`\delta `$ $`=`$ $`2M_Z^2{\displaystyle \frac{v_1v_2}{v^2}}{\displaystyle \frac{1}{\mu }}={\displaystyle \frac{M_Z^2\mathrm{sin}2\beta }{\mu }}\sqrt{1{\displaystyle \frac{v_\nu ^2}{v^2\mathrm{cos}^2\beta }}},`$ (15)
$`ϵ_i`$ $`=`$ $`{\displaystyle \frac{M_Z}{v}}\left(v_{\nu _i}{\displaystyle \frac{\mu _i}{\mu }}v_1\right).`$ (16)
From the above, only the combination $`\nu _l(ϵ_1\nu _1+ϵ_2\nu _2+ϵ_3\nu _3)/ϵ`$, with $`ϵ^2=ϵ_1^2+ϵ_2^2+ϵ_3^2`$, mixes with the gauginos. This state will have an effective vanishing lepton number and the other two orthogonal combinations decouple from the neutralino mass matrix. In this case, only the eigenstate
$$\nu _l^{}=\nu _l+\frac{sϵ}{M_1}\stackrel{~}{B}\frac{cϵ}{M_2}\stackrel{~}{W}_3,$$
(17)
gets a seesaw mass, i.e.
$$m_{\nu _l^{}}=ϵ^2\left(\frac{s^2}{M_1}+\frac{c^2}{M_2}\right),$$
(18)
whereas the other two neutrinos remain massless. They may get masses through one-loop radiative corrections from the usual trilinear R-parity violating terms which we have not yet considered.
The two gauginos mix with the neutrino $`\nu _l`$ and form mass eigenstates given by
$`\stackrel{~}{B}^{}`$ $`=`$ $`\stackrel{~}{B}+{\displaystyle \frac{sc\delta }{M_1M_2}}\stackrel{~}{W}_3{\displaystyle \frac{sϵ}{M_1}}\nu _l,`$ (19)
$`\stackrel{~}{W}_3^{}`$ $`=`$ $`\stackrel{~}{W}_3{\displaystyle \frac{sc\delta }{M_1M_2}}\stackrel{~}{B}+{\displaystyle \frac{cϵ}{M_2}}\nu _l.`$ (20)
The physical states $`\stackrel{~}{B}^{}`$ and $`\stackrel{~}{W}_3^{}`$ now contain $`\nu _l`$. This gives the main feature of R-parity violation, which is the decay of the lightest neutralino. By virtue of their $`\nu _l`$ components, both neutralinos will now decay into a lepton or an antilepton and a weak gauge boson, such as $`\stackrel{~}{W}_3^{}l^{}W^+`$ and $`l^+W^{}`$, thus violating lepton number. Since the mixing of the neutralinos may also have $`CP`$ violation through the complex gaugino masses (thus making $`\delta `$ complex), a lepton asymmetry may be generated from these decays. However, the amount of asymmetry thus generated is several orders of magnitude too small because it has to be much less than $`(ϵ/M_{1,2})^2`$, which is of order $`m_{\nu _l^{}}/M_{1,2}`$, i.e. $`<5\times 10^{13}`$ if $`m_{\nu _l^{}}<0.05`$ eV and $`M_{1,2}>100`$ GeV. In addition, the out-of-equilibrium condition on the decay width of the lightest neutralino imposes an upper bound on $`(ϵ/M_{1,2})^2`$ which is independent of $`m_{\nu _l^{}}`$, and that also results in an asymmetry very much less than $`10^{10}`$.
We now consider the R-parity violating trilinear couplings, i.e. $`\lambda `$ and $`\lambda ^{}`$ of Eq. (8). Since the particles involved should have masses at most equal to the supersymmetry breaking scale, i.e. a few TeV, their $`L`$ violation together with the $`B+L`$ violation by sphalerons would erase any primordial $`B`$ or $`L`$ asymmetry of the Universe . To avoid such a possibility, we may reduce $`\lambda `$ and $`\lambda ^{}`$ to less than about $`10^7`$, but a typical minimum value of $`10^4`$ is required for realistic neutrino masses in one-loop order . Hence it appears that the MSSM with R-parity violation is not only unsuitable for leptogenesis, it is also a destroyer of any lepton or baryon asymmetry which may have been created by some other means before the electroweak phase transition.
## 4 Leptogenesis from Neutralino Decay
From the discussion of the previous section we observe that for a leptogenesis mechanism to be successful in the MSSM with R-parity violation, two requirements have to be fulfilled. First we must use lepton-number violating terms which are not constrained by neutrino masses. Second we must satisfy the out-of-equilibrium condition for the lightest neutralino in such a way that the asymmetry is not automatically suppressed. More explicitly, we will consider the possibility that the heavier neutralino does not satisfy the out-of-equilibrium condition and decays very quickly, but the lighter neutralino decays very slowly and satisfies the out-of-equilibrium condition. Since the asymmetry comes from the interference of the one-loop $`CP`$ violating contribution of the heavier neutralino, it is then unsuppressed. We will demonstrate explicitly in the following how this scenario may be realized.
We assume first that $`M_1>M_2`$, so that the bino $`\stackrel{~}{B}`$ is heavier than the wino $`\stackrel{~}{W}_3`$. While the former couples to both $`\overline{e}_{i_L}\stackrel{~}{e}_{i_L}`$ and $`\overline{e}_i^c\stackrel{~}{e}_i^c`$, the latter couples only to $`\overline{e}_{i_L}\stackrel{~}{e}_{i_L}`$, because the $`e_i^c`$ are singlets under $`SU(2)_L`$. Since R-parity is violated, one combination of the $`\stackrel{~}{e}_{i_L}`$ and another of the $`\stackrel{~}{e}_i^c`$ mix with the charged Higgs boson of the supersymmetric standard model: $`h^\pm =h_2^\pm \mathrm{cos}\beta +h_1^\pm \mathrm{sin}\beta `$. Let us denote them by $`\stackrel{~}{l}_L`$ and $`\stackrel{~}{l}^c`$ respectively. Their corresponding leptons are of course $`l_L`$ and $`l^c`$. Hence both $`\stackrel{~}{B}^{}`$ and $`\stackrel{~}{W}_3^{}`$ may decay into $`l^{}h^\pm `$.
We assume next that the $`\stackrel{~}{l}_L`$ mixing with $`h^{}`$ is negligible, so that the only relevant coupling is that of $`\stackrel{~}{B}`$ to $`\overline{l}^ch^+`$. Hence $`\stackrel{~}{W}_3^{}`$ decay (into $`l^{}h^\pm `$) is suppressed because it may only do so through the small component of $`\stackrel{~}{B}`$ that it contains, assuming of course that all charged sleptons are heavier than $`\stackrel{~}{B}`$ or $`\stackrel{~}{W}_3`$.
With this choice that the heavier neutralino $`\stackrel{~}{B}^{}`$ decays quickly and the lighter neutralino $`\stackrel{~}{W}_3^{}`$ decays much more slowly, we now envisage the following leptogenesis scenario. At temperatures well above $`T=M_{SUSY}`$, there are fast lepton-number and R-parity violating interactions, which will wash out any $`L`$ or $`B`$ asymmetry of the Universe in the presence of sphalerons. This will be the case even at temperatures around $`M_1`$, when $`\stackrel{~}{W}_3^{}`$ interactions violate $`L_i`$ as well as $`B3L_i`$ for $`i=e,\mu ,\tau `$ . We assume here that all other supersymmetric particles are heavier than the neutralinos, so that at temperatures below $`M_1`$ we need only consider the interactions of $`\stackrel{~}{B}^{}`$ and $`\stackrel{~}{W}_3^{}`$. In Figure 1 we show the lepton-number violating processes (a) $`\stackrel{~}{B}^{}l_R^\pm h^{}`$, where we have adopted the more conventional notation of an outgoing $`l_R`$ in place of an incoming $`l^c`$. These processes are certainly still fast and there can be no $`L`$ asymmetry. At temperatures far below the mass of the heavier neutralino, the $`\stackrel{~}{B}^{}`$ interactions are suppressed and we need only consider those of $`\stackrel{~}{W}_3^{}`$. With our assumptions, the lepton-number violating processes (b) $`\stackrel{~}{W}_3^{}l_R^\pm h^{}`$ are slow and will satisfy the out-of-equilibrium condition for generating a lepton asymmetry of the Universe. Specifically, it comes from the interference of this tree-level diagram with the one-loop (c) self-energy and (d) vertex diagrams. Since the unsuppressed lepton-number violating couplings of $`\stackrel{~}{B}^{}`$ are involved, a realistic lepton asymmetry may be generated. It is then converted by the still active sphalerons into the present observed baryon asymmetry of the Universe. In this scenario the mass of $`\stackrel{~}{W}_3^{}`$ also has to be small enough so that the scattering processes mediated by the heavier $`\stackrel{~}{B}^{}`$ are negligible at temperature below $`M_2`$ when the asymmetry is produced.
We start with the well-known interaction of $`\stackrel{~}{B}`$ with $`l`$ and $`\stackrel{~}{l}_R`$ given by
$$\frac{e\sqrt{2}}{\mathrm{cos}\theta _W}[\overline{l}\left(\frac{1\gamma _5}{2}\right)\stackrel{~}{B}\stackrel{~}{l}_R+H.c.].$$
(21)
We then allow $`\stackrel{~}{l}_R`$ to mix with $`h^{}`$, and $`\stackrel{~}{B}`$ to mix with $`\stackrel{~}{W}_3`$, so that the interaction of the physical state $`\stackrel{~}{W}_3^{}`$ of Eq. (20) with $`l`$ and $`h^\pm `$ is given by
$$\left(\frac{sc\xi \delta r}{M_1M_2}\right)\left(\frac{e\sqrt{2}}{\mathrm{cos}\theta _W}\right)[\overline{l}\left(\frac{1\gamma _5}{2}\right)\stackrel{~}{W}_3^{}h^{}+H.c.],$$
(22)
where $`\xi `$ represents the $`\stackrel{~}{l}_Rh^{}`$ mixing and is assumed real, but the parameter $`\delta `$ of Eq. (15) is complex. We have also inserted a correction factor $`r=(1+M_2/\mu \mathrm{sin}2\beta )/(1M_2^2/\mu ^2)`$ for finite values of $`M_2/\mu `$. The origin of a nontrivial $`CP`$ phase in the above is from the $`2\times 2`$ Majorana mass matrix spanning $`\stackrel{~}{B}`$ and $`\stackrel{~}{W}_3`$, with complex $`M_1`$ and $`M_2`$. It is independent of the phase of $`\mu `$ and contributes negligibly to the neutron electric dipole moment because the magnitude of $`\delta `$ is very small. \[Note that the usual assumption of $`CP`$ violation in supersymmetric models is that $`M_1`$ and $`M_2`$ have a common phase, in which case the phase of $`\delta `$ would be equal to the phase of $`\mu `$.\]
The decay width of the bino is then
$$\mathrm{\Gamma }_{\stackrel{~}{B}^{}}=\mathrm{\Gamma }(\stackrel{~}{B}^{}l^+h^{})+\mathrm{\Gamma }(\stackrel{~}{B}^{}l^{}h^+)=\frac{1}{4\pi }\xi ^2\frac{e^2}{c^2}\frac{(M_{\stackrel{~}{B}^{}}^2m_h^2)^2}{M_{\stackrel{~}{B}^{}}^3},$$
(23)
while that of the wino is
$$\mathrm{\Gamma }_{\stackrel{~}{W}_3^{}}=\mathrm{\Gamma }(\stackrel{~}{W}_3^{}l^+h^{})+\mathrm{\Gamma }(\stackrel{~}{W}_3^{}l^{}h^+)=\frac{1}{4\pi }\xi ^2\left(\frac{es|\delta |r}{M_1M_2}\right)^2\frac{(M_{\stackrel{~}{W}_3^{}}^2m_h^2)^2}{M_{\stackrel{~}{W}_3^{}}^3}.$$
(24)
Using Eqs. (21) and (22), we calculate the interference between the tree-level and self-energy + vertex diagrams of Figure 1 and obtain the following asymmetry from the decay of $`\stackrel{~}{W}_3^{}`$:
$$ϵ=\frac{\mathrm{\Gamma }(\stackrel{~}{W}_3^{}l^+h^{})\mathrm{\Gamma }(\stackrel{~}{W}_3^{}l^{}h^+)}{\mathrm{\Gamma }_{\stackrel{~}{W}_3^{}}}=\frac{\alpha \xi ^2}{2\mathrm{cos}^2\theta _W}\frac{Im\delta ^2}{|\delta |^2}\left(1\frac{m_h^2}{M_{\stackrel{~}{W}_3^{}}^2}\right)^2\frac{x^{1/2}g(x)}{(1x)},$$
(25)
where $`x=M_{\stackrel{~}{W}_3^{}}^2/M_{\stackrel{~}{B}^{}}^2`$ and
$$g(x)=1+\frac{2(1x)}{x}\left[\left(\frac{1+x}{x}\right)\mathrm{ln}(1+x)1\right].$$
(26)
If the $`\stackrel{~}{W}_3^{}`$ interactions satisfy the out-of-equilibrium condition, then a lepton asymmetry may be generated from the above decay asymmetry. Note that in the above expression for $`ϵ`$, the parameter $`\delta `$ appears only in the combination $`\mathrm{Im}\delta ^2/|\delta |^2`$, which may be of order one. If the absolute value of $`\delta `$ is small, it slows down the decay rate of $`\stackrel{~}{W}_3^{}`$ and a departure from equilibrium may be achieved without affecting the amount of decay asymmetry generated in the process.
At the time this lepton asymmetry is generated, if the sphaleron interactions are still in equilibrium, they will convert it into a baryon asymmetry of the Universe . If the electroweak phase transition is strongly first-order, the sphaleron interactions freeze out at the critical temperature. Lattice simulations suggest that for a Higgs mass of around $`m_H70`$ GeV, the critical temperature is around $`T_c150`$ GeV . Higher values of $`m_H`$ will increase the critical temperature, but the increase is slower than linear. For example, for $`m_H150`$ GeV, the critical temperature could go up to $`T_c250`$ GeV.
For a second-order or weakly first-order phase transition<sup>1</sup><sup>1</sup>1 Note that we do not require the electroweak phase transition to be first-order for satisfying the out-of-equilibrium condition. See for example Ref. ., the sphaleron interactions freeze out at a temperature lower than the critical temperature. After the electroweak phase transition ($`T<T_c`$), the sphaleron transition rate is given by
$$\mathrm{\Gamma }_{sph}(T)=(2.2\times 10^4\kappa )\frac{[2M_W(T)]^7}{[4\pi \alpha _WT]^3}e^{E_{sph(T)}/T},$$
(27)
where
$$M_W(T)=\frac{1}{2}g_2v(T)=\frac{1}{2}g_2v(T=0)\left(1\frac{T^2}{T_c^2}\right)^{1/2},$$
(28)
and the free energy of the sphaleron is $`E_{sph}(T)(2M_W(T)/\alpha (W))B(m_h/M_W)`$, with $`B(0)=1.52`$, $`B(\mathrm{})=2.72`$ and $`\kappa =e^{3.6}`$ . In this case, the sphalerons freeze out at a temperature $`T_{out}`$ which is the temperature at which their interaction strength equals the expansion rate of the universe,
$$\mathrm{\Gamma }_{sph}(T_{out})=H(T_{out})=1.7\sqrt{g_{}}\frac{T_{out}^2}{M_{Pl}}.$$
(29)
For a critical temperature of about $`T_c250`$ GeV, the freeze-out temperature comes out to be around $`T_{out}200`$ GeV.
These discussions indicate that as long as the lepton asymmetry is generated at a temperature above, say 200 GeV, it will be converted to a baryon asymmetry of the Universe. Since the sphaleron interactions grow exponentially fast, they can convert a lepton asymmetry to a baryon asymmetry \[a $`(BL)`$ asymmetry to be precise\] by the time the temperature drops by only a few GeV. In the next section we discuss how the decay asymmetry of the neutralinos becomes a lepton asymmetry of the Universe.
## 5 Boltzmann Equations
We now solve the Boltzmann equations to estimate the amount of lepton asymmetry created after the decays of the neutralinos. When the decay of the $`\stackrel{~}{W}_3^{}`$ satisfies the out-of-equilibrium condition, i.e. when the decay rate is slower than the expansion rate of the Universe, the generated asymmetry is of the order of the decay asymmetry given in Eq. (25). This argument could replace the details of solving the Boltzmann equations for an order-of-magnitude estimate of the asymmetry in many scenarios. However, in the present case there are other constraints and depleting factors, and we need to solve the Boltzmann equations explicitly for a reliable estimate.
If the $`\stackrel{~}{W}_3^{}`$ decay rate is much less than the expansion rate of the Universe, the generated lepton asymmetry is the same as the decay asymmetry. In other words, the out-of-equilibrium condition reads
$$K_{\stackrel{~}{W}_3^{}}=\frac{\mathrm{\Gamma }_{\stackrel{~}{W}_3^{}}}{H(M_{\stackrel{~}{W}_3^{}})}1,$$
(30)
where $`H(T)`$ is the Hubble constant at the temperature $`T`$ and is given by
$$H(T)=\sqrt{\frac{4\pi ^3g_{}}{45}}\frac{T^2}{M_{Planck}},$$
(31)
with $`g_{}`$ the number of massless degrees of freedom which we take equal to $`106.75`$ and $`M_{Pl}10^{18}`$ GeV is the Planck scale. If this condition is satisfied, the lepton asymmetry is given by $`n_L=n_ln_{\overline{l}}ϵ/g_{}`$. But in practice, when $`K1`$, there is no time for the asymmetry to grow to its maximum value before the sphaleron transitions are over. So we need to study the case $`K1`$. Furthermore, a reasonable amount of asymmetry cannot be obtained unless the inverse decay and the scattering from bino exchange have rates lower than the expansion rate of the Universe. All these effects result in the further diminution of the lepton asymmetry and we need to solve the Boltzmann equations to take care of them properly.
At temperatures $`T<M_2`$, the decays of $`\stackrel{~}{W}_3^{}`$ given in Eq. (24) start generating an asymmetry. At this time there are important damping contributions coming from the inverse decays of $`\stackrel{~}{W}_3^{}`$ and $`\stackrel{~}{B}^{}`$ as well as the scattering processes $`l^\pm +h^{}\stackrel{~}{B}^{}l^{}+h^\pm `$. As we will see, the last two processes are especially important because $`\stackrel{~}{B}^{}`$ tends to remain in equilibrium and its presence washes out the created lepton asymmetry from the $`\stackrel{~}{W}_3^{}`$ decays. The reason is that its interactions are strong enough so that the Boltzmann exponential suppression of its number density may not be sufficient to compensate its large inverse decay and scattering cross sections. The effect of the scattering $`l^\pm +h^{}\stackrel{~}{W}_3^{}l^{}+h^\pm `$ is on the other hand negligible because it is suppressed by a factor of $`[(sc\delta r)/(M_1M_2)]^2`$ with respect to the scattering $`l^\pm +h^{}\stackrel{~}{B}^{}l^{}+h^\pm `$. Neglecting this term and defining the variable $`zM_{\stackrel{~}{W}_3}/T`$, the Boltzmann equations are then:
$`{\displaystyle \frac{dX_{\stackrel{~}{W}_3^{}}}{dz}}`$ $`=`$ $`\gamma _{\stackrel{~}{W}_3^{}}^{eq}{\displaystyle \frac{z}{sH(M_{\stackrel{~}{W}_3^{}})}}\left({\displaystyle \frac{X_{\stackrel{~}{W}_3^{}}}{X_{\stackrel{~}{W}_3^{}}^{eq}}}1\right),`$ (32)
$`{\displaystyle \frac{dX_L}{dz}}`$ $`=`$ $`\gamma _{\stackrel{~}{W}_3^{}}^{eq}{\displaystyle \frac{z}{sH(M_{\stackrel{~}{W}_3^{}})}}\left[\epsilon \left({\displaystyle \frac{X_{\stackrel{~}{W}_3^{}}}{X_{\stackrel{~}{W}_3^{}}^{eq}}}1\right){\displaystyle \frac{1}{2}}{\displaystyle \frac{X_L}{X_\gamma }}\right]`$ (33)
$`{\displaystyle \frac{z}{sH(M_{\stackrel{~}{B}^{}})}}\left({\displaystyle \frac{M_{\stackrel{~}{B}^{}}}{M_{\stackrel{~}{W}_3^{}}}}\right)^2\left[\gamma _{\stackrel{~}{B}^{}}^{eq}{\displaystyle \frac{1}{2}}{\displaystyle \frac{X_L}{X_\gamma }}+2{\displaystyle \frac{X_L}{X_\gamma }}\gamma _{scatt.}^{eq}\right],`$
where we have defined the number densities per comoving volume $`X_i=n_i/s`$ in terms of the number densities of particles “i” and
$$s=g_{}\frac{2\pi ^2}{45}T^3$$
(34)
is the entropy density. The equilibrium distributions of the number densities are given by the Maxwell-Boltzmann statistics:
$`n_{\stackrel{~}{W}_3^{}}`$ $`=`$ $`g_{\stackrel{~}{W}_3^{}}{\displaystyle \frac{M_{\stackrel{~}{W}_3^{}}^2}{2\pi ^2}}TK_2(M_{\stackrel{~}{W}_3^{}}/T),`$ (35)
$`n_{\stackrel{~}{B}^{}}`$ $`=`$ $`g_{\stackrel{~}{B}^{}}{\displaystyle \frac{M_{\stackrel{~}{B}^{}}^2}{2\pi ^2}}TK_2(M_{\stackrel{~}{B}^{}}/T),`$ (36)
$`n_\gamma `$ $`=`$ $`{\displaystyle \frac{g_\gamma T^3}{\pi ^2}},`$ (37)
where $`g_{\stackrel{~}{W}_3^{}}=1`$, $`g_{\stackrel{~}{B}^{}}=1`$, and $`g_\gamma =2`$ are the numbers of degrees of freedom of $`\stackrel{~}{W}_3^{}`$, $`\stackrel{~}{B}^{}`$, and the photon respectively.
The quantities $`\gamma _{\stackrel{~}{W}_3^{}}^{eq}`$ and $`\gamma _{\stackrel{~}{B}^{}}^{eq}`$ are the reaction densities for the decays and inverse decays of $`\stackrel{~}{W}_3^{}`$ and $`\stackrel{~}{B}^{}`$:
$`\gamma _{\stackrel{~}{W}_3^{}}^{eq}`$ $`=`$ $`n_{\stackrel{~}{W}_3^{}}^{eq}{\displaystyle \frac{K_1(M_{\stackrel{~}{W}_3^{}}/T)}{K_2(M_{\stackrel{~}{W}_3^{}}/T)}}\mathrm{\Gamma }_{\stackrel{~}{W}_3^{}},`$ (38)
$`\gamma _{\stackrel{~}{B}^{}}^{eq}`$ $`=`$ $`n_{\stackrel{~}{B}^{}}^{eq}{\displaystyle \frac{K_1(M_{\stackrel{~}{B}^{}}/T)}{K_2(M_{\stackrel{~}{B}^{}}/T)}}\mathrm{\Gamma }_{\stackrel{~}{B}^{}},`$ (39)
$`K_1`$ and $`K_2`$ being the usual modified Bessel functions. The reaction density for the scattering is given by
$$\gamma _{scatt.}^{eq}=\frac{T}{64\pi ^4}_{(m_h+m_l)^2}^{\mathrm{}}𝑑s\widehat{\sigma }_{\stackrel{~}{B}^{}}(s)\sqrt{s}K_1(\sqrt{s}/T),$$
(40)
where $`\widehat{\sigma }_{\stackrel{~}{B}^{}}`$ is the reduced cross section and is given by $`2[s(m_h+m_l)^2][s(m_hm_l)^2]\sigma _{\stackrel{~}{B}^{}}/s\mathrm{\hspace{0.17em}2}s\sigma _{\stackrel{~}{B}^{}}`$. The cross section $`\sigma _{\stackrel{~}{B}^{}}`$ does not contain the contribution of the on-mass-shell bino (which is already taken into account in the decay and inverse decay terms). This is achieved by replacing the usual propagator $`1/(sm^2+i\mathrm{\Gamma }m)`$ with the off-mass-shell propagator :
$$D_s^1=\frac{sm^2}{(sm^2)^2+\mathrm{\Gamma }^2m^2}.$$
(41)
The cross section $`\sigma _{\stackrel{~}{B}^{}}`$ contains the s- and t-channel contributions together with their interference terms and is given by
$`\sigma _{\stackrel{~}{B}^{}}\sigma (l^\pm h^{}\stackrel{~}{B}^{}l^{}h^\pm )=`$
$`{\displaystyle \frac{1}{8\pi s^2}}\left({\displaystyle \frac{e^2\xi ^2}{cos^2\theta _W}}\right)^2m_{\stackrel{~}{B}^{}}^2\left[{\displaystyle \frac{s^2}{D_s^2}}+{\displaystyle \frac{4s}{D_s}}+{\displaystyle \frac{2s}{M_{\stackrel{~}{B}^{}}^2}}\left(2+4{\displaystyle \frac{s+M_{\stackrel{~}{B}^{}}^2}{D_s}}\right)\mathrm{ln}\left(1+{\displaystyle \frac{s}{M_{\stackrel{~}{B}^{}}^2}}\right)\right]`$ (42)
In Eq. (40) the integral is dominated by the s-channel contribution in the resonance region and to a very good approximation, $`\gamma _{scatt.}^{eq}`$ reduces to
$$\gamma _{scatt.}^{eq}=\frac{T}{512\pi ^4}\frac{M_{\stackrel{~}{B}^{}}^4}{\mathrm{\Gamma }_{\stackrel{~}{B}^{}}}K_1(M_{\stackrel{~}{B}^{}}/T)\frac{e^4\xi ^4}{cos^4\theta _W}.$$
(43)
From Eqs. (34) to (43), we find the Boltzmann equations, i.e. (32) and (33), to be given by
$`{\displaystyle \frac{dX_{\stackrel{~}{W}_3^{}}}{dz}}`$ $`=`$ $`zK_{\stackrel{~}{W}_3^{}}{\displaystyle \frac{K_1(z)}{K_2(z)}}\left(X_{\stackrel{~}{W}_3^{}}X_{\stackrel{~}{W}_3^{}}^{eq}\right),`$
$`{\displaystyle \frac{dX_L}{dz}}`$ $`=`$ $`zK_{\stackrel{~}{W}_3^{}}{\displaystyle \frac{K_1(z)}{K_2(z)}}\left[\epsilon (X_{\stackrel{~}{W}_3^{}}X_{\stackrel{~}{W}_3^{}}^{eq}){\displaystyle \frac{1}{2}}{\displaystyle \frac{X_{\stackrel{~}{W}_3}}{X_\gamma }}X_L\right]`$ (44)
$`z\left({\displaystyle \frac{M_{\stackrel{~}{B}^{}}}{M_{\stackrel{~}{W}_3^{}}}}\right)^2K_{\stackrel{~}{B}^{}}\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{K_1(zM_{\stackrel{~}{B}^{}}/M_{\stackrel{~}{W}_3^{}})}{K_2(zM_{\stackrel{~}{B}^{}}/M_{\stackrel{~}{W}_3^{}})}}{\displaystyle \frac{X_{\stackrel{~}{B}^{}}}{X_\gamma }}X_L+2{\displaystyle \frac{X_L}{X_\gamma }}{\displaystyle \frac{\gamma _{scatt.}^{eq}}{s\mathrm{\Gamma }_{\stackrel{~}{B}^{}}}}\right],`$
with
$$K_{\stackrel{~}{B}^{}}=\mathrm{\Gamma }_{\stackrel{~}{B}^{}}/H(M_{\stackrel{~}{B}^{}}),$$
(45)
which gives the strength of lepton-number violation in the decays of $`\stackrel{~}{B}^{}`$.
If we now ignore the inverse decay and scattering processes, we can simplify the problem by requiring the out-of-equilibrium condition to be
$$K_{\stackrel{~}{W}_3^{}}<1.$$
(46)
With the terms proportional to $`K_{\stackrel{~}{B}^{}}`$ this condition is necessary but not sufficient. In the present scenario for a large asymmetry we also require $`K_{\stackrel{~}{B}^{}}>1`$. Indeed, $`K_{\stackrel{~}{B}^{}}`$ is larger than $`K_{\stackrel{~}{W}_3^{}}`$ by a factor $`R_K=K_{\stackrel{~}{B}^{}}/K_{\stackrel{~}{W}_3^{}}[(sc\delta r)/(M_1M_2)]^2(M_{\stackrel{~}{W}_3^{}}/M_{\stackrel{~}{B}^{}})`$, which is larger than one by several orders of magnitude. Therefore $`\stackrel{~}{B}^{}`$ remains in equilibrium and the $`\stackrel{~}{B}^{}`$ damping terms, due to its inverse decay and scattering, dominate over the $`\stackrel{~}{W}_3^{}`$ inverse decay damping term as long as the Boltzmann suppression factor in the $`\stackrel{~}{B}^{}`$ equilibrium distribution has not compensated the large value of $`R_K`$. For example with the set of parameters $`M_{\stackrel{~}{W}_3^{}}=2`$ TeV, $`M_{\stackrel{~}{B}^{}}=3`$ TeV, $`\mathrm{sin}2\beta =0.5`$, $`\xi =2\times 10^3`$, $`\mu =5`$ TeV used in Ref. , we obtain $`K_{\stackrel{~}{W}_3^{}}=0.63`$ and $`K_{\stackrel{~}{B}^{}}=7.8\times 10^5`$, and the $`\stackrel{~}{B}^{}`$ damping terms dominate over those of $`\stackrel{~}{W}_3^{}`$ as long as the temperature is above $`65`$ GeV (the former differs from the latter by a factor $`R_K(M_{\stackrel{~}{B}^{}}/M_{\stackrel{~}{W}_3^{}})^3e^{(M_{\stackrel{~}{B}^{}}M_{\stackrel{~}{W}_3^{}})/T}`$). In this case the inverse decay and scattering of $`\stackrel{~}{B}^{}`$ cause a considerable wash-out of the asymmetry because $`\stackrel{~}{W}_3^{}`$ has mostly decayed away already at temperatures well above $``$ 65 GeV. This is illustrated in Figure 2 showing the effects of various terms in the Boltzmann equations.
To avoid this wash-out, the value of $`M_{\stackrel{~}{B}^{}}`$ has to be larger in order that the $`\stackrel{~}{B}^{}`$ number density is further suppressed at temperatures below $`M_{\stackrel{~}{W}_3^{}}`$ when the asymmetry is produced. Varying the parameters of these $`\stackrel{~}{B}^{}`$ damping terms, it appears difficult to induce a sufficiently large asymmetry of order $`10^{10}`$ for $`M_{\stackrel{~}{B}^{}}`$ below 4 TeV. Two typical situations for which a large asymmetry is produced are for example:
$`M_{\stackrel{~}{B}^{}}=6\mathrm{TeV},M_{\stackrel{~}{W}_3^{}}=3.5\mathrm{TeV},\xi =5\times 10^3,`$
$`\mu =10\mathrm{TeV},\mathrm{sin}2\beta =0.10,m_h=200\mathrm{GeV},`$ (47)
and
$`M_{\stackrel{~}{B}^{}}=5\mathrm{TeV},M_{\stackrel{~}{W}_3^{}}=2\mathrm{TeV},\xi =5\times 10^3,`$
$`\mu =7.5\mathrm{TeV},\mathrm{sin}2\beta =0.05,m_h=200\mathrm{GeV},`$ (48)
for which we have $`K_{\stackrel{~}{W}_3^{}}=0.02`$, $`K_{\stackrel{~}{B}^{}}^{}=2.4\times 10^6`$ and $`K_{\stackrel{~}{W}_3^{}}=0.02`$, $`K_{\stackrel{~}{B}^{}}^{}=2.9\times 10^6`$ respectively. At $`T=M_Z`$ the leptonic asymmetry produced is $`X_L=1.0\times 10^{10}`$ with the parameters of Eqs. (47) and $`X_L=1.2\times 10^{10}`$ with those of Eq. (48). Figures 3 and 4 show the evolution of the asymmetry in these two cases. As can be seen from these figures, the damping effects of the inverse decay of $`\stackrel{~}{W}_3^{}`$ and of the scattering are small<sup>2</sup><sup>2</sup>2 Note that in the Boltzmann equations we neglected the damping contributions of the scatterings $`l^\pm l^\pm h^\pm h^\pm `$ mediated by a neutralino in the t-channel. Their effect is negligible within the ranges of the parameters we consider.. The damping effect from the inverse decay of $`\stackrel{~}{B}^{}`$ is however not small and reduces the asymmetry by a factor of 2 to 4 by washing out all the asymmetry produced above $`T300400`$ GeV.
A large asymmetry of order $`10^{10}`$ is produced provided $`M_{\stackrel{~}{B}^{}}`$ is of the order 4 TeV or more. A low value of $`\mathrm{sin}2\beta `$ below $`0.30`$ is generally necessary. Values of $`\xi `$ around $`35\times 10^3`$, of $`\mu `$ around $`510`$ TeV, and of $`M_{\stackrel{~}{W}_3^{}}`$ from 1 TeV to $`2M_{\stackrel{~}{B}^{}}/3`$ are also preferred.
## 6 Charged Scalar Mass Matrix
The mechanism we propose for leptogenesis requires the decay of $`\stackrel{~}{B}^{}`$ to be fast, while that of $`\stackrel{~}{W}_3^{}`$ is very slow. This is achieved by requiring $`l_R^\pm h^{}`$ to be the main decay mode and $`l_L^\pm h^{}`$ to be negligible. Hence $`\stackrel{~}{l}_R`$ must mix with $`h^{}`$ readily, so that $`\stackrel{~}{B}^{}`$ could decay directly, but $`\stackrel{~}{W}_3^{}`$ could decay only through its small $`\stackrel{~}{B}`$ component. We now consider the charged scalar mass matrix which determines this mixing. As shown in the following, our present scenario requires one more new ingredient, i.e. the presence of nonholomorphic soft terms.
The value of the $`\tau _Rh^+`$ mixing parameter $`\xi `$ is governed by the charged scalar mass matrix which follows from the quadratic terms in the Lagrangian:
$$\mathrm{\Phi }^{}_{S^\pm }^2\mathrm{\Phi },$$
(49)
with $`\mathrm{\Phi }=[h_1^{},h_2^+,\stackrel{~}{e}_{i_L}^{},\stackrel{~}{e}_i^c]^T`$. In the case where all $`N`$ constants are put to zero and for parameters satisfying the various constraints from the Boltzmann equations (see previous section), it is difficult to generate a sufficiently large value of the $`\stackrel{~}{l}^ch^{}`$ mixing parameter $`\xi `$ (see Appendix A). Thus we will ignore the $`\mu _i`$ terms and the associated vaccuum expectation values $`v_{\nu _i}`$ in the following. To induce a large $`\stackrel{~}{e}^ch^{}`$ mixing, we introduce the nonholomorphic terms of Eqs. (11) and (12). The mass matrix is then given by
$$_{S^\pm }^2=\left(\begin{array}{cccc}_{h_1^{}h_1^+}^2& _{h_1^{}h_2^+}^2& _{h_1^{}\stackrel{~}{e}_{i_L}^{}{}_{}{}^{}}^2& _{h_1^{}\stackrel{~}{e}_i^c}^2\\ _{h_2^{}h_1^+}^2& _{h_2^{}h_2^+}^2& _{h_2^{}\stackrel{~}{e}_{i_L}^{}{}_{}{}^{}}^2& _{h_2^{}\stackrel{~}{e}_i^c}^2\\ _{\stackrel{~}{e}_{j_L}h_1^+}^2& _{\stackrel{~}{e}_{j_L}h_2^+}^2& _{\stackrel{~}{e}_{j_L}\stackrel{~}{e}_{i_L}^{}}^2& _{\stackrel{~}{e}_{j_L}\stackrel{~}{e}_i^c}^2\\ _{\stackrel{~}{e}_j^ch_1^+}^2& _{\stackrel{~}{e}_j^ch_2^+}^2& _{\stackrel{~}{e}_j^c\stackrel{~}{e}_{i_L}^{}}^2& _{\stackrel{~}{e}_j^c\stackrel{~}{e}_i^c}^2\end{array}\right),$$
(50)
where
$`_{h_1^{}h_1^+}^2`$ $`=`$ $`{\displaystyle \frac{g^2}{4}}v_2^2B{\displaystyle \frac{v_2}{v_1}},`$ (51)
$`_{h_1^{}h_2^+}^2`$ $`=`$ $`{\displaystyle \frac{g^2}{4}}v_1v_2B,`$ (52)
$`_{h_1^{}\stackrel{~}{e}_{i_L}^{}}^2`$ $`=`$ $`N_i^B,`$ (53)
$`_{h_1^{}\stackrel{~}{e}_i^c}^2`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}N_i^ev_2,`$ (54)
$`_{h_2^{}h_2^+}^2`$ $`=`$ $`{\displaystyle \frac{g^2}{4}}v_1^2B{\displaystyle \frac{v_1}{v_2}},`$ (55)
$`_{h_2^{}\stackrel{~}{e}_{i_L}^{}}^2`$ $`=`$ $`B_i^{},`$ (56)
$`_{h_2^{}\stackrel{~}{e}_i^c}^2`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}N_i^ev_1,`$ (57)
$`_{\stackrel{~}{e}_{j_L}\stackrel{~}{e}_{i_L}^{}}^2`$ $`=`$ $`(M_L^2)_{ij}{\displaystyle \frac{1}{8}}(g^2g^2)(v_1^2v_2^2)\delta _{ij}+{\displaystyle \frac{1}{2}}f_{jk}^ef_{ik}^ev_1^2,`$ (58)
$`_{\stackrel{~}{e}_{j_L}\stackrel{~}{e}_i^c}^2`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}f_{ji}^e\mu v_2+{\displaystyle \frac{1}{\sqrt{2}}}(A_ef_e)_{ji}v_1+{\displaystyle \frac{1}{\sqrt{2}}}N_{ji}^ev_2,`$ (59)
$`_{\stackrel{~}{e}_j^c\stackrel{~}{e}_i^c}^2`$ $`=`$ $`(M_{\stackrel{~}{e}^c}^2)_{ji}{\displaystyle \frac{g^2}{4}}(v_1^2v_2^2)\delta _{ij}+{\displaystyle \frac{1}{2}}f_{ki}^ef_{kj}^ev_1^2.`$ (60)
In the following we assume that $`N_i^B`$ and $`B_i^{}`$ are negligible so that they do not induce a large mixing of the left-handed charged sleptons with the charged Higgs boson. We are then left with the $`N_i^e`$ and $`N_{ji}^e`$ terms. Going to the basis of the physical charged Higgs boson $`h^+`$ and of the Goldstone boson $`G^+`$, the latter decouples and in the basis $`[h^+,\stackrel{~}{e}_{i_L}^{},\stackrel{~}{e}_i^c]`$ we obtain the mass matrix
$$_{S^\pm }^2=\left(\begin{array}{ccc}m_W^22\frac{B}{\mathrm{sin}2\beta }& 0& \frac{1}{\sqrt{2}}N_i^ev\\ 0& _{\stackrel{~}{e}_{j_L}\stackrel{~}{e}_{i_L}^{}}^2& _{\stackrel{~}{e}_{j_L}\stackrel{~}{e}_i^c}^2\\ \frac{1}{\sqrt{2}}N_i^ev& _{\stackrel{~}{e}_j^c\stackrel{~}{e}_{i_L}^{}}^2& _{\stackrel{~}{e}_j^c\stackrel{~}{e}_i^c}^2\end{array}\right).$$
(61)
We observe that only one combination of the charged right-handed sleptons mixes with the charged Higgs boson:
$$\stackrel{~}{l}^c=\frac{N_1^e\stackrel{~}{e}_1^c+N_2^e\stackrel{~}{e}_2^c+N_3^e\stackrel{~}{e}_3^c}{N^e},$$
(62)
with $`(N^e)^2=(N_1^e)^2+(N_2^e)^2+(N_3^e)^2`$. In the Lagrangian the mass term which couples the charged Higgs boson with the sleptons reduces therefore to the following single term:
$$\frac{1}{\sqrt{2}}vN^eh^+\stackrel{~}{l}^c+h.c.$$
(63)
In the case of one family, we obtain to a very good approximation:
$`\xi `$ $`=`$ $`{\displaystyle \frac{\frac{1}{\sqrt{2}}N^ev}{m_{h^+}^2M_{\stackrel{~}{e}_L\stackrel{~}{e}_L^{}}^2\frac{\left(M_{\stackrel{~}{e}_L\stackrel{~}{e}^c}^2\right)^2}{m_{h^+}^2M_{\stackrel{~}{e}_L\stackrel{~}{e}_L^{}}^2}}},`$ (64)
$`\xi ^{}`$ $`=`$ $`{\displaystyle \frac{M_{\stackrel{~}{e}_L\stackrel{~}{e}^c}^2}{m_{h^+}^2M_{\stackrel{~}{e}_L\stackrel{~}{e}_L^{}}^2}}\xi ,`$ (65)
where $`\xi ^{}`$ is the $`h^+\stackrel{~}{e}_L`$ mixing and
$$m_{h^+}^2=m_W^22\frac{B}{\mathrm{sin}2\beta }.$$
(66)
As required in section 3, the mixing $`\xi ^{}`$ has to be much smaller than the mixing $`\xi `$ in order to avoid having the transition $`\stackrel{~}{W}_3^{}\stackrel{~}{e}_L^\pm e_L^{}h^\pm e_L^{}`$. This requires
$$|\xi ^{}|<\left|\xi \frac{2s^2\delta r}{M_1M_2}\right|$$
(67)
which for the parameters of Eq. (47) implies that $`\xi ^{}<7\times 10^5\xi `$. Our mechanism requires also that $`m_{h^+}<M_{\stackrel{~}{W}_3^{}}`$ and that the mass of any charged slepton is larger than $`M_{\stackrel{~}{W}_3^{}}`$. We then have
$$\xi \frac{\frac{1}{\sqrt{2}}N^ev}{m_{h^+}^2M_{\stackrel{~}{e}_L\stackrel{~}{e}_L^{}}^2},$$
(68)
$$\left|m_e\mu \mathrm{tan}\beta +A_em_e+\frac{1}{\sqrt{2}}N_ev_2\right|<M_{\stackrel{~}{e}_L\stackrel{~}{e}_L^{}}^2\left|\frac{2s^2\delta r}{M_1M_2}\right|,$$
(69)
which for the set of parameters given in Eq. (47) and for $`M_{\stackrel{~}{e}_L\stackrel{~}{e}_L^{}}^2M_{\stackrel{~}{B}^{}}^2`$ gives $`N^e1`$ TeV and $`|m_e\mu \mathrm{tan}\beta +A_em_e+\frac{1}{\sqrt{2}}N_ev_2|<(50\mathrm{GeV})^2`$. The latter condition requires $`m_e<0.02`$ GeV, $`|A_em_e|<(50\mathrm{GeV})^2`$, and $`|N_e|<15`$ GeV, or a cancellation of the three terms together. In the case where there is no cancellation between these terms the charged slepton which mixs with the charged Higgs boson must have predominantly an electron or $`\mu `$ flavor. In the case of a cancellation between these terms, all flavors are possible. For other sets of parameters which lead to large asymmetries, these numerical bounds can be relaxed easily by a factor of 2 to 4. In the more general case of three families, similar constraints and relations are obtained.
## 7 Two-Loop Neutrino Mass
It is interesting to note that in addition to inducing a lepton asymmetry, the nonholomorphic terms $`N_i^e`$ could also generate a neutrino mass. Since lepton number is violated at most by one unit in each term, the neutrino mass should include at least two lepton-violating vertices in a loop diagram.
There exists a one-loop diagram (Fig. 5), which contributes to the sneutrino “Majorana” mass. In general, the sneutrinos could have diagonal lepton number conserving masses, i.e. $`\stackrel{~}{\nu }^{}\stackrel{~}{\nu }`$. But there can be also lepton-number violating mass terms, i.e. $`\stackrel{~}{\nu }\stackrel{~}{\nu }`$ . In the present model, the nonholomorphic terms give rise to a lepton-number violating sneutrino-antisneutrino mixing term, i.e. $`\frac{1}{2}\delta m_{\stackrel{~}{\nu }_{ij}}^2\stackrel{~}{\nu }_i\stackrel{~}{\nu }_j+h.c.`$. In the case of one family, we get
$$\delta m_{\stackrel{~}{\nu }}^2\frac{1}{8\pi ^2}\frac{\mu ^2\xi ^2}{v^2}m_l^2.$$
(70)
This lepton-number violating sneutrino mass can then induce a Majorana neutrino mass .
In the present case the one-loop diagram of Fig. 6 gives a neutrino mass
$`m_\nu `$ $``$ $`{\displaystyle \frac{1}{32\pi ^2}}{\displaystyle \frac{e^2}{\mathrm{sin}^2\theta _W}}\delta m_{\stackrel{~}{\nu }}^2M_{\stackrel{~}{W}_3^{}}{\displaystyle \frac{M_{\stackrel{~}{\nu }}^2M_{\stackrel{~}{W}_3^{}}^2M_{\stackrel{~}{W}_3^{}}^2\mathrm{ln}\left(M_{\stackrel{~}{\nu }}^2/M_{\stackrel{~}{W}_3^{}}^2\right)}{\left(M_{\stackrel{~}{\nu }}^2M_{\stackrel{~}{W}_3^{}}^2\right)^2}}`$ (71)
$``$ $`{\displaystyle \frac{1}{256\pi ^4}}{\displaystyle \frac{e^2}{\mathrm{sin}^2\theta _W}}\mu ^2{\displaystyle \frac{m_l^2}{v^2}}\xi ^2M_{\stackrel{~}{W}_3^{}}{\displaystyle \frac{M_{\stackrel{~}{\nu }}^2M_{\stackrel{~}{W}_3^{}}^2M_{\stackrel{~}{W}_3^{}}^2\mathrm{ln}\left(M_{\stackrel{~}{\nu }}^2/M_{\stackrel{~}{W}_3^{}}^2\right)}{\left(M_{\stackrel{~}{\nu }}^2M_{\stackrel{~}{W}_3^{}}^2\right)^2}}.`$
In the case where the lepton $`l`$ which mixes with $`h^+`$ is essentially $`\tau `$, we get
$`m_\nu `$ $`=`$ $`{\displaystyle \frac{1}{256\pi ^4}}{\displaystyle \frac{e^2}{\mathrm{sin}^2\theta _W}}\mu ^2{\displaystyle \frac{m_\tau ^2}{v^2}}\xi ^2M_{\stackrel{~}{W}_3^{}}{\displaystyle \frac{M_{\stackrel{~}{\nu }}^2M_{\stackrel{~}{W}_3^{}}^2M_{\stackrel{~}{W}_3^{}}^2\mathrm{ln}\left(M_{\stackrel{~}{\nu }}^2/M_{\stackrel{~}{W}_3^{}}^2\right)}{\left(M_{\stackrel{~}{\nu }}^2M_{\stackrel{~}{W}_3^{}}^2\right)^2}},`$ (72)
which has the correct order of magnitude. For example with the parameters of Eq. (47) and taking $`M_{\stackrel{~}{\nu }}M_{\stackrel{~}{B}^{}}`$ we get $`m_{\nu _\tau }0.1`$ eV. The value of $`\xi `$ we need for having the right order of magnitude for the asymmetry is therefore also the one we need to have a neutrino mass, in agreement with the present data on atmospheric neutrinos. In Eq. (71), the factor $`\delta m_{\stackrel{~}{\nu }}^2`$ appears because there is GIM (Glashow-Iliopoulos-Maiani) suppression from summing over all possible neutral slepton eigenstates in the loop. In Fig. 5 the two-point functions of the form $`f(m^2,m^2,p)(i/\pi ^2)d^4k(k^2m^2)^1((k+p)^2m^2)^1`$ have been (roughly) approximated by $`1`$ while in Fig. 6 the two-point functions have been calculated explicitly (as required by the fact that for these diagrams, a GIM suppression mechanism is operative). This can also be understood from another point of view. Since the diagonal terms of the sneutrino mass come from the lepton-number conserving interactions, they should not contribute to the Majorana mass of a neutrino. Only the lepton-number violating sneutrino mass, which is the mass-squared difference, should contribute to the Majorana neutrino mass. This makes the neutrino mass proportional to the mass-squared difference after GIM cancellation. Note that combining both one-loop diagrams, a two-loop diagram is obtained which is similar to the diagram proposed in Ref. with a different lepton-number violating soft term.
In the case of three families, the induced neutrino mass terms involving $`\nu _e`$ are suppressed by the small value of the electron mass with respect to the $`\mu `$ or $`\tau `$ mass. Therefore, unless $`N_e^e`$ is much larger than $`N_\tau ^e`$ and $`N_\mu ^e`$, $`\nu _e`$ essentially decouples and acquires a very small mass; we get $`m_{\nu _e}10^8`$ eV or less. In this case, the mass matrix of the sneutrinos in the $`\mu \tau `$ sector is of the form:
$$\frac{1}{2}\mathrm{\Phi }_{\stackrel{~}{\nu }}^{}_{\stackrel{~}{\nu }}^2\mathrm{\Phi }_{\stackrel{~}{\nu }},$$
(73)
with $`\mathrm{\Phi }_{\stackrel{~}{\nu }}=(\stackrel{~}{\nu }_\mu ,\stackrel{~}{\nu }_\tau ,\stackrel{~}{\nu }_\mu ^{},\stackrel{~}{\nu }_\tau ^{})^T`$ and
$$_{\stackrel{~}{\nu }}^2=\left(\begin{array}{cccc}M_{L_\mu }^2& 0& \delta m_{\stackrel{~}{\nu }_{\mu \mu }}^2& \delta m_{\stackrel{~}{\nu }_{\mu \tau }}^2\\ 0& M_{L_\tau }^2& \delta m_{\stackrel{~}{\nu }_{\mu \tau }}^2& \delta m_{\stackrel{~}{\nu }_{\tau \tau }}^2\\ \delta m_{\stackrel{~}{\nu }_{\mu \mu }}^2& \delta m_{\stackrel{~}{\nu }_{\mu \tau }}^2& M_{L_\mu }^2& 0\\ \delta m_{\stackrel{~}{\nu }_{\mu \tau }}^2& \delta m_{\stackrel{~}{\nu }_{\tau \tau }}^2& 0& M_{L_\tau }^2\end{array}\right),$$
(74)
where for simplicity we have assumed in Eq. (9) a diagonal matrix $`M_L^2=diag(M_{L_\mu }^2,M_{L_\tau }^2)`$ and with
$$\delta m_{\stackrel{~}{\nu }_{ij}}^2\frac{1}{8\pi ^2}\frac{\mu ^2\xi ^2}{v^2}m_{l_i}m_{l_j}\frac{N_i^eN_j^e}{(N^e)^2}.$$
(75)
From the mass matrix of Eq. (74), the diagrams of Fig. 6 induce then the following neutrino mass term
$$\frac{1}{2}\mathrm{\Psi }_\nu ^{}_\nu \mathrm{\Psi }_\nu ,$$
(76)
with $`\mathrm{\Psi }_\nu =(\nu _\mu ,\nu _\mu ^{},\nu _\tau ,\nu _\tau ^{})^T`$ and
$$_\nu A\left(\begin{array}{cccc}0& \mathrm{\Delta }_{\mu \mu }& 0& \mathrm{\Delta }_{\mu \tau }\\ \mathrm{\Delta }_{\mu \mu }& 0& \mathrm{\Delta }_{\mu \tau }& 0\\ 0& \mathrm{\Delta }_{\mu \tau }& 0& \mathrm{\Delta }_{\tau \tau }\\ \mathrm{\Delta }_{\mu \tau }& 0& \mathrm{\Delta }_{\tau \tau }& 0\end{array}\right),$$
(77)
where
$`\mathrm{\Delta }_{ii}`$ $`=`$ $`M_{\stackrel{~}{W}_3^{}}\left(f(M_{L_i}^2+\delta m_{\stackrel{~}{\nu }_{ii}}^2,M_{\stackrel{~}{W}_3^{}}^2,0)f(M_{L_i}^2\delta m_{\stackrel{~}{\nu }_{ii}}^2,M_{\stackrel{~}{W}_3^{}}^2,0)\right)`$ (78)
$``$ $`2\delta m_{\stackrel{~}{\nu }_{ii}}^2M_{\stackrel{~}{W}_3^{}}{\displaystyle \frac{M_{L_i}^2M_{\stackrel{~}{W}_3^{}}^2M_{\stackrel{~}{W}_3^{}}^2\mathrm{ln}\left(M_{L_i}^2/M_{\stackrel{~}{W}_3^{}}^2\right)}{\left(M_{L_i}^2M_{\stackrel{~}{W}_3^{}}^2\right)^2}},`$
$`\mathrm{\Delta }_{\mu \tau }`$ $`=`$ $`{\displaystyle \frac{2\delta m_{\stackrel{~}{\nu }_{\mu \tau }}^2}{M_{L_\tau }^2M_{L_\mu }^2}}M_{\stackrel{~}{W}_3^{}}\left(f(M_{L_\tau }^2,M_{\stackrel{~}{W}_3^{}}^2,0)f(M_{L_\mu }^2,M_{\stackrel{~}{W}_3^{}}^2,0)\right)`$ (79)
$``$ $`2\delta m_{\stackrel{~}{\nu }_{\mu \tau }}^2M_{\stackrel{~}{W}_3^{}}(M_{L_\tau }^2M_{L_\mu }^2)^1(M_{L_\tau }^2M_{\stackrel{~}{W}_3^{}}^2)^1(M_{L_\mu }^2M_{\stackrel{~}{W}_3^{}}^2)^1`$
$`\times `$ $`\left[M_{L_\mu }^2M_{\stackrel{~}{W}_3^{}}^2\mathrm{ln}\left({\displaystyle \frac{M_{L_\mu }^2}{M_{\stackrel{~}{W}_3^{}}^2}}\right)+M_{L_\tau }^2M_{\stackrel{~}{W}_3^{}}^2\mathrm{ln}\left({\displaystyle \frac{M_{\stackrel{~}{W}_3^{}}^2}{M_{L_\tau }^2}}\right)+M_{L_\tau }^2M_{L_\mu }^2\mathrm{ln}\left({\displaystyle \frac{M_{L_\tau }^2}{M_{L_\mu }^2}}\right)\right],`$
and
$$A=\frac{1}{64\pi ^2}\frac{e^2}{\mathrm{sin}^2\theta _W}.$$
(80)
This matrix can lead easily to a maximal mixing between the $`\mu `$ and $`\tau `$ neutrinos. This will be the case in particular if $`\mathrm{\Delta }_{\mu \mu }\mathrm{\Delta }_{\tau \tau }`$ which implies
$$\frac{m_\tau N_\tau ^e}{M_{L_\tau }^2}\frac{m_\mu N_\mu ^e}{M_{L_\mu }^2}.$$
(81)
In addition it can be seen easily that in the limit where $`M_{L_\mu }^2=M_{L_\tau }^2`$, the determinant of the neutrino mass matrix vanishes, leading to a large hierarchy of masses (as required by atmospheric and solar neutrino experiments, taking into account the fact that the mass of $`\nu _e`$ is below $`10^8`$ eV in the present scenario). For example with the parameters of Eq. (47) and taking in addition $`N_\mu ^e=14N_\tau ^e`$, $`M_{L_\mu }M_{\stackrel{~}{B}^{}}=6`$ TeV, and $`M_{L_\tau }7.5`$ TeV, we obtain one neutrino with a mass $`10^3`$ eV and one with a mass $`10^5`$ eV in addition to the electron neutrino with a mass below $`10^8`$ eV. In this case the mixing between the $`\mu `$ and $`\tau `$ flavors is large ($`\mathrm{sin}2\alpha =0.99`$) while that of the electron flavor with the two other flavors is very much suppressed. Note that the values of the lepton-number violating mass terms $`\delta m_{\stackrel{~}{\nu }_{ij}}^2`$ induced by Fig. 5 are several orders of magnitude below the phenomenological bounds $`\delta m_e<350`$ MeV, $`\delta m_\mu <50`$ GeV, and $`\delta m_\tau <450`$ GeV obtained for $`M_{SUSY}1`$ TeV in Ref. .
In summary, from the above qualitative estimate, we observe that realistic neutrino masses could be accomodated easily in the present scenario, in agreement with atmospheric and solar neutrino experiments. A large mixing and a hierarchy of neutrino masses appear rather naturally. A more quantitative estimate would require an explicit calculation of the two-loop integrals involved, but since there are still many free parameters, it will not add much to our understanding in any case.
## 8 Conclusion
We have studied a model of leptogenesis in a R-parity violating supersymmetric model. The lightest neutralino $`\stackrel{~}{W}_3^{}`$ is assumed to be mostly the $`SU(2)`$ gaugino but its decay into $`l^\pm h^{}`$ is suppressed because the required $`\stackrel{~}{l}_Lh^+`$ mixing is negligible. On the other hand, $`\stackrel{~}{W}_3^{}`$ has a small component of $`\stackrel{~}{B}`$, the $`U(1)`$ gaugino, which decays readily because the required $`\stackrel{~}{l}_Rh^+`$ mixing is of order $`10^3`$ from the presence of nonholomorphic R-parity violating soft terms in the Lagrangian. The decay asymmetry of $`\stackrel{~}{W}_3^{}`$ is then evolved into a lepton asymmetry of the Universe by solving the Boltzmann equations in detail numerically. We demonstrate how each term in the equations affects the eventual outcome of the proposed scenario. The charged scalar mass matrix and the neutralino sector are discussed in detail. A realistic scenario of radiative neutrino mass generation in two loops is presented, which originates from the same lepton-number violating nonholomorphic terms.
Acknowledgements
This work was supported in part by the U.S. Department of Energy under Grant No. DE-FG03-94ER40837. One of us (U.S.) acknowledges the hospitality of the University of California at Riverside where this work was completed.
## Appendix A Complete Charged Scalar Mass Matrix
From Eqs. (7) to (12), neglecting small terms of order $`(\lambda _{ijk})^2`$ and $`\lambda _{ijk}\lambda _{}^{}{}_{lmn}{}^{}`$, the complete charged scalar mass matrix is given by Eq. (50) with (see also Refs. for the holomorphic part):
$`_{h_1^{}h_1^+}^2`$ $`=`$ $`{\displaystyle \frac{g^2}{4}}(v_2^2v_\nu ^2)B{\displaystyle \frac{v_2}{v_1}}\mu \mu _i{\displaystyle \frac{v_{\nu _i}}{v_1}}+{\displaystyle \frac{1}{2}}f_{ij}^ef_{kj}^ev_{\nu _i}v_{\nu _k}N_i^B{\displaystyle \frac{v_{\nu _i}}{v_1}},`$ (82)
$`_{h_1^{}h_2^+}^2`$ $`=`$ $`{\displaystyle \frac{g^2}{4}}v_1v_2B,`$ (83)
$`_{h_1^{}\stackrel{~}{e}_{i_L}^{}{}_{}{}^{}}^2`$ $`=`$ $`{\displaystyle \frac{g^2}{4}}v_1v_{\nu _i}+\mu \mu _i{\displaystyle \frac{f_{kj}^ef_{ij}^e}{2}}v_1v_{\nu _k}+{\displaystyle \frac{1}{2}}f_{lj}^e(\lambda _{ikj}\lambda _{kij})v_{\nu _l}v_{\nu _k}+N_i^B,`$ (84)
$`_{h_1^{}\stackrel{~}{e}_i^c}^2`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}f_{ji}^e\mu _jv_2{\displaystyle \frac{1}{\sqrt{2}}}(A_ef_e)_{ji}v_{\nu _j}+{\displaystyle \frac{1}{\sqrt{2}}}N_i^ev_2,`$ (85)
$`_{h_2^{}h_2^+}^2`$ $`=`$ $`{\displaystyle \frac{g^2}{4}}(v_1^2+v_\nu ^2)B{\displaystyle \frac{v_1}{v_2}}B_i^{}{\displaystyle \frac{v_{\nu _i}}{v_2}},`$ (86)
$`_{h_2^{}\stackrel{~}{e}_{i_L}^{}{}_{}{}^{}}^2`$ $`=`$ $`{\displaystyle \frac{g^2}{4}}v_2v_{\nu _i}B_i^{},`$ (87)
$`_{h_2^{}\stackrel{~}{e}_i^c}^2`$ $`=`$ $`+{\displaystyle \frac{1}{\sqrt{2}}}f_{ji}^e(\mu v_{\nu _j}\mu _jv_1)+{\displaystyle \frac{1}{\sqrt{2}}}\lambda _{kji}(\mu _kv_{\nu _j}\mu _jv_{\nu _k})`$ (88)
$`+{\displaystyle \frac{1}{\sqrt{2}}}N_{ji}^ev_{\nu _j}+{\displaystyle \frac{1}{\sqrt{2}}}N_i^ev_1,`$
$`_{\stackrel{~}{e}_{j_L}\stackrel{~}{e}_{i_L}^{}}^2`$ $`=`$ $`(M_L^2)_{ij}{\displaystyle \frac{1}{8}}(g^2g^2)(v_1^2v_2^2+v_\nu ^2)\delta _{ij}+{\displaystyle \frac{g^2}{4}}v_{\nu _i}v_{\nu _j}+{\displaystyle \frac{1}{2}}f_{jk}^ef_{ik}^ev_1^2+\mu _j\mu _i`$ (89)
$`+{\displaystyle \frac{1}{2}}f_{jl}^ev_1v_{\nu _k}(\lambda _{kil}\lambda _{ikl})+{\displaystyle \frac{1}{2}}f_{il}^ev_1v_{\nu _k}(\lambda _{kjl}\lambda _{jkl}),`$
$`_{\stackrel{~}{e}_{j_L}\stackrel{~}{e}_i^c}^2`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}f_{ji}^e\mu v_2+{\displaystyle \frac{1}{\sqrt{2}}}(A_ef_e)_{ji}v_1+{\displaystyle \frac{1}{\sqrt{2}}}(A_{kji}^eA_{jki}^e)v_{\nu _k}`$ (90)
$`+{\displaystyle \frac{1}{\sqrt{2}}}\mu _k(\lambda _{kji}\lambda _{jki})v_2+{\displaystyle \frac{1}{\sqrt{2}}}N_{ji}^ev_2,`$
$`_{\stackrel{~}{e}_j^c\stackrel{~}{e}_i^c}^2`$ $`=`$ $`(M_{\stackrel{~}{e}^c}^2)_{ji}{\displaystyle \frac{g^2}{4}}(v_1^2v_2^2+v_\nu ^2)\delta _{ij}+{\displaystyle \frac{1}{2}}f_{ki}^ef_{lj}^ev_{\nu _k}v_{\nu _l}+{\displaystyle \frac{1}{2}}f_{ki}^ef_{kj}^ev_1^2`$ (91)
$`+{\displaystyle \frac{1}{2}}f_{kj}^e(\lambda _{lki}\lambda _{kli})v_1v_{\nu _l}+{\displaystyle \frac{1}{2}}f_{ki}^e(\lambda _{lkj}\lambda _{klj})v_1v_{\nu _l}.`$
In Eqs. (82) and (86), the tadpoles conditions for the $`h_1^0`$ and $`h_2^0`$ fields have been used:
$`\left(m_{H_1}^2+\mu ^2+{\displaystyle \frac{1}{8}}(g^2+g_{}^{}{}_{}{}^{2})(v_1^2v_2^2+v_\nu ^2)\right)v_1+(\mu \mu _i+N_i^B)v_{\nu _i}=\mathrm{\hspace{0.17em}\hspace{0.17em}0},`$ (92)
$`\left(m_{H_2}^2+\mu ^2+\mu _i^2{\displaystyle \frac{1}{8}}(g^2+g_{}^{}{}_{}{}^{2})(v_1^2v_2^2+v_\nu ^2)\right)v_2+Bv_1+B_i^{}v_{\nu _i}=\mathrm{\hspace{0.17em}\hspace{0.17em}0}.`$ (93)
In Eqs. (82) to (91), $`v_{\nu _i}`$ are given by the corresponding tadpole conditions for the $`\stackrel{~}{\nu _i}`$ fields:
$$\left((M_L^2)_{ji}+\mu _i\mu _j+\frac{1}{8}(g^2+g_{}^{}{}_{}{}^{2})(v_1^2v_2^2+v_\nu ^2)\delta _{ij}\right)v_{\nu _j}+B_i^{}v_2+N_i^Bv_1+\mu \mu _iv_1=\mathrm{\hspace{0.17em}\hspace{0.17em}0}.$$
(94)
Without nonholomorphic terms, it is difficult to obtain a large $`\stackrel{~}{e}_Rh^+`$ mixing without generating a large $`\stackrel{~}{e}_Lh^+`$ mixing as well (which would induce an undesirably large $`\stackrel{~}{W}_3^{}h^+e_L`$ decay rate) or without requiring very fine tuning between the values of $`v_{\nu _i}`$ and $`\mu _i`$. In the case of one family (putting all indices equal), this can be seen easily. First, the $`_{h_2^{}\stackrel{~}{e}^c}^2`$ matrix element is proportional to the neutrino mass and hence very small. Second, the $`_{h_1^{}\stackrel{~}{e}^c}^2`$ matrix element, neglecting a small term proportional to the neutrino mass, is proportional to the $`_{\stackrel{~}{e}_L\stackrel{~}{e}^c}^2`$ matrix element. Hence it can be shown easily that it is not possible to have a sufficiently large $`_{h_1^{}\stackrel{~}{e}^c}^2`$ matrix element (inducing $`\xi `$ of order $`10^3`$) together with a sufficiently small $`\stackrel{~}{e}_Lh^{}`$ mixing . The latter mixing gets a contribution $`\xi (\mu /\mu _l)[_{\stackrel{~}{e}_L\stackrel{~}{e}^c}^2/max(_{\stackrel{~}{e}_L\stackrel{~}{e}_L^{}}^2,_{\stackrel{~}{e}^c\stackrel{~}{e}^c}^2)]`$. Now, in the case of three families, due to the $`A^e`$ terms in $`_{\stackrel{~}{e}_L\stackrel{~}{e}^c}^2`$, both matrix elements are not any more proportional and the $`\stackrel{~}{e}_L^ih^+`$ mixings can be made as small as necessary independently of the value of $`\xi `$. However, for values of $`\stackrel{~}{e}_L`$ and $`\stackrel{~}{e}^c`$ masses of the order 4 TeV or more (see section 4), a value of $`\xi `$ around $`10^3`$ requires that the $`_{h_1^{}\stackrel{~}{e}^c}^2`$ matrix element is of order $`10^3\times (`$4 TeV$`)^2`$ (125 GeV)<sup>2</sup> which implies very large values of $`\mu _l`$ and $`v_{\nu _l}`$. Hence extreme fine tuning between the values of $`\mu _i`$ and $`v_{\nu _i}`$ is needed to obtain a small enough neutrino mass in Eq. (18).
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# Superradiant light scattering and grating formation in cold atomic vapours
## I Introduction
Recent experiments by S. Inouye et al. at MIT have demonstrated the formation of atomic matter waves in a cigar-shaped Bose-Einstein Condensate (BEC) pumped by an off-resonant laser beam, together with highly directional scattering of light along the major axis of the condensate. This emission has been interpreted as superradiant Rayleigh scattering, and some theoretical work describing this experiment has been recently published . In particular, the work of Moore and Meystre describes the Rayleigh scattering in a BEC using a model which extends the Collective Atomic Recoil Laser (CARL) model originally proposed by Bonifacio et al. to include a quantum-mechanical description of the centre-of-mass motion of the atoms in the condensate . The conclusions of ref. were that the original CARL theory, which treats the atomic centre-of-mass motion classically, fails when the temperature of the atomic sample is below the recoil temperature $`T_R=\mathrm{}\omega _r/k_B`$, where $`\omega _r=\mathrm{}|\stackrel{}{q}|^2/2m`$ is the recoil frequency, $`m`$ is the atomic mass, $`\stackrel{}{q}=\stackrel{}{k}\stackrel{}{k}_s`$ is the difference between the pump and scattered wavevectors and $`k_B`$ is Boltzmann’s constant. However, the cubic dispersion relation derived in ref. reduces to that of the original semiclassical CARL model for large atomic densities: More specifically, the quantum corrections to the classical motion are negligibly small when the CARL parameter $`\rho `$ in the free electron laser (FEL) limit , roughly interpreted as the average number of photons scattered per atom, is greater than one. This suggests that a fully quantum-mechanical description of the atomic centre-of-mass motion may not be necessary in order to describe the main experimental results of ref., i.e. the temporal evolution of the scattered light intensity and the spatial grating in the condensate. We are aware that a semiclassical theory is necessarily limited in its description of the radiation statistics and the quantum degenerate nature of the condensate, which require a full quantum analysis. Nevertheless, we consider the semiclassical approach useful in order to give an intuitive description of the physical mechanism underlying the observed effects. We stress however that in spite of its simplicity, the semiclassical model produces good quantitative agreement with the experimental results of ref..
## II Model
The model described is bidimensional and semiclassical. We represent the cigar-shaped atomic sample as an ellipsoid with length $`L`$ and diameter $`W`$, where $`LW`$ as shown in Fig.1. The sample is exposed to a classical plane wave radiation electric field $`\stackrel{}{E}_0(y,t)=\widehat{x}_0e^{ik(yct)}+c.c.`$, polarised along the $`\widehat{x}`$ axis and incident along the axis $`\widehat{y}`$, with $`_0`$ real and constant and where $`k=\omega /c`$. We assume that the scattered radiation consists of two radiation pulses propagating along the $`\widehat{z}`$ axis, with electric fields polarised as the incident field:
$$\stackrel{}{E}(z,t)=\widehat{x}[_1(z,t)e^{ik(zct)}+_2(z,t)e^{ik(z+ct)}+c.c],$$
(1)
where $`_{1,2}(z,t)`$ are slowly varying complex amplitudes. The dominance of the scattering along the $`\pm \widehat{z}`$ axis over that in other directions is due to the geometry of the atomic sample.
The atomic sample is described as a collisionless gas of atoms, each with two internal energy levels. The internal evolution of each atom is described by the density matrix elements $`\rho _{mn}`$ ($`m,n=1,2`$) for the lower, (1), and upper, (2), levels. The off-diagonal elements $`\rho _{12}=\rho _{21}^{}`$ describe the dipole moment induced by the radiation fields via the relation $`\stackrel{}{d}=\widehat{x}\mu (\rho _{12}+c.c)`$, where $`\mu `$ is the dipole matrix element. The diagonal elements $`\rho _{11}`$ and $`\rho _{22}`$ describe the probability of an atom being in the lower or in the upper level, respectively. The off-diagonal elements may be described conveniently as a sum of three polarisation waves:
$$\rho _{12}=S_0e^{ik(yct)}+S_1e^{ik(zct)}+S_2e^{ik(z+ct)}.$$
(2)
The dipole moment of each atom contributes to the macroscopic polarisation of the atomic sample described by $`\stackrel{}{P}=n(\stackrel{}{x})\stackrel{}{d}`$, where $`n(\stackrel{}{x})`$ is the atomic density. This polarisation is a source for the radiation field via Maxwell’s wave equation which yields, in the usual Slowly Varying Envelope Approximation (SVEA),
$$\left(\frac{_1}{t}+c\frac{_1}{z}\right)e^{ikz}+\left(\frac{_2}{t}c\frac{_2}{z}\right)e^{ikz}=\frac{i\omega \mu }{2ϵ_0}n(\stackrel{}{x})\left\{S_0e^{iky}+S_1e^{ikz}+S_2e^{ikz}\right\},$$
(3)
where we have neglected the terms proportional to $`e^{\pm 2i\omega t}`$. We assume that the atomic sample can be described as a collection of $`N`$ point particles with positions $`\stackrel{}{x}_j`$, so that $`n(\stackrel{}{x})=_{j=1}^N\delta ^{(3)}(\stackrel{}{x}\stackrel{}{x}_j)`$. Multiplying both sides by $`e^{ikz}`$ and integrating over the $`\widehat{z}`$ axis from $`z\mathrm{\Delta }z/2`$ to $`z+\mathrm{\Delta }z/2`$, where $`\mathrm{\Delta }z=\lambda /2`$, Eq.(3) yields
$$\left(\frac{_{1,2}}{t}\pm c\frac{_{1,2}}{z}\right)\mathrm{\Delta }z=\frac{i\omega \mu }{2ϵ_0}\left\{S_0e^{ik(yz_j)}+S_{1,2}+S_{2,1}e^{2ikz_j}\right\}\delta (xx_j)\delta (yy_j),$$
(4)
where the upper sign corresponds to the first subscript and we have assumed the field amplitudes $`_{1,2}`$ are spatially slowly varying over $`\mathrm{\Delta }z`$. Assuming also that $`_{1,2}`$ are independent of $`x`$ and $`y`$, we can integrate on the plane $`(x,y)`$ over the section $`A=\pi W^2/4`$ of the condensate, so that Eq.(4) becomes
$$\left(\frac{_{1,2}}{t}\pm c\frac{_{1,2}}{z}\right)=\frac{i\omega \mu \overline{n}}{2ϵ_0}S_0e^{ik(yz)}+S_{1,2}+S_{2,1}e^{2ikz},$$
(5)
where $`\overline{n}=N/A\mathrm{\Delta }z`$ is the average density and $`..=(1/N)_{j=1}^N(..)_j`$.
In this model the atomic centre-of mass motion is treated classically, with each atom described as a point particle with a given position and momentum. The radiation fields drive the centre-of-mass motion of the atoms via the force
$`\stackrel{}{F}=(0,\stackrel{}{d}{\displaystyle \frac{(\stackrel{}{E}_0+\stackrel{}{E})}{y}},\stackrel{}{d}{\displaystyle \frac{(\stackrel{}{E}_0+\stackrel{}{E})}{z}}).`$
Neglecting the fast-varying temporal terms, the equations for the atomic velocity components are:
$`m{\displaystyle \frac{dv_y}{dt}}`$ $`=`$ $`ik\mu _0[S_0^{}+S_1^{}e^{ik(yz)}+S_2^{}e^{ik(y+z)}c.c]`$ (6)
$`m{\displaystyle \frac{dv_z}{dt}}`$ $`=`$ $`ik\mu \{S_1^{}_1S_2^{}_2+S_2^{}_1e^{2ikz}S_1^{}_2e^{2ikz}+S_0^{}[_1e^{ik(zy)}_2e^{ik(y+z)}]c.c.\}.`$ (7)
We assume that the detuning $`\delta =\omega \omega _a`$ between the optical fields and the atomic resonance is much larger than the natural linewidth of the atomic transition, $`\gamma `$, so that the atoms always remain in their lower internal energy states ($`\rho _{11}1`$ and $`\rho _{22}0`$). Moreover, assuming that the scattering time scale is much longer than the relaxation time $`\gamma ^1`$, we can adiabatically eliminate the atomic polarisations, i.e. $`S_k=i(\mu /\mathrm{})_k/(\gamma +i\delta )\mathrm{\Omega }_k/2\delta `$, where $`\mathrm{\Omega }_k=2\mu _k/\mathrm{}`$, $`|\mathrm{\Omega }_k|`$ is the Rabi frequency for the field $`k`$ and $`k=0,1,2`$. With these approximations, Eqs.(6) and (7) yield:
$`m{\displaystyle \frac{dv_y}{dt}}`$ $`=`$ $`i\mathrm{}k(\mathrm{\Omega }_0/4\delta )[\mathrm{\Omega }_1e^{ik(zy)}\mathrm{\Omega }_2^{}e^{ik(z+y)}c.c]`$ (8)
$`m{\displaystyle \frac{dv_z}{dt}}`$ $`=`$ $`i\mathrm{}k(\mathrm{\Omega }_0/4\delta )[\mathrm{\Omega }_1e^{ik(zy)}+\mathrm{\Omega }_2^{}e^{ik(z+y)}c.c]+i(\mathrm{}k/2\delta )[\mathrm{\Omega }_1\mathrm{\Omega }_2^{}e^{2ikz}c.c].`$ (9)
It is seen that the interference between the pump and the scattered fields forms two bidimensional periodic potentials $`V_{1,2}(y,z)|_0_{1,2}|\mathrm{cos}[k(zy)\pm \varphi _{1,2}]`$ in the plane $`(\widehat{y},\widehat{z})`$, where $`\varphi _{1,2}`$ are the phases of the complex amplitudes $`_{1,2}`$. A weaker 1D potential $`V_3(z)|_1_2|\mathrm{cos}[2kz+\varphi _1\varphi _2]`$ forms along the $`\widehat{z}`$ axis due to the interference of the two counterpropagating scattered fields. If the pump intensity is large enough, we can assume $`_0_{1,2}`$ and neglect the ponderomotive potential $`V_3`$. Then, Eqs.(8),(9) and (5) can be conveniently written in the following dimensionless form :
$`{\displaystyle \frac{d\theta _{1,2}}{d\overline{t}}}`$ $`=`$ $`p_{1,2},`$ (10)
$`{\displaystyle \frac{dp_{1,2}}{d\overline{t}}}`$ $`=`$ $`[A_{1,2}e^{\pm i\theta _{1,2}}+c.c],`$ (11)
$`{\displaystyle \frac{A_{1,2}}{\overline{t}}}`$ $`\pm `$ $`{\displaystyle \frac{A_{1,2}}{\overline{z}}}=e^{i\theta _{1,2}}`$ (12)
where $`\theta _{1,2}=k(zy)`$, $`p_{1,2}=(m/\mathrm{}k\rho )(v_zv_y)`$ and $`A_{1,2}=2i(ϵ_0/\mathrm{}\omega \overline{n}\rho )^{1/2}_{1,2}`$ are scaled atomic position, atomic momentum and field amplitude variables respectively. The dimensionless time and space coordinates, $`\overline{t}=\omega _r\rho t`$ and $`\overline{z}=\omega _r\rho z/c`$, are scaled in terms of the collective recoil bandwidth, $`\rho \omega _r`$, where $`\omega _r=\mathrm{}k^2/m`$ is the single-atom recoil frequency and $`\rho =(\mathrm{\Omega }_0/2\delta )^{2/3}(\omega \mu ^2\overline{n}/ϵ_0\omega _r^2\mathrm{})^{1/3}`$ is the dimensionless CARL parameter . At $`\overline{t}=0`$, the atoms are assumed to be randomly distributed in position and have zero momentum, and the amplitudes of the scattered fields are set to zero.
## III Analysis
In this simple model the two scattered fields are uncoupled and symmetric. For each field (1,2) individually, Eqs.(10)-(12) are formally identical to those which describe pulse propagation in a high gain free electron laser (FEL) . It is already known that they admit a self-similar solution of the form $`A_{1,2}(\overline{z},\overline{t})=\pm \overline{z}𝒜(u)`$, where $`u=\sqrt{|\overline{z}|}(\overline{t}\overline{z})`$ and $`𝒜(u)`$ is the solution of a set of ordinary differential equations . This self-similar solution describes the superradiant emission of radiation pulses whose duration decreases in proportion to the fourth root of the peak intensity. The pulse shape can be approximated by a hyperbolic secant function, followed by some non-linear ‘ringing’, similar to that which occurs in superfluorescence from inverted two-level atoms .
A simpler model can be obtained by approximating the spatial derivative in the field equation (12) by a damping term i.e.
$$\frac{dA_{1,2}}{d\overline{t}}=e^{i\theta _{1,2}}\kappa A_{1,2}$$
(13)
where $`\kappa =c/2\omega _r\rho L`$ and $`L/c`$ is the transit time of the photon along the major axis of the condensate. In this approximation, the finite interaction time due to the escape of radiation from the atomic sample is represented by an incoherent decay of the field amplitude in the sample at a rate $`c/2L`$, half the inverse of the radiation ‘lifetime’ in the atomic sample. A more general treatment where the radiation is scattered in a direction making an angle $`\psi `$ with respect to the $`\widehat{z}`$ axis should give $`\kappa (c/2\omega _r\rho )[|\mathrm{sin}\psi |/W+|\mathrm{cos}\psi |/L]`$ . As $`LW`$, the radiation is least strongly damped along the major axis of the sample.
An approximate solution to eqs.(10),(11) and (13) can be found assuming $`\kappa 1`$ and adiabatically eliminating the field variables, i.e. $`A_{1,2}\kappa ^1\mathrm{exp}[i\theta _{1,2}]`$. In this limit, the rate of change of the average scaled momentum is $`(d/d\overline{t})p_{1,2}=2\kappa |A_{1,2}|^2`$. A third-order analysis of the equations in the mean-field limit (i.e. with radiation propagation modelled by the damping term) gives the following approximate solution:
$$|A_{1,2}|^2\frac{1}{2\kappa ^2}\mathrm{sech}^2\left[(\overline{t}\overline{t}_D)/\sqrt{2\kappa }\right]$$
(14)
and
$$p_{1,2}\sqrt{\frac{2}{\kappa }}\left\{1+\mathrm{tanh}\left[(\overline{t}\overline{t}_D)/\sqrt{2\kappa }\right]\right\},$$
(15)
where $`\overline{t}_D=\sqrt{2\kappa }\mathrm{ln}(|b_0|/\sqrt{2})`$ is the delay time of the peak and $`b_0=\mathrm{exp}[i\theta _{1,2}(\overline{t}=0)]`$ is the initial bunching, which can be assumed to be $`1/\sqrt{N}`$ for a condensate of $`N`$ atoms. In the linear regime the exponential gain is
$$G=\omega _r\rho \sqrt{2/\kappa }=(3\gamma /\delta )\sqrt{(2I_0N/m\omega )(\lambda ^2/A)},$$
(16)
whereas the peak value of the scattered intensity is
$$I_{peak}=(\gamma /\delta )^2[(3/2\pi )(\lambda ^2/A)N]^2I_0,$$
(17)
where $`I_0=2cϵ_0|_0|^2`$ is the pump intensity, $`A`$ is the cross sectional area of the condensate and $`\gamma =\mu ^2k^3/6\pi \mathrm{}ϵ_0`$ is the natural decay rate of the atomic transition.
## IV Comparison with the MIT experiment
In the MIT experiment, a sodium BEC was exposed to a single off-resonant laser pulse red-detuned by $`\delta /2\pi =1.7`$ GHz from the $`3S_{1/2}3P_{3/2}`$ transition, with $`\lambda =0.589\mu \text{m}`$ and natural width $`\gamma =0.31\times 10^8`$ s<sup>-1</sup>. The recoil frequency is $`\omega _r=3\times 10^5`$ s<sup>-1</sup>. We assume that the condensate had a diameter of $`20\mu \text{ m}`$ and a length of $`200\mu \text{m}`$, approximately $`N=5\times 10^5`$ atoms participate in the emission of a scattered radiation pulse. The dimensionless parameters are $`\rho 44\times I_0^{1/3}`$ and $`\kappa 5.5\times 10^4\times I_0^{1/3}`$, where $`I_0`$ is the pump intensity in mW/cm<sup>2</sup>. As $`I_0>1`$ and consequently $`\rho 1`$, the results of ref. indicate that quantum effects due to atomic diffraction should be negligibly small for this experiment, even though $`TT_R`$. As $`\kappa 1`$, the atoms emit two superradiant pulses along the major axis of the condensate. The gain is approximately $`G82\times \sqrt{I_0}`$, where $`G`$ is given in ms<sup>-1</sup>, whereas the peak occurs after a time $`t_D=\mathrm{ln}(2N)/G(170/\sqrt{I_0})\mu \text{s}`$, in good agreement with the measured values of ref.. Furthermore, from Eq.(15) the modulus of the average atomic velocity is $`v=(\mathrm{}k/m)\rho |p_{1,2}|(\lambda /2\pi )G0.7\sqrt{I_0}\text{cm s}^1`$. Fig. 2 shows the temporal evolution of the main peak of the scattered intensity, as given by the approximate formula (14), for the parameters of the MIT experiment and three different values of the incident intensity, 3.8 mW/cm<sup>2</sup> (solid line), 2.4 mW/cm<sup>2</sup> (dashed line) and 1.4 mW/cm<sup>2</sup> (dotted line).
In addition to the temporal evolution of the scattered radiation pulses, there are other predictions of this semiclassical model which are consistent with the results of the MIT experiment: Firstly, superradiance is observable only if the Doppler broadening of the atomic resonance is sufficiently small that $`\sigma t_{SR}1`$, where $`\sigma `$ is the rms spread of the gaussian spectral distribution and $`t_{SR}1/G`$ is the superradiant time . The observed spectral width of the Bragg resonance of the BEC of approximately $`5`$ kHz (corresponding to a velocity spread of few mm s$`{}_{}{}^{}1`$) yields $`\sigma t_{SR}0.16\times I_0^{1/2}`$. We observe that, using $`\sigma =k(k_BT/m)^{1/2}`$, a temperature of only $`1\mu K`$ (approximately the BEC transition temperature for the MIT experiment) would increase the frequency spread by a factor of $`15`$, enough to destroy the superradiant emission. This explains why superradiant emission was observed only at the extremely low temperatures below the threshold for Bose-Einstein condensation . Secondly, superradiant emission parallel and antiparallel to the $`\widehat{z}`$ axis induces an average atomic velocity $`\stackrel{}{v}_{1,2}(G/k)[\widehat{y}\widehat{z}]`$, respectively at 45<sup>0</sup> degree with respect to the negative (positive) direction of the $`\widehat{z}`$ axis, as observed in the MIT experiment. We assume the existence of two distinct families of atoms interacting with the two independent superradiant pulses $`A_1`$ and $`A_2`$. However different orders of atomic velocity, i.e. $`\stackrel{}{v}_{m,n}=m\stackrel{}{v}_1+n\stackrel{}{v}_2`$, with $`m,n`$ integers, have also been observed. More precisely, the orders $`(2,0)`$, $`(1,1)`$, $`(0,2)`$, $`(2,1)`$ and $`(1,2)`$ other than the usual $`(1,0)`$ and $`(0,1)`$, have been clearly observed in the experiment after increasing the exposure time to the laser source and letting the atomic cloud expand ballistically after the interaction . The formation of this momentum grating can be explained as a sequential superradiant scattering process in which the atoms emit $`m`$ pulses along the positive $`\widehat{z}`$ and $`n`$ pulses along the negative $`\widehat{z}`$ axis, acquiring a total recoil velocity $`\stackrel{}{v}_{m,n}`$. The extremely narrow resonance line allows the atoms to emit up to three sequential superradiantly scattered pulses before $`\sigma t1`$, which is consistent with the observation of the atomic momentum distribution.
## V Conclusions
In conclusion, we have presented a semiclassical model describing the superradiant Rayleigh scattering from a Bose-Einstein condensate observed in Ref. . The model is much simpler than those previously used to explain the results of as the atomic centre-of-mass motion is treated classically. The evolution of the scattered intensity and the atomic motion due to recoil as calculated from this simple model are in quantitative agreement with the experimental results. The fact that quantum centre-of-mass effects such as atomic diffraction are negligible is a consequence of the high density of the condensate. In our model the BEC is essentially described as a collisionless Doppler-free atomic gas. The results presented here suggest that together with its high density, the most important property of the condensate with regard to superradiant light scattering is its very low temperature rather than its quantum degenerate nature. In this respect the situation is similar to that of ultraslow propagation of light in a BEC . Subsequent observations of ultraslow propagation in a hot vapour demonstrated that the significant property of the BEC was that it was a Doppler-free optical medium rather than a quantum degenerate one.
## ACKNOWLEDGMENTS
The authors would like to thank the Royal Society of Edinburgh for support of G.R.M.R. and the EPSRC for support of B.M<sup>c</sup>N.
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# Phase space localization of chaotic eigenstates: Violating ergodicity
## I Introduction
For a bounded, classically chaotic system, ergodicity is defined with respect to the energy surface, the only available invariant space of finite measure. In an extension developed just over twenty years ago, the consequences of ergodicity for the eigenstates of a corresponding quantum system were conjectured to give rise to a locally, Gaussian random behavior . Shortly thereafter, work ensued on defining the concept of eigenstate ergodicity within a more rigorous framework . Some of the paradoxes and peculiarities have been recently explored as well . One expression of eigenstate ergodicity is that a typical eigenstate would fluctuate over the energy surface, but otherwise be featureless, in an appropriate pseudo-phase-space representation such as the Wigner transform representation . Any statistically significant deviation from ergodicity in individual eigenstates is termed phase space localization.
It came as a surprise when Heller discovered eigenstates were “scarred” by short, unstable periodic orbits . A great deal of theoretical and numerical research followed , and experiments also . In fact, scarring is just one of the means by which phase space localization can exist in the eigenstates of such systems. Another means would be localizing effects due to transport barriers such as cantori or broken separatrices . Despite these studies and the semiclassical construction of an eigenstate , the properties of individual eigenstates remain somewhat a mystery.
Individual eigenfunctions may not be physically very relevant in many situations, especially those involving a high density of states. In this case, groups of states contribute towards localization in ways that may be understood with available semiclassical theories. One simple and important quantity where this could arise is the time average of an observable as this is a weighted sum of several (in principle all) eigenfunctions. In the Heisenberg picture, where the operator is evolving in time, the expectation value of the observable could be measured with any state. Phase space localization features would be especially evident if this state were chosen to be a wave packet well localized in such spaces.
One of us has already studied the correlation between level velocities and wavefunction intensities in connection with localization . This can be directly connected to the issues raised above, and our treatment thus extends the previous work. The present paper is a companion to a study of similar problems in continuous Hamiltonian systems (as opposed to Hamiltonian maps, i.e. discrete time systems) and billiards . The methods used in these two papers complement each other and the results in the present paper are detailed as the systems studied are much simpler.
In a system with a non-degenerate spectrum the time average of an observable $`\widehat{A}`$ in state $`|\alpha `$ is
$$\widehat{A}(t)_t=\underset{n}{}|\alpha |\psi _n|^2\psi _n|\widehat{A}|\psi _n,$$
(1)
where $`|\psi _n`$ are the eigenstates of the Hamiltonian $`\widehat{H}`$. If the system depends upon a parameter $`\lambda `$ which varies continuously, an energy level “velocity”, $`E_n/\lambda `$, for the $`n^{th}`$ level can be defined (velocity is a bit of a misnomer for it is actually just a slope - we are not evolving the system parameter in time). The level velocity-intensity correlation is identical to the time average above if we identify $`\widehat{H}/\lambda `$ with $`\widehat{A}`$, as from the so called Hellmann-Feynman theorem (for instance ):
$$\psi _n|\frac{\widehat{H}}{\lambda }|\psi _n=\frac{E_n}{\lambda }.$$
(2)
It must be noted that while we call the operator average Eq. (1) a “correlation” it is not the true correlation that is obtained by dividing out the rms values of the wavefunction intensities and the operator expectation value (as defined in ). In other words we are going to study the covariance rather than the correlation. This is followed in this paper for two reasons; first, dividing out these quantities does not retain the meaning of the time average of an observable and second, the root mean square of the wavefunction intensities which is essentially the inverse participation ratio in phase space is itself a fairly complex quantity reflecting on phase-space localization.
## II Correlation for the quantum bakers map
### A Generalities
We first formulate in general the quantity to be studied and the approach to its semiclassical evaluation. The classical dynamical systems are discrete maps on the dimensionless unit two-torus whose cyclical coordinates are denoted $`(q,p)`$. The quantum kinematics is set in a space of dimension $`N`$ where this is related to the scaled Planck constant as $`N=1/h`$, and the classical limit is the large $`N`$ limit. The quantum dynamics is specified by a unitary operator $`U`$ (quantum map) that propagates states by one discrete time step. The quantum stationary states are the eigensolutions of this operator. The $`N`$ eigenfunctions and eigenangles are denoted by $`\{|\psi _i,\varphi _i;i=0,\mathrm{},N1\}`$. The eigenvalues lie on the unit circle and are members of the set $`\{\mathrm{exp}(i\varphi _i)i=0,\mathrm{},N1\}`$.
The central object of interest is:
$$C_A(\alpha )=\underset{i=0}{\overset{N1}{}}|\alpha |\psi _i|^2\psi _i|\widehat{A}|\psi _i.$$
(3)
The operator $`\widehat{A}`$ is a Hermitian operator and either describes an observable or a perturbation allowing one to follow the levels’ motions continuously. The state $`|\alpha `$ represents a wave packet on the two-torus that is well localized in the $`(q,p)`$ coordinates . We will be interested in the quantum effects over and above the classical limit and we will require that the operator is traceless. Otherwise we will need to subtract the uncorrelated product of the averages of the eigenfunctions (unity) and the trace of the operator. This immediately also implies that the correlation according to Random Matrix Theory (RMT) is zero as well. The ensemble average of $`C_A(\alpha )`$ will wash out random oscillations that are a characteristic of the Gaussian distributed eigenfunctions of the random matrices. Specific localization properties that we will discuss are then not part of the RMT models of quantized chaotic systems. In the framework of level velocities we are considering the situation where the average level velocity is zero, i.e., there is no net drift of the levels.
As we noted in the introduction the correlation $`C_A(\alpha )`$ is simply the time average of the observable:
$$C_A(\alpha )=\alpha |U^n\widehat{A}U^n|\alpha _n$$
(4)
where the large angular brackets denotes the time $`(n)`$ average. This requires that degeneracies do not exist and we assume that this is the case as we are primarily interested in quantized chaotic systems. A physically less transparent identity that is nevertheless useful in subsequent evaluations is:
$$C_A(\alpha )=\alpha |U^n|\alpha \text{Tr}\left(\widehat{A}U^n\right)_n$$
(5)
This may be written more symmetrically as
$$C_A(\alpha )=\text{Tr}\left(|\alpha \alpha |U^n\right)\text{Tr}\left(\widehat{A}U^n\right)_n.$$
(6)
Thus, the correlation is a sort of time evolved average correlation between the two operators $`\widehat{A}`$ and $`|\alpha \alpha |`$. The semiclassical expressions for these are however different as complications arise from the classical limit of $`|\alpha \alpha |`$ which would be varying over scales of $`\mathrm{}`$ that govern the validity of the stationary phase approximations. However, we may anticipate, based on the last form, that the semiclassical expression would be roughly the correlations of the classical limits of these two operators .
### B Semiclassical Evaluation
The bakers map is a very attractive system to study the quantities discussed in the introduction. The classical dynamics is particularly simple (it is sometimes referred to as the “harmonic oscillator of chaos”). A simple quantization is due to Balazs and Voros (where a discussion of the classical dynamics may also found). As a model of quantum chaos it shows many generic features including the one central to this study namely scarring localization of eigenfunctions . There are detailed semiclassical theories that have been verified substantially . We neglect certain anomalous features of the quantum bakers map that would eventually show up in the classical limit. This is reasonable in the range of scaled Planck constant values we have used in the following.
We use the second time averaged expression, Eq. (5), for the correlation. We do not repeat here details of the quantization of the bakers map or the semiclassical theories of this operator except note that we use the anti-periodic boundary conditions as stipulated by Saraceno in order to retain fully the classical symmetries.
The semiclassical theory of the bakers map deals with the powers of the propagator. The trace of $`U^n`$, the time $`n`$ propagator, has been written in the canonical form of a sum over classical hyperbolic periodic orbits with the phases being actions and the amplitudes relating to the linear stability of the orbits. The complications with Maslov phases is absent here . Also, the semiclassical expressions have been derived for matrix elements of the time $`n`$ propagator in the wave packets basis . The time domain dominates the study of the quantum maps, the Fourier transform to the spectrum being done exactly. Our approach to the correlation is then naturally built in the time domain. The situation is different in the case of Hamiltonian time independent flows where the energy domain is very useful.
We use the semiclassical expression for the propagator diagonal matrix elements derived in :
$`\alpha |U^n|\alpha `$ $``$ $`{\displaystyle \underset{\gamma }{}}{\displaystyle \frac{\mathrm{exp}(iS_\gamma /\mathrm{})}{\sqrt{\mathrm{cosh}(\lambda n)}}}{\displaystyle \underset{j}{}}\mathrm{exp}[{\displaystyle \frac{\mathrm{cosh}(n\lambda )1}{2\mathrm{cosh}(n\lambda )\mathrm{}}}`$ (8)
$`\times (\delta q^2+\delta p^2){\displaystyle \frac{i\delta q\delta p}{\mathrm{}}}\mathrm{tanh}(n\lambda )].`$
Here $`\gamma `$ labels periodic orbits of period $`n`$ including repetitions. The Lyapunov exponent is $`\lambda `$ which is $`\mathrm{ln}(2)`$ for the usual bakers map (corresponding to the $`(1/2,1/2)`$ partition and Bernoulli process). Also $`\mathrm{}=h/(2\pi )=1/(2\pi N)`$, $`\delta q=q_jq_\alpha `$ and a similar relation for $`p`$. The position of $`j`$-th periodic point on the periodic orbit $`\gamma `$ is $`(q_j,p_j)`$. The centroids of the wave packets, assumed circular Gaussians, are $`(q_\alpha ,p_\alpha )`$. The choice of type of wave packets is not crucial for the features we seek. We note that the simplicity of this expression for the propagator derives from the simplicity of the classical bakers map, especially the fact that the stable and unstable manifolds are everywhere aligned with the $`(q,p)`$ axes. That Eq. (8) happens to be a periodic orbit sum differs from the similar treatment for billiards as found in where such sums are treated as homoclinic orbit sums. Note however, that the local linearity of the bakers map renders the two approaches (periodic orbit, homoclinic orbit) equivalent.
A generalization of the trace formula for the propagator is given below that is easily derived by the usual procedure employed for the propagator itself . Such a formula was derived in for the case of Hamiltonian flows in the energy domain. We make the simplifying assumption that the operator $`\widehat{A}`$ is diagonal in the position representation (we could treat the case of $`\widehat{A}`$ being diagonal in momentum alone as well). This avoids the problem of a Weyl-Wigner association of operators to functions on the torus. The quantum operator $`\widehat{A}`$ under this simplifying assumption has an obvious classical limit and associated function which is denoted by $`A(q)`$. The other major assumption used in deriving the formula below is that it does not vary on scales comparable to or smaller than $`\mathrm{}`$.
Thus we derive:
$$\text{Tr}\left(\widehat{A}U^n\right)\underset{\gamma }{}\frac{\mathrm{exp}(iS_\gamma /\mathrm{})}{2\mathrm{sinh}(n\lambda /2)}\underset{j}{}A(q_j).$$
(9)
The index $`j`$ again labels points along the periodic orbit $`\gamma `$. The sum over the periodic orbit is the analogue of the integral of the Weyl transform over a primitive periodic orbit in the Hamiltonian flow case . The special case $`\widehat{A}=I`$ the identity corresponds to the usual trace formula . Note that we have written the sums above as being over periodic orbits, while the trace formulas have been often written as sums over fixed points.
The first step is to multiply the two semiclassical periodic orbit sums in Eq. (8) and in Eq. (9). Since there is a time average, $`n`$ is assumed large enough, but not too large (so that these expansions retain some accuracy). All hyperbolic functions are approximated by their dominant exponential dependences. The diagonal approximation and the uniformity principle is used as well.
$`C_A(\alpha )`$ $`=`$ $`{\displaystyle \underset{\gamma }{}}\sqrt{2}\mathrm{exp}(n\lambda ){\displaystyle \underset{T}{}}\left({\displaystyle \underset{j}{}}F(q_j,p_j)\right)`$ (11)
$`\times \left({\displaystyle \underset{j}{}}A(Tq_j,Tp_j)\right)_n`$
Here we have taken a more general dependence for $`A`$ (including the possibility of momentum dependence). $`T`$ represents elements of the symmetry group of the system including time-reversal symmetry and including, of course, unity. These symmetries imply in general, though not as a rule, distinct (for $`TI`$) orbits with identical actions. One assumes that the overwhelming number of action degeneracies are due to such symmetries.
The function $`F`$ is the approximated Gaussian:
$$F(q_j,p_j)=\mathrm{exp}\left[\frac{1}{2\mathrm{}}(\delta q^2+\delta p^2)\frac{i\delta q\delta p}{\mathrm{}}\right]$$
(12)
Using $`\lambda =\mathrm{ln}(2)`$ and the fact that there are approximately $`2^n/n`$ orbits of period $`n`$, one finds
$$C_A(\alpha )=\sqrt{2}\underset{T}{}\underset{l=M}{\overset{M}{}}\stackrel{~}{C}_T(l)$$
(13)
where $`\stackrel{~}{C}(l)`$ is a classical $`l`$-step correlation:
$$\stackrel{~}{C}_T(l)=\frac{1}{n}\underset{j=1}{\overset{n}{}}F(q_j,p_j)A(Tq_{j+l},Tp_{j+l}).$$
(14)
The time average is taken over a typical orbit. We abandon any specific periodic orbit and appeal to ergodicity, taking $`n`$ and also $`M`$ as practically infinite. This is with the assumption that such correlations will decay with time $`l`$. In fact, below we calculate such correlations explicitly and display the decay. Note that $`\stackrel{~}{C}_T(l)\stackrel{~}{C}_T(l)`$ in general. Although these are classical correlations, in the sense that $`q_j,p_j`$ represent a classical orbit, $`\mathrm{}`$ appears as a parameter in them through $`F`$. Further, using the ergodic principle we can replace time averages in $`\stackrel{~}{C}_T(l)`$ by appropriate phase space averages:
$$\stackrel{~}{C}_T(l)=𝑑q𝑑pF(q,p)A(Tf^l(q,p),Tg^l(q,p))$$
(15)
where we have used the fact that the total phase space volume (area) is unity, and $`f^l(q,p)=q_l`$, $`g^l(q,p)=p_l`$ are the classical $`l`$-step integrated mappings.
### C Special case and verifications
We first consider the case that $`\widehat{A}=A_0(\widehat{T}_p+\widehat{T}_p^{})/2`$, where $`\widehat{T}_p`$ is the unitary single-step momentum translation operator that is diagonal in the position representation and $`A_0`$ is a constant real number. This implies that the associated function is $`A(q)=A_0\mathrm{cos}(2\pi q)`$. Below we consider $`A_0=1`$ as the strength of the perturbation. The elements of $`T`$, apart from the identity ($`I`$), are time-reversal ($`TR`$) symmetry and parity ($`P`$). Time reversal in the bakers map is $`(T(q)=p,T(p)=q)`$ followed by backward iteration, while parity is the transformation $`(T(q)=1q,T(p)=1p)`$.
We begin with the evaluation of the forward correlation $`(l0)`$ corresponding to $`T=I`$.
$$\stackrel{~}{C}_I(l)=_{\mathrm{}}^{\mathrm{}}𝑑q𝑑pF(q,p)\mathrm{cos}(2\pi 2^lq).$$
(16)
This follows from the equality:
$$f^l(q)=2^lq(\text{mod}\mathrm{\hspace{0.17em}1})$$
(17)
for the bakers map. The limits of the integrals can be extended to the entire plane as long as the centroid of the weighting factor $`(q_\alpha ,p_\alpha )`$ is far enough away from the edges of the unit phase space square that the Gaussian tails are small there. The integral is elementary, and using $`h=1/N`$ one gets:
$$\stackrel{~}{C}_I(l0)=\frac{1}{\sqrt{2}N}\mathrm{exp}(2^{2l}\pi /(2N))\mathrm{cos}(2\pi 2^lq_\alpha ).$$
(18)
This explicit expression shows the super-exponential decrease with time $`l`$ in the correlation coefficients. It is interesting to note that the logarithmic time scale which sets an important quantum-classical correspondence scale of divergence for chaotic systems, here $`\tau =1/\lambda \mathrm{ln}(1/2\pi \mathrm{})=\mathrm{ln}(N)/\mathrm{ln}(2)`$, enters the correlation decay. In fact, the correlations are significant to precisely half the log-time. We anticipate this feature to hold in general, including autonomous Hamiltonian systems.
Since $`g^l(p)=\mathrm{\hspace{0.17em}2}^lp(\text{mod }\mathrm{\hspace{0.17em}1})`$, for $`(l0)`$, the time-reversed backward correlation $`(l0)`$ is
$$\stackrel{~}{C}_{TR}(l0)=\frac{1}{\sqrt{2}N}\mathrm{exp}(2^{2l}\pi /(2N))\mathrm{cos}(2\pi 2^lp_\alpha )$$
(19)
which also decays super-exponentially and is responsible for the $`(qp)`$ symmetry in the final correlation.
Next we turn to the other, apparently more curious, correlations: the backward identity correlations and the forward time-reversed one. As an example of a backward $`(l0)`$ identity correlation consider $`l=1`$:
$`f^1(q)=\{\begin{array}{cc}q/2\hfill & \text{for }p<1/2\hfill \\ (q+1)/2\hfill & \text{for }p>1/2\hfill \end{array}`$
Therefore
$`\stackrel{~}{C}_I(1)`$ $`=`$ $`{\displaystyle _0^1}𝑑q{\displaystyle _0^{1/2}}𝑑pF(q,p)\mathrm{cos}(\pi q)`$ (21)
$`{\displaystyle _0^1}𝑑q{\displaystyle _{1/2}^1}𝑑pF(q,p)\mathrm{cos}(\pi q)`$
In fact, since $`\mathrm{cos}(\pi q)`$ vanishes at 1/2, there is no discontinuity in the full integral, but it is more difficult to evaluate (and to approximate). If one were to take the upper limits of the $`p`$ integrals to be infinity, there would be errors at $`p=1/2`$. However, this is not terribly damaging, and tolerating a small discontinuity at this point due to this approximation leads to:
$$\stackrel{~}{C}_I(1)=\pm \frac{1}{\sqrt{2}N}\mathrm{exp}(\pi /(8N))\mathrm{cos}(\pi q_\alpha )$$
(22)
the sign depending on if $`p_\alpha <1/2`$ or if $`p_\alpha >1/2`$ respectively. The time-reversed, forward correlation, $`\stackrel{~}{C}_{TR}(1)`$, is the same as this except for interchanging the roles of $`q_\alpha `$ and $`p_\alpha `$.
The generalization of this to higher times is (take $`l>0`$ below):
$$\stackrel{~}{C}_I(l)=\underset{\nu =0}{\overset{2^l1}{}}_0^1𝑑q_{\nu /2^l}^{(\nu +1)/2^l}𝑑pF(q,p)\mathrm{cos}(2\pi (q+\overline{\nu })/2^l)$$
(23)
where $`\nu `$ represents a partition of the bakers map at time $`l`$, and $`\overline{\nu }`$ results from the bit-reversal of the binary expansion of $`\nu `$. The momentum gets exponentially partitioned with time, and it precludes going beyond the log-time here as well (like the forward correlation), although there is apparently no super-exponential decrease here. Indeed if we evaluate the above after neglecting finite limits in each of the $`p`$ integrals above, so that we would have $`2^l`$ discontinuities at time $`l`$, we get:
$$\stackrel{~}{C}_I(l)=\frac{1}{\sqrt{2}N}\mathrm{exp}(\pi /(2^{(2l+1)}N))\mathrm{cos}(2\pi (q_\alpha +\overline{\nu })/2^l)$$
(24)
depending on if $`p_\alpha `$ lies in the interval $`(\nu /2^l,(\nu +1)/2^l)`$. So that for $`l`$ large and $`N`$ fixed, the exponential goes to unity; effectively, for large $`N`$ and any $`l`$, the exponential can be replaced by unity. Even the $`q_\alpha `$-dependent part of the argument in the function (cos) itself is tending to vanish, so that the integral seems to give the area of the Gaussian ($`h`$). The approximation of putting all $`p`$ limits to infinity makes sense only if the Gaussian state is well within a zone of the partition and this is necessarily violated at half the log-time. So the approximate expression of Eq. (24) breaks down beyond $`\tau /2`$. This lack of a super-exponential cutoff as seen with the previous correlations considered is due to two special conditions. First, the argument of the cosine has no $`p`$-dependence. Second, all the stable manifolds are perfectly parallel to the $`p`$ axis. We would recover super-exponential decay in all the correlations if the operator, $`\widehat{A}`$, being considered was a constant function along neither the stable nor the unstable manifold. In this sense, we have chosen a maximally difficult operator with which to test the semiclassical theory, though it simplifies the quantum calculations.
As before, $`\stackrel{~}{C}_{TR}(l)(q,p)=\stackrel{~}{C}_I(l)(p,q)`$. Parity symmetry is benign and leads to an overall multiplication by a factor of 2. Thus, the final semiclassical expression for the full correlation for the quantum bakers map is:
$`C_A(\alpha )`$ $`=`$ $`{\displaystyle \frac{2}{N}}[{\displaystyle \underset{l=0}{\overset{T_1}{}}}\mathrm{exp}(2^{2l}\pi /(2N))\mathrm{cos}(2\pi 2^lq_\alpha )`$ (28)
$`+{\displaystyle \underset{l=1}{\overset{T_2}{}}}\mathrm{exp}(\pi /(2^{2l+1}N)){\displaystyle \underset{\nu =0}{\overset{2^l1}{}}}(\mathrm{cos}(2\pi (q_\alpha +\overline{\nu })/2^l)`$
$`\times \mathrm{\Theta }(p_\alpha \nu /2^l)\mathrm{\Theta }((\nu +1)/2^lp_\alpha ))]`$
$`+(q_\alpha p_\alpha ).`$
where $`T_1`$ can be infinite but it is sufficient to stop just beyond half the log-time. As just discussed, $`T_2`$, is more problematic here, and we do not have an expression to use beyond $`\tau /2`$. $`\mathrm{\Theta }`$ is the Heavyside step function that is zero if the argument is negative and unity otherwise. The correlation is of the order $`1/N`$ or $`\mathrm{}`$. If one were to divide by the number of states in Eq. (3) so that it is a true average, this quantity would decrease as $`1/N^2`$ or $`\mathrm{}^2`$.
For the case of $`N=100`$, we compare in Fig. (1) the full quantum correlation given by Eq. (3) with the final semiclassical evaluation given by Eq. (28). The absolute value of the correlation function is contoured and superposed on a grey scale. Figure (1a) shows the quantum calculation for the full phase space. In other words, the intensity (value) of each point, $`(q,p)`$, on the plot represents the $`C_A(\alpha )`$-calculation for a wave packet centered at $`(q_\alpha =q,p_\alpha =p)`$. The first sum in Eq. (28) (over $`T_1`$ terms) is a smooth function, and it also displays an additional symmetry about $`1/2`$ in both canonical variables separately. This extra symmetry is broken by the second sum (over $`T_2`$ terms). Figure 1(b) compares the semiclassical formula to the exact quantum calculation. We have taken eight “forward” correlations (excluding zero), i.e. $`T_1=8`$, while we have only taken two “backward” correlations, i.e. $`T_2=2`$. This is because it appears that the approximations that go into the latter expressions lead to non-uniformly converging quantities and it works better to stop at a earlier point in the series. The (artificial) discontinuities at $`1/2`$ and $`1/4`$ are seen prominently in the semiclassical results. Otherwise, it turns out that the semiclassical approximation captures many fine-scale features of the correlations, some of which will be discussed below. Figures (2a,b) are for specific one dimensional sections of the same quantities. The agreement is very good.
### D Classical features in the correlation
A strong (positive) correlation is indicated at the classical fixed points $`(0,0)`$ and $`(1,1)`$, with the rest of the significant correlations being negative. They are dominated by several classical structures as illustrated in Fig. (3). Here the $`N`$ value used is 200, and superposed on the significant contour features are the following classical orbits:
* the period-2 orbit at $`(1/3,2/3)`$, $`(2/3,1/3)`$ is by far the most prominent structure. This is shown in Fig. (3a) by two circular dots. Also, we can look at these structures closely through 1-d slices. In Fig. (2a), the correlation is seen to be large and negative at $`(q_\alpha =.66)`$. The period-2 structure is dominating the landscape;
* next in importance is the primary homoclinic orbit to the period 2 orbit in i), $`(1/3,1/3)`$, which goes to $`(2/3,1/6)`$, quickly gets into the region of the period two orbit and is difficult to resolve. The parity and time-reversal symmetric image points are also indicated. It turns out that there is an infinite set of periodic orbits which approximate this orbit more and more closely. Its effects may be present simultaneously, and indistinguishable from the homoclinic orbit itself . The relevant family (set) is denoted by $`(001)_{01}`$ which is based on a complete binary coding of the orbits . For example, the first few periodic orbits of the family are associated with the binary codes $`(00101)`$, $`(0010101)`$, and $`(001010101)`$. They are also shown in Fig. (3b), including the symmetric image points. In the 1-d slice of Fig. (2a) we see this orbit as well;
* there is an infinite number of orbits homoclinic to the period-2 orbit. They become increasingly more complicated. The next associated periodic orbit family, $`(0011)_{01}`$, is shown in Fig. (3c), including the symmetric image points. This family was noted by Saraceno to scar eigenfunctions . Also shown in this figure is the period-4 along the diagonal lines: $`(3/5,3/5)(1/5,4/5)(2/5,2/5)(4/5,1/5)`$. Figure (4) shows sections at $`p_\alpha =3/5,4/5,2/5`$ to highlight this orbit. In Figs. (4a) $`p_\alpha =3/5`$ and has a local minimum at $`q_\alpha =3/5`$; (b) $`p_\alpha =4/5`$ and has a local minimum at $`q_\alpha =1/5`$; and (c) $`p_\alpha =2/5`$ and has a local minimum at $`q_\alpha =2/5`$. These are marked, to indicate location along $`q_\alpha `$ by filled circles. The other minima are due to competing nearby structures of the period-2 orbit and its principal homoclinic excursion;
* the orbit homoclinic to the period-4 orbit included in Fig. (3c) with the initial condition $`(1/5,2/5)`$ (and its symmetric partners) is shown in Fig. (3d); and
* points, such as $`(0,1/4)`$, which are homoclinic to the fixed points $`(0,0)`$, $`(1,1)`$ also show prominently.
That these structures are in a sense invariant, i.e. not specific to $`N=100`$ is shown in Fig. (5a,b) where the correlation (absolute value) is shown for $`N=128`$ and $`200`$ respectively. The phase-space resolution of the correlation is increasing with $`N`$, while the overall magnitude is decreasing as $`1/N`$. The peculiar properties of the quantum bakers map for N equaling a power of two is tested by $`N=128`$. Here the correlation is “cleaner” and the stable and unstable manifolds at 1/4, 1/2, and 3/4 of the fixed points are clearly visible. The peaks are well enunciated as well. Both Figs. (5a,b) have contours up to 2/3 peak height, so a direct comparison is meaningful. Higher $`N`$ values show more clearly the secondary homoclinic orbit to the period-2 orbit.
We may compare these structures with the inverse participation ratio defined as:
$$P(\alpha )=\underset{i=0}{\overset{N1}{}}\left|\alpha |\psi _i\right|^4.$$
(29)
It is illustrated in Fig. (6). It shows marked enhancements at the period-2 and period-4 (along the symmetry lines) orbits, and closer examination reveals all orbits up to period-4 are present and one orbit of period-6 along the symmetry lines (the diagonals); see Ref. for a more detailed discussion.
### E General operators and selective enhancements
The results so far have dealt with the special case $`A(q)=\mathrm{cos}(2\pi q)`$. It seems natural to suspect that the structures highlighted in the correlation are dependent on the choice of the operator. This turns out to be true, and we show here how this works in the bakers map. We reemphasize though that were the eigenstates behaving ergodically, the correlations would have been consistent with zero to within statistical uncertainties independent of the choice of the operator. In this sense, a complete view of the extent to which the eigenstates manifest phase space localization properties comes only from considering both the full phase plane of wave packets and enough operators to span roughly the space of possible perturbations of the energy surface. The flexibility of operator choice does provide a means to enhance selectively particular features of interest supposing one had a specific localization question in mind. As an illustration, note that localization about the period-3 orbit barely appeared in the contour plot of Fig. (3), and yet, we show below that it can be made to show up prominently with other operators.
Since the case $`A(q)=\mathrm{sin}(2\pi nq)`$ has vanishing correlations for any integer $`n`$ due to symmetry, the other cases of interest are the higher harmonics of the cosine. Therefore consider:
$$A(q)=\mathrm{cos}(2\pi nq).$$
(30)
If $`n=2^m`$ for some positive integer $`m`$, a rather remarkable scaling property of the quantum bakers map is revealed that is actually implicit in the way the bakers map was originally quantized in . Semiclassically, the correlations are identical to the case $`m=0`$. For example, consider $`A(q)=\mathrm{cos}(4\pi q)`$. Then the one-step back classical correlation becomes identical to the zero-th order correlation corresponding to $`A(q)=\mathrm{cos}(2\pi q)`$. The correlations all shift by $`m`$ in the sense that $`C(l)C(l+m)`$. Thus, there is a kind of scale invariance in the correlation like classical fractals, although this is not self-similarity in the same curve. Quantum calculations reflect this invariance to a remarkable degree as seen in Fig. (7) where the $`N=200`$ and $`p_\alpha =1/3`$ case is shown.
Other harmonics do weight differently the same localization effects (classical structures). In Fig. (8), $`N`$ and $`p_\alpha `$ are taken the same as in Fig. (7). The cases $`n=1,3,5`$ are all very different from each other, but note that the case $`n=6`$ almost coincides with $`n=3`$ for the same reason that powers of two harmonics are nearly same. Thus only operators of odd harmonics give the possibility of providing new or unique information about the nonergodicity in the eigenstates. The period-2 orbit localization is accentuated at $`n=3`$, since for $`n=3m`$ where $`m`$ is a positive integer, $`\mathrm{cos}(6\pi mq)`$ has a maximum of $`+1`$ at $`q=2/3`$, whereas for all other integers $`n`$, $`\mathrm{cos}(2\pi nq)=1/2`$ at the same point. In short the perturbation (or measurement) is more significant at the location of the period-2 orbit for $`n=3`$. On the other hand, the case $`n=5`$ is similar to the fundamental harmonic case at $`q=2/3`$ where the perturbation is also equal to $`1/2`$.
The case $`n=7`$ is interesting as $`\mathrm{cos}(14\pi q)`$ has a maximum at $`q=1/7`$ which coincides with a period-3 orbit at $`(1/7,4/7)`$. In Fig. (9), we see the correlation ($`N=100)`$ corresponding to this operator and the dominant structure is this period-3 orbit and its symmetric partner. Also visible are the stable and unstable manifolds of these orbits. In fact, it is the multiples of the $`2^m1`$ harmonics which selectively highlights the period $`m`$ orbits.
Summarizing then, the correlations reflect that the bakers map eigenstates are not ergodic, and manifest strongly phase space localization properties. There do not exist transport barriers such as cantori or diffusive dynamics in the bakers map, so whatever localization that exists should be due to scarring by the short periodic orbits. This is confirmed in the examples shown with connections to their homoclinic orbits illustrated as well. The perturbation or observable determines the regions of phase space that will light up in the correlation measure. A semiclassical theory predicts reasonably well many of these structures. The correlation is semiclassically written as a sum of classical correlations that are super-exponentially cut off after about half the log-time scale.
## III The standard map
### A The map and the mixed phase space regime
The standard map (a review is found in ) has many complications that can arise in more generic models and we turn to their study. It is also an area-preserving, two-dimensional map of the cylinder onto itself that may be wrapped on a torus. We will consider identical settings of the phase space and Hilbert space as for the bakers map discussed above. The standard map has a parameter that controls the degree of chaos and thus we can study the effect of regular regions in phase space, i.e. the generic case of mixed dynamics.
The classical standard map is given by the recursion
$`q_{i+1}`$ $`=`$ $`(q_i+p_{i+1})\text{mod}(1)`$ (31)
$`p_{i+1}`$ $`=`$ $`(p_i(k/2\pi )\mathrm{sin}(2\pi q_i))\text{mod}(1),`$ (32)
where $`i`$ is the discrete time. The parameter $`k`$ is of principal interest and it controls the degree of chaos in the map. Classically speaking, an almost complete transition to ergodicity and mixing is attained above values of $`k5`$, while the last rotational KAM torus breaks around $`k.971`$.
The quantum map in the discrete position basis is given by
$`n|U|n^{}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{iN}}}\mathrm{exp}\left(i\pi (nn^{})^2/N\right)`$ (34)
$`\times \mathrm{exp}\left(i{\displaystyle \frac{kN}{2\pi }}\mathrm{cos}(2\pi (n+a)/N)\right).`$
The parameter to be varied will be the “kicking strength” $`k`$, while the phase $`a=1/2`$ for maximal quantum symmetries, and $`n,n^{}=\mathrm{\hspace{0.17em}0},\mathrm{},N1`$.
We use the unitary operator and evaluate the correlation as in Eq. (3) with $`A(q)=\mathrm{cos}(2\pi q)`$ here as well. This corresponds exactly to the level velocity induced by a change in the parameter $`k`$. In Fig. (10) is shown the absolute value of the correlation for various values of parameter $`k`$. Case (a) corresponds to $`k=.1\times (2\pi )`$ and is dominated by the KAM curves as the perturbation has not yet led to significant chaos. Highlighted is the fixed point resonance region at the origin that is initially stable. An unstable point is located at the point $`(1/2,0\mathrm{or}1)`$. The separatrix or the stable and the unstable manifolds of this point are aligned along the local ridges seen in the correlation. Also the period-2 resonance region is visible. Higher resolution not shown here, corresponding to higher values of $`N`$ reveal weakly the period-3 resonance as well. Case (b) corresponds to $`k=.3\times (2\pi )`$ while (c) and (d) to $`k=.9\times (2\pi )`$ and $`k=2.3\times (2\pi )`$ respectively. We note the gradual destruction of the KAM tori and the emergence of structures that are dominated by hyperbolic orbits. A more detailed classical-quantum correspondence is however not attempted here.
These contour plots do not reveal the difference in the magnitude between the correlations in the stable and unstable regions. In Fig. (11), we have plotted the correlation at the origin $`(0,0)`$, which is also a fixed point, as a function of the parameter. The value $`k/(2\pi )1`$ corresponds to a transition to complete classical chaos and is reflected in this plot as erratic and small oscillations. The large correlation in the mixed phase space regime arises from the non-ergodic nature of the classical dynamics. The classical fixed point loses stability at $`k^{}/(2\pi )=4/(2\pi ).63`$ and this is roughly the region at which the correlation starts to dip away from unity toward lower values.
The gross features and principal $`\mathrm{}`$ behavior in this regime is easy to derive in terms of purely classical correlations as follows:
$`C_A(\alpha )`$ $`=`$ $`\text{Tr}\left(|\alpha \alpha |\widehat{A}(n)\right)_n`$ (35)
$`=`$ $`{\displaystyle 𝑑q𝑑p[|\alpha \alpha |]_W[\widehat{A}(n)]_W}_n`$ (36)
where $`[.]_W`$ is the Weyl-Wigner transform of the operator in the brackets and $`\widehat{A}(n)`$ is the operator after a time $`n`$. Without worrying about the toral nature of the phase space and the Weyl-Wigner transforms, we treat the problem as in a plane. This is justified by the use of localized, Gaussian wave packets. Otherwise, we could imagine that the Wigner transform of the projector would follow from an infinite series of Gaussian states that is equivalent to discretizing the Gaussian. We use a normalized, “circular” Gaussian of width $`\sqrt{\mathrm{}}`$. The approximation comes in when we replace $`[\widehat{A}(n)]_W`$ by $`A(f^n(q,p))`$ where the latter is the classical function evaluated at the classically iterated point $`q_n=f^n(q,p)`$. We expect this approximation to be valid in the case of regular dynamics over a much longer time scale than found with chaotic dynamics. To a good approximation,
$`C_A(\alpha )`$ $``$ $`{\displaystyle }dqdp\left({\displaystyle \frac{1}{\pi \mathrm{}}}\right)\mathrm{exp}[{\displaystyle \frac{1}{\mathrm{}}}((qq_\alpha )^2`$ (38)
$`+(pp_\alpha )^2)]A\left(f^n(q,p)\right)_n`$
As intuitively expected, there is no principal $`\mathrm{}`$ dependence in the correlation since there is a non-zero classical limit. At $`\mathrm{}=0`$, we could replace the Gaussian forms by $`\delta `$functions and would get simply $`C_A(\alpha )=A(f^n(q_\alpha ,p_\alpha ))_n`$.
This however vanishes as the classical system becomes more ergodic and is no more capable of predicting the correlation. Higher-order corrections are needed. It is in this regime that we studied the bakers map and found that the correlation has a principal part that scales (almost) as $`\mathrm{}`$ and classical correlations based on periodic orbits predict the localization features that arise out of quantum interference.
We return to Fig. (11) to remark on some of these properties. Notice that the simple estimate of Eq. (38) performs very well, even as the phase space is becoming increasingly chaotic. It is quite unexpected that the oscillations after the onset of full mixing (around $`k/(2\pi )=1`$) should follow this estimate. However, after the transition to chaos the classical estimate will depend on the times over which the averaging is done and as this increases the estimate would vanish.
### B Chaotic regime
We attempt in some measure a semiclassical theory for the correlation in the chaotic regime along the lines adopted for the quantum bakers map. Of the two ingredients in Eq. (5) one of them remains the same, namely Eq. (9). However the diagonal elements of the propagator in Eq. (8) have to be generalized.
In a semiclassical expression for the matrix elements of the propagator as a homoclinic orbit sum is given. Although this was derived with the example of the billiard in mind, it can be interpreted as a generalization of Eq. (8) for area-preserving, two-dimensional maps. We, however, interpret the sum not as a homoclinic orbit sum, but as a periodic orbit sum. To each homoclinic orbit there is a neighboring periodic orbit that we will use instead. This will form the points around which the expansions are carried out and the result is identical to that in . Thus we write
$$\alpha |U^n|\alpha \underset{\gamma }{}\mathrm{exp}\left(iS_\gamma /\mathrm{}i\pi \nu /2\right)\underset{j}{}B_j$$
(39)
where
$`B_j`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{A_0}}}\mathrm{exp}\{{\displaystyle \frac{1}{2\mathrm{}A_0}}[(\text{tr}2)(\delta q^2+\delta p^2)`$ (41)
$`+\mathrm{\hspace{0.17em}2}i(m_{21}\delta p^2m_{12}\delta q^2+\delta q\delta p(m_{22}m_{11}))]\}.`$
$`B_j`$ generalizes the Gaussian form (including the prefactor) in Eq. (8). Again $`j`$ labels points along the periodic orbit and $`\delta q=q_jq_\alpha `$, $`\delta p=p_jp_\alpha `$ are as before deviations from the centroid of the wave packet. The two dimensional matrix elements, $`m_{ij}`$, are the elements of the stability matrix at the periodic point $`j`$ along the periodic orbit $`\gamma `$. The deviations $`\delta q`$ and $`\delta p`$ after $`n`$ iterations of the map are given by:
$$\left(\begin{array}{c}\delta p_n\\ \delta q_n\end{array}\right)=\left(\begin{array}{cc}m_{11}\hfill & m_{12}\hfill \\ m_{21}\hfill & m_{22}\hfill \end{array}\right)\left(\begin{array}{c}\delta p\\ \delta q\end{array}\right)$$
(42)
The invariant is the trace of this matrix that is denoted tr. While $`A_0=m_{11}+m_{22}+i(m_{21}m_{12})`$, $`\nu `$ is a phase that will not play a crucial role below. In the case of the bakers map $`m_{12}=m_{21}=0`$ and $`m_{11}=2^n,m_{22}=2^n`$ uniformly at all points in phase space, as well as $`\nu =0`$. On substitution of this in Eq. (39) we get Eq. (8).
The dependence on individual matrix elements of the stability matrix complicates the use of this formula in general. However we note that the Gaussian is effectively cutting off periodic points that are not close to $`\alpha `$ and therefore we may take the $`m_{ij}`$ elements to be the stability matrix at this point. In the chaotic regime each of the matrix elements grow exponentially with time $`n`$. Thus we have that $`\mathrm{exp}(\lambda n)m_{ij}`$ constant, where $`\lambda `$ is the Lyapunov exponent. We call this saturated constant $`m_{ij}`$ as well. Below we will assume that the exponential growth has been factored out of these elements. Also we use $`\mathrm{exp}(\lambda n)A_0a_0`$. The terms inside the exponential function in Eq. (41) saturate in time $`n`$ while the prefactor goes as $`\mathrm{exp}(\lambda n/2)`$. It follows then that
$$B_j\sqrt{\frac{2}{a_0}}\mathrm{exp}(\lambda n/2)F(q_j,p_j)$$
(43)
where $`F(q_j,p_j)`$ is
$`F(q_j,p_j)`$ $`=`$ $`\mathrm{exp}\{{\displaystyle \frac{1}{2\mathrm{}a_0}}[(\delta q^2+\delta p^2)`$ (44)
$`+`$ $`2i(m_{21}\delta p^2m_{12}\delta q^2+\delta q\delta p(m_{22}m_{11}))]\}.`$ (45)
Here the $`m_{ij}`$ elements already have the exponential behavior factored out. For example in the case of the bakers map $`m_{22}=1`$ while all the other elements are zero and this gives consistently the approximated Gaussian form in Eq. (12). Further steps are identical to the case of the bakers map and leads to the generalization of Eq. (13):
$$C_A(\alpha )=\sqrt{\frac{2}{a_0}}\underset{T}{}\underset{l=M}{\overset{M}{}}\stackrel{~}{C}_T(l),$$
(46)
where the classical correlations are calculated as in Eq. (14) with the function $`F`$ being that in Eq. (44). We may then expect all the principal conclusions from the study of the bakers map to be carried over, principally the decrease in the correlation as $`\mathrm{}`$, the correlations being cut off after half the log-time scale, and the effects of classical orbits.
More detailed analysis in the lines of the special case discussed in case of the bakers map will run into the following difficulties. First, the $`m_{ij}`$ elements will depend on $`\alpha `$ in general. Exceptions are uniformly hyperbolic systems such as the cat or sawtooth maps (and, of course, the bakers map). A second difficulty is that the correlations have to be evaluated to half the log-time while classically iterating the map (analytically) over such times is often not possible. The classical correlations that arise in the study of rms values of level velocities involved correlations that exponentially decreased in time while here we are likely to get generalizations of forms such as in Eq. (18) that will require us to go up to log-times. We have calculated the correlations for times 1, -1, and 2 but will not display them as they are by themselves not very useful. A third problem with this form of the generalization is that it is not explicitly real.
We have used Eq. (46) and for the $`m_{ij}`$ used either those calculated at one point in phase space (such as the origin) or in fact assumed those that are relevant for the bakers map. While fine structures are not reproduced, the general features are captured equally well in both these approaches. To illustrate the quality of the approximation we again look at the correlation at the point $`(0,0)`$ as a function of $`k`$ in Fig. (12) (as in the previous figure). The solid line is the semiclassical prediction based upon using the same $`m_{ij}`$ values at all values of $`k`$. It is seen that even with these (over) simplifications the semiclassical expressions capture much of the oscillations with the parameter and the magnitude.
## IV Summary and conclusions
We have studied the details of phase space localization present in the quantum time evolution of operators. This was related to a measure of localization involving the correlation between the level velocities and wavefunction intensities. While individual quantum states show well known interesting scars of classical orbits, groups of states weighted appropriately provide both a convenient and important quantity to study semiclassically. We were interested principally in those features whose origins were quantum mechanical. The quantities studied had both a vanishing classical limit as well as vanishing RMT averages.
We studied simple maps as a way to understand the general features that will appear. We found that the operator dictated to a large extent which parts of phase space will display prominent localization features and further that these localization features are often related to classical periodic orbits and their homoclinic structures. The time average of the operator for wave packets was explicitly related semiclassically to classical correlations. These were shown to be cut off super-exponentially after half the log-time in the quantum bakers map. Thus the localization features in quantum systems associated with scars were reproduced using long (periodic) orbits but short time correlations. The localization would disappear in the classical limit as the magnitude of the quantum correlations or time averages are proportional to (scaled) $`\mathrm{}`$.
General systems were approached using the quantum standard map and complications that would arise were discussed. Also the case of mixed phase space was seen to be well reproduced by a simple classical argument. The generalization to Hamiltonian systems contains many of the features and structures are also (not surprisingly) present in this case.
We gratefully acknowledge support from the National Science Foundation under Grant No. NSF-PHY-9800106 and the Office of Naval Research under Grant No. N00014-98-1-0079.
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# On the Inclusive Determination of |𝑽_{𝒖𝒃}| from the Lepton Invariant Mass Spectrum
## 1 Introduction
A precise knowledge of the magnitude of the Cabibbo–Kobayashi–Maskawa (CKM) matrix element $`V_{ub}`$ is of great importance to the study of CP violation at the $`B`$ factories. Whereas the main target of CP-asymmetry measurements is to fix the angles of the unitarity triangle, the length of the sides of this triangle are determined by the smallest CKM elements, $`|V_{ub}|`$ and $`|V_{td}|`$. The measurement of $`|V_{ub}|`$, which neither involves CP violation nor rare loop processes, appears to be the simplest step in the overall determination of the unitarity triangle. Yet, at present this measurement is limited by uncomfortably large theoretical uncertainties.
$`|V_{ub}|`$ can be determined most directly from semileptonic $`B`$ decays into charmless final states. The theoretical interpretation of exclusive decays such as $`B\pi l\nu `$ or $`B\rho l\nu `$ is limited by the necessity to predict the $`B\pi `$ or $`B\rho `$ transition form factors, which parameterize the complicated hadronic interactions relevant to these decays. The inclusive decays $`BXl\nu `$ admit a cleaner theoretical analysis based on a heavy-quark expansion . However, an obstacle is that experimentally it is necessary to impose restrictive cuts to suppress the background from $`BX_cl\nu `$ decays (i.e., decays into final states with charm). Accounting for such cuts theoretically is difficult and usually introduces significant uncertainties.
The first determination of $`|V_{ub}|`$ from inclusive decays was based on a measurement of the charged-lepton energy spectrum close to the endpoint region, which is kinematically forbidden for decays with a charm hadron in the final state. This restriction eliminates about 90% of the $`BX_ul\nu `$ signal, and it is difficult to calculate reliably the small fraction of the remaining events. The reason is that in this portion of phase space the conventional heavy-quark expansion breaks down and must be replaced by a twist expansion, in which an infinite tower of local operators is resummed into a shape function describing the light-cone momentum distribution of the $`b`$ quark inside a $`B`$ meson . A promising strategy is to infer the fraction of semileptonic decays in the endpoint region from a study of the photon spectrum in $`BX_s\gamma `$ decays, since the leading nonperturbative effects in the two decays are described by the same shape function . Recently, it has been emphasized that a cut $`M_X<M_D`$ on the hadronic invariant mass would provide for a better discrimination between $`BX_ul\nu `$ and $`BX_cl\nu `$ decays . Such a cut eliminates the charm background, while affecting only about 20% of the signal events. Despite of this advantage, however, it is difficult to calculate precisely what this fraction is .
Bauer, Ligeti and Luke (BLL) have pointed out that the situation may be better if a discrimination based on a cut on the lepton invariant mass $`\sqrt{q^2}`$ is employed . Requiring $`q^2>(M_BM_D)^2`$ eliminates the charm background, while containing about 20% of the signal events. This fraction is much less than in the case of a cut on hadronic invariant mass, but the inclusive rate with a lepton invariant mass cut offers the advantage of being calculable using a conventional heavy-quark expansion, without a need to resum an infinite series of nonperturbative corrections. BLL find that the fraction of $`BX_ul\nu `$ events with $`q^2>q_0^2`$, where $`q_0^2(M_BM_D)^2`$ is required to eliminate the charm background, is given by
$$F(q_0^2)=(1+\widehat{q}_0^2)(1\widehat{q}_0^2)^3+\stackrel{~}{X}(\widehat{q}_0^2)\frac{\alpha _s(m_b)}{\pi }\left(\widehat{q}_0^22\widehat{q}_0^6+\widehat{q}_0^8\right)\frac{12\lambda _2}{m_b^2}+\mathrm{},$$
(1)
where $`\widehat{q}_0^2=q_0^2/m_b^2`$, and the dots represent corrections of higher order in the expansion in powers of $`\alpha _s(m_b)`$ and $`\mathrm{\Lambda }/m_b`$ (with $`\mathrm{\Lambda }`$ a characteristic hadronic scale). The function $`\stackrel{~}{X}(\widehat{q}_0^2)`$ can be obtained from results presented in . The nonperturbative parameter $`\lambda _2=\frac{1}{4}(M_B^{}^2M_B^2)0.12`$ GeV<sup>2</sup> was introduced in . As it stands, eq. (1) seems to provide a solid theoretical basis for a systematic analysis of the partial decay rate in an expansion in logarithms and powers of $`\mathrm{\Lambda }/m_b`$. (This leaves aside the fact that the lepton invariant mass cut eliminates about 80% of all $`BX_ul\nu `$ events, and therefore one may worry that violations of quark–hadron duality may be larger than for the total inclusive semileptonic rate.)
The purpose of this work is to analyze the structure of the result (1) in more detail. We show that the relevant mass scale $`\mu _c`$ controlling the size of corrections in the heavy-quark expansion is less than or of order the charm-quark mass, rather than the heavier $`b`$-quark mass. This is so because the largest values of the hadronic invariant mass and energy accessible are of order the charm mass or less. Although the ratio $`m_c/m_b`$ is usually taken as a constant when discussing the heavy-quark limit, it is well known that the convergence of heavy-quark expansions at the charm scale can be poor. This makes our observation relevant. We suggest that an appropriate framework in which to investigate the leading corrections to the heavy-quark limit in the present case is a modified version of the heavy-quark expansion to which we refer as “hybrid expansion”. The idea is that in the kinematic region where $`M_BE_XM_X\mathrm{\Lambda }`$ one can perform a two-step expansion in the ratios $`E_X/M_B`$ and $`\mathrm{\Lambda }/E_X`$. A similar strategy has been used in applications of the heavy-quark effective theory (HQET) to resum the so-called “hybrid logarithms” $`\alpha _s\mathrm{ln}(m_c/m_b)`$ arising in current-induced $`bc`$ transitions using renormalization-group (RG) equations . In the context of inclusive decays, an approach similar in spirit to our proposal was suggested first by Mannel . We use the hybrid expansion to obtain a RG-improved expression for the perturbative corrections to the quantity $`F(q_0^2)`$ at next-to-leading order (NLO), as well as to estimate the size of higher-order power corrections omitted in (1).
## 2 Structure of the hadronic tensor
The strong-interaction dynamics relevant to inclusive $`B`$ decays is encoded in a hadronic tensor defined as the forward matrix element of the time-ordered product of two weak currents between $`B`$-meson states. Although the variable of prime interest to our discussion is the lepton invariant mass, the physics of the hadronic tensor is most naturally described in terms of the hadronic invariant mass and energy, $`M_X`$ and $`E_X`$. These variables are related by $`q^2=M_B^22M_BE_X+M_X^2`$, and thus the restriction $`q^2>q_0^2`$ implies
$$M_XE_X\frac{M_B^2+M_X^2q_0^2}{2M_B},M_\pi ^2M_X^2(M_B\sqrt{q_0^2})^2.$$
(2)
With the optimal choice $`q_0^2=(M_BM_D)^2`$ this gives $`M_\pi M_XM_D`$ and $`M_XE_XM_D\frac{1}{2}(M_D^2M_X^2)/M_B`$. Both variables vary between $`M_\pi `$ and $`M_D`$. If the cutoff $`q_0^2`$ is chosen higher, as may be required for experimental reasons, their maximal values become less than $`M_D`$.
The corresponding variables entering a partonic description of inclusive decay rates are the parton invariant mass and energy, $`\sqrt{p^2}`$ and $`vp`$, where $`p`$ is related to the lepton momentum $`q`$ by $`p=m_bvq`$, and $`v`$ is the velocity of the $`B`$ meson. The phase space for the dimensionless variables $`\widehat{p}^2=p^2/m_b^2`$ and $`z=2vp/m_b`$ is
$$2\sqrt{\widehat{p}^2}z1\widehat{q}_0^2+\widehat{p}^2,0\widehat{p}^2(1\sqrt{\widehat{q}_0^2})^2.$$
(3)
For the purpose of our discussion here $`m_b`$ is the pole mass of the $`b`$ quark. Alternative mass definitions will be discussed in more detail later. The optimal value of $`\widehat{q}_0^2`$ is $`\widehat{q}_0^2=(M_BM_D)^2/m_b^2(1\frac{m_c}{m_b})^2`$. In this case the largest values of the parton variables are $`z_{\mathrm{max}}=2\frac{m_c}{m_b}`$ and $`\widehat{p}_{\mathrm{max}}^2=(\frac{m_c}{m_b})^2`$. Without a restriction on $`q^2`$, the phase space for these variables would be such that $`z_{\mathrm{max}}=2`$ and $`\widehat{p}_{\mathrm{max}}^2=1`$. In other words, the lepton invariant mass cut restricts both variables to a region where they are parametrically suppressed, such that $`vp`$ and $`\sqrt{p^2}`$ are at most of order $`m_c`$. Besides $`m_b`$, the charm-quark mass is then a relevant scale that determines the magnitude of strong-interaction effects in the heavy-quark expansion.
To make the parametric suppression noted above more explicit, we introduce a characteristic scale $`\mu _c`$ and an associated small expansion parameter $`ϵ`$ by
$$\mu _c=\frac{m_b^2q_0^2}{2m_b}=O(m_c),ϵ=\frac{1\widehat{q}_0^2}{2}=\frac{\mu _c}{m_b},$$
(4)
and rewrite the fraction $`F(q_0^2)`$ of $`BX_ul\nu `$ events with $`q^2>q_0^2`$ as
$$F(ϵ)=16ϵ^3(1ϵ)\left[C_0(ϵ)\frac{3}{2}C_2(ϵ)\frac{\lambda _2}{\mu _c^2}+O[(\mathrm{\Lambda }/\mu _c)^3]\right],$$
(5)
where $`C_n(ϵ)=1+O(ϵ,\alpha _s)`$ are short-distance coefficients. From (1) we find
$`C_0(ϵ)`$ $`=`$ $`1+{\displaystyle \frac{\alpha _s}{3\pi }}\left[\left(6\mathrm{ln}2ϵ{\displaystyle \frac{3}{2}}{\displaystyle \frac{2\pi ^2}{3}}\right)+ϵ\left(8\mathrm{ln}2ϵ+{\displaystyle \frac{73}{6}}\right)+O(ϵ^2)\right]+O(\alpha _s^2),`$
$`C_2(ϵ)`$ $`=`$ $`{\displaystyle \frac{18ϵ^2+8ϵ^3}{1ϵ}}+O(\alpha _s).`$ (6)
The definition of $`ϵ`$ and $`\mu _c`$, as well as the explicit form of the coefficients $`C_n(ϵ)`$, depend on the definition of the heavy-quark mass. The above result for $`C_0(ϵ)`$ refers to the pole mass. Later we will introduce a more suitable mass definition and give an exact NLO expression for $`C_0(ϵ)`$ including the higher-order terms in $`ϵ`$.
It is apparent from (5) that the leading power correction proportional to $`\lambda _2`$ is of order $`(\mathrm{\Lambda }/\mu _c)^2`$. It is not difficult to see that, in the presence of a lepton invariant mass cut, also in higher orders the power corrections scale like $`(\mathrm{\Lambda }/\mu _c)^n`$. For simplicity of the argument we work to leading order in $`\alpha _s`$, where only the tree diagram shown in Figure 1 contributes to the hadronic tensor. In momentum space, the $`u`$-quark propagator gives a contribution
$$\frac{(p+k)^\mu }{p^2+2pk+k^2},$$
(7)
where $`k=O(\mathrm{\Lambda })`$ is the residual momentum of the heavy quark inside the $`B`$ meson . Roughly speaking, the heavy-quark expansion is obtained by replacing the residual momentum with a covariant derivative, $`k^\mu iD^\mu `$, thereby introducing the soft interactions of the $`u`$-quark jet with the background field of the light degrees of freedom in the $`B`$ meson. Let us discuss how the three terms in the denominator of the propagator scale in the different kinematic regions relevant to the determination of $`|V_{ub}|`$. Whereas the $`k^2\mathrm{\Lambda }^2`$ term is always suppressed, the relative magnitude of the two other terms depends on the kinematic region considered. For the regions of the large charged-lepton energy or low hadronic invariant mass, the terms $`p^2\mathrm{max}[\mathrm{\Lambda }m_b,m_c^2]`$ and $`pk\mathrm{\Lambda }m_b`$ are of the same magnitude, and it is thus necessary to resum terms of the form $`(pk/p^2)^n`$ to all orders in the heavy-quark expansion. This leads to a twist expansion, where these terms are absorbed into a nonperturbative shape function . In contrast, for large lepton invariant mass, $`pk\mathrm{\Lambda }m_c`$ is parametrically suppressed with respect to $`p^2m_c^2`$ by a power of $`\mathrm{\Lambda }/m_c`$.
These general observations can also be derived from the explicit expression for the differential decay rate expressed in terms of the variables $`z`$ and $`\widehat{p}^2`$, normalized to the total decay rate. At NLO in the heavy-quark expansion we find
$`{\displaystyle \frac{1}{\mathrm{\Gamma }}}{\displaystyle \frac{\text{d}^2\mathrm{\Gamma }}{\text{d}z\text{d}\widehat{p}^2}}`$ $`=`$ $`2z^2(32z)\delta (\widehat{p}^2)+{\displaystyle \frac{\alpha _s}{3\pi }}E(z,\widehat{p}^2)`$ (8)
$`\delta (\widehat{p}^2)\left[{\displaystyle \frac{z}{3}}(3627z16z^2){\displaystyle \frac{\lambda _1}{m_b^2}}+(12+12z63z^2+8z^3){\displaystyle \frac{\lambda _2}{m_b^2}}\right]`$
$`\delta ^{}(\widehat{p}^2)\left[{\displaystyle \frac{z^2}{3}}(18+3z14z^2){\displaystyle \frac{\lambda _1}{m_b^2}}+z^2(6+3z10z^2){\displaystyle \frac{\lambda _2}{m_b^2}}\right]`$
$`\delta ^{\prime \prime }(\widehat{p}^2){\displaystyle \frac{z^4}{3}}(32z){\displaystyle \frac{\lambda _1}{m_b^2}}+\mathrm{},`$
where $`\lambda _1`$ is a nonperturbative parameter related to the average kinetic energy of the $`b`$ quark inside the $`B`$ meson , and the function $`E(z,\widehat{p}^2)`$ gives the perturbative correction calculated in , which also includes the correction to the total decay rate appearing in the denominator on the left-hand side. The power corrections in (8) have been derived using the results of . In the kinematic region where $`z=O(ϵ)`$ and $`\widehat{p}^2=O(ϵ^2)`$, it is instructive to change variables from $`(z,\widehat{p}^2)`$ to $`(z,\xi )`$, where $`\xi =4\widehat{p}^2/z^2=p^2/(vp)^2[0,1]`$ is related to the parton velocity in the $`B`$ rest frame. The kinematic range for these variables is
$$0\xi 1,0z\frac{2}{\xi }(1\sqrt{12ϵ\xi })=2ϵ+ϵ^2\xi +O(ϵ^3).$$
(9)
The variable $`\xi `$ is of order unity irrespective of the lepton invariant mass cut. Thus, after the transformation the only parametrically small quantity is $`z=O(ϵ)`$. In terms of the new variables, the double-differential decay rate turns into an expansion in powers of $`\mathrm{\Lambda }/(zm_b)`$, and the perturbative corrections contain single logarithms of $`z`$. Integrating the double-differential rate over $`z`$, and keeping only terms of leading order in $`ϵ`$, we obtain
$`{\displaystyle \frac{1}{\mathrm{\Gamma }}}{\displaystyle \frac{\text{d}\mathrm{\Gamma }}{\text{d}\xi }}`$ $`=`$ $`16ϵ^3\{\delta (\xi )[1{\displaystyle \frac{\alpha _s}{3\pi }}(6\mathrm{ln}{\displaystyle \frac{\mu _c}{m_b}}+8\mathrm{ln}2(\mathrm{ln}21)+{\displaystyle \frac{31}{2}})]`$ (10)
$`{\displaystyle \frac{\alpha _s}{3\pi }}\left[{\displaystyle \frac{4\mathrm{ln}\xi }{\xi }}+{\displaystyle \frac{1}{\xi }}\left(7\sqrt{1\xi }8\mathrm{ln}(1+\sqrt{1\xi })\right)\right]_+`$
$`{\displaystyle \frac{3\lambda _2}{2\mu _c^2}}\delta (\xi ){\displaystyle \frac{3(\lambda _1+\lambda _2)}{\mu _c^2}}\delta ^{}(\xi ){\displaystyle \frac{2\lambda _1}{\mu _c^2}}\delta ^{\prime \prime }(\xi )+\mathrm{}\}+O(ϵ^4).`$
The integral of this expression over $`\xi `$ reproduces the leading terms in $`ϵ`$ in (5). The result for the power corrections shows that indeed $`\mu _c=O(m_c)`$ is the characteristic scale of the hybrid expansion. Perturbative logarithms of $`\mu _c/m_b`$ appear because the intrinsic scale of the hadronic tensor is smaller than the mass of the external $`b`$ quarks, suggesting that the appropriate scale to evaluate the running coupling $`\alpha _s`$ is significantly less than $`m_b`$. This is in accordance with the observation that the physical scale derived using the Brodsky–Lepage–Mackenzie (BLM) scale-setting prescription strongly decreases with increasing $`q^2`$. We will come back to the question of scale setting in the next section.
For later purposes, it will be useful to have yet another way of reproducing the result (5). To this end, we represent the fraction $`F(q_0^2)`$ as a contour integral in the complex $`\widehat{p}^2`$ plane. Such a representation exists because the hadronic tensor is an analytic function in the complex plane apart from discontinuities located on the positive real $`\widehat{p}^2`$ axis, as illustrated in Figure 2. We write
$$F(q_0^2)=\frac{i}{2\pi }\underset{|\widehat{p}^2|=\widehat{p}_{\mathrm{max}}^2}{}\text{d}\widehat{p}^2T(\widehat{p}^2,\widehat{q}_0^2),$$
(11)
where $`\widehat{p}_{\mathrm{max}}^2=(1\sqrt{\widehat{q}_0^2})^2`$, and the correlator $`T(\widehat{p}^2,\widehat{q}_0^2)`$ can be obtained using dispersion relations . Eliminating $`\widehat{q}_0^2`$ in favor of $`ϵ`$, and parameterizing $`\widehat{p}^2=\widehat{p}_{\mathrm{max}}^2e^{i\phi }`$ on the contour of integration, the result can be written in the form
$$F(ϵ)=16ϵ^3(1ϵ)\underset{0}{\overset{2\pi }{}}\frac{\text{d}\phi }{2\pi }t(e^{i\phi },ϵ),$$
(12)
where the function $`t(e^{i\phi },ϵ)`$ contains the corrections to the heavy-quark limit. We obtain
$`t(e^{i\phi },ϵ)`$ $`=`$ $`C_0(ϵ)+{\displaystyle \frac{\alpha _s}{3\pi }}Y(e^{i\phi },ϵ){\displaystyle \frac{3}{2}}C_2(ϵ){\displaystyle \frac{\lambda _2}{\mu _c^2}}+\left[D_1(ϵ){\displaystyle \frac{\lambda _1}{\mu _c^2}}+D_2(ϵ){\displaystyle \frac{\lambda _2}{\mu _c^2}}\right]e^{i\phi }`$ (13)
$`{\displaystyle \frac{4}{5}}E_1(ϵ){\displaystyle \frac{\lambda _1}{\mu _c^2}}e^{2i\phi }+O[(\mathrm{\Lambda }/\mu _c)^3].`$
The function $`Y(e^{i\phi },ϵ)`$ with
$$Y(x,ϵ)=\frac{10}{3}x\frac{184}{9}+\left(\frac{29}{3}\frac{5}{3}x\right)u\mathrm{ln}\frac{u+1}{u1}2\mathrm{ln}^2\frac{u+1}{u1}+O(ϵ)$$
(14)
and $`u=\sqrt{1x}`$ is defined such that its contour integral vanishes. The coefficients $`D_i(ϵ)`$ and $`E_1(ϵ)`$ are equal to 1 in the limit $`ϵ0`$. Their explicit expressions are
$`D_1(ϵ)`$ $`=`$ $`{\displaystyle \frac{ϵ^2}{(1ϵ)(1\sqrt{12ϵ})^2}}\left(1{\displaystyle \frac{7}{4}}ϵ+{\displaystyle \frac{4}{5}}ϵ^2\right)+O(\alpha _s),`$
$`D_2(ϵ)`$ $`=`$ $`{\displaystyle \frac{ϵ^2}{(1ϵ)(1\sqrt{12ϵ})^2}}\left(1+{\displaystyle \frac{3}{4}}ϵ4ϵ^2\right)+O(\alpha _s),`$
$`E_1(ϵ)`$ $`=`$ $`{\displaystyle \frac{ϵ^4}{(1ϵ)(1\sqrt{12ϵ})^4}}\left(1{\displaystyle \frac{10}{9}}ϵ\right)+O(\alpha _s).`$ (15)
Equations (10) and (13) allow us to study the behavior of the leading corrections in the heavy-quark expansion in more detail. In particular, we will utilize them to estimate unknown, higher-order power corrections. We will, however, first investigate the perturbative corrections in more detail, using the hybrid expansion as a tool to perform a systematic RG improvement of the one-loop expressions in (1) and (2).
## 3 RG improvement and definition of the heavy-quark mass
Based on the observation that the event fraction $`F(q_0^2)`$ in (1) receives very small corrections of order $`\beta _0\alpha _s^2`$, where $`\beta _0`$ is the first coefficient of the $`\beta `$-function, BLL have argued that the perturbative uncertainty in their prediction is negligible . The purpose of this section is to critically reanalyze the perturbative uncertainty in the calculation of this quantity. We first note that in the present case it is misleading to associate the size of $`\beta _0\alpha _s^2`$ corrections with a physical BLM scale in the process. The total semileptonic rate and the lepton invariant mass spectrum receive very large corrections of order $`\beta _0\alpha _s^2`$, corresponding to very low BLM scales . The fraction $`F(ϵ)`$ is defined as the ratio of the partially integrated lepton spectrum and the total rate. If both of these quantities have low physical scales, the same must be true for their ratio. In our opinion the small $`\beta _0\alpha _s^2`$ correction observed in is thus due to an accidental cancellation and does not bare any physical significance.
Here we follow a different strategy to estimate the potential importance of higher-order effects. We have argued that a useful framework in which to analyze the fraction $`F(ϵ)`$ is provided by a hybrid expansion, in which the physics associated with the three mass scales $`m_b\mu _c\mathrm{\Lambda }`$ is disentangled. Since the intermediate scale $`\mu _c`$ is only about 1 GeV or less, and since the running of the strong coupling in the region between $`m_b`$ and 1 GeV is significant, we expect important higher-order perturbative corrections resolving the scale ambiguity. At one-loop order, this is indicated by the presence of logarithms of $`ϵ`$ in (2). At any given order in an expansion in powers of $`ϵ`$ the contributions associated with the two couplings $`\alpha _s(m_b)`$ and $`\alpha _s(\mu _c)`$ can be separated by solving RG equations in the hybrid expansion at NLO. The residual scale ambiguity left after RG improvement provides an estimate of the perturbative uncertainty in the result. Unfortunately, to perform this program one must compute, at every order in $`ϵ`$, the two-loop anomalous dimensions of a new tower of higher-dimensional operators. At present, these anomalous dimensions are known only for the operators entering at the leading order in $`ϵ`$, although partial results exist for the operators relevant to the $`O(ϵ)`$ terms .
We now discuss the RG improvement at leading order in $`ϵ`$ in detail. The first step in the construction of the hybrid expansion is to expand the weak currents in the definition of the hadronic tensor in terms of operators of the HQET, with the result
$$\overline{q}\gamma _\mu (1\gamma _5)bC_1\left(\frac{m_b}{\mu }\right)\overline{q}\gamma _\mu (1\gamma _5)h_v+C_2\left(\frac{m_b}{\mu }\right)\overline{q}v_\mu (1+\gamma _5)h_v+O(1/m_b),$$
(16)
where $`v_\mu `$ is the $`B`$-meson velocity, $`h_v`$ are the velocity-dependent fields of the HQET, and $`\mu `$ is the scale at which the operators are renormalized. The RG-improved expressions for the Wilson coefficients $`C_i(m_b/\mu )`$ are known at NLO. The terms of order $`1/m_b`$ in (16) would contribute at order $`ϵ`$ in the hybrid expansion and can be neglected for the discussion of the leading terms. It is important that all dependence on the $`b`$-quark mass is explicit in (16). We now insert this result into the time-ordered product of currents in the hadronic tensor and perform an operator product expansion of the current product. This is an expansion in logarithms and powers of $`\mathrm{\Lambda }/\mu _c`$. The scale $`m_b`$ does not appear in the matrix elements of the hybrid expansion. We find that at leading order in $`ϵ`$ only the HQET current product proportional to $`C_1^2(m_b/\mu )`$ contributes. Using the known NLO expression for the coefficient $`C_1(m_b/\mu )`$ , we obtain
$`C_0(ϵ)`$ $`=`$ $`\left({\displaystyle \frac{\alpha _s(\mu )}{\alpha _s(m_b)}}\right)^{\frac{4}{\beta _0}}[1+{\displaystyle \frac{2\alpha _s(m_b)}{\pi }}(Z_{\mathrm{hl}}{\displaystyle \frac{25}{12}}+{\displaystyle \frac{\pi ^2}{3}})`$ (17)
$`{\displaystyle \frac{2\alpha _s(\mu )}{\pi }}(Z_{\mathrm{hl}}\mathrm{ln}{\displaystyle \frac{\mu }{2\mu _c}}{\displaystyle \frac{11}{6}}+{\displaystyle \frac{4\pi ^2}{9}})]+O(ϵ)`$
$``$ $`\left({\displaystyle \frac{\alpha _s(\mu )}{\alpha _s(m_b)}}\right)^{\frac{12}{25}}\left[10.71{\displaystyle \frac{\alpha _s(m_b)}{\pi }}{\displaystyle \frac{\alpha _s(\mu )}{\pi }}\left(3.372\mathrm{ln}{\displaystyle \frac{\mu }{\mu _c}}\right)\right]+O(ϵ).`$
Here $`\beta _0=\frac{25}{3}`$ and $`Z_{\mathrm{hl}}=\frac{9403}{7500}\frac{7\pi ^2}{225}`$ are perturbative coefficients evaluated for $`n_f=4`$ light quark flavors, as is appropriate for a scale of order $`m_c`$. Note that the NLO corrections proportional to the coupling $`\alpha _s(m_b)`$ contain the corrections to the total semileptonic rate. The above result for $`C_0(ϵ)`$ gives the RG-improved form of the leading term in $`ϵ`$ in the one-loop expression in (2). This result is formally scale independent at NLO. The renormalization scale $`\mu `$ should be chosen of order $`\mu _c`$ in order to avoid large logarithms in the hybrid expansion. (The inappropriate choice $`\mu =m_b`$ would reproduce the one-loop result with $`\alpha _s`$ evaluated at $`m_b`$.)
The dashed line in Figure 3 shows the perturbative prediction for the fraction $`F(ϵ)`$ at leading order in $`ϵ`$, as a function of the renormalization scale. Here and below we use the two-loop expression for the running coupling constant normalized such that $`\alpha _s(M_Z)=0.118`$. In accordance with relation (19) below we use $`m_b=5.0`$ GeV for the pole mass, noting that this value (and thus the normalization of the dashed curve) has a large uncertainty. Two observations are important. First, the scale dependence of the dashed curve is significant, and hence the result obtained with an appropriate choice of scale $`\mu \mu _c`$ is much lower than that obtained with the naive choice $`\mu =m_b`$. Secondly, perturbation theory in the on-shell scheme breaks down at a scale not much less than the appropriate scale $`\mu \mu _c1`$ GeV. We conclude that in the on-shell scheme there is a large perturbative uncertainty in the calculation of the coefficient $`C_0(ϵ)`$, which is not apparent from the naive one-loop result. This conclusion is in contrast with the assumption made by BLL, that the perturbative uncertainty is negligible .
The breakdown of perturbation theory at a scale of order $`\mu _c`$ can be traced back to the large coefficient of the NLO correction proportional to $`\alpha _s(\mu )`$ in (17). We will now show that the size of this coefficient can be reduced significantly by adopting a more appropriate definition of the heavy-quark mass. So far we have worked in the on-shell scheme, where the mass is defined as the pole in the renormalized quark propagator, $`m_b=m_b^{\mathrm{pole}}`$. Since the pole mass is affected by IR renormalon ambiguities , it is better to eliminate it from the final expressions for inclusive decay rates. If a new mass definition $`m_b^{}`$ is introduced via $`m_b^{\mathrm{pole}}=Z_mm_b^{}`$, it follows from (5) that
$$C_0^{}(ϵ^{})=C_0^{\mathrm{pole}}(ϵ^{})\left[1+\frac{Z_m1}{ϵ^{}}\frac{(12ϵ^{})(34ϵ^{})}{1ϵ^{}}+\mathrm{}\right],$$
(18)
where the prime on $`ϵ`$ indicates that this parameter in sensitive to the definition of $`m_b`$. We observe that a multiplicative redefinition of $`m_b`$ with $`Z_m(\alpha _s)=1+O(\alpha _s)`$, such as the relation between the pole mass and the running mass defined in the $`\overline{\mathrm{MS}}`$ scheme, is not appropriate in our case, since it would lead to a contribution to $`C_0^{}(ϵ^{})`$ that is enhanced by a factor of $`\alpha _s/ϵ^{}`$. This would upset the power counting in the hybrid expansion. We suggest instead to work with a short-distance mass subtracted at a scale $`\mu _f`$ of order $`\mu _c`$, which is the natural scale of our problem.
It is well known that the convergence of the perturbative series for near on-shell problems in heavy-quark physics can be largely improved by introducing low-scale subtracted quark masses, which have the generic property that they differ from the pole mass by an amount proportional to the subtraction scale $`\mu _f`$. Several such mass definitions exist and have been applied to various processes . To illustrate the point, we use the potential-subtracted (PS) mass $`m_b^{\mathrm{PS}}(\mu _f)`$ introduced by Beneke and evaluate it at the scale $`\mu _f=\mu _c^{\mathrm{PS}}`$.<sup>1</sup><sup>1</sup>1We could instead evaluate the PS mass at a scale $`\mu _f=\xi \mu _c`$ with $`\xi =O(1)`$, however this would lead to more complicated expressions. Varying $`\xi `$ between 1 and 2 leads to a variation of the results by an amount similar to the perturbative uncertainty estimated later in this section. At NLO, the relation between the pole mass and the PS mass reads
$$m_b^{\mathrm{pole}}=m_b^{\mathrm{PS}}(\mu _c^{\mathrm{PS}})+\mu _c^{\mathrm{PS}}\frac{4\alpha _s(\mu )}{3\pi }\left\{1+\frac{\alpha _s(\mu )}{2\pi }\left[\beta _0\left(\mathrm{ln}\frac{\mu }{\mu _c^{\mathrm{PS}}}+\frac{11}{6}\right)4\right]+\mathrm{}\right\},$$
(19)
which is formally independent of the scale $`\mu `$ at which the coupling is renormalized. Note that the difference between the two mass definitions is a perturbative series multiplying the scale $`\mu _c^{\mathrm{PS}}=ϵ^{\mathrm{PS}}m_b^{\mathrm{PS}}`$. At NLO in $`\alpha _s`$, it then follows from (18) that
$`C_0^{\mathrm{PS}}(ϵ^{\mathrm{PS}})`$ $`=`$ $`C_0^{\mathrm{pole}}(ϵ^{\mathrm{PS}})\left[1+{\displaystyle \frac{4\alpha _s(\mu )}{3\pi }}{\displaystyle \frac{\mu _c^{\mathrm{PS}}}{ϵ^{\mathrm{PS}}m_b^{\mathrm{PS}}}}{\displaystyle \frac{(12ϵ^{\mathrm{PS}})(34ϵ^{\mathrm{PS}})}{1ϵ^{\mathrm{PS}}}}\right]`$ (20)
$`=`$ $`C_0^{\mathrm{pole}}(ϵ^{\mathrm{PS}})\left[1+{\displaystyle \frac{4\alpha _s(\mu )}{\pi }}+O(ϵ^{\mathrm{PS}})\right].`$
From now on we will use the PS mass $`m_bm_b^{\mathrm{PS}}(\mu _c^{\mathrm{PS}})`$ in all our equations and omit the label “PS” on the quantities $`m_b`$, $`\mu _c`$ and $`ϵ`$. Using (17) and adding the extra contribution proportional to $`\alpha _s(\mu )`$, we obtain
$$C_0^{\mathrm{PS}}(ϵ)\left(\frac{\alpha _s(\mu )}{\alpha _s(m_b)}\right)^{\frac{12}{25}}\left[10.71\frac{\alpha _s(m_b)}{\pi }+\frac{\alpha _s(\mu )}{\pi }\left(0.63+2\mathrm{ln}\frac{\mu }{\mu _c}\right)\right]+O(ϵ).$$
(21)
The introduction of the PS mass has much reduced the size of the NLO correction. The result for the leading contribution to $`C_0(ϵ)`$ in the PS scheme is shown by the solid line in Figure 3. It exhibits a better stability that in the on-shell scheme, and it is stable down to lower values of the renormalization scale.
The value of the PS mass at the scale $`\mu _2=2`$ GeV has been determined from a sum-rule analysis of the $`b\overline{b}`$ production cross section near threshold, with the result $`m_b^{\mathrm{PS}}(2\mathrm{GeV})=(4.59\pm 0.08)`$ GeV (corresponding to $`\overline{m}_b(m_b)=(4.25\pm 0.08)`$ GeV in the $`\overline{\mathrm{MS}}`$ scheme) . At NLO, we can use relation (19) to convert this into a value of the PS mass at the scale $`\mu _c`$. This gives the implicit equation
$`m_b^{\mathrm{PS}}(\mu _c)`$ $`=`$ $`m_b^{\mathrm{PS}}(\mu _2)+\mu _c{\displaystyle \frac{4\alpha _s(\mu _2)}{3\pi }}\left[\left({\displaystyle \frac{\mu _2}{\mu _c}}1\right)\left(1+{\displaystyle \frac{203}{36}}{\displaystyle \frac{\alpha _s(\mu _2)}{\pi }}\right){\displaystyle \frac{25}{6}}{\displaystyle \frac{\alpha _s(\mu _2)}{\pi }}\mathrm{ln}{\displaystyle \frac{\mu _2}{\mu _c}}\right]`$ (22)
$`\stackrel{!}{=}`$ $`\mu _c+\sqrt{\mu _c^2+q_0^2},`$
from which we determine the scale $`\mu _c`$ and then the mass $`m_b=m_b^{\mathrm{PS}}(\mu _c)`$. For example, we find $`m_b4.73`$ GeV, $`\mu _c1.13`$ GeV, $`ϵ0.24`$ for $`q_0^2=(M_BM_D)^2`$, and $`m_b4.79`$ GeV, $`\mu _c0.83`$ GeV, $`ϵ0.17`$ for $`q_0^2=15`$ GeV<sup>2</sup>.
In (21) we have obtained a RG-improved expression for the short-distance coefficient at leading order in $`ϵ`$. It is at present not possible to extend this analysis to higher orders in the hybrid expansion, since the corresponding two-loop anomalous dimensions of higher-dimensional operators are unknown. However, since in the PS scheme the leading term in $`ϵ`$ gives the dominant contribution to $`C_0(ϵ)`$, we expect that the unresolved scale ambiguity in the higher-order terms does not introduce a large uncertainty. Our final expression for the short-distance coefficient at NLO is
$$C_0^{\mathrm{PS}}(ϵ)\left(\frac{\alpha _s(\mu )}{\alpha _s(m_b)}\right)^{\frac{12}{25}}\left[10.71\frac{\alpha _s(m_b)}{\pi }+\frac{\alpha _s(\mu )}{\pi }\left(0.63+2\mathrm{ln}\frac{\mu }{\mu _c}\right)+\frac{ϵ\overline{\alpha }_s}{\pi }G(ϵ)\right],$$
(23)
where the scale in $`\overline{\alpha }_s`$ in the $`O(ϵ)`$ term remains undetermined. The exact result for the function $`G(ϵ)`$ in the PS scheme is
$`G(ϵ)`$ $`=`$ $`{\displaystyle \frac{1}{12(1ϵ)ϵ^4}}\left[L_2(12ϵ){\displaystyle \frac{\pi ^2}{6}}\left({\displaystyle \frac{13}{12}}8ϵ+28ϵ^2{\displaystyle \frac{128}{3}}ϵ^3+20ϵ^4\right)\mathrm{ln}(12ϵ)\right]`$ (24)
$`{\displaystyle \frac{1}{6(1ϵ)ϵ^3}}\left[{\displaystyle \frac{1}{12}}{\displaystyle \frac{89}{12}}ϵ+21ϵ^2+40ϵ^364ϵ^4+\left(1+ϵ+{\displaystyle \frac{4}{3}}ϵ^2+2ϵ^3\right)\mathrm{ln}2ϵ\right]`$
$`+{\displaystyle \frac{4}{3ϵ}}\left[L_2(2ϵ)L_2(12ϵ)+{\displaystyle \frac{\pi ^2}{6}}\right]`$
$`=`$ $`{\displaystyle \frac{8}{3}}\mathrm{ln}2ϵ{\displaystyle \frac{95}{18}}+ϵ\left({\displaystyle \frac{32}{15}}\mathrm{ln}2ϵ+{\displaystyle \frac{1337}{450}}\right)+O(ϵ^2).`$
The left-hand plot in Figure 4 shows $`C_0(ϵ)`$ for the optimal choice $`q_0^2=(M_BM_D)^2`$ as a function of the renormalization scale. The width of the band reflects the sensitivity of the result to the value of the coupling $`\overline{\alpha }_s`$ associated with the $`O(ϵ)`$ terms in (23), which we vary between $`\alpha _s(m_b)`$ and $`\alpha _s(\mu )`$. To estimate the residual scale dependence we vary $`\mu `$ between the values $`\mu _c`$ and $`2\mu _c`$. For lower values the perturbative expansion diverges, since the running coupling $`\alpha _s(\mu )`$ strongly increases below $`\mu 1`$ GeV. For comparison, we mention that the naive perturbative analysis with fixed scale $`\mu =m_b`$ adopted in would give the much smaller value $`C_0(ϵ)1.16`$ at minimal $`q_0^211.6`$ GeV<sup>2</sup>. The fact that we find larger QCD corrections will have important implications for the extraction of $`|V_{ub}|`$. The right-hand plot in Figure 4 shows $`C_0(ϵ)`$ as a function of $`q_0^2`$. The width of the band represents the total scale uncertainty, estimated by variation of $`\mu `$ and $`\overline{\alpha }_s`$ as described above.
An independent way to estimate the uncertainty in the value of the perturbative coefficient $`C_0(ϵ)`$ is based on the contour representation (12) for the fraction $`F(ϵ)`$. As we have just discussed, the value of $`C_0(ϵ)`$ depends on the definition of the heavy-quark mass. However, the variation of the $`O(\alpha _s)`$ correction in (13) along the circle in the complex momentum plane is independent of mass redefinitions. We may thus take the $`\phi `$-variation of the one-loop correction, given by $`\alpha _s(\mu _c)/3\pi `$ times the variation of the real part of the function $`Y(e^{i\phi },ϵ)`$ in (14), as a typical size of an $`O(\alpha _s)`$ correction in the problem at hand. For an asymptotic series, the value of that correction provides an estimate for the magnitude of unknown higher-order corrections. The dashed lines in the right-hand plot in Figure 4 show this variation as an error band applied to the central values of $`C_0(ϵ)`$. This independent evaluation of higher-order effects is in good agreement with our previous estimate of the perturbative uncertainty, giving us confidence that this estimate is a realistic one.
## 4 Higher-order power corrections
Uncertainties enter the theoretical prediction for the fraction $`F(ϵ)`$ also at the level of power correction. First, there are unknown $`O(\alpha _s)`$ corrections to the Wilson coefficient $`C_2(ϵ)`$ multiplying the term proportional to $`\lambda _2/\mu _c^2`$ in (5). To estimate their effect, we replace the bracket $`[\mathrm{}]`$ in this equation with $`C_0(ϵ)[1\frac{3}{2}C_2(ϵ)\lambda _2/\mu _c^2+\mathrm{}]`$, which amounts to multiplying the tree-level coefficient $`C_2(ϵ)`$ in the original expression with $`C_0(ϵ)`$. At $`q_0=(M_BM_D)^2`$, the difference is a $`5\%`$ effect. Potentially more important are higher-order power corrections scaling as $`(\mathrm{\Lambda }/\mu _c)^3`$. The operator matrix elements contributing at third order in the heavy-quark expansion can be identified , but little is known about their actual size. Naive dimensional analysis suggests that, with a typical hadronic scale $`\mathrm{\Lambda }0.5`$ GeV, a third-order power correction could be of order $`(\mathrm{\Lambda }/\mu _c)^30.09`$ for $`q_0^2=(M_BM_D)^2`$ and $`(\mathrm{\Lambda }/\mu _c)^30.22`$ for $`q_0^2=15`$ GeV<sup>2</sup>, but clearly these are rough estimates which must be taken with caution.
We will attempt to extract as much information about power corrections as possible from the formulae derived in Section 2 for the $`BXl\nu `$ decay rate and spectra in the presence of a lepton invariant mass cut. We start with the quantity $`F(ϵ)`$ itself, which as shown in (5) receives a moderate second-order power correction proportional to the parameter $`\lambda _20.12`$ GeV<sup>2</sup>. The dashed line in Figure 5 shows an estimate of the unknown $`(\mathrm{\Lambda }/\mu _c)^3`$ correction, obtained by raising this second-order term to the power 3/2. To address the question to what extent this is a conservative estimate of a “generic” higher-order correction, we focus on the differential spectrum in (10) and on the contour representation in (12). The function $`F(ϵ)`$ is obtained from these results by performing integrals over $`\xi `$ or over the contour in the complex plane, respectively. However, the differential distributions contain additional information about power corrections, which is not seen after the integrations are performed. We first discuss the case of the contour integral in (12), taking the point of view that except for the region of small $`\phi `$ the magnitude of $`t(e^{i\phi },ϵ)`$ can be used to estimate the “generic” size of corrections to the heavy-quark limit. This is so because on any point on the circle far away from the real, positive $`\widehat{p}^2`$ axis, the function $`t(e^{i\phi },ϵ)`$ admits an operator product expansion in a series of local operators, whose contributions scale like powers of $`\mathrm{\Lambda }/\mu _c`$. The average over the circle determines the corrections to the function $`F(q_0^2)`$. However, the finer details of the distribution on the circle become relevant, e.g., when in a real experiment events with different hadronic masses and energies are weighted by different efficiencies. In other words, the various terms proportional to $`\lambda _1/\mu _c^2`$ and $`\lambda _2/\mu _c^2`$ in (13) are as valid as estimate of a second-order power correction as is the $`\lambda _2/\mu _c^2`$ term in (5). Specifically, we calculate the average value of the modulus of the power corrections in (13) on the circle in the complex plane, and then we raise this number to the power 3/2 to obtain an estimate of a “generic” $`(\mathrm{\Lambda }/\mu _c)^3`$ correction. For our numerical analysis we use $`\lambda _1=(0.30\pm 0.15)`$ GeV<sup>2</sup>, which is in the ball park of recent determinations . The result is shown by the dark band in Figure 5, whose width reflects the sensitivity to the value of $`\lambda _1`$. If we were to consider larger values of $`|\lambda _1|`$ the upper limit of the band would increase. Another estimate of power corrections can be obtained from the coefficients of the various $`\delta `$-function terms in (10). If we take one third of the geometric average of the three coefficients, and raise the result to the power 3/2, we obtain the light band shown in the figure.
The above analysis shows that, as expected, there is a large uncertainty in the estimate of higher-order power corrections. We do not claim that these corrections are likely to be as large as indicated by the upper limit of the light band in Figure 5, however, as we have shown this would indeed be possible without introducing any unnaturally large coefficients or parameters. Keeping this caveat in mind, we will from now on use the upper limit of the dark band in the figure as our estimate of third-order power corrections. Numerically, this estimate is close to the one obtained by BLL .
## 5 Phenomenological implications and summary
The proposal of BLL is to use the theoretical calculation of the fraction $`F(q_0^2)`$ to obtain a model-independent determination of the CKM matrix element $`|V_{ub}|`$ with controlled and small theoretical uncertainty . To this end, one uses the relation $`\mathrm{Br}(BX_ul\nu )|_{q^2>q_0^2}=F(q_0^2)\mathrm{\Gamma }(BX_ul\nu )\tau _B`$, where $`\tau _B`$ is the $`B`$-meson lifetime, and $`\mathrm{\Gamma }(BX_ul\nu )`$ is the total semileptonic decay rate into charmless final states. This rate can be calculated with high accuracy in terms of a low-scale subtracted $`b`$-quark mass, including perturbative corrections of order $`[\alpha _s(m_b)]^2`$ and power corrections of order $`(\mathrm{\Lambda }/m_b)^2`$. For our purposes, we use the PS mass defined at the scale $`\mu _2=2`$ GeV. Then the expression for the total rate is
$$\mathrm{\Gamma }(BX_ul\nu )=\frac{G_F^2|V_{ub}|^2[m_b^{\mathrm{PS}}(\mu _2)]^5}{192\pi ^3}(1+\delta _{\mathrm{pert}}+\delta _{\mathrm{power}}),$$
(25)
where $`\delta _{\mathrm{pert}}0.04`$ at two-loop order, and $`\delta _{\mathrm{power}}=\frac{\lambda _19\lambda _2}{2m_b^2}0.03`$. The small uncertainties in these two quantities are negligible for our numerical analysis below. Using these results, we obtain the master formula
$$|V_{ub}|=2.96\times 10^3\left[\frac{\mathrm{Br}(BXl\nu )|_{q^2>q_0^2}}{10^3F^{}(q_0^2)}\frac{1.6\mathrm{ps}}{\tau _B}\right]^{1/2},$$
(26)
where all theoretical uncertainties are contained in the function
$$F^{}(q_0^2)=\left(\frac{m_b^{\mathrm{PS}}(\mu _2)}{4.59\mathrm{GeV}}\right)^5F(q_0^2).$$
(27)
The mass dependence due to factor $`[m_b^{\mathrm{PS}}(\mu _2)]^5`$ from the total decay rate is positively correlated with the mass dependence of the function $`F(q_0^2)`$. As a result, our predictions for the function $`F^{}(q_0^2)`$ become extremely sensitive to the value of the $`b`$-quark mass. For practical purposes, this dependence can be parameterized as
$$F^{}(q_0^2)\left(\frac{m_b^{\mathrm{PS}}(\mu _2)}{4.59\mathrm{GeV}}\right)^{\mathrm{\Delta }(q_0^2)},$$
(28)
where
$$\mathrm{\Delta }(q_0^2)10+\frac{q_0^2(M_BM_D)^2}{1\text{GeV}^2}.$$
(29)
In Table 1, we show our final results for the quantity $`F^{}(q_0^2)`$ and its theoretical uncertainties (as estimated above) for some representative values of $`q_0^2`$. For comparison, we note that BLL obtained the values $`F^{}((M_BM_D)^2)=0.169\pm 0.016`$ and $`F^{}(15\mathrm{GeV}^2)=0.061\pm 0.013`$, where the dominant theoretical error was assumed to be due to higher-order power corrections. Our central values are significantly higher because of the larger perturbative correction obtained after RG improvement. Note that our error estimates are about 2.5 times as large as those quoted by BLL. The difference between the central values of the two calculations is about $`1\sigma `$ of our errors, and about $`2\sigma `$ of their errors.
In Figure 6 we show a graphical representation of the fraction $`F^{}(q_0^2)`$ and its total theoretical uncertainty. This result, together with the master formula (26), provides the theoretical basis for the determination of $`|V_{ub}|`$. The right-hand plot in the figure shows the fractional theoretical uncertainty in the result for $`|V_{ub}|`$. Although our error estimates are more pessimistic than those presented by BLL, we still conclude that their method provides a very promising route for a precise determination of $`|V_{ub}|`$. For a realistic cut on the lepton invariant mass in the vicinity of $`q_0^212.5`$ GeV<sup>2</sup>, which is about 1 GeV<sup>2</sup> above the optimal value, the theoretical uncertainty in $`|V_{ub}|`$ is close to 10%. A determination with such an accuracy would be a significant improvement with respect to the present knowledge of this important parameter. We believe it would also be more reliable than a future determination obtained by combining the partial decay rates in the endpoint regions of $`BX_s\gamma `$ and $`BX_ul\nu `$ decay spectra , which is limited by uncontrollable power corrections of first order in $`\mathrm{\Lambda }/m_b`$ that violate the factorization of soft and collinear singularities.
According to Table 1, the dominant sources of theoretical uncertainty in the extraction of $`|V_{ub}|`$ are associated with the sensitivity to the value of the $`b`$-quark mass and with unknown, higher-order power corrections. Whereas it is not obvious how one should obtain a reliable value for the power corrections, the precision in the value of the $`b`$-quark mass can presumably be improved by reducing the theoretical uncertainties in the analysis of $`(b\overline{b})`$ bound states.<sup>2</sup><sup>2</sup>2We stress, however, that using the so-called Upsilon mass defined as one half of the mass of the $`\mathrm{{\rm Y}}(1S)`$ bottomonium state does not eliminate the uncertainty associated with the variation of the $`b`$-quark mass. As discussed in , this choice obscures the presence of an unknown nonperturbative contribution to the bound-state mass, which is neglected in the perturbative expression of the $`B`$-meson decay rate in terms of the Upsilon mass. In other words, in such a scheme the value of $`m_b`$ is known (by definition) with very high precision, but for consistency the uncertainty shown in the third column in Table 1 must then be added to the other theoretical uncertainties. In addition, it would be possible to reduce the perturbative uncertainty in the calculation in two ways, by calculating the exact $`O(\alpha _s^2)`$ corrections to the fraction $`F(q_0^2)`$ (the two-loop corrections to the total decay rate are known ), and by computing the two-loop anomalous dimensions of the operators contributing at $`O(ϵ)`$ in the hybrid expansion. Both calculations are technically feasible and should be done.
In summary, we have analyzed the structure of the heavy-quark expansion for the inclusive, semileptonic $`BXl\nu `$ decay rate with a lepton invariant mass cut $`q^2q_0^2`$. This expansion is characterized by a hadronic scale $`\mu _c<m_c`$ determined by the value of $`q_0^2`$. Because $`m_b\mu _c\mathrm{\Lambda }`$, the heavy-quark expansion can be organized as a combined (hybrid) expansion in two small mass ratios. The physics associated with the two large scales $`m_b`$ and $`\mu _c`$ is disentangled using the HQET, whereas the physics on the scale $`\mu _c`$ can be separated from long-distance physics associated with $`\mathrm{\Lambda }`$ utilizing an operator product expansion. We have used this formalism to obtain a RG-improved expression for the leading short-distance coefficient in the heavy-quark expansion at NLO. The summation of large logarithms in the hybrid expansion turns out to be important and strongly enhances the overall size of the perturbative correction. We have also emphasized that in order to obtain a stable perturbative prediction it is important to eliminate the $`b`$-quark pole mass in favor of a low-scale subtracted quark mass, such as the PS mass. Finally, we have presented several independent estimates of higher-order power corrections in the heavy-quark expansion, which at present do not permit a rigorous treatment. We find that with realistic values of the lepton invariant mass cut the overall theoretical uncertainty in the extraction of $`|V_{ub}|`$ is about 10%, which is larger than previously estimated but still significantly less than the current uncertainty in this parameter.
###### Acknowledgments.
I am grateful to Martin Beneke, Alex Kagan, Zoltan Ligeti and Mark Wise for useful discussions. This work was supported in part by the National Science Foundation.
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# Reconstructing the ZOA from Galaxy Peculiar Velocities
## 1. Introduction
Extinction due to the galactic plane obscures about $`25\%`$ of the optically visible universe. In order to account for the Local Group motion relative to the Cosmic Microwave Background (CMB), the flow of galaxies in the Great-Attractor (hereafter, GA) area and other similar phenomena, the full distribution of matter, especially in the local universe, is essential. Direct measurement of the distribution of matter/galaxies requires extensive, tedious and dedicated observational programs in all the available electromagnetic wave-bands (see review by Kraan-Korteweg & Lahav 2000 and references therein).
However, a complementary approach for studying the universe behind the Milky Way is to use the available dynamical data, e.g., galaxy peculiar velocity catalogs, together with statistical reconstruction methods, e.g., Wiener filtering (WF), in order to uncover the mass density distribution, with resolution scale $`\stackrel{>}{}500\mathrm{km}\mathrm{s}^1`$, hence, singling out the dynamically most significant structures.
Peculiar velocities of galaxies enable a direct and reliable measurement of the underlying mass distribution, under the natural assumption that galaxies are unbiased tracers of the large-scale, gravitationally induced, velocity field. Furthermore, since peculiar velocities are non-local and have contributions from different scales and different regions, analysis of the peculiar velocity field provides information on regions not covered by the data, e.g., the ZOA (Kolatt et al. 1995; Zaroubi et al. 1999), and on scales larger than the sampled regions (Hoffman et al. 2000).
Kolatt et al. (1995) were the first to attempt to reconstruct the ZOA from galaxy peculiar velocity data, where they used the POTENT method (Bertschinger & Dekel 1989; Dekel 1994) to reconstruct the mass-density distribution, within a sphere of radius $`8000\mathrm{km}\mathrm{s}^1`$, from the Mark III galaxy peculiar velocity catalog (Willick et al. 1997) with $`1200\mathrm{km}\mathrm{s}^1`$ resolution. Their study has resulted in predicting that the mass distribution of the Great Attractor peaks precisely at the center of the ZOA at a distance of $`4500\mathrm{km}\mathrm{s}^1`$. Zaroubi et al. (1999) have used the same peculiar velocity catalog to Wiener reconstruct the mass density distribution within $`8000\mathrm{km}\mathrm{s}^1`$ sphere. Their main conclusions were consistent with those of Kolatt et al. (1995).
The common use of WF is for straightforward noise suppression, but it can be easily generalized to achieve two further goals: to reconstruct the density field from the observed radial velocities and to interpolate or extrapolate the reconstruction to regions of poor sampling (for review, see Zaroubi et al. 1995 & Hoffman 2000). The later aspect is of special use for the reconstruction of the ZOA. In this work the WF is used to reconstruct the mass density distribution within $`8000\mathrm{km}\mathrm{s}^1`$ sphere from the largest galaxy peculiar velocity catalog yet. This catalog is a combination of the SFI and the ENEAR peculiar galaxy catalog.
The WF approach has been already applied to the IRAS two and three-dimensional galaxy distribution (Lahav et al. 1994), the IRAS three-dimensional redshift distortions (Fisher et al. 1995; Webster et al. 1997), the COBE/DMR cosmic microwave background mapping (Bunn et al. 1994), and to galaxy peculiar velocity catalogs of Mark III and ENEAR (Zaroubi et al. 1999; 2000a).
The outline of this paper is as follows. In § 2 we briefly describe the peculiar velocity data used in the present analysis. The method of Wiener reconstruction from peculiar velocity data is introduced in § 3, and the results of its application to the SEcat data set are presented in § 4. The paper concludes with a general discussion (§ 5).
## 2. The Data Sets
The ENEAR catalog have been extracted from the all-sky ENEAR redshift survey comprising about 1600 galaxies. Individual galaxy distances were estimated from a direct $`D_n\sigma `$ template relation derived by combining all the available cluster data, corrected for incompleteness and associated diameter-bias. From the observed scatter of the template relation the estimated fractional error in the inferred distance of a galaxy is $`\mathrm{\Delta }0.19`$, nearly independent of the velocity dispersion. An objective grouping procedure has been applied to the data in order to lower the inhomogeneous Malmquist bias before correction and to avoid strong non-linear effects (in particular large velocities of galaxies in clusters). The final catalog consists of about $`750`$ objects.
The SFI catalog of peculiar velocities of galaxies (Giovanelli et al. 1999), contains about 1300 field spiral galaxies with Tully-Fisher distances. After the grouping procedure the final dataset consists of distances, radial peculiar velocities and errors for $`1250`$ objects, ranging from individual field galaxies to rich clusters.
The combined catalog, SEcat, consists of $`2000`$ objects, uniformly covering, apart from the ZOA region, the local universe up to distance of $`6000\mathrm{km}\mathrm{s}^1`$. The error in the distance of the objects measured with $`D_n\sigma `$, namely objects from the ENEAR catalog, are assumed to have two contributions, the first is the usual $`D_n\sigma `$ distance proportional errors. The second is a constant error of $`250\mathrm{km}\mathrm{s}^1`$ that accounts for the non-linear velocities of galaxies in the high density environment in which early-type galaxies preferentially reside.
The inferred distances are corrected for the homogeneous and inhomogeneous Malmquist bias (for details see Freudling et al. 1999; da Costa et al. 2000b). The latter was estimated using the PSCz density field (Branchini et al. 1999), corrected for the effects of peculiar velocities, using the expressions given by Willick et al. (1997). In this calculation, a correction for redshift limit of the sample is included.
Finally the results are compared with the mass-density reconstruction from the Mark III catalog (Willick et al. 1997). This catalog, consists of more than 3400 galaxies, has been compiled from several data sets of spirals and elliptical/S0 galaxies with distances inferred by the forward Tully-Fisher and $`D_n\sigma `$ distance indicators. These data were re-calibrated and self-consistently put together as a homogeneous catalog for velocity analysis. The catalog provides radial velocities and inferred distances with errors on the order of $`1721\%`$ of the distance per galaxy. After grouping, the catalog contains $`1200`$ objects. The sampling covers the whole sky outside the ZOA, but with an anisotropic and non-uniform density that is a strong function of distance. The good sampling typically ranges out to $`6000\mathrm{km}\mathrm{s}^1`$ but it may be limited to only $`4000\mathrm{km}\mathrm{s}^1`$ in some directions or extend beyond $`8000\mathrm{km}\mathrm{s}^1`$ in other directions. The inhomogeniety of the Mark III sampling together with the complicated calibration proceedure employed to obtain the final catalog led many authors to question its reliability (e.g., Davis et al. 1996).
## 3. Wiener Filter
Here we limit the description of the WF to the actual application of the method to the case of radial velocity data. The data for the WF analysis are given as a set of observed radial peculiar velocities $`u_i^o`$ sampled at positions $`𝐫_i`$ with estimated errors $`ϵ_i`$ that are assumed to be uncorrelated. The peculiar velocities are assumed to be corrected for systematic errors such as Malmquist bias. The observed velocities are thus related to the true underlying velocity field $`𝐯(𝐫)`$, or its radial component $`u_i`$ at $`𝐫_i`$, via
$$u_i^o=𝐯(𝐫_i)\widehat{𝐫}_i+ϵ_iu_i+ϵ_i.$$
(1)
We assume that the peculiar velocity field $`𝐯(𝐫)`$ and the density fluctuation field $`\delta (𝐫)`$ are related via linear gravitational-instability theory, $`\delta =f(\mathrm{\Omega })^1𝐯`$, where $`f(\mathrm{\Omega })\mathrm{\Omega }^{0.6}`$ and $`\mathrm{\Omega }`$ is the mean universal density parameter. Under the assumption of a specific theoretical prior for the power spectrum $`P(k)`$ of the underlying density field, we can write the WF minimum-variance estimator of the fields as
$$𝐯^{\mathrm{WF}}(𝐫)=<𝐯(𝐫)u_i^o><u_i^ou_j^o>^1u_j^o$$
(2)
and
$$\delta ^{\mathrm{WF}}(𝐫)=<\delta (𝐫)u_i^o><u_i^ou_j^o>^1u_j^o.$$
(3)
In these equations $`<\mathrm{}>`$ denotes an ensemble average. The assumption that linear theory is valid on all scales enables us to estimate, given the power spectrum, the ensemble average quantities appearing in Eqs. (2) & (3) . The reader is referred to Zaroubi et al. (1999) for the explicit mathematical formulae used in the calculation. We choose to reconstruct the density field with a finite Gaussian smoothing of radius $`900\mathrm{km}\mathrm{s}^1`$.
## 4. Results
First we compare the density reconstruction from the Mark III, SFI and ENEAR catalogs. Figure 1 shows the reconstructed mass-density distribution for each catalog, smoothed with a $`900\mathrm{km}\mathrm{s}^1`$ Gaussian, on a spherical shell at $`4000\mathrm{km}\mathrm{s}^1`$ distance. The assumed power spectrum used in the reconstruction from Mark III, SFI & ENEAR has been determined through maximum likelihood analysis by Zaroubi et al. (1997), Freudling et al. (1999) and Zaroubi et al. (2000a), respectively.
In the three maps the existence of the GA supercluster on the left and the Perseus-Pisces (P-P) supercluster to the right, $`(l,b)(135{}_{}{}^{},30{}_{}{}^{})`$ is evident. However, the ENEAR map, relative to the two others, shows a more localized GA and P-P. In fact, since the ENEAR catalog measures velocities of early-type galaxies preferentially residing in high density environments, this difference is expected. Conversely, smaller overdensities for the GA and P-P shown in the SFI density reconstruction map is due to the tendency of spiral galaxies to reside in the field. The insufficient sampling of the P-P supercluster in the Mark III catalog renders its recovered density smaller than that of the GA.
The big structure, centered at $`(l,b)(200{}_{}{}^{},30{}_{}{}^{})`$, appearing in the SFI density reconstruction and in the correspondent SEcat spherical shell (see Figure 3), does not have a counterpart in the ENEAR density map and has much lower density in the Mark III reconstruction. Comparison with the IRAS 1.2-Jy redshift survey density reconstruction (Webster et al. 1997) shows that the peak location of this structure coincides with the position of the massive cluster N1600. The IRAS 1.2-Jy $`4000\mathrm{km}\mathrm{s}^1`$ shell further shows the existence of several other clusters, i.e., Cancer, Camelopardalis, $`𝐂_\beta `$, $`𝐂_\gamma `$, $`𝐂_\delta `$ and P-P that can account for this huge concentration seen extended from $`l180{}_{}{}^{}220^{}`$ and centered around the ZOA.
Obviously, the ENEAR and SFI catalogs complement, therefore we combined them to one catalog, SEcat. Figure 2 shows the maps of the density field, recovered rfom the SEcat catalog, in 4 different slices using a Gaussian smoothing of $`900\mathrm{km}\mathrm{s}^1`$. In the Supergalactic plane slice (upper left) the main features of our local universe can be easily identified, including the GA and the P-P superclusters at the left and right parts of the map respectively; the Local supercluster appears at the center of the map. The SGY=0 slice (upper right) coincides roughly with the plane obscured by the ZOA. Another two slices are at SGZ$`=\pm 4000\mathrm{km}\mathrm{s}^1`$ are also shown.
Figure 3 shows the Aitoff projection of the WF SEcat reconstructed density in Galactic coordinates, evaluated across shells at various distances. The structures in these maps match closely those seen in similar reconstruction from the IRAS 1.2-Jy redshift catalog (Webster et al. 1997). The $`4000\mathrm{km}\mathrm{s}^1`$ shell has been discussed earlier, the other shells will be discussed in detail elsewhere (Zaroubi et al. 2000b)
The velocity field along the Supergalactic plane is presented in Fig. 4, showing the existence of two convergence regions which roughly coincide with the locations of the GA and PP.
## 5. Discussion
In spite of the high level of extinction due to the galactic plane, galaxy peculiar velocities together with statistical reconstruction techniques, e.g., WF, present a very useful tool for mapping the ZOA. This approach complements the very challenging task of directly mapping the universe behind the Milky Way.
In this contribution we have showed that the WF method could indeed, within the resolution limit, faithfully reconstruct the nearby universe including regions masked by the Galactic plane. Several issues still need to be addressed as the details of the reconstructed maps can vary from catalog to catalog depending on the distance indicator, i.e., TF vs. $`D_n\sigma `$ the sampling, noise properties and systematic differences, e.g., calibration.
In attempting to combine the ENEAR and SFI data-sets, one needs to ensure their consistency. Indeed various indications, e.g., calibration, zero point, measured bulk flow (da Costa et al. 2000a & 2000b) support the assumed compatibility of ENEAR and SFI, enabling their combination to one new catalog of spiral and elliptical galaxies, SEcat. These issues will be explored in detail in a forthcoming paper (Zaroubi et al. 2000b).
### Acknowledgments.
I would like to thank M. Bernardi, L.N. da Costa and my long term collaborator Y. Hoffman for their contribution to this work and C. Cress for her helpful comments on the manuscrip. The contribution of the ENEAR team is gratefully acknowledged. The financial support of the Deutsche Forschungsgemeinschaft (DFG) is acknowledged.
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# Squeezed Condensates submitted to Phys. Rev. A
## 1 Introduction
Like electromagnetic waves, atomic matter waves can to a large extent be manipulated in space. Interference phenomena have been observed, and a large number of high precision atom optics measurements have been carried out . One of the potential future applications of atomic Bose-Einstein condensates is as coherent matter wave sources for interferometric time and frequency standards, detection of inertial effects, and a host of related technological tasks. The spatial and temporal coherence of condensates has been verified, and coherent amplification of matter waves has been demonstrated to establish the close analogy with laser and maser sources of light.
For high precision purposes, however, it was realized long time ago, that so-called non-classical states of light may be more useful than classical field states . It is therefore natural to consider the production of such ’non-classical’ states of atoms as well. The term ’non-classical’ is ambiguous here, since what is very non-classical for light, e.g. a number state, can seem perfectly classical for atoms, and we shall instead use the term quantum correlated states. Due to the collisional interaction between atoms in a Bose condensate, this system contains a non-linearity, in equivalence with the Kerr-effect for light, and this collisional interaction has already been observed in a matter wave analog of four-wave mixing in non-linear optics. There have been several proposals to utilize collisional effects to produce certain quantum correlated states, such as Schrödinger cat-like states , and to observe the effect of the non-linearity on the dynamics of the condensate . It is clear from these studies that, even though the collisional interaction can be controlled experimentally it is not easy to control precisely the quantum features of interest; the problem is a genuine multi-mode one due to the spatial degrees of freedom, and one has to pay attention to the role of non-condensed atoms, as well.
In this Paper, we suggest to follow instead the procedure of quantum optics for production of squeezed light: we suggest to implement an ’atomic OPO’. The optical parametric oscillator (OPO) is a device where light is down converted, so that a single pump photon at frequency $`2\omega `$ is converted into two photons, each of frequency $`\omega `$ (degenerate case). The process is the inverse of second harmonic generation (SHG), and in practical experiments, one often sees a strong field of frequency $`\omega `$, which is first frequency doubled in the SHG crystal, and the high frequency field is subsequently down-converted in a similar non-linear crystal, now working as an OPO.
The matter wave analogue of the SHG process is one in which the atoms are combined by photo-association into di-atomic molecules. This process has been analyzed theoretically , and it was recently demonstrated experimentally, that part of an atomic condensate may be converted by stimulated Raman transitions into molecules this way . We suggest to follow the analogy with the OPO quite directly, i.e., to apply laser fields to drive the Raman transition from the molecular state back to free atoms (photo-dissociation), and we note, that the quantum correlations appear due to the fact that only even number states for the atomic component will be present in the sample, because atoms are created in pairs.
The dynamics of coupled atomic and molecular degenerate gasses has been studied in some detail with particular focus on the spatial and temporal dependence of the mean field dynamics, e.g., the appearance of solitons . We focus on the “quantum optics” aspects, i.e., the atom number distribution, and the consequences for atom counting experiments.
The organization of this paper is as follows: In Sec. 2, we present our proposal, and we derive two theoretical methods, one approximate and one exact, to analyze the dynamics of the system. In Sec. 3 we present numerical results for various relevant quantities. In Sec. 4 we analyze the use of our quantum correlated atoms as squeezed input in two-state Rabi-oscillations with a large, “classical”, condensate. In Sec. 5, we discuss the results and we briefly indicate some alternative ideas for the production and application of quantum correlated condensates.
## 2 The model
### 2.1 Trapped cold atoms in second quantization
It is convenient to describe dilute gasses of cold atoms in a trap in the second quantized formalism. In this description the Hamiltonian for the system at low temperatures is usually taken to be:
$$H_0^{3\mathrm{D}}=𝑑\stackrel{}{r}\left\{\widehat{\psi }^{}(\stackrel{}{r})h_1\widehat{\psi }(\stackrel{}{r})+\frac{g}{2}\widehat{\psi }^{}(\stackrel{}{r})\widehat{\psi }^{}(\stackrel{}{r})\widehat{\psi }(\stackrel{}{r})\widehat{\psi }(\stackrel{}{r})\right\}.$$
(1)
The atomic field operators $`\widehat{\psi }(\stackrel{}{r})`$ and $`\widehat{\psi }^{}(\stackrel{}{r})`$ annihilate and create an atom at position $`\stackrel{}{r}`$ and they satisfy the equal time commutation relations
$$[\widehat{\psi }(\stackrel{}{r}),\widehat{\psi }^{}(\stackrel{}{r}^{})]=\delta ^3(\stackrel{}{r}\stackrel{}{r}^{})$$
(2)
The single particle Hamiltonian $`h_1`$ is the one appropriate for a single atom of mass $`m`$ in the external trapping potential $`V_{ext}`$
$$h_1=\frac{\mathrm{}^2}{2m}\stackrel{}{}^2+V_{ext}(\stackrel{}{r}).$$
(3)
The $`g`$ term in Eq.(1) describes the collisional two body interaction of the atoms. The simple contact form is an approximation appropriate for cold, dilute gases . The value of the strength parameter is $`g=4\pi \mathrm{}a_s/m`$ where $`a_s`$ is the s-wave scattering length.
In most experiments the trapping potential can be approximated well by a three dimensional harmonic oscillator characterized by three frequencies $`\omega _x`$, $`\omega _y`$ and $`\omega _z`$. If one of these frequencies is significantly smaller than the other two the atoms will form an elongated cloud and we describe such an effectively one dimensional system by a single frequency $`\omega `$ .
We shall work with dimensionless equations and choose to measure time in units of $`\omega ^1`$, lengths in units of $`a_0=\sqrt{\mathrm{}/\omega m}`$, and energy in units of $`\mathrm{}\omega `$. We are then left with
$$H_0=𝑑x\left\{\widehat{\psi }^{}(x)h_1\widehat{\psi }(x)+\frac{g}{2}\widehat{\psi }^{}(x)\widehat{\psi }^{}(x)\widehat{\psi }(x)\widehat{\psi }(x)\right\}$$
(4)
where
$$h_1=\frac{1}{2}\frac{^2}{x^2}+\frac{1}{2}x^2.$$
(5)
and where $`g=4\pi a_s/a_0`$ is a dimensionally correct, one dimensional interaction strength .
### 2.2 Photodissociation from a molecular condensate
Starting from a molecular condensate (prepared by photo-association or by other means) we can imagine to photo-dissociate or “down convert” the molecules to pairs of free atoms. A suitable Hamiltonian to describe photo-dissociation is
$$H_{PD}=\frac{1}{2}𝑑x\left\{b(x,x^{},t)\widehat{\psi }^{}(x)\widehat{\psi }^{}(x^{})+b^{}(x,x^{},t)\widehat{\psi }(x)\widehat{\psi }(x^{})\right\}.$$
(6)
This Hamiltonian clearly creates and annihilates atoms in pairs. A full description should also include the molecules such that each atomic pair creation would be accompanied by the annihilation of a molecule. Here we will assume the molecular condensate to be large and the number of molecules removed to be small. In that case it is reasonable to describe the molecular condensate by a time-independent c-number field. The function $`b`$ in Eq.(6) is the product of this field and of the coupling to a position dependent laser field which is also assumed to be classical. To represent the position dependence of the molecular condensate and of the laser field, the relative wavefunction of the pair of atoms, and the time dependence of the laser field, we use the ansatz
$$b(x,x^{},t)=\frac{B}{2\pi \sigma _r\sigma _{cm}}\mathrm{exp}\left(2i\mathrm{\Delta }t\right)\mathrm{exp}\left(\frac{1}{2}\frac{(xx^{})^2}{\sigma _r^2}\right)\mathrm{exp}\left(\frac{1}{2}\frac{(\frac{1}{2}(x+x^{}))^2}{\sigma _{cm}^2}\right)$$
(7)
where $`\sigma _r\sigma _{cm}a_0`$.
### 2.3 Operator equations of motion
A successful tool for the description of a Bose condensate is the Gross-Pitaevskii equation (GPE). It can be obtained from the Heisenberg equation of motion for the atomic field $`\widehat{\psi }(x,t)`$ by taking average values and by replacing the mean value of an operator product by the product of mean values. In our case the Heisenberg equation reads:
$$i\frac{\widehat{\psi }(x,t)}{t}=\left(\frac{1}{2}\frac{}{x}+\frac{1}{2}x^2+g\widehat{\psi }^{}(x,t)\widehat{\psi }(x,t)\right)\widehat{\psi }(x,t)+𝑑x^{}b(x,x^{},t)\widehat{\psi }^{}(x^{},t).$$
(8)
The last term in this equation is due to the exchange of atom pairs with the molecular condensate via photo-dissociation.
It is easy to see that using the average of Eq.(8) has shortcomings when we try to describe photo-dissociation: If we start from the atomic vacuum defined by
$`\widehat{\psi }(x,0)|0`$ $`=`$ $`0`$ (9)
the incoupling term on the right hand side has vanishing average value. Therefore $`\widehat{\psi }(x,t)`$ will stay zero also at all later times and no useful information can be extracted.
To gain knowledge of the state created we therefore proceed to study expressions quadratic in the field operators. To shorten notation we define
$`\widehat{R}(x,y,t)`$ $``$ $`\widehat{\psi }^{}(x,t)\widehat{\psi }(y,t)`$ (10)
$`\widehat{S}(x,y,t)`$ $``$ $`\widehat{\psi }(x,t)\widehat{\psi }(y,t).`$ (11)
For these operators we get the following Heisenberg equations of motion:
$`i{\displaystyle \frac{\widehat{R}(x,y,t)}{t}}`$ $`=`$ $`\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{^2}{x^2}}{\displaystyle \frac{1}{2}}x^2{\displaystyle \frac{1}{2}}{\displaystyle \frac{^2}{y^2}}+{\displaystyle \frac{1}{2}}y^2\right)\widehat{R}(x,y,t)`$ (12)
$`+g\left(\widehat{R}(y,y,t)\widehat{R}(x,x,t)\right)\widehat{R}(x,y,t)`$
$`+{\displaystyle 𝑑z\left\{b(z,y,t)\widehat{S}^{}(x,z,t)b^{}(x,z,t)\widehat{S}(z,y,t)\right\}}`$
and
$`i{\displaystyle \frac{\widehat{S}(x,y,t)}{t}}`$ $`=`$ $`\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{^2}{x^2}}+{\displaystyle \frac{1}{2}}x^2{\displaystyle \frac{1}{2}}{\displaystyle \frac{^2}{y^2}}+{\displaystyle \frac{1}{2}}y^2\right)\widehat{S}(x,y,t)`$ (13)
$`+g\left(\delta (xy)+\widehat{R}(x,x,t)+\widehat{R}(y,y,t)\right)\widehat{S}(x,y,t)`$
$`+{\displaystyle 𝑑z\left\{b(z,y,t)\widehat{R}(z,x)+b(x,z,t)\widehat{R}(z,y)\right\}}+b(x,y,t).`$
Note that $`b(x,y,t)`$ now appears as an inhomogeneous source term in the $`\widehat{S}`$ equation. This guarantees a nontrivial behaviour when we take averages. Note also that when taking averages we have a problem with the interaction terms which will couple the second order expectations to fourth order expectations. These fourth order terms have to be factorized in some approximate way to obtain a closed set of equations.
### 2.4 c-number equations
#### 2.4.1 Exact equations for $`g=0`$
When $`g=0`$, Eq.(12) and Eq.(13) can be reduced to two coupled linear equations for the moments
$$R(x,y,t)\widehat{R}(x,y,t),S(x,y,t)\widehat{S}(x,y,t).$$
(14)
They read:
$`i{\displaystyle \frac{R(x,y,t)}{t}}`$ $`=`$ $`\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{^2}{x^2}}{\displaystyle \frac{1}{2}}x^2{\displaystyle \frac{1}{2}}{\displaystyle \frac{^2}{y^2}}+{\displaystyle \frac{1}{2}}y^2\right)R(x,y,t)`$ (15)
$`+{\displaystyle 𝑑z\left\{b(z,y,t)S^{}(x,z,t)b^{}(x,z,t)S(z,y,t)\right\}}`$
$`i{\displaystyle \frac{S(x,y,t)}{t}}`$ $`=`$ $`\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{^2}{x^2}}+{\displaystyle \frac{1}{2}}x^2{\displaystyle \frac{1}{2}}{\displaystyle \frac{^2}{y^2}}+{\displaystyle \frac{1}{2}}y^2\right)S(x,y,t)`$ (16)
$`+{\displaystyle 𝑑z\left\{b(z,y,t)R(z,x)+b(x,z,t)R(z,y)\right\}}`$
$`+b(x,y,t).`$
Moreover, $`R`$ and $`S`$ uniquely determine all higher order expectation values. This can be seen in a number of ways. One is to note that the Wigner distribution is a multi-dimensional gaussian distribution fully characterized by its second order moments . Another follows from the observation that when $`g=0`$, the Heisenberg equation of motion (8) implies that $`\widehat{\psi }(x,t)`$ can at all times be expressed as a linear combination of the initial values
$`\widehat{\psi }(x,t)`$ $`=`$ $`{\displaystyle 𝑑y\left\{f(x,y,t)\widehat{\psi }(y,0)+g(x,y,t)\widehat{\psi }^{}(y,0)\right\}}.`$ (17)
Eq.(17), its hermitian conjugate and the fact that our system starts in the vacuum state then suggest the following scheme for calculation of any operator product at arbritrary $`t`$: Use the commutation relation (2) to move all $`\widehat{\psi }(x,0)`$ to the right of any $`\widehat{\psi }^{}(y,0)`$ (normal ordering). Of all the terms produced in this process only the ones consisting entirely of c-numbers are nonzero as the vacuum expectation of any normal ordered product of operators vanishes in the vacuum state. To evalute the c-number terms we formally need to calculate integrals of products of the $`f`$ and $`g`$ functions of Eq.(17). It is, however, not difficult to see that these integrals factorize and that the factors are exactly the ones involved in calculating expectation values of products of only two field operators. The end result is that the average of any operator product is replaced by a sum of all possible factorizations into two-operator expectations
$`\widehat{\psi }^{}(x_1)\widehat{\psi }^{}(x_2)\widehat{\psi }(x_3)\widehat{\psi }(x_4)`$ $`=`$ $`\widehat{\psi }^{}(x_1)\widehat{\psi }^{}(x_2)\widehat{\psi }(x_3)\widehat{\psi }(x_4)`$ (18)
$`+\widehat{\psi }^{}(x_1)\widehat{\psi }(x_3)\widehat{\psi }^{}(x_2)\widehat{\psi }(x_4)`$
$`+\widehat{\psi }^{}(x_1)\widehat{\psi }(x_4)\widehat{\psi }^{}(x_2)\widehat{\psi }(x_3).`$
This is a simple version of Wick’s theorem .
#### 2.4.2 Approximate equations for $`g0`$
When $`g0`$ we have to include the interaction term of Eqs.(12,13) in Eqs.(15,16). Unfortunately, the decomposition Eq.(17) is no longer exact, and there is no simple way to reduce the mean values of four-operator products to products of $`R`$ and $`S`$. Rather than the simple replacement, e.g., $`\widehat{R}(y,y)\widehat{R}(x,y)R(y,y)R(x,y)`$, we choose to apply the Wick prescription as this is correct to lowest order. We then get
$`g\widehat{R}(y,y,t)\widehat{R}(x,y,t)\widehat{R}(x,x,t)\widehat{R}(x,y,t)`$ $``$ $`g\left(S^{}(x,y,t)S(y,y,t)S(x,y,t)S^{}(x,x,t)\right)`$ (19)
$`+2g\left(R(y,y,t)R(x,x,t)\right)R(x,y,t)`$
$`g\left(\widehat{R}(x,x,t)+\widehat{R}(y,y,t)+\delta (xy)\right)\widehat{S}(x,y,t)`$ $``$ $`2g\left(R(y,y,t)+R(x,x,t)\right)S(x,y,t)`$ (20)
$`+g\left(R(x,y,t)S(x,x,t)+R^{}(x,y,t)S(y,y,t)\right)`$
$`+g\delta (xy)S(x,y,t)`$
When these expressions are inserted into Eqs.(15,16), we arrive at the equations we want to solve numerically. We use a split-step approach where the kinetic energy is treated by a Fourier method. The remaining terms are dealt with by a fourth order Runge-Kutta scheme. In this one-dimensional problem the equations are quite manageable.
### 2.5 The positive $`P`$ pseudo-probability distribution
Pseudo-probablity distributions (PPD’s) are well-established tools in quantum mechanics. The most well known of the distributions are the Glauber-Sudarshan P-function and the Wigner function but especially in quantum optics a number of other distributions have also been useful. Common to all the PPD’s is that they provide the expectation values of properly ordered operator products as weighted c-number averages. The *positive $`P`$ distribution* ($`P_+`$) that we will be using here gives the expection values of normally ordered products by replacing $`\widehat{\psi }`$ by a c-number function $`\psi _1`$ and $`\widehat{\psi }^{}`$ by a c-number function $`\psi _2`$, e.g.:
$$\widehat{\psi }^{}(x,t)\widehat{\psi }(y,t)=d[\psi _1]d[\psi _2]\psi _2(x)\psi _1(y)P_+[\psi _1,\psi _2,t].$$
(21)
Note that $`P_+`$ is the joint distribution of two spatial functions and it is therefore an immensely complicated functional in general. It satisfies a multi-dimensional Fokker-Planck equation which, however, opens the door to a Monte Carlo sampling as we can translate the Fokker-Planck equation for the distribution to Langevin equations for stochastic realizations of $`\psi _1(x,t)`$ and $`\psi _2(x,t)`$. These equations resemble the GPE but they are coupled and they contain noise terms. A derivation of the equations without incoupling is given in (see also ) and the inclusion of incoupling is straightforward. In the notation of stochastic differential equations the equations can be written
$`id\psi _1(x)`$ $`=`$ $`\left(h_1\psi _1(x)+g\psi _2(x)\psi _1(x)\psi _1(x)\right)dt`$ (23)
$`+{\displaystyle b(x,x^{})\psi _2(x^{})𝑑x^{}}+dW_1(x)`$
$`id\psi _2(x)`$ $`=`$ $`\left(h_1\psi _2(x)+g\psi _1(x)\psi _2(x)\psi _2(x)\right)dt`$ (25)
$`+{\displaystyle b^{}(x,x^{})\psi _1(x^{})𝑑x^{}}+dW_2(x)`$
where $`h_1`$ is still defined in Eq.(3) and the noise terms are gaussian and given by
$`dW_{1,2}(x,t)`$ $`=`$ $`0`$ (26)
$`dW_1(x,t)dW_2(x^{},t^{})`$ $`=`$ $`0`$ (27)
$`dW_1(x,t)dW_1(x^{},t^{})`$ $`=`$ $`idt\left(b(x,x^{},t)+g\psi _1(x,t)\delta (xx^{})\right)\delta (tt^{})`$ (28)
$`dW_2(x,t)dW_2(x^{},t^{})`$ $`=`$ $`idt\left(b^{}(x,x^{},t)+g\psi _1(x,t)\delta (xx^{})\right)\delta (tt^{}).`$ (29)
By numerically simulating Eqs.(23,25) we are able to calculate expectation values of arbitrary, normally ordered field operator products with the only approximation that the results are subject to sampling errors due to the use of finite ensembles. We describe in the appendix our procedure to synthesize the noise $`dW`$ in our simulations.
The crucial drawback of the method is the well know sudden divergence in some of the unphysical moments of the $`P_+`$ distribution. Unphysical moments exist because the translation from operator products to products of $`\psi _1`$ and $`\psi _2`$ never involves $`\psi _1^{}`$ and $`\psi _2^{}`$. This leaves some room for $`P_+`$ to behave badly and unfortunately it exploits this freedom. In the wavefunction realizations, some wavefunctions diverge or they make very large excursions that are difficult to follow numerically and which make a devastating impact on the sampling error. The $`P_+`$ Monte Carlo method is therefore limited to short times where only few atoms have been created and nonlinear effects are still small. This suits our purpose, since we are only interested in short time dynamics, and we shall trust the result produced by the Langevin equations as long as none of the wavefunctions in the ensemble have escaped the region where we have confidence in our integration algorithm.
## 3 Results
In this section we show results for some of the quantities of interest that we are able to calculate in our model. Although the main new feature lie in the quantum correlations we first show a very classical quantity, namely the density profile. We then proceed to look at the eigenvalues of the one-particle density matrix. The largest of these eigenvalues defines the condensate fraction and the corresponding eigenvector is the condensate wavefunction. Finally we turn to a two-body quantity, the second order correlation function.
### 3.1 The density profile and the number of atoms
The atomic density is given by the diagonal elements of the one-body density operator in the position representation, that is
$$\rho (x,x)=\widehat{\psi }^{}(x)\widehat{\psi }(x)=R(x,x)=\overline{\psi _2(x)\psi _1(x)},$$
(30)
where the overbar in the last expression denotes the average over many realizations of the stochastic $`\psi _1(x,t)`$ and $`\psi _2(x,t)`$. In Fig. 1 is shown a typical plot of this profile at $`\omega t=2.4`$. It has the characteristic gaussian shape of the harmonic oscillator ground state and as we shall see in Sec. 3.2 a large fraction of the atoms indeed occupy a common wavefunction close to this state. The R&S equations (15,16) with the interaction terms (19,20) give results in excellent agreement with the $`P_+`$ simulations.
The total number of photo-dissociated atoms is obtained as the trace of the one-body density-operator or, according to Eq.(30), simply as
$$N=𝑑x\rho (x,x).$$
(31)
In Fig. 2 this number is shown as a function of time for $`g=0`$ and for $`g=0.01`$. The agreement between the R&S equations and the $`P_+`$ method is seen to be quite good.
### 3.2 The condensate fraction and wavefunction
By diagonalizing the one-body density matrix we obtain an orthonormal basis of single-particle states. The eigenvalues correspond to the populations of these states and in the case of a condensate one of these eigenvalues dominates, i.e., most particles occupy the same state. A dynamical picture of the condensation process is the Bose-enhancement of the scattering into the most occupied state. Here we expect a similar effect to take place. At first several states of the system are occupied by the atoms created. As the number of atoms grows the stimulated character of the creation becomes more important and a mode competition results in one mode being preferentially occupied.
In Fig.3 we show the condensate fraction, i.e. the ratio of the largest eigenvalue of the one-body density matrix to the sum of the eigenvalues for different values of the interaction strenght $`g`$. It is seen that as expected the condensate fraction is in general an increasing function of time. The effects of interactions are rather small at these low atom numbers. Note that unlike studies of stationary condensates at $`T=0`$, where interactions are responsible for the breakdown of a simple product state ansatz for the system and the existence of atoms outside the condensate, our incoupling by itself produces atoms both in the condensate and outside the condensate. In fact, our calculations show that the second-largest eigenvalue accounts for most of the atoms which are not in the condensate.
As for the condensate wavefunction we see an interesting phenomenon: Although the density profile associated with the condensed part of the one-body density matrix is close to that of the trap ground state, the condensate wavefunction is in fact not stationary. The atoms have condensed into a state more resembling a squeezed state<sup>1</sup><sup>1</sup>1This position-momentum squeezing should not be confused with the the atom-field squeezing discussed later. and if the incoupling is stopped the wavefunction widths show an oscillating behavior. In Fig. 4 we show $`\widehat{x}^2`$ and $`\widehat{p}^2`$ of the condensate wavefunction as a function of time. We see that at $`\omega t=2.4`$ when the incoupling is stopped, the wavefunction is too wide in momentum space as compared to the ground state of the trap ($`\widehat{p}^2>1/2`$).
One way to avoid this oscillation is to apply $`\delta `$-kick cooling to the system. This procedure is efficient if there is a linear correlation between position and momentum. In the original suggestion the correlation between position and momentum is brought about by free expansion, but an examination of the condensate wavefunction shows that we have a similar correlation here. The idea is to apply a tight, harmonic trapping potential for a short time interval. If this interval is so short that any changes in position can be ignored, the effect is simply a momentum kick also varying linearly with position. Ideally, this kick brings all the particles to rest. In Fig. 4 we demonstrate that the procedure is effective in our problem. The results are shown for $`g=0`$, but a similar reduction is achieved for non-vanishing $`g`$.
### 3.3 The second order correlation function $`g^{(2)}(x,y)`$
More detailed information about the quantum state of the system is desired and available, and a natural quantity to consider is the second order correlation function
$$g^{(2)}(x,y)\frac{\widehat{\psi }^{}(x)\widehat{\psi }^{}(y)\widehat{\psi }(y)\widehat{\psi }(x)}{\widehat{\psi }^{}(x)\widehat{\psi }(x)\widehat{\psi }^{}(y)\widehat{\psi }(y)}.$$
(32)
It measures the probability to find two atoms at positions $`x`$ and $`y`$ normalized to the single-particle densities at the two positions. For a thermal state $`g^{(2)}=2`$ while for a coherent state $`g^{(2)}=1`$.
As $`g^{(2)}`$ involves the expectation value of a product of four field operators we are faced with similar factorization problem as when we derived the R&S equations. Again we will resort to the Wick prescription although it should be realized that this is only exact for states obtained with $`g=0`$.
We have already evaluated the expectation in the enumerator of Eq.(32) in terms of $`R`$ and $`S`$ in the gaussian case. This was done in Eq.(18) and we get for the second order correlation function
$$g^{(2)}(x,y)=1+\frac{\left|R(x,y)\right|^2+\left|S(x,y)\right|^2}{R(x,x)R(y,y)}.$$
(33)
In contrast with the R&S equations the $`P_+`$ method has no problems handling expectation values like the enumerator of Eq.(32), and $`g^{(2)}`$ can be determined exactly up to sampling errors. At relatively short times and low atom numbers we have therefore an excellent tool to obtain exact results even for $`g0`$.
In Fig. 5 we show a plot of $`g^{(2)}(0,0)`$ as a function of time for various values of $`g`$. The central value slightly above 3 indicates a strong bunching effect where two atoms are more likely to be found close together than in a coherent or a thermal state. This result can be compared with the analytical expression for a single mode squeezed state, generated by the Hamiltonian $`\beta \left(\widehat{a}^2+\widehat{a}^{}^2\right)`$:
$$g^{(2)}=\frac{\widehat{a}^{}\widehat{a}^{}\widehat{a}\widehat{a}}{\widehat{a}^{}\widehat{a}^2}=3+\frac{1}{\widehat{a}^{}\widehat{a}}$$
(34)
In the figure we plot both results of the R&S equations using Eq.(33) and the exact $`P_+`$ results. Good agreement is found between the two approaches until $`\omega t2.5`$ and hereafter the R&S equations fail to capture a decrease in the value of $`g^{(2)}`$. This decrease indicates a threshold effect that we will discus in the next section.
### 3.4 Threshold effect
In the semi-classical treatments of the laser and of the parametric oscillator, one identifies a threshold in the stationary balance between gain and loss; the fields shift from fluctuations around zero to fluctuations around finite intensities . Above threshold these optical systems have smaller relative fluctuations of the intensity, and it is natural to expect a similar threshold behaviour in our model. There is a seemingly important difference between our model and the optical systems in the fact that we do not have an explicit dissipative mechanism. It has been known for a long time, and it has been demonstrated explicitly for a large number of physical systems, however, that quite generically, the interactions in many-body systems lead to ergodicity of eigenstates and a dynamical relaxation without coupling to an external bath. For a recent review, see .
The R&S equations with their underlying assumption of a gaussian Wigner distribution, centered around vanishing atomic field, are clearly unable to describe correctly the system around and above threshold, but the $`P_+`$ simulations are exact, and the discrepancy between the two methods is thus most likely explained by a threshold effect. To investigate more closely the threshold hypothesis, we show in Fig. 6 scatter plots of $`\psi _2(0)\psi _1(0)`$ at $`\omega t=2.4`$, $`3.0`$ and $`3.6`$ for a situation with $`g=0.02`$. From photodetection theory we can deduce the following expression for the atom number distribution
$$P_n(t)=\overline{\frac{[\psi _2(x,t)\psi _1(x,t)𝑑x]^n}{n!}\mathrm{exp}\left(\psi _2(x,t)\psi _1(x,t)𝑑x\right)}.$$
(35)
This expression, however, exhausts the statistical precision of the $`P_+`$ method, since it involves higher moments of the simulated amplitudes, which yield larger and larger fluctuations. Instead, we heuristically present histograms in the figures of how the real parts, $`\mathrm{Re}(\psi _2(0)\psi _1(0))`$ are distributed, and these histograms are in good qualitative agreement with our picture of a bifurcation of the solution when we reach threshold in the process. It is seen how the distribution at $`\omega t=2.4`$ is strongly peaked at zero with an exponential tail along the the real axis. At $`\omega t=3.0`$ this tail extends to larger values, and a shoulder around 120 atoms/$`a_0`$ starts to appear, and at $`\omega t=3.6`$ a second maximum has developed. This explains the lowering of $`g^{(2)}`$ seen in Fig. 5.
## 4 Application of a squeezed condensate
In this section we analyse a possible application of the state created in our model. We show how its peculiar statistical properties can be utilized to produce precise measurements. Our suggestion is in direct analogy with the use of a *squeezed vacuum* in quantum optics experiments with beam-splitters. When a beam in such an experiment is incident on a 50/50 beam-splitter and is split in two, the analysis of the noise properties (quantum fluctuations) of the two daughter beams depends crucially on realizing that the incoming beam is not only split at the beam-splitter, it is actually mixed with vacuum coming in from the back-side of the mirror. By replacing this vacuum by a squeezed state we may control the statistical properties of the daughter beams.
The matter wave analogue of the beam-splitter could in our case be a laser pulse which is able to coherently change the internal state of atoms. Such a pulse can drive each atom into a superposition of two internal states.
Suppose now that we have two internal states of the atoms, state $`a`$ and state $`b`$, and that we apply a $`\pi /2`$-pulse. We then have in the Heisenberg picture
$`\widehat{\psi }_a`$ $``$ $`\widehat{\psi }_a^{}={\displaystyle \frac{1}{\sqrt{2}}}\left(\widehat{\psi }_a+\widehat{\psi }_b\right)`$ (36)
$`\widehat{\psi }_a`$ $``$ $`\widehat{\psi }_b^{}={\displaystyle \frac{1}{\sqrt{2}}}\left(\widehat{\psi }_a\widehat{\psi }_b\right).`$ (37)
The total number operators of atoms in state $`a`$ and state $`b`$ after the pulse are thus given by
$`\widehat{N}_a^{}`$ $`=`$ $`{\displaystyle 𝑑x\widehat{\psi ^{}}_a^{}(x)\widehat{\psi }_a^{}(x)}`$ (38)
$`=`$ $`{\displaystyle \frac{1}{2}}\widehat{N}_a+{\displaystyle \frac{1}{2}}\widehat{N}_b+{\displaystyle \frac{1}{2}}{\displaystyle 𝑑x\left\{\widehat{\psi }_a^{}(x)\widehat{\psi }_b(x)+\widehat{\psi }_b^{}(x)\widehat{\psi }_a(x)\right\}}`$
$`\widehat{N}_b^{}`$ $`=`$ $`{\displaystyle 𝑑x\widehat{\psi ^{}}_b^{}(x)\widehat{\psi }_b^{}(x)}`$ (39)
$`=`$ $`{\displaystyle \frac{1}{2}}\widehat{N}_a+{\displaystyle \frac{1}{2}}\widehat{N}_b{\displaystyle \frac{1}{2}}{\displaystyle 𝑑x\left\{\widehat{\psi }_a^{}(x)\widehat{\psi }_b(x)+\widehat{\psi }_b^{}(x)\widehat{\psi }_a(x)\right\}}.`$
We will now concentrate on the *difference* in the number of atoms in the two states,
$`\widehat{N}_a^{}\widehat{N}_b^{}`$ $`=`$ $`{\displaystyle 𝑑x\left\{\widehat{\psi }_a^{}(x)\widehat{\psi }_b(x)+\widehat{\psi }_b^{}(x)\widehat{\psi }_a(x)\right\}}.`$ (40)
We will also let $`\widehat{\psi }_a`$ initially describe a large condensate while $`\widehat{\psi }_b`$ describes our photodissociated state. That is, we assume that the photodissociation is to the internal state $`b`$, while at the time of the $`\pi /2`$-pulse these atoms are overlapped with a large normal condensate in the internal state $`a`$. The large condensate is assumed to be in a coherent state with a definite phase , e.g.
$`\widehat{\psi }_a(x)|`$ $`=`$ $`e^{i\theta }\varphi _a(x)|`$ (41)
$`|\widehat{\psi }_a^{}(x)`$ $`=`$ $`e^{i\theta }\varphi _a(x)|`$ (42)
with $`\varphi _a`$ real. The mean number of atoms is given by
$$N_a=\widehat{N}_a=𝑑x\varphi _a^2(x).$$
(43)
Using Eq.(40) we find
$$\widehat{N}_a^{}\widehat{N}_b^{}=0$$
(44)
and
$$(\widehat{N}_a^{}\widehat{N}_b^{})^2=N_a+N_b+dxdy\varphi _a(x)\varphi _a(y)\mathrm{Re}[R_b(x,y,)+e^{2i\theta }S_b(x,y,)].$$
(45)
Ordinary vacuum in the $`b`$-state is the special case with $`S_b=R_b=N_b=0`$ and we note that the typical imbalance of populations is $`\sqrt{\mathrm{Var}(N_a^{}N_b^{})}=\sqrt{N_a}`$. Now, we use the nontrivial state created by photodissociation as squeezed vacuum. We imagine to have experimental control over $`\theta `$ and choose this phase optimally in order to reduce $`\mathrm{Var}(N_a^{}N_b^{})`$.
In Fig. 7 we plot this minimum value of $`\mathrm{Var}(N_a^{}N_b^{})/N_a`$ as a function of the time of production of the squeezed condensate. It is clearly seen how the noise is rather quickly suppressed almost perfectly. In the particular case shown $`N_a`$ was taken to be $`10^3`$ and $`\varphi _a`$ was of gaussian shape. The time-dependent number of atoms in the $`b`$-condensate can be approximately read out of Fig. 2. The width of $`\varphi _a`$ was chosen to be the same as the equilibrium width of the $`b`$-state atoms. It is amazing, but of course already well-known in quantum optics, that this supression can take place even though the number of atoms initially in the $`b`$-state is very small compared to the average number of atoms in the large condensate in the $`a`$-state
At times beyond $`\omega t2`$ the noise suppression is lost in the exact $`P_+`$ simulations. We recall that for $`g=0`$, the R&S equations are exact, and it is thus natural to ascribe the discrepancy between the two methods to the interactions, which, by analogy with the Kerr-effect in optics, cause a deformation of the gaussian state. In Fig. 8 we show a scatter plot of the amplitudes of the projections, $`\psi _{1a}`$, of 3000 $`P_+`$ realizations of $`\psi _1`$ on $`\varphi _a`$ at $`\omega t=2.4`$ for $`g=0.00`$ and for $`g=0.02`$. Such plots should be interpreted with care, as in general the *pair* $`(\psi _1,\psi _2)`$ is needed to calculate all normal-ordered expectations. The Kerr-effect, however, influences the $`S`$-function, and to calculate $`\widehat{\psi }^2`$ we only need to average $`\psi _1^2`$. In our case the relevant contributions to this average can thus be depicted by a scatter plot of $`\psi _{1a}`$ alone.
We see in Fig. 8a a gaussian state represented by points which form an ellipse-like structure. If the distibution had no prefered direction in phase-space $`\widehat{\psi }^2=\overline{\psi _1^2}`$ would vanish, but this is clearly not the case for our squeezed state. In Fig. 8b the the interaction modifies the phase accumulation of points with large values of $`\left|\psi _1\right|`$ and deforms the distribution to an S-like shape, and the mean value of $`\psi _1^2`$ is reduced. If the state is centered around a finite field amplitude, the intensity dependent phase shift transforms a circular distribution into a bean-shaped one, and in that case the Kerr effect actually produces squeezing .
## 5 Discussion
In this paper we have studied the matter wave analogue of an optical parametric oscillator. We have developed a description based on the correlations of the atomic field and we have supplemented this by the $`P_+`$ method. When interactions are taken into account both treatments were limited to relatively short times and low numbers of atoms. We showed how the state can be used as squeezed vacuum to improve the statistics of experiments with larger condensates.
On the theoretical side, it is of course interesting to remedy the breakdown of the R&S method as we know that for large condensates an even simpler description, the Gross-Pitaevskii equation, is very successful. To bridge the gap between the two regimes which are both describable in simple terms would be interesting and could be useful also in other scenarios where special quantum states of a BEC are created. In such work it is very likely that time dependent Monte Carlo methods, based for example on the positive P distribution, $`P_+`$, will play a significant role.
Let us finish this paper with a brief mentioning of other proposals for the preparation of quantum correlated atomic condensates. As mentioned in the introduction, the collisional interaction has been proposed as an agent for the preparation of quantum correlated states, such as Schrödinger cat-like states . The states prepared by our proposal are substantially less ’non-classical’, but provided the existence of a molecular condensate, we believe that the experimental requirements of our proposal are more easy to meet.
It was suggested some time ago to prepare spin-squeezed states of non-degenerate atomic gasses and of trapped ions , and concrete proposals were made, involving coupling of the atoms via their interaction with a quantized field mode , or their center-of-mass motion . More recently, it was suggested, and experimentally proven possible, to induce transitions between atomic states in an optically thick sample by absorption of non-classical light, and thereby to transfer quantum correlations from the light field to the atoms . Ingredients of these proposals become even more powerful with cavities around the samples . Very promising, recent ideas are based on quantum non-demolition (QND) measurements of, e.g., quadratures or populations in the atomic sample by refraction of light beams. Such detection suffices to establish quantum correlations in the sample, even though the detection is done with classical field states . One may simply use the state following the QND measurement together with the known classical outcome of the measurement to perform subsequent high-precision experiments, or one may consider feed-back loops, in which the system is driven towards a specific quantum correlated state.
All of these proposals are also relevant for degenerate gasses. The QND methods are probably closest to real implementation, since phase contrast methods, already used for imaging of condensates in many experiments, only have to be carried out with properly chosen parameters to lead to quantum correlated atomic states. But all methods are interesting, and they may illuminate various aspects of the condensate dynamics, as illustrated, e.g. by the threshold phenomenon and the ’Kerr-effect’ suppression of squeezing in the proposal in this paper.
## Synthesis of correlated noise
In order to numerically simulate Eqs.(23,25) we discretize time and space, and we synthesize noise terms, $`dW_{1,2}(x_n,t_i)`$, that obey discretized versions of Eqs.(26-29). It is well-known how independent (pseudo) random numbers from different distributions can be created. For example a Gaussian distribution can be created starting from uniformly distributed numbers via the Box-Müller method, i.e. we know how to produce $`\{dU_{1,2}(x_n,t_i)\}`$ so that
$`dU_{1,2}(x_n,t_i)`$ $`=`$ $`0`$ (46)
$`dU_\alpha (x_n,t_i)dU_\beta (x_m,t_j)`$ $`=`$ $`\delta _{\alpha \beta }\delta _{nm}\delta _{ij}.`$ (47)
We see that the correlation functions, Eqs.(28,29), contain two terms: one from the interaction and one from the incoupling. These can be treated separately if we split the noise in two independent contributions
$$dW_{1,2}(x_n,t_i)=dW_{1,2}^g(x_n,t_i)+dW_{1,2}^b(x_n,t_i)$$
(48)
Due to the contact form of the interaction, the corresponding noise term poses no difficulties; we simply choose
$$dW_{1,2}^g(x_n,t_i)=\sqrt{\frac{\pm ig\psi (x_i)dt}{dx}}dU_{1,2}^g(x_n,t_i),$$
(49)
where $`dU_{1,2}^g`$ is chosen with the properties (46,47). The incoupling term is created by multiplication and convolution of uncorrelated noise $`dU_{1,2}^b`$ of the form (46,47) with suitable Gaussian functions:
$$dW_{1,2}^b(x_n,t_i)=𝒩\mathrm{exp}\left(\pm it_i\mathrm{\Delta }\right)\mathrm{exp}\left(\frac{x_n^2}{2\sigma _a^2}\right)\underset{n^{}}{}dxdU_{1,2}^b(x_n^{},t_i)\mathrm{exp}\left(\frac{(x_nx_n^{})^2}{2\sigma _b^2}\right)$$
(50)
It turns out that choosing
$$\sigma _a^2=2\sigma _{cm}^2,\sigma _b^2=\frac{2\sigma _{cm}^2\sigma _r^2}{4\sigma _{cm}^2\sigma _r^2},𝒩=\sqrt{\frac{\pm idtB}{dx\sqrt{\pi }\sigma _r\sigma _{cm}\sigma _b}}$$
(51)
is sufficient to fulfill Eq.(26) with $`b`$ given by Eq.(7). $`\sigma _r`$ and $`\sigma _b`$ are rather small, and in practice the sum in Eq.(50) only needs to involve a few terms.
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# HIIphot: AUTOMATED PHOTOMETRY OF H II REGIONS APPLIED TO M51
## 1 Introduction
H II regions are an effective optical tracer of ongoing massive star formation. Even a single early-type star produces enough Lyman continuum photons to ionize a quantity of gas sufficient to produce readily detectable recombination lines in galaxies at distances of many Mpc. Observations in the Balmer lines (predominantly H$`\alpha `$) and other nebular emission lines (e.g. \[N II\] $`\lambda \lambda `$ 6548,6584, \[S II\] $`\lambda \lambda `$ 6717,6731, and \[O III\] $`\lambda `$ 5007) enable estimation of physical conditions within discrete H II regions (Evans & Dopita (1985), Osterbrock (1989) ($`AGN^2`$), Ferland, et al. (1998)). Conversely, measurement of the H II region luminosity function can indicate global patterns in the process of star cluster formation (Kennicutt, Edgar & Hodge (1989) (hereafter KEH89), Banfi, Rampazzo, Chincarini & Henry (1993), Feinstein (1997), Wyder, Hodge & Skelton (1997)). An excellent review of the field is provided by Oey & Clarke (1998).
We sought a fully-automated technique for determination of the positions, fluxes, and sizes of H II regions in galaxies. Such a tool is crucial for efficient and reproducible characterization of their star formation properties, especially if meaningful intercomparison between datasets is a primary goal. Another advantage of an automated approach over conventional measurement by visual inspection is that many systematic effects, like that of reduced spatial resolution for more distant galaxies or of differences in limiting sensitivity, can be quantitatively ascertained.
Well-resolved images of spiral and irregular galaxies reveal that massive star clusters often form in close proximity, usually leading to a confused and rather inhomogeneous distribution of ionized gas. This inherent clumpiness makes automated photometry difficult. Photometric accuracy can be further hindered by the presence of a variable background of diffuse ionized gas (DIG). In complex environments of this nature it is practically impossible to cleanly separate the contribution of neighboring extended objects to the observed surface brightness distribution. Methods of correcting for source overlap in the special case of stellar photometry (e.g. DAOPHOT, Stetson (1987); ALLFRAME, Stetson (1994)) cannot be applied without accurate models for the intrinsic structure of every source. At present it is infeasible to construct a comprehensive set of models spanning the observed properties of resolved H II regions in nearby galaxies. Any automated photometric procedure optimized for H II regions must consequently provide adaptivity to the actual source morphology. This can be accomplished using an iterative approach to “grow” sources from an initial guess at the shape. We have developed HIIphot, a user-friendly procedure written in IDL <sup>1</sup><sup>1</sup>1For information on the Interactive Data Language (IDL), see http://www.rsinc.com. which employs such a method.
McCall, Straker & Uomoto (1996) demonstrated the potential of an automated photometry procedure for H II regions based on a simple iterative growth mechanism. Their method, called percent-of-peak photometry (PPP), involved growth from local maxima down to a constant fraction of the difference between each peak and its local background. In the grand-design spiral NGC 3398 and the flocculent spiral NGC 4414, PPP successfully reproduced the LF obtained through standard fixed-threshold photometry (FTP). Kingsburgh & McCall (1998) have recently applied PPP during their analysis of four nearby dwarf galaxies. Unfortunately, McCall and collaborators are unable to recover more than a small percentage of the observed flux for even the brightest sources when using PPP. This stems from the fact that they can only grow to 70% of peak before the automated method becomes susceptible to rapid growth and merging of adjacent regions. Our present research was inspired by the desire to overcome this limitation using criteria to carefully regulate growth in saddle points between neighboring regions. Also, we hoped to recover all the observed flux rather than only that contributed by each region’s brightest pixels by defining larger “seeds” with a better match in shape to the source structure (rather than growing from a single pixel). McCall et al. argued that the flux detected using PPP was directly proportional the total source flux, but their line of reasoning assumed an idealized Stromgren sphere geometry for all regions. It is difficult to imagine that this assumption could be satisfied in general.
This paper is organized in the following manner. Section 2 presents the concepts and algorithms employed within HIIphot. Section 3 contains a very brief description of the M51 dataset used for illustrating the capabilities of the procedure. Section 4 presents the population of H II regions in M51, including a new, more sensitive luminosity function (LF). Finally, we conclude in Section 5 with a summary and view toward the near future.
## 2 HIIphot procedure
HIIphot is a completely automated method for photometry of H II regions. Our algorithm is sufficiently general to work well for distant galaxies, but provides the most substantial benefit during analysis of narrowband images of complicated, highly-resolved systems. Below is an explicit description of the procedure.
### 2.1 Initial detection of sources
Following the discussion of astronomical object recognition recently presented by Thilker, Braun & Walterbos (1998), hereafter TBW98, Thilker (1999); and Mashchenko, Thilker & Braun (1999), hereafter MTB99, we recognize that an ideal technique for decomposition of narrowband images into individual objects might employ: (1) calculation of projected physical models describing all anticipated source morphologies, (2) cross-correlation of image data with each model to find tentative matches, and finally (3) pruning of the composite detection list to correct for multiple detections of the same source. Regrettably, this direct approach is not currently viable due to a lack of sufficient computing power and a comprehensive set of models. One practical alternative might be to select a set of sufficiently diverse empirical models and evaluate the degree to which they match the data at some limited set of sky positions. HIIphot employs this strategy, followed by an iterative growing procedure to permit departures from the idealized models.
The HIIphot collection of empirical models includes six basic morphologies, each considered at various sizes and with major-to-minor axis ratios ranging from one to two, stepping by 0.25. We permit different position angles, sampling with an increment of 15°. In each morphology, the predicted surface brightness of the radially symmetric (base) model is computed as:
$$f(r)=exp\frac{(rr_0)^2}{2\sigma ^2}.$$
(1)
Figure 1 shows a model-center cross-cut for every morphological class. Each profile has been normalized to unit peak brightness. We include Gaussians by setting r<sub>0</sub> = 0, whereas ring models of varied “shell thickness” (relative to the ring diameter) are generated by taking r$`{}_{0}{}^{}>0`$ and adopting various $`\sigma /r_0`$ ratios. Specific choices of $`\sigma /r_0`$ were selected in order to sample thin rings, thick rings, and centrally depressed structures. Note that in Fig. 1 we varied r<sub>0</sub> with the intent of producing sources having the same characteristic size. As mentioned above, each radially symmetric base model is stretched and rotated in numerous ways for comparison with the data.
The essential challenge when incorporating these parameterized “guess morphologies” into HIIphot is finding a way to limit the number of sky positions at which any model must be compared with the data. Because we use at least 100 stretched/rotated variants for each base model of a given size, together with typically 50 base model sizes, it is prohibitive to compute a cross-correlation between each model and the entire image plane (as implemented by TBW98 for the case of H I shells). Instead, we determine a list of “tentative match” sky coordinates for each model by tabulating significant local maxima in the convolution of the data with an appropriately-sized circular Gaussian. In this manner, we detect structures with dissimilar morphology but having about the same size in a single pass. Our technique works because we only look for H II regions of a given characteristic size on images that have been smoothed to remove source structure on smaller scales. Even the most well-resolved ring (for instance) will end up looking like a Gaussian after some degree of spatial smoothing. We use “lowered” Gaussian kernels (as also employed in DAOPHOT, Stetson (1987)) as a means of removing slowly varying background structure from the galaxy image during our multi-resolution, convolution-based procedure. Each Gaussian kernel was truncated at a radius of 1.5$`\sigma `$ and offset with a constant in order to provide an integral over the kernel of zero. We tabulate solely those convolution maxima which have peaks exceeding a $`5\sigma `$ threshold. Variance associated with random fluctuations in the convolution of each Gaussian kernel with our data is measured in a user-selected sky region. Modest flat-fielding errors are not problematic due to our use of a lowered Gaussian as the convolution kernel.
After compiling a list of tentative centroid positions for sources of each characteristic size, direct comparison between a set of stretched/rotated models (Gaussians and rings) and the data is accomplished by calculation of a noise-corrected version of Pearson’s linear correlation coefficient, $`\rho `$. (As described in detail by MTB99, this statistic allows robust estimation of “goodness of fit” and completeness. Note that $`\rho `$ is invariant under linear scaling of the data, so the flux of a region and the level of its local background are irrelevant. Only the best match (highest $`\rho `$) model together with it’s value of $`\rho `$ are retained for each tentative source. The entire list of tentative sources is sorted by the $`\rho `$ value of the entries. We then employ a cutoff, $`\rho _{crit}`$, in the correlation coefficient in order to retain only the best matches. For this paper we adopted $`\rho _{crit}=0.25`$, although the median value was $`0.75`$. Remember that so far we are only creating a ranked list of possible detections.
### 2.2 HIIphot footprint and seed definition
Having this sorted list of potential detections, we next eliminate multiple detections associated with the same observed emission. This is accomplished by defining “footprints” in the image for each source. Beginning with the highest ranked detection, we loop over all regions allowing each one to “claim” pixels of the input image. Each detection is allowed to place a footprint if the following conditions are satisfied: (1) the associated model centroid has not been claimed, (2) 90% of the data flux inside the model’s 20% isophotal boundary remains unclaimed, and (3) the detection’s signal-to-noise is greater than $`\frac{S}{N}_{crit}`$. (See Section 2.6 for a detailed discussion of signal-to-noise in the context of HIIphot. We introduce a formal analysis based on uncertainties associated with the independent line+continuum and continuum images, rather than merely the continuum-subtracted image.) Regions satisfying these conditions take as a footprint all unclaimed pixels within the 20% isophotal level of their best-match model. Effectively, our footprint convention allows simultaneous rejection of multiple detections (naturally leaving only the best-match model) and introduces a buffer between neighboring regions. One can think of this procedure as a detailed fitting process in which all sources are compared with a finite number of relatively well-matched models.
Due to line-of-sight projection, one should anticipate overlapping H II regions in most galaxies (even if perfectly face-on, due to the finite disk thickness). Note that our methodology makes it possible to separately detect and analyze sources even when they have complete spatial overlap if their morphologies are sufficiently distinct. If a compact, highly significant source first places a footprint in an area containing many surface brightness enhancements with a range in size, the probability is substantial that a larger, less significant source will overlap the initial detection. This large detection will be allowed into the catalog provided the compact region does not contain more than 10% of the observed flux within the big model’s 20% isophotal boundary (presuming the other standard conditions for a footprint are also met). The HIIphot procedure naturally treats partially overlapping and fully overlapping detections in this manner.
Figure 2 illustrates the complete HIIphot procedure by showing the same image section within a continuum subtracted H$`\alpha `$ image of M51 at various stages of the processing. In particular, Fig. 2a shows our HIIphot footprints. The image data has been scaled logarithmically to keep from saturating the inner portions of the galaxy. Scaling is identical in each panel so as to facilitate comparison between panels 2c–2f. All marked regions are associated with a convolution peak ($`5\sigma `$ or better) at either original or somewhat degraded resolution. Fig. 3 shows two small subsections of Fig. 2c (see description further below) with linear scaling in order to demonstrate the significance of low surface brightness detections which are difficult to appreciate in Fig. 2. Not all of the detections shown in Fig. 3 are used in the construction of our H II region LF, as we demand that every “photometric source” have a final signal-to-noise ratio in excess of five. Nevertheless, all detections plotted are thought to be genuine, having been originally discovered using HIIphot and later confirmed by visual inspection of individual continuum-subtracted images (before CR-rejection) viewed at various resolutions. Recall that we never make use of (and draw no conclusions from) these intrinsically questionable detections. Essentially they should be considered candidates, until deeper observations become available.
Because our empirical models are only a first order approximation to actual source structure, footprints often contain pixels that are not bright enough to justifiably remain in the final boundary of the region. We account for this by rejecting all pixels which fall outside a “bounding isophote” defined by 50% of the median data value found within each footprint (where the median is measured relative to an estimated local background). We call these trimmed footprints “seeds” since they are composed of only those pixels destined to belong to a region, but do not yet reflect changes associated with the iterative growth procedure. Our procedure ensures that all seed boundaries follow isophotal contours within footprint boundaries, although the specific cutoff varies depending on the distribution of pixel intensities within any given H II region. Notice that this conservative approach makes it possible for ring-like footprints to reject pixels which fall within the object’s central surface brightness depression. Fig. 2b shows the HIIphot seeds for our subsection of M51. In practice, our particular implementation of the “seed” convention is motivated by the following arguments: (1) multi-pixel seeds provide a “head start” for the iterative growth process, making it easier to reliably separate adjacent regions, (2) defining seeds as a data-regulated subset of model footprints allows the first true excursion of region boundaries from our set of empirical models.
### 2.3 Iterative growth of detections
Given a set of seed pixels associated with every H II region in a galaxy, it might seem a simple matter to iteratively add pixels to each region until reaching the maximum extent of all nebulae. In fact, the implementation of a well-behaved iterative growing algorithm is far from trivial and there is no established convention for determining the “edge” of an H II region. McCall, Straker & Uomoto (1996) encountered difficulty in growing their sources to isophotal cutoffs fainter than about 70% of the local peak. Potential inhomogeneity in the diffuse background level and crowding of regions having remarkably different flux conspire to make the PPP method less suitable except in a limited set of well-behaved circumstances. HIIphot attempts to carefully control the rate of growth in saddle points between regions by introducing a slowly declining threshold which determines the set of pixels considered for growth during a given iteration. Pixels having values below this global threshold are ignored until later iterations. In this way, neighboring regions approach their saddle point at an equal rate no matter what the difference in peak value or total counts between sources.
Iterative growth commences by setting the global threshold for pixel consideration equal to the highest bounding isophote and is reduced by 0.02 dex before each subsequent iteration. Regions as a whole are considered for growth only if the median value of the pixels in a seed’s “exterior perimeter” exceeds the slowly declining threshold. This implies that only the seed having the highest bounding isophote is considered during the first iteration. Any time that more than one region is allowed to grow during an iteration, HIIphot cycles through the active regions in order of decreasing correlation coefficient, $`\rho `$. Qualified pixels (lying above the global threshold) which are adjacent to or diagonal from any pixel already belonging to the region being augmented are potentially added to the source if they are not claimed by other regions. That is, pixels from the exterior perimeter of a region can be added if they are bright enough. We also require that at least 50% of the perimeter pixels are added during any given iteration. If this is not the case, we postpone growth until the global threshold declines further, so as to simultaneously add most of an entire isophotal ring. Growth for a particular region continues in this manner until either: (1) the observed surface brightness profile flattens sufficiently, or (2) no more qualified pixels can ever be reached due to being surrounded by other regions or because of the intrinsic data values. Note that regions can “stall” for many iterations and do not immediately cease growing just because neighbor pixels cannot be considered (as a result of the global threshold). In other words, our iterative procedure amounts to carefully adding lower isophotal contours to all qualified regions after specifically accounting for a slightly unequal start brought about by our adaptive definition of seed boundaries.
Figure 4 shows schematic representations of a hypothetical source being considered for iterative growth. In panel (4a), the dark shaded pixels belong to the region’s interior perimeter set during active iteration $`n`$, whereas lighter colored pixels compose the exterior perimeter group. The current region boundary is indicated with a heavy solid line. Median values for the interior/exterior perimeter sets will be used to determine if the surface brightness profile has flattened sufficiently in order to stop further growth (in subsequent iterations). Some of the lightly shaded pixels have been marked with a circle. These exterior perimeter pixels have values above the global threshold and will be added to the region during iteration $`n`$. Note that more than 50% of the lightly shaded pixels fall into this category. If this had not been the case, growth for this region would stall until the HIIphot global threshold declined enough to allow a majority of the exterior perimeter pixels to augment the region. Panel (4b) is similar to Panel (4a), except that it has been drawn for the following iteration, $`n+1`$.
The question of how to determine whether a surface brightness profile has “flattened” is somewhat difficult to treat on anything other than pragmatic grounds. Presently there is no established connection between the rate of surface brightness decline and specific physical conditions within an H II nebula. Originally we demanded that regions grow until the difference in median values between interior and exterior perimeter pixel sets indicated the surface brightness profile was no longer declining. This choice resulted in very large H II regions and significant bumping of adjacent regions, since in crowded fields brightness profiles rarely flatten out before encountering a neighbor. Our reason for requiring that the surface brightness profile flatten completely was that we sought to fairly treat all regions, regardless of their environment. Using this procedure we effectively determined groups of pixels most plausibly associated with the same ionizing source. That is, for this methodology, our H II “regions” included compact cores and related diffuse emission (DIG). Although interesting in its own right, this non-conventional definition of an H II region makes it difficult to compare the current results with previous work and we sought a more flexible alternative.
In the end, we elected to permit an array of different stopping points ranging from very little growth to nearly the generous “flat result” described above. This amounted to adopting a series of cutoffs in terminal surface brightness slope, \[10, 4, 2, 1.5, and 1\] EM/pc, then running the growth procedure from beginning to end for each. Notice that the specific cutoff values given here are only appropriate if the calibrated narrowband data are expressed in the conventional units of EM, cm<sup>-6</sup> pc, and must be rescaled for any other case. In section 4.2, we show that this approach allows us to directly address systematic uncertainty in H II region fluxes (and the resulting luminosity function) associated with our decision to stop growth at a given point. Figs. 2c and 2e present images of the M51 subsection with H II region boundaries marked for two different values of the terminal surface brightness slope. In Fig. 2c growth has stopped at a slope of 10 EM/pc, leaving a substantial fraction of diffuse emission possibly associated with discrete H II regions remaining outside the HIIphot boundaries. Fig. 2e shows the result for growth continuing until surface brightness profiles flatten to 1 EM/pc. Notice how isolated regions do eventually stop growing on their own, while the crowded central area has effectively been subdivided into numerous chunks (each plausibly associated with an embedded ionizing source). The behavior of the iterative growth procedure is further illustrated in Fig. 5, which shows another section of M51 at several stages of growth. For reference, all the results presented in this paper are based on a terminal profile slope of 1.5 EM/pc, midway between the degree of growth shown in panels (e) and (f) of Fig. 5.
One important advantage of adopting terminal surface brightness slopes (specified in physical units) is that our procedure is relatively robust to changes in signal-to-noise. One can think of other criteria for stopping growth that are not as reliable. For instance, we initially tried to quantify the “stopping point” for H II region growth in terms of various critical multiples of the formal error in the dimensionless surface brightness slope. This procedure appeared promising when analyzing our basic dataset, but was found to introduce substantial bias in the definition of H II region boundaries during experiments in which the signal-to-noise was globally reduced by factors ranging from 2 through 5. In short, as the test images were made noisier, growth stopped progressively sooner despite the fact that the underlying observed surface brightness profiles were no different. Our adopted procedure is substantially more well-behaved under these circumstances and generates luminosity functions which are statistically indistinguishable at all but the lowest luminosities. Sadly, loss of low luminosity sources is unavoidable with degraded signal-to-noise no matter how regions are grown.
### 2.4 Correction for underlying emission
For H II regions embedded in a background of diffuse ionized gas it is important to accurately estimate the DIG flux contribution to the observed counts within a region’s boundary. In past studies, most authors have gauged the background contribution by interactively selecting one or more positions near each H II region they thought to be representative of the level underlying the source. Our method works as follows: (1) after final region boundaries are available, we define as “background” pixels all those unclaimed pixels within a projected distance of 250 pc from the boundary of an H II region, (2) next we select a uniformly-spaced set of “control points” to represent these background data, only accepting those which are at least 75% surrounded by other background pixels within a circular domain of 250 pc diameter, (3) we then compute the median value of all background pixels within the domain of each control point, and (4) we finally compute a surface fit to these median values. Our surface-fitting procedure generates a low-order solution on small scales by interpolating between the 3 nearest control points at every position, but in a global sense the product is a very high-order surface. The result is essentially an image of the diffuse emission present in the original data and therefore represents an excellent means of quantifying the diffuse fraction in galaxies (e.g. Hoopes, Walterbos & Greenwalt (1996)). Note that we compute a different surface-fit for each requested version of the region boundaries, as the degree of iterative growth will influence pixel membership in background annuli and therefore the estimated background level for each emission line source. Figs. 2d and 2f show the diffuse background for the growth states illustrated in Figs. 2c and 2e, respectively. Note how the level of diffuse emission is estimated to be substantially lower in the second case.
### 2.5 HIIphot data products
The output of HIIphot consists of several images and one catalog detailing properties of all detected regions. The catalog tabulates the following quantities among others: ID#, right ascension, declination, pixel position, number of pixels contained by the region, effective FWHM, major axis FWHM, axial ratio, position angle, total flux after correction for background emission, $`1\sigma `$ uncertainty in total flux after correction, and the peak surface brightness inside the region. See Section 2.6 for a description of how total corrected flux and its error are calculated. Note that right ascension, declination, pixel position, and FWHM values refer to the best-fitting empirical model associated with each region, so from before region growth.
The images produced by HIIphot include: (1) a copy of the continuum-subtracted line image with the various “after growth” boundaries marked, (2) several versions of the background surface fit (corresponding to different levels of growth), and (3) integer maps delineating the position and extent of each footprint, seed, and grown region. These integer maps can be used for supplementary analysis if identically gridded images at different wavebands are available. Among the most obvious applications are computation of line ratios or equivalent width for emission line objects. Finally, the HIIphot user has the option of dumping postage stamp collages depicting each source in the catalog.
### 2.6 Flux determination and signal-to-noise in HIIphot
The background-corrected emission line flux of an H II region is computed using the continuum image ($`C`$), the line+continuum image ($`L`$), and our HIIphot surface fit to diffuse background emission remaining in the continuum-subtracted line image after growth of sources ($`D`$). In this derivation we assume that $`C`$, $`L`$, and $`D`$ remain in ADUs. Additionally, we require that no sky background has been subtracted from either $`C`$ or $`L`$. This is essential if photon noise is to be properly modeled during estimation of signal-to-noise. For region $`i`$ (composed of pixels $`j=1\mathrm{}n_i`$) the background corrected emission line flux, $`F_i`$, is calculated as:
$$F_i=(L_{ij}S_iC_{ij})D_{ij},$$
(2)
where $`S_i`$ is the continuum scaling factor appropriate for region $`i`$. In practice we hold $`S_i`$ constant for all regions. The formal $`1\sigma `$ uncertainty, $`\delta F_i`$, associated with $`F_i`$ is given by the quadratic sum of standard deviations associated with individual terms of Eq. 2:
$$\delta F_i=\frac{(\delta L_{ij}^{}{}_{}{}^{2}+S_{i}^{}{}_{}{}^{2}\delta C_{ij}^{}{}_{}{}^{2}+\delta D_{ij}^{}{}_{}{}^{2})^{1/2}}{g}$$
(3)
As written here, the units of $`F_i`$ and $`\delta F_i`$ are ADUs, while $`g`$ is the gain in terms of electrons per ADU. To convert into physically meaningful units we multiply by an appropriate calibration factor. We assume that $`C`$ and $`L`$ are both essentially sky-noise limited, implying the following relations:
$$\delta L_{sky}=\frac{1}{g}\left(\frac{gX_LL_{sky}}{n_L}\right)^{1/2},\mathrm{and}$$
(4)
$$\delta C_{sky}=\frac{1}{g}\left(\frac{gX_CC_{sky}}{n_C}\right)^{1/2}.$$
(5)
In these expressions, $`n_L`$ and $`n_C`$ are the number of images (assumed to have comparable exposure) combined to create $`L`$ and $`C`$, respectively. Multiplicative factors $`X_L`$ and $`X_C`$ have values near unity or slightly higher in order to account for the possibility that read-noise may still make a small contribution to the noise budget in the HIIphot sky region. $`\delta L_{sky},\delta C_{sky},L_{sky},`$ and $`C_{sky}`$ are each measured within the sky region of the input images, implying appropriate values for $`X_L`$ and $`X_C`$. This information then constrains a photon noise model for the data, as we can fold $`X_L`$ and $`X_C`$ together with $`g`$ to represent an effective gain ($`g_L=gX_L`$ and $`g_C=gX_C`$) for $`L`$ and $`C`$.
Next, we estimate the level of noise per pixel in the brighter, interesting portions of $`L`$ and $`C`$. The relevant equations are:
$$\delta L_{ij}^{}{}_{}{}^{2}=\frac{g_LL_{ij}}{n_L}e^{},\mathrm{and}$$
(6)
$$\delta C_{ij}^{}{}_{}{}^{2}=\frac{g_CC_{ij}}{n_C}e^{}.$$
(7)
Eqns. 6 and 7 specify most of the terms in Eq. 3. Variance in the diffuse background level, $`\delta D_{ij}^{}{}_{}{}^{2}`$, is determined empirically on the basis of the measured standard deviation near control points used during the background fitting process. Note that we convert such measurements to electrons before computation of $`\delta F_i`$.
Two comments must be put forward at this point. Our assumption of sky-noise limited imagery is a conservative approach. By adopting Eqs. 47, we guarantee that $`\delta L_{ij}`$ and $`\delta C_{ij}`$ will always be predicted accurately or overestimated. If readnoise contributes substantially to the standard deviation of pixel values in the user-selected sky region, the measured values of $`g_L`$ and $`g_C`$ insure that it will contribute an identical fraction of the estimated noise-budget for pixels that are substantially brighter (due to observed emission from the object of interest). In reality, this is not the case – detector readnoise is independent of the observed pixel intensity. Overprediction of error terms $`\delta L_{ij}^{}{}_{}{}^{2}`$ and $`\delta C_{ij}^{}{}_{}{}^{2}`$ is the unavoidable consequence. Furthermore, we argue that our procedure for evaluating the source signal-to-noise ratio ($`F_i/\delta F_i`$), is more accurate than the traditional method based only on continuum-subtracted data, especially in the limit of bright continuum emission. Previous studies have usually gauged the standard deviation per pixel in one or more selected “sky areas,” then added this term in quadrature based on the number of pixels in a region. This implies that their estimated signal-to-noise is independent of the original observed datavalues. Identical sources located in various positions on top of a variable background of (continuum or line) emission will be assigned identical signal-to-noise, even though the true uncertainty increases for sources embedded in a bright background.
## 3 Narrowband observations of M51
We obtained narrowband images of the M51/NGC5195 system as part of a separate project concerning diffuse ionized gas (DIG) in spiral galaxies (Greenawalt (1998)). These data were obtained in 1992 March using the No.–1 0.9 m telescope at Kitt Peak. Three H$`\alpha `$ \+ \[N II\] and two \[S II\] images of 20 min integration were recorded in addition to a set of offband continuum exposures. The bandpass of each filter was approximately 75Å. Complete details concerning our observations and image reduction can be found in Greenawalt, Walterbos, Thilker & Hoopes (1998). For the present analysis a slightly different flux calibration was used. The continuum-subtracted image originally presented in Greenawalt et al. was calibrated by comparison with the photometric data of van der Hulst, Kennicutt, Crane & Rots (1988). An identical procedure was employed by Rand (1992) in a detailed study of 616 M51 H II regions. Because we sought to compare our photometry directly with Rand’s, we bootstrapped to his flux scale using a sample of 10 bright, moderately isolated H II regions. The magnitude of this recalibration amounted only to $``$10%, most likely attributable to the use of different regions by Rand and Greenawalt.
The 1$`\sigma `$ noise of our continuum-subtracted line image is at an H$`\alpha `$ emission measure (EM) of 9.9 pc cm<sup>-6</sup> at $`1.8\mathrm{}`$ FWHM resolution. This corresponds to a surface brightness of $`2.0\times 10^{17}`$ erg cm<sup>-2</sup> s<sup>-1</sup> arcsec<sup>-2</sup>, or 3.5 Rayleighs. Our noise implies a limiting (5$`\sigma `$) point source flux of $`3.6\times 10^{16}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, or equivalently an H$`\alpha `$ luminosity of $`3.9\times 10^{36}`$ erg s<sup>-1</sup>, neglecting any background confusion. For this calculation and the analysis below, we have assumed a distance of 9.6 Mpc to M51 (Sandage & Tammann (1974)). At this distance, the $`1.8\mathrm{}`$ seeing prevailing during our observations corresponds to a linear resolution of 84 pc. No correction was attempted for extinction, in order to facilitate comparison with earlier H II region surveys. In reality, van der Hulst, Kennicutt, Crane & Rots (1988) report on extinction toward a large sample of M51’s most luminous H II regions, finding about 2 visual magnitudes in most cases. This should be kept in mind when interpreting our results.
## 4 Results
We detected 1618 H II regions in the field of view of our observations, excluding sources located in M51’s interacting companion NGC5195. Of the total sample, 1229 regions were classified as “photometric-quality” detections having $`\frac{S}{N}_{final}5`$. Only these 1229 H II regions have been considered in the analysis described below. Fig. 6 displays our continuum-subtracted H$`\alpha `$ image, with all source boundaries marked. Note that the extent of each H II region has been determined using a terminal surface brightness slope of 1.5 EM/pc. All results in the rest of the paper correspond to this choice, unless stated otherwise. The image has been logarithmically scaled in order to preserve contrast over a wide range in surface brightness. Notice the hand-drawn loop surrounding NGC 5195. It indicates the region specifically excluded from our M51 HIIphot survey.
### 4.1 Overall comparison with R92
In this section we compare our H II region catalog with the list compiled by Rand (1992), hereafter R92, based on visual inspection and classification. We highlight the overall correspondence between our results and those of R92, but also describe variations attributable to procedural differences. We include a look at catalog completeness as a function of luminosity and morphological properties.
The most straightforward comparison between both catalogs is the total number of detected H II regions. Our detection list encompasses the entire galaxy, even the outer disk and confused nuclear portions not considered by Rand, suggesting we should find more sources than previously reported. However, other competing factors need to be considered as well:
(1) During the definition of footprints, HIIphot considers if a set of neighboring peaks is best described as a collection of individual regions or should be grouped into one or more composite aggregations. R92 always classified each neighboring peak as a separate source.
(2) HIIphot is perfectly consistent during the evaluation of marginal detections. Our estimate of signal-to-noise for each detection is evaluated on the basis of both the line+continuum and continuum datasets rather than just the continuum-subtracted line image (see Section 2.6).
(3) The catalogs were generated using images having different intrinsic sensitivity. The 1$`\sigma `$ noise in our data was $`2.0\times 10^{17}`$ erg cm<sup>-2</sup> s<sup>-1</sup> arcsec<sup>-2</sup>, whereas Rand’s imagery went down to $`1\times 10^{17}`$ erg cm<sup>-2</sup> s<sup>-1</sup> arcsec<sup>-2</sup> (EM = 5 pc cm<sup>-6</sup>), when evaluated at comparable spatial resolution.
(4) Contamination by emission-line objects other than traditional H II regions is a concern in our catalog. As a fundamental part of the R92 source selection process, each tentative detection was individually checked in a number of ways. Rand demanded that every source be centrally peaked, have a limited degree of circular symmetry, and possess a peak flux exceeding the background by more than 50%. HIIphot does not use any of these criteria. This means that our procedure is much more likely to result in a catalog containing “objects” such as localized enhancements in the widespread curtain of DIG, whether they be ionized on the spot by an embedded OB field star or by leakage from an H II region in a neighboring part of the galaxy, possibly hundreds of pc away. We elected to accept detections of this sort for three reasons: (1) they are intrinsically interesting, (2) they have little effect on the derived slope of the H II region luminosity function, and (3) eliminating such detections from the catalog would either require human intervention or an extra a priori constraint on the properties of detected H II regions. Note that contamination by planetary nebulae is equally unlikely in the R92 and HIIphot catalogs, since they would be too faint at a distance of 9.6 Mpc. Vassiliadis, Dopita, Morgan & Bell (1992) show that in the Magellanic Clouds there are no planetary nebulae with log L$`{}_{H\alpha }{}^{}>36`$, suggesting that both H II region catalogs are probably uncontaminated by PNe.
In the range of galactocentric radius examined by Rand (1 kpc $`<R_g<`$ 15 kpc), he reported 616 individual H II regions. HIIphot detected 1184 H II complexes with fluxes $`5\sigma `$ in the same area. Some of these objects are composed of multiple components. Although the total number of regions tabulated by HIIphot is more than reported by R92, we find that the agreement is substantially better in the range log L$`{}_{H\alpha }{}^{}>38.5`$. Rand detected 67 H II regions of this luminosity class, whereas we have 80.
The astute reader might ask how many of the regions detected by HIIphot are exactly the same sources described by R92. We explicitly checked this issue, finding that HIIphot misses only 2 of the 616 regions of R92. At the position of these two sources, we inspected our data and found no evidence for a significant detection. Note that our assessment of correspondence between the R92 source list and the HIIphot catalog was completed by way of manual inspection. During this process we demanded not only positional agreement, but also similar size between detections considered as being one in the same. It is important to note that due to the different methods of photometry it is unlikely that most sources had exactly the same effective boundary.
Figures 7 and 8 present a comparison of the measured luminosities for the regions in common to both our catalog and the R92 source list. Fig. 7 shows that there is a clear correlation (having slope $`1`$) between the luminosities measured by Rand and HIIphot. A limited number of H II regions fall substantially below the main cloud of datapoints. These sources most likely represent cases in which HIIphot broke a single R92 detection into one principal H II region and a small number of fainter peripherial sources. Figure 8 more clearly indicates a very subtle systematic change in the ratio of HIIphot/R92 luminosities as a function of source strength. We find that HIIphot tends to return slightly higher flux levels for very faint sources in comparison to the measurements of R92. This trend only begins in earnest at luminosities well below Rand’s estimated completeness limit (at L$`{}_{H\alpha }{}^{}10^{37.6}`$ erg s<sup>-1</sup>). As shown in the next section, this systematic bias and the inherent scatter in Fig. 7 has very little (if any) influence on the H II region LF.
So what are the detections “missed” by R92? In the area surveyed by Rand, we find that $``$80% of the regions picked up by HIIphot (but not listed in R92) have an H$`\alpha `$ luminosity less than $`10^{37.5}`$ erg s<sup>-1</sup>. Morphologically it is clear why most of them were not included in R92 – often they are rather diffuse and/or elongated. Sometimes these faint sources fall in the bright halo of a more significant H II region. In this case, interactive inspection of our continuum-subtracted line image typically reveals a relatively compact source superimposed on the brighter neighbor.
### 4.2 The H II region luminosity function
Figure 9 presents a comparison of the R92 differential luminosity function and our HIIphot result for exactly the same sources. Above their turnover points, both luminosity functions can be modeled as power law distribution (as first suggested by KEH89). Using only the data for regions in bins with log L$`{}_{H\alpha }{}^{}37.6`$ (those thought to be complete in R92) and assuming the standard functional form,
$$dN(L)=AL^\alpha dL,$$
(8)
we find that $`\alpha _{\mathrm{H}\mathrm{II}\mathrm{𝑝ℎ𝑜𝑡}}=1.60\pm 0.07`$, whereas $`\alpha _{\mathrm{R92}}=1.66\pm 0.07`$. (For these fits we assumed simple counting statistics in order to assign a variance to each value of the LF. The weights used during $`\chi ^2`$ minimization were inversely proportional to the variance of each bin, essentially giving the most influence to bins with the highest number of detections and reducing the influence of under-populated bins. This procedure is appropriate as long as no bins suffering from incompleteness are included in the fit.) The fact that there is no significant difference between the results plotted in Fig. 9 indicates that our HIIphot flux measurement technique does not introduce bias into the observed H II region luminosity functions.
Restoring the regions ignored for our comparison with R92, Fig. 10 presents an observed LF for all 1229 of the M51 sources detected by HIIphot with $`S/N_{final}5`$. The weighted power-law slope for this distribution is $`\alpha _{\mathrm{H}\mathrm{II}\mathrm{𝑝ℎ𝑜𝑡}}=1.75\pm 0.06`$, including only bins for which log L$`{}_{H\alpha }{}^{}37.6`$. Notice the break in the power-law at a luminosity of 10<sup>38.9</sup> erg s<sup>-1</sup>. Our results confirm that M51 has a Type II LF, as originally defined by KEH89. For the purpose of comparison, Fig. 10 also presents the differential luminosity function from R92 - including the 2 sources not detected by HIIphot (this explains the slight difference with respect to Fig. 9). Rand estimated that his LF was complete down to log L = 37.6. In the text below, we carefully address incompleteness in the HIIphot catalog. In any case, Fig. 10 shows that our observed LF is marginally steeper than reported in R92. This probably reflects the enhanced sensitivity of our procedure at low luminosities and for relatively diffuse H II regions.
Note that our H II LF is subject to a systematic uncertainty associated with our choice of when to stop growing regions. In particular, the observed LF slope varies substantially if growth is stopped early or allowed to continue until region surface brightness profiles are more nearly flat. The LF slope quoted above ($`\alpha =1.75`$) was obtained by growing H II regions until a terminal surface brightness slope of 1.5 EM/pc. If we had instead adopted a 1 EM/pc cutoff, the LF slope would have been shallower ($`\alpha 1.7`$). In the case of minimal growth, the LF slope approaches $`\alpha =1.9`$. We view this systematic difference as more of a change in the definition of an H II region, rather than uncertainty in our nominal result. The key point is that this “bias” can be properly addressed in a study of a sample of spirals by adopting the same procedure for all galaxies.
### 4.3 Investigation of incompleteness and blending effects
It is important to assess incompleteness effects. We investigated systematic trends such as the loss of faint, or bright but relatively diffuse, H II regions using simulations in which artificial sources were added to our original images. These altered data were subsequently reprocessed using HIIphot.
Blending due to limited spatial resolution can induce catalog incompleteness in crowded environments and flatten the observed LF at the faint end. H II regions tend to be inherently clumpy in terms of their spatial distribution, especially along spiral arms. This implies that the distribution of artificial H II regions should not be uniform across an image, but instead that additional H II regions should be placed preferentially in areas having recent star formation. Two ways to achieve this result are described below:
(1) Select a representative subset of detected regions in a galaxy and permit limited random walks away from actual tabulated positions, adding an artificial H II region in each slightly-randomized location.
(2) Use our HIIphot surface fit to the diffuse emission throughout a galaxy as a weighting function for the probability of placing an artificial H II region in any particular spot.
We elected to use the second method, as it provided more flexibility. For comparison, we also produced simulated images in which we distributed artificial sources at random.
We sought to reflect the intrinsic morphological diversity of H II regions in the prescription employed to generate artificial H II regions, rather than just adding unresolved sources of varied luminosity. Our incompleteness testing procedure allowed 3 types of simulated H II region: (1) small elliptical Gaussians, FWHM$`{}_{eff}{}^{}=100`$ pc, of varied axial ratio and position angle; (2) large elliptical Gaussians, FWHM$`{}_{eff}{}^{}=200`$ pc, of varied axial ratio and position angle; and (3) background-subtracted copies of actual H II regions extracted from the original data (mean FWHM<sub>eff</sub> = 134 pc), scaled down to varying flux levels. In all cases, photon noise was added to each source before adding it into the line+continuum image being modified. Note that all direct image modification was performed on line+continuum images, since they are the relevant “observable” data. Afterwards, continuum subtraction was performed to compute a modified continuum-free line image. HIIphot was run using the modified continuum-subtracted line image, the modified line+continuum image, and the original continuum image.
Fig. 10 suggests that our observed M51 H II LF might begin suffering from incompleteness for sources as bright as $`10^{37.037.5}`$ erg s<sup>-1</sup> (just above the turnover point). We chose to insert 100 artificial regions of each type at 5 distinct luminosity values, log$`(L_{H\alpha })=36.2,36.6,37.0,37.4`$, and $`37.8`$. Obviously each combination of source type, luminosity, and spatial distribution was investigated during a separate trial. The goal was to reliably constrain the true (“corrected”) H II LF down to log$`(L_{H\alpha })36.6`$.
Because actual (and simulated) H II regions have a spatial extent defined by irregular boundaries, one cannot simply inter-compare center positions for each detection in order to determine if simulated H II regions have been recovered. Each simulated source was assigned a code indicating whether it was: (1) recovered cleanly, (2) recovered as a blend, (3) essentially unrecovered, but partially blended with one other region, (4) essentially unrecovered, but partially blended with multiple regions, or (5) completely unrecovered. Sources were considered to be “recovered” (codes 1 and 2) if a single detection boundary encompassed pixels that contained at least 2/3 of the inserted region’s footprint flux (above the 20% isophote of the simulated source), otherwise the source was labeled “unrecovered” (codes 3, 4, and 5). For successfully recovered sources, cleanliness of recovery was judged by the fraction of total detection flux contributed by the simulated H II region. If a simulated region contributed at least 50% of the total flux in a detection, then it was assumed to be a clean recovery (code 1). Otherwise, blended recovery was indicated (code 2). For unrecovered synthetic sources, we evaluated the number of neighboring detections claiming at least one pixel of the unrecovered source footprint. If no pixels belonging within the simulated region’s 20% isophote were part of an HIIphot detection, then the artificial source was considered completely unrecovered (code 5). Likewise, if one and only one HIIphot detection claimed a pixel belonging to the unrecovered source footprint, the synthetic source was labeled essentially unrecovered, single blend (code 3). If more than one detection claimed a synthetic footprint pixel, then code 4 (multiple blend) was indicated. As a tool for determining the dependence of recovery statistics and photometric accuracy on variations in the local environment, we also classified the degree of crowding in the vicinity of each simulated H II region.
Before discussing the results of our completeness testing procedure, we note that the simulations provide a way to quantify the accuracy of our photometry as a function of luminosity. Because we know the exact flux of all simulated regions added to an image, we can determine the standard deviation of flux measurements for cleanly recovered sources. We examined the distribution of fractional flux discrepancy, defined as $`(F_{observed}F_{true})/F_{true}`$, for each cleanly recovered artificial source without close neighbors. We find that the standard deviation of fractional flux discrepancy increases with decreasing source luminosity (as expected), ranging from 0.1 for log$`(L_{H\alpha })=37.4`$ up to 0.3 for log$`(L_{H\alpha })=36.6`$. Fractional flux discrepancy values were negligible for log$`(L_{H\alpha })37.8`$. The measurement scatter is significantly reduced for small sources. Furthermore, the median value of fractional flux discrepancy is very near zero for all but our faintest artificial sources.
Tables 1 and 2 present the end results of our completeness testing procedure for M51. Specifically, we have tabulated the percentage of simulated detections falling in each of the five recovery categories for all region types and luminosity values. Table 1 indicates the values for source placement via weighting the distribution of artificial sources to regions of star formation, while Table 2 shows what was recovered for the uniform source distribution.
As expected the simulations indicate our H II LF begins to be substantially incomplete by log L = 37.4 (about $`L_{H\alpha }/L_{H\alpha ,RMS}=32`$). At this luminosity, nearly one third of the “actual” variety synthetic H II regions could not be recovered by HIIphot (see Table 1). We do find that small sources are easiest to recover. Large Gaussians were much more susceptible to blending with one or more sources. Actual regions appear to be intermediate – harder to recover than 100 pc Gaussians ($`1.2\times `$ PSF FWHM), but significantly easier than 200 pc Gaussians ($`2.4\times `$ PSF FWHM). These statements hold for both the weighted and uniform source distributions. Table 1, which shows the results for our weighted distribution tests, is most appropriate for the galaxy at large. However, the uniform distribution recovery statistics should be used when looking at completeness issues in uncrowded regions.
The results of our completeness testing procedure allowed us to perform Monte Carlo simulations designed to gauge systematic bias due to blending of faint, indistinct regions with brighter sources in observed luminosity functions. A separate paper will discuss the detailed findings of this investigation in a more general context. However, we were able to show that for the M51 completeness statistics (presented in Tables 1 and 2) the slope of the luminosity function above the low luminosity turnover was rather insensitive to “upward contamination” (see R92) potentially brought about via blending.
This result is actually somewhat of a coincidence related to the specific observed power law slope of the M51 LF. For intrinsically steeper luminosity functions, having $`\alpha 2.0`$ for instance, blending can lead to a shift in the turnover point (to higher L) and create an artificial hump at slightly higher L (in excess of the true number of sources per bin). Shallower LFs than M51 are less susceptible to blending effects, but suffer severely from non-detection of low luminosity regions. For such systems, the turnover of the LF becomes rather broad and fitting a power law slope to bins just above the turnover leads to an underestimate of the true $`\alpha `$ (that is, we are fooled to think the LF is shallower than it actually is). As stated above, the M51 power law slope ($`\alpha =1.75`$) is just shallow enough to avoid severe upward contamination, but not yet flat enough to substantially change the histogram character near the turnover point. Consequently, we conclude the observed M51 LF slope is rather robust to systematic bias and suspect that the true (unobservable) LF slope falls within the quoted uncertainty range for $`\alpha `$.
### 4.4 Distance-related effects on the observed LF
Systematic application of HIIphot to a large sample of galaxies will be able to address the effects of limited spatial resolution and sensitivity on observed H II region luminosity functions. Both of these observational characteristics are directly related to the distance for the object of interest. We can gauge the intrinsic bias related to limited resolution and sensitivity by deriving two LFs for each galaxy, one at the actual distance of the observed system and another using data which has been degraded to make the observations appear as if the galaxy was at the distance of our most-removed system. Although not really an issue in the present context, since we are studying a single galaxy which is already moderately distant (9.6 Mpc), this section has been included to demonstrate the technique and show that it is rather easy to realize given the HIIphot procedure.
We adopted a conservative procedure for generating image sets corresponding to the same galaxy at various distances. Instead of merely convolving the continuum-subtracted line image, then regridding, and adding noise (as is typically done), we independently transform the line+continuum and continuum images, only then creating the continuum-subtracted result. This procedure is required to accurately keep track of the photon statistics associated with continuum emission underlying H II regions. Neglecting this “hidden” noise may result in an overestimate of sensitivity when mimicking the effects of increased distance to a system.
The following step-by-step summary explicitly outlines our procedure:
(1) Select a “blank sky” region within the field of view of the continuum-subtracted line image.
(2) Determine the median level and standard deviation of this sky region in both the line+continuum and continuum images.
(3) Subtract the respective sky level from the line+continuum and continuum images.
(4) Determine the total number of galaxy counts in each image.
(5) Convolve with an appropriate Gaussian kernel (having peak of unity) in order to reduce the spatial resolution in both images.
(6) Regrid the convolved line+continuum and continuum images to a scale which results in the same PSF as the original data.
(7) Scale down the number of counts in each regridded image to be consistent with the totals determined in Step 4. That is, $`F_{total,new}=F_{total,orig}(\frac{D_{orig}}{D_{new}})^2`$, where $`D_{orig}`$ and $`D_{new}`$ are the original and increased distances to the galaxy, respectively.
(8) Add the appropriate sky level back into each image.
(9) Based on the assumption that the original data were sky-noise-limited, add photon noise according to a model derived from our blank sky region (see Step 2). This model insures that the magnitude of simulated noise is higher in bright parts of the image.
(10) Using the distant images created in Steps 1-9, perform continuum subtraction to compute the line only image.
For the present demonstration, we generated datasets corresponding to the appearance of M51 at distances of 15, 30, and 45 Mpc (having PSF FWHM $``$ 130, 260, and 390 pc respectively). Fig. 11 shows a comparison of the H II LFs obtained by running HIIphot on these degraded data. The dotted line traces the actual observed M51 LF, presented earlier in Fig. 10.
Two effects are rather striking. The completeness limit at low luminosities increases in a smooth but dramatic fashion. Moving from 9.6 Mpc to 15 Mpc, the rapid loss of faint, isolated point sources takes place and our incompleteness limit (in this case judged by the LF turnover) rises slightly faster than one might expect according to the inverse square law. This effect is mitigated as the galaxy gets even more distant. Perhaps blending allows a small fraction of adjacent weak sources to be recovered as single (brighter) objects. The second striking effect shown in Fig. 11 is the influence of blending on the slope of the LF. It is clear that the LFs tabulated for 30 and 45 Mpc have substantially shallower power-law slopes than the 9.6 Mpc LF over a limited range of luminosity. Indeed, the best-fit LF slope ranges from $`\alpha =1.75\pm 0.06`$ for the original data, to $`\alpha =1.22\pm 0.08`$ for the case of M51 at 45 Mpc (fitting only sources with log L$`{}_{H\alpha }{}^{}>38.0`$). This effect is brought about by blending of H II regions (as spatial resolution is degraded) with some help from catalog incompleteness. Blending also explains the increased apparent luminosity of the brightest H II regions as the galaxy becomes more distant, although this effect is not illustrated by Fig. 11 (due to the choice of bin size).
It is worth noting that above a limiting luminosity, the original and degraded H II LFs are essentially identical within the errors. For this example, the completeness tests of Section 4.3 imply that all versions of the M51 data in Fig. 11 are complete above log L $``$ 38.6. In the few histogram bins above this limit, minimal difference between the various LFs is apparent.
### 4.5 Comparison of arm & interarm regions
The results of R92 included a demonstration of changes in H II region properties for those sources located in interarm gaps. We can classify arm/interarm status based on masking of the diffuse background image produced by HIIphot. Using this technique we confirm the difference in LF slope observed by Rand for arm versus interarm H II region populations.
A simple way to designate H II regions as belonging in the arm or interarm populations relies on masking of the HIIphot surface fit to the diffuse emission remaining after definition of H II region boundaries. These images typically show very conspicuous spiral structure. We experimented with several isophotal cutoffs to obtain a boundary that closely resembled that of R92. Although the present goal was to see if we could develop a masking technique to efficiently reproduce the classification scheme of Rand (who carefully subdivided the entire sample of H II regions on the basis of spiral arm morphology), a more appropriate characterization of our new method would be one in which the isophotal mask is used to segregate regions on the basis of their local star forming environment. Under the assumption that more DIG is found in areas of enhanced recent star formation, the “arm” H II regions identified by our mask could be thought of as sources that lie within especially active star forming areas of the galaxy. In the end, a cutoff at an H$`\alpha `$ EM of 30 pc cm<sup>-6</sup> worked well in both contexts for M51 (with $`i=20\mathrm{°}`$).
Fig. 12 shows the arm and interarm LFs created using the mask described above. The straight lines are weighted power-law fits to the data for H II regions brighter than log L = 37.6 and 37.0, for arm and interarm respectively. Our simulations of the previous section suggest that the catalog of interarm sources is complete to lower luminosities than the general population. This was the motivation behind choosing different lower limits for arm and interarm power-law fits. We find that there is a difference in slope between the two populations. The spiral arm population is best-fit with $`\alpha =1.72\pm 0.06`$, whereas the interarm regions have a much steeper power-law slope given by $`\alpha =1.96\pm 0.15`$. The brightest H II regions are found almost exclusively within the spiral arms. Only two interarm H II regions in M51 are more luminous than L = 10<sup>38.1</sup> erg s<sup>-1</sup>.
### 4.6 Correlation between H$`\alpha `$ luminosity and H II region size
We find that there is a correlation between the H$`\alpha `$ luminosity of a region and its projected surface area (PSA). Fig. 13 shows a plot of log L versus log PSA. Our data is best-fit by a line of slope 1.71, substantially higher than the predicted value of 1.5 for a classical (radiation bounded) Stromgren sphere of constant density. The scatter about the fit is rather large, especially for small H II regions. Note that we have chosen to present the correlation between log L and log PSA, rather than log of H II region effective radius (r<sub>eff</sub>), because projected surface area is more directly related to our observations in the case of sources having irregular shape. We suspect that the slightly steeper than expected slope in Fig. 13 could be related to clumping within H II regions.
### 4.7 Characteristics of the DIG
The diffuse fraction in spiral galaxies remains of substantial interest for studies of ISM morphology and energetics. Defined as the ratio of DIG H$`\alpha `$ luminosity to total H$`\alpha `$ luminosity (Walterbos & Braun (1994)), the diffuse fraction has been estimated in a number of ways by different authors. The most common techniques are based on isophotal masking (e.g. Ferguson, Wyse, Gallagher & Hunter (1996), Hoopes, Walterbos & Greenwalt (1996), Wang, Heckman & Lehnert (1999)), although authors usually disagree on precise methodology. Classification of DIG has also been accomplished using explicit identification of traditional H II regions (Walterbos & Braun (1994)) and using maps of H$`\alpha `$ equivalent width (Veilleux, Cecil & Bland-Hawthorn (1995)). It is remarkable that the results obtained using diverse methods are quite similar, with a diffuse fraction of $`0.4\pm 0.1`$ being common for spiral galaxies (Greenawalt (1998)). Nevertheless, the diversity of methods employed makes it difficult to compare results in a detailed manner and accurately address relative uncertainties.
There are a few obvious drawbacks to each of the techniques previously used to estimate the diffuse fraction. In particular, inspection of the Fig. 2 in Ferguson, Wyse, Gallagher & Hunter (1996), Fig. 4b in Hoopes, Walterbos & Greenwalt (1996), and Fig. 8 in Greenawalt, Walterbos, Thilker & Hoopes (1998) reveals many instances of faint but highly localized H$`\alpha `$ emission being lumped into the DIG. These sources could be compact H II regions or even planetary nebulae. Attributing the flux of these faint discrete sources to DIG tends to artificially boost the diffuse fraction by a small (perhaps insignificant) amount. Secondly, several authors have pointed out that the total DIG luminosity should receive a contribution from locations in which H II regions are projected onto a slowly varying, diffuse background. The most commonly adopted solution is to assume that pixels occupied by an H II region each contribute the mean DIG intensity when totaling up DIG. This is undoubtedly an underestimate, as H II regions often have prominent DIG haloes, implying that the DIG superimposed on H II regions will typically be brighter than average. HIIphot addresses both of these problems, because it first measures flux associated with all discrete emission line sources and then individually estimates a background level for each region.
We calculate the diffuse fraction by independently totaling: (1) F<sub>HII</sub>, the background-corrected flux associated with all detected H II regions (except those with $`\frac{S}{N}_{final}<5`$) , and (2) F<sub>tot</sub>, the flux of the entire image. The diffuse fraction is then given by (F$`{}_{tot}{}^{}`$F<sub>HII</sub>)/F<sub>tot</sub>. This is the method of Walterbos & Braun (1994), but accomplished in a repeatable automated manner. By computing the diffuse fraction for various requested stopping-points during the iterative growth process, we can accurately constrain the diffuse fraction and also place an upper limit on the amount of DIG ionized in the field, apparently unrelated to classical H II regions.
Using region boundaries established by our nominal 1.5 EM/pc terminal surface brightness slope, we find that the diffuse fraction for M51 is $`0.45\pm 0.01`$. The uncertainty quoted here only accounts for the possibility of variation in the sky background. Other uncertainties such as those associated with continuum subtraction and growth termination criteria will also play a role, as will flat-fielding errors. In fact, as described below, these factors may actually dominate the diffuse fraction uncertainty.
Just over half of the observed H$`\alpha `$ emission from M51 can be unambiguously associated with classical H II regions. Fig. 14 presents portions of our HIIphot surface fit to control points located in the diffuse emission not overlapped by H II regions. In this plot, we have only shown the diffuse background surface fit for pixels covered by an H II region - all other areas show the original data. Panels (d) and (f) of Fig. 2 present the entire smoothly varying surface fit for comparison. Note that there is still substantial spatial correlation between areas of bright DIG and obvious concentrations of H II regions. Indeed, bright rims around a significant number of H II regions can be seen in Fig. 14. Taken together, these facts seem to imply that we have not yet recovered all the H$`\alpha `$ emission which is powered by Lyman continuum photon sources inside traditional H II regions. It is entirely possible that our 1.5 EM/pc growth limit remains too conservative and should really be lowered if we want to accurately characterize the massive stars in H II regions.
Figs. 2e and 5f represent the case in which H II regions were grown to encompass all apparently associated emission down to the sensitivity limit of the current data (by adopting a 1 EM/pc stopping point during the iterative procedure). Estimating the diffuse fraction with this set of region boundaries leads us to conclude that at most 38% of M51’s H$`\alpha `$ luminosity might originate via ionization by some mechanism other than leakage of Lyman continuum photons from H II regions. We do not mean to say that the conventional diffuse fraction is 0.38, but instead that a substantial fraction of the DIG emission in M51 cannot be plausibly tied to specific H II regions with the current data. This emission is still somewhat spatially correlated with the local density of H II regions, but ionization by “field” sources such as OB stars not in associations (Hoopes et al. 1999, in prep) may be largely responsible.
Our determination of the nominal diffuse fraction for M51 is clearly subject to systematic uncertainties related to our choice of the terminal surface brightness slope and uncertainty in the determination of the scale factor used during continuum-subtraction. Both of these can be empirically gauged. By computing the diffuse fraction immediately after growth commences (with a 10 EM/pc cutoff, see Fig. 5c), we obtain a hard upper limit of 0.68. At the very least, 32% of the H$`\alpha `$ emission from M51 is contained within the cores of classical H II regions. Systematic changes related to error in continuum-subtraction are easily measured by producing new versions of the line-only image then recomputing the background surface-fit and H II region fluxes. We generated “test” H$`\alpha `$ images by varying the continuum-subtraction scale factor $`\pm `$ 3% (1 $`\sigma `$) from our best-guess value. In the case of 1.5 EM/pc nominal growth boundaries, this resulted in diffuse fractions of 0.40 and 0.49, respectively for increased and decreased continuum emission.
We note in passing that the diffuse fraction is also potentially influenced by extinction variations across the face of a galaxy. The optical depth towards H II regions is probably elevated with respect to field DIG. Unfortunately, correcting for this systematic error would be rather difficult even in the case of measured Balmer decrements, given the unavoidable uncertainty in the geometry of emitting and absorbing volumes.
## 5 Summary
We have developed a new IDL procedure, which we have designated HIIphot, which is capable of performing fully-automated, repeatable photometry of H II regions. The procedure can detect and accurately characterize faint sources embedded in crowded fields, even in the presence of a substantially inhomogeneous, diffuse H$`\alpha `$ background.
In this paper we have applied HIIphot to the analysis of the grand-design spiral M51, studied previously by R92 and KEH89. Our results are in general agreement with these authors, although we detect more than twice the number of H II regions described by R92. In total, we find 1229 sources above $`5\sigma `$ having luminosity greater than about 10<sup>36.1</sup> erg s<sup>-1</sup>. The LF obtained from this catalog of $`5\sigma `$ sources is reasonably well fit by a power law distribution having $`\alpha `$ = 1.75 $`\pm `$ 0.06, below a break in the observed number of regions near log L = 38.9. This break confirms the earlier classification of M51 as exhibiting a Type II LF.
In the near future, we plan to apply HIIphot for the analysis of an extensive galaxy sample for which high-quality, sensitive narrowband observations already exist. The sample will contain substantially more galaxies than observed by KEH89. Given the HIIphot code, it will be trivial to “reobserve” each of the galaxies at a common distance in order to inter-compare LFs in the absence of bias associated with different degrees of blending due to limited resolution and sensitivity. As a predecessor to this large study, we present HIIphot results for a smaller sample of 11 spirals in Thilker et al. (2000, Paper II). Therein we also develop a procedure for fitting HII LFs with predictions from population synthesis models of star cluster formation and evolution.
DAT gratefully acknowledges the support and encouragement of RB and RW, his dissertation advisors. DAT further acknowledges the congenial staff of NFRA for their hospitality during many collaborative trips to work with RB. Veronica Fierro has also been of great help to the authors, finding bugs in our code via repeated trial-and-error throughout the development stage of HIIphot. DAT has been funded through the NASA Graduate Student Researcher Program (NGT-51640) and by NSF grant AST9617014 to RAMW. The HIIphot IDL source code and explanatory documentation will soon be available by request from DAT. Contact dthilker@nrao.edu for details.
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# 1 Acknowledgements
## 1 Acknowledgements
This work is supported by grants provided by Fundação para a Ciência e a Tecnologia, PRAXIS/C/FIS/12247/1998, PESO/P/PRO/15127/1999 and NATO ”Outreach” Cooperation Program.
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# Notes on Infinite Layer Quantum Hall Systems
## I Introduction
The quantum Hall effect arises from physics that is special to two dimensions. Nevertheless, it was demonstrated in an early experiment that it survives a small amount of three dimensionality; more precisely, that an infinite layer system with interlayer tunneling weak compared to the single layer (mobility) gap would continue to exhibit dissipationless transport and a quantized Hall resistance per layer. The quantized Hall phases in the organics also rely on the same effect. In the presence of tunneling the chiral edge states which exist in each layer hybridize and yield a family of “one and a half” dimensional phases that live on the surfaces of the three dimensional systems and exhibit interesting transport in the direction transverse to the layers.
Following the pioneering theoretical work, a flurry of work has focused on the integer Hall effect in infinite layers systems with a particular emphasis on the properties of the chiral metal that forms at their surface in the quantum Hall phases and some of this has found support in experiments. Recently, interaction effects and magnetic order at $`\nu =1`$ (per layer) have been discussed as well. All of this work has antecedents in work on multi-component systems such as double layer devices, but the limit of infinitely many layers sharpens quantitative differences into qualitatively new features, e.g., a different universality class for the transitions out of the quantum Hall phases.
In this paper we extend the study of three dimensional quantum Hall phases in two directions. First, we offer a systematic analysis of the simplest fractional quantum Hall phases in infinite layer systems with a special focus on the phases that exhibit interlayer correlations. Here we find several novel bulk properties, most notably a non-trivial structure for the quasiparticles. This part of our work builds on the pioneering and early work of Qiu, Joynt and MacDonald (QJM), who studied the energetics of multilayer states, and noted the irrational partitioning of the quasiparticle charge. The second axis of our work is the analysis of the “one and a half” dimensional phases that arise at the surfaces of fractional quantum Hall phases. In this regard we report a general analysis of the surface conduction by combining renormalization group arguments (which have strong consequences for the ground state structure in chiral systems) and a conductivity sum rule. We also carry out computations of the surface conductivity of the disordered system in a weak tunneling expansion for a large class of states. Together, these show that the surfaces of states with marginal or irrelevant electron tunneling are semi-metallic, in that they are perfect insulators at all frequencies in the clean limit but begin to conduct at finite temperature and disorder. Some of the results derived here were announced previously in a companion paper.
This last theme, of analyzing the surface phases, is interesting from a more theoretical standpoint in that the surface phases are “half” a dimension up from one dimension where exact solutions are often possible. We have not made very much progress on this front. What we have accomplished is to solve the case of bilayers in some generality , find special solutions for three and four layers, and recast the infinite layer problems in algebraic and non-standard fermionized formulations that appear potentially fruitful. We report these here in the hope that some readers will find them useful in carrying this program to completion. We should also note that the construction of more complex bulk states than those considered in this paper could well lead to a richer class of such problems.
The outline of this paper is as follows. We begin in Section II with a discussion of the bulk properties of multilayer fractional states, in particular the novel features of their quasihole excitations which include a non-trivial internal charge structure and irrational statistics. In Section III we discuss the edge theory of general multilayer states in the presence of nearest-neighbor single-electron tunneling. We specialize to the case where the tunneling is marginal in Section IV, first considering the clean limit and then adding disorder and interactions. In Section V we extend our analysis to states where tunneling is irrelevant. We conclude in Section VI by presenting a method for fermionizing the multilayer edge theory via the addition of auxiliary degrees of freedom. The appendices contain a discussion of an exactly soluble model problem with marginal tunneling (App. A), exact solutions for the spectra of two edge theories for the special case of four layers (App. B), and some technical details including Klein factors (App. C), and the proof of an assertion made in the main text (App. D).
## II Bulk Properties
Consider a system which consists of $`N`$ parallel layers of 2DEGs in a strong perpendicular magnetic field. We assume there is a confining potential which restricts the electrons in each layer to a region with the topology of a disc. The phase diagram of this system was first investigated by Qiu, Joynt, and MacDonald (QJM) . At low electron densities they found a variety of solid phases, while at higher densities they found the ground state to be an incompressible fluid described by the $`N`$-layer generalization of the Laughlin wavefunction:
$$\mathrm{\Psi }_0(\{z_{i,\alpha }\})=\underset{i=1}{\overset{N}{}}\underset{\alpha <\beta }{\overset{N_i}{}}(z_{i\alpha }z_{i\beta })^{K_{ii}}\underset{i<j}{\overset{N}{}}\underset{\alpha =1}{\overset{N_i}{}}\underset{\beta =1}{\overset{N_j}{}}(z_{i\alpha }z_{j\beta })^{K_{ij}}e^{_{\alpha ,i}|z_{\alpha ,i}|^2/4}.$$
(1)
Here $`z_{i\alpha }`$ is the coordinate of electron $`\alpha `$ in layer $`i`$, and $`N_i`$ is the number of electrons in layer $`i`$. The exponents are specified by a symmetric, $`N\times N`$ matrix $`K`$. The diagonal elements of this matrix determine the electron correlations within each layer and the off-diagonal elements specify the correlations between layers. For the wavefunction to describe electrons the diagonal elements of $`K`$ must be odd integers. The filling factor in layer $`i`$ is given by
$$\nu _i=\underset{j=1}{\overset{N}{}}(K^1)_{ij}.$$
(2)
Note that it is possible to construct more complicated multilayer states by beginning with the states in Eq. (1) and carrying out a generalization of the single-layer Haldane-Halperin hierarchy construction. For example, at the first level of the hierarchy the effective $`K`$ matrix is given by
$$K_{(1)}=K+AP^1,$$
(3)
where $`A`$ is a diagonal $`N\times N`$ matrix with elements $`\pm 1`$, and $`P`$ is a symmetric $`N\times N`$ matrix with even integer elements along the diagonal and integer elements everywhere else. The sign of the element $`A_{jj}`$ determines whether the excitations in layer $`j`$ are quasiholes (+1) or quasielectrons ($`1`$) and the matrix $`P`$ specifies the (bosonic) multilayer state into which these excitations condense. Using the matrix $`K_{(1)}`$ in Eq. (2) gives the filling factors at the first level of the hierarchy.
We will restrict ourselves to the case where the matrix $`K`$ is tridiagonal
$$K_{ij}=m\delta _{ij}+n(\delta _{i,j1}+\delta _{i,j+1}),$$
(4)
where $`m`$ and $`n`$ are nonnegative integers. Since we are interested in the limit of large $`N`$, we will impose periodic boundary conditions along the direction perpendicular to the layers by the identification $`N+11`$. One of the difficulties encountered when working without periodic boundary conditions is discussed briefly below. The standard convention is to refer to the state whose $`K`$ matrix is of the form (4) as the “$`nmn`$” state. Using Eqns. (2) and (4), the filling factor per layer in the $`nmn`$ state is
$$\nu =\frac{1}{m+2n}.$$
(5)
From this expression we see that there are multiple states with the same filling factor per layer. For example, the states 050 and 131 both have filling factor $`\nu =1/5`$ per layer. QJM found that as the interlayer separation $`d`$ was decreased there is a phase transition between the 050 state in which there are no interlayer correlations and the 131 state in which electrons in neighboring layers are correlated.
Note that Eq. (5) applies to the case $`N\mathrm{}`$ or to the case of finite $`N`$ with periodic boundary conditions. For a physically realizable multilayer system, i.e., one with $`N`$ finite but without periodic boundary conditions, one finds that the filling factor $`\nu _j`$ is not independent of the layer index $`j`$. If one solves Eq. (2) for finite $`N`$ without periodic boundary conditions, using the $`K`$ matrix given in Eq. (4) one finds
$$\nu _j=\frac{1}{m+2n}\left[1\frac{y^j+y^{N+1j}}{1+y^{N+1}}\right],$$
(6)
where
$$y=\frac{m}{2n}+\sqrt{\left(\frac{m}{2n}\right)^21}.$$
(7)
From this form we find that for large $`N`$ the filling factor near the middle of the multilayer stack is given approximately by Eq. (5), but exhibits oscillations about this value as one approaches the end layers. Such non-uniformities in the electron density would cost electrostatic energy, and could be mitigated by the formation of quasihole excitations near the outermost layers. Therefore, one would expect that the true ground state for a finite multilayer with interlayer correlations would be determined by balancing the creation energy of quasiholes against the electrostatic energy of the density oscillations. It is interesting to note that to construct a state in a finite stack with uniform filling by carrying out the hierarchy construction in only the outermost layers one must go to an infinite level in the hierarchy. Henceforth we will avoid these complications by working with periodic boundary conditions, arguing that in the limit of large $`N`$ they can safely be ignored.
### A Quasihole Excitations
In the last section we learned that for a range of electron densities the ground state of the multilayer system is an incompressible state described by the wavefunction (1). It is natural to ask about the properties of the excitations of this state. By analogy with the single-layer quantum Hall effect we write the wavefunction for a single quasihole at point $`\xi `$ in layer $`j`$ as:
$$\mathrm{\Psi }_{(j,\xi )}(\{z_{i\alpha }\})=\underset{\alpha =1}{\overset{N_j}{}}(z_{j\alpha }\xi )\mathrm{\Psi }_0(\{z_{i\alpha }\}).$$
(8)
If one interprets $`|\mathrm{\Psi }_{(j,\xi )}|^2e^{\beta V(\{z\})}`$ as the Boltzmann factor for a classical 2D generalized Coulomb plasma at inverse temperature $`\beta `$ with an impurity at $`\xi `$, the perfect screening conditions give:
$$K_{ik}Q_k^{(j)}=\delta _{ij},$$
(9)
where $`Q_k^{(j)}=(K^1)_{kj}`$ is the local deviation of the charge in layer $`k`$ from its ground state value, due to the quasihole in layer $`j`$. The total charge of the quasihole is
$$Q=\underset{k}{}Q_k^{(j)}=\underset{k}{}(K^1)_{jk}=\nu _j=\frac{1}{m+2n},$$
(10)
where we have used Eq. (2). Thus, the relation between the total quasihole charge and the filling factor is the same as in the single-layer quantum Hall effect. However, if there are interlayer correlations, i.e., $`n0`$, the quasihole charge is spread over many layers. Indeed, in the limit $`N\mathrm{}`$ the individual charges are irrational:
$$Q_k^{(j)}=\frac{1}{\sqrt{m^24n^2}}\left(\frac{\sqrt{m^24n^2}m}{2n}\right)^{|kj|}.$$
(11)
The same Coulomb plasma analogy, combined with the mean field approximation, allows us to find the spatial distribution of the extra charge excited by the hole. The Debye screening equation for an average “potential” $`V_i^{(j,\xi )}(z)`$ acting on an electron at the point $`z`$ in layer $`i`$ can be written as
$$\beta ^2V_i^{(j,\xi )}=4\pi \delta _{ij}\delta (z\xi )2+4\pi \underset{k}{}K_{ik}n_k,$$
(12)
where the single-particle average charge density
$$n_i=\left(2\pi _jK_{ij}\right)^1e^{\beta V_i}$$
(13)
is chosen so that at $`V_i=0`$ the unperturbed density would be restored. Clearly, the screening in each successive layer is limited by the Debye screening length which is on the order of the magnetic length ($`l=1`$ in chosen units). Therefore, as the amount of charge induced by the hole is reduced with the separation along the stack, this charge also spreads over a wider and wider area, further reducing the charge density. The shape of the induced charge distribution can be found by linearizing the density (13) in Eq. (12) and solving the resulting set of linear partial differential equations by Fourier transformation. In particular, for the r.m.s. spread of the distribution created in layer $`j`$ by the quasihole in layer $`i`$ we obtain
$$r_{|ij|}^2=2I_2(|ji|)/I_1(|ji|),\mathrm{where}I_p(j)_0^{2\pi }𝑑q\frac{\mathrm{cos}(qj)}{K_q^p},$$
(14)
and $`K_q=m+2n\mathrm{cos}q`$ is the Fourier transform of Eq. (4). Specifically, for the 131 state we find $`r_j^2=4\left(3+\sqrt{5}|j|\right)/5`$. If we were to color in the perturbed charge density, the hole would produce a characteristic “hourglass” shape.
To compute the statistics of these excitations we follow the method of Arovas, Schrieffer, and Wilczek and consider a state with two quasiholes, one at point $`\xi `$ in layer $`j`$ and one at point $`\eta `$ in layer $`k`$:
$$\mathrm{\Psi }_{(j,\xi ;k,\eta )}=\underset{\alpha =1}{\overset{N_j}{}}(z_{j\alpha }\xi )\underset{\beta =1}{\overset{N_k}{}}(z_{k\beta }\eta )\mathrm{\Psi }_0.$$
(15)
The Berry phase for moving the quasihole at $`(k,\eta )`$ around a closed loop of radius $`R`$ containing the quasihole at $`(j,\xi )`$ is
$`\gamma _B`$ $`=`$ $`i{\displaystyle _\xi }𝑑\eta \mathrm{\Psi }_{(j,\xi ;k,\eta )}|_\eta \mathrm{\Psi }_{(j,\xi ;k,\eta )}`$ (16)
$`=`$ $`i{\displaystyle _\xi }𝑑\eta {\displaystyle d^2z\frac{1}{\eta z}\mathrm{\Psi }_{(j,\xi ;k,\eta )}|\underset{\beta =1}{\overset{N_k}{}}\delta ^{(2)}(zz_{k\beta })|\mathrm{\Psi }_{(j,\xi ;k,\eta )}}.`$ (17)
The expectation value appearing in this equation is just the electron density at the point $`z`$ in layer $`k`$ for the two-quasihole state. We can write this as
$$\mathrm{\Psi }_{(j,\xi ;k,\eta )}|\underset{\beta =1}{\overset{N_k}{}}\delta ^{(2)}(zz_{k\beta })|\mathrm{\Psi }_{(j,\xi ;k,\eta )}=\rho _k^{(0)}(z)+\delta \rho _k^{(k,\eta )}(z)+\delta \rho _k^{(j,\xi )}(z),$$
(18)
where $`\rho _k^{(0)}(z)`$ is the density in the ground state, $`\delta \rho _k^{(k,\eta )}(z)`$ is the change in density due to the quasihole at $`(k,\eta )`$ and $`\delta \rho _k^{(j,\xi )}(z)`$ is the change in density due to the quasihole at $`(j,\xi )`$. If we substitute Eq. (18) into Eq. (17), the $`\rho _k^{(0)}`$ term gives the Aharonov-Bohm phase, the $`\delta \rho _k^{(k,\eta )}`$ term gives zero by symmetry, and hence
$$\gamma _B^{(\mathrm{stat})}=id^2z\delta \rho _k^{(j,\xi )}(z)_\xi 𝑑\eta \frac{1}{\eta z}=2\pi Q_k^{(j)},$$
(19)
where once again $`Q_k^{(j)}`$ is the charge deviation in layer $`k`$ due to a quasihole in layer $`j`$, which we found above to be equal to $`(K^1)_{kj}`$, see Eq. (9). The statistical angle for interchanging two quasiholes in the same layer is
$$\mathrm{\Delta }\gamma =\frac{1}{2}\gamma _B^{(\mathrm{stat})}=\pi (K^1)_{jj}=\frac{\pi }{\sqrt{m^24n^2}}.$$
(20)
We find the statistical angle of the quasiholes is an irrational multiple of $`\pi `$ in the infinite layer system.
## III Edge Theory
In this section we consider the edge theory of the multilayer $`nmn`$ state in the presence of interlayer single-electron tunneling. The low-energy effective Hamiltonian of the edge theory in the absence of tunneling is:
$$_0=_{L/2}^{L/2}𝑑x\frac{1}{4\pi }V_{ij}:_xu_i_xu_j:,$$
(21)
where $`V`$ is a symmetric, positive definite, $`N\times N`$ matrix which depends on the interactions and confining potentials at the edge, and we take the $`x`$-axis along the edges of length $`L`$. The $`N`$ chiral bosons appearing in this Hamiltonian obey the equal-time commutation relations
$$[u_i(x),u_j(x^{})]=i\pi K_{ij}\mathrm{sgn}(xx^{}).$$
(22)
The electron charge density and the electron creation operator in layer $`i`$ are given by
$$\rho _i(x)=\frac{1}{2\pi }(K^1)_{ij}_xu_j(x),\mathrm{\Psi }_i^{}(x)e^{iu_i(x)}.$$
(23)
Note that because of the factor of $`K^1`$ in the expression for the density operator (23) there is a non-trivial relation between the matrix $`V`$ appearing in the Hamiltonian and the matrix determining the density-density couplings, which we will call $`U`$. If we rewrite the Hamiltonian (21) as
$$_0=_{L/2}^{L/2}𝑑xU_{ij}:\rho _i(x)\rho _j(x):,$$
(24)
which serves as our definition of $`U`$, we find by equating these two forms that
$$V=\frac{1}{\pi }K^1UK^1.$$
(25)
We assume that the system is translationally invariant along the direction perpendicular to the layers, which we shall take to be the $`z`$-axis. We shall often refer to this as the “vertical” direction. Specifically, we assume the matrices $`K_{ij}`$ and $`U_{ij}`$ depend only on the difference $`|ij|`$. The tridiagonal form of $`K`$ given above in Eq. (4) certainly obeys this constraint, provided we recall that we are assuming periodic boundary conditions along $`z`$. From Eq. (25) we see that if both $`K`$ and $`U`$ are translationally invariant, so is $`V`$. Hence to diagonalize $`K`$ and $`V`$ we perform a discrete Fourier transformation in the $`z`$ direction, defining:
$$u_q(x)=\frac{1}{\sqrt{N}}\underset{j=1}{\overset{N}{}}e^{iqj}u_j(x),$$
(26)
where
$$q\left\{2\pi n/N\right|n=0,1,\mathrm{},N1\},$$
(27)
and we identify $`q`$ and $`q+2\pi k`$ for any integer $`k`$.
Using transformation (26), the commutation relations (22) become
$$[u_q(x),u_q^{}(x^{})]=i\pi \delta _{q+q^{},0}K_q\mathrm{sgn}(xx^{}),$$
(28)
where
$$K_qm+2n\mathrm{cos}q,$$
(29)
are the eigenvalues of the matrix (4). We will largely restrict our analysis to those states for which all the edge modes move in the same direction, i.e., the maximally chiral case. This requires that $`K`$ be positive definite. From Eqns. (27) and (29) this implies $`n<m/2`$.
Finally, we rescale the fields, defining
$$\varphi _q(x)=\frac{1}{\sqrt{K_q}}u_q(x),$$
(30)
which have conventionally normalized commutation relations
$$[\varphi _q(x),\varphi _q^{}(x^{})]=i\pi \delta _{q+q^{},0}\mathrm{sgn}(xx^{}).$$
(31)
Using the transformations (26) and (30) to express the Hamiltonian (21) in terms of the fields $`\varphi _q`$ we find
$$_0=_{L/2}^{L/2}𝑑x\frac{1}{4\pi }v_q:_x\varphi _q_x\varphi _q:.$$
(32)
The velocities $`v_q`$ are given by $`v_q=V_qK_q`$ where $`V_q`$ are the eigenvalues of the $`V`$ matrix. We now specialize to the case where the density-density coupling matrix is of the form
$$U_{ij}=\pi v\delta _{ij}+\pi g(\delta _{i,j1}+\delta _{i,j+1}),$$
(33)
where $`v`$ is a velocity determined by the confining potential at the edge and $`g`$ parameterizes the nearest-neighbor density-density interactions between the layers. Using the relation between $`U`$ and $`V`$ (25) we find that the eigenmode velocities appearing in the Hamiltonian (32) are
$$v_q=\frac{v+2g\mathrm{cos}q}{m+2n\mathrm{cos}q}.$$
(34)
Next we consider adding interlayer single-electron tunneling to the multilayer edge theory. The electron annihilation operator in layer $`i`$ is given in Eq. (23) in terms of the original bosonic field $`u_i(x)`$. Using the transformations (26) and (30) the tunneling operator between layers $`j`$ and $`j+1`$ can be written
$$\lambda \mathrm{\Psi }_j(x)\mathrm{\Psi }_{j+1}^{}(x)+\mathrm{h}.\mathrm{c}.\lambda e^{i\alpha _{jq}\varphi _q(x)}+\mathrm{h}.\mathrm{c}.,$$
(35)
where $`\lambda `$ is the tunneling amplitude and we have defined the coefficients
$$\alpha _{jq}\frac{1}{\sqrt{N}}e^{iqj}(1e^{iq})\sqrt{K_q}.$$
(36)
As defined in Eq. (23) the electron operators in different layers do not obey proper anticommutation relations. In Appendix C we demonstrate that this can be remedied without significantly altering the form of the tunneling operator (35).
The lowest-order perturbative RG flow of the tunneling amplitude $`\lambda `$, assumed to be uniform along the edge, is controlled by the scaling dimension, $`\delta _\lambda `$, of the tunneling operator (35):
$$\frac{d\lambda }{d\mathrm{}}=\left(2\delta _\lambda \right)\lambda ,$$
(37)
where the short-distance cutoff increases as $`\mathrm{}`$ increases. Since the fields $`\varphi _q`$ obey canonically normalized commutation relations, (31), we can read off the scaling dimension of the tunneling operator from Eqns. (29), (35), and (36):
$$\delta _\lambda =\frac{1}{2}\underset{q}{}\alpha _{jq}\alpha _{jq}=mn,$$
(38)
where there is no sum on $`j`$. Using Eqns. (37) and (38) we see that tunneling is relevant for $`m<n+2`$, irrelevant for $`m>n+2`$, and marginal for $`m=n+2`$.
If we combine this result with the condition of maximum chirality, $`n<m/2`$, we obtain the diagram in Fig. 1. The only maximally-chiral state for which the tunneling is relevant is 010, i.e., uncorrelated $`\nu =1`$ layers. There are two maximally-chiral states for which tunneling is marginal: 131 and the bosonic state 020. For all other maximally-chiral multilayer states interlayer electron tunneling is irrelevant.
There are two bosonic states, 121 with relevant tunneling and 242 with marginal tunneling, which lie on the boundary of the maximally-chiral region. The corresponding $`K`$ matrices have zero eigenvalues in the limit of an infinite number of layers. The complication this zero eigenvalue adds is the possibility of “soft” modes which are not strictly confined to the edge. The 121 state is the first of a sequence of “pyramid” states: 121, 12321, 1234321,$`\mathrm{}`$, where in an obvious extension of the $`nmn`$ notation the state $`onmno`$ has next-nearest-neighbor correlations specified by the exponent $`o`$, in addition to nearest-neighbor ($`n`$) and intralayer correlations ($`m`$), and similarly for the $`ponmnop`$ states, etc. All of the pyramid states lie on the boundary of the maximally chiral region (i.e., their $`K`$ matrices are positive semi-definite) and have relevant tunneling. Were it not for the complication from the soft modes, their edge theories would be exactly soluble by fermionization as discussed in Section VI.
The full Hamiltonian of the multilayer edge theory with tunneling can be written using Eqns. (32), (35), and (38):
$$=_{L/2}^{L/2}dx[\frac{1}{4\pi }v_q:_x\varphi _q_x\varphi _q:+\underset{j}{}\frac{\lambda }{(2\pi a)^{\delta _\lambda }}(e^{i\alpha _{jq}\varphi _q(x)}+\mathrm{h}.\mathrm{c}.)],$$
(39)
where $`a`$ is a short distance cutoff. In the following sections we will be concerned with the $`z`$-axis conductivity of this theory, so we will need the form of the current operator along $`z`$. If $`\rho _j^{(2D)}(x,t)`$ is the two-dimensional charge density operator in layer $`j`$, the continuity equation can be written
$$\frac{}{t}\rho _j^{(2D)}(x,t)=\frac{1}{d}\left(𝒥_{j1,j}^z(x,t)𝒥_{j,j+1}^z(x,t)\right)_x𝒥_j^x(x,t),$$
(40)
where $`𝒥_j^x`$ is the current density along the edge of layer $`j`$, $`𝒥_{j,l}^z`$ is the current density flowing from layer $`j`$ to layer $`l`$, and $`d`$ is the layer separation. Introducing the discrete Fourier transforms
$$\rho ^{(2D)}(k,q,t)=d\underset{j}{}𝑑xe^{iqj}e^{ikx}\rho _j^{(2D)}(x,t),$$
(41)
with an analogous definition for $`𝒥^x(k,q,t)`$, and
$$𝒥^z(k,q,t)=d\underset{j}{}𝑑xe^{iq(j1/2)}e^{ikx}𝒥_{j1,j}^z(x,t),$$
(42)
where the allowed values of the dimensionless $`z`$-momentum $`q`$ are given in Eq. (27), the continuity equation (40) becomes
$$\frac{}{t}\rho ^{(2D)}(k,q,t)=\frac{2i}{d}\mathrm{sin}(q/2)𝒥^z(k,q,t)ik𝒥^x(k,q,t).$$
(43)
Note $`\rho _j^{(2D)}(x,t)(1/d)\rho _j(x,t)`$, where $`\rho _j`$ is given in Eq. (23). Therefore, at $`k=0`$ we can write Eq. (43) as
$$𝒥^z(k=0,q,t)=\frac{d}{2\mathrm{sin}(q/2)}[H,\rho ^{(2D)}(k=0,q,t)].$$
(44)
Evaluating the commutator with the help of Eqns. (22), (23), and (35), we find
$$I^z(t)𝒥^z(k=0,q=0,t)=\frac{i\lambda d}{(2\pi a)^{\delta _\lambda }}\underset{j}{}_{L/2}^{L/2}dx[e^{i\alpha _{jq^{}}\varphi _q^{}(x,t)}\mathrm{h}.\mathrm{c}.].$$
(45)
This expression will be used to calculate the $`z`$-axis conductivity of the multilayer system.
## IV Marginal Tunneling Cases
We know there are only two maximally-chiral states with marginal tunneling: 131 and the bosonic state 020. In this section we begin by showing that in the absence of disorder these states exhibit insulating behavior in the vertical direction, i.e., $`\sigma ^{zz}(\omega )`$ is identically zero for all frequencies at $`T=0`$. We will then find that adding disorder to the edge theory increases the vertical conductivity, and therefore the surface behaves like a “chiral semi-metal”. We conclude this section by considering how the transport is modified by nearest-neighbor density-density interactions.
### A Clean Limit
Our discussion of the system in the absence of disorder will proceed in two steps. First, we derive a sum rule which relates the frequency integral of the conductivity to the expectation value of the tunneling term in the Hamiltonian. Next we prove that there exists a finite range of values for the tunneling amplitude $`\lambda `$, and the interaction strength $`g`$, for which the ground state in the presence of tunneling and interactions is identical to the ground state in the absence of tunneling and interactions. We then combine these two results to arrive at our conclusions about the zero temperature conductivity. We emphasize that this argument is non-perturbative in $`\lambda `$ and $`g`$. This is an important point since we do not have a quadratic form of the Hamiltonian for the multilayer state with marginal tunneling.
The conductivity sum rule is derived by considering the double commutator of the Hamiltonian with the charge density operator. Specifically, we calculate:
$$[[,\rho ^{(2D)}(k=0,q,t)],\rho ^{(2D)}(k=0,q,t)],$$
(46)
where the Hamiltonian is given in Eq. (39), and from Eqns. (23), (26), (30), (41), and $`\rho _j^{(2D)}(x,t)=(1/d)\rho _j(x,t)`$ we have
$$\rho ^{(2D)}(k=0,q,t)=\frac{1}{2\pi }\sqrt{N}_{L/2}^{L/2}𝑑x\frac{1}{\sqrt{K_q}}_x\varphi _q(x,t).$$
(47)
Using the commutation relations for the fields $`\varphi _q(x)`$ (31) it is a straightforward exercise to evaluate the commutator (46) with the help of Eqns. (39) and (47). One finds:
$$[[,\rho ^{(2D)}(k=0,q,t)],\rho ^{(2D)}(k=0,q,t)]=4\mathrm{sin}^2(q/2)_\lambda ,$$
(48)
where $`_\lambda `$ is the tunneling part of the Hamiltonian.
We next relate the expectation value of this double commutator to the integrated conductivity. The continuity equation (40) introduced in the previous section can be written
$$\rho ^{(2D)}(k=0,q,t)=\frac{2}{d}\mathrm{sin}(q/2)\frac{d\omega }{2\pi }𝑑t^{}\frac{e^{i\omega (tt^{})}}{\omega }𝒥^z(k=0,q,t).$$
(49)
Suppressing the argument $`k=0`$, using $`[,\rho ]=i\rho /t`$, and taking the expectation value of Eq. (49) we find
$$[[,\rho _j^{(2D)}(q,t)],\rho _j^{(2D)}(q,t)]=\frac{4}{d^2}\mathrm{sin}^2(q/2)\frac{d\omega }{2\pi }𝑑te^{i\omega t}[𝒥^z(q,0),𝒥^z(q,t)].$$
(50)
If we take the expectation value of Eq. (48), and use it to eliminate the l.h.s. of Eq. (50) we find in the limit $`q0`$:
$$_\lambda =\frac{4}{d^2}\frac{d\omega }{2\pi }𝑑te^{i\omega t}[I^z(0),I^z(t)],$$
(51)
where we have used $`I^z(t)=𝒥^z(q=0,t)`$. The Kubo formula relates the commutator appearing in this expression to the $`z`$-axis conductivity:
$$\mathrm{}e\sigma ^{zz}(\omega )=\frac{1}{2\omega }\frac{1}{NLd}𝑑te^{i\omega t}[I^z(t),I^z(0)].$$
(52)
Thus we arrive at the final form of the sum rule
$$𝑑\omega \mathrm{}e\sigma ^{zz}(\omega )=\frac{\pi d}{NL}_\lambda .$$
(53)
This is an exact relation valid at any temperature. It applies to all multilayer states, even in the presence of disorder, regardless of the relevancy of the tunneling term.
Next we turn our attention to proving the stability of the ground state in the presence of interlayer tunneling and interactions. We will describe in detail the case of the bosonic 020 state. A proof along the same lines can be constructed for the 131 state. We begin by writing the Hamiltonian of the 020 edge theory in terms of the original radius $`R=1`$ Bose fields $`u_j(x)`$:
$$_{020}=_{L/2}^{L/2}dx\underset{j}{}[\frac{v}{16\pi }:(_xu_j)^2:+\frac{g}{8\pi }_xu_j_xu_{j+1}+\frac{\lambda }{(2\pi a)^2}(e^{i(u_ju_{j+1})}+\mathrm{h}.\mathrm{c}.)].$$
(54)
Instead of performing the discrete Fourier transformation along $`z`$ (26) we simply rescale the fields, defining $`u_j(x)=\sqrt{2}\phi _j(x)`$, where $`\phi _j(x)`$ are radius $`R=1/\sqrt{2}`$ chiral bosons. We can then define a set of $`N`$ $`\widehat{su}(2)_1`$ Kac-Moody (KM) currents:
$$J_j^z(x)=\frac{1}{2\pi \sqrt{2}}_x\phi _j(x),J_j^\pm (x)=J_j^x\pm iJ_j^y=\frac{1}{2\pi a}e^{i\sqrt{2}\phi _j(x)},$$
(55)
in terms of which the Hamiltonian (54) can be written
$`_{020}={\displaystyle _{L/2}^{L/2}}dx[{\displaystyle \frac{\pi v}{6}}:\left[𝐉_j(x)\right]^2:+\pi gJ_j^z(x)J_{j+1}^z(x)`$ (56)
$`+\lambda (J_j^{}(x)J_{j+1}^+(x)+J_{j+1}^{}(x)J_j^+(x))],`$ (57)
where we have employed the identity
$$_{L/2}^{L/2}dx:\left[J^z(x)\right]^2:=_{L/2}^{L/2}dx\frac{1}{3}:\left[𝐉(x)\right]^2:.$$
(58)
Consider the theory without interactions, $`g=0`$, and without tunneling $`\lambda =0`$. The zero-energy ground state of this Hamiltonian, $`|0`$, is formed by taking the tensor product of the highest-weight state of the spin-0 irreducible representation of each $`\widehat{su}(2)_1`$ algebra. Thus, the ground state satisfies
$$J_{j,n}^a|0=0,\mathrm{for}j;a=x,y,z;n0,$$
(59)
where the Fourier components of the currents are defined by
$$J_{j,n}^a_{L/2}^{L/2}𝑑xJ_j^a(x)e^{2\pi inx/L},n\mathrm{𝖹𝖹},$$
(60)
and obey the algebra
$$[J_{i,n}^a,J_{j,m}^b]=\frac{n}{2}\delta _{ij}\delta ^{ab}\delta _{n+m,0}+i\delta _{ij}ϵ^{abc}J_{j,n+m}^c.$$
(61)
Using Eqns. (57) and (60) we can write the Hamiltonian as
$$_{020}=[\frac{\pi v}{6L}:J_{j,n}^aJ_{j,n}^a:+\frac{\pi g}{L}J_{j,n}^zJ_{j+1,n}^z+\frac{\lambda }{L}(J_{j,n}^{}J_{j+1,n}^++J_{j+1,n}^{}J_{j,n}^+)].$$
(62)
For later reference we also record the expression for the current operator in terms of the KM generators
$$I^z=\frac{i\lambda d}{L}\left[J_{j,n}^{}J_{j+1,n}^+J_{j+1,n}^{}J_{j,n}^+\right].$$
(63)
From Eqns. (59) and (62) we see that not only is $`|0`$ the zero-energy ground state of the Hamiltonian with $`g=\lambda =0`$, it is a zero-energy eigenstate of the Hamiltonian for all $`g`$ and $`\lambda `$:
$$_{020}|0=0.$$
(64)
Of course, this does not mean that $`|0`$ is the ground state of the Hamiltonian with non-zero $`g`$ and $`\lambda `$. To establish this we now demonstrate that for a range of $`g`$ and $`\lambda `$ the Hamiltonian is positive semi-definite. For $`\lambda >0`$, $`g>0`$, we rewrite Eq. (57) as
$`_{020}`$ $`=`$ $`{\displaystyle _{L/2}^{L/2}}dx\{{\displaystyle \frac{\pi g}{2}}:(J_j^z+J_{j+1}^z)^2:+{\displaystyle \frac{\lambda }{2}}[:(J_j^{}+J_{j+1}^{})(J_j^++J_{j+1}^+):+(J^+J^{})]`$ (66)
$`+({\displaystyle \frac{\pi v}{2}}\pi g):\left[J_j^z(x)\right]^2:\lambda [:J_j^{}(x)J_j^+(x):+:J_j^+(x)J_j^{}(x):]\}.`$
The two terms in the first line are clearly non-negatively defined, while the terms in the second line can be shown to be non-negatively defined for $`v>2g+8\lambda /\pi `$ using Eq. (58) and the basic identity
$$\left(𝐉_j\right)^2=\left(J_j^z\right)^2+\frac{1}{2}(J_j^+J_j^{}+J_j^{}J_j^+).$$
(67)
An analogous construction can be done for $`\lambda <0`$ or $`g<0`$, and we conclude that the Hamiltonian (57) is positive semi-definite provided
$$2|g|+8|\lambda |/\pi <v.$$
(68)
This is also the sufficient condition for the stability of the original ground state, as follows from our previous result that the Hamiltonian annihilates the original ground state $`|0`$.
Note that the normal ordering does not present a complication because it amounts to subtracting off the unit operator times the expectation value of the Hamiltonian in the state $`|0`$, which is zero. Therefore, since the Hamiltonian is positive semi-definite and we know the state $`|0`$ is a zero-energy eigenstate, it follows that $`|0`$ is the ground state of the Hamiltonian provided Eq. (68) is satisfied.
We can now combine the sum rule with our knowledge of the exact ground state to say something about the conductivity. At zero temperature the expectation value on the r.h.s. of the sum rule (53) is a ground state expectation value. Since we know the ground state $`|0`$ is unchanged by tunneling and interactions and $`0|_1|0=0`$ we can conclude
$$𝑑\omega \mathrm{}e\sigma ^{zz}(\omega )=0.$$
(69)
The real part of $`\sigma ^{zz}(\omega )`$ is dissipative and hence it cannot be negative. Therefore, from Eq. (69) we have
$$\sigma ^{zz}(\omega )=0.$$
(70)
We reiterate that this is an exact statement for the clean 020 multilayer at $`T=0`$ provided $`g`$ and $`\lambda `$ satisfy Eq. (68). As we mentioned above, a similar result can be shown to hold for the 131 state, the other maximally-chiral state with marginal tunneling. The condition for the stability of the ground state of the 131 edge theory is also given by Eq. (68). In the next section we will explore what happens when we add disorder to these systems.
### B Disordered Case
In the previous section we found that the clean multilayer with marginal tunneling exhibits insulating behavior. In this section we study the marginal tunneling case in the presence of disorder. We perform a perturbative calculation of the vertical conductivity via the Kubo formula, finding that disorder increases the conductivity.
In Section III we discussed the general edge theory of clean multilayer systems. Our first task is to add disorder to this model. We will then use the Kubo formula to find the $`z`$-axis conductivity.
The disorder we consider is a random scalar potential $`V_j(x)`$ in each layer $`j`$, which couples to the Bose fields via the charge density
$$_V=_{L/2}^{L/2}𝑑xV_j(x)\rho _j(x)=_{L/2}^{L/2}𝑑x\frac{1}{2\pi }V_j(x)K_{jk}_xu_k(x).$$
(71)
We assume $`V_j(x)`$ is a Gaussian random variable, uncorrelated between different layers, but for now we will not specify how it is correlated within the same layer.
Using the Fourier transformation (26) and rescaling (30) of the Bose fields we can write Eq. (71) as
$$_V=_{L/2}^{L/2}𝑑x\frac{1}{2\pi }\frac{1}{\sqrt{K_q}}V_q_x\varphi _q(x),$$
(72)
where
$$V_q(x)=\frac{1}{\sqrt{N}}\underset{j}{}e^{iqj}V_j(x).$$
(73)
Recalling the form of the clean multilayer edge theory from Section III, our full Hamiltonian is
$``$ $`=`$ $`{\displaystyle _{L/2}^{L/2}}dx[{\displaystyle \frac{1}{4\pi }}v_q:_x\varphi _q_x\varphi _q:+{\displaystyle \underset{j}{}}{\displaystyle \frac{\lambda }{(2\pi a)^2}}(e^{i\alpha _{jq}\varphi _q(x)}+\mathrm{h}.\mathrm{c}.)`$ (75)
$`+{\displaystyle \frac{1}{2\pi }}{\displaystyle \frac{1}{\sqrt{K_q}}}V_q_x\varphi _q(x)].`$
We move the disorder term into the phase of the tunneling amplitude by shifting the bosons:
$$\varphi _q\varphi _q+\frac{1}{v_q\sqrt{K_q}}^x𝑑yV_q(y).$$
(76)
If we define the phases
$$\gamma _j(x)\frac{1}{\sqrt{N}}\underset{q}{}e^{iqj}(1e^{iq})\frac{1}{v_q}^x𝑑yV_q(y),$$
(77)
then with Eq. (76) the Hamiltonian (75) is
$$=_{L/2}^{L/2}dx[\frac{1}{4\pi }v_q:_x\varphi _q_x\varphi _q:+\frac{\lambda }{(2\pi a)^2}\underset{j}{}(e^{i\gamma _j(x)}e^{i\alpha _{jq}\varphi _q(x)}+\mathrm{h}.\mathrm{c}.)],$$
(78)
and the current operator (45) becomes
$$I^z=\frac{i\lambda d}{(2\pi a)^2}\underset{j}{}_{L/2}^{L/2}dx[e^{i\gamma _j(x)}e^{i\alpha _{jq}\varphi _q(x)}\mathrm{h}.\mathrm{c}.].$$
(79)
We are now in a position to calculate the $`z`$-axis conductivity. We will first evaluate the Matsubara two-point function of the current operator (79). Since Eq. (78) is a theory of $`N`$ interacting chiral bosons, the best we can do is a perturbative calculation in the tunneling amplitude $`\lambda `$. Although we will work perturbatively in $`\lambda `$, we will be treating the disorder exactly. We thus expect the $`𝒪(\lambda ^2)`$ result to be reliable, since the disorder essentially randomizes the phase of the electrons between tunneling events. More formally, if the tunneling amplitude is a coordinate-dependent, delta-correlated random variable, then the RG flow (37) is modified to
$$\frac{d\mathrm{\Delta }}{d\mathrm{}}=(32\delta _\lambda )\mathrm{\Delta },$$
(80)
where $`\mathrm{\Delta }`$ is the variance of the tunneling. Thus for $`\delta _\lambda =2`$ we would conclude that the tunneling is irrelevant. The combination $`\lambda e^{i\gamma _j(x)}`$ that appears in the Hamiltonian (78) is in general not delta-correlated, as we will discuss below, but since the case of $`x`$-independent tunneling is marginal we expect any departure from uniformity makes the tunneling irrelevant and therefore justifies a perturbative expansion in $`\lambda `$.
Using the expression for the current operator (79) we evaluate the Matsubara two-point function
$`𝖢(\tau )`$ $`=`$ $`T_\tau I^z(\tau )I^z(0)`$ (81)
$`=`$ $`{\displaystyle \frac{d^2\lambda ^2}{L^4}}{\displaystyle }dxdx^{}{\displaystyle \underset{j,k}{}}T_\tau (e^{i\gamma _j(x)}:e^{i\alpha _{jq}\varphi _q(\tau ,x)}:e^{i\gamma _j(x)}:e^{i\alpha _{jq}\varphi _q(\tau ,x)}:)`$ (83)
$`\times (e^{i\gamma _k(x^{})}:e^{i\alpha _{kq}\varphi _q(0,x^{})}:e^{i\gamma _k(x^{})}:e^{i\alpha _{kq}\varphi _q(0,x^{})}:),`$
where we have normal-ordered the vertex operators. Note that because the current operator (79) has a factor of $`\lambda `$ out front, there is an explicit factor of $`\lambda ^2`$ in Eq. (83), and therefore to find this function to order $`\lambda ^2`$ we can evaluate the expectation value using the $`\lambda =0`$ Hamiltonian. One finds that only the cross terms in the product of the current operators give non-vanishing contributions, and only when $`j=k`$. The terms have the form
$`T_\tau :e^{i\alpha _{jq}\varphi _q(\tau ,x)}::e^{i\alpha _{jq}\varphi _q(0,x^{})}:`$ (85)
$`={\displaystyle \underset{q}{}}{\displaystyle \frac{L^{\alpha _q^2}}{\left[(2\beta iv_q)\mathrm{sinh}\left(\frac{\pi }{\beta v_q}(xx^{}+iv_q\tau +ia\mathrm{sgn}\tau )\right)\right]^{\alpha _q^2}}},`$
where $`\alpha _q^2\alpha _{jq}\alpha _{jq}`$, with no sum on $`j`$. This is a very complicated function; a great simplification occurs if the eigenmode velocities, $`v_q`$, are identical for all $`q`$. From the expression for $`v_q`$ (34) we see that this occurs when $`g=nv/m`$. For the two multilayer states with marginal tunneling, 020 and 131, this condition gives $`g=0`$ and $`g=v/3`$, respectively. We will refer to the case where $`v_q`$ is independent of $`q`$ as the “solvable point.” For 020 it corresponds to no density-density couplings between neighboring layers while for 131 it occurs at a non-zero value of the density-density interaction strength. Note that for more general forms of the matrices $`K`$ and $`U`$, i.e., not tridiagonal, the solvable point corresponds to $`UK`$.
Working at the solvable point, writing $`v_qv_0`$, and using our knowledge of the scaling dimension of the tunneling operator (38),
$$\underset{q}{}\alpha _q^2=\underset{q}{}\alpha _{jq}\alpha _{jq}=2\delta _\lambda =4,$$
(86)
we find upon disorder averaging Eq. (83)
$$\overline{𝖢}(\tau )=\frac{2d^2\lambda ^2}{(2\beta v_0)^4}\underset{j}{}𝑑x𝑑x^{}\frac{W_j(x,x^{})}{\left[\mathrm{sinh}\left(\frac{\pi }{\beta v_0}(xx^{}+iv_0\tau +ia\mathrm{sgn}\tau )\right)\right]^4},$$
(87)
where we have defined the function
$$W_j(x,x^{})\overline{e^{i[\gamma _j(x)\gamma _j(x^{})]}}.$$
(88)
If we take the correlation function of the scalar potential to be
$$\overline{V_j(x)V_k(x^{})}=\delta _{jk}Z(xx^{}),$$
(89)
then using the definition of the phase $`\gamma _j(x)`$ (77) we find that the function $`W_j(x,x^{})`$ can be written
$$W_j(x,x^{})=\mathrm{exp}\left[\frac{1}{v_0^2}_x^x^{}𝑑y_1_x^x^{}𝑑y_2Z(y_1y_2)\right],$$
(90)
Note that $`W_j(x,x^{})W(xx^{})`$ is independent of $`j`$ and depends only on the magnitude of the relative coordinate $`|xx^{}|`$. Therefore the sum over $`j`$ in Eq. (87) just gives a factor of the number of layers, $`N`$, and we can replace the integration over $`x`$ and $`x^{}`$ by an integration over the average coordinate, which gives a factor of $`L`$, and an integration over the relative coordinate $`y=xx^{}`$. If we define the Fourier transform of the function $`W(x)`$:
$$\stackrel{~}{W}(k)𝑑xe^{ikx}W(x),$$
(91)
then we can write Eq. (87) as
$$\overline{𝖢}(\tau )=\frac{2d^2NL\lambda ^2}{(2\beta v_0)^4}\frac{dk}{2\pi }\stackrel{~}{W}(k)𝑑y\frac{e^{iky}}{\left[\mathrm{sinh}\left(\frac{\pi }{\beta v_0}(y+iv_0\tau +ia\mathrm{sgn}\tau )\right)\right]^4}$$
(92)
The factor of $`L`$ in this expression will cancel against the factor of $`1/L`$ in the Kubo formula (52), and in anticipation of this we have extended the range of the $`y`$ integral to include the entire real axis since we are interested in the result in the limit where the system size $`L\mathrm{}`$. The integration over $`y`$ can now be performed by the method of residues. The integrand has fourth-order poles on the imaginary $`y`$-axis at the points $`y_n=i[v_0(\beta n\tau )a\mathrm{sgn}\tau ]`$ for integer $`n`$, and is exponentially small in the upper (lower) half plane for $`k>0`$ ($`k<0`$). One finds
$$𝑑y\frac{e^{iky}}{\left[\mathrm{sinh}\left(\frac{\pi }{\beta v_0}(y+iv_0\tau +ia\mathrm{sgn}\tau )\right)\right]^4}=\frac{\pi }{3}\left[4\left(\frac{\beta v_0k}{\pi }\right)+\left(\frac{\beta v_0k}{\pi }\right)^3\right]\frac{e^{v_0k\tau }}{e^{\beta v_0k}1}.$$
(93)
We now Fourier transform Eq. (92) with respect to the imaginary time $`\tau `$,
$`\overline{𝖢}(\omega _n)`$ $`=`$ $`{\displaystyle _0^\beta }𝑑\tau e^{i\omega _n\tau }\overline{𝖢}(\tau )`$ (94)
$`=`$ $`{\displaystyle \frac{d^2\lambda ^2NL}{48\pi ^4v_0^4}}{\displaystyle 𝑑p\stackrel{~}{W}(p/v_0)\frac{p^3+4\pi ^2p/\beta ^2}{i\omega _n+p}},`$ (95)
where we have changed the integration variable to $`pv_0k`$. Analytically continuing ($`i\omega _n\omega +i0^+`$) we find
$`\overline{𝒞}(\omega )`$ $`=`$ $`i{\displaystyle _0^{\mathrm{}}}𝑑te^{i\omega t}\overline{[I^z(t),I^z(0)]}`$ (96)
$`=`$ $`{\displaystyle \frac{d^2\lambda ^2NL}{48\pi ^4v_0^4}}{\displaystyle 𝑑p\stackrel{~}{W}(p/v_0)\frac{p^3+4\pi ^2p/\beta ^2}{\omega +p+i0^+}}.`$ (97)
Using this result in the Kubo formula (52) with the help of the identity
$$\frac{1}{x\pm i0^+}=𝒫\frac{1}{x}i\pi \delta (x),$$
(98)
and the observation that $`\stackrel{~}{W}(k)`$ is real since $`W(x)`$ is an even function of $`x`$, we arrive finally at
$$\mathrm{}e\overline{\sigma ^{zz}}(\omega )=\frac{\lambda ^2d}{48\pi ^3v_0^4}(\omega ^2+4\pi ^2T^2)\stackrel{~}{W}(\omega /v_0).$$
(99)
This is our central result. It is the $`𝒪(\lambda ^2)`$ term in the vertical conductivity at the solvable point for both the 020 and 131 states at a finite temperature $`T`$, including potential disorder which enters through the function $`\stackrel{~}{W}(k)`$. The imaginary part of the conductivity can of course be found from Eq. (99) by the usual dispersion relations.
Let us first consider the limit of a clean system, $`V_j(x)0`$. From Eqns. (89) and (90), we see that in this case $`W(x)=1`$ and therefore $`\stackrel{~}{W}(k)=2\pi \delta (k)`$. Hence from Eq. (99) we find
$$\mathrm{}e\sigma ^{zz}(\omega )=\frac{\lambda ^2d}{6v_0^3}T^2\delta (\omega ).$$
(100)
This is consistent with our result in Section IV A that the $`z`$-axis conductivity in the clean system vanishes at $`T=0`$. An interesting aspect of this result is that it is proportional to $`\delta (\omega )`$, i.e., the conductivity vanishes at all $`\omega 0`$, and the DC conductivity is infinite. One question to ask is whether these features are a consequence of the fact that Eq. (100) is only the first non-vanishing term in a perturbative expansion in $`\lambda `$.
Although we cannot definitively answer the question about the finiteness of the DC conductivity beyond $`𝒪(\lambda ^2)`$, we can prove that at higher orders in $`\lambda `$ the conductivity cannot vanish at all $`\omega 0`$. First we note that in the clean system the current operator commutes with the Hamiltonian in the absence of tunneling: $`[I^z,_0]=0`$. For example, for the 020 state this can be seen from Eqns. (62) and (63), using the commutation relations for the Kac-Moody currents given in Eq. (61). Since the current operator commutes with $`_0`$, the expectation value of the current-current commutator evaluated using the $`\lambda =0`$ Hamiltonian vanishes, and therefore by the Kubo formula (52) the vertical conductivity must vanish for all non-zero frequencies. This explains why the conductivity is zero except at $`\omega =0`$ in Eq. (100). By the same reasoning we see that if the current operator commutes with the full Hamiltonian, $`[I^z,]=0`$, then $`\sigma ^{zz}(\omega )`$ would be zero for all $`\omega 0`$ to all orders in $`\lambda `$. The converse of this is also true, i.e., if $`\sigma ^{zz}(\omega )`$ is zero for all $`\omega 0`$ then $`[I^z,]=0`$. This is proven in Appendix D. For the multilayer edge theory with marginal tunneling: $`[I^z,]0`$. For the 020 case this can be established using Eqns. (62) and (63). Therefore, since the current operator does not commute with the full Hamiltonian, we can conclude that $`\sigma ^{zz}(\omega )`$ cannot be zero for all $`\omega 0`$ for the clean 020 or 131 multilayers. Hence, we would expect that the $`\delta (\omega )`$ factor in Eq. (100) would be broadened by the higher-order corrections in $`\lambda `$. Recalling that $`\lambda `$ is dimensionless and the conductivity should not depend on the sign of $`\lambda `$, one plausible scenario is:
$$\mathrm{}e\sigma ^{zz}(\omega )=\frac{\lambda ^2d}{6v_0^3}T\delta (\omega /T)\frac{\lambda ^2d}{6v_0^3}T\frac{\lambda ^2}{(\omega /T)^2+𝒪(\lambda ^4)}.$$
(101)
This conjectured form has several interesting features. First, we find that the DC conductivity goes to zero with temperature as $`T`$. We expect this to be a generic feature of the exact result in $`\lambda `$ for the clean system, provided the DC conductivity is finite. This differs from the disordered system (99) where we find the DC conductivity vanishes as $`T^2`$. In addition, note that the $`\omega 0`$ limit of Eq. (101) is independent of the tunneling amplitude $`\lambda `$.
We now return to the result for $`\overline{\sigma ^{zz}}(\omega )`$ in the presence of disorder (99). If the disorder is delta-correlated, the function $`Z(x)`$ defined in Eq. (89) is
$$Z(x)=\mathrm{\Delta }\delta (x),$$
(102)
where $`\mathrm{\Delta }`$ is the variance of the disorder. For this case $`W(x)=\mathrm{exp}(\mathrm{\Delta }|x|/v_0^2)`$, which has a Fourier transform $`\stackrel{~}{W}(k)=(2\mathrm{\Delta }/v_0^2)/(k^2+(\mathrm{\Delta }/v_0^2)^2)`$, and hence (99) yields
$$\mathrm{}e\overline{\sigma ^{zz}}(\omega )=\frac{\lambda ^2d\mathrm{\Delta }}{24\pi ^3v_0^4}\frac{\omega ^2+4\pi ^2T^2}{\omega ^2+(\mathrm{\Delta }/v_0)^2}.$$
(103)
Note that at zero temperature the conductivity vanishes in the DC limit. This conclusion is independent of the details of the disorder and holds as long as $`\stackrel{~}{W}(0)`$ is finite. We also see that the conductivity approaches a constant as $`\omega \mathrm{}`$. This second feature is a consequence of taking the disorder to be delta-correlated. We see that $`\overline{\sigma ^{zz}}(\omega )`$ goes to a constant at large $`\omega `$ because of the slow fall off of $`\stackrel{~}{W}(k)`$, which is in turn caused by the non-analyticity of $`W(x)`$ at $`x=0`$. We shall now investigate how Eq. (103) is modified if the disorder has a finite correlation length.
As a first step, we determine the asymptotic form of the disorder-averaged phase, $`W(x)`$, at small and large $`x`$ when the scalar potential disorder has a finite correlation length. We return to the expression for the correlation function of the random potential (89) and write
$$Z(x)\mathrm{\Delta }^2A(x/\mathrm{}_0),$$
(104)
where $`\mathrm{\Delta }`$ is a parameter with the dimensions of energy that sets the strength of the disorder, $`\mathrm{}_0`$ is the correlation length of the disorder, and $`A(z)`$ is a dimensionless function which we take to have the properties: $`A(0)=1`$, $`A(z)=A(z)`$, and
$$_0^{\mathrm{}}𝑑zA(z)=1.$$
(105)
Now consider the disorder-averaged phase, $`W(x)`$, given in Eq. (90). For $`|x|\mathrm{}_0`$, i.e., for distances small compared with the correlation length of the random potential, we can replace $`Z(y_1y_2)`$ by its value at $`y_1y_2=0`$ and obtain:
$$W(\mathrm{}_0z)e^{(\mathrm{\Delta }\mathrm{}_0z/v_0)^2}1\left(\frac{\mathrm{\Delta }\mathrm{}_0z}{v_0}\right)^2,\mathrm{for}|z|1.$$
(106)
In the opposite limit $`|x|\mathrm{}_0`$ we change integration variables in Eq. (90) to $`y_c=(y_1+y_2)/2`$ and $`y_r=y_1y_2`$:
$$W(x)=\mathrm{exp}\left[\frac{1}{v_0^2}\left(_{|x|}^0𝑑y_r(|x|+y_r)Z(y_r)+_0^{|x|}𝑑y_r(|x|y_r)Z(y_r)\right)\right].$$
(107)
Since $`|x|/\mathrm{}_01`$ we can extend the integration over $`y_r`$ to infinity with small error. Using the symmetry of $`A`$ and defining
$$\zeta =_0^{\mathrm{}}𝑑zzA(z),$$
(108)
we then find
$$W(\mathrm{}_0z)\mathrm{exp}\left[2\left(\frac{\mathrm{\Delta }\mathrm{}_0}{v_0}\right)^2(|z|\zeta )\right],\mathrm{for}|z|1.$$
(109)
We see that at short distances $`W(x)`$ is parabolic (106) while at long distances it decays exponentially (109). The presence of a finite correlation length for the random potential does not change the long distance behavior of $`W(x)`$, but it does “round out” the non-analyticity at $`x=0`$ present in the case of delta-correlated disorder. Generally, this is sufficient to make the function rapidly decaying as $`\omega \mathrm{}`$.
Having found the generic behavior of the function $`W(x)`$, we now choose a specific form:
$$W(x)=\frac{1}{\mathrm{cosh}(\xi x)},$$
(110)
characterized by the parameter $`\xi `$ which has the dimensions of an energy. This choice of $`W(x)`$ has several appealing features. First, it has the correct asymptotic forms at large and small values of $`x`$. Indeed, if we match the asymptotics of Eq. (110) to Eqns. (106) and (109) we find
$$\frac{\mathrm{\Delta }}{v_0}=\xi ,\mathrm{}_0=\frac{1}{2\xi }.$$
(111)
We see that for the choice of $`W(x)`$ in Eq. (110), the parameter $`\xi `$ is proportional to the strength of the disorder and inversely proportional to the correlation length. Hence this form allows us to explore the effect of a finite correlation length in the region $`\mathrm{\Delta }\mathrm{}_01`$.
A second advantage of Eq. (110) is that its Fourier transform can be evaluated in closed form. A straightforward computation gives
$$\stackrel{~}{W}(k)=𝑑x\frac{e^{ikx}}{\mathrm{cosh}(\xi x)}=\frac{\pi /\xi }{\mathrm{cosh}\left(\pi k/2\xi \right)}.$$
(112)
Using this result in the general form above for the conductivity (99) we find
$$\mathrm{}e\overline{\sigma ^{zz}}(\omega )=\frac{\lambda ^2d}{48\pi ^2v_0^4}\frac{\omega ^2+4\pi ^2T^2}{\xi \mathrm{cosh}\left(\pi \omega /2\xi v_0\right)}.$$
(113)
Comparing Eq. (113) with the expression for the conductivity in the case of delta-correlated disorder (103) we see that the finite correlation length gives a conductivity which exponentially decays to zero as $`\omega \mathrm{}`$ rather than approaching a nonzero value. This result for $`\overline{\sigma ^{zz}}`$ (113) is shown in Fig. 2 at various temperatures. For large temperatures the conductivity is peaked around $`\omega =0`$, while at smaller temperatures the maximum occurs at a finite frequency.
### C Effect of Interactions
In the preceding section we calculated the order $`\lambda ^2`$ term in the $`z`$-axis conductivity of the multilayer states with marginal tunneling, 020 and 131, at the solvable point where the eigenmode velocities $`v_q`$ are independent of $`q`$. In this section we discuss how the zero-temperature result is modified away from the solvable point.
Since we will only consider $`T=0`$ in this section we shall work directly in real time $`t`$. The real-time zero-temperature analog of the expression for the Matsubara two-point function of the current operator (87) is
$$\overline{𝒞}(t)=\frac{2id^2\lambda ^2NL}{(2\pi )^4}𝑑yW(y)\underset{q}{}\frac{1}{(yv_qt+ia\mathrm{sgn}t)^{\alpha _q^2}},$$
(114)
where from the definitions (77) and (88) we see that the expression for $`W(x)`$ in Eq. (90) is changed to
$$W(x)=\mathrm{exp}\left[\left(\frac{2}{N}\underset{q}{}\frac{\mathrm{sin}^2(q/2)}{v_q^2}\right)_0^x𝑑y_1_0^x𝑑y_2Z(y_1y_2)\right].$$
(115)
The corresponding retarded correlation function is
$$\overline{𝒞^R}(t)=\frac{2id^2\lambda ^2NL}{(2\pi )^4}\theta (t)𝑑yW(y)\left[\underset{q}{}\frac{1}{(yv_qt+ia)^{\alpha _q^2}}\underset{q}{}\frac{1}{(yv_qtia)^{\alpha _q^2}}\right].$$
(116)
We would now like to perform the $`y`$-integration. From Eq. (36) we see that the exponents appearing in Eq. (116) are
$$\alpha _q^2=\frac{2}{N}(3\mathrm{cos}q2\mathrm{cos}^2q),$$
(117)
and therefore, as a function of complex $`y`$, the first integrand in Eq. (116) has $`N`$ branch point singularities at the points $`v_qtia`$ and similarly for the second integrand at $`v_qt+ia`$. This is a complicated singularity structure, but observe that in the first term these branch points are all below the real axis, while in the second term they are all above the real axis. Hence, since $`_q\alpha _q^2`$ is an integer, we can certainly choose branch cuts that lie entirely on one side of the real axis for each term. When we considered the solvable point in the previous section we were able to evaluate the conductivity for any disorder $`W(y)`$. We shall not attempt the same feat away from the solvable point. The evaluation of the $`y`$-integral in Eq. (116) is greatly simplified if $`W(y)`$ is a meromorphic function of $`y`$. One such function was discussed in the previous section (110), and for the remainder of this section we specialize to this form.
Using $`W(y)=1/\mathrm{cosh}(\xi y)`$ in Eq. (116), we can close the first term in the upper-half plane and the second term in the lower-half plane, avoiding all the complicated branch point singularities, and picking up the residues of the poles of $`W(y)`$. We find
$`\overline{𝒞^R}(t)`$ $`=`$ $`i\theta (t)\overline{[I^z(t),I^z(0)]}`$ (118)
$`=`$ $`{\displaystyle \frac{2id^2\lambda ^2NL}{(2\pi )^3}}\theta (t)\xi ^3{\displaystyle \underset{r=\mathrm{}}{\overset{\mathrm{}}{}}}(1)^r{\displaystyle \underset{q}{}}{\displaystyle \frac{1}{[v_q\xi ti\pi (r+1/2)]^{\alpha _q^2}}}.`$ (119)
Therefore, from the Kubo formula (52) we arrive at
$$\mathrm{}e\overline{\sigma ^{zz}}(\omega )=\frac{\lambda ^2d}{(2\pi )^3}\frac{\xi ^3}{\omega }\underset{r=\mathrm{}}{\overset{\mathrm{}}{}}(1)^r𝑑te^{i\omega t}\underset{q}{}\frac{1}{[v_q\xi ti\pi (r+1/2)]^{\alpha _q^2}}.$$
(120)
This is certainly far from a closed form result, but it is a convenient form for extracting the high and low frequency asymptotics of the conductivity. At large (positive) $`\omega `$ the asymptotic form is determined by the closest singularity to the real axis in the upper half of the complex $`t`$ plane. From Eq. (120) this clearly occurs a finite distance away from the real-$`t`$ axis and we find
$$\mathrm{}e\overline{\sigma ^{zz}}(\omega )\stackrel{\omega \mathrm{}}{}\mathrm{exp}\left[\frac{\pi \omega }{2\xi v_{\mathrm{max}}}\right],$$
(121)
where $`v_{\mathrm{max}}`$ is the maximum of $`v_q`$ over $`q`$. For the 020 state we see from the expression for $`v_q`$ (34) that at the solvable point ($`g=0`$) $`v_0=v/2`$ while away from the solvable point ($`g>0`$) $`v_{\mathrm{max}}=v/2+g`$. Comparing Eq. (121) with the asymptotic form of Eq. (113):
$$\mathrm{}e\overline{\sigma ^{zz}}(\omega )\stackrel{\omega \mathrm{}}{}\mathrm{exp}\left[\frac{\pi \omega }{2\xi v_0}\right],$$
(122)
we find that one effect of interactions at $`T=0`$ is to soften the exponential decay of the conductivity at large frequencies.
Returning to Eq. (120), we now consider the behavior at small frequencies. If we assume $`\omega >0`$ we can close the $`t`$ integral in the upper-half plane because the integrand is exponentially small there. Hence all the terms with $`r<0`$ give no contribution. For each $`r0`$ there remains $`N`$ branch point singularities in the upper-half plane enclosed by the contour. Since we are interested in small $`\omega `$ we stretch the integration contour into a large circle so that everywhere along the contour $`|t|`$ is very large. We can then expand the integrand in powers of $`1/t`$ and integrate term-by-term:
$`{\displaystyle 𝑑te^{i\omega t}\underset{q}{}\frac{1}{[v_q\xi ti\pi (r+1/2)]^{\alpha _q^2}}}`$ (125)
$`={\displaystyle 𝑑te^{i\omega t}\left(\underset{q}{}\frac{1}{v_q^{\alpha _q^2}}\right)\frac{1}{(\xi t)^4}\left[1+\frac{i\pi (r+1/2)}{\xi t}\underset{q}{}\frac{\alpha _q^2}{v_q}+𝒪\left(\frac{1}{t^2}\right)\right]}`$
$`={\displaystyle \frac{\pi }{3}}\left({\displaystyle \underset{q}{}}{\displaystyle \frac{1}{v_q^{\alpha _q^2}}}\right){\displaystyle \frac{\omega ^3}{\xi ^4}}\left[1{\displaystyle \frac{\pi }{4}}{\displaystyle \frac{\omega }{\xi }}(r+1/2){\displaystyle \underset{q}{}}{\displaystyle \frac{\alpha _q^2}{v_q}}+𝒪(\omega ^2)\right].`$
If we retain only the lowest term in $`\omega `$ in this expansion, we note that upon substituting this into Eq. (120) we would encounter the divergent sum
$$\underset{r=0}{\overset{\mathrm{}}{}}(1)^r.$$
(126)
By considering the solvable point, where $`v_q=v_0`$ independent of $`q`$, it is clear that the proper regularization of this sum is
$$\underset{r=0}{\overset{\mathrm{}}{}}(1)^r=\underset{\delta 0}{lim}\underset{r=0}{\overset{\mathrm{}}{}}(e^\delta )^r=\underset{\delta 0}{lim}\frac{1}{1+e^\delta }=\frac{1}{2},$$
(127)
and thus from Eqns. (120), (125), and (127) we have
$$\mathrm{}e\overline{\sigma ^{zz}}(\omega )\stackrel{\omega 0}{}\frac{\lambda ^2d}{48\pi ^2}\left(\underset{q}{}\frac{1}{v_q^{\alpha _q^2}}\right)\frac{\omega ^2}{\xi }.$$
(128)
Comparing this with the small frequency asymptotics of Eq. (113) at $`T=0`$:
$$\mathrm{}e\overline{\sigma ^{zz}}(\omega )\stackrel{\omega 0}{}\frac{\lambda ^2d}{48\pi ^2}\frac{1}{v_0^4}\frac{\omega ^2}{\xi },$$
(129)
we see that the interactions do not modify the result that the conductivity goes to zero in the DC limit at zero temperature. While the power of $`\omega `$ is the same in Eqns. (128) and (129), the interactions do modify the prefactor.
To summarize, we have considered how interactions away from the solvable point modify the vertical conductivity at $`T=0`$. It should be noted that in arriving at Eq. (120) the interactions were treated exactly. We find that at large frequencies the interactions soften the decay of the conductivity, but it remains exponential, while at small frequencies the interactions do not change the frequency exponent of the conductivity. We remark that the difficulty encountered when attempting to apply the method presented in this section to investigate the effect of interactions at a finite temperature is that in the finite-$`T`$ version of Eq. (116) there are complicated branch point singularities on both sides of the real-$`y`$ axis, see Eq. (85). Nevertheless, we believe that even at a finite temperature the asymptotics of $`\sigma ^{zz}(\omega )`$ are not modified by interactions.
## V Irrelevant Tunneling
In the previous section we have been concerned exclusively with the case of marginal tunneling, namely the 020 and 131 multilayer states. The calculation presented in Section IV B can readily be extended to the more general case of the $`nmn`$ state. Essentially the only change in the computation is that the scaling dimension of the tunneling operator, $`\delta _\lambda =mn`$, differs from its marginal value of 2. Therefore, the poles in the complex $`y`$-plane for the integrand in Eq. (93) are now of order $`2(mn)`$. One finds that for the solvable point $`g=nv/m`$:
$$\mathrm{}e\overline{\sigma _{nmn}^{zz}}(\omega )=\frac{\lambda ^2d}{(2\pi )^{2\delta _\lambda 1}v_0^{2\delta _\lambda }}\frac{\stackrel{~}{W}(\omega /v_0)}{(2\delta _\lambda 1)!}\underset{k=1}{\overset{\delta _\lambda 1}{}}[\omega ^2+(2\pi kT)^2].$$
(130)
For the case of relevant tunneling, $`\delta _\lambda =mn=1`$, the product over $`k`$ is absent in Eq. (130), and if we use $`\stackrel{~}{W}(\omega /v_0)=(2\mathrm{\Delta })/(\omega ^2+(\mathrm{\Delta }/v_0)^2)`$ appropriate for delta-correlated disorder we reproduce the behavior of the chiral metal.
In the remaining cases, $`\delta _\lambda 2`$, we see at $`T=0`$ the conductivity vanishes as $`\omega 0`$ as $`\omega ^{2(\delta _\lambda 1)}`$, and at $`\omega =0`$ it vanishes with temperature as $`T^{2(\delta _\lambda 1)}`$. In the absence of interlayer correlations, i.e., for the $`0m0`$ state, this means the DC conductivity obeys $`\sigma ^{zz}T^{2m2}`$. This disagrees with the result of Balents and Fisher, who find by a scaling argument that $`\sigma ^{zz}T^{2m3}`$. While it is unclear whether Balents and Fisher have in mind the clean system or the disordered system, we believe the reconciliation of this discrepancy is that, as we found for the marginal case, the temperature scaling of the DC conductivity is different in the clean and disordered systems. If we consider the clean limit of Eq. (130) by writing $`\stackrel{~}{W}(\omega /v_0)=(v_0/T)\delta (\omega /T)`$, and then conjecture that the interactions broaden the delta function, we find:
$$\mathrm{}e\overline{\sigma _{nmn}^{zz}}(\omega )\frac{\lambda ^2d}{T}\frac{g}{(\omega /T)^2+𝒪(g^2)}\underset{k=1}{\overset{\delta _\lambda 1}{}}[\omega ^2+(2\pi kT)^2],$$
(131)
where $`g`$ is the dimensionless interaction strength. The broadening of the delta function by interactions is suggested by the fact that the nearest-neighbor density-density interaction Hamiltonian does not commute with the current operator $`I^z`$. For the case of the $`0m0`$ state the DC conductivity of Eq. (131) goes to zero as $`T^{2m3}`$, in agreement with Balents and Fisher.
Finally, returning to the disordered case (130), we see that for the 050 state the DC conductivity scales as $`T^8`$, while for 131 it scales as $`T^2`$ (99). Recall from Eq. (5) that both of these states occur at the same electron density, $`\nu =1/5`$ per layer. In their numerical work, QJM found a phase transition between these states as the layer separation $`d`$ was varied. The significant difference in the temperature scaling of the DC vertical conductivity in these states suggests a way to experimentally detect this phase transition.
## VI Fermionization
In certain special cases the multilayer edge theory can be solved exactly. The case of $`N=2`$ layers is discussed extensively elsewhere, and examples of exact solutions for $`N=4`$ layers are given in Appendix B. A potentially useful first step to solving the edge theory for arbitrary $`N`$ is to find a fermionic representation. In this section we discuss the fermionization of the edge theory for an arbitrary number of layers $`N`$. In some cases the procedure involves the introduction of auxiliary degrees of freedom as discussed in Ref. .
We begin with the observation that the scaling dimension (38) of the tunneling operator is always an integer. This underlies the fact that the chiral Hamiltonian (39) can be fermionized exactly, with the tunneling term expressed as a product of Fermi operators whose number equals twice the scaling dimension.
The usual mapping between chiral radius-one bosons and fermions requires independent bosonic fields. Even though we are interested in correlated states with non-trivial commutators (22), these can be readily diagonalized by a linear transformation. For example, for the 131 state we can write
$$u_j=\phi _{j1/2}+\phi _j+\phi _{j+1/2},j=1,2,\mathrm{},N$$
(132)
where the $`2N`$ bosons $`\phi _j`$ with integer and half-integer indices have the usual commutation relations $`[\phi _i(x),\phi _j(x^{})]=i\pi \delta _{ij}\mathrm{sgn}(xx^{})`$. Clearly, Eq. (132) is not a canonical transformation because the number of degrees of freedom has been doubled. This formal difficulty can be circumvented if we introduce an independent set of auxiliary bosons with half-integer indices $`\stackrel{~}{u}_{i+1/2}`$, $`i=1,2,\mathrm{},N`$ with identical commutation relations,
$$[\stackrel{~}{u}_{i+1/2}(x),\stackrel{~}{u}_{j+1/2}(x^{})]=i\pi K_{ij}\mathrm{sgn}(xx^{}),[u_i(x),\stackrel{~}{u}_{j+1/2}(x^{})]=0.$$
(133)
Then, we can further write for the 131 state
$$\stackrel{~}{u}_{j+1/2}=\phi _j+\phi _{j+1/2}\phi _{j+1},$$
(134)
and the transformation $`u_i,\stackrel{~}{u}_{i+1/2}\phi _j`$ becomes canonical.
At the end of a calculation the auxiliary degrees of freedom need to be projected out. Such a projection has been explicitly carried out by us in Ref. for the simpler case of correlated quantum Hall bilayers. However, since the auxiliary degrees of freedom are independent of the physical ones, all quantum-mechanical averages involving physical quantities will automatically be correct, independent of the Hamiltonian chosen for the additional bosons $`\stackrel{~}{u}_{i+1/2}`$. In the following we select this Hamiltonian in the form (21), i.e., identical to that for the physical bosonic fields sans the tunneling term. In this case, the transformations (132) and (134) give
$`_{131}`$ $`=`$ $`{\displaystyle _{L/2}^{L/2}}dx\{{\displaystyle \frac{1}{4\pi }}{\displaystyle \underset{i,j=1}{\overset{N}{}}}\stackrel{~}{V}_{ij}(:_x\phi _i_x\phi _j:+:_x\phi _{i+1/2}_x\phi _{j+1/2}:)`$
$`+{\displaystyle \underset{j}{}}[{\displaystyle \frac{\lambda }{(2\pi a)^2}}:e^{i\left(\phi _{j1/2}+\phi _j\phi _{j+1}\phi _{j+3/2}\right)}:+\mathrm{h}.\mathrm{c}.]\},`$
where the modified interaction matrix is given by the matrix product $`\stackrel{~}{V}_{ij}K_{il}V_{lj}`$. Now we can use the standard fermionization prescription:
$$\psi _{j+\alpha }=\frac{1}{\sqrt{2\pi a}}e^{i\phi _{j+\alpha }},$$
(135)
where $`\alpha =0,1/2`$. In particular, in the fermionic representation the original electron annihilation operator is written as a product of three fermion operators, $`\mathrm{\Psi }_j\psi _{j1/2}\psi _j\psi _{j+1/2}`$, while the complete fermionized Hamiltonian for the surface of the 131 system becomes
$`_{131}`$ $`=`$ $`{\displaystyle _{L/2}^{L/2}}dx\{{\displaystyle \underset{j,\alpha }{}}\stackrel{~}{V}_{jj}:\psi _{j+\alpha }^{}i_x\psi _{j+\alpha }:+{\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha ,ij}{}}\stackrel{~}{V}_{ij}:\psi _{i+\alpha }^{}\psi _{i+\alpha }::\psi _{j+\alpha }^{}\psi _{j+\alpha }:`$ (137)
$`+{\displaystyle \underset{j}{}}[\lambda \psi _{j1/2}^{}\psi _j^{}\psi _{j+1}\psi _{j+3/2}+\mathrm{h}.\mathrm{c}.]\}.`$
A curious feature of this transformation is that the introduced auxiliary degrees of freedom can be thought of as located at the layers sandwiched between the physical layers—this is the reason for assigning half-integer indices to them.
Similar fermion representations also exist for other correlated multilayer states. For the bosonic state 020, the interlayer correlations are absent, and the commutation relations can be diagonalized by the transformation
$$u_j=\phi _j+\phi _j,$$
(138)
where the auxiliary degrees of freedom labeled by the pseudospin index are located in the same layers. The fundamental boson operator can be written as $`\mathrm{\Psi }_j=\psi _j\psi _j`$, and the full Hamiltonian is
$`_{020}`$ $`=`$ $`{\displaystyle _{L/2}^{L/2}}dx\{{\displaystyle \underset{j,\sigma }{}}\stackrel{~}{V}_{jj}:\psi _{j\sigma }^{}i_x\psi _{j\sigma }:+{\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma ,ij}{}}\stackrel{~}{V}_{ij}:\psi _{i\sigma }^{}\psi _{i\sigma }::\psi _{j\sigma }^{}\psi _{j\sigma }:`$ (140)
$`+{\displaystyle \underset{j}{}}[\lambda \psi _j^{}\psi _j^{}\psi _{j+1}\psi _{j+1}+\mathrm{h}.\mathrm{c}.]\},\sigma =,.`$
Despite the similarity with its counterpart for the 131 state \[Eq. (137)\], this Hamiltonian involves subtly different correlations.
In the case of 020 state, one can also avoid doubling the number of degrees of freedom by using the alternative fermionization introduced in Ref. . Specifically, the neighboring layers are combined into pairs,
$$u_{2n+1}=\phi _{2n+1}+\phi _{2n+2},u_{2n+2}=\phi _{2n+1}\phi _{2n+2},$$
(141)
and the tunneling operators are fermionized alternatingly as anomalous two- or four-fermion combinations,
$$\frac{1}{(2\pi a)^2}e^{i(u_1u_2)}=\psi _2^{}i_x\psi _2^{},\frac{1}{(2\pi a)^2}e^{i(u_2u_3)}=\psi _1^{}\psi _2\psi _3\psi _4,\mathrm{},$$
(142)
where $`\psi _j`$ is given in Eq. (135). In cases where the number of layers $`N`$ is small, e.g., $`N=2,3,4`$, alternative fermionizations exists for the 020 and 131 states, which in some cases allow exact solutions, see Ref. and Appendix B.
The fermionization is also very simple for the pyramid states. As an example consider the 12321 state where the commutation relations are satisfied if we write
$$u_j=\phi _{j1}+\phi _j+\phi _{j+1},$$
(143)
and the corresponding tunneling operator is bilinear,
$$\frac{1}{2\pi a}e^{i(u_ju_{j+1})}=\psi _{j1}^{}\psi _{j+1}.$$
(144)
The edge theory with tunneling is quadratic and therefore exactly soluble, except for the complication of the soft modes mentioned in Section III.
## VII Summary
We have investigated the fractional quantum Hall effect in infinite layer systems. In the bulk we find quasiparticles with several unusual (irrational) features and at the edge we find surface phases which can be considered “chiral semi-metals” when the interlayer tunneling is not relevant. In addition we have presented fermionizations of the edge theory which may prove useful in finding exact solutions.
## VIII Acknowledgments
We would like to acknowledge support by an NSF Graduate Research Fellowship (JDN), DOE Grant DE-FG02-90ER40542 (LPP) as well as NSF grant No. DMR-99-78074, US-Israel BSF grant No. 9600294, and fellowships from the A. P. Sloan Foundation and the David and Lucille Packard Foundation (SLS). LPP and SLS thank the Aspen Center for Physics for its hospitality during the completion of part of this work.
## A Model problem with marginal tunneling
Here we consider a model problem with marginal interlayer tunneling. The Hamiltonian is chosen in analogy with the Hamiltonian (78),
$$=_{L/2}^{L/2}dx\{v\psi _j^{}i_x\psi _j+\frac{1}{2}[\lambda e^{i\gamma _j(x)}(\psi _j^{}i_x\psi _{j+1}i_x\psi _j^{}\psi _{j+1})+\mathrm{h}.\mathrm{c}.]\},$$
(A1)
with the marginal tunneling operator in parenthesis of a form reminiscent of the two-fermion-fermionized version of the marginal tunneling operator (142). The great advantage of the Hamiltonian (A1) is its linearity, which allows us to calculate the conductivity for this model exactly. In order to do this, we perform a gauge transformation $`\psi _je^{i\zeta _j(x)}\psi _j`$, with $`\zeta _{j+1}\zeta _j\gamma _j`$. This gives, in obvious matrix notation,
$$=_{L/2}^{L/2}𝑑x\left[\mathrm{\Psi }^{}\mathrm{\Lambda }i_x\mathrm{\Psi }+\frac{1}{2}\mathrm{\Psi }^{}\left(\mathrm{\Lambda }\widehat{V}+\widehat{V}\mathrm{\Lambda }\right)\mathrm{\Psi }\right],$$
(A2)
where $`\mathrm{\Psi }(\psi _1\psi _2\mathrm{})^T`$, the tri-diagonal matrix $`\mathrm{\Lambda }_{ll^{}}v\delta _{ll^{}}+\lambda (\delta _{l,l^{}+1}+\delta _{l+1,l^{}})`$ is coordinate-independent, while $`\widehat{V}\mathrm{diag}(V_1,V_2,\mathrm{})`$ with $`V_j_x\zeta _j`$. The corresponding transverse current operator has the form
$$I^z=_{L/2}^{L/2}𝑑x\mathrm{\Psi }^{}\left[\widehat{t}i_x+\frac{1}{2}\left(\widehat{t}\widehat{V}+\widehat{V}\widehat{t}\right)\right]\mathrm{\Psi },$$
(A3)
where $`\widehat{t}_{ij}i\lambda (\delta _{i,j+1}\delta _{i+1,j})`$.
The formal solution of this model (for a given realization of disorder) can be written in the limit of infinite system size ($`L\mathrm{}`$) using the transfer matrix approach described in Appendix A of Ref. . Using the Kubo formula for the vertical conductance, we obtain
$$G^{zz}(\omega )=\frac{L}{Nd}\sigma ^{zz}(\omega )=\frac{1}{2N\omega }\underset{k}{}_{L/2}^{L/2}𝑑x\left(n_kn_{k+\omega }\right)\mathrm{Tr}\left[\stackrel{~}{𝒱}(x)S_k(x,y)\stackrel{~}{𝒱}(y)S_{k+\omega }(y,x)\right],$$
(A4)
where $`n_k=(e^{\beta vk}+1)^1`$, the matrix in the vertex
$$\stackrel{~}{𝒱}(x)\mathrm{\Lambda }^{1/2}\left[\widehat{t}i_x+\frac{1}{2}\left(\widehat{t}\widehat{V}+\widehat{V}\widehat{t}\right)\right]\mathrm{\Lambda }^{1/2},$$
(A5)
and the transfer matrix
$$S_k(a,b)=T_x\mathrm{exp}\left(i_a^b𝑑x\left[k\mathrm{\Lambda }^1\frac{1}{2}\mathrm{\Lambda }^{1/2}\left(\widehat{t}\widehat{V}+\widehat{V}\widehat{t}\right)\mathrm{\Lambda }^{1/2}\right]\right).$$
(A6)
To compute the disorder averages we assume that $`V_j(x)`$ are independent delta-correlated Gaussian random variables with zero average, $`\overline{V_i(x)V_j(y)}=\mathrm{\Delta }\delta _{ij}\delta (xy)`$. Then, a calculation at a finite temperature $`T`$ gives
$$\overline{G^{zz}}=\frac{L}{24}\left(\omega ^2+4\pi ^2T^2\right)_\pi ^\pi \frac{dq}{2\pi }\frac{t^2(q)}{\mathrm{\Lambda }^3(q)}\frac{2\stackrel{~}{\mathrm{\Delta }}(q)}{\omega ^2+\stackrel{~}{\mathrm{\Delta }}^2(q)},$$
(A7)
where $`\mathrm{\Lambda }(q)=v+2\lambda \mathrm{cos}q`$, $`t(q)=2\lambda \mathrm{sin}q`$, and
$$\stackrel{~}{\mathrm{\Delta }}(q)=\frac{\mathrm{\Delta }}{4}\left[2\mathrm{\Lambda }(q)+1+\frac{\mathrm{\Lambda }^2(q)}{\sqrt{v^2\lambda ^2}}\right].$$
As in the physical marginal cases, 020 and 131, at zero temperature and in the absence of disorder the conductivity vanishes. In the leading order in $`\lambda `$ Eq. (A7) is identical to the result for 020 with delta-correlated disorder (103). However, in the clean limit we find $`G^{zz}\delta (\omega )`$ for the model problem, which is not the correct behavior for the 020 and 131 states. The difference arises because, in contrast to 020 and 131, in the model problem the current operator $`I^z`$ (A3) commutes with the full Hamiltonian (A2).
## B Exact Solutions for $`N=4`$
In this appendix we describe exact solutions for the spectrum of the edge theory for the special case of $`N=4`$ layers with periodic boundary conditions. We begin with a discussion of the 010 state with both tunneling and interactions and then consider the 020 state. We remark that similar solutions exist for the 010, 020, and 131 states in $`N=3`$ layers.
### 1 010 State
The Hamiltonian for the 010 state is
$$_{010}^{N=4}=_{L/2}^{L/2}dx[\frac{v}{4\pi }:(_xu_i)^2:+\frac{g}{2\pi }:_xu_i_xu_{i+1}:+\frac{\lambda }{2\pi a}(e^{i(u_iu_{i+1})}+\mathrm{h}.\mathrm{c}.)],$$
(B1)
where $`u_5u_1`$. A straightforward fermionization $`\psi _i=e^{iu_i(x)}/\sqrt{2\pi a}`$ would yield a quartic Hamiltonian because of the density-density interaction term. We thus perform the orthogonal transformation $`u_i(x)=\mathrm{\Lambda }_{ij}\varphi _j(x)`$ where
$$\mathrm{\Lambda }=\frac{1}{2}\left(\begin{array}{cccc}1& 1& 1& 1\\ 1& 1& 1& 1\\ 1& 1& 1& 1\\ 1& 1& 1& 1\end{array}\right).$$
(B2)
The Hamiltonian (B1) then reads
$`_{010}^{N=4}`$ $`=`$ $`{\displaystyle _{L/2}^{L/2}}dx[{\displaystyle \frac{1}{4\pi }}v_i:(_x\varphi _i)^2:`$ (B4)
$`+{\displaystyle \frac{\lambda }{2\pi a}}(e^{i(\varphi _2\varphi _3)}+e^{i(\varphi _3\varphi _4)}+e^{i(\varphi _2+\varphi _3)}+e^{i(\varphi _3+\varphi _4)}+\mathrm{h}.\mathrm{c})],`$
where $`v_{1,3}=v\pm 2g`$ and $`v_2=v_4=v`$. The original fields, $`u_i`$, are all radius $`R=1`$ chiral bosons. We restrict ourselves to the sector where the total topological charge of the fields $`u_i`$ is constant, which corresponds to the conservation of the total fermion number in the fermionic form of (B1). Then the requirement that the topological charges of the original ($`u_i`$) and transformed ($`\varphi _i`$) bosons be integral is consistent with taking the fields $`\varphi _i`$ to all have unit radius as well. We can therefore fermionize according to $`\eta _i(x)=e^{i\varphi _i(x)}/\sqrt{2\pi a}`$ and then express the fermions in terms of their Fourier components $`\eta _i(x)=_ke^{ikx}c_{ik}/\sqrt{L}`$. In terms of these mode operators the Hamiltonian (B4) is
$$_{010}^{N=4}=\underset{k}{}[v_ikc_{ik}^{}c_{ik}+\lambda (c_{2k}^{}c_{3k}+c_{3k}^{}c_{4k}+c_{2k}c_{3k}+c_{3k}^{}c_{4k}^{}+\mathrm{h}.\mathrm{c}.)],$$
(B5)
where the momentum $`k`$ takes values in $`(2\pi /L)(\mathrm{𝖹𝖹}+1/2)`$. This form of the edge theory is quadratic and hence exactly soluble via a Bogoliubov transformation. We find:
$$_{010}^{N=4}=\underset{k}{}ϵ_i(k)b_{ik}^{}b_{ik},$$
(B6)
where $`\{b_{ik},b_{jk^{}}^{}\}=\delta _{ij}\delta _{kk^{}}`$ are four independent branches of fermionic oscillators with dispersions
$$\begin{array}{cc}ϵ_1(k)=(v+2g)k,& ϵ_3(k)=(vg)k\sqrt{(gk)^2+4\lambda ^2},\\ ϵ_2(k)=vk,& ϵ_4(k)=(vg)k+\sqrt{(gk)^2+4\lambda ^2}\end{array}.$$
(B7)
We find that two of the branches develop curvature if and only if both $`g`$ and $`\lambda `$ are nonzero.
### 2 020 State
Now consider the 020 state whose bosonic Hamiltonian is given above in Eq. (54). The fields $`u_i`$ all have unit radius, but their commutators are not conventionally normalized, see Eq. (22). From our solution above for the 010 state we know that we can perform an orthogonal transformation (B2) and still have four radius one bosons $`\varphi _i`$. If we then define rescaled fields $`\theta _i(x)\varphi _i(x)/\sqrt{2}`$ the 020 Hamiltonian becomes
$`_{020}^{N=4}`$ $`=`$ $`{\displaystyle _{L/2}^{L/2}}dx[{\displaystyle \frac{v_i}{8\pi }}:(_x\theta _i)^2:`$ (B9)
$`+{\displaystyle \frac{\lambda }{(2\pi a)^2}}(e^{i\sqrt{2}(\theta _2\theta _3)}+e^{i\sqrt{2}(\theta _3\theta _4)}+e^{i\sqrt{2}(\theta _2+\theta _3)}+e^{i\sqrt{2}(\theta _3+\theta _4)}+\mathrm{h}.\mathrm{c}.)],`$
where $`\theta _i`$ are radius $`R=1/\sqrt{2}`$ chiral bosons with conventionally normalized commutators, and the velocities $`v_i`$ are the same as those given above in the solution of the 010 state. We can define a quartet of $`\widehat{su}(2)_1`$ KM currents as in Eq. (55). The Hamiltonian becomes
$$_{020}^{N=4}=_{L/2}^{L/2}dx[\frac{\pi v_i}{3}:[𝐉_i(x)]^2:+4\lambda (J_2^xJ_3^x+J_3^xJ_4^x)],$$
(B10)
where we have again used the identity (58). Next, we define a new set of currents $`\stackrel{~}{J}_i^a(x)=R^{ab}J_i^b(x)`$, which also obey independent $`\widehat{su}(2)_1`$ algebras, where $`R^{ab}\delta ^{az}\delta ^{bx}+\delta ^{ay}\delta ^{by}\delta ^{ax}\delta ^{bz}`$. This rotation leaves the first term in the Hamiltonian (B10) unchanged and in the second term gives $`J_i^x\stackrel{~}{J}_i^z`$. If we express the rotated KM currents, $`\stackrel{~}{J}_i^a`$, in terms of four new radius $`R=1/\sqrt{2}`$ chiral bosons $`\stackrel{~}{\theta _i}`$ via Eq. (55), we find
$$_{020}^{N=4}=_{L/2}^{L/2}𝑑x\frac{1}{4\pi }M_{ij}:_x\stackrel{~}{\theta }_i_x\stackrel{~}{\theta }_j:,$$
(B11)
where
$$M\left(\begin{array}{cccc}v/2+g& 0& 0& 0\\ 0& v/2& \lambda /\pi & 0\\ 0& \lambda /\pi & v/2g& \lambda /\pi \\ 0& 0& \lambda /\pi & v/2\end{array}\right).$$
(B12)
We have succeeded in finding a quadratic representation of the Hamiltonian which can be readily diagonalized by performing an orthogonal transformation from the fields $`\stackrel{~}{\theta _i}`$ to new fields $`\vartheta _i`$. The final form of the Hamiltonian is
$$_{020}^{N=4}=_{L/2}^{L/2}𝑑x\frac{1}{4\pi }\stackrel{~}{v}_i:(_x\vartheta _i)^2:,$$
(B13)
where the velocities are
$$\begin{array}{cc}\stackrel{~}{v}_1=v/2+g& \stackrel{~}{v}_3=v/2g/2+\sqrt{(g/2)^2+2(\lambda /\pi )^2}\\ \stackrel{~}{v}_2=v/2& \stackrel{~}{v}_4=v/2g/2\sqrt{(g/2)^2+2(\lambda /\pi )^2}\end{array}.$$
(B14)
The exact spectrum of the 020 state contains four independent free chiral bosons whose velocities are renormalized by both interactions and tunneling.
## C Klein Factors
In this appendix we discuss Klein factors which are needed to produce the proper anticommutation relations between different species of fermions. One place where Klein factors are needed is in the definition of the physical electron operators at the edges of each layer. We modify the definition of the electron operator at the edge of layer $`i`$ (23) to read
$$\mathrm{\Psi }_i(x)e^{iA_{ij}N_j}e^{iu_i(x)},$$
(C1)
where $`A`$ is an $`N\times N`$ matrix and $`N_j`$ is the topological charge of the Bose field $`u_j(x)`$ defined by
$$N_j=\frac{1}{2\pi }_{L/2}^{L/2}𝑑x_xu_j(x).$$
(C2)
Using the commutation relations of the chiral bosons (22) one can readily show
$$\mathrm{\Psi }_i(x)\mathrm{\Psi }_j(x^{})=\mathrm{\Psi }_j(x^{})\mathrm{\Psi }_i(x)e^{i\pi K_{ij}\mathrm{sgn}(xx^{})}e^{i(A_{ik}K_{kj}A_{jk}K_{ki})},$$
(C3)
and the combination appearing in the tunneling operator (35) can be written
$$\mathrm{\Psi }_i(x)\mathrm{\Psi }_{i+1}^{}(x)e^{i(A_{ik}A_{i+1,k})N_k}e^{i[u_i(x)u_{i+1}(x)]}e^{i(A_{i+1,k}K_{k,i+1}A_{i+1,k}K_{ki})}.$$
(C4)
If we write $`A=\pi BK^1`$, for some matrix $`B`$, then we have
$`\mathrm{\Psi }_i(x)\mathrm{\Psi }_j(x^{})`$ $`=`$ $`\mathrm{\Psi }_j(x^{})\mathrm{\Psi }_i(x)e^{i\pi K_{ij}}e^{i\pi (B_{ij}B_{ji})},`$ (C5)
$`\mathrm{\Psi }_i(x)\mathrm{\Psi }_{i+1}^{}(x)`$ $``$ $`e^{i[u_i(x)u_{i+1}(x)]}e^{i\pi (B_{i+1,i+1}B_{i+1,i})}e^{i\pi (B_{il}B_{i+1,l})(K^1)_{lk}N_k}.`$ (C6)
From Eq. (C5) we see that for the electron operators in different layers to anticommute we must have $`(B_{ij}B_{ji})\mathrm{𝖹𝖹}_{\mathrm{even}}`$ if $`K_{ij}`$ is odd and $`(B_{ij}B_{ji})\mathrm{𝖹𝖹}_{\mathrm{odd}}`$ if $`K_{ij}`$ is even. From Eq. (C6) we see that if the elements of the matrix $`B`$ are integers and $`(B_{il}B_{i+1,l})(K^1)_{lk}`$ is an integer for all $`i`$ and $`k`$, then the tunneling operator will be the same as in the absence of the Klein factors up to a sign which depends on the topological charge sector.
As a concrete example, consider the multilayer 131 state. If the number of layers $`N`$ is such that $`\mathrm{det}(K)`$ is odd (i.e., $`N\mathrm{mod}\mathrm{\hspace{0.17em}3}2`$) then the choice
$$B_{ij}=\{\begin{array}{cc}\mathrm{det}(K)& \mathrm{for}ji2\\ 0& \mathrm{otherwise}\end{array}$$
(C7)
gives $`\{\mathrm{\Psi }_i(x),\mathrm{\Psi }_j(x^{})\}=0`$ for all $`i,j`$ and modifies the tunneling term by a factor of $`\pm 1`$. Note that in correlation functions, such as the current two-point function used to compute the conductivity (83), the tunneling operator is always accompanied by its hermitian conjugate and therefore factors such as $`e^{i\pi (B_{il}B_{i+1,l})(K^1)_{lk}N_k}`$ will cancel.
In the above discussion we have constructed Klein factors out of the topological charge operators of the Bose fields in order to give the correct anticommutation relations between the physical electron operators without significantly modifying the tunneling Hamiltonian. When considering the fermionization of the multilayer edge theory, as discussed in Section VI, we have a different goal. In addition to ensuring that the electron operators anticommute we want the different species in the fermionized Hamiltonian to obey proper anticommutation relations. In this case an alternative approach to the one described above proves advantageous. This procedure is best described by considering a specific case, for example the 131 state. One can introduce $`2N`$ unitary operators $`F_{i+\alpha }`$, where $`\alpha =0,1/2`$ and $`i=1,\mathrm{},N`$, by formally enlarging the Hilbert space. These operators commute with the Bose fields $`u_i`$ and obey
$$\{F_{i+\alpha },F_{j+\alpha ^{}}\}=\{F_{i+\alpha }^{},F_{j+\alpha ^{}}^{}\}=\{F_{i+\alpha },F_{j+\alpha ^{}}^{}\}=0,$$
(C8)
for $`i+\alpha j+\alpha ^{}`$. The electron operators are modified according to
$$\mathrm{\Psi }_i(x)F_{i1/2}F_iF_{i+1/2}e^{iu_i(x)}.$$
(C9)
It is a straightforward exercise to demonstrate that Eqns. (C8) and (C9) imply that electron operators in different layers properly anticommute. The tunneling operator is now
$$\mathrm{\Psi }_j(x)\mathrm{\Psi }_{j+1}^{}(x)e^{i[u_j(x)u_{j+1}(x)]}F_{j1/2}F_jF_{j+1}^{}F_{j+3/2}^{}.$$
(C10)
If we then follow the procedure described in Section VI with the modification that the new fermion species each absorb a factor of $`F`$:
$$\psi _{j+\alpha }\frac{1}{\sqrt{2\pi a}}F_{j+\alpha }e^{i\phi _{j+\alpha }(x)},$$
(C11)
we arrive at the same fermionized Hamiltonian (137), but now with the different fermion species properly anticommuting.
## D Proof that $`\mathrm{}e\sigma ^{zz}(\omega )\delta (\omega )[I^z,H]=0`$
In this appendix we prove that the real part of the $`z`$-axis conductivity is proportional to a delta function in frequency if and only if the current operator commutes with the full Hamiltonian of the edge theory. The basic connection is of course the Kubo formula (52):
$$\mathrm{}e\sigma ^{zz}(\omega )=\frac{1}{2\omega }\frac{1}{NLd}_{\mathrm{}}^{\mathrm{}}𝑑te^{i\omega t}[I^z(t),I^z(0)].$$
(D1)
One of the directions is trivial: if $`[I^z,H]=0`$, then $`I^z(t)`$ is independent of time, and hence the commutator in Eq. (D1) vanishes identically implying that for all $`\omega 0`$, $`\mathrm{}e\sigma ^{zz}(\omega )=0`$ and therefore $`\mathrm{}e\sigma ^{zz}\delta (\omega )`$.
We now prove the converse, i.e., we assume $`\mathrm{}e\sigma ^{zz}\delta (\omega )`$. Inverting the Fourier transform in Eq. (D1) we find
$$[I^z(t),I^z(0)]=2NLd_{\mathrm{}}^{\mathrm{}}𝑑\omega e^{i\omega t}\omega \mathrm{}e\sigma ^{zz}(\omega ).$$
(D2)
Using our assumption in this equation gives
$$[I^z(t),I^z(0)]=\underset{n}{}e^{\beta E_n}n|[I^z(t),I^z(0)]|n=0,$$
(D3)
where $`|n`$ and $`E_n`$ are eigenstates and eigenvalues of $`H`$. Note that the sum in Eq. (D3) is a linear combination of terms $`n|[I^z(t),I^z(0)]|n`$ with non-negative coefficients which depend on $`\beta `$. Since the sum must vanish for arbitrary $`\beta `$ we conclude that
$$n|[I^z(t),I^z(0)]|n=0,n.$$
(D4)
Writing $`I^z(t)=\mathrm{exp}(iHt)I^z(0)\mathrm{exp}(iHt)`$ and inserting a summation over a complete set of states $`\{|m\}`$, Eq. (D4) is equivalent to
$$\underset{m}{}\mathrm{sin}[(E_nE_m)t]|n|I^z(0)|m|^2=0,n,t.$$
(D5)
The quantities $`|n|I^z(0)|m|^2`$ are clearly non-negative. Since the sum in Eq. (D5) must vanish for all $`n`$ and at all times $`t`$ we can conclude
$$|n|I^z(0)|m|^2=0,nm.$$
(D6)
This equation says that the current has no off-diagonal matrix elements in the basis of energy eigenstates, i.e., the current operator $`I^z`$ is diagonal in the basis of eigenstates of $`H`$. Since this means $`I^z`$ and $`H`$ can be simultaneously diagonalized it implies that they commute: $`[I^z(0),H]=0`$, completing the proof.
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# I. Introduction
## I. Introduction
The canonical commutation relations (henceforth the CCR) were initially introduced in 1927 by Dirac as generalizations of Heisenberg’s commutation relations in order to discuss radiation theory . Since then, the CCR have proven to be of central importance in the quantum description of bosonic systems , i.e. systems of identical particles satisfying Bose-Einstein statistics, such as the photons Dirac considered. The most elementary formulation of the CCR for $`n`$ degrees of freedom is
$$[q_j,p_k]q_jp_kp_kq_j=i\delta _{jk}\text{1I},j,k\{1,\mathrm{},n\},$$
where $`\delta _{jk}`$ is the Kronecker delta, 1I is the identity operator on the complex Hilbert space $`𝒦`$, and $`q_j,p_k`$ are symmetric operators on $`𝒦`$. It is well-known that the operators $`q_j,p_k`$ must be unbounded, and so the CCR are to be understood on a suitable dense subspace of $`𝒦`$. Under certain circumstances, these operators may be exponentiated to obtain unitaries
$$U_k(a)=e^{iap_k}\text{and}V_j(a)=e^{iaq_j},$$
with $`a\text{I R}`$, satisfying
$`U_k(a)V_j(b)`$ $`=e^{iab}V_j(b)U_k(a),j=k,`$
$`U_k(a)V_j(b)`$ $`=V_j(b)U_k(a),jk,`$
for all $`a,b\text{I R}`$, $`j,k=1,\mathrm{},n`$. These relations constitute the Weyl form of the CCR for $`n`$ degrees of freedom. These unitaries generate a $`C^{}`$-algebra on $`𝒦`$. A representation of (the Weyl form of) the CCR is a $`C^{}`$-homomorphism preserving the Weyl relations. A remarkable early result about the CCR is the Stone-von Neumann uniqueness theorem : all irreducible representations of the CCR for $`n`$ degrees of freedom are unitarily equivalent.
When rigorous mathematical analysis of the CCR in the case of infinitely many degrees of freedom — the case of relevance to quantum statistical mechanics and quantum field theory — began in the 1950’s, it was quickly realized that there are uncountably infinitely many unitarily inequivalent irreducible representations of the CCR in this case and that the choice of proper representation is crucial in any physical application. This last point deserves emphasis. It has become clear from rigorous study of concrete models in constructive quantum field theory that bosonic systems with identical kinematics but physically distinct dynamics require inequivalent representations of the CCR. Roughly speaking, the kinematical aspects determine the choice of CCR-algebra, whereas the dynamics fix the choice of the representation of the given CCR-algebra in which to make the relevant, perturbation-free computations. (It is also believed — and proven in a number of indicative special cases — that perturbation series in one representation provide divergent and at best asymptotic approximations to the physically relevant quantities in another, unitarily inequivalent representation.) For an overview of these matters and further references, the reader is referred to .
Many classes of representations of the CCR for infinitely many degrees of freedom have been rigorously studied in the literature - we mention, in particular, the Fock , coherent , quasifree (or symplectic), infinite product , and, more recently, quadratic and polynomial representations . In this paper, we wish to construct and prove certain mathematically and physically relevant properties of a broad class of representations of the CCR containing all of the above (except possibly the infinite product representations) as special cases.
We suggest heuristically the nature of these representations with the following easily describable examples. Let $`\{q_k,p_k\}_{k=1}^{\mathrm{}}`$ be a system of densely defined operators acting on the complex Hilbert space $`𝒦`$ and satisfying
$$[q_j,p_k]=i\delta _{jk}\text{1I}$$
on a suitable dense subset of $`𝒦`$, and for the standard annihilation and creation operators $`a_k\frac{1}{\sqrt{2}}(q_k+ip_k)`$ and $`a_k^{}\frac{1}{\sqrt{2}}(q_kip_k)`$ let there be a vector $`\mathrm{\Omega }𝒦`$ such that $`a_k\mathrm{\Omega }=0`$ for all $`k\text{I N}`$. In other words, we consider the Fock representation<sup>1</sup><sup>1</sup>1Often called the Fock-Cook representation, since it was in that this representation was given a mathematically rigorous and relativistically covariant form. of a bosonic system with infinitely many degrees of freedom <sup>2</sup><sup>2</sup>2in the so-called basis-dependent formulation — We discuss the basis-independent formulation and discuss their relation in Chapter II. (see ). In this setting, the coherent, resp. quasifree, quadratic or polynomial canonical transformations can be written as
$`(coherent)q_kq_k,`$ $`p_kp_k+\lambda _k\text{1I},`$
$`(quasifree)q_kq_k,`$ $`p_kp_k+\underset{k_1}{{\displaystyle }}\lambda _{kk_1}q_{k_1},`$
$`(quadratic)q_kq_k,`$ $`p_kp_k+\underset{k_1,k_2}{{\displaystyle }}\lambda _{kk_1k_2}:q_{k_1}q_{k_2}:,`$
$`(polynomial)q_kq_k,`$ $`p_kp_k+\underset{k_1,\mathrm{},k_n}{{\displaystyle }}\lambda _{kk_1\mathrm{}k_n}:q_{k_1}\mathrm{}q_{k_n}:,`$
with all coefficients totally symmetric in the indices.<sup>3</sup><sup>3</sup>3In point of fact, the polynomial representations of the CCR constructed in were constructed globally and not via their infinitesimal generators as is done in this paper. The canonical transformations of general degree we construct in this paper are of the form
$$q_kq_k,p_kp_k+\underset{m,k_1,\mathrm{},k_m}{}\lambda _{kk_1\mathrm{}k_m}:q_{k_1}\mathrm{}q_{k_m}:.$$
$`1.1`$
We emphasize that this generalization opens up a much larger class of representations of the CCR — representations, which involve an extremely singular perturbation of the original Fock representation. Aside from the mathematical interest in this extension of the representation theory of an important class of $`C^{}`$-algebras associated with the CCR, there are advantages for the quantum theory of bosonic systems with infinitely many degrees of freedom. As alluded to earlier, such representations allow rigorous treatment of interacting systems with extremely singular “bare” interaction. One may therefore handle systems with new classes of dynamics (for a suggestive treatment of the connection between representation and dynamics, see the discussion in ). This will entail that the divergences which typically arise in perturbation theory can be avoided, since the very need for a perturbation expansion from the “wrong” representation is obviated. Furthermore, as briefly shown in the final section of , but which is well-known to theoretical physicists, the Hamilton operator, which represents the total energy of the system, can often be simplified and in certain situations even be diagonalized by such canonical transformations. This leads to significant computational advantages. We postpone addressing these issues in detail to a later publication.
In Chapter II, the general setting of this work will be more precisely specified. In particular, we define and collect necessary facts about the $`C^{}`$-algebras associated with the CCR and about their Fock representations. The new transformations of general degree will be constructed in Section 3.1, and those which lead to regular representations of the CCR will be characterized in a number of different manners — see Theorems 3.1.5 and 3.1.6. We generalize a result in by showing that the general degree transformations yielding regular representations of the CCR can be brought into the transparent form given in equation (1.1), which we call the standard form of such transformations. This fact proves to be useful in the proof of some of the later results.
Convenient sufficient conditions assuring the irreducibility of such representations are given in Section 3.2 — see Theorem 3.2.2. However, we show in Section 3.2 that there do exist such unrestricted order representations which are reducible. Chapter IV is given over to a proof that the class of representations constructed in Section 3.1 subsumes the previously studied classes discussed above. In particular, it is explicitly shown that every (irreducible) quasifree or coherent representation of the CCR is unitarily equivalent to one of our unrestricted order representations (Theorem 4.5), since it is already evident that the quadratic and polynomial representations are included. In the process, we provide a new characterization of quasifree states on CCR-algebras (Proposition 4.4).
In Chapter V we restrict our attention to the computationally more manageable canonical transformations of finite degree and prove necessary and sufficient conditions for such representations to be quasi-equivalent to Fock (Theorem 5.6), coherent or quasifree (Theorem 5.10 and Corollary 5.11) representations. For our purposes, the primary interest of these results is in the violation of the identified necessary and sufficient conditions — for then one has representations of bosonic systems of infinitely many degrees of freedom which describe new physics, i.e. which model bosonic systems that cannot be described by the earlier classes of representations. However, in this paper we do not try to study the new physical content of these representations. As an aside, the previously known conditions for unitary equivalence of quasifree states follow as a special case of the results in Chapter V — see Theorem 5.12. The significantly less transparent conditions for unitary equivalence with the quadratic representations constructed in will not be given here, though they are known to us. (See the end of Chapter V for a brief discussion.)
As this paper is an extension of and employs a number of the results and arguments from that paper, we shall maintain the same notational conventions, as detailed in the next chapter. Early versions of some of the results presented in this paper appeared in the Diplomarbeit .
## II. Notation and General Setting
We begin with an arbitrary real nondegenerate symplectic space $`(H,\sigma )`$ with an associated regular Weyl system $`(𝒦,W(f))`$ consisting of a complex Hilbert space $`𝒦`$ and a mapping $`W:H𝒰(𝒦)`$ from $`H`$ into the group $`𝒰(𝒦)`$ of unitary operators on $`𝒦`$ which satisfies the following axioms :
$$W(f)W(g)=e^{i\sigma (f,g)/2}W(f+g),f,gH,$$
$`2.1`$
$$W(f)^{}=W(f),fH,$$
$`2.2`$
and
$$\text{I R}tW(tf)(𝒦)\text{is weakly continuous for all}fH.$$
$`2.3`$
($`(𝒦)`$ denotes the set of all bounded linear operators $`A:𝒦𝒦`$.) Condition (2.3) entails that the map $`\text{I R}tW(tf)(𝒦)`$ is actually strongly continuous, hence by Stone’s Theorem one knows that for each $`fH`$ there exists a self-adjoint operator $`\mathrm{\Phi }(f)`$ on $`𝒦`$ such that $`W(tf)=e^{it\mathrm{\Phi }(f)}`$ for all $`t\text{I R}`$ and by (2.1) the map $`f\mathrm{\Phi }(f)`$ is (real) linear. In fact, there exists a dense domain of vectors $`D_W𝒦`$ which is a core of and left invariant by every $`\mathrm{\Phi }(f)`$ ; it is on this domain that the linearity just mentioned can be verified. On this domain one also verifies that the generators satisfy the CCR:
$$\mathrm{\Phi }(f)\mathrm{\Phi }(g)\mathrm{\Phi }(g)\mathrm{\Phi }(f)=i\sigma (f,g)\text{1I},f,gH.$$
$`2.4`$
We therefore also call the Weyl system $`(𝒦,W(f))`$ and its associated generators as above a regular representation of the CCR over $`(H,\sigma )`$. In the physical literature such infinitesimal generators $`\mathrm{\Phi }(f)`$ are called field operators, and we shall use this language in the following.
We shall denote by $`𝒜(H,\sigma )`$ the $`C^{}`$-algebra on $`𝒦`$ generated by the operators $`\{W(f)fH\}`$. As the notation indicates, the algebra $`𝒜(H,\sigma )`$ does not depend on the choice of representation of the Weyl operators $`\{W(f)fH\}`$ ( or Theorem 5.2.8 in ). $`𝒜(H,\sigma )`$ is a simple $`C^{}`$-algebra<sup>4</sup><sup>4</sup>4The fact that the CCR-algebra is simple can be seen as the correct generalization of the Stone-von Neumann uniqueness theorem to the case of infinitely many degrees of freedom. Indeed, since $`𝒜(H,\sigma )`$ is simple, all of its representations are isomorphic. When $`H`$ is finite-dimensional, this isomorphism is unitarily implementable, entailing the result in . and is nonseparable if $`H`$ is infinite-dimensional . There are, in fact, many different $`C^{}`$-algebras one can associate with the CCR (see for a discussion of some of the alternatives), and the one we have chosen is minimal in the sense of set containment ; but, for practical purposes the choice is immaterial, since one is generally interested in a von Neumann algebra which is ‘generated’ by the $`C^{}`$-algebra, and all the $`C^{}`$-algebras discussed in (realized concretely on a given representation space) have the same weak closure. In the following it shall be understood that one choice of $`\sigma `$ has been fixed, and we shall write $`𝒜(H)`$ instead of $`𝒜(H,\sigma )`$.
For any real linear map $`J:HH`$ satisfying $`\sigma (Jf,Jg)=\sigma (f,g)`$, $`\sigma (Jf,f)>0`$ ($`f0`$), and $`J^2=\text{1I}`$, one can introduce a complex structure on $`H`$ as follows <sup>5</sup><sup>5</sup>5It should be mentioned that such a $`J`$ does not necessarily exist for arbitrary choice of symplectic space $`(H,\sigma )`$ . : $`(\alpha +i\beta )f\alpha f+\beta Jf`$, for all $`fH`$, and $`\alpha ,\beta \text{I R}`$. Moreover, $`f,g_{}\sigma (Jf,g)+i\sigma (f,g)`$ defines a scalar product on $`H`$ such that $`(H,,_{})`$ is a complex preHilbert space , the completion of which we shall denote by $``$. If, on the other hand, one begins with a complex Hilbert space $``$ with scalar product $`,_{}`$, then with $`\sigma (f,g)=\mathrm{}mf,g_{}`$, $`(,\sigma )`$ is a real symplectic space of the sort with which we began, and with $`f,g\mathrm{}ef,g_{}`$, then $`(,,)`$ is a real Hilbert space.
With a choice of a $`\sigma `$-admissible complex structure $`J`$ on $`(H,\sigma )`$, there is an important representation of $`𝒜(H)`$ called the Fock representation. This is given as the GNS-representation $`(𝒦,\pi _J,\mathrm{\Omega })`$ associated with the state $`\omega _J`$ determined on $`𝒜(H)`$ by
$$\omega _J(W(f))e^{\sigma (Jf,f)/4}.$$
$`2.5`$
Given such a representation, one can define the following ‘annihilation’ and ‘creation’ operators:
$$a(f)\frac{1}{\sqrt{2}}(\mathrm{\Phi }(f)+i\mathrm{\Phi }(Jf)),a^{}(f)\frac{1}{\sqrt{2}}(\mathrm{\Phi }(f)i\mathrm{\Phi }(Jf)),$$
$`2.6`$
where $`\pi _J(W(tf))=e^{it\mathrm{\Phi }(f)}`$. One has then $`a(f)\mathrm{\Omega }=0`$ for all $`fH`$.
For the purposes of this paper, we shall assume that a choice of $`\sigma `$-admissible complex structure $`J`$ has been made on $`(H,\sigma )`$ and held fixed, so we have the complex one-particle space $``$ and a corresponding Fock representation. Since as sets $`H=`$, we shall distinguish notationally the vector $`f`$ viewed as an element of the real Hilbert space $`H`$ from the same vector, denoted as $`\stackrel{~}{f}`$, viewed as an element of the complex Hilbert space $``$. If $`\{\stackrel{~}{e}_k\}_{k\text{I N}}`$ forms an orthonormal basis in $``$, then the set $`\{e_k,Je_k\}_{k\text{I N}}`$ forms a symplectic orthonormal system in $`H`$, in particular
$`\sigma (e_k,e_l)=\sigma (Je_k,Je_l)=0`$ $`,\sigma (e_k,Je_l)=\delta _{kl},`$
$`e_k,e_l=Je_k,Je_l=\delta _{kl}`$ $`,e_k,Je_l=0.`$
In this paper, whenever a choice of symplectic orthonormal basis $`\{e_k,Je_k\}_{k\text{I N}}`$ has been made, the Hilbert subspace of $`H`$ generated by $`\{e_k\}_{k\text{I N}}`$ will be denoted by $`V`$.
With the choice of complex structure made as above, the corresponding Fock state $`\omega _J:𝒜(H)\text{ }\text{C}`$ now satisfies
$$\omega _J(W(f))=e^{f^2/4},$$
and the associated GNS-space may be represented by the symmetric Fock space $`_+()`$. We recall that the Fock space $`()=_{n=0}^{\mathrm{}}^n`$ ($`^0=\text{ }\text{C}`$, $`^n`$ is the $`n`$-fold tensor product of $``$ with itself), and that $`_+()`$ is the totally symmetric subspace $`_{n=0}^{\mathrm{}}P_+^n`$ of $`()`$, where $`P_+`$ is the projection
$$P_+(\stackrel{~}{f}_1\stackrel{~}{f}_2\mathrm{}\stackrel{~}{f}_n)=\frac{1}{n!}\underset{\pi }{}\stackrel{~}{f}_{\pi (1)}\stackrel{~}{f}_{\pi (2)}\mathrm{}\stackrel{~}{f}_{\pi (n)}$$
($`\stackrel{~}{f}_i`$, $`\pi S_n`$, the group of permutations on the set $`\{1,2,\mathrm{},n\}`$). The projection operator $`P_n:()^n`$ projects onto the n-particle subspace. The vector $`\mathrm{\Omega }(1,0,0,\mathrm{})_+()`$ is the Fock vacuum. $`_0`$ is the finite-particle subspace, i.e. the linear span of the ranges of $`\{P_n\}_{n\text{I N}_0}`$. The elements of $`_0`$ will be called finite-particle vectors. As is well-known, the GNS-representation for the Fock state may be identified with this Fock space representation.
With $`a(\stackrel{~}{f})`$ the usual annihilation operator in $`_+()`$ for $`\stackrel{~}{f}`$ and the adjoint operator $`a^{}(\stackrel{~}{f})`$ (an extension of) the corresponding creation operator, then the linear self-adjoint operator $`\mathrm{\Phi }_S(\stackrel{~}{f})\frac{1}{\sqrt{2}}(\overline{a(\stackrel{~}{f})+a^{}(\stackrel{~}{f})})`$ (the bar denotes the closure of the operator) is called the Segal field operator in $`_+()`$. If we also view $`\stackrel{~}{f}`$ as an element of $`H`$, then we have $`\mathrm{\Phi }_S(\stackrel{~}{f})=\mathrm{\Phi }(f)`$, after identifying $`𝒜(H)`$ with $`\pi _J(𝒜(H))`$ (since they are isomorphic). $`_0`$ is a core for $`\mathrm{\Phi }(f)`$, and for $`\phi _0`$, $`\mathrm{\Phi }(f)`$ satisfies the bound $`\mathrm{\Phi }(f)\phi \sqrt{2}\sqrt{n_\phi +1}f\phi `$, where $`n_\phi `$ equals the smallest $`n\text{I N}`$ such that $`P_N\phi =0`$ for all $`N>n`$. If $`\{f_n\}_{n=1}^{\mathrm{}}`$ converges in $`H`$ to $`f`$, then for every $`\phi _0`$ the sequence $`\{\mathrm{\Phi }(f_n)\phi \}_{n=1}^{\mathrm{}}`$ converges in $`_+()`$ to $`\mathrm{\Phi }(f)\phi `$.
To make a notational connection to the discussion in Chapter I and in keeping with common harmonic oscillator conventions, if $`\{e_k,Je_k\}`$ is a symplectic orthonormal basis in $`H`$, then the ‘position operator’ and the ‘momentum operator’ corresponding to the $`k`$-th degree of freedom are given by
$$q_k\mathrm{\Phi }_S(\stackrel{~}{e}_k)=\mathrm{\Phi }(e_k),p_k\mathrm{\Phi }_S(i\stackrel{~}{e}_k)=\mathrm{\Phi }(Je_k).$$
To elucidate what is meant in the following by linear canonical transformations, we recall two well-studied classes of representations of the CCR — the coherent and the quasifree representations. If $`(𝒦,W(f))`$ is a representation of $`𝒜(H)`$ and $`l:H\text{I R}`$ is a linear map, then
$$\pi _l(W(f))\widehat{W}(f)e^{il(f)}W(f),fH,$$
determines a representation of $`𝒜(H)`$ generally called a coherent representation if $`(𝒦,W(f))`$ is a Fock representation (see, e.g. ). This leads to the relationship $`\widehat{\mathrm{\Phi }}(f)=\mathrm{\Phi }(f)+l(f)\text{1I},fH`$, between the generators of the representations. One may equivalently start with a Fock state $`\omega _J`$ on $`𝒜(H)`$ with associated representation $`(𝒦,\pi _J)`$ and define a coherent state $`\omega _l`$ by
$$\omega _l(W(f))\omega _J(W(f))e^{il(f)},fH.$$
Of course, the GNS representation of $`𝒜(H)`$ corresponding to $`\omega _l`$ is given on $`𝒦`$ by
$$\pi _l(W(f))e^{il(f)}\pi _J(W(f)),fH.$$
Quasifree (or symplectic) representations can be obtained from a given Fock representation of the CCR-algebra $`𝒜(H)`$ by a Bogoliubov transformation
$`\mathrm{\Phi }_T(f)`$ $`\mathrm{\Phi }(Tf)=\mathrm{\Phi }(f)+\mathrm{\Phi }((T\text{1I})f),`$
$`(\text{or equivalently}\pi _T(W(f))`$ $`\pi _J(W(Tf)),`$
using a symplectic operator $`T`$, i.e. one leaving the symplectic form invariant: $`\sigma (Tf,Tg)=\sigma (f,g)`$, for all $`f,gH`$. The canonical transformations associated with the coherent and quasifree representations constitute the inhomogeneous linear group of canonical transformations, studied by Shale and Berezin , among many others.
A crucial technical tool in this paper is the use of $`Q`$-space techniques. Let $`\{e_k,Je_k\}`$ be a symplectic orthonormal basis. Moreover, let $`x=(x_1,x_2,\mathrm{})`$ be a point in $`Q\times _{k=1}^{\mathrm{}}\text{I R}`$, and $`\mathrm{\Sigma }`$ be the $`\sigma `$-algebra generated by the cylinder sets of $`Q`$ with Lebesgue measurable base. Then $`\mu =_{k=1}^{\mathrm{}}\mu _k`$, where each $`\mu _k`$ is the Gaussian measure $`d\mu _k=\pi ^{\frac{1}{2}}e^{x_k^2}dx_k`$, is a probability measure on $`(Q,\mathrm{\Sigma })`$. It is well-known that there exists a unitary map $`S`$ of $`_+()`$ onto $`L^2(Q,d\mu )`$ such that
$`S\mathrm{\Omega }_0=1`$ $`\text{and}SP_+(\stackrel{~}{e}_{k_1}\stackrel{~}{e}_{k_2}\mathrm{}\stackrel{~}{e}_{k_r})=(r!)^{\frac{1}{2}}(\sqrt{2})^r:x_{k_1}x_{k_2}\mathrm{}x_{k_r}:,`$
$`S\mathrm{\Phi }(e_k)S^1`$ $`q_k=x_k\text{1I}=\text{(multiplication by)}x_k,\text{and}`$
$`S\mathrm{\Phi }(Je_k)S^1`$ $`p_k={\displaystyle \frac{1}{i}}{\displaystyle \frac{}{x_k}}+ix_k\text{1I},`$
where the operator equations are understood to hold on the dense set $`S_0`$. We shall drop the symbol $`S`$, when the identification is clear.
Let $`\underset{¯}{k}=\{k_1,k_2,\mathrm{},k_r\}`$ be a multiple index in $`\text{I N}^r=\times _{j=1}^r\text{I N}`$, $`\{\lambda _{\underset{¯}{k}}\}_{\underset{¯}{k}\text{I N}^r}`$ a sequence of real numbers, totally symmetric in $`\underset{¯}{k}`$ , $`_{\underset{¯}{k}}\lambda _{\underset{¯}{k}}^2<\mathrm{}`$, and
$`I_n=\{\underset{¯}{k}\mathrm{max}\{k_1,k_2,\mathrm{},k_r\}n\}`$. Then
$$f_n^{(r)}=\underset{\underset{¯}{k}I_n}{}\lambda _{\underset{¯}{k}}e_{k_1}e_{k_2}\mathrm{}e_{k_r}H^r,$$
with $`f_n^{(r)}^2=_{\underset{¯}{k}I_n}\lambda _{\underset{¯}{k}}^2`$. Since with $`m<n`$, we have $`f_n^{(r)}f_m^{(r)}^2=_{\underset{¯}{k}I_n\backslash I_m}\lambda _{\underset{¯}{k}}^2`$ and $`_{\underset{¯}{k}I_n}\lambda _{\underset{¯}{k}}^2`$ converges with $`n\mathrm{}`$, it follows that
$$f_n^{(r)}f^{(r)}=\underset{\underset{¯}{k}}{}\lambda _{\underset{¯}{k}}e_{k_1}e_{k_2}\mathrm{}e_{k_r}\text{in}H^r.$$
Consider the sequence of operators
$$A(f_n^{(r)})=\underset{\underset{¯}{k}I_n}{}\lambda _{\underset{¯}{k}}:\mathrm{\Phi }(e_{k_1})\mathrm{\Phi }(e_{k_2})\mathrm{}\mathrm{\Phi }(e_{k_r}):,$$
the $`Q`$–space realization of which is given by
$$A_n=A(f_n^{(r)})=\underset{\underset{¯}{k}I_n}{}\lambda _{\underset{¯}{k}}:x_{k_1}x_{k_2}\mathrm{}x_{k_r}:\text{1I}.$$
Note that since it is a polynomial, $`A_n=A_n(x)`$ is in $`L^2(Q,d\mu )`$ with
$`A_n^2`$ $`=2^rr!{\displaystyle \underset{\underset{¯}{k}I_n}{}}\lambda _{\underset{¯}{k}}^2=2^rr!f_n^{(r)}^2,`$
$`\text{and}A_nA_m^2`$ $`=2^rr!f_n^{(r)}f_m^{(r)}^2`$
( cf. the proof of Lemma I.18 of ). Therefore, $`A_n(x)`$ converges in $`L^2(Q,d\mu )`$, and we shall call the a.e.-defined limit $`A(x)=_{\underset{¯}{k}}\lambda _{\underset{¯}{k}}:x_{k_1}x_{k_2}\mathrm{}x_{k_r}:`$, which up to a factor of $`\sqrt{r!}(\sqrt{2})^r`$ corresponds to $`f^{(r)}`$ .
The advantage of the $`Q`$–space formulation is that all functions of the elements of $`\{\mathrm{\Phi }(e_k)k\text{I N}\}`$ become multiplication operators on $`L^2(Q,d\mu )`$. $`A_n(x)`$ and $`A(x)`$ are measurable, real-valued functions on $`Q`$ which are finite almost everywhere with respect to $`\mu `$. So with $`D(A)\{\phi A\phi L^2(Q,d\mu )\}`$ (similarly for $`D(A_n)`$), $`(A_n\phi )(x)=A_n(x)\phi (x)`$ and $`(A\phi )(x)=A(x)\phi (x)`$ are self-adjoint operators ( cf. , VIII.3 Proposition 1 ). Thus, for every $`f^{(r)}P_+V^r`$, $`A(f^{(r)})`$ represents a well-defined self-adjoint multiplication operator.
## III. A General Class of Representations of the CCR
In Section 3.1 we shall use $`Q`$-space techniques to construct our general class of representations of the CCR, as motivated above. These techniques were already employed to establish some of the results proven in about quadratic representations of the CCR. Our representations of general degree of the CCR will be defined in terms of linear maps $`\mathrm{\Lambda }`$ from the one-particle space $`H`$ into $`Q`$-space itself. We shall give a number of characterizations of those $`\mathrm{\Lambda }`$ which yield representations of the CCR, i.e. which determine canonical transformations of the field operators. It will be shown in Section 3.2 that when such $`\mathrm{\Lambda }`$ are bounded, then the resulting representation is irreducible. We shall also explain how unbounded $`\mathrm{\Lambda }`$ can lead to reducible representations.
###### 3.1 Canonical Transformations of Arbitrary Degree
We shall use the basic facts that the set $`\{1\}\{:x_{k_1}\mathrm{}x_{k_m}:k_1,\mathrm{},k_m,m\text{I N}\}`$ is an orthogonal basis in $`L^2(Q,d\mu )`$ and that
$$\underset{k_1,\mathrm{},k_m}{}\lambda _{k_1\mathrm{}k_m}:x_{k_1}\mathrm{}x_{k_m}:_2^2=\underset{k_1,\mathrm{},k_m}{}\lambda _{k_1\mathrm{}k_m}^2\frac{m!}{2^m},$$
$`\mathrm{3.1.1}`$
for arbitrary $`\lambda _{k_1\mathrm{}k_m}\text{I R}`$ symmetric in the indices $`k_1\mathrm{}k_m`$ (see Section 4.3 in ). Thus, standard arguments entail the following lemma.
###### Lemma 3.1.1
Let $`F=\underset{k_1,\mathrm{},k_m}{}c_{k_1\mathrm{}k_m}:x_{k_1}\mathrm{}x_{k_m}:L^2(Q,d\mu )`$ with complex numbers $`c_{k_1\mathrm{}k_m}`$ symmetric in the indices $`k_1\mathrm{}k_m`$. Then
$$F,:x_{k_1}\mathrm{}x_{k_m}:=c_{k_1\mathrm{}k_m}\frac{m!}{2^m},$$
and $`F`$ is real-valued (a.e. $`\mu `$) if and only if $`c_{k_1\mathrm{}k_m}\text{I R}`$ for all $`k_1,\mathrm{},k_m\text{I N}`$.
We also state a well-known fact about the $`Q`$-space representation of Fock space. Recall that for $`fV`$, $`x(f)\text{1I}=\mathrm{\Phi }(f)`$, given our stated convention of dropping the unitary $`S`$. The number operator $`N`$ on Fock space is a self-adjoint operator satisfying $`NP_n\phi =nP_n\phi `$, for all $`\phi ()`$ and commutes with $`P_+`$.
###### Lemma 3.1.2
The set $`𝒢`$ equal to the linear span of $`\{e^{ix(f)}fV\}`$ is a core for the number operator $`N`$ and for $`\mathrm{\Phi }(g)`$, given any $`gH`$.
In the next proposition we relate the operators $`\mathrm{\Phi }(f)`$ and $`\mathrm{\Phi }(f)+F\text{1I}`$, for any $`FL^2(Q,d\mu )`$, via a unitary transformation.
###### Proposition 3.1.3
If $`FL^2(Q,d\mu )`$ is real-valued ($`\mu `$ a.e.), then the operator $`\mathrm{\Phi }(f)+F\text{1I}`$, is essentially self-adjoint on $`L^{\mathrm{}}(Q,d\mu )`$, when $`fV`$, and on $`e^{iG}𝒢`$, when $`fV`$ ($`G`$ will be defined shortly). Furthermore, if $`fV`$, then the closure of the corresponding operator is given in terms of the self-adjoint $`\mathrm{\Phi }(f)`$ by
$$\overline{\mathrm{\Phi }(f)+F\text{1I}}=e^{iG}\mathrm{\Phi }(f)e^{iG},$$
$`\mathrm{3.1.2}`$
for any choice of ($`\mu `$ a.e.) real-valued $`GL^2(Q,d\mu )`$ such that $`G/x(f_2)=F`$, where $`f_2`$ is determined uniquely by $`f`$ by the decomposition $`f=f_1+Jf_2`$, $`f_1,f_2V`$.
###### Demonstration Proof
By Lemma 4.3.1 in , $`L^{\mathrm{}}(Q,d\mu )`$ is a core for the corresponding multiplication operator for every $`FL^2(Q,d\mu )`$. If $`fV`$, then the operator $`\mathrm{\Phi }(f)+F\text{1I}`$ is symmetric, $`D(\mathrm{\Phi }(f)+F\text{1I})=D(\mathrm{\Phi }(f))D(F\text{1I})`$, and, on $`L^{\mathrm{}}(Q,d\mu )L^2(Q,d\mu )`$, it is equal to the operator corresponding to multiplication by the $`L^2`$-function $`x(f)+F`$, i.e. $`\mathrm{\Phi }(f)+F\text{1I}`$ is essentially self-adjoint on $`L^{\mathrm{}}(Q,d\mu )`$.
One may assume that $`fV`$ and therefore $`f=Je_1+v`$ for a suitable $`vV`$ (after choosing the basis $`\{e_kk\text{I N}\}`$ appropriately). There exist suitable $`G_n`$, $`n\text{I N}\{0\}`$, in the subspace of $`L^2(Q,d\mu )`$, which includes the constant functions and the set $`\{:x_{k_1}\mathrm{}x_{k_m}:k_1,\mathrm{},k_m2\}`$, such that
$$F=\underset{n=0}{\overset{}{}}:x_1^n:G_n.$$
Set $`G=_{n=0}^{\mathrm{}}\frac{1}{n+1}:x_1^{n+1}:G_n`$. The series converges with respect to the $`L^2`$-norm, i.e. $`GL^2(Q,d\mu )`$, since
$`G_2^2`$ $`=\underset{𝑛}{{\displaystyle }}{\displaystyle \frac{1}{(n+1)^2}}:x_1^{n+1}:_2^2G_n_2^2`$
$`\stackrel{(\mathrm{3.1.1})}{=}\underset{𝑛}{{\displaystyle }}{\displaystyle \frac{1}{(n+1)^2}}2^{n1}(n+1)!G_n_2^2`$
$`\underset{𝑛}{{\displaystyle }}\mathrm{\hspace{0.17em}2}^nn!G_n_2^2`$
$`=\underset{𝑛}{{\displaystyle }}:x_1^n:_2^2G_n_2^2`$
$`=F_2^2<\mathrm{},`$
and $`G,G_n`$, $`n\text{I N}`$, are ($`\mu `$-a.e.) real-valued, by Lemma 3.1.1.
Note that $`:x_1^k:/x_1=k:x_1^{k1}:`$ is a polynomial of degree $`k1`$ with leading term $`kx_1^{k1}`$. It is evident that $`𝒢L^{\mathrm{}}(Q,d\mu )`$, but one has, furthermore, $`e^{iG}𝒢D(\mathrm{\Phi }(f))`$ and $`\mathrm{\Phi }(f)e^{iG}\phi =e^{iG}(\mathrm{\Phi }(f)F\text{1I})\phi `$, for all $`\phi 𝒢`$, since
$$e^{i_{n=0}^m\frac{1}{n+1}:x_1^{n+1}:G_n}e^{iG}$$
strongly as $`m\mathrm{}`$ (see the proof of Lemma 4.3.1 in ) and
$`\mathrm{\Phi }(f)e^{i_{n=0}^m\frac{1}{n+1}:x_1^{n+1}:G_n}\phi `$ $`=({\displaystyle \frac{}{ix_1}}+ix_1+x(v))e^{i_{n=0}^m\frac{1}{n+1}:x_1^{n+1}:G_n}\phi `$
$`=e^{i_{n=0}^m\frac{1}{n+1}:x_1^{n+1}:G_n}(\mathrm{\Phi }(f)\underset{n=0}{\overset{𝑚}{{\displaystyle }}}:x_1^n:G_n)\phi `$
$`e^{iG}(\mathrm{\Phi }(f)F\text{1I})\phi `$
with respect to the $`L^2`$-norm. Therefore,
$`e^{iG}(\mathrm{\Phi }(f)+F\text{1I})e^{iG}\phi `$ $`=e^{iG}\mathrm{\Phi }(f)e^{iG}\phi +F\phi `$
$`=(\mathrm{\Phi }(f)F\text{1I})\phi +F\phi `$
$`=\mathrm{\Phi }(f)\phi ,`$
for any $`\phi 𝒢`$. Since $`𝒢`$ is a core for the self-adjoint operator $`\mathrm{\Phi }(f)`$, the symmetric operator $`\mathrm{\Phi }(f)+F\text{1I}`$ is essentially self-adjoint on $`e^{iG}𝒢`$ and (3.1.2) holds.
To see the truth of the final assertion of the proposition, note that if $`G_0L^2(Q,d\mu )`$ satisfies $`\frac{}{x_1}G_0=F`$, then $`GG_0`$ is an element of the subspace generated by $`\{:x_{k_1}\mathrm{}x_{k_m}:k_1,\mathrm{},k_m2\}L^2(Q,d\mu )`$, i.e. $`\mathrm{\Phi }(f)`$ commutes with $`e^{i(GG_0)}`$. $`\overline{)}`$
We therefore see that for any real-linear, densely defined $`\mathrm{\Lambda }:HL^2(Q,d\mu )`$ we obtain self-adjoint transforms of the field operators $`\mathrm{\Phi }(f)`$:
$$\mathrm{\Phi }_\mathrm{\Lambda }(f)\overline{\mathrm{\Phi }(f)+\mathrm{\Lambda }f\text{1I}},fD(\mathrm{\Lambda }),$$
where the closure is understood to be taken on $`D(\mathrm{\Phi }(f))D(\mathrm{\Lambda }f\text{1I})`$. It is, of course, not true in general that the operators $`\{\mathrm{\Phi }_\mathrm{\Lambda }(f)fD(\mathrm{\Lambda })\}`$ form a representation of the CCR. We shall concentrate upon those which do.
###### Definition 3.1.4
Let $``$ be the set of all real-linear, densely defined maps from $`H`$ to $`L^2(Q,d\mu )`$, and let $`_{CCR}`$ be the subset of $``$ consisting of elements $`\mathrm{\Lambda }`$ such that
$$\pi _\mathrm{\Lambda }(W(f))e^{i\mathrm{\Phi }_\mathrm{\Lambda }(f)},fD(\mathrm{\Lambda }),$$
defines a regular representation $`(\pi _\mathrm{\Lambda },L^2(Q,d\mu ))`$ of the CCR-algebra $`𝒜(D(\mathrm{\Lambda }))`$.
Note that, by Corollary 4.1.2 in , this definition generalizes the one made in .
Hence, each $`\mathrm{\Lambda }_{CCR}`$ induces a canonical transformation on the quantum fields $`\mathrm{\Phi }(f)\mathrm{\Phi }_\mathrm{\Lambda }(f)`$, which itself is exponentiable to yield a regular representation of the algebra $`𝒜(D(\mathrm{\Lambda }))`$. We now wish to characterize the members of this set $`_{CCR}`$ and shall do so in more than one way. Let $`P_n:L^2(Q,d\mu )L^2(Q,d\mu )`$, $`n\text{I N}`$, be the orthogonal projection onto the subspace of $`n`$-particle vectors and $`P:HH`$ be the orthogonal projection onto the subspace $`JV`$.
###### Theorem 3.1.5
Let $`\mathrm{\Lambda }`$. Then $`\mathrm{\Lambda }_{CCR}`$ if and only if $`\mathrm{\Lambda }h`$ is a ($`\mu `$ a.e.) real-valued function, for each $`hD(\mathrm{\Lambda })`$, and
$$\mathrm{\Lambda }f,a^{}(Pg)\psi =\mathrm{\Lambda }g,a^{}(Pf)\psi ,$$
$`\mathrm{3.1.3}`$
for arbitrary $`f,gD(\mathrm{\Lambda })`$ and $`\psi D(a^{}(Pf))D(a^{}(Pg))`$. The assertion still holds if (3.1.3) is replaced by
$$a(Pf)P_n\mathrm{\Lambda }g=a(Pg)P_n\mathrm{\Lambda }f,$$
$`\mathrm{3.1.4}`$
for all $`f,gD(\mathrm{\Lambda })`$ and $`n\text{I N}`$.
###### Demonstration Proof
Assume that $`\mathrm{\Lambda }_{CCR}`$. Then because the field operator $`\mathrm{\Phi }_\mathrm{\Lambda }(h)`$ must be self-adjoint, the function $`\mathrm{\Lambda }h`$ must be real-valued ($`\mu `$ a.e.). Equations (3.1.3) and (3.1.4) are trivial if $`f,gV`$, so that $`Pf=0=Pg`$. Hence, one may assume that $`f=Je_1+v_1`$ and $`g=c_1Je_1+c_2Je_2+v_2`$ for suitable $`c_1,c_2\text{I R}`$ and $`v_1,v_2V`$, after choosing the basis $`\{e_k,Je_kk\text{I N}\}`$ appropriately. Differentiating the equation
$$e^{it\mathrm{\Phi }_\mathrm{\Lambda }(f)}\phi ,e^{is\mathrm{\Phi }_\mathrm{\Lambda }(g)}\psi =e^{its\sigma (f,g)}e^{is\mathrm{\Phi }_\mathrm{\Lambda }(g)}\phi ,e^{it\mathrm{\Phi }_\mathrm{\Lambda }(f)}\psi $$
with respect to $`s`$ and $`t`$ and evaluating at $`t=0=s`$, one can conclude, using Theorem VIII.7 in that
$$\mathrm{\Phi }_\mathrm{\Lambda }(f)\phi ,\mathrm{\Phi }_\mathrm{\Lambda }(g)\psi =\mathrm{\Phi }_\mathrm{\Lambda }(g)\phi ,\mathrm{\Phi }_\mathrm{\Lambda }(f)\psi +i\sigma (f,g)\phi ,\psi ,$$
for any $`\phi ,\psi 𝒢D(\mathrm{\Phi }_\mathrm{\Lambda }(f))D(\mathrm{\Phi }_\mathrm{\Lambda }(g))`$. Since $`\mathrm{\Lambda }f`$ and $`\mathrm{\Lambda }g`$ are $`\mu `$ a.e. real-valued, one has
$`0`$ $`=\mathrm{\Phi }_\mathrm{\Lambda }(f)\mathrm{\Omega },\mathrm{\Phi }_\mathrm{\Lambda }(g)\psi \mathrm{\Phi }_\mathrm{\Lambda }(g)\mathrm{\Omega },\mathrm{\Phi }_\mathrm{\Lambda }(f)\psi i\sigma (f,g)\mathrm{\Omega },\psi `$
$`=\mathrm{\Phi }(f)\mathrm{\Omega },\mathrm{\Phi }(g)\psi +\mathrm{\Lambda }f\mathrm{\Omega },\mathrm{\Phi }(g)\psi +\mathrm{\Phi }(f)\mathrm{\Omega },\mathrm{\Lambda }g\psi +\mathrm{\Lambda }f\mathrm{\Omega },\mathrm{\Lambda }g\psi `$
$`\mathrm{\Phi }(g)\mathrm{\Omega },\mathrm{\Phi }(f)\psi \mathrm{\Lambda }g\mathrm{\Omega },\mathrm{\Phi }(f)\psi \mathrm{\Phi }(g)\mathrm{\Omega },\mathrm{\Lambda }f\psi \mathrm{\Lambda }g\mathrm{\Omega },\mathrm{\Lambda }f\psi `$
$`i\sigma (f,g)\mathrm{\Omega },\psi `$
$`=\mathrm{\Phi }(f)\mathrm{\Omega },\mathrm{\Phi }(g)\psi \mathrm{\Phi }(g)\mathrm{\Omega },\mathrm{\Phi }(f)\psi i\sigma (f,g)\mathrm{\Omega },\psi `$
$`+\mathrm{\Lambda }f\mathrm{\Omega },\mathrm{\Phi }(g)\psi +\mathrm{\Phi }(f)\mathrm{\Omega },\mathrm{\Lambda }g\psi \mathrm{\Phi }(g)\mathrm{\Omega },\mathrm{\Lambda }f\psi \mathrm{\Lambda }g\mathrm{\Omega },\mathrm{\Phi }(f)\psi `$
$`+\mathrm{\Lambda }f\mathrm{\Omega },\mathrm{\Lambda }g\psi \mathrm{\Lambda }g\mathrm{\Omega },\mathrm{\Lambda }f\psi `$
$`=\mathrm{\Lambda }f\mathrm{\Omega },\mathrm{\Phi }(g)\psi +\mathrm{\Phi }(f)\mathrm{\Omega },\mathrm{\Lambda }g\psi \mathrm{\Phi }(g)\mathrm{\Omega },\mathrm{\Lambda }f\psi \mathrm{\Lambda }g\mathrm{\Omega },\mathrm{\Phi }(f)\psi `$
$`=\mathrm{\Lambda }f,(\mathrm{\Phi }(g)\overline{\mathrm{\Phi }(g)\mathrm{\Omega }}\text{1I})\psi \mathrm{\Lambda }g,(\mathrm{\Phi }(f)\overline{\mathrm{\Phi }(f)\mathrm{\Omega }}\text{1I})\psi .`$
By using
$`(\mathrm{\Phi }(f)\overline{\mathrm{\Phi }(f)\mathrm{\Omega }}\text{1I})\psi `$ $`=(\mathrm{\Phi }(Je_1)\overline{\mathrm{\Phi }(Je_1)\mathrm{\Omega }}\text{1I})\psi =(\mathrm{\Phi }(Je_1)+ix_1)\psi `$
$`=(\mathrm{\Phi }(Je_1)+i\mathrm{\Phi }(e_1))\psi =\sqrt{2}a^{}(Je_1)\psi `$
$`=\sqrt{2}a^{}(Pf)\psi `$
and the similar equality $`(\mathrm{\Phi }(g)\overline{\mathrm{\Phi }(g)\mathrm{\Omega }}\text{1I})\psi =\sqrt{2}a^{}(Pg)\psi `$, one can conclude that (3.1.3) is fulfilled for $`\psi 𝒢`$.
The next step is to prove (3.1.3) for polynomials $`\psi L^2(Q,d\mu )`$. Since $`𝒢`$ is a core for $`N`$ (Lemma 3.1.2), for arbitrary $`k_1,\mathrm{},k_m\text{I N}`$ there exists a sequence $`\{\psi _n\}𝒢`$ which converges in $`L^2(Q,d\mu )`$ to $`:x_{k_1}\mathrm{}x_{k_m}:`$ such that also the sequence $`\{N\psi _n\}`$ converges in $`L^2(Q,d\mu )`$ to $`N:x_{k_1}\mathrm{}x_{k_m}:`$. For arbitrary $`hH`$ one has
$`a^{}(h)`$ $`(\psi _n:x_{k_1}\mathrm{}x_{k_m}:)^2`$
$`h^2(\psi _n:x_{k_1}\mathrm{}x_{k_m}:),(N+1)(\psi _n:x_{k_1}\mathrm{}x_{k_m}:)0,`$
as $`n\mathrm{}`$, so that (3.1.3) holds for $`\psi =:x_{k_1}\mathrm{}x_{k_m}:`$. It follows that
$`\mathrm{\Lambda }f,a^{}(Pg)\psi `$ $`=\underset{n=0}{\overset{}{{\displaystyle }}}P_n\mathrm{\Lambda }f,P_na^{}(Pg)\psi =\underset{n=1}{\overset{}{{\displaystyle }}}P_n\mathrm{\Lambda }f,a^{}(Pg)P_{n1}\psi `$ $`\mathrm{3.1.5}`$
$`=\underset{n=1}{\overset{}{{\displaystyle }}}\mathrm{\Lambda }f,a^{}(Pg)P_{n1}\psi =\underset{n=1}{\overset{}{{\displaystyle }}}\mathrm{\Lambda }g,a^{}(Pf)P_{n1}\psi `$
$`=\mathrm{\Lambda }g,a^{}(Pf)\psi ,`$
for $`\psi D(a^{}(Pf))D(a^{}(Pg))`$. The chain of equalities (3.1.5) establishes the asserted equivalence of (3.1.3) and (3.1.4).
Assume now that (3.1.3) holds for arbitrary $`f,gD(\mathrm{\Lambda })`$ and that $`\mathrm{\Lambda }h`$ is $`\mu `$ a.e. real-valued, for each $`hD(\mathrm{\Lambda })`$. The assertion in this direction will be established in part by appealing to another characterization of $`\mathrm{\Lambda }_{CCR}`$ appearing in Theorem 3.1.6, which will be presented subsequently. In particular, here it will be shown that there exists an extension $`\mathrm{\Lambda }^{}`$ of $`\mathrm{\Lambda }`$ containing $`V`$ in its domain of definition and satisfying the hypothesis of Theorem 3.1.6. It will then follow from Theorem 3.1.6 that both $`\mathrm{\Lambda }^{}`$ and $`\mathrm{\Lambda }`$ are contained in $`_{CCR}`$.
Let $`gVD(\mathrm{\Lambda })`$ be arbitrary but fixed. The set $`\{PffD(\mathrm{\Lambda })\}`$ is dense in $`JV`$, so the set $`\{a^{}(Pf)\psi \psi D(a^{}(Pf))D(a^{}(Pg)),fD(\mathrm{\Lambda })\}`$ is dense in the orthogonal complement of the set $`\{\mathrm{\Omega }\}`$. Therefore, $`\mathrm{\Lambda }g`$ must be a multiple of $`\mathrm{\Omega }`$, i.e. a constant function, whenever $`Pg=0`$, i.e. whenever $`gV`$ (use the hypothesis (3.1.3)). Thus, the restriction of $`\mathrm{\Lambda }`$ to $`VD(\mathrm{\Lambda })`$ determines a linear form $`\mathrm{}:VD(\mathrm{\Lambda })\text{I R}`$. But $`\mathrm{}`$ has a linear extension $`\mathrm{}^{}`$ to $`V`$, so $`\mathrm{\Lambda }`$ has a linear extension $`\mathrm{\Lambda }^{}`$ to $`D(\mathrm{\Lambda })+V`$ such that $`\mathrm{\Lambda }^{}V=\mathrm{}^{}`$ and such that (3.1.3) holds with $`\mathrm{\Lambda }`$ replaced by $`\mathrm{\Lambda }^{}`$ (both sides of (3.1.3) under this replacement are equal to zero for the additional vectors $`gV`$ and constant $`\mathrm{\Lambda }g`$). It shall be established that this extension $`\mathrm{\Lambda }^{}`$ fulfills the hypothesis of Theorem 3.1.6.
For arbitrary $`f_1,\mathrm{},f_mJD(\mathrm{\Lambda }^{})V`$, one has
$$a^{}(f_2):x(f_3)\mathrm{}x(f_m):=\sqrt{2}:x(f_2)x(f_3)\mathrm{}x(f_m):,$$
so that
$`\mathrm{\Lambda }^{}Jf_1,:x(f_2)\mathrm{}x(f_m):`$ $`={\displaystyle \frac{1}{\sqrt{2}}}\mathrm{\Lambda }^{}Jf_1,a^{}(f_2):x(f_3)\mathrm{}x(f_m):`$
$`={\displaystyle \frac{i}{\sqrt{2}}}\mathrm{\Lambda }^{}Jf_1,a^{}(Jf_2):x(f_3)\mathrm{}x(f_m):`$
$`\stackrel{(\mathrm{3.1.3})}{=}{\displaystyle \frac{i}{\sqrt{2}}}\mathrm{\Lambda }^{}Jf_2,a^{}(Jf_1):x(f_3)\mathrm{}x(f_m):`$
$`=\mathrm{\Lambda }^{}Jf_2,:x(f_1)x(f_3)\mathrm{}x(f_m):,`$
and the conditions (i)-(iii) of Theorem 3.1.6 are satisfied for $`\mathrm{\Lambda }^{}`$. $`\overline{)}`$
We now give the announced second characterization of $`\mathrm{\Lambda }_{CCR}`$.
###### Theorem 3.1.6
If $`\mathrm{\Lambda }_{CCR}`$, then there exists an extension $`\mathrm{\Lambda }^{}_{CCR}`$ of $`\mathrm{\Lambda }`$ with $`VD(\mathrm{\Lambda }^{})`$. Moreover, if $`VD(\mathrm{\Lambda })`$, then $`\mathrm{\Lambda }_{CCR}`$ is equivalent to the following three conditions:
(i) The functions $`\mathrm{\Lambda }h`$ are $`\mu `$ a.e. real-valued, for all $`hD(\mathrm{\Lambda })`$.
(ii) For arbitrary $`f_1,\mathrm{},f_mJD(\mathrm{\Lambda })V`$,
$$\mathrm{\Lambda }Jf_1,:x(f_2)x(f_3)\mathrm{}x(f_m):=\mathrm{\Lambda }Jf_2,:x(f_1)x(f_3)\mathrm{}x(f_m):.$$
$`\mathrm{3.1.6}`$
(iii) $`\mathrm{\Lambda }f`$ is a (real) constant function for all $`fV`$.
###### Demonstration Proof
The proof of Theorem 3.1.5 shows that $`\mathrm{\Lambda }_{CCR}`$ implies that there exists an extension $`\mathrm{\Lambda }^{}`$ of $`\mathrm{\Lambda }`$ which fulfills conditions (i)-(iii) and $`VD(\mathrm{\Lambda }^{})`$. It therefore remains only to prove that $`VD(\mathrm{\Lambda })`$ and (i)-(iii) imply $`\mathrm{\Lambda }_{CCR}`$. By applying a suitable coherent transformation, it may be assumed that $`\mathrm{\Lambda }h,\mathrm{\Omega }=0`$, for all $`hD(\mathrm{\Lambda })`$. By choosing a suitable basis $`\{e_kk\text{I N}\}`$ of $`V`$, it may also be assumed that $`\{Je_kk\text{I N}\}D(\mathrm{\Lambda })`$.
Now let $`f,gD(\mathrm{\Lambda })`$ be arbitrary but fixed. One can choose the basis of $`V`$ such that $`f=c_1Je_1+v_1`$, $`g=c_2Je_1+c_3Je_2+v_2`$, with $`c_1,c_2,c_3\text{I R}`$ and $`v_1,v_2V`$. There are suitable $`\lambda _{k_1\mathrm{}k_m}\text{I R}`$ ($`m,k_1,\mathrm{},k_m\text{I N}`$) such that $`\lambda _{k_1\mathrm{}k_m}`$ is symmetric in $`k_2,\mathrm{},k_m`$ and such that
$$\mathrm{\Lambda }Je_{k_1}=\underset{𝑚}{}\underset{k_2,\mathrm{},k_m}{}\lambda _{k_1\mathrm{}k_m}:x_{k_2}\mathrm{}x_{k_m}:,$$
for arbitrary $`k_1\text{I N}`$. The $`\lambda _{k_1\mathrm{}k_m}`$ are symmetric in the indices according to Lemma 3.1.1 and (3.1.6):
$`\lambda _{k_1k_2k_3\mathrm{}k_m}`$ $`={\displaystyle \frac{2^m}{m!}}\mathrm{\Lambda }Je_{k_1},:x_{k_2}x_{k_3}\mathrm{}x_{k_m}:`$
$`={\displaystyle \frac{2^m}{m!}}\mathrm{\Lambda }Je_{k_2},:x_{k_1}x_{k_3}\mathrm{}x_{k_m}:`$
$`=\lambda _{k_2k_1k_3\mathrm{}k_m}.`$
The function defined by
$$G=\underset{𝑚}{}\underset{\{1,2\}\{k_1,\mathrm{},k_m\}\mathrm{}}{}\frac{\lambda _{k_1\mathrm{}k_m}}{m}:x_{k_1}\mathrm{}x_{k_m}:,$$
is an $`L^2`$-function, since by (3.1.1) one has
$`G_2^2`$ $`=\underset{𝑚}{{\displaystyle }}\underset{\{1,2\}\{k_1,\mathrm{},k_m\}\mathrm{}}{{\displaystyle }}{\displaystyle \frac{\lambda _{k_1\mathrm{}k_m}^2}{m^2}}{\displaystyle \frac{m!}{2^m}}`$
$`\underset{𝑚}{{\displaystyle }}\underset{1\{k_1,\mathrm{},k_m\}}{{\displaystyle }}{\displaystyle \frac{\lambda _{k_1\mathrm{}k_m}^2}{m^2}}{\displaystyle \frac{m!}{2^m}}+\underset{𝑚}{{\displaystyle }}\underset{2\{k_1,\mathrm{},k_m\}}{{\displaystyle }}{\displaystyle \frac{\lambda _{k_1\mathrm{}k_m}^2}{m^2}}{\displaystyle \frac{m!}{2^m}}`$
$`\underset{𝑚}{{\displaystyle }}\underset{l_2,\mathrm{},l_m}{{\displaystyle }}m{\displaystyle \frac{\lambda _{1l_2\mathrm{}l_m}^2}{m^2}}{\displaystyle \frac{m!}{2^m}}+\underset{𝑚}{{\displaystyle }}\underset{l_2,\mathrm{},l_m}{{\displaystyle }}m{\displaystyle \frac{\lambda _{2l_2\mathrm{}l_m}^2}{m^2}}{\displaystyle \frac{m!}{2^m}}`$
$`=\underset{𝑚}{{\displaystyle }}\underset{l_2,\mathrm{},l_m}{{\displaystyle }}\lambda _{1l_2\mathrm{}l_m}^2{\displaystyle \frac{(m1)!}{2^m}}+\underset{𝑚}{{\displaystyle }}\underset{l_2,\mathrm{},l_m}{{\displaystyle }}\lambda _{2l_2\mathrm{}l_m}^2{\displaystyle \frac{(m1)!}{2^m}}`$
$`={\displaystyle \frac{1}{2}}\underset{𝑚}{{\displaystyle }}\underset{l_2,\mathrm{},l_m}{{\displaystyle }}\lambda _{1l_2\mathrm{}l_m}:x_{l_2}\mathrm{}x_{l_m}:_2^2`$
$`+{\displaystyle \frac{1}{2}}\underset{𝑚}{{\displaystyle }}\underset{l_2,\mathrm{},l_m}{{\displaystyle }}\lambda _{2l_2\mathrm{}l_m}:x_{l_2}\mathrm{}x_{l_m}:_2^2`$
$`={\displaystyle \frac{1}{2}}\mathrm{\Lambda }Je_1_2^2+{\displaystyle \frac{1}{2}}\mathrm{\Lambda }Je_2_2^2<\mathrm{}.`$
Then
$`{\displaystyle \frac{}{x_1}}G`$ $`=\underset{m_0\mathrm{}}{lim}\underset{m=1}{\overset{m_0}{{\displaystyle }}}\underset{j=1}{\overset{𝑚}{{\displaystyle }}}\underset{k_i,ij}{{\displaystyle }}{\displaystyle \frac{\lambda _{k_1\mathrm{}k_{j1}1k_{j+1}\mathrm{}k_m}}{m}}:x_{k_1}\mathrm{}x_{k_{j1}}x_{k_{j+1}}\mathrm{}x_{k_m}:`$
$`=\underset{𝑚}{{\displaystyle }}\underset{l_2,\mathrm{},l_m}{{\displaystyle }}\lambda _{1l_2\mathrm{}l_m}:x_{l_2}\mathrm{}x_{l_m}:`$
$`=\mathrm{\Lambda }Je_1.`$
But $`\mathrm{\Lambda }v_1`$ is a constant function and $`\mathrm{\Lambda }v_1,\mathrm{\Omega }=0`$, so $`\mathrm{\Lambda }v_1=0`$ and $`c_1\frac{}{x_1}G=\mathrm{\Lambda }f`$. Proposition 3.1.3 then entails
$$\mathrm{\Phi }_\mathrm{\Lambda }(f)=e^{iG}\mathrm{\Phi }(f)e^{iG}.$$
Similarly, one has
$$\frac{}{x_2}G=\mathrm{\Lambda }Je_2,$$
and, since
$`{\displaystyle \frac{}{x(c_2e_1+c_3e_2)}}G`$ $`=c_2{\displaystyle \frac{}{x_1}}G+c_3{\displaystyle \frac{}{x_2}}G`$
$`=c_2\mathrm{\Lambda }Je_1+c_3\mathrm{\Lambda }Je_2`$
$`=\mathrm{\Lambda }g`$
(recall that $`\frac{}{x(c_2e_1+c_3e_2)}=\sqrt{2}(a(c_2e_1+c_3e_2))=\sqrt{2}(c_2a(e_1)+c_3a(e_2))=c_2\frac{}{x_1}+c_3\frac{}{x_2}`$ and use $`\mathrm{\Lambda }v_2=0`$, as well), Proposition 3.1.3 also implies
$$\mathrm{\Phi }_\mathrm{\Lambda }(g)=e^{iG}\mathrm{\Phi }(g)e^{iG}.$$
Hence, one has $`\mathrm{\Phi }_\mathrm{\Lambda }(f+g)=e^{iG}\mathrm{\Phi }(f+g)e^{iG}`$, which implies
$`e^{i\mathrm{\Phi }_\mathrm{\Lambda }(f)}e^{i\mathrm{\Phi }_\mathrm{\Lambda }(g)}`$ $`=e^{iG}e^{i\mathrm{\Phi }(f)}e^{i\mathrm{\Phi }(g)}e^{iG}`$
$`=e^{iG}e^{\frac{i}{2}\sigma (f,g)}e^{i\mathrm{\Phi }(f+g)}e^{iG}`$
$`=e^{\frac{i}{2}\sigma (f,g)}e^{i\mathrm{\Phi }_\mathrm{\Lambda }(f+g)},`$
and the proof is complete. $`\overline{)}`$
It follows from Theorem 3.1.5 that each $`\mathrm{\Lambda }_{CCR}`$ is of the form $`\mathrm{\Lambda }=\mathrm{\Lambda }_l+\mathrm{\Lambda }_q`$, where $`\mathrm{\Lambda }_l_{CCR}`$ is linear, that is to say, the associated field operators $`\mathrm{\Phi }_{\mathrm{\Lambda }_l}(f)`$, $`fD(\mathrm{\Lambda }_l)`$, are of the form $`\mathrm{\Phi }(g)+c\text{1I}`$, for suitable $`gH`$ and $`c\text{I R}`$. In other words, the transformation $`\mathrm{\Phi }(f)\mathrm{\Phi }_{\mathrm{\Lambda }_l}(f)`$ is one of the inhomogeneous linear canonical transformations alluded to in Chapter II. It then follows that, with a suitable choice of linear $`\mathrm{\Lambda }_l`$, the operator $`\mathrm{\Lambda }_q_{CCR}`$ satisfies $`P_0\mathrm{\Lambda }_qf=0=P_1\mathrm{\Lambda }_qf`$, for any $`fD(\mathrm{\Lambda }_q)`$. The set of such operators $`\mathrm{\Lambda }_q`$ will be denoted by $`_{CCR}^q`$. The superscript $`q`$ is chosen because the degree of such transformations is quadratic or higher. The structure of the linear elements of $`_{CCR}`$ will be discussed in Chapter IV. Here we shall consider the elements of $`_{CCR}^q`$. The following result generalizes Proposition 3.3.4 in from the quadratic case to this general setting.
###### Proposition 3.1.7
Let $`\mathrm{\Lambda }_{CCR}^q`$. Then there exists a unique maximal extension $`\mathrm{\Lambda }_{max}_{CCR}^q`$ of $`\mathrm{\Lambda }`$.
###### Demonstration Proof
By Theorem 3.1.6, it may be assumed that $`VD(\mathrm{\Lambda })`$, and by suitably choosing the basis $`\{e_kk\text{I N}\}`$ of $`V`$, it may also be assumed that $`\{e_k,Je_kk\text{I N}\}D(\mathrm{\Lambda })`$. If $`f=f_1+Jf_2D(\mathrm{\Lambda })`$ with $`f_1,f_2V`$, then Lemma 3.1.1 and Theorem 3.1.6 entail
$`\mathrm{\Lambda }f`$ $`=\underset{m,k_1,\mathrm{},k_m}{{\displaystyle }}{\displaystyle \frac{2^m}{m!}}\mathrm{\Lambda }Jf_2,:x_{k_1}x_{k_2}\mathrm{}x_{k_m}::x_{k_1}\mathrm{}x_{k_m}:`$
$`=\underset{m,k_1,\mathrm{},k_m}{{\displaystyle }}{\displaystyle \frac{2^m}{m!}}\mathrm{\Lambda }Je_{k_1},:x(f_2)x_{k_2}\mathrm{}x_{k_m}::x_{k_1}\mathrm{}x_{k_m}:.`$
If $`\mathrm{\Lambda }^{}_{CCR}^q`$ is an extension of $`\mathrm{\Lambda }`$ and $`f=f_1+Jf_2D(\mathrm{\Lambda }^{})`$, with $`f_1,f_2V`$, then
$`\mathrm{}`$ $`>\mathrm{\Lambda }^{}f^2`$
$`=\underset{m,k_1,\mathrm{},k_m}{{\displaystyle }}{\displaystyle \frac{2^m}{m!}}\mathrm{\Lambda }Je_{k_1},:x(f_2)x_{k_2}\mathrm{}x_{k_m}::x_{k_1}\mathrm{}x_{k_m}:^2`$
$`\stackrel{(\mathrm{3.1.1})}{=}\underset{m,k_1,\mathrm{},k_m}{{\displaystyle }}{\displaystyle \frac{2^m}{m!}}|\mathrm{\Lambda }Je_{k_1},:x(f_2)x_{k_2}\mathrm{}x_{k_m}:|^2.`$
Hence, define $`\mathrm{\Lambda }_{max}`$ by
$$D(\mathrm{\Lambda }_{max})=\{f+Jgf,gV,\underset{m,k_1,\mathrm{},k_m}{}\frac{2^m}{m!}|\mathrm{\Lambda }Je_{k_1},:x(g)x_{k_2}\mathrm{}x_{k_m}:|^2<\mathrm{}\}$$
and
$$\mathrm{\Lambda }_{max}f=\underset{m,k_1,\mathrm{},k_m}{}\frac{2^m}{m!}\mathrm{\Lambda }Je_{k_1},:x(f_2)x_{k_2}\mathrm{}x_{k_m}::x_{k_1}\mathrm{}x_{k_m}:,$$
for $`f=f_1+Jf_2D(\mathrm{\Lambda }_{max})`$, $`f_1,f_2V`$. Since all other assertions are now clear, it remains only to show that $`\mathrm{\Lambda }_{max}_{CCR}`$. But for $`f_1,\mathrm{},f_mJD(\mathrm{\Lambda }_{max})V`$ (with $`f_l=_kc_{lk}e_k`$), one sees
$`\mathrm{\Lambda }_{max}Jf_1,:x(f_2)\mathrm{}x(f_m):`$
$`=\underset{k_2,\mathrm{},k_m}{{\displaystyle }}c_{2k_2}\mathrm{}c_{mk_m}\mathrm{\Lambda }_{max}Jf_1,:x(e_{k_2})\mathrm{}x(e_{k_m}):`$
$`=\underset{k_2,\mathrm{},k_m}{{\displaystyle }}c_{2k_2}\mathrm{}c_{mk_m}\mathrm{\Lambda }Je_{k_m},:x(f_1)x(e_{k_2})\mathrm{}x(e_{k_{m1}}):`$
$`=\underset{k_m}{{\displaystyle }}c_{mk_m}\mathrm{\Lambda }Je_{k_m},:x(f_1)x(f_2)\mathrm{}x(f_{m1}):`$
$`=\underset{k_m}{{\displaystyle }}c_{mk_m}\mathrm{\Lambda }Je_{k_m},:x(f_2)x(f_1)\mathrm{}x(f_{m1}):`$
$`=\mathrm{\Lambda }_{max}Jf_2,:x(f_1)x(f_3)\mathrm{}x(f_m):.`$
Theorem 3.1.6 then implies $`\mathrm{\Lambda }_{max}_{CCR}`$. $`\overline{)}`$
We then can use this result to show, as in Section 3.3 of for the quadratic case, that any $`\mathrm{\Lambda }_{CCR}^q`$ has a particular form, which leads to a convenient “standard” form for the corresponding field operators $`\mathrm{\Phi }_\mathrm{\Lambda }(f)`$.
###### Proposition 3.1.8
For each $`\mathrm{\Lambda }_{CCR}^q`$ there exist an orthonormal basis $`\{e_kk\text{I N}\}`$ of $`V`$ and real numbers $`\lambda _{kk_1\mathrm{}k_m}`$ totally symmetric in the indices such that all of the following conditions are satisfied:
$$\underset{m,k_1,\mathrm{},k_m}{}\lambda _{kk_1\mathrm{}k_m}^2\frac{m!}{2^m}<\mathrm{},$$
for any $`k\text{I N}`$; moreover, if $`\mathrm{\Lambda }^{}_{CCR}^q`$ is defined on the linear span of $`\{e_k,Je_kk\text{I N}\}`$ by $`\mathrm{\Lambda }^{}e_k=0`$ and $`\mathrm{\Lambda }^{}Je_k=_{m,k_1,\mathrm{},k_m}\lambda _{kk_1\mathrm{}k_m}:x_{k_1}\mathrm{}x_{k_m}:`$, then $`\mathrm{\Lambda }\mathrm{\Lambda }_{max}^{}`$.
###### Demonstration Proof
Once again, it may be assumed that $`VD(\mathrm{\Lambda })`$ and $`\{Je_kk\text{I N}\}D(\mathrm{\Lambda })`$. Set
$$\lambda _{kk_1\mathrm{}k_m}=\frac{2^m}{m!}\mathrm{\Lambda }Je_k,:x_{k_1}\mathrm{}x_{k_m}:.$$
Then, by Lemma 3.1.1, one finds
$`\mathrm{}`$ $`>\mathrm{\Lambda }Je_k^2`$
$`=\underset{m,k_1,\mathrm{},k_m}{{\displaystyle }}\lambda _{kk_1\mathrm{}k_m}:x_{k_1}\mathrm{}x_{k_m}:^2`$
$`=\underset{m,k_1,\mathrm{},k_m}{{\displaystyle }}\lambda _{kk_1\mathrm{}k_m}^2{\displaystyle \frac{m!}{2^m}}.`$
Proposition 3.1.7 then completes the proof. $`\overline{)}`$
Since $`\mathrm{\Lambda }_{max}`$ uniquely exists, one may consider $`\mathrm{\Lambda }_{CCR}^q`$ as being defined on a symplectic orthonormal basis $`\{e_k,Je_k\}_{k\text{I N}}`$ such that $`\mathrm{\Lambda }e_k=0`$ and $`\mathrm{\Lambda }Je_k=_{m,k_1,\mathrm{},k_m}\lambda _{kk_1\mathrm{}k_m}:x_{k_1}\mathrm{}x_{k_m}:`$. Then one has
$$q_k\mathrm{\Phi }(e_k)\mathrm{\Phi }_\mathrm{\Lambda }(e_k)=\mathrm{\Phi }(e_k)=q_k,$$
and
$`p_k\mathrm{\Phi }(Je_k)\mathrm{\Phi }_\mathrm{\Lambda }(Je_k)`$ $`=\overline{p_k+{\displaystyle \underset{m,k_1,\mathrm{},k_m}{}}\lambda _{kk_1\mathrm{}k_m}:x_{k_1}\mathrm{}x_{k_m}:\text{1I}}`$
$`=\overline{p_k+{\displaystyle \underset{m,k_1,\mathrm{},k_m}{}}\lambda _{kk_1\mathrm{}k_m}:q_{k_1}\mathrm{}q_{k_m}:}.`$
In short, the standard form of a canonical transformation of arbitrary degree is that given in equation (1.1). This standard form, along with being physically more transparent, was useful in the special case of quadratic transformations in to establish results concerning the unitary equivalence of such representations with the Fock representation. Though we do not prove such results here for transformations of arbitrary degree, we shall give necessary and sufficient conditions for unitary, resp. quasi-, equivalence between Fock, coherent, and quasifree representations and representations of finite degree in Chapters IV and V.
###### 3.2 Irreducibility of the Representation
In this section we shall show that for bounded $`\mathrm{\Lambda }_{CCR}`$, the corresponding representation $`\pi _\mathrm{\Lambda }`$ of the CCR is irreducible. If $`\mathrm{\Lambda }`$ is unbounded, then it can occur that $`\pi _\mathrm{\Lambda }`$ is reducible, as we shall explain. We begin with a technical lemma concerning the closability and continuity properties of $`\mathrm{\Lambda }_{CCR}`$ such that $`P_0\mathrm{\Lambda }`$ is the zero operator, i.e. such that the range of $`\mathrm{\Lambda }`$ is orthogonal to the vacuum vector $`\mathrm{\Omega }`$. Note that this is true of each $`\mathrm{\Lambda }_{CCR}`$, up to a coherent transformation, i.e. a transformation of degree zero.
###### Lemma 3.2.1
Let $`\mathrm{\Lambda }_{CCR}`$ with $`P_0\mathrm{\Lambda }0`$. Then $`\mathrm{\Lambda }`$ is closable and $`\overline{\mathrm{\Lambda }}_{CCR}`$. Furthermore, for any sequence $`\{f_n\}_{n\text{I N}}D(\mathrm{\Lambda })`$ such that $`f_nfD(\overline{\mathrm{\Lambda }})`$ and $`\mathrm{\Lambda }f_n\overline{\mathrm{\Lambda }}f`$ as $`n\mathrm{}`$, then the operators $`\{e^{i\mathrm{\Phi }_\mathrm{\Lambda }(f_n)}\}_{n\text{I N}}`$ converge strongly to $`e^{i\mathrm{\Phi }_{\overline{\mathrm{\Lambda }}}(f)}`$ as $`n\mathrm{}`$.
###### Demonstration Proof
It will first be shown that such $`\mathrm{\Lambda }`$ are closable. Once again, it may be assumed that $`VD(\mathrm{\Lambda })`$ and $`\{Je_kk\text{I N}\}D(\mathrm{\Lambda })`$. Let $`\{g_n\}`$ in $`D(\mathrm{\Lambda })`$ be a sequence such that $`g_n=h_n+Jh_n^{}`$, with $`h_n,h_n^{}V`$, $`g_n0`$ and $`\mathrm{\Lambda }g_nF`$, for some $`FL^2(Q,d\mu )`$. Then for arbitrary $`m,k_1,\mathrm{},k_m\text{I N}`$, Theorem 3.1.6 implies
$`\mathrm{\Lambda }g_n,:x_{k_1}\mathrm{}x_{k_m}:`$ $`=\mathrm{\Lambda }Jh_n^{},:x_{k_1}x_{k_2}\mathrm{}x_{k_m}:`$
$`=\mathrm{\Lambda }Je_{k_1},:x(h_n^{})x_{k_2}\mathrm{}x_{k_m}:`$
$`0,`$
as $`n\mathrm{}`$. Hence, $`F=0`$ and $`\mathrm{\Lambda }`$ is closable. Since $`\overline{\mathrm{\Lambda }}`$ fulfills (3.1.4), Theorem 3.1.5 entails $`\overline{\mathrm{\Lambda }}_{CCR}`$.
In addressing the final assertion in the lemma, one may assume that $`VD(\mathrm{\Lambda })`$, $`\mathrm{\Lambda }=\overline{\mathrm{\Lambda }}`$, $`f=Je_1+v`$, with $`vV`$, and
$$\mathrm{\Lambda }f=\underset{𝑚}{}\underset{l=0}{\overset{}{}}\underset{k_1,\mathrm{},k_m2}{}c_{lk_1\mathrm{}k_m}:x_1^lx_{k_1}\mathrm{}x_{k_m}:,$$
for suitable $`c_{lk_1\mathrm{}k_m}`$, symmetric in the indices $`k_1\mathrm{}k_m`$. Set
$$G_n=\underset{m+ln}{\underset{2k_1,\mathrm{},k_mn}{}}c_{lk_1\mathrm{}k_m}\frac{1}{l+1}:x_1^{l+1}x_{k_1}\mathrm{}x_{k_m}:$$
and
$$G=\underset{m,l,k_1,\mathrm{},k_m}{}c_{lk_1\mathrm{}k_m}\frac{1}{l+1}:x_1^{l+1}x_{k_1}\mathrm{}x_{k_m}:.$$
Note that the proof of Proposition 3.1.3 implies that $`GL^2(Q,d\mu )`$. Then one has
$`\mathrm{\Phi }_\mathrm{\Lambda }(f_n)e^{iG_m}\phi `$ $`=\mathrm{\Phi }(f_n)e^{iG_m}\phi +\mathrm{\Lambda }f_ne^{iG_m}\phi `$
$`\mathrm{\Phi }(f)e^{iG_m}\phi +\mathrm{\Lambda }fe^{iG_m}\phi `$
$`=\mathrm{\Phi }_\mathrm{\Lambda }(f)e^{iG_m}\phi ,`$
for $`\phi 𝒢`$. But the set $`\{e^{iG_m}\phi m\text{I N},\phi 𝒢\}`$ is contained in the domain of the strong graph limit of the sequence $`\{\mathrm{\Phi }_\mathrm{\Lambda }(f_n)\}`$, and this strong graph limit is a symmetric and closed operator (see, e.g., Theorem VIII.27 in ). Furthermore, this strong graph limit acts upon the elements of this set as $`\mathrm{\Phi }_\mathrm{\Lambda }(f)`$. As in the proof of Proposition 3.1.3, one may conclude that $`e^{iG_m}\phi `$ converges to $`e^{iG}\phi `$ as $`m\mathrm{}`$ and, since $`\frac{}{x_1}G_m`$ converges to $`\mathrm{\Lambda }f`$ in the $`L^2`$-norm,
$$\mathrm{\Phi }_\mathrm{\Lambda }(f)e^{iG_m}\phi =e^{iG_m}(\mathrm{\Phi }(f)+\mathrm{\Lambda }f\text{1I}\frac{}{x_1}G_m\text{1I})\phi e^{iG}\mathrm{\Phi }(f)\phi .$$
In other words, the set $`e^{iG}𝒢`$ is contained in the domain of the strong graph limit of the sequence $`\{\mathrm{\Phi }_\mathrm{\Lambda }(f_n)\}`$, and it acts upon this set as $`\mathrm{\Phi }_\mathrm{\Lambda }(f)=e^{iG}\mathrm{\Phi }(f)e^{iG}`$. But $`𝒢`$ is a core for $`\mathrm{\Phi }(f)`$, so it follows that $`\mathrm{\Phi }_\mathrm{\Lambda }(f)`$ is, in fact, the strong graph limit of $`\{\mathrm{\Phi }_\mathrm{\Lambda }(f_n)\}`$. The final assertion of the lemma then follows from Theorems VIII.21 and VIII.26 in . $`\overline{)}`$
We now provide a sufficient condition on $`\mathrm{\Lambda }`$ which entails that $`\pi _\mathrm{\Lambda }`$ is irreducible.
###### Theorem 3.2.2
Let $`\mathrm{\Lambda }_{CCR}`$. If a dense subset of $`V`$ is contained in $`D(\overline{\mathrm{\Lambda }P_0\mathrm{\Lambda }})`$, then $`\pi _\mathrm{\Lambda }`$ is irreducible. In particular, $`\pi _\mathrm{\Lambda }`$ is irreducible for bounded $`\mathrm{\Lambda }P_0\mathrm{\Lambda }`$, resp. for bounded $`\mathrm{\Lambda }`$.
###### Demonstration Proof
Let $`T:L^2(Q,d\mu )L^2(Q,d\mu )`$ be bounded and assume
$$[T,e^{i\mathrm{\Phi }_\mathrm{\Lambda }(g)}]=0,$$
for any $`gD(\mathrm{\Lambda })`$. Then
$$[T,e^{i\mathrm{\Phi }(g)}]=0,$$
$`\mathrm{3.2.1}`$
for $`g`$ in a dense subset of $`V`$. But Lemma 3.2.1 implies that equation (3.2.1) holds for all $`gV`$. Since $`\{e^{ix(f)}fV\}`$ generates the maximally abelian von Neumann algebra $`L^{\mathrm{}}(Q,d\mu )`$, $`T`$ may be identified with an element of $`L^{\mathrm{}}(Q,d\mu )`$ (or, more accurately, with the corresponding multiplication operator).
Now, Proposition 3.1.3 implies that for any $`gD(\mathrm{\Lambda })V`$ there exists a $`GL^2(Q,d\mu )`$ such that
$$e^{i\mathrm{\Phi }_\mathrm{\Lambda }(g)}=e^{iG}e^{i\mathrm{\Phi }(g)}e^{iG}.$$
But $`[T,e^{iG}]=0`$ entails equation (3.2.1) for $`gD(\mathrm{\Lambda })V`$ and, therefore, by Lemma 3.2.1, also for any $`gH`$. Since any Fock representation is irreducible, it follows that $`T`$ is a multiple of the identity. $`\overline{)}`$
We wish to show that there do exist $`\mathrm{\Lambda }_{CCR}`$ such that $`\pi _\mathrm{\Lambda }`$ is reducible. To set this up properly, we first prove the following lemma.
###### Lemma 3.2.3
If $`𝒦`$ is a Hilbert space and $`T:𝒦D(T)𝒦`$ is an unbounded densely defined operator, then there exists a densely defined $`ST`$ such that $`\overline{R(S)}𝒦`$.
###### Demonstration Proof
The equality $`D(T^{})=𝒦`$ would imply the boundedness of $`T^{}`$ and therefore of $`T`$. Thus, there exists an $`f𝒦D(T^{})`$, with which one may define the operator $`S`$ as the restriction of $`T`$ to
$$D(S)=\{gD(T)f,Tg=0\}.$$
$`fD(T^{})`$ entails the existence of a unit vector $`f_nD(T)`$ such that
$$\underset{n\mathrm{}}{lim}|f,Tf_n|\mathrm{}.$$
It may be assumed that $`f_nD(S)`$. Let $`gD(T)`$ be arbitrary. Then
$$g\frac{f,Tg}{f,Tf_n}f_nD(S)\text{and}g\frac{f,Tg}{f,Tf_n}f_ng,$$
as $`n\mathrm{}`$. Hence, $`g\overline{D(S)}`$. Since $`D(T)`$ is dense in $`𝒦`$, this establishes that $`\overline{D(S)}=𝒦`$. $`\overline{)}`$
We can now show that there exists a linear $`\mathrm{\Lambda }_{CCR}`$ such that $`\pi _\mathrm{\Lambda }`$ is reducible. In fact, to each linear element $`\mathrm{\Lambda }`$ of $`_{CCR}`$ there corresponds a symplectic transformation, which is unbounded if $`\mathrm{\Lambda }`$ is unbounded (see ). According to Lemma 3.2.3, we may restrict the domain of any unbounded symplectic transformation to a set which is still dense in $`H`$ in such a manner that the range of the restriction is not dense in $`H`$. However, the representation induced by a symplectic transformation is irreducible if and only if the range of the symplectic operator is dense in $`H`$, as we shall see in the next chapter (Lemma 4.1).
We also wish to point out that Corollary 4.1.2 in is false as stated; in particular, the claim of irreducibility does not follow. As discussed in , the argument sketch given in tacitly assumed that $`\overline{\mathrm{\Lambda }}`$ is defined on a proper standard basis, which certainly follows if $`\overline{\mathrm{\Lambda }}=\mathrm{\Lambda }_{max}`$, and hence also if $`\mathrm{\Lambda }`$ is bounded, but which is not true in general. Florig also provides an example of a reducible quadratic representation (see Section 2.3 in ).
## IV. Quasifree States and Linear Canonical Transformations
In this chapter we shall restrict our attention to quasifree states on $`𝒜(H)`$ and the associated representations. The notion of quasifree state was introduced by D.W. Robinson in his study of the ground state of the Bose gas. It was shown in that such (pure) states can be obtained by Bogoliubov transformations of a Fock state, hence making it clear that the class of representations (commonly called symplectic representations) studied by Segal and Shale , among others, essentially coincided with the quasifree representations. We wish to show that the GNS representation of any pure quasifree state is unitarily equivalent to one of the representations $`\pi _\mathrm{\Lambda }`$ constructed in the previous chapter, for a suitable choice of $`VH`$ and a linear $`\mathrm{\Lambda }_{CCR}`$. The polynomial representations constructed globally under certain boundedness restrictions in and the quadratic representations of are clearly included among the representations of finite degree (special cases of the class constructed in Chapter III) discussed in more detail in Chapter V. Since the coherent representations are special cases of pure quasifree representations (see Proposition 4.4) and are therefore also subsumed in the class of representations presented in Chapter III, we see that our methods serve to unify the approaches to these various classes of representations, as well as to extend them to arbitrary degree.
We recall that if $`\omega _J`$ is a Fock state on $`𝒜(H)`$ with associated representation $`(𝒦,\pi _J)`$, then a coherent state $`\omega _l`$ is given by
$$\omega _l(W(f))\omega _J(W(f))e^{il(f)},fH.$$
The GNS representation of $`𝒜(H)`$ corresponding to $`\omega _l`$ is given on $`𝒦`$ by
$$\pi _l(W(f))e^{il(f)}\pi _J(W(f)),fH.$$
From Theorem 3.1 of (but see also ), it follows that the representations $`\pi _l`$ and $`\pi _J`$ are unitarily equivalent if and only if the map $`l:H\text{I R}`$ is bounded. And if a quasifree representation is obtained from a given Fock representation of $`𝒜(H)`$ by
$$\mathrm{\Phi }_T(f)\mathrm{\Phi }(Tf)=\mathrm{\Phi }(f)+\mathrm{\Phi }((T\text{1I})f),$$
using a symplectic operator $`T`$, it is known that the representations $`\pi _J`$ and $`\pi _T`$ are unitarily equivalent if and only if the operator $`\text{1I}|T|`$ is Hilbert-Schmidt. (See Theorem 3.2 in for a basis-dependent formulation of this result.)
Given a dense subspace $`H_0H`$, we shall show that all the quasifree states on $`𝒜(H_0)`$ can be obtained from the Fock state $`\omega _J`$ by symplectic transformations $`\mathrm{\Lambda }_{CCR}`$. To begin, we consider pure quasifree states. It is known that, in our terminology, a pure quasifree state is, up to a coherent transformation, a Fock state . In order to be more precise, we need to introduce some notation.
Let $`\omega _F^{}`$ be a Fock state on $`𝒜(H_0)`$, so there exists an associated scalar product $`s^{}`$ on $`H_0`$ , with respect to which the completion of $`H_0`$ will be denoted by $`H^{}`$ (and the induced scalar product on $`H^{}`$ will again be called $`s^{}`$). The scalar product $`s^{}`$ is such that the symplectic form $`\sigma `$ determining $`𝒜(H)`$ is continuous with respect to $`s^{}`$, when restricted to $`H_0`$ . Hence, the restriction of $`\sigma `$ to $`H_0`$ extends uniquely to a nondegenerate symplectic bilinear form $`\sigma ^{}`$ on $`H^{}`$. Moreover, there exists an operator $`J^{}:H^{}H^{}`$ which induces a complex structure on $`H^{}`$ , so that, in particular,
$$s^{}(f,g)=\sigma ^{}(J^{}f,g),$$
for all $`f,gH^{}`$. Since $`H_0`$ is dense in $`H`$, resp. $`H^{}`$, with respect to $`s`$, resp. $`s^{}`$, and $`\sigma `$ and $`\sigma ^{}`$ are nondegenerate, we may assume $`f=f^{}HH^{}`$, whenever
$$\sigma (f,h)=\sigma ^{}(f^{},h)fH,f^{}H^{},hH_0.$$
$`4.1`$
There is no loss of generality, since we do not assume that $`\sigma (f_1,f_2)=\sigma ^{}(f_1,f_2)`$ for $`f_1,f_2HH^{}`$.
The existence of a symplectic $`T:HH`$ with $`s^{}(f,f)=s(Tf,Tf)`$, for all $`fH`$, was already proven for $`H=H_0=H^{}`$ in . We shall need to generalize this result. First, we characterize symplectic maps and irreducible symplectic representations.
###### Lemma 4.1
An operator $`T:HD(T)H`$ is symplectic (with respect to $`\sigma `$), i.e. $`\sigma (Tf,Tg)=\sigma (f,g)`$ for all $`f,gD(T)`$, if and only if
$$JT^1JT^{}$$
(it is not assumed here that $`T^1`$, resp. $`T^{}`$, is necessarily densely defined). A self-adjoint operator $`T`$ is symplectic if and only if $`JT^1J=T`$. For symplectic $`T`$, the representation $`\pi _T`$ defined by
$$\pi _T(W(f))=e^{i\mathrm{\Phi }(Tf)},fD(T),$$
is irreducible if and only if $`\overline{R(T)}=H`$.
###### Demonstration Proof
Let $`T`$ be symplectic and $`gD(T)`$. Then one has for $`fD(JT^1J)`$
$`s(JT^1Jf,g)`$ $`=\sigma (T^1Jf,g)=\sigma (Jf,Tg)`$
$`=s(f,Tg)=s(T^{}f,g).`$
Hence, $`JT^1JT^{}`$. The converse follows from the equalities
$`\sigma (f,g)`$ $`=s(JT^1Tf,g)=s(T^{}JTf,g)`$
$`=s(JTf,Tg)=\sigma (Tf,Tg),`$
for all $`f,gD(T)`$. If $`T`$ is symplectic and self-adjoint, then $`T^1`$ is densely defined ($`T`$’s null space is trivial), so, by the above result, $`T`$ is a self-adjoint extension of $`JT^1J=J(T^{})^1J=(JT^1J)^{}`$, and thus $`T=JT^1J`$.
Turning to the characterization of irreducible symplectic representations, if $`gR(T)^{}`$, then $`e^{i\mathrm{\Phi }(Jg)}\pi _T(𝒜(R(T)))^{}`$, so that $`\pi _T`$ is reducible if $`\overline{R(T)}H`$. Assume now that $`\overline{R(T)}=H`$. If the sequence $`\{f_n\}_{n\text{I N}}H`$ converges to $`fH`$, then $`\{e^{i\mathrm{\Phi }(f_n)}\}_{n\text{I N}}`$ converges to $`e^{i\mathrm{\Phi }(f)}`$ strongly (use Lemma 3.1.1 with $`\mathrm{\Lambda }=0`$). Thus, an element of the commutant of $`\pi _T(𝒜(R(T)))`$ must commute with the elements of $`\pi _J(𝒜(H))`$, which is itself a Fock representation and, hence, irreducible. $`\overline{)}`$
With this in hand, we can now generalize the mentioned result of Manuceau and Verbeure.
###### Proposition 4.2
Given the above-established notation, there exists a subspace $`H_1H_0`$ of $`HH^{}`$ such that the following conditions are fulfilled. If one defines an operator $`K`$ by $`KJ^{}`$ and $`D(K)=\{fH_1J^{}fH_1\}`$, then $`JK:HD(JK)H`$ is a symplectic (with respect to $`\sigma `$) positive self-adjoint (with respect to $`s`$) operator, and $`T=(JK)^{1/2}`$ is a symplectic transformation with $`D(T)=H_1`$ and
$$s^{}(f,g)=s(Tf,Tg),f,gH_1.$$
###### Demonstration Proof
By using Zorn’s Lemma, it is easy to show that there exists a subspace $`H_1H_0`$ of $`HH^{}`$ such that $`H_1`$ is maximal with the property
$$\sigma (f,g)=\sigma ^{}(f,g),f,gH_1.$$
$`4.2`$
Consider the restriction $`s_{H_1}^{}`$ of the positive quadratic form $`s^{}`$ determined by the form core $`Q(s_{H_1}^{})=H_1`$. Because of the aforesaid maximality, $`H_1`$ is closed with respect to the norm
$$f=\sqrt{s^{}(f,f)+s(f,f)},fHH^{}.$$
The quadratic form $`s_{H_1}^{}`$ determines a self-adjoint operator $`A:HD(A)H`$ (use, e.g. Theorem VIII.15 in ), and the closure of $`D(A)`$ with respect to the above norm $``$ is $`H_1`$. Hence,
$$\sigma (f,JAg)=s(f,Ag)=s_{H_1}^{}(f,g)=\sigma ^{}(f,J^{}g),f,gD(A)H_1$$
is also true for $`fH_1`$. But the equality $`\sigma (f,JAg)=\sigma ^{}(f,J^{}g)`$, for any $`fH_1H_0`$ entails the equality $`JAg=J^{}gHH^{}`$ (see (4.1)) and, thus, by the maximality of $`H_1`$, the equality $`JAg=J^{}gH_1`$, for $`gD(A)`$. According to the definition of $`K`$, one has $`Kg=J^{}g`$, for $`gD(K)`$, hence $`Ag=JJ^{}g=JKg`$, for $`gD(A)`$. From (4.2) one sees
$`s(f,JKg)`$ $`=\sigma (f,Kg)=\sigma ^{}(f,Kg)`$
$`=\sigma ^{}(Kf,g)=\sigma (Kf,g)=s(JKf,g),`$
for $`f,gD(JK)H_1`$ (so $`Kf,KgH_1`$), thus the operator $`JKA`$ is symmetric. Hence, $`JK=A`$ is positive and self-adjoint, and one can define the positive self-adjoint operator $`T=(JK)^{1/2}`$. The equality $`s^{}(f,g)=s(Tf,Tg)`$, which holds for all $`f,gD(A)`$, is therefore still true for $`f,gD(T)=H_1`$, which is the closure of $`D(A)`$ with respect to the norm $``$.
The operator $`JK`$ is symplectic with respect to $`\sigma `$, since
$$\sigma (JKf,JKg)=\sigma (Kf,Kg)=\sigma ^{}(Kf,Kg)=\sigma ^{}(f,g)=\sigma (f,g),$$
for $`f,gD(JK)H_1`$ (so $`Kf,KgH_1`$).
It remains to prove that $`T`$ is symplectic. From and Lemma 4.1, one can decompose $`H`$ and $`JK`$ as follows:
$$H=UJU,JK=LJL^1J,$$
with $`L:UU`$ self-adjoint and $`0L\text{1I}`$. Therefore, $`T=L^{1/2}JL^{1/2}J=JT^1J`$ is symplectic, by Lemma 4.1. $`\overline{)}`$
This permits us to characterize pure quasifree states. (If the proof is not yet clear, then read the first few lines of the proof of Proposition 4.4.)
###### Corollary 4.3
Let $`H_0`$ be a dense subspace of $`H`$. Each pure quasifree state $`\omega `$ on $`𝒜(H_0)`$ has a characteristic function of the form
$$\omega (W(f))=e^{il(f)\frac{s(Tf,Tf)}{4}},fH_0,$$
for some linear form $`l:H_0\text{I R}`$ and a symplectic positive self-adjoint operator $`T:H_0D(T)H`$.
With this result, we can characterize general quasifree states.
###### Proposition 4.4
Let $`H_0`$ be a dense subspace of $`H`$. Each quasifree state $`\omega `$ on $`𝒜(H_0)`$ has a characteristic function of the form
$$\omega (W(f))=e^{il(f)\frac{s(Tf,Tf)}{4}},fH_0,$$
for some linear form $`l:H_0\text{I R}`$ and a symplectic operator $`T:H_0D(T)H`$.
###### Demonstration Proof
By , $`\omega `$ has a characteristic function of the form
$$\omega (W(f))=e^{il(f)\frac{s^{}(f,f)}{4}},fH_0,$$
for some scalar product $`s^{}`$ on a Hilbert space $`MH_0`$. The symplectic form $`\sigma `$ can be continuously (with respect to $`s^{}`$) extended to a bilinear form $`\sigma ^{}`$ on $`M`$ (see inequality (2) in ). Let $`P:MM`$ be the orthogonal projection onto the closure of $`\{fM\sigma (f,g)=0,gM\}`$ and set $`Q=\text{1I}P`$. The inequality (see, once again, (2) in )
$$|\sigma ^{}(f,g)|^2=|\sigma ^{}(Qf,Qg)|^2s^{}(Qf,Qf)s^{}(Qg,Qg),$$
for all $`f,gM`$, implies that one can define a state $`\omega ^{}`$ on $`𝒜(H_0)`$ by
$$\omega ^{}(W(f))=e^{s^{}(Qf,Qf)/4},$$
for all $`fH_0`$, by Proposition 10 in . According to Proposition 11 in the same paper, the state $`\omega ^{}`$ is also primary, since the restriction of $`\sigma ^{}`$ to $`QM\times QM`$ is nondegenerate. The discussion in Section IV of also implies the existence of a scalar product $`s_0`$ on $`H_0`$ associated with a pure state on $`𝒜(H_0)`$ such that
$$s_0(f,f)s^{}(Qf,Qf),$$
for all $`fH_0`$. Thus, there exists a symplectic operator $`S:H_0H`$ such that
$$s(Sf,Sf)=s_0(f,f)s^{}(Qf,Qf)s^{}(f,f),$$
for all $`fH_0`$, using Corollary 4.3. Furthermore, from Theorem VIII.15 in one has
$$s^{}(f,g)=s(A^{\frac{1}{2}}f,A^{\frac{1}{2}}g),f,gD(A^{\frac{1}{2}})H_0,$$
for a suitable self-adjoint $`A`$. But then $`s(A^{\frac{1}{2}},A^{\frac{1}{2}})s(S,S)`$ is a positive quadratic form on $`H_0`$, so, by appealing once again to Theorem VIII.15 in , there exists a positive self-adjoint operator $`B:HD(B)H`$ such that
$$s(Bf,Bf)=s(A^{\frac{1}{2}}f,A^{\frac{1}{2}}f)s(Sf,Sf),$$
$`4.3`$
for all $`fH_0`$.
Now define isometries (with respect to $`s`$) $`U,V:HH`$ by
$$Ue_k=e_{3k},UJe_k=Je_{3k}\text{and}Ve_k=e_{3k+1},VJe_k=Je_{3k+2},k\text{I N},$$
and set $`T=US+VB`$. Then the equalities (by definition, $`U`$ commutes with $`J`$)
$`\sigma (Tf,Tg)`$ $`=\sigma ((US+VB)f,(US+VB)g)=\sigma (USf,USg)`$
$`=\sigma (Sf,Sg)=\sigma (f,g),`$
for all $`f,gD(T)`$, entail that $`T`$ is symplectic. The claim then follows after noting that (4.3) implies
$$s(Tf,Tf)=s(Sf,Sf)+s(Bf,Bf)=s(A^{\frac{1}{2}}f,A^{\frac{1}{2}}f)=s^{}(f,f),$$
for all $`fH_0`$. $`\overline{)}`$
It is now also clear that coherent representations are special cases of pure quasifree representations. This permits us to prove the result announced at the beginning of the chapter.
###### Theorem 4.5
Let $`\pi `$ be the GNS-representation associated to a pure quasifree state on $`𝒜(H_0)`$, where $`H_0`$ is a dense subspace of $`H`$. For a suitable choice of $`VH`$, there exists a linear $`\mathrm{\Lambda }_{CCR}`$ such that $`\pi `$ is unitarily equivalent to $`\pi _\mathrm{\Lambda }`$.
###### Demonstration Proof
From the above discussion, it may be assumed that there exist a positive self-adjoint operator $`T:HD(T)H`$ with $`H_0D(T)`$ and a linear form $`l:H_0\text{I R}`$ such that
$$\pi (W(f))=e^{i\mathrm{\Phi }(Tf)+il(f)},$$
for all $`fH_0`$. From the proof of Proposition 4.2 it is clear that $`H`$ and $`T`$ decompose as
$$H=UJU,T=AJA^1J,$$
with $`A:UU`$ self-adjoint and satisfying $`0A\text{1I}`$. Set
$$U_1=\{(\text{1I}+JA)\phi \phi U\}H.$$
It is easy to see that $`JU_1U_1^{}`$. With $`\phi U`$, the inclusions
$$\phi =(\text{1I}+JA)\frac{1}{\text{1I}+A^2}\phi J(\text{1I}+JA)\frac{A}{\text{1I}+A^2}\phi U_1JU_1$$
and
$$J\phi J(U_1JU_1)=U_1JU_1$$
imply $`H=U_1JU_1`$. The restriction of $`T`$ to $`U_1`$ is then an isometry, since
$$T(\text{1I}+JA)\phi =(A+(JA^1J)JA)\phi =J(JA+\text{1I})\phi ,$$
for all $`\phi U`$, and
$`T(\text{1I}+JA)\phi ^2`$ $`=J(JA+\text{1I})\phi ^2=JA\phi ^2+\phi ^2`$
$`=(\text{1I}+JA)\phi ^2.`$
One can similarly prove that $`H=\overline{TU_1}J\overline{TU_1}`$, so there exist a unitary (considering $``$ instead of $`H`$) $`W:HH`$ such that $`WT`$ is the identity on $`U_1`$. The equalities
$$0=\sigma (WTf,WTg)\sigma (f,g)=\sigma (f,(WT\text{1I})g),$$
for $`fU_1`$, $`gD(T)`$, entail $`R(WT\text{1I})U_1`$. Choose now $`V=U_1`$ and define $`\mathrm{\Lambda }:H_0L^2(Q,d\mu )`$ by $`\mathrm{\Lambda }f=l(f)+x((WT\text{1I})f)`$ for any $`fH_0`$. Then $`\mathrm{\Lambda }_{CCR}`$ and $`\mathrm{\Phi }_\mathrm{\Lambda }(f)=\mathrm{\Phi }(WTf)+l(f)\text{1I}`$. According to , since $`|W|\text{1I}=0`$ is a Hilbert-Schmidt operator, there exists a unitary $`L`$ such that $`L\mathrm{\Phi }(f)L^{}=\mathrm{\Phi }(Wf)`$, for all $`fH`$. In particular, one has $`L\mathrm{\Phi }(Tf)L^{}=\mathrm{\Phi }(WTf)`$, for all $`fD(T)`$. It is therefore clear that the representation $`\pi _\mathrm{\Lambda }`$ corresponding to $`\mathrm{\Phi }_\mathrm{\Lambda }(f)`$ is unitarily equivalent to the representation $`\pi `$ given by
$$\pi (W(f))=e^{i\mathrm{\Phi }(Tf)+il(f)},fH_0.$$
$`\overline{)}`$
To close this chapter, we give a characterization of our linear canonical transformations. Recall that $`P:HH`$ is the orthogonal projection onto the subspace $`JV`$.
###### Proposition 4.6
To every linear $`\mathrm{\Lambda }_{CCR}`$ there corresponds a linear form $`l:D(\mathrm{\Lambda })\text{I R}`$ and a symmetric operator $`S:VD(S)=JPD(\mathrm{\Lambda })V`$ such that
$$\mathrm{\Lambda }f=x(SJPf)+l(f),$$
$`4.4`$
for every $`fD(\mathrm{\Lambda })`$. Moreover, each such pair $`(l,S)`$ defines a linear $`\mathrm{\Lambda }_{CCR}`$ in this manner, with $`D(\mathrm{\Lambda })=D(SJP)`$.
###### Demonstration Proof
Let $`\mathrm{\Lambda }_{CCR}`$. From Theorem 3.1.6 it may be assumed that $`VD(\mathrm{\Lambda })`$ and that $`\mathrm{\Lambda }`$ has the form given in (4.4). Thus, one sees that
$$\mathrm{\Phi }_\mathrm{\Lambda }(f)=\mathrm{\Phi }(f)+\mathrm{\Lambda }f\text{1I}=\mathrm{\Phi }(f)+(x(SJPf)+l(f))\text{1I}=\mathrm{\Phi }(f+SJPf)+l(f)\text{1I},$$
since $`R(S)V`$. But since $`\mathrm{\Lambda }_{CCR}`$, these operators must satisfy the CCR; hence, the operator $`\text{1I}+SJP`$ must be symplectic. The resultant equalities
$`\sigma (f,g)`$ $`=\sigma (f+SJPf,g+SJPg)`$ $`4.5`$
$`=\sigma (f,g)+\sigma (f,SJPg)+\sigma (SJPf,g)`$
$`=\sigma (f,g)+\sigma (Pf,SJPg)+\sigma (SJPf,g)`$
$`=\sigma (f,g)+s(JPf,SJPg)s(SJPf,JPg),`$
for all $`f,gD(\mathrm{\Lambda })`$, imply $`s(JPf,SJPg)=s(SJPf,JPg)`$, in other words, $`S`$ is a symmetric operator. The same computation (4.5) shows that if $`S`$ is symmetric, then (4.4) defines an element $`\mathrm{\Lambda }_{CCR}`$. $`\overline{)}`$
## V. Canonical Transformations of Finite Degree
In this chapter we restrict our attention to the computationally simpler canonical transformations of arbitrary but finite degree.<sup>6</sup><sup>6</sup>6These are the counterparts in our approach to the polynomial representations of . Note, however, that due to the boundedness assumptions made in , which do not need to be made here, we shall be discussing a larger class of representations than does . In any case, the questions treated below are not addressed in — at the cost of the additional technicalities involved in working infinitesimally with representations of the CCR, one gains a more detailed computational power than one apparently can attain when working globally from the outset. We provide necessary and sufficient conditions on the mapping $`\mathrm{\Lambda }`$ of finite degree so that the associated representation $`\pi _\mathrm{\Lambda }`$ of the CCR is unitarily equivalent to a Fock, a coherent or a quasifree representation. As we show, these results contain the previously known conditions for the unitary equivalence of irreducible quasifree representations. The case of unitary equivalence with a quadratic representation is briefly indicated at the end of the chapter. We emphasize that when the conditions isolated in this chapter are violated, one has a representation of the CCR which can describe bosonic systems with infinitely many degrees of freedom manifesting physics different from that describable by the Fock, coherent, quasifree or quadratic representations.
###### Definition 5.1
Let $`_{CCR}^{(n)}`$ denote the set of elements $`\mathrm{\Lambda }`$ of $`_{CCR}`$ such that
$$\mathrm{\Lambda }=\underset{i=0}{\overset{𝑛}{}}P_i\mathrm{\Lambda }\underset{i=0}{\overset{n1}{}}P_i\mathrm{\Lambda }.$$
Such elements and the corresponding canonical transformations and representations will be said to be of degree $`n`$.
Using the equality (3.1.1) given above and the estimate (4.3.5) given in , the following lemma can be easily proven.
###### Lemma 5.2
There exist real constants $`C_{lm}>0`$ such that
$$\phi _l\phi _m_2C_{lm}\phi _l_2\phi _m_2,$$
for all $`\phi _lR(\underset{jl}{}P_j)`$ and $`\phi _mR(\underset{jm}{}P_j)`$. Therefore, the finite-particle vectors are contained in the domain of the field operators $`\mathrm{\Phi }_\mathrm{\Lambda }(f)`$, for all $`fD(\mathrm{\Lambda })`$, whenever $`\mathrm{\Lambda }_{CCR}^{(n)}`$.
Another straightforward fact we shall need is given in the next lemma.
###### Lemma 5.3
Let $`\mathrm{\Lambda }_{CCR}^{(n)}`$. Then one has
$$P_m\mathrm{\Phi }_\mathrm{\Lambda }(f)P_m\mathrm{\Phi }_\mathrm{\Lambda }(f)\underset{kn+m+1}{}P_k,$$
for any $`m\text{I N}\{0\}`$ and $`fD(\mathrm{\Lambda })`$.
###### Demonstration Proof
For arbitrary $`\phi D(\mathrm{\Phi }_\mathrm{\Lambda }(f))`$ and $`\psi L^2(Q,d\mu )`$, one sees from Lemma 5.2 that
$`\psi ,P_m\mathrm{\Phi }_\mathrm{\Lambda }(f)\phi `$ $`=\mathrm{\Phi }_\mathrm{\Lambda }(f)P_m\psi ,\phi `$
$`=\underset{k\mathrm{max}\{n+m,m+1\}}{{\displaystyle }}P_k\mathrm{\Phi }_\mathrm{\Lambda }(f)P_m\psi ,\phi `$
$`=\psi ,P_m\mathrm{\Phi }_\mathrm{\Lambda }(f)\underset{kn+m+1}{{\displaystyle }}P_k\phi .`$
$`\overline{)}`$
We give a characterization of the existence of strong graph limits of sequences of field operators in representations of degree $`n`$.
###### Lemma 5.4
Let $`\mathrm{\Lambda }_{CCR}^{(n)}`$ and $`\{f_m\}_{m\text{I N}}D(\mathrm{\Lambda })`$. There exist vectors $`g_m,h_mV`$ such that $`f_m=g_m+Jh_m`$, for every $`m\text{I N}`$. The strong graph limit of the sequence of operators $`\{\mathrm{\Phi }_\mathrm{\Lambda }(f_m)\}`$ exists if and only if the sequences $`\{h_m\}`$ and $`\{x(g_m)+\mathrm{\Lambda }f_m\}`$ converge in their respective Hilbert spaces. If this strong graph limit exists and the sequence $`\{f_m\}`$ converges, then also the sequence $`\{\mathrm{\Lambda }f_m\}`$ converges.
###### Demonstration Proof
Assume that the strong graph limit indicated exists (the other direction can be proven using the argument of Lemma 3.2.1). Since this limit is densely defined, for arbitrary $`ϵ>0`$ there exists a sequence $`\{\phi _m\}_{m\text{I N}}`$ such that both it and the sequence $`\{\mathrm{\Phi }_\mathrm{\Lambda }(f_m)\phi _m\}_{m\text{I N}}`$ converge and such that $`\phi _m\mathrm{\Omega }_2<ϵ`$, for all $`m\text{I N}`$. Hence, the sequence $`\{P_0\phi _m\}_{m\text{I N}}`$ converges and one has $`P_0\phi _m>1ϵ`$. One may therefore choose the sequence $`\{\phi _m\}_{m\text{I N}}`$ such that, in addition, one has $`P_0\phi _m=\mathrm{\Omega }`$, for all $`m\text{I N}`$. There also exists a real constant $`C`$ such that $`CC_{kl}`$, for $`k,l2n+2`$, where the constants $`C_{kl}`$ are those evoked in Lemma 5.2, and thus
$$\mathrm{\Phi }(f)\phi _2Cf\phi _2,$$
for arbitrary $`\phi `$ in the range of the projection $`_{l2n+2}P_l`$ and $`fH`$. One also has the following estimate, using previously established notation:
$`\mathrm{\Phi }_\mathrm{\Lambda }(f_m)\mathrm{\Omega }_2`$ $`=\mathrm{\Phi }(f_m)\mathrm{\Omega }+\mathrm{\Lambda }f_m_2`$ $`5.1`$
$`=\mathrm{\Phi }(Jh_m)\mathrm{\Omega }+x(g_m)+\mathrm{\Lambda }f_m_2`$
$`=ix(h_m)+x(g_m)+\mathrm{\Lambda }f_m_2`$
$`=(x(h_m)_2^2+x(g_m)+\mathrm{\Lambda }f_m_2^2)^{1/2}`$
$`\stackrel{(\mathrm{3.1.1})}{=}({\displaystyle \frac{1}{2}}h_m_2^2+x(g_m)+\mathrm{\Lambda }f_m_2^2)^{1/2}`$
$`{\displaystyle \frac{1}{2}}\mathrm{max}\{h_m,x(g_m)+\mathrm{\Lambda }f_m_2\}`$
$`{\displaystyle \frac{1}{4}}(h_m+x(g_m)+\mathrm{\Lambda }f_m_2).`$
Of course, also the sequence $`\{_{j=0}^nP_j\mathrm{\Phi }_\mathrm{\Lambda }(f_m)\phi _m\}_{m\text{I N}}`$ converges. Now consider the estimate (obtained using Lemma 5.3)
$`\underset{j=0}{\overset{n+1}{{\displaystyle }}}P_j\mathrm{\Phi }_\mathrm{\Lambda }(f_m)\phi _m\mathrm{\Phi }_\mathrm{\Lambda }(f_m)\mathrm{\Omega }_2=\underset{j=0}{\overset{n+1}{{\displaystyle }}}P_j\mathrm{\Phi }_\mathrm{\Lambda }(f_m)\underset{l=1}{\overset{2n+2}{{\displaystyle }}}P_l\phi _m_2`$
$`\underset{j=0}{\overset{n+1}{{\displaystyle }}}P_j\mathrm{\Phi }(Jh_m)\underset{l=1}{\overset{2n+2}{{\displaystyle }}}P_l\phi _m_2+\underset{j=0}{\overset{n+1}{{\displaystyle }}}P_j(x(g_m)+\mathrm{\Lambda }f_m)\underset{l=1}{\overset{2n+2}{{\displaystyle }}}P_l\phi _m_2`$
$`\mathrm{\Phi }(Jh_m)\underset{l=1}{\overset{2n+2}{{\displaystyle }}}P_l\phi _m_2+(x(g_m)+\mathrm{\Lambda }f_m)\underset{l=1}{\overset{2n+2}{{\displaystyle }}}P_l\phi _m_2`$
$`C(h_m+x(g_m)+\mathrm{\Lambda }f_m_2)\underset{l=1}{\overset{2n+2}{{\displaystyle }}}P_l\phi _m_2`$
$`C(h_m+x(g_m)+\mathrm{\Lambda }f_m_2)ϵ`$
$`\stackrel{(5.1)}{}\mathrm{\Phi }_\mathrm{\Lambda }(f_m)\mathrm{\Omega }_2\mathrm{\hspace{0.17em}4}Cϵ.`$
From this estimate, one sees that the boundedness of the sequence
$`\{\underset{j=0}{\overset{n+1}{}}P_j\mathrm{\Phi }_\mathrm{\Lambda }(f_m)\phi _m\}_{m\text{I N}}`$ entails the boundedness of the sequence $`\{\mathrm{\Phi }_\mathrm{\Lambda }(f_m)\mathrm{\Omega }_2\}_{m\text{I N}}`$. Thus, for any $`\delta >0`$, there exist convergent sequences $`\{\phi _m\}_{m\text{I N}}`$ and
$`\{\underset{j=0}{\overset{n+1}{}}P_j\mathrm{\Phi }_\mathrm{\Lambda }(f_m)\phi _m\}_{m\text{I N}}`$ such that
$$\underset{j=0}{\overset{n+1}{}}P_j\mathrm{\Phi }_\mathrm{\Lambda }(f_m)\phi _m\mathrm{\Phi }_\mathrm{\Lambda }(f_m)\mathrm{\Omega }_2<\delta ,$$
for all $`m\text{I N}`$. This implies the convergence of the sequence $`\{\mathrm{\Phi }_\mathrm{\Lambda }(f_m)\mathrm{\Omega }\}`$. It may therefore be concluded that the sequences $`\{h_m\}_{m\text{I N}}`$ and $`\{x(g_m)+\mathrm{\Lambda }f_m\}_{m\text{I N}}`$ converge (use (5.1) with $`f_m`$ replaced by $`f_{m_1}f_{m_2}`$). $`\overline{)}`$
It is easy to see that the following lemma is true. We simply record it here for later reference.
###### Lemma 5.5
Let $`\pi `$ be a representation of a $`C^{}`$-algebra $`𝒜`$, $`\{A_n\}`$ a sequence of elements of $`𝒜`$, and $`k`$ a cardinal number. The strong graph limit of the operator sequence $`\{\pi (A_n)\}`$ exists and is densely defined if and only if the strong graph limit of the operator sequence $`\{k\pi (A_n)\}`$ exists and is densely defined, where $`k\pi `$ is the direct sum of $`k`$ copies of $`\pi `$.
We can finally prove our characterization of the quasi-equivalence of a representation $`\pi _\mathrm{\Lambda }`$ of degree $`n`$ with the original Fock representation $`\pi _J`$. This generalizes Theorem 4.1 in , which was restricted to the case $`n=2`$. Of course, if $`\pi _\mathrm{\Lambda }`$ is irreducible (see Theorem 3.2.2), then quasi-equivalence implies unitary equivalence.
###### Theorem 5.6
Let $`\mathrm{\Lambda }_{CCR}^{(n)}`$. The representation $`\pi _\mathrm{\Lambda }`$ is quasi-equivalent to the restriction of the Fock representation $`\pi _J`$ to $`𝒜(D(\mathrm{\Lambda }))`$ if and only if the operator $`\overline{\mathrm{\Lambda }}`$ is Hilbert-Schmidt.
###### Demonstration Proof
Assume that the representations $`\pi _\mathrm{\Lambda }`$ and $`\pi _J_{𝒜(D(\mathrm{\Lambda }))}`$ are quasi-equivalent. Let $`fH`$ and $`\{f_n\}_{n\text{I N}}D(\mathrm{\Lambda })`$ be arbitrary with $`\{f_n\}_{n\text{I N}}`$ converging to $`f`$. The convergence of the sequence $`\{\mathrm{\Phi }(f_n)\phi \}_{n\text{I N}}`$ to $`\mathrm{\Phi }(f)\phi `$ for finite-particle vectors $`\phi `$ entails that $`\mathrm{\Phi }(f)`$ is the strong graph limit of $`\{\mathrm{\Phi }(f_n)\}_{n\text{I N}}`$ (see Theorem VIII.27 in ). As quasi-equivalence is the same as unitary equivalence up to multiplicity (see e.g. Theorem 2.4.26 in ), Lemma 5.5 implies that the strong graph limit of $`\{\mathrm{\Phi }_\mathrm{\Lambda }(f_n)\}_{n\text{I N}}`$ exists, so that, by Lemma 5.4, the sequence $`\{\mathrm{\Lambda }f_n\}_{n\text{I N}}`$ converges. This entails that $`\mathrm{\Lambda }`$ is bounded, and thus, by Lemma 3.2.1, one may assume that $`D(\mathrm{\Lambda })=H`$. According to Theorem 5.2.14 in there must exist a dense subset $`𝒦`$ of $`L^2(Q,d\mu )`$ such that
$$\underset{k=1}{\overset{}{}}(\mathrm{\Phi }_\mathrm{\Lambda }(e_k)+i\mathrm{\Phi }_\mathrm{\Lambda }(Je_k))\phi _2^2<\mathrm{},$$
$`5.2`$
for all $`\phi 𝒦`$, since a densely defined number operator exists. Thus, one must have
$`\underset{k=1}{\overset{}{{\displaystyle }}}\underset{l=0}{\overset{𝑛}{{\displaystyle }}}P_l(\mathrm{\Phi }_\mathrm{\Lambda }(e_k)+i\mathrm{\Phi }_\mathrm{\Lambda }(Je_k))\phi _2^2`$
$`=\underset{k=1}{\overset{}{{\displaystyle }}}\underset{l=0}{\overset{𝑛}{{\displaystyle }}}P_l(\mathrm{\Phi }_\mathrm{\Lambda }(e_k)+i\mathrm{\Phi }_\mathrm{\Lambda }(Je_k))\underset{m=0}{\overset{2n+1}{{\displaystyle }}}P_m\phi _2^2`$
$`<\mathrm{},`$
and, since the finite-particle vectors are in the domain of the number operator of the Fock representation $`\pi _J`$,
$$\underset{k=1}{\overset{}{}}\underset{l=0}{\overset{𝑛}{}}P_l(\mathrm{\Phi }_\mathrm{\Lambda }(e_k)+i\mathrm{\Phi }_\mathrm{\Lambda }(Je_k)\sqrt{2}a(e_k))\underset{m=0}{\overset{2n+1}{}}P_m\phi _2^2<\mathrm{}.$$
$`5.3`$
For each $`ϵ>0`$ there exists a vector $`\phi _ϵ𝒦`$ such that $`\phi _ϵ\mathrm{\Omega }_2<ϵ`$; and set
$$C\underset{r,s2n+1}{\mathrm{max}}C_{rs}$$
(see Lemma 5.2). Then, since $`(\mathrm{\Lambda }e_k+i\mathrm{\Lambda }Je_k)P_0\phi _ϵ_2=(\mathrm{\Lambda }e_k+i\mathrm{\Lambda }Je_k)_2P_0\phi _ϵ_2`$ ($`P_0\phi _ϵ`$ is just a constant function), one has the estimate
$`\underset{l=0}{\overset{𝑛}{{\displaystyle }}}`$ $`P_l(\mathrm{\Phi }_\mathrm{\Lambda }(e_k)+i\mathrm{\Phi }_\mathrm{\Lambda }(Je_k)\sqrt{2}a(e_k))\underset{m=0}{\overset{2n+1}{{\displaystyle }}}P_m\phi _ϵ_2`$ $`5.4`$
$`=\underset{l=0}{\overset{𝑛}{{\displaystyle }}}P_l(\mathrm{\Lambda }e_k+i\mathrm{\Lambda }Je_k)\underset{m=0}{\overset{2n+1}{{\displaystyle }}}P_m\phi _ϵ_2`$
$`\underset{l=0}{\overset{𝑛}{{\displaystyle }}}P_l(\mathrm{\Lambda }e_k+i\mathrm{\Lambda }Je_k)P_0\phi _ϵ_2\underset{m=1}{\overset{2n+1}{{\displaystyle }}}\underset{l=0}{\overset{𝑛}{{\displaystyle }}}P_l(\mathrm{\Lambda }e_k+i\mathrm{\Lambda }Je_k)P_m\phi _ϵ_2`$
$`(\mathrm{\Lambda }e_k+i\mathrm{\Lambda }Je_k)P_0\phi _ϵ_2\underset{m=1}{\overset{2n+1}{{\displaystyle }}}(\mathrm{\Lambda }e_k+i\mathrm{\Lambda }Je_k)P_m\phi _ϵ_2`$
$`\mathrm{\Lambda }e_k+i\mathrm{\Lambda }Je_k_2P_0\phi _ϵ_2\underset{m=1}{\overset{2n+1}{{\displaystyle }}}C\mathrm{\Lambda }e_k+i\mathrm{\Lambda }Je_k_2P_m\phi _ϵ_2`$
$`(1ϵ)\mathrm{\Lambda }e_k+i\mathrm{\Lambda }Je_k_2(2n+1)Cϵ\mathrm{\Lambda }e_k+i\mathrm{\Lambda }Je_k_2`$
$`=(1ϵ(2n+1)Cϵ)\mathrm{\Lambda }e_k+i\mathrm{\Lambda }Je_k_2.`$
Choosing $`ϵ>0`$ such that $`1ϵ(2n+1)Cϵ\frac{1}{2}`$, one sees that
$$\underset{𝑘}{}\mathrm{\Lambda }e_k+i\mathrm{\Lambda }Je_k_2^2<\mathrm{},$$
so that $`\mathrm{\Lambda }`$ is a Hilbert-Schmidt operator ($`\mathrm{\Lambda }e_k`$ and $`\mathrm{\Lambda }Je_k`$ are a.e. real-valued):
$$\underset{𝑘}{}(\mathrm{\Lambda }e_k_2^2+\mathrm{\Lambda }Je_k_2^2)=\underset{𝑘}{}\mathrm{\Lambda }e_k+i\mathrm{\Lambda }Je_k_2^2.$$
The asserted sufficiency of the condition will now be proven. Without loss of generality, one may assume that $`\mathrm{\Lambda }=\overline{\mathrm{\Lambda }}`$ and $`P_0\mathrm{\Lambda }=0`$. There exist real constants $`\lambda _{kk_1\mathrm{}k_m}`$ symmetric in the indices, such that (using (3.1.1))
$$\mathrm{\Lambda }Je_k=\underset{m,k_1,\mathrm{},k_m}{}\lambda _{kk_1\mathrm{}k_m}:x_{k_1}\mathrm{}x_{k_m}:$$
and
$$\underset{𝑘}{}\underset{m,k_1,\mathrm{},k_m}{}\lambda _{kk_1\mathrm{}k_m}^2\frac{m!}{2^m}=\underset{𝑘}{}\mathrm{\Lambda }Je_k_2^2<\mathrm{}.$$
Setting
$$G=\underset{m,k_1,\mathrm{},k_m}{}\frac{\lambda _{k_1\mathrm{}k_m}}{m}:x_{k_1}\mathrm{}x_{k_m}:L^2(Q,d\mu ),$$
one sees that
$$e^{iG}\mathrm{\Phi }(f)e^{iG}=\mathrm{\Phi }(f),$$
for all $`fV`$, as well as
$$e^{iG}\mathrm{\Phi }(Je_k)e^{iG}=\mathrm{\Phi }_\mathrm{\Lambda }(Je_k),$$
for all $`k\text{I N}`$, by Proposition 3.1.3. This then demonstrates that
$$e^{iG}\mathrm{\Phi }(g)e^{iG}=\mathrm{\Phi }_\mathrm{\Lambda }(g),$$
for all $`g`$ in the linear span of $`\{Je_kk\text{I N}\}`$. Lemma 3.2.1 completes the proof. $`\overline{)}`$
And next we give a characterization of the quasi-equivalence of a representation of degree $`n`$ with a quasifree representation.
###### Theorem 5.7
Let $`\mathrm{\Lambda }_{CCR}^{(n)}`$ and $`\pi `$ be a GNS-representation of a pure quasifree state on $`𝒜(D(\mathrm{\Lambda }))`$. There exists a symplectic operator $`K:D(\mathrm{\Lambda })H`$ such that $`\mathrm{\Phi }_{P_1\mathrm{\Lambda }}(f)=\mathrm{\Phi }(Kf)`$, for all $`fD(\mathrm{\Lambda })`$. The representations $`\pi _\mathrm{\Lambda }`$ and $`\pi `$ are quasi-equivalent if and only if the following conditions are fulfilled:
(i) $`\pi _{(P_0+P_1)\mathrm{\Lambda }}`$ and $`\pi `$ are quasi-equivalent.
(ii) The closure of the restriction of $`(\mathrm{\Lambda }P_0\mathrm{\Lambda }P_1\mathrm{\Lambda })_{max}`$ to the range of $`K`$ is Hilbert-Schmidt.
###### Demonstration Proof
By Lemma 4.1 and Corollary 4.4, it may be assumed that $`\pi `$ is of the form
$$\pi (W(f))=e^{i(\mathrm{\Phi }(Tf)+l(f))},fD(\mathrm{\Lambda }),$$
for a symplectic operator $`T:D(\mathrm{\Lambda })H`$ with dense range and a linear form $`l:D(\mathrm{\Lambda })\text{I R}`$. After a coherent transformation, one may assume that $`l`$ is the zero mapping. Let the representations $`\pi `$ and $`\pi _\mathrm{\Lambda }`$ be quasi-equivalent. Then the restriction of $`\pi _J`$ to $`𝒜(R(T))`$ and the representation $`\pi ^{}`$, defined by
$$\pi ^{}(W(f))=e^{i\mathrm{\Phi }_\mathrm{\Lambda }(T^1f)},fR(T),$$
are quasi-equivalent. The linear part of the field operators associated with $`\pi ^{}`$ is, up to constants, equal to $`\mathrm{\Phi }_{P_1\mathrm{\Lambda }}(T^1f)`$. There exists a symplectic transformation $`S:R(T)H`$ such that $`\mathrm{\Phi }(Sf)=\mathrm{\Phi }_{P_1\mathrm{\Lambda }}(T^1f)`$. As in the proof of Theorem 5.6, one may extend $`\pi ^{}`$ to a representation of $`𝒜(H)`$ which is quasi-equivalent to $`\pi _J`$. This extension will also be denoted by $`\pi ^{}`$. The operator $`\mathrm{\Lambda }T^1`$ is bounded; consequently $`S`$ is bounded, as well. It may be assumed that $`S=\overline{S}`$.
Proceeding as in the proof of Theorem 5.6, one needs to consider the linear part of the annihilation operators associated with $`\pi ^{}`$ separately. Furthermore, the operators $`a(e_k)`$ in equation (5.3) must be replaced by somewhat different terms, discussed below. Let $`\mathrm{\Phi }^{}(f)`$, $`fH`$, denote the field operators associated with $`\pi ^{}`$. The linear part of the associated annihilation operators $`\frac{1}{\sqrt{2}}(\mathrm{\Phi }^{}(e_k)+i\mathrm{\Phi }^{}(Je_k))`$ is given by
$`{\displaystyle \frac{1}{\sqrt{2}}}(\mathrm{\Phi }(Se_k)+i\mathrm{\Phi }(SJe_k))`$ $`={\displaystyle \frac{1}{2}}(a(Se_k)+a^{}(Se_k)+ia(SJe_k)+ia^{}(SJe_k))`$ $`5.5`$
$`={\displaystyle \frac{1}{2}}(a((SJSJ)e_k)+a^{}((S+JSJ)e_k)).`$
The operators $`SJSJ`$ and $`J`$ commute, so that, if $`SJSJ=UA`$ is the polar decomposition, then $`U`$ and $`J`$ commute. Moreover, since $`R(A)=`$, one has $`U^{}U=\text{1I}`$. Therefore, (5.3) is still fulfilled if $`a(e_k)`$ is replaced by $`a(Ue_k)`$.
The bound
$$(SJSJ2U)f(S+JSJ)f,$$
$`5.6`$
will be proven for every $`fH`$. Note that for an arbitrary unit vector $`fH`$, one has
$`Af^2`$ $`=(SJSJ)f^2`$
$`=Sf^2+JSJf^22Sf,JSJf`$
$`=Sf^2+SJf^2+2\sigma (Sf,SJf)`$
$`=Sf^2+SJf^2+2\sigma (f,Jf)`$
$`=Sf^2+SJf^2+2.`$
Similarly, it follows that
$$(S+JSJ)f^2=Sf^2+SJf^22.$$
From the bound
$$1=|\sigma (SJf,Sf)|SJfSf$$
it follows that
$$Af^2Sf^2+\frac{1}{Sf^2}+24,$$
and $`A2\text{1I}`$. This proves (5.6), as $`A2\text{1I}`$ and $`U^{}U=\text{1I}`$ imply
$`(SJSJ2U)f^2`$ $`=(A2\text{1I})f^2`$
$`=f,(A^24A+4\text{1I})f`$
$`f,(A^28\text{1I}+4\text{1I})f`$
$`=(SJSJ)f^24`$
$`=(S+JSJ)f^2.`$
One can then proceed as in the proof of Theorem 5.6 to find suitable Hilbert-Schmidt conditions - one must simply modify (5.5) by the term $`a(Ue_k)`$, resp., replace $`a(e_k)`$ in (5.3) by $`a(Ue_k)`$. In the counterpart to (5.4) one finds the additional terms
$`a^{}((S+JSJ)e_k)P_0\phi _ϵ_2`$ $`\stackrel{(\mathrm{3.1.1})}{=}(S+JSJ)e_kP_0\phi _ϵ_2,`$
$`a^{}((S+JSJ)e_k)\underset{l=1}{\overset{2n+1}{{\displaystyle }}}P_l\phi _ϵ_2`$ $`Cϵ(S+JSJ)e_k,`$
and
$`a((SJSJ2U)`$ $`e_k)\underset{l=1}{\overset{2n+1}{{\displaystyle }}}P_l\phi _ϵ_2`$
$`C(SJSJ2U)e_k\underset{l=1}{\overset{2n+1}{{\displaystyle }}}P_l\phi _ϵ_2`$
$`\stackrel{(5.6)}{}C(S+JSJ)e_kϵ.`$
It may therefore be concluded that $`S+JSJ`$ is a Hilbert-Schmidt operator. Furthermore, the constant part of the field operator $`\mathrm{\Phi }^{}(f)`$ is a bounded linear form, so by Lemma 5.8 it may be assumed without loss of generality that it is trivial. Thus, the representations $`\pi `$ and $`\pi _{P_0\mathrm{\Lambda }+P_1\mathrm{\Lambda }}`$ are quasi-equivalent. The closure of the operator $`(\mathrm{\Lambda }P_0\mathrm{\Lambda }P_1\mathrm{\Lambda })T^1`$ is also Hilbert-Schmidt. That the mapping $`S^1:R(S)D(S)=H`$ is bounded is implied by the bound
$$SfSf|\sigma (Sf,SJf)|=f^2,$$
for every $`fH`$, as well as
$$Sf\frac{1}{S}f.$$
Hence, recalling $`S=\overline{KT^1}`$, the operator
$`\overline{(\mathrm{\Lambda }P_0\mathrm{\Lambda }P_1\mathrm{\Lambda })T^1}S^1`$ $`=\overline{(\mathrm{\Lambda }P_0\mathrm{\Lambda }P_1\mathrm{\Lambda })T^1TK^1}`$
$`=\overline{(\mathrm{\Lambda }P_0\mathrm{\Lambda }P_1\mathrm{\Lambda })_{max}PK^1}`$
$`=\overline{(\mathrm{\Lambda }P_0\mathrm{\Lambda }P_1\mathrm{\Lambda })_{max}R(K)}`$
is Hilbert-Schmidt.
To prove that the stated Hilbert-Schmidt conditions imply the desired quasi-equivalence, it is sufficient to show that the representation admits a densely defined number operator, in other words, by Theorem 5.2.14 in it suffices to show that for any finite-particle vector $`\phi =_{l=0}^mP_l\phi `$, and any orthonormal basis $`\{f_k,Jf_kk\text{I N}\}`$ of $`H`$, one has
$`\underset{𝑘}{{\displaystyle }}`$ $`(\mathrm{\Phi }^{}(f_k)+i\mathrm{\Phi }^{}(Jf_k))\phi \sqrt{2}a(Uf_k)\phi _2^2`$
$`D^{(m)}(S+JSJ_{HS}^2+\overline{(\mathrm{\Lambda }P_0\mathrm{\Lambda }P_1\mathrm{\Lambda })_{max}R(K)}_{HS}^2)\phi _2^2`$
$`<\mathrm{},`$
for a constant $`D^{(m)}\text{I R}`$ which does not depend on the choice of the basis (the Hilbert-Schmidt norm of an operator is denoted by $`_{HS}`$). This can be done by using arguments already employed above (see Lemma 5.2 and equations (5.5) and (5.6)). It follows then that
$$\underset{𝑘}{}(\mathrm{\Phi }^{}(f_k)+i\mathrm{\Phi }^{}(Jf_k))\phi ^2$$
is bounded by a finite real number, which only depends on $`m`$, $`\phi _2`$, $`S+JSJ_{HS}^2`$, and $`\overline{(\mathrm{\Lambda }P_0\mathrm{\Lambda }P_1\mathrm{\Lambda })_{max}R(K)}_{HS}^2`$. $`\overline{)}`$
Taken together, Theorems 5.6 and 5.7 give conditions so that our higher order representations are inequivalent to the well-studied class of representations associated with inhomogeneous linear canonical transformations, and hence so that they can describe bosonic systems with different physics.
###### Lemma 5.8
Let $`\pi `$ be a GNS-representation of a quasifree state $`\omega `$ on $`𝒜(H_0)`$, where $`H_0`$ is a dense subspace of $``$. There exists a linear form $`k:H_0\text{I R}`$ and a scalar product $`s^{}`$ on $`H_0`$ such that
$$\omega (W(f))=e^{ik(f)\frac{s^{}(f,f)}{4}},$$
for all $`fH_0`$ . If $`l:H_0\text{I R}`$ is a linear form, then the representation $`\pi _l`$ of $`𝒜(H_0)`$ defined by $`\pi _l(W(f))=e^{il(f)}\pi (W(f))`$ is quasi-equivalent to $`\pi `$ if and only if $`l`$ is bounded with respect to $`s^{}`$.
###### Demonstration Proof
By Corollary 4.4, it may be assumed that the linear form $`k`$ is trivial and that $`\pi (W(f))=e^{i\mathrm{\Phi }(Tf)}`$, $`fH_0`$, for a symplectic mapping $`T:H_0H`$. Suppose the indicated representations are quasi-equivalent. The convergence $`f_n\stackrel{s^{}}{}0`$ implies the convergence $`Tf_n\stackrel{𝑠}{}0`$, so that the strong graph limit of $`\{\mathrm{\Phi }(Tf_n)\}`$ is $`0`$. The assumed quasi-equivalence implies also that $`0`$ is the strong graph limit of the sequence $`\{\mathrm{\Phi }(Tf_n)+l(f_n)\}`$. Appealing to Lemma 5.4, one concludes that also $`\{l(f_n)\}`$ converges to $`0`$.
Assume now that the linear form $`l`$ is bounded with respect to $`s^{}`$, so that $`l(f)=s^{}(g,f)`$ for some element $`g`$ in the completion of $`H_0`$ with respect to $`s^{}`$. There exists a sequence $`\{g_n\}H_0`$ converging to $`g`$ with respect to $`s^{}`$; hence the sequence $`\{Tg_n\}`$ converges to some element $`hH`$ with respect to $`s`$. Then, for $`fH_0`$, one sees
$`e^{i\mathrm{\Phi }(Jh)}\mathrm{\Phi }(Tf)e^{i\mathrm{\Phi }(Jh)}`$ $`=\mathrm{\Phi }(Tf)\sigma (Jh,Tf)\text{1I}`$
$`=\mathrm{\Phi }(Tf)+s(h,Tf)\text{1I}`$
$`=\mathrm{\Phi }(Tf)+\underset{𝑛}{lim}s(Tg_n,Tf)\text{1I}`$
$`=\mathrm{\Phi }(Tf)+s^{}(g,f)\text{1I}`$
$`=\mathrm{\Phi }(Tf)+l(f)\text{1I}.`$
$`\overline{)}`$
We provide an extension to Proposition 3.1.7 for irreducible representations of finite degree.
###### Proposition 5.9
Let $`\mathrm{\Lambda }_{CCR}^q_{CCR}^{(n)}`$ determine an irreducible representation $`\pi _\mathrm{\Lambda }`$ of $`𝒜(D(\mathrm{\Lambda }))`$. If $`l:H_0\text{I R}`$ is a linear form such that the representation $`\pi _l`$ of $`𝒜(H_0)`$ defined by $`\pi _l(W(f))=e^{il(f)}\pi _\mathrm{\Lambda }(W(f))`$ is unitarily equivalent to $`\pi _\mathrm{\Lambda }`$, then $`l`$ has the form $`l(f)=\sigma (g,f)`$, for all $`fD(\mathrm{\Lambda })`$, for some $`gD(\mathrm{\Lambda }_{max})`$. Therefore, the set $`D(\mathrm{\Lambda }_{max})`$ is the maximal extension of the test function space $`D(\mathrm{\Lambda })`$.
###### Demonstration Proof
Let $`U:L^2(Q,d\mu )L^2(Q,d\mu )`$ be unitary such that
$$U^{}\mathrm{\Phi }_\mathrm{\Lambda }(f)U=\mathrm{\Phi }_\mathrm{\Lambda }(f)+l(f)\text{1I},$$
for all $`fD(\mathrm{\Lambda })`$. Then the operator $`\pi _{\mathrm{\Lambda }_{max}}(W(g))U^{}\pi _{\mathrm{\Lambda }_{max}}(W(g))U`$ commutes with the elements of $`\pi _\mathrm{\Lambda }(𝒜(D(\mathrm{\Lambda })))`$, for all $`gD(\mathrm{\Lambda }_{max})`$. By the assumed irreducibility, it must be a multiple of the identity. Thus, for each $`gD(\mathrm{\Lambda }_{max})`$ one has
$$U^{}\pi _{\mathrm{\Lambda }_{max}}(W(g))U=e^{ic_g}\pi _{\mathrm{\Lambda }_{max}}(W(g)),$$
for suitable $`c_g\text{I R}`$.
It shall next be shown that $`e^{ic_g}=e^{il^{}(g)}`$, for a suitable linear form $`l^{}:D(\mathrm{\Lambda }_{max})\text{I R}`$. But the equality $`e^{ic_{t_1g}}e^{ic_{t_2g}}=e^{ic_{(t_1+t_2)g}}`$ implies $`e^{ic_{tg}}=e^{itk}`$, for all $`t\text{I R}`$ and a suitable $`k\text{I R}`$. Thus, it has been shown that
$$U^{}\mathrm{\Phi }_{\mathrm{\Lambda }_{max}}(tg)U=\mathrm{\Phi }_{\mathrm{\Lambda }_{max}}(tg)+tk\text{1I},$$
in other words,
$$U^{}\pi _{\mathrm{\Lambda }_{max}}(W(f))U=e^{il^{}(f)}\pi _{\mathrm{\Lambda }_{max}}(W(f)),$$
for all $`fD(\mathrm{\Lambda }_{max})`$, with $`l^{}:D(\mathrm{\Lambda }_{max})\text{I R}`$ linear. Hence, there exists an extension $`l^{}:D(\mathrm{\Lambda }_{max})\text{I R}`$ of $`l`$ such that the representation $`\pi _l^{}`$ of $`𝒜(D(\mathrm{\Lambda }_{max}))`$ defined by $`\pi _l^{}(W(f))=e^{il^{}(f)}\pi _{\mathrm{\Lambda }_{max}}(W(f))`$, for $`fD(\mathrm{\Lambda }_{max})`$, is unitarily equivalent to $`\pi _{\mathrm{\Lambda }_{max}}`$. Therefore, it may be assumed that $`l=l^{}`$ and $`\mathrm{\Lambda }=\mathrm{\Lambda }_{max}`$.
Let $`\{e_k,Je_kk\text{I N}\}D(\mathrm{\Lambda })`$ be an orthonormal symplectic basis in $`H`$. For each $`fD((P_m\mathrm{\Lambda })_{max})`$, there exist real constants $`\lambda _{k_1\mathrm{}k_m}(f)`$, which are totally symmetric in the indices, such that
$$(P_m\mathrm{\Lambda })_{max}f=\underset{k_1,\mathrm{},k_m}{}\lambda _{k_1\mathrm{}k_m}(f):x_{k_1}\mathrm{}x_{k_m}:.$$
If $`\{f_n\}`$ is a convergent sequence in $`V`$, then the strong graph limit of $`\{\mathrm{\Phi }(f_n)\}`$ exists, so that also the strong graph limit of $`\{\mathrm{\Phi }(f_n)+l(f_n)\}`$ exists. Therefore, the restriction of $`l`$ to $`V`$ is bounded. Set $`l_kl(e_k)`$, $`h_kl_kJe_kH`$ and
$$A_l(f)\mathrm{\Phi }_\mathrm{\Lambda }(f)\underset{𝑚}{}\underset{k_1,\mathrm{},k_ml}{}\lambda _{k_1\mathrm{}k_m}(f):x_{k_1}\mathrm{}x_{k_m}:\text{1I},$$
with $`D(A_l(f))=D(\mathrm{\Phi }_\mathrm{\Lambda }(f))(_{k_1,\mathrm{},k_ml}D(x_{k_1}\mathrm{}x_{k_m}\text{1I}))`$. Then the operator $`A_l(f)`$ is symmetric and the sequence $`\{A_l(f)\phi \}_{l\text{I N}}`$ converges to $`\mathrm{\Phi }(f)\phi `$ for every $`\phi 𝒢`$. Since $`𝒢`$ is a core for $`\mathrm{\Phi }(f)`$ (Lemma 3.1.2), it follows from Theorem VIII.27 of that $`\mathrm{\Phi }(f)`$ is the strong graph limit of $`\{A_l(f)\}_{l\text{I N}}`$.
With $`U^{}x_kU=x_k+l_k`$ and $`p_l(f)`$ suitable polynomials of degree $`n2`$, one sees therefore that the strong graph limit of
$`U^{}`$ $`A_l(f)U`$ $`5.7`$
$`=U^{}\mathrm{\Phi }_\mathrm{\Lambda }(f)U\underset{m=2}{\overset{𝑛}{{\displaystyle }}}\underset{k_1,\mathrm{},k_ml}{{\displaystyle }}\lambda _{k_1\mathrm{}k_m}(f)U^{}:x_{k_1}\mathrm{}x_{k_m}:U`$
$`=\mathrm{\Phi }_\mathrm{\Lambda }(f)+(l(f)+p_l(f)\underset{k_1,\mathrm{},k_{n1}l}{{\displaystyle }}\lambda _{k_1\mathrm{}k_{n1}}(f):x_{k_1}\mathrm{}x_{k_{n1}}:)\text{1I}`$
$`\underset{k_1,\mathrm{},k_nl}{{\displaystyle }}(\lambda _{k_1\mathrm{}k_n}(f):x_{k_1}\mathrm{}x_{k_n}:+\lambda _{k_1\mathrm{}k_n}(f)nl_{k_n}:x_{k_1}\mathrm{}x_{k_{n1}}:)\text{1I}`$
is $`U^{}\mathrm{\Phi }(f)U`$. According to the proof of Lemma 5.4, the existence of the indicated strong graph limit entails the convergence of
$`p_l(f)`$ $`\underset{k_1,\mathrm{},k_{n1}l}{{\displaystyle }}\lambda _{k_1\mathrm{}k_{n1}}(f):x_{k_1}\mathrm{}x_{k_{n1}}:`$
$`\underset{k_1,\mathrm{},k_nl}{{\displaystyle }}(\lambda _{k_1\mathrm{}k_n}(f):x_{k_1}\mathrm{}x_{k_n}:+\lambda _{k_1\mathrm{}k_n}(f)nl_{k_n}:x_{k_1}\mathrm{}x_{k_{n1}}:).`$
Following the argument of Lemma 3.1, it can also be shown that the strong graph limit of $`\{U^{}A_l(f)U\}_{l\text{I N}}`$ is of the form $`\overline{\mathrm{\Phi }(f)+F\text{1I}}`$, for a suitable $`FR(_{k=0}^{n1}P_k)L^2(Q,d\mu )`$. Hence, there exists a $`\mathrm{\Lambda }^{}_{CCR}^{(n1)}`$ such that
$$U^{}\mathrm{\Phi }(f)U=\mathrm{\Phi }_\mathrm{\Lambda }^{}(f),$$
for all $`fD(\mathrm{\Lambda })=D(\mathrm{\Lambda }^{})`$.
According to Theorem 5.6, $`\mathrm{\Lambda }^{}`$ is a Hilbert-Schmidt operator. Hence one sees that
$`\mathrm{}`$ $`>\underset{𝑘}{{\displaystyle }}P_{n1}\mathrm{\Lambda }^{}(Je_k)_2^2`$
$`=\underset{𝑘}{{\displaystyle }}\underset{k_1,\mathrm{},k_n}{{\displaystyle }}\lambda _{k_1\mathrm{}k_n}(Je_k)l_{k_n}n:x_{k_1}\mathrm{}x_{k_{n1}}:_2^2`$
$`\stackrel{(\mathrm{3.1.1})}{=}\underset{𝑘}{{\displaystyle }}\underset{k_1,\mathrm{},k_{n1}}{{\displaystyle }}(\underset{k_n}{{\displaystyle }}\lambda _{k_1\mathrm{}k_n}(Je_k)l_{k_n}n)^2{\displaystyle \frac{(n1)!}{2^{n1}}}`$
$`=n^2{\displaystyle \frac{(n1)!}{2^{n1}}}\underset{𝑘}{{\displaystyle }}\underset{k_1,\mathrm{},k_{n1}}{{\displaystyle }}(\underset{k_n}{{\displaystyle }}\lambda _{k_1\mathrm{}k_n}(Je_k)l_{k_n})^2`$
$`\stackrel{(\text{Lemma}\mathrm{\hspace{0.17em}3.1.1})}{=}n{\displaystyle \frac{n!}{2^{n1}}}\underset{𝑘}{{\displaystyle }}\underset{k_1,\mathrm{},k_{n1}}{{\displaystyle }}(\underset{k_n}{{\displaystyle }}{\displaystyle \frac{2^n}{n!}}\mathrm{\Lambda }Je_k,:x_{k_1}\mathrm{}x_{k_{n1}}x_{k_n}:l_{k_n})^2`$
$`=2n\underset{k,k_1,\mathrm{},k_{n1}}{{\displaystyle }}{\displaystyle \frac{2^n}{n!}}\mathrm{\Lambda }Je_k,:x_{k_1}\mathrm{}x_{k_{n1}}x(Jh):^2,`$
which entails $`hD((P_n\mathrm{\Lambda })_{max})`$ (see the proof of Proposition 3.1.7).
There exists a $`GL^2(Q,d\mu )R(P_{n+1})`$ such that $`\mathrm{\Phi }_{(P_n\mathrm{\Lambda })_{max}}(h)=e^{iG}\mathrm{\Phi }(h)e^{iG}`$. By replacing $`U`$ by $`e^{iG}Ue^{iG}`$ and $`\mathrm{\Phi }_\mathrm{\Lambda }(f)`$ by $`e^{iG}\mathrm{\Phi }_\mathrm{\Lambda }(f)e^{iG}`$, which is the same as replacing $`\mathrm{\Lambda }`$ by another element of $`_{CCR}`$ differing from $`\mathrm{\Lambda }`$ only in the component of degree $`n`$, it may be assumed that $`(P_n\mathrm{\Lambda })_{max}h=0`$ and hence that the range of the mapping $`(P_n\mathrm{\Lambda })_{max}`$ lies in the subspace of $`L^2(Q,d\mu )`$ generated by
$$\{:x(f_1)\mathrm{}x(f_n):\sigma (h,f_1)=\mathrm{}=\sigma (h,f_n)=0,f_1,\mathrm{},f_nV\}$$
(use (3.1.3) and the fact that $`a(h)(P_n\mathrm{\Lambda })_{max}f=0`$, for $`fD((P_n\mathrm{\Lambda })_{max})`$). The operator $`e^{i\mathrm{\Phi }(h)}`$ commutes with these elements. From the definition of $`h`$, the adjoint actions of $`e^{i\mathrm{\Phi }(h)}`$ and $`U`$ on the field operators $`\mathrm{\Phi }_\mathrm{\Lambda }(f)=\mathrm{\Phi }(f)`$, $`fV`$, are identical, inducing the same coherent canonical transformation. Hence, $`e^{i\mathrm{\Phi }(h)}U`$ commutes with $`e^{ix(f)}`$, for all $`fV`$, and thereby may be identified with the multiplication operator corresponding to some suitable element of $`L^{\mathrm{}}(Q,d\mu )`$, as in the proof of Theorem 3.2.2. Thus, one has
$$U^{}((P_n\mathrm{\Lambda })_{max}f)U=(P_n\mathrm{\Lambda })_{max}f,$$
for all $`fD((P_n\mathrm{\Lambda })_{max})`$, and it therefore follows that
$$U^{}\mathrm{\Phi }_{(\mathrm{\Lambda }P_n\mathrm{\Lambda })_{max}}(f)U=\mathrm{\Phi }_{(\mathrm{\Lambda }P_n\mathrm{\Lambda })_{max}}(f)+l(f)\text{1I},$$
for all $`fD(\mathrm{\Lambda }_{max})`$.
Next, one can consider $`\mathrm{\Lambda }P_n\mathrm{\Lambda }`$ instead of $`\mathrm{\Lambda }`$ and prove $`hD((P_{n1}\mathrm{\Lambda })_{max})`$. Repeating this process finitely many times, one concludes that $`hD(\mathrm{\Lambda }_{max})`$ and
$$U^{}\mathrm{\Phi }(f)U=\mathrm{\Phi }(f)+l(f)\text{1I},$$
for all $`fD(\mathrm{\Lambda }_{max})`$, with $`U`$ a suitable unitary. As before, the boundedness of $`l`$ follows, and thus, again, the existence of a $`gH`$ such that $`l(f)=\sigma (g,f)`$, for all $`fH`$, is assured. But then the equalities
$$\sigma (g,e_k)=l(e_k)=l_k=\sigma (h,e_k),$$
for all $`k\text{I N}`$, imply that $`g+hV`$ and, thus, $`gD(\mathrm{\Lambda })`$, since $`VD(\mathrm{\Lambda })`$ and $`hD(\mathrm{\Lambda })`$. $`\overline{)}`$
###### Theorem 5.10
Let $`\mathrm{\Lambda }_{CCR}^{(n)}`$ and $`\pi `$ be a GNS-representation of a quasifree state on $`𝒜(D(\mathrm{\Lambda }))`$. Assume that $`\pi _\mathrm{\Lambda }`$ and $`\pi _{P_1\mathrm{\Lambda }}`$ are irreducible (see Theorem 3.2.2). The representations $`\pi _\mathrm{\Lambda }`$ and $`\pi `$ are quasi-equivalent if and only if the following conditions are fulfilled:
(i) $`\pi _{(P_0+P_1)\mathrm{\Lambda }}`$ and $`\pi `$ are quasi-equivalent.
(ii) The closure of the operator $`(\text{1I}P_0P_1)\mathrm{\Lambda }`$ is Hilbert-Schmidt.
###### Demonstration Proof
($``$) Since $`\pi _{P_1\mathrm{\Lambda }}`$ is irreducible and since quasi-equivalence is an equivalence relation, after applying a suitable linear transformation, it may be assumed that $`P_0\mathrm{\Lambda }0`$ and $`P_1\mathrm{\Lambda }0`$, and then Theorem 5.7 yields the quasi-equivalence of $`\pi `$ and $`\pi _\mathrm{\Lambda }`$.
($``$) Now let $`\pi `$ and $`\pi _\mathrm{\Lambda }`$ be quasi-equivalent. Let $`s^{}`$ be the scalar product associated to $`\pi `$ and $`H^{}`$ be the completion of $`D(\mathrm{\Lambda })`$ with respect to $`s^{}`$, as above. By Corollary 4.4, it may further be assumed that
$$\pi (W(f))=e^{il(f)}e^{i\mathrm{\Phi }(Tf)},fD(\mathrm{\Lambda }),$$
for a symplectic $`T:D(\mathrm{\Lambda })H`$ and a linear form $`l:D(\mathrm{\Lambda })\text{I R}`$. By applying a suitable coherent transformation, one has $`l0`$. By Lemma 5.5, if a sequence $`\{f_n\}`$ in $`D(\mathrm{\Lambda })`$ converges with respect to $`s^{}`$, then the strong graph limit of $`\{\mathrm{\Phi }(Tf_n)\}`$ exists, as does the strong graph limit of $`\{\mathrm{\Phi }_\mathrm{\Lambda }(f_n)\}`$. From Lemma 5.4 and $`P_1\mathrm{\Lambda }0`$, one notes that the sequences $`\{f_n\}`$ and $`\{\mathrm{\Lambda }f_n\}`$ also converge, so it may be assumed that $`D(\mathrm{\Lambda })`$ is closed with respect to $`s^{}`$, i.e. $`T`$ is closed. Theorem 5.9 entails that each coherent transformation of $`\pi _\mathrm{\Lambda }`$ inducing a quasi-equivalent representation is of the form
$$\mathrm{\Phi }_\mathrm{\Lambda }(f)\mathrm{\Phi }_\mathrm{\Lambda }(f)+\sigma (f,g)\text{1I},$$
for all $`fD(\mathrm{\Lambda })`$ and some suitable $`gH`$. From Lemma 5.8 one concludes that such coherent transformations of $`\pi `$ are of the form
$$\mathrm{\Phi }(Tf)\mathrm{\Phi }(Tf)+s^{}(f,g)\text{1I},$$
for all $`fD(\mathrm{\Lambda })`$ and some suitable $`gH`$. Thus, the assumed quasi-equivalence implies that, for arbitrary $`fD(\mathrm{\Lambda })`$ and arbitrary but fixed $`gD(\mathrm{\Lambda })`$, one has
$$s(Tf,Tg)=s^{}(f,g)=\sigma (f,g^{})=s(f,Jg^{}),$$
for a suitable $`g^{}H`$. Hence, $`TgD(T^{})`$, for all $`gD(T)=D(\mathrm{\Lambda })`$, which implies consecutively the boundedness of $`T`$ (since $`T`$ is closed) and then of $`\mathrm{\Lambda }`$.
Let $`T_V=U|T_V|`$ be the polar decomposition of the restriction $`T_V`$. By considering the quasifree quasi-equivalent representations $`\pi 𝒜(V)`$ and $`\pi _\mathrm{\Lambda }𝒜(V)`$, one may conclude that $`|T_V|\text{1I}`$ is a Hilbert-Schmidt operator. Since $`T`$ is symplectic, one has $`\sigma (f,g)=0`$, for all $`f,gR(U)`$, so that $`JR(U)R(U)^{}`$. Let $`P:HH`$, resp. $`Q:HH`$, be the orthogonal projection onto $`R(U)`$, resp. $`JR(U)`$. Note that $`|T_V|`$ is invertible, since for any $`0fH`$ there exists a $`gH`$ such that $`\sigma (Tf,Tg)=\sigma (f,g)0`$. Hence, one can define the operator $`S`$ on $`H`$ by
$$S=U|T_V|^1U^{}JU|T_V|U^{}J+(\text{1I}PQ).$$
By Lemma 4.1, $`S=JS^1J`$ is symplectic. Furthermore, the operator $`|S|\text{1I}=S\text{1I}`$ is Hilbert-Schmidt, so we may replace $`T`$ by $`ST`$, since quasi-equivalence is an equivalence relation. Note that $`T_V=U`$ is an isometry. By replacing $`T`$ with $`WT`$ for a suitable $`W`$, which is unitary when viewed as an operator on $``$, it may further be assumed that $`R(U)V`$.
There exist real constants $`\lambda _{k_1\mathrm{}k_m}(f)`$, symmetric in the indices, such that
$$\mathrm{\Lambda }f=\underset{m,k_1,\mathrm{},k_m}{}\lambda _{k_1\mathrm{}k_m}(f):x_{k_1}\mathrm{}x_{k_m}:,fH.$$
As in the proof of Theorem 5.8, one applies the inverse transformation of $`\mathrm{\Lambda }`$ to $`\pi _\mathrm{\Lambda }`$ and $`\pi `$. One then uses the strong graph limit of the sequence
$$\{\mathrm{\Phi }_\mathrm{\Lambda }(f)\underset{𝑚}{}\underset{k_1,\mathrm{},k_ml}{}\lambda _{k_1\mathrm{}k_m}(f):x_{k_1}\mathrm{}x_{k_m}:\text{1I}\}_{l\text{I N}},$$
which is $`\mathrm{\Phi }(f)`$, to show that $`\pi _J`$ is quasi-equivalent to a representation $`\pi ^{}`$ defined by
$$\pi ^{}(W(f))=e^{i\mathrm{\Phi }^{}(f)},fH,$$
where (note that $`x(Ue_k)=\mathrm{\Phi }(Te_k)`$)
$$\mathrm{\Phi }^{}(f)=\overline{\mathrm{\Phi }(Tf)\underset{m,k_1,\mathrm{},k_m}{}\lambda _{k_1\mathrm{}k_m}(f):x(Ue_{k_1})\mathrm{}x(Ue_{k_m}):\text{1I}}.$$
$`5.8`$
As in the proof of Theorem 5.6, one can derive Hilbert-Schmidt conditions. In equation (5.3), one only has to replace $`a(e_k)`$ by $`a(Ue_k)`$, since the linear term of the annihilation operators $`\frac{1}{\sqrt{2}}(\mathrm{\Phi }^{}(e_k)+i\mathrm{\Phi }^{}(Je_k))`$ is, by (5.8),
$`{\displaystyle \frac{1}{\sqrt{2}}}(\mathrm{\Phi }(Te_k)+i\mathrm{\Phi }(TJe_k))`$ $`={\displaystyle \frac{1}{\sqrt{2}}}(\mathrm{\Phi }(Ue_k)+i\mathrm{\Phi }(TJe_k))`$
$`=a(Ue_k)+{\displaystyle \frac{i}{\sqrt{2}}}\mathrm{\Phi }((TJJU)e_k).`$
In the counterpart to (5.4), one obtains the additional terms
$$\mathrm{\Phi }((TJJU)e_k)P_0\phi _ϵ_2\stackrel{(\mathrm{3.1.1})}{=}\frac{1}{\sqrt{2}}(TJJU)e_kP_0\phi _ϵ_2,$$
and
$`\mathrm{\Phi }((TJJU)e_k)\underset{l=1}{\overset{2n+1}{{\displaystyle }}}P_l\phi _ϵ_2`$ $`C(TJJU)e_k\underset{l=1}{\overset{2n+1}{{\displaystyle }}}P_l\phi _ϵ_2`$
$`Cϵ(TJJU)e_k.`$
One may therefore conclude that the closure of the mapping $`\mathrm{\Lambda }(P_0+P_1)\mathrm{\Lambda }`$ is Hilbert-Schmidt, thus implying the quasi-equivalence of $`\pi _\mathrm{\Lambda }`$ and $`\pi _{(P_0+P_1)\mathrm{\Lambda }}`$ and hence the quasi-equivalence of $`\pi `$ and $`\pi _{(P_0+P_1)\mathrm{\Lambda }}`$. $`\overline{)}`$
A further immediate consequence is the following result.
###### Corollary 5.11
Let $`\mathrm{\Lambda }_{CCR}^q_{CCR}^{(n)}`$ determine an irreducible representation $`\pi _\mathrm{\Lambda }`$ of $`𝒜(D(\mathrm{\Lambda }))`$. $`\pi _\mathrm{\Lambda }`$ is quasi-equivalent to a GNS-representation of a quasifree state if and only if $`\overline{\mathrm{\Lambda }}`$ is a Hilbert-Schmidt operator. If $`\mathrm{\Lambda }_{CCR}^{(n)}`$ is bounded, then $`\pi _\mathrm{\Lambda }`$ is quasi-equivalent to a GNS-representation of a quasifree state if and only if $`\overline{\mathrm{\Lambda }P_1\mathrm{\Lambda }}`$ is a Hilbert-Schmidt operator.
###### Demonstration Proof
Under the given hypothesis, one has $`P_1\mathrm{\Lambda }0`$. Thus, with Theorem 3.2.2, the stated assertions follow at once from the previous Theorem. $`\overline{)}`$
The results in this paper can be used to provide an alternative proof to a well-known criterion for the unitary equivalence of two pure quasifree states . Recall that pure quasifree states are Fock.
###### Theorem 5.12
Let $`\pi =\pi _J`$ and $`\pi ^{}=\pi _J^{}`$ be two pure quasifree states on $`𝒜(H_0)`$, where $`H_0`$ is a dense subspace of $`H`$. Then they are unitarily equivalent if and only if $`JJ^{}`$ is Hilbert-Schmidt with respect to $`s`$, the scalar product on $`H`$ associated with $`\pi _J`$ (equivalently, Hilbert-Schmidt with respect to $`s^{}`$).
###### Demonstration Proof
Using notation already established in Chapter IV, it follows easily from Lemma 5.4 that if two Fock states $`\omega _F=\omega _J`$ and $`\omega _F^{}=\omega _J^{}`$ on $`𝒜(H_0)`$ are unitarily equivalent then the associated scalar products $`s`$ and $`s^{}`$ are equivalent - this entails $`H=H^{}`$. By Proposition 4.4, Theorems 4.5 and 5.7, the operator (employing the notation of Proposition 4.2)
$$(|T|\text{1I})(|T|+\text{1I})=|T|^2\text{1I}=JK\text{1I}=JJ^{}\text{1I},$$
resp. $`J^{}J=J(JJ^{}\text{1I})`$, is Hilbert-Schmidt with respect to $`s`$, resp. $`s^{}`$. $`\overline{)}`$
As in Theorem 5.12, one can use the results of Chapters IV and V, particularly Theorems 4.5 and 5.10 and Proposition 4.4, to give an alternative proof of the criteria characterizing the quasi-equivalence of quasifree states (see for increasingly general results), but we shall not give the details here.
In are given necessary and sufficient conditions so that two irreducible quadratic representations (see ) are unitarily equivalent. One could generalize that result in order to obtain necessary and sufficient conditions so that an irreducible quadratic representation and an irreducible representation of finite degree are unitarily equivalent. But, as these conditions are not particularly transparent, we shall not present them here.
Finally, we mention that in this paper and in the previous ones, , the choice of the complex structure, and thus the choice of the Fock representation, has been held fixed. It is therefore of interest to point out that in the unitary equivalence of two quadratic representations constructed from different Fock representations is characterized. Those arguments can also be generalized to the case of the representations of finite degree discussed in this chapter, but once again the conditions which emerge are not particularly edifying.
## Acknowledgements
The research program of which this paper is a continuation would not have taken place without Georg Reents’ initial impetus and participation. In addition, early versions of some of the results of this paper appeared in , a Diplomarbeit carried out under Dr. Reents’ direction. We have also benefitted from an exchange with Prof. Paul Robinson concerning our differing approaches to these problems.
## REFERENCES
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# PECULIARITIES OF INCLINED FORCE AFFECTION ON HOMOGENEOUS ONE-DIMENSIONAL ELASTIC LUMPED-PARAMETERS LINE
## 1 Introduction
In papers have been considered longitudinal vibrations in one-dimensional elastic lumped-parameters line. However in such line vibrations of more general form are possible, if we attribute the idea of one-dimensionality only to general shape of line, not to degree of freedom of vibrating elements of line.
This paper investigates the peculiarities of exact analytic solutions for this type of problems.
## 2 Vibrations in semi-infinite elastic line under inclined external force affecting on start of line
Regarding stated above, supposing two degrees of vibration freedom for elastic line elements, the model can be presented in form shown in fig.1. In case when amplitude of vibration is small (linear vibration), this model can be described by two sets of equations - with respect to $`x`$\- and $`y`$-projections of external force affection correspondingly:
$$\{\begin{array}{c}m\frac{d^2\mathrm{\Delta }_1}{dt^2}=F\left(t\right)\mathrm{cos}\alpha +s\left(\mathrm{\Delta }_2\mathrm{\Delta }_1\right)\\ m\frac{d^2\mathrm{\Delta }_2}{dt^2}=s\left(\mathrm{\Delta }_3+\mathrm{\Delta }_12\mathrm{\Delta }_2\right)\\ \mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}.\\ m\frac{d^2\mathrm{\Delta }_n}{dt^2}=s\left(\mathrm{\Delta }_{n+1}+\mathrm{\Delta }_{n1}2\mathrm{\Delta }_n\right)\\ \mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\end{array}$$
(1)
$$\{\begin{array}{c}m\frac{d^2y_1}{dt^2}=F\left(t\right)\mathrm{sin}\alpha +s\left(y_2y_1\right)\\ m\frac{d^2\mathrm{\Delta }y}{dt^2}=s\left(y_3+y_12y_2\right)\\ \mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}.\\ m\frac{d^2y_n}{dt^2}=s\left(y_{n+1}+y_{n1}2y_n\right)\\ \mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\end{array}$$
(2)
where $`\alpha `$ is angle of external force inclination to the axe of line.
Each of these sets of equations is similar to ones having been investigated in . consequently, we can write at once exact analytic solutions for every of them.
For $`x`$-component of vibration: periodic regime $`\left(\beta <1\right)`$
$$\mathrm{\Delta }_n=j\frac{F_0\mathrm{cos}\alpha }{\omega \sqrt{sm}}e^{j\left[\omega t\left(2n1\right)\tau \right]}$$
(3)
aperiodic one $`\left(\beta >1\right)`$
$$\mathrm{\Delta }_n=\left(1\right)^n\frac{F_0\mathrm{cos}\alpha }{\omega \sqrt{sm}}\gamma ^{2n1}e^{j\omega t}$$
(4)
critical one $`\left(\beta =1\right)`$
$$\mathrm{\Delta }_n=\left(1\right)^n\frac{F_0\mathrm{cos}\alpha }{2s}e^{j\omega t}$$
(5)
For $`y`$-component of vibration we obtain correspondingly: periodic regime $`\left(\beta <1\right)`$
$$y_n=j\frac{F_0\mathrm{sin}\alpha }{\omega \sqrt{sm}}e^{j\left[\omega t\left(2n1\right)\tau \right]}$$
(6)
aperiodic one $`\left(\beta >1\right)`$
$$y_n=\left(1\right)^n\frac{F_0\mathrm{sin}\alpha }{\omega \sqrt{sm}}\gamma ^{2n1}e^{j\omega t}$$
(7)
and critical one $`\left(\beta =1\right)`$
$$y_n=\left(1\right)^n\frac{F_0\mathrm{sin}\alpha }{2s}e^{j\omega t}$$
(8)
As a result of superposition, there forms an inclined wave propagating in positive direction of axe $`x`$; this is confirmed by diagram of vibration shown in fig.2.
It is characteristic that inclined pattern of vibration remains as under unforced vibrations in lumped-parameters line as under limiting process to distributed-parameters line.
Really, basing on results presented in , solution, e.g. for unforced vibration, has the following form: for $`x`$-component of vibration
$$\mathrm{\Delta }_n=\frac{X_k\mathrm{cos}\left(2i1\right)\tau }{\mathrm{cos}\left(2k1\right)\tau }e^{j\omega t}$$
(9)
for $`y`$-component
$$y_n=\frac{Y_k\mathrm{cos}\left(2i1\right)\tau }{\mathrm{cos}\left(2k1\right)\tau }e^{j\omega t}$$
(10)
where $`X_k`$ and $`Y_k`$ are $`x`$\- and $`y`$-components of vibration amplitude of $`k`$th element which parameters are specified; $`k`$ is number of element which vibration is specified.
In case of limiting process to distributed-parameters line, we can present
$$\rho =\frac{m}{a};s=\frac{T}{a};n=\frac{x_0}{a}$$
where $`\rho `$ is density; $`T`$ is tension in line; $`x_0`$ is distance from start of line to the point of rest of investigated element of line; $`m`$ is mass of element of line. With it solutions (3)$`÷`$(8) transform to the set of equations
$$\{\begin{array}{c}x=j\frac{F_0\mathrm{cos}\alpha }{\omega \sqrt{T\rho }}e^{j\omega \left(tx_0\sqrt{\rho /T}\right)}+x_0\\ y=j\frac{F_0\mathrm{cos}\alpha }{\omega \sqrt{T\rho }}e^{j\omega \left(tx_0\sqrt{\rho /T}\right)}\end{array}$$
(11)
Obtained set of equations describes paramertically the inclined wave propagating along the axe $`x`$, as shown in fig.3.
We can see of carried out investigation that inclined vibrations arise far from always as a consequence of nonlinear processes in elastic system, as it was supposed before. Inclined waves can arise quite naturally under affection of force inclined to the direction of wave process propagation. And this conclusion can be quite simply extended to a most wide spectrum of vibration process.
## 3 Elements-of-line motion trajectory
Paying attention to a separate element motion trajectory, we can see easy, this trajectory has form of ellipse circumscribed around the point of rest of element. And inclination of wave forms, at the cost of shift phase of motion along element-to-element elliptic trajectories. Presented structure of vibration is well-known in physics, particularly in wave processes in unbounded volumes of liquid. ”In a wave, motion of liquid is non-stationary. So trajectories of separate particles are far from coinciding with lines of current in time. They have absolutely other form. Under small amplitudes they are circumferences in a great approximation. We find these circular trajectories as on surface as in depth of liquid. Only in the most upper layers the diameters of circular ways are the most large” \[3, pp.300-301\].
Indeed, vibration processes in space have their peculiarities. None the less, it is characteristic that basic regularities are run down already in one-dimensional model. It also follows of obtained solutions that ellipsoidal pattern of vibrations of elements remains as in critical as in aperiodic regimes. Consequently, in last case in the line forms compound wave fast-decaying along the line, and this is one more peculiarity that exact analytic solutions demonstrate.
## 4 On new class of functions being the solution of wave equation
Foregoing generalizations can be extended also to solution of wave equation in the whole.
It is known that differential equation of hyperbolic type
$$\frac{^2y(x,t)}{x^2}=\frac{k^2}{\omega ^2}\frac{^2x(x,t)}{t^2}=0$$
(12)
has general solution \[5, p.300\]
$$\mathrm{\Phi }(x,t)=\mathrm{\Phi }_1\left(kx\omega t\right)+\mathrm{\Phi }_2\left(kx+\omega t\right)$$
(13)
where $`c=\omega /k`$ is velocity of wave propagation, i.e. in the form of two explicit functions with respect to $`\left(xct\right)`$ and $`\left(x+ct\right)`$ correspondingly. Till now this solution was considered the only and complete, due to theorem of uniqueness of solution of differential equation. None the less, there exists one more class of functions being the solution of differential equation (12) but not taken into account by solution (13). We can present general form of this class of functions in the form
$$y(x,t)=\mathrm{\Phi }_1\left(kx\omega t+\psi _1\left(y\right)\right)+\mathrm{\Phi }_2\left(kx+\omega t+\psi _2\left(y\right)\right)$$
(14)
where $`\psi _1\left(y\right)`$and $`\psi _2\left(y\right)`$are some twice-differentiable functions. In other words, given solution (14) belongs to the class of implicit functions whose regularities of behavior and technique of differentiation and integration essentially differ from such for explicit functions. Important that, while for explicit functions definite systematization of differential equations has been created and for definite class of these equations the regularities and schemes to obtain solutions have been defined, for implicit functions all these developments are absent. Naturally, for today correspondence of expression (14) to differential equation (12) can be checked only by the most simple way - by straight substitution (14) into (12).
For it, on the grounds of known laws of implicit functions differentiation, find first and second particular derivatives of expression (14) with respect to $`x`$ and $`t`$. To simplify calculation, consider a half of right part of expression (14)
$$y(x,t)=\mathrm{\Phi }_1\left(kx\omega t+\psi _1\left(y\right)\right)=\mathrm{\Phi }_1\left(A\right)$$
(15)
where $`A=kx\omega t+\psi _1\left(y\right)`$. First derivatives have the form
$$\frac{y}{x}=\frac{k\frac{d\mathrm{\Phi }_1}{dA}}{1\frac{d\psi _1}{dy}\frac{d\mathrm{\Phi }_1}{dA}};\frac{y}{t}=\frac{\omega \frac{d\mathrm{\Phi }_1}{dA}}{1\frac{d\psi _1}{dy}\frac{d\mathrm{\Phi }_1}{dA}}$$
(16)
Second derivatives after transformation and substitution of expressions (16) take form
$`{\displaystyle \frac{^2y}{x^2}}`$ $`=`$ $`{\displaystyle \frac{k^2\left[\frac{d^2\mathrm{\Phi }_1}{dA^2}+\frac{d^2\psi _1}{dy^2}\left(\frac{d\mathrm{\Phi }_1}{dA}\right)^3\right]}{\left(1\frac{d^2\psi _1}{dy^2}\frac{d\mathrm{\Phi }_1}{dA}\right)^3}}`$ (17)
$`{\displaystyle \frac{^2y}{t^2}}`$ $`=`$ $`{\displaystyle \frac{\omega ^2\left[\frac{d^2\mathrm{\Phi }_1}{dA^2}+\frac{d^2\psi _1}{dy^2}\left(\frac{d\mathrm{\Phi }_1}{dA}\right)^3\right]}{\left(1\frac{d^2\psi _1}{dy^2}\frac{d\mathrm{\Phi }_1}{dA}\right)^3}}`$
Substituting (17) into (12), we obtain required. Similarly we can prove the correspondence of second part of expression (14) to equation (12).
Thus, solution (17) defines a whole class of implicit functions satisfying the linear wave equation. And presence of new class of functions being the solution of equation (12) does not violate a least the theorem of uniqueness of solution of differential equation, because under definite conditions
$$\psi _1\left(y\right)0;\psi _2\left(y\right)0$$
expression (14) degenerates into (13). hereby it is proved that solution known before is a particular case of more general class of functions.
The found class of implicit functions defines nonlinear wave; its degree of deformation depends on form of functions $`\psi _1\left(y\right)`$ and $`\psi _2\left(y\right)`$. For example, in particular case of expression (14) (see fig.4)
$$y=c\mathrm{sin}\left(kx\omega t+y\mathrm{cot}\alpha \right)$$
(18)
solution of wave equation (12) describes progressive wave propagating along axe $`x`$ and inclined by angle $`\alpha `$, what completely corresponds with inclined vibration in one-dimensional line investigated above.
## 5 Conclusions
1. Analyzing solutions for semi-infinite model on free end of which affects a force inclined to the axe, we have found that as a result of this affection, in the line propagate inclined waves described by implicit function.
2. Solutions of wave equation in the form of implicit functions are generalizing for known solution being superposition of running waves.
3. We have ascertained that under affection of inclined force the elements of line follow elliptic trajectories.
## 6 Symbols
$`F\left(t\right)`$ is external force affecting on the line; $`F_0`$ is amplitude of external force; $`T`$ is tension in the line; $`X_i`$, $`Y_i`$ are $`x`$\- and $`y`$-components of vibration amplitude of $`k`$th element whose parameters of vibration are given; $`a`$ is distance between the elements of line; $`f`$ is frequency of vibration in the line; $`i`$, $`k`$, $`n`$ are indexes; $`m`$ is mass of element of line; $`k`$ is wave number; $`s`$ is stiffness coefficient of line; $`t`$ is time parameter; $`x_0`$ is distance from start of line to the point of rest of last element of line; $`y_i`$ is displacement of $`i`$th element of elastic line in vertical plane; $`\mathrm{\Delta }_i`$ is instantaneous longitudinal displacement of $`k`$th element of line; $`\alpha `$ is angle of affecting force inclination to the axe of elastic line; $`\beta `$, $`\gamma _+`$, $`\gamma _{}`$, $`\tau `$ are parameters of line; $`\omega `$ is circular frequency of affecting force.
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# Field induced long-range-ordering in an 𝑆=1 quasi-one-dimensional Heisenberg antiferromagnet
## I Introduction
Many lower-dimensional antiferromagnets show a singlet ground state and an energy gap in their excitation spectrum originating from quantum many body effects. The first indication of the quantum energy gap came out in 1983. Haldane conjectured that the excitation spectrum of one-dimensional (1D) Heisenberg antiferromagnets (HAFs) with integer spin quantum number ($`S`$) has an energy gap between the ground and first excited states, while the corresponding system with half-odd integer $`S`$ has no energy gap. This conjecture has been tested both theoretically and experimentally. It is now generally accepted that the quantum energy gap (Haldane gap) does exist in 1D HAFs with integer $`S`$.
In real materials, there always exists an interaction between spin chains. We call this class of materials a quasi-one-dimensional (Q1D) magnet. Usually Q1D antiferromagnets exhibit a long-range ordering (LRO) at finite temperature due to the inter-chain coupling ($`J^{}`$). The compound Ni(C<sub>2</sub>H<sub>8</sub>N<sub>2</sub>)<sub>2</sub>NO<sub>2</sub>(ClO<sub>4</sub>), abbreviated NENP, is a typical example of $`S`$=1 Q1D HAFs and does not show any indication of LRO down to 300 $`\mu `$K, whereas the $`S`$=1 Q1D HAF CsNiCl<sub>3</sub> shows LRO at $``$4.7 K. The different behavior between the two compounds comes from the difference in the strength of $`J^{}`$. How robust is the Haldane phase in a Q1D HAF against perturbation? Sakai and Takahashi studied theoretically the ground state properties of an $`S`$=1 Q1D HAF with a single-ion anisotropy of the form $`DS_z^2`$. They showed that the Haldane disordered phase exists for $`zJ^{}/J`$$``$0.05 in a rather wide range of $`D`$ values, where $`z`$ is the number of neighboring chains and $`J`$ is the intra-chain exchange interaction.
As was demonstrated experimentally, strong magnetic fields destroy the Haldane gap and the system recovers magnetism. Then, we expect a magnetic ordering to occur in an $`S`$=1 Q1D HAF under high fields and at low temperatures. The heat capacity measurement made on NENP in magnetic fields did not reveal any indication of LRO. One of the reasons for this is that an external magnetic field applied along the chain axis of NENP induces a staggered field on the Ni<sup>2+</sup> sites because of the presence of two crystallographically inequivalent sites for Ni<sup>2+</sup>. This staggered field causes a small energy gap near the transition field from the disordered to the magnetized states and thus prevents the occurrence of LRO at low temperatures.
We have reported the observation of a field-induced LRO in the $`S`$=1 Q1D HAF Ni(C<sub>5</sub>H<sub>14</sub>N<sub>2</sub>)<sub>2</sub>N<sub>3</sub>(ClO<sub>4</sub>), abbreviated NDMAZ, and in Ni(C<sub>5</sub>H<sub>14</sub>N<sub>2</sub>)<sub>2</sub>N<sub>3</sub>(PF<sub>6</sub>), abbreviated NDMAP. Since there is only one site for Ni<sup>2+</sup> in NDMAZ and MAMAP, these compounds are ideal for studying the field-induced LRO. From a heat capacity ($`C_p`$) measurement on a single crystal sample of NDMAZ, we have observed an anomaly at about 0.6 K and 12 T which indicated that a magnetic ordering occurred there. Because of limitations in our calorimeter, we were unable to follow how the position of the anomaly in $`C_p`$ changes with temperature ($`T`$) and magnetic field ($`H`$). We then tried to synthesize the $`S`$=1 Q1D HAF, NDMAP, with a weaker $`J`$ than NDMAZ so that LRO is expected to be induced at a lower field. We have measured $`T`$ and $`H`$ dependence of $`C_p`$ in a single crystal sample of NDMAP and constructed the $`H`$-$`T`$ phase diagram. We found an anisotropy in the phase boundary curve separating the paramagnetic and LRO phases, and we presented a qualitative interpretation for this anisotropy. Electron spin resonance (ESR) measurements made on a single crystal of NDMAP gave further evidence for the existence of the field-induced LRO phase and of the anisotropy in the phase boundary curve.
In this paper we report detailed results of heat capacity and magnetization measurements made on NDMAP. We also report a theoretical analysis of the $`H`$-$`T`$ phase diagram.
The format of this paper is as follows. In Sec.II we present the relevant background information and details. The experimental results are given in Sec.III. In Sec.IV, a theoretical consideration is given on the $`H`$-$`T`$ phase diagram. The last section (Sec.V) is devoted to discussion and conclusions.
## II Preliminary Details
The compound Ni(C<sub>5</sub>H<sub>14</sub>N<sub>2</sub>)<sub>2</sub>N<sub>3</sub>(PF<sub>6</sub>) (NDMAP) has the orthorhombic structure with the space group $`Pnmn`$ shown in Fig. 1. The lattice parameters are, $`a`$= 18.046 $`\AA `$, $`b`$= 8.7050 $`\AA `$ and $`c`$= 6.139 $`\AA `$. The structure consists of Ni(C<sub>5</sub>H<sub>14</sub>N<sub>2</sub>)<sub>2</sub>N<sub>3</sub> chains along the $`c`$ axis. These chains are well separated from each other by PF<sub>6</sub> molecules. All the Ni<sup>2+</sup> sites in a chain are equivalent, $`i.e.`$, only one site exists for Ni<sup>2+</sup>.
From the analysis of the magnetic susceptibility data, the following values are obtained; $`J/k_B`$= 30.0 K, $`D/J`$= 0.3, $`g_{}`$= 2.10 and $`g_{}`$= 2.17, where $`g_{}`$ and $`g_{}`$ are the $`g`$ values parallel and perpendicular to the chain $`c`$ axis, respectively. Neutron inelastic scattering measurements were done on single crystals of deuterated NDMAP at $`T`$= 1.4 K . From the analysis of the data, the values of exchange and anisotropy parameters are determined to be, $`J`$= 2.28 meV (= 26.5 K), $`J_x^{}`$= 3.5$`\times `$10<sup>-4</sup> meV (= 4.1$`\times `$10<sup>-3</sup> K), $`J_y^{}`$= 1.8$`\times `$10<sup>-3</sup> meV (= 2.1$`\times `$10<sup>-2</sup> K) and $`D`$= 0.70 meV (= 8.1 K), where $`J_x^{}`$ and $`J_y^{}`$ are the inter-chain exchange interactions along the $`a`$ and $`b`$ axes, respectively.
The single crystals of NDMAP used in this study were grown from an aqueous solution of NaN<sub>3</sub>, Ni(NO<sub>3</sub>)<sub>2</sub>$``$6H<sub>2</sub>O and 1,3-diamino-2,2-dimethylpropane. After filtration, KPF<sub>6</sub> was added to the solution. Well shaped blue single crystals up to 5mm$`\times `$5mm$`\times `$20mm were obtained after several weeks. Fully deuterated single crystals of NDMAP (NDMAP-d<sub>28</sub>) were grown from a D<sub>2</sub>O solution of the same ingredients except that a deuterated 1,3-diamino-2,2-dimethylpropane was used.
The single crystals thus obtained were checked by a four circle X-ray diffractometer. We confirmed that the lattice parameters of our crystal are almost identical with those reported before. We measured $`C_p`$ of NDMAP and NDMAP-d<sub>28</sub> and found that the $`H`$-$`T`$ phase diagrams of the two systems are essentially the same. This means that the magnetic parameters in NDMAP and NDMAP-d<sub>28</sub> are not different.
Heat capacity measurements were performed with a MagLab<sup>HC</sup> microcalorimeter (Oxford Instruments, UK). The temperature and magnetic field ranges accessible with this calorimeter are, 0.45 K$``$$`T`$$``$ 200 K and 0$``$$`H`$$``$12 T. Magnetization measurements were done with a MagLab<sup>VSM</sup> vibrating-sample-magnetometer (Oxford Instruments, UK). The temperature and magnetic field ranges available with this magnetometer are, 1.5 K$``$$`T`$$``$300 K and 0$``$$`H`$$``$12 T.
## III Experimental Results
### A Heat Capacity
Figure 2 shows the temperature dependence of the heat capacity of NDMAP, including the contribution of the lattice measured in zero field. The data are well expressed by the following equation,
$$C_p=aT+bT^3$$
(1)
in the temperature range between 2 and 5 K with $`a`$= 0.109 and $`b`$= 0.00653. We use hereafter this $`bT^3`$ term to subtract the contribution of lattice heat capacity as has been done by many authors.
We show in Figs. 3(a)-(c) the temperature dependence of magnetic heat capacity ($`C_m`$) of NDMAP, after subtracting the lattice heat capacity, in magnetic fields applied parallel to the $`a`$, $`b`$, and $`c`$ axes, respectively. In all field directions, we see an anomaly in $`C_m`$ at finite fields above a critical value. This anomaly signals that a magnetic ordering occurs there.
One of the advantages of our calorimeter is that a field dependent $`C_p`$ can be measured under a constant temperature. Strictly speaking, we need to change temperature to measure $`C_p`$. However, the temperature increment necessary for the measurement is 0.5 - 1 % of the temperature we set so that temperature change during the measurement may be considered as small. Figures 4(a)-(c) show such ”field scan” data measured at several temperatures for the magnetic field directions parallel to the $`a`$, $`b`$, and $`c`$ axes, respectively. In addition to the sharp peak at the field denoted by $`H_{\mathrm{LRO}}`$, a broad feature is seen around the field named as $`H_\mathrm{c}`$. Here, $`H_{\mathrm{LRO}}`$ means the position of $`H`$ at which the field-induced LRO occurs for a given $`T`$.
Combining all the information obtained from the heat capacity measurements, both of the ”temperature scan” and ”field scan” procedures, we present the $`H`$-$`T`$ phase diagram of NDMAP in Fig. 5. In addition to the anisotropic phase boundary (curve (A)) separating the disordered and LRO phases, we have another boundary (curve (B)) separating the Haldane and the disordered phases which is also anisotropic. The two curves (A) and (B) seem to merge at a finite $`H`$ when extrapolated to $`T`$= 0 K, for respective field directions.
### B Magnetization
Figures 6(a)-(c) show the temperature dependence of susceptibility (magnetization divided by applied magnetic field, $`M/H`$) in NDMAP measured in magnetic fields applied parallel to the $`a`$, $`b`$, and $`c`$ axes, respectively. The behavior of the susceptibility at $`H`$= 1 T is reminiscent of that taken at a much lower field ($`H`$= 0.01 T); a broad peak around 35 K and a steep decrease in susceptibility with decreasing temperature below about 20 K. On increasing $`H`$, $`M/H`$ does not extrapolate to zero with $`T`$$``$0. This behavior of $`M/H`$ is similar to the one observed in an $`S`$=$`\frac{1}{2}`$ 1D HAF and indicates that a transition from the gapped to a gapless phase occurs at a higher field. On increasing $`H`$ further, $`M/H`$ shows a minimum and an up turn at low temperatures. The insets of Figs. 6(a)-(c) show the low temperature part of the data. We see in the inset of Figs. 6(a) and (b) that $`M/H`$ becomes almost temperature independent below a temperature whose value is field dependent. The temperature independent susceptibility reminds us of the perpendicular susceptibility ($`\chi _{}`$) of an anisotropic antiferromagnet below the Néel temperature ($`T_\mathrm{N}`$). We show below that we are actually observing $`\chi _{}`$ in this compound. Because the sign of the single-ion anisotropy term ($`DS_z^2`$) is positive in this compound, spins in the ordered phase are expected to lie in a plane perpendicular to the $`c`$ axis (the quantization axis of the $`D`$ term is taken parallel to the $`c`$ axis). The anisotropy in the $`c`$ plane of the form $`E(S_x^2S_y^2)`$ is very small. Therefore, when $`H`$ is applied along the $`a`$ or $`b`$ axes, spins point perpendicularly to $`H`$ in the $`c`$ plane keeping an antiferromagnetic arrangement, thus giving $`\chi _{}`$.
We plot in Fig. 5 the transition points obtained from the $`M/H`$ data shown in Fig. 6(b). These points are defined as the temperatures where $`d(M/H)/dT`$ shows a minimum for a given $`H`$. We see that the transition points determined from the heat capacity and magnetization measurements agree well with each other.
## IV Theoretical Consideration on the Phase Diagram
In this section, we try to reproduce the $`H`$-$`T`$ phase diagram observed in NDMAP, exhibiting an interesting behavior of $`T_\mathrm{N}`$ as a function of $`H`$: At low temperatures, there occurs LRO only above a certain critical field and $`T_\mathrm{N}`$ shows an increase with increasing fields, the rate of which depends on the direction of $`H`$. We focus our attention especially on the physics behind the phase diagram. To this end, we adopt the mean-field approximation for the interchain interaction, which is known to work quite well except for fields in the vicinity of the critical field. In the following, we use the energy unit $`J`$=1, so that $`t`$, $`𝐡`$, $`j^{}`$ and $`d`$ are renormalized quantities of $`T`$, $`𝐇`$, $`J^{}`$ and $`D`$, respectively. According to the mean-field theory, the renormalized Néel temperature $`t_\mathrm{N}(𝐡)`$ as a function of the renormalized field $`𝐡`$ is given by the solution satisfying the following equation:
$$1/j^{}z=\chi _{\mathrm{st}}(t_\mathrm{N}(𝐡);𝐡),$$
(2)
where $`\chi _{\mathrm{st}}(t;𝐡)`$ is the staggered susceptibility for the one-dimensional magnetic system.
Then, we calculate the staggered susceptibility of the magnetic chain by means of the quantum transfer matrix method combined with the finite-temperature density matrix renormalization group. As was mentioned in the previous section, Ni<sup>2+</sup> spins in NDMAP has an easy-plane anisotropy, and thus the magnetic chain is well described by the following Hamiltonian:
$$=J[\underset{n}{}\{𝐒_n𝐒_{n+1}+d(S_n^z)^2\}\underset{n}{}g\mu _\mathrm{B}𝐡𝐒_n],$$
(3)
where $`𝐒_n`$ represents the spin-$`1`$ operator at the $`n`$th site. Neglecting the small anisotropy in the $`c`$ plane, we consider the following two cases for the field $`𝐡`$ applied (i) along the $`z`$($`c`$)-axis (perpendicular to the easy plane) and (ii) along the $`y`$-axis (in the easy plane).
It is noted here that the situation is quite different between the two cases, since the former field reserves the axial symmetry around the $`c`$-axis while the latter breaks it. The nature of quantum fluctuations and hence the staggered susceptibility depends crucially on the symmetry of the system. In the following, we consider the two cases separately.
(i) The field applied perpendicularly to the easy plane ($`𝐡=(0,0,h)`$)
As was mentioned, below the critical field $`h_\mathrm{c}`$ the nonmagnetic Haldane phase called a quantum disordered phase is the ground state of the system, while above it the so called Tomonaga-Luttinger liquid state becomes the ground state. Although the former has an excitation gap to the triplet state, the latter has a gapless excitation spectrum and hence is critical. This criticality of the Tomonaga-Luttinger liquid is characterized by the critical exponent $`\eta `$ defined by the divergence of the staggered susceptibility at zero temperature:
$$\chi _{\mathrm{st}}(t;h)=At^{(2\eta )},$$
(4)
where $`A`$ is a constant, which scarcely depends on $`𝐡`$ in our calculation. Note that in the classical system $`\eta =0`$. We may safely use Eq. (4) at low temperatures well below the temperature at which the susceptibility is maximum ($``$35 K). In Fig. 7, $`\eta `$’s, estimated from the numerical calculations for the temperature region $`0.1<t<1.0`$, are shown by the solid circles, each of which has an error bar of $`\pm 0.05`$. The data have been shifted in fields so that the theoretical critical field coincides with the one observed. The solid curve represents the phenomenological relation between $`\eta `$ and $`h`$: $`\eta =0.3\mathrm{exp}\{\beta (hh_\mathrm{c})\}+0.2`$ with $`\beta =0.5`$ and $`h_\mathrm{c}=0.2`$.
Now, we reproduce the phase diagram, using Eqs. (2) and (4) with $`J=26.5`$K. In Fig. 8 we show the theoretical result by the solid curve and the experimental data by the solid circles. Here, we adjusted the theoretical curve to reproduce the experimental point $`T_\mathrm{N}=0.92`$ K at 11 T. From Figs. 7 and 8, we see that the increase of $`T_\mathrm{N}`$ with $`H`$ is a consequence of the decrease in $`\eta `$ with $`H`$. Remembering that $`\eta `$ measures the degree of quantum fluctuations, we can say that the quantum fluctuation out of the easy plane is reduced by $`H`$ so that the Néel state becomes more stable. In contrast to this quantum system, the field dependence of the constant $`A`$ is an only source of the field dependence of $`T_\mathrm{N}`$ in the classical system, being very mild. We show in Fig. 7 $`\eta `$ estimated from the phase diagram (Fig. 5) by the open circles, which follows also the phenomenological relation but $`\beta =1.0`$. Considering the large error bars in $`\eta `$ of our estimations, we do not think the discrepancy in $`\beta `$ so seriously. Observation of the field dependence of $`\eta `$ by other methods is desired.
In stronger fields, spins cant in the field direction and hence $`A`$ decreases seriously. Thus, the phase boundary curve closes at the upper critical field, where the magnetic moment saturates. From the phase diagram, we estimate the interchain coupling $`j^{}z`$ to be $`1.2\times 10^3`$, using the value $`A=2.0`$. This value is favorably compared with that obtained from the neutron inelastic scattering experiments ( 2$`j_x^{}`$+2$`j_y^{}`$=$`1.9\times 10^3)`$.
(ii) The field applied in the easy plane ($`𝐡=(0,h,0)`$)
In this case, above $`h_\mathrm{c}`$, the ground state has the Néel order and a gap opens again in the excitation spectrum because of the symmetry breaking field. This situation is quite different from the former case. The field dependence of $`T_\mathrm{N}`$ reminds us of the soliton scenario in the classical system. Remembering the form of the staggered susceptibility in the classical system, we postulate the following in this case:
$$\chi _{\mathrm{st}}(t;h)=(B/t^2)\mathrm{exp}\{\alpha (hh_\mathrm{c})/t\},$$
(5)
where $`B`$ is a constant. This form is confirmed by our numerical calculations, in the temperature region $`0.1<t<1.0`$, with the coefficient $`\alpha `$, a little less than unity. Although $`\alpha `$ may represent quantum effects in the formation energy of soliton, we assume, for simplicity, the classical value 1 for it.
Now, we again reproduce the phase diagram for this case using Eqs. (2) and (5). The theoretical curves are adjusted as before using the experimental point $`T_\mathrm{N}=2.2`$ K at 12 T and the critical field value $`H_{\mathrm{CF}}=5.7`$ T for $`𝐡`$ parallel to the $`a`$-axis, and $`T_\mathrm{N}=2.7`$ K at 12 T and $`H_{\mathrm{CF}}=5.4`$ T for $`𝐡`$ parallel to the $`b`$-axis, respectively. The different values of $`H_{\mathrm{CF}}`$ are due to an in-plane anisotropy, being neglected in this paper. The phase boundary curves are shown by the solid curve for $`𝐡`$ parallel to the $`a`$-axis and by the dotted curve for $`𝐡`$ parallel to the $`b`$-axis with the corresponding experimental points, respectively, by the solid and the open circles in Fig. 9. The agreement between theory and experiment is satisfactory. We mention that the soliton scenario still works in our quantum system: The symmetry breaking field yields the uniaxial symmetry and hence the soliton is a dominant source for the fluctuations in this system. Since the soliton formation energy, i.e. the gap, increases with $`hh_\mathrm{c}`$ and hence the staggered susceptibility increases exponentially at low temperatures, $`T_\mathrm{N}`$ shows a rapid increase with $`H`$, the rate of which is marked contrast with the former case. Although the soliton scenario is effective in our case, more sophisticated study is required to establish further a quantum analogue of the soliton in classical spin chains.
The phase boundary curve closes also at the upper critical field as in the former case. We estimate the interchain coupling $`j^{}z`$ to be $`1.4\times 10^3`$, using the value $`B=0.6`$. The value $`j^{}`$ estimated again agrees with the value observed.
## V Discussion and Conclusions
In Sec. III, we have presented detailed results of heat capacity and magnetization measurements on NDMAP from which we have constructed the $`H`$-$`T`$ phase diagram shown in Fig. 5. We have been successful in explaining theoretically the phase boundary curve separating the paramagnetic and LRO phases in Sec. IV using the values of intra-chain exchange interaction and anisotropy constants determined from the neutron inelastic scattering measurements. The inter-chain exchange interactions estimated theoretically are close to those obtained from the neutron experiment.
We discuss the lower field boundary separating the Haldane and paramagnetic phases (curve (B) in Fig. 5). We argue below that the anomaly in $`C_m`$ observed along this curve is due to the field dependence of the first excited triplet. We analyzed the low temperature part of the magnetic heat capacity data using a two-level system model, with a singlet ground state and the lowest state of the excited triplet with an energy difference ($`\mathrm{\Delta }_{}(H)`$), which gives a Schottky type anomaly. We show in Fig. 10, $`\mathrm{\Delta }_{}(H_x)`$, $`\mathrm{\Delta }_{}(H_y)`$ and $`\mathrm{\Delta }_{}(H_z)`$ thus obtained for the field directions parallel to the $`a(x)`$, $`b(y)`$ and $`c(z)`$ axes, respectively. The solid curves in Fig. 10 represent the theoretical energy level as a function of $`H`$. Here, we used the value determined from the neutron inelastic scattering experiment for the energy gap in respective field directions at $`H`$$``$0. We see in this figure that the agreement between theory and experiment is satisfactory.
Finally, we discuss the temperature dependence of $`M/H`$ (Figs. 6(a)-(c)). We calculated $`M/H`$ as a function of $`T`$ using the quantum transfer matrix method with a density-matrix renormalization group technique. We compare theory and experiment in Fig. 11(a) for selected values of $`H`$ parallel to the $`a`$ axis. Here, we used $`J/k_\mathrm{B}`$=26.5 K obtained from the neutron inelastic scattering study and $`g`$=2.14 determined from the ESR measurement. We see a good agreement between theory and experiment without any adjustable parameters. We have obtained the value $`J/k_\mathrm{B}`$=30.0 K from a fitting of the theory with the susceptibility data at high temperature range above about 40 K. Because the lattice parameters change with temperature, it is not surprizing if the exchange interaction constant determined at low temperatures is different from that at high temperatures. Figure 11(b) shows the case when $`H`$ is applied along the $`c`$ axis. Since no ESR data are available along this direction, we assumed the value 2.05 for $`g`$. The agreement between theory and experiment is not as good as in Fig. 11(a). Further study is necessary to clarify this point.
In conclusion, we have reported detailed results of heat capacity and magnetization measurements on a single crystal sample of the $`S`$=1 Q1D HAF, NDMAP. From these results, we constructed the $`H`$-$`T`$ phase diagram which exhibits the quantum disordered Haldane, field-induced LRO and thermally disordered paramagnetic phases. The phase boundary curve separating the paramagnetic and LRO phases is anisotropic; the increase of $`T_\mathrm{N}`$ with $`H`$ along the $`a`$ and $`b`$ axes is more rapid than that along the $`c`$ axis. We calculated $`T_\mathrm{N}`$ as a function of $`H`$ by taking into account the inter chain coupling as the form of a mean-field. We first evaluated numerically the staggered susceptibility of the $`S`$=1 1D antiferromagnet for $`H`$ applied perpendicularly to the easy plane, which shows, at low temperatures, a typical divergence of the Tomonaga-Luttinger liquid with the critical exponent $`\eta `$. Then, we got a satisfactory agreement with the experimental results using the exchange and anisotropy constants obtained from the neutron scattering experiment. It is interesting to note that the transition temperature $`T_\mathrm{N}`$ is governed by the critical exponent $`\eta `$ of the Tomonaga-Luttinger liquid. On the other hand, for $`H`$ applied in the easy plane, we invoked the soliton scenario and got again a satisfactory agreement with experiment.
## Acknowledgements
This work was partially supported by a Grant-in-Aid for Scientific Research from the Japanese Ministry of Education, Science, Sports and Culture. Z. H. was supported by the Research Fellowships of the Japan Society for the Promotion of Science for Young Scientists. The computation in this work has been done using the facilities of the Supercomputer Center, ISSP, University of Tokyo.
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# Theoretical Calculations of Atomic Data for Spectroscopy
## 1. Introduction
Spectroscopic data for atoms and ions fall into the category of fundamental scientific research and this by itself may be good justification for experimental and theoretical work. The data are also essential to the analysis of spectra of laboratory and astronomical plasmas. These applications have been as important motivation as, if not greater than, the fundamental physics For several decades, theoretical atomic physicists have devoted great effort to the development of methods and computation of vast amounts of data, particularly that needed for the study of optical spectra. On the other hand, the progress of the atomic data sets for X-ray spectroscopy has been much slower, due in part to the scarcity of high quality X-ray spectra to motivate the work. At present, the new generation of observatories (like Chandra and XMM) and the incoming facilities with the Constellation-X program are starting to provide high resolution, high signal-to-noise spectra for which the need of improved atomic data is not just granted but urgent. But the production of high quality atomic data for X-ray spectroscopy is far from simple or routine. The high energy conditions of the emitting plasmas demand the study of large numbers of energy levels as well as inner-shell and Auger processes, Thus, current theoretical methods and computational tools encounter severe difficulties in producing high accuracy data.
The present review of theoretical atomic physics is aimed at astronomers and spectroscopists in general who need to understand the limitations of atomic data. In that sense, the present paper tries to give a general, but relatively simple, overview of the most commonly used approximations and methods for the production of spectroscopic data. The review pays particular attention to the the calculation of atomic structure and the spectral processes of photoionization and electron impact excitation. Other processes like collisional ionization and proton impact excitation, although relevant, are not discussed in much detail due to lack of space. For discussions on electron-ion recombination see papers by Nahar and Savin in this volume. Also, the present review emphasizes inner-shell processes that are most relevant to X-ray spectroscopy.
From the point of view of the kinds of atomic data needed for the analysis of spectra there are two general atomic physics problems to be solved:
| 1) Atomic structure | \- Energy Levels |
| --- | --- |
| | \- Radiative transition rates |
| | \- Autoionization rates |
| 2) Scattering | \- ion + photon |
| | \- ion + charged particle |
The atomic structure problem is concerned with the computation of energy levels and rates of spontaneous transition rates among levels of the same ion (transition rates or $`A`$-values), and from autoionizing levels of one ion to levels of the next ionized species (autoionization rates). The scattering problem, relevant to the formation of spectra of warm and hot plasmas, has to do with all different processes that can occur after a collision of an ion with a photon or with a charged particle. The most important physical processes that occur after ion scattering are illustrated in Figure 1.
## 2. Atomic Structure Calculations
From the point of view of quantum mechanics the physical representation of ions and spectral processes comes down to the fundamental problem of solving the Schrödinger equation
$$H\mathrm{\Psi }_i=E_i\mathrm{\Psi }_i.$$
(1)
Here, $`i`$ represents the set of quantum numbers necessary to describe the system, $`\mathrm{\Psi }_i`$ are the wavefunctions of the system, $`E_i`$ are the eigenvalues, and $`H`$ is the Hamiltonian which, for a non-relativistic system, can be written as
$$H=\underset{i=1}{\overset{N}{}}\frac{p_i}{2m_e}\underset{i=1}{\overset{N}{}}\frac{Ze^2}{r_i}+\underset{ij}{}\frac{e^2}{|r_ir_j|},$$
(2)
where the first term is the sum of the kinetic energies of all electrons, the second term is the potential energy due to Coulomb attraction of all electrons by the nucleus, the third term is the energy due to electrostatic repulsion between the electrons, and $`N`$ is the total number of electrons in the system. The previous equation neglects the spin of the electron, which can be considered separately in non-relativistic systems.
The presence of the two-electron operators $`1/|r_ir_j|`$ makes it impossible to obtain exact solutions to the Schrödinger equation for the $`N`$-electron system. On the other hand, the two-electron operator should never be neglected as it is often comparable in magnitude to the Coulomb attraction term or one-electron operator. One can show
$$\frac{twoelectronterm}{oneelectronterm}\xi \frac{1}{2}\frac{N(N1)}{ZN}.$$
(3)
For a neutral ion ($`Z=N`$), $`1/4\xi 1/2`$.
Approximate methods for solving Equation (1) replace the two-electron terms by approximate one-electron potentials to give an effective Hamiltonian of the form
$$H^{eff}=\underset{i}{\overset{N}{}}H_i^{eff}=\underset{i=1}{\overset{N}{}}\left[\frac{1}{2}\frac{p_i^2}{m_e}+\frac{Ze^2}{r_i}V_i^{eff}(r_i)\right].$$
(4)
The exact eigenfunctions of $`H^{eff}`$ can now be constructed from products of single-electron eigenfunctions of $`H_i^{eff}`$. Since $`H`$ commutes with angular momentum and spin operators ($`L^2,L_z,S^2,S_z`$) the eigenfunctions of $`H^{eff}`$ are required to be eigenfunctions of these operators as well.
For highly ionized ions ($`ZN`$), the inter-electronic repulsions are only small perturbations relative to the much stronger nuclear central potential. Then, it is a good approximation to choose $`V_i^{eff}`$ as a central field potential. If this can be achieved the orbital wavefunction may be written explicitly as products of spherical harmonics, spin functions, and radial functions. The advantage of using central field orbitals cannot be overestimated. The effect is to reduce the computational problem from solving $`N`$ coupled integro-partial-differential equations (one for each orbital) in four variables to $`N`$ coupled integro-differential equations in a single radial variable.
In general $`V^{eff}`$ does not need to be of central field type, but the advantages of a central potential are so great that essentially all methods for atomic structure calculations use spherically averaged potentials even for neutral and open-shell systems. There are several techniques regularly used in the atomic structure calculations. The most important are: use of model potentials, methods based on the Hartree-Fock theory, semiempirical methods, perturbation methods, and the R-matrix method in the close coupling formalism. These methods, with the exception of the R-matrix, are discussed below followed by a discussion of two particularly important physical effects, i.e. configuration interaction and relativistic effects. The R-matrix method will be described in Section 3.5.
### 2.1. Model potentials in atomic structure
Alkali atoms and ions are good systems for testing model potentials since, by having a single valence electron, their analysis leads to a single equation for the one-electron wavefunction. For these ions one would expect the potential for large $`r`$ to approach
$$V^{eff}(r)=\frac{(ZN_c)}{r},$$
(5)
where $`N_c`$ is the number of core electrons. The solution of the Schrödinger equation based on this potential is discussed by Bates and Damgaard (1949). In spite of the simplicity of this potential it gives accurate oscillator strengths in many cases, but it is not always reliable (Bromander et al. 1978).
A large variety of model potentials have been constructed from the assumptions that they behave as shown in Equation (5) for $`r>R`$, for some suitable $`R`$, while for $`r`$ approaching 0 the potential could be infinite, equal to a finite boundary value, or even zero. For a review of model potentials in atomic structure see Hibbert (1982).
A generally applicable potential is the modified Thomas-Fermi potential by Eissner and Nussbaumer (1969). This potential is implemented in the widely used computer program SUPERSTRUCTURE (Eissner, Jones, and Nussbaumer 1974). This program uses a Thomas-Fermi-Dirac type of central potential to generate one-electron orbitals.
This potential differs from that in Equation (5) not only in the added sophistication of the model, but also in that this one depends on the angular momentum of the valence electrons. SUPERSTRUCTURE is quite efficient and can provide relatively accurate results ($`1\%`$ for energy levels and $`10\%`$ for oscillator strengths). Recently the code has been extended to calculate autoionization and dielectronic recombiantion data (AUTOLSJ by Bely-Dubau 1982 and AUTOSTRUCTURE by Badnell 1985, unpublished) has extended the code to calculate autoionization rates and dielectronic recombination.
### 2.2. Atomic structure calculations base on the Hartree-Fock Formalism
In contrast to the local central-field potentials, the Hartree-Fock (HF) method entails the computation of the one-electron orbitals in the non-local potential (direct and exchange) generated from electronic orbitals in a self-consistent manner using the variational principle. More extensive discussion of the HF method can be found in Hartree (1957), Slater (1960), Froese Fischer (1977), Cohen and McEachran (1980).
Early HF calculations indicated that in going from the ground state of beryllium $`1s^22s^2{}_{}{}^{1}S`$ to the first excited states $`1s^22s2p^{1.3}P^o`$, the $`1s`$ orbital remains almost unchanged. This led Fock (1933) to realize that accurate wavefunctions could be found by varying only valence orbitals while keeping the orbitals for the core fixed. This approach is known as the frozen core (FC) approximation.
In spite of the success of the FC approximation, this is not always directly applicable, such as in the case of configurations containing equivalent electrons. The simplest example of such configurations is the ground state of helium $`1s^2{}_{}{}^{1}S`$, which has two equivalent electrons, both described by a single radial function $`P_{1s}(r)`$. Yet, if one identifies one of the electron as “core” and the other as “valence” they require different radial functions. Evidently, these two $`1s`$-orbitals cannot be made orthogonal. The use of nonorthogonal orbitals to solve the atomic structure problem leads to equations much more complex than otherwise. Calculations with such nonorthogonal orbitals have been carried out only for a few systems (e.g. Pratt 1956; Froese 1966; Jucys 1967).
The FC approximation yields goods results for simple systems, but for the mayority of cases it is necessary to include electronic configuration interactions like in the multiconfiguration Hartree-Fock (MCHF) or in the superposition of configurations (SOC) methods. For example, for boron-like cores a possible two-configuration core wave function is
$$\mathrm{\Psi }(^1S)=c_1\mathrm{\Psi }_1(1s^22s^2{}_{}{}^{1}S)+c_2\mathrm{\Psi }_2(1s^22p^2{}_{}{}^{1}S).$$
(6)
Clementi and Veillard (1965) showed that including $`\mathrm{\Phi }(1s^22p^2{}_{}{}^{1}S)`$ accounts for most of the correlation energy error for all values of $`Z`$.
The MCHF (e.g. Froese Fischer 1977) method computes orthonormal orbitals self-consistently in an iterative fashion for every choice of the $`\{c_i\}`$ coefficients. Then, these coefficients are varied and the radial wavefunctions are recomputed until sufficient degree of convergence is achieved. This makes the method potentially very accurate, but computationally lengthy.
Based on the MCHF technique, Hibbert (1975) developed the SOC method in which analytic radial functions depending on variational parameters are used. Core wave functions of the kind shown in Equation (7) are used and the $`\{c_i\}`$ coefficients are determined for a given choice of the radial wavefunction parameters. Then, the radial wavefunctions are changed and the coefficients determined again until the description of the atom converges satisfactorily. The SOC method is somewhat more efficient than the MCHF technique and is able to provide similarly accurate results. The main implementation of the SOC method is in the computer program CIV3 (Hibbert 1975).
Another approach used to improve calculations of atomic structure arises from the realization that a single valence electron may polarize the spherically symmetric core, including a dipole moment at the nucleus of the form $`\alpha /r^2`$, where $`\alpha `$ is the dipole polarizability of the core (Hartree 1957). This gives rise to an additional attractive field with a long-range $`r`$ dependence which is experienced by the valence electron, and at large distances may be represented by a “polarization potential”
$$V_{pol}=\alpha /2r^4.$$
(7)
This potential is singular at $`r=0`$, but it can be modified to avoid this singularity as (Biermann and Trefftz 1953)
$$V_{pol}=(\alpha /2r^4)[1exp(x^p)],x=r/\rho .$$
(8)
The use of a $`V_{pol}`$ in the orbital equations for valence electrons can be formally justified (Caves and Dalgarno 1972), and it has been shown (e.g. Cohen and McEachran 1980, and references therein) to provide a very satisfactory description of simple spectra. However, the method has the disadvantage that the choice of the parameters $`p`$ and $`\rho `$ in Equation (9) remains completely ad hoc.
### 2.3. Semiempirical methods
Semiempirical methods try to compute the atomic structure of ions by solving simplified forms of the HF equations. One of the first examples of this approach is the Hartree-Fock-Slater (HFS) scheme, which consists in replacing the two-electron non-local exchange terms of the Hamiltonian by the statistical potential function
$$V_{xs}(r)=\frac{3}{2}\left(\frac{24}{\pi }\rho \right)^{1/3},$$
(9)
where $`\rho `$ is the local electron density $`\rho (r)`$ of the atom.
Later, Cowan (1967) modified the HFS approach to write the total potential of the ion as
$$V^i(r)=2Z/r+V_e(r)+V_{xs}(r),$$
(10)
where $`V_e`$ is the potential energy for the density $`\rho \rho _i`$ of electrons other than $`i`$ and $`V_{xs}`$ is an exchange term of the form
$$V_{xs}=k_1f(r)\left[\frac{\rho ^{}}{\rho ^{}+k_2/(nl)}\right]\left(\frac{\rho ^{}}{\rho }\right)\left(\frac{24}{\pi }\rho \right)^{1/3},$$
(11)
where $`k_1`$ and $`k_2`$ are constants, $`\rho ^{}`$ is the electron density less the densities of the electron $`i`$ and of the electron with which it is paired, and
$$f(r)=\{\begin{array}{cc}1\hfill & rr_0\hfill \\ 1+k(1r/r_0)\hfill & r<r_0\hfill \end{array}$$
(12)
for suitable $`k`$ and $`r_0`$.
This definition for the atomic potential is mostly empirically motivated and requires preconceived wave functions with which to construct the electron density functions $`\rho `$ and $`\rho ^{}`$. Once these functions are created the atomic structure equations are solved in an iterative fashion. Then, the accuracy of the results is assessed upon the agreement between obtained energy levels and oscillator strengths with a given sample of experimental values. The method has the advantage of being quite efficient, but it requires a lot of care in the construction of the initial electron density distribution. Further, it is difficult to estimate the accuracy of any given calculation except by the observed agreement of a limited sample of data with experimental values.
The Cowan code has been widely used by many researchers like Kurucz (e.g. Kurucz 1988; Kurucz and Peytremann 1975), who computed millions of energy levels and oscillators strengths for most ions of astrophysical interest.
### 2.4. Perturbation treatments
Perhaps the best known perturbation treatment is the Z-expansion method which was first introduced by Hylleraas (1930). The main idea behind this method is that one can rewrite the $`N`$-electron atom Hamiltonian of Equation (2) in terms of units of length of $`Z`$ atomic units (au) and unit of energy of $`Z^2`$ au (thus different units for each $`N`$-electron ion) to obtain
$$H=\left[\frac{1}{2}\underset{i}{\overset{N}{}}\frac{p^2}{m_e}\underset{i}{\overset{N}{}}\frac{1}{r}\right]+Z^1\left[\underset{ii}{}\frac{1}{|r_ir_j|}\right].$$
(13)
Then, for highly ionized systems ($`ZN`$) one can divide this Hamiltonian as indicated by the brackets into a one-electron zero order Hamiltonian ($`H_0`$) and a two-electron perturbation ($`H_1`$) with the expansion parameter $`Z^1`$. Perturbation theory leads to the expansion for the energy and wavefunction of the form
$$E=E_0Z^2+E_1Z+E_2+E_3Z^1+\mathrm{}$$
(14)
$$\mathrm{\Psi }=\mathrm{\Psi }_0Z^{3/2}+\mathrm{\Psi }_1Z^{1/2}+\mathrm{\Psi }_2Z^{1/2}+\mathrm{}$$
(15)
Here, $`H_0`$ is the sum of hydrogenic Hamiltonians, so $`\mathrm{\Psi }_0`$ is a linear combination of Slater determinants of hydrogenic orbitals. Subsequent orders of the Hamiltonian and wavefunctions are found as in standard perturbation theory.
The Z-expansion method is conceptually simple, but computationally lengthy when trying to go beyond the first perturbation term. One advantage of the method is that a single calculation provides results for a whole isoelectronic sequence, but accuracy is normally restricted to highly ionized ions. Another problem occurs in cases of strong configuration interaction (see Section 2.5) where states are labeled by the single configuration with the largest mixing coefficient, thus zero order mixing does not represent the physical state being considered.
Several extensions of the Z-expansion method have been developed (see Cohen 1988, Crossley 1969, and references therein) and a lot of data for astrophysical applications has been published in recent years by Vainshtein, Safronova, and collaborators (e.g. Vainshtein and Safronova 1980; Safronova et al. 1998).
### 2.5. Additional considerations
#### Configuration interaction
Configuration interaction (CI) can be seen as a way to correct for Hartree’s single-electron orbital approximation that pictures each electron moving individually in the field of the nucleus screened by the other electrons. Each electron is described by a single wavefunction and the whole atom is described by the Slater determinant of these wavefunction (e.g. Slater 1960). This model admits the familiar configuration description of, for example, the ground state of Be $`1s^22s^2(^1S)`$. This kind of description is often not very accurate. Much better wavefunctions may be obtained from linear combinations of single configuration wavefunctions of the same total angular momentum and spin symmetry (Condon and Shortley 1935). This approach is named CI.
Going back to the case of the ground state of Be, a good CI description may be
$$\mathrm{\Psi }(^1S)=\alpha \mathrm{\Psi }(1s^22s^2{}_{}{}^{1}S)+\beta \mathrm{\Psi }(1s^22p^2{}_{}{}^{1}S),$$
(16)
where $`\alpha `$ and $`\beta `$ are the so-called mixing coefficients.
Such CI representation has important effects on atomic quantities like oscillator strengths. For example, in the transition probability of the resonant transition $`1s^22s^2(^1S)1s^22s2p(^1P^o)`$. In the single configuration model the transition probability is determined by the matrix element $`<2s|r|2p>`$. On the other hand, in the two-configuration description of the ground $`{}_{}{}^{1}S`$ state, the matrix element becomes $`(\alpha \pm \beta )<2s|r|2p>`$. Depending on the relative signs of the mixing coefficients and the $`\pm `$ sign, the transition may be strengthened or weakened, but which way the value will go is generally impossible to know without a full solution of the atomic structure problem.
CI is very important in the representation of the majority of atomic systems and is a standard capability of codes like SUPERSTRUCTURE, CIV3, Cowan’s code, and Froese Fischer’s MCHF code.
#### Relativistic Effects
So far we have neglected relativistic effects in atomic structure calculations and we have considered oscillator strengths for dipole allowed transitions under LS coupling. In order to consider forbidden transitions it is necessary to use $`jj`$ (or at least intermediate) coupling, and this requires some treatment of relativistic effects. For heavy atoms and ions relativistic effects are important even for allowed transitions. There are two ways of treating relativistic effects: by the addition of Breit-Pauli operators to non-relativistic equations, or by the fully relativistic Dirac formalism. The Breit-Pauli operators are seven and each one accounts for a specific physical effect (Bethe and Salpeter 1957). These operators are: (1) the mass operator which gives the correction due to the relativistic variation of mass with velocity; (2) the Darwin term which is characteristic of the Dirac theory but has no obvious physical interpretation; (3) the spin-orbit coupling term which arises from the interaction of spin and orbital magnetic moments of each electron; (4) the spin-other-orbit term which is due to the interaction between the spin of one electron with the orbital magnetic moment of another electron; (5) and (6) are spin-spin coupling terms that describe the interactions between the spin magnetic moments of pairs of electrons; (7) the orbit-orbit coupling term which accounts for the interactions between the orbital magnetic moments of pairs of electrons.
The expectation value of the orbit-orbit coupling term cannot be calculated accurately in the one-electron approximation. Moreover, experience suggests that this term is rather small for all ions of interest to the present discussion (Ufford and Callen 1958; Cohen and McEachran 1980). The mass and Darwin terms contribute to the overall energy shift, but do not break LS coupling. The spin-spin, orbit-orbit, and spin-other-orbit operators serve to split energy terms into fine structure levels.
Most current codes for atomic structure calculations such as SUPERSTRUCTURE and CIV3 use the Breit-Pauli approximation to account for relativistic effects. A more precise treatment of relativistic effects using the Dirac formalism has been implemented in the computer package GRASP (Dyall et al. 1989).
## 3. Scattering Calculations
In considering the problems of photoionization and excitation or ionization by electron impact it is customary to think of the residual ion as the “target” or “core” with $`N`$ electrons and the incoming/outgoing (electron impact/photoionization) as “free electron”. Thus, the wavefunctions of the $`(N+1)`$-electron system can be expanded in terms of products of wavefunctions of the core ($`\varphi _i`$) and those of the electron ($`\theta _i`$), i.e.
$$\mathrm{\Psi }=\underset{i}{}\varphi _i(x_1,\mathrm{},x_N)\theta _i(x_{N+1}).$$
(17)
Substitution of (18) in Equation (2) yields the system of coupled equations
$$(^2+k_i^2)\theta _i(x)+\underset{i^{}}{}V_{ii^{}}\theta _i(x)=0$$
(18)
where
$$V_{ii^{}}(x)=\frac{2Z}{r_1}\delta _{ii^{}}+\varphi ^{}(x_1,\mathrm{},x_N)\times \underset{n=1}{\overset{N}{}}\frac{2}{r_{N+1,n}}\varphi (x_1,\mathrm{},x_N)dx_1\mathrm{}dx_N$$
(19)
and $`k_i^2`$ is defined by
$$E=E_i(N)+k_i^2.$$
(20)
It is important to notice that these equations have to be solved for every value of the energy ($`k^2`$) and of the total angular momentum ($`L`$) of the $`(N+1)`$-electron system. This last condition motivates the so-called partial waves expansion in which all states of definite angular momentum of the free electron are considered separately. Furthermore, the calculated collision strength is divided into partial waves as
$$\mathrm{\Omega }=\underset{l=0}{\overset{\mathrm{}}{}}\mathrm{\Omega }_l$$
(21)
There are several methods to solve the scattering problem and calculate the cross sections necessary for practical applications. Below is a discussion of the most widely used methods ordered according to their level of sophistication, starting with the simplest method. This discussion is followed by a study of several physical effects that affect the cross sections. The techniques are: the central field approximation, used in photoionization calculations; the Gaunt factor and the Coulomb-Born approximation, used for electron impact excitation; and the Distorted Wave and the R-matrix methods, used for both photoionization and electron impact excitation and ionization, as well as atomic structure.
### 3.1. The central field approximation
This is the simplest approximation used to solve the multi-electron problem. As discussed above, in the context of atomic structure calculations, it is assumed that inter-electronic couplings are small and can be treated as small perturbations with respect to an effective central field potential that includes the attraction by the nucleus less the screening by the core electrons. A further substantial simplification of the problem results by treating subshells of equivalent electrons (having the same values of $`n`$ and $`l`$, but different values of $`m`$ and $`s`$) by means of a single radial function. This has the effect of reducing the number of coupled radial equations to only one for each subshell.
Central field type calculations of photoionization cross sections have been carried out by Reilman and Manson (1979), who adopted Hartree-Slater wave functions (Herman and Skillman 1963), and Verner et al. (1993), using the Dirac-Slater potential (Slater 1960; Band et al. 1979). These cross sections are reasonable accurate at high photon energies, but often give poor results near the threshold, particularly where CI is strong. For example, in the photoionization of Fe I the central field cross section from the threshold to the $`3d`$ subshell is underestimated by more than three orders of magnitude, as seen in Figure 2 (Bautista and Pradhan 1995). Similar errors of up to orders of magnitude are found in Fe II through Fe V (Bautista and Pradhan 1998).
Other problems with central field cross sections are that they neglect resonances and, because only degenerate subshells of equivalent electrons rather than individual energy levels are considered, the utility of these cross sections in detailed spectral modeling is limited.
### 3.2. The Gaunt factor
Burgess, Seaton, and Van Regemorter (Burgess 1961; Seaton 1962, Van Regemorter 1962) suggested an approximate formula to obtain near threshold collision strengths for optically allowed transitions. The formula is based on the Kramers approximation and the Gaunt factor, $`g`$, and is usually accurate within a factor of a few. Later, the formula was modified to replace $`g`$, which is a varying function with energy, by an empirical parameter $`\overline{g}0.2`$. The formula is
$$\mathrm{\Omega }(i,i^{})=\frac{8\pi }{3\sqrt{(}3)}S(i,i^{})\overline{g},$$
(22)
where $`S(i,i^{})`$ is the line strength, which relates to the absorption oscillator strength as $`g_i^{}f_{i^{}i}=(2/3)(E_i^{}E_i)S(i,i^{})`$
At high energies $`\overline{g}`$ increases logarithmically and the collision strength takes the form (Seaton 1962)
$$\mathrm{\Omega }(i.i^{})=\frac{4}{3}S(i,i^{})\mathrm{ln}[4k^2/(r_0\mathrm{\Delta }E)^2],$$
(23)
where $`\mathrm{\Delta }E=E_i^{}E_i`$, and $`r_0`$ is the distance of closest approach given by
$$\frac{2z}{r_0}\frac{l(l+1)}{r_o^2}+k^2=0.$$
(24)
### 3.3. The Coulomb-Born approximation
In the Coulomb-Born (CB) approximation $`V_{ii^{}}`$ is replaced by its asymptotic form at large $`r`$, $`2z/r`$, where $`z=ZN+1`$. In doing this the solutions to the scattering problem are known analytically. The CB approximation is best for highly charged ions and transitions not affected by channel coupling and large partial waves (e.g. Seaton 1975; Van Regemorter 1960).
### 3.4. The Distorted Wave approximation
The Distorted Wave (DW) method assumes the coupling between different target states to be weak ($`V_{ii^{}}=0`$, for $`ii^{}`$ in Equation 19). Then, the system of coupled equations is reduced to
$$\{(^2+k_i^2)+V_{ii}\}\theta _i=0.$$
(25)
Some of the most refined forms of the method were developed by Saraph, Seaton, and Shemming (1969) and Eissner and Seaton (1972) including exchange, which is important in many cases, like in intercombination transitions. In addition, correlation functions, $`\chi _i`$, which are wavefunctions of the $`(N+1)`$-electron system, were introduced to modify the wavefunctions as
$$\mathrm{\Psi }=\underset{i}{}\varphi _i(x_1,\mathrm{},x_N)\theta _i(x_{N+1})+\underset{j=1}{}c_j\chi _j(x_1,\mathrm{},x_{N+1}),$$
(26)
where the coefficients $`c_j`$ are determined variationally. This form of expansion of the wavefunction is known as the close-coupling (CC) expansion (Seaton 1953).
The DW method usually gives accurate collision strengths in the absence of resonances for systems more than a few times ionized (see paper by Bhatia in this volume). Hershkowitz and Seaton (1973) showed that the DW method could also provide information about bound states of the $`(N+1)`$-electron system, and hence obtain resonance structures. However, for detailed calculations of cross sections including resonances it is usually preferred to use the R-matrix method.
### 3.5. The R-matrix method
The R-matrix method is the most sophisticated of the techniques discussed here and is also the most accurate. The method takes into account nearly all of the physical effects that contribute to cross sections for astrophysical applications and is applicable to all kinds of ions, from neutral to highly ionized species. Of course, with the increased complexity of the calculations the R-matrix method can be computationally very intensive.
The R-matrix theory starts by dividing the configuration space by a sphere of radius $`a`$ centered on the target nucleus. In the internal region, $`r<a`$, where $`r`$ is the relative coordinate of the free electron, electron exchange and correlation between the scattered electron and the target are important. Thus a CC expansion like that in Equation (27) is adopted for the system. In the external region, $`r>a`$, electron exchange between the free electron and the target can be neglected if $`a`$ is large enough to contain the charge distribution of the target. Then the scattered electron moves in the long-range multipole potentials of the target ion. The internal and external regions are linked by the R-matrix on the boundary, $`r=a`$.
Following the theory from Burke et al. (1971), and Burke and Robb (1975), the total wavefunction $`\mathrm{\Psi }`$ in the inner region for any energy $`E`$ can be written in terms of the basis states set $`\{\psi _k\}`$ as
$$\mathrm{\Psi }=\underset{k}{}A_{Ek}\psi _k$$
(27)
where the $`\psi _k`$ functions are energy independent and are expanded as shown in Equation (27), and the energy dependence is carried through the $`A_{Ek}`$ coefficients.
Then, the $`R`$ matrix is defined as
$$R_{ij}(E)=\frac{1}{2a}\underset{k}{}\frac{w_{ik}(a)w_{jk}(a)}{E_kE},$$
(28)
where
$$\frac{1}{r}w_{ik}(r)=<\chi _i|\psi _k>$$
(29)
This $`R`$ matrix is the basic solution of the electron-scattering problem as it allows one to determine the atomic structure of the $`(N+1)`$ system, the collision strengths, and photoionization cross sections. The $`R`$ matrix, the amplitude of the wavefunctions at the $`r=a`$ boundary, and the poles $`E_k`$ are obtained from the eigenvalues and eigenvectors of the Hamiltonian matrix. Further, quantitative results for physical parameters of interest are obtained by matching the solutions with those in the outer region which are known analytically.
One important point to note is that in the R-matrix method the solution inner region is obtained only once, then cross sections for any number of energy points are readily available. Furthermore, although the basic computations in this method are lengthly, it turns out to be a very efficient technique as it provides results for large numbers of points which allow complex resonance structures in the cross sections to be delineated.
The R-matrix method has been implemented in the RMATRX package of codes (Berrington, Eissner, and Norrington 1995) which has been widely used by several groups like the Opacity Project, that carried out extensive atomic structure and photoionization calculations of astrophysically important ions, and the Iron Project, and the RmaX Project (see papers by Pradhan and Berrington in this volume).
### 3.6. Important physical effects
There are several specific physical effects that should be taken into account when calculating or evaluating cross sections for photoionization or collisional excitation/ionization. These are: target CI, resonances, convergence of the wavefunction expansion, convergence of partial waves expansion, relativistic effects, and radiation damping.
#### Configuration interaction
The first requirement in any scattering calculation is a good representation of the target, i.e. accurate wavefunctions for the target ion. Such representation usually requires the inclusion of CI in the atomic structure model, which as we saw in Section 2.5.1 affects the calculated energy levels and oscillator strengths of the ion. Furthermore, comparisons between calculated and experimental energies and oscillator strengths of the target are important indicators of the quality of the target representation and the overall accuracy of the obtained cross sections.
#### Resonances
Resonances are an important part of collisional excitation/ionization and photoionization cross sections. Physically, resonances occur when the incoming particle (electron or photon) with just the right kinetic energy gets trapped into an autoionizing state<sup>1</sup><sup>1</sup>1Autoionizing states are compound states of the (electron+ion) system located above the ionization potential. These states result from the excitation of two or more electrons of the system or by excitation of inner-shell electrons. of the $`(N+1)`$-electron system. Then, as the electron remains trapped for a time before autoionization occurs, the time delay yields a phase shift in the wavefunction that manifests itself in sharp peaks or throughs in the cross sections (resonances).
Resonances appear as Rydberg series converging onto the various excitation thresholds of the target. In complex ions with many levels close in energy the series of resonances often overlap and interference effects can occur which makes the resonances inseparable from the background cross section. An example of this is the photoionization cross section of the ground state of Fe I (Figure 2). The R-matrix cross section (Bautista and Pradhan 1995) given by the solid line is over three orders of magnitude greater than the central field results without resonances. Similar discrepancies with respect to central field photoionization cross section are found for all iron ions up to Fe V (Bautista and Pradhan 1998). Like in the case of photoionization near the inner $`3d`$ subshell of Fe I, resonances are a general phenomenon near inner-shell ionization thresholds. Figure 3 shows the photoionization cross section of Fe XV from the first ionization threshold to just above the K-shell threshold (Bautista 2000a). Complex resonances structures are seen converging onto both L and K-shell thresholds. The anhancement of the cross section due to resonances near the L-shell threshold was first pointed out by Hanque and Pradhan (1999) These resonance can change the appearance of the thresholds from sharp edges to more complex and continuous transitions. The resonances also enhance the photoionization rates. For example, assuming a power law continuum of the kind $`FE^\alpha `$ for $`E`$ from 500 to 700 Ry, the resonances near the K-shell threshold of Fe XV enhance the ionization rate of this ion by factors of 1.7 for $`\alpha =1`$, 1.8 for $`\alpha =2`$, and 2.9 for $`\alpha =3`$.
Resonances in the collision strengths are known to enhance the excitation rates by up to several factors in the case of valence electron excitation. In the case of inner-shell excitation, Bautista (2000b) carried out R-matrix calculations for Fe XVI and compared the results with earlier DW results. In the region free of resonances the agreement is generally very good, but huge numbers of near threshold narrow resonances are found from the R-matrix calculation (see Figure 4). These resonances enhance the effective collision strengths by up to three order of magnitude, as shown in Table 1.
#### Convergence of the partial waves expansion
It is an standard approach to expand the collision strengths in partial waves from 0 to infinity from every possible value of the angular momentum of the free electron. Clearly, infinite expansion terms cannot be computed in practice and one takes only the lowest dominant partial waves and extrapolate through infinity. The convergence of the partial waves expansion is usually rapid for forbidden transitions, but it can be quite slow for allowed transitions. Also, the number of partial waves needed for convergence increases with increasing energy of the free electron. Thus, the convergence of the partial waves expansion becomes a difficult practical problem, and a possible source of error in the collision strengths for highly ionized systems for which very high collision energies need to be considered (e.g. Burke and Seaton 1986; Chidichimo 1988, 1989; Eissner et al. 1999)
#### Convergence of the close coupling expansion
The general form of the CC expansion for the radial wavefunction $`\mathrm{\Psi }`$ in terms of a $`N`$-electron target basis $`\chi _i`$ and the scattering electron function $`\theta _i`$ is
$$\mathrm{\Psi }=\underset{i}{\overset{\mathrm{}}{}}\chi _i\theta _i+_ϵ\chi _ϵ\theta _ϵ.$$
(30)
In the DW and the standard R-matrix approach, however, the integral term that accounts for the target continuum is neglected or replaced by a discrete sum over bound correlation functions (Equation 27) and the sum over the infinite number of target states is truncated to a small number of strongly coupled states. This approximation is usually good when considering valence electron excitations among the lowest energy levels of the ion. However, for excitations to highly excited levels and inner-shell excitations the convergence of the CC expansion must be looked at in great detail. For example, Sawey and Berrington (1993) showed that in order to obtain accurate collision strengths for the $`n=4`$ levels of He I one must include target states in the CC expansion up to at least $`n=5`$. Otherwise, if the CC expansion is too small the obtained collision strengths tend to be severely overestimated.
The use of pseudo-orbitals to accelerate the convergence of CC expansions was first introduced by Burke, Gallagher, and Geltman (1969). Recant developments are the convergent close-coupling (CCC) and R-matrix with pseudo-states (RMPS) methods (see the review by Gorczyca et al. in this volume).
#### Relativistic effects
If relativistic corrections are small, one can carry out the scattering calculation in LS coupling and then perform an algebraic transformation of the $`K`$ (reactance) matrices to fine structure (Saraph 1972, Luo and Pradhan 1990).
As the relativistic effects become significant one must allow for the fine structure splitting of the target by using the so-called term-coupling-coefficients (TCCs) to diagonalize the Hamiltonian with the inclusion of relativistic corrections.
For more complete and proper treatament of relativistic effects it is necessary to use either the Breit-Pauli or the Dirac formulations. The Breit-Pauli operators have been implemented into the RMATRX package of codes by the Iron Project (Hummer et al. 1993). One limitation with relativistic calculations is that by splitting the structure of the core into fine structure the size of the computation often exceeds the capacity of modern vectorized supercomputers. Perhaps, this problem may be solved in the near future with the use of massively parallel computers.
#### Radiation damping
Radiation damping of resonances can be an important effect for highly charged ions. When the radiative de-excitation rates of autoionizing states are comparable to the autoionization rates the resonances associated with these states become damped. In the case of photoionization damping of resonances occurs because as the target is photo-excited to an autoionizing state radiative de-excitation competes with autoionization. In collisional excitation damping of resonances indicates a competition between dielectronic recombination and autoionization.
Radiation damping in electron-ion scattering were studied by Pradhan (1981, 1983a, 1983b) using the branching ratios between autoionization and radiative rates and more recently by Zhang and Pradhan (1995) employing the detailed Bell and Seaton (1985) theory of dielectronic recombination. Radiation damping was found to be important in reducing the strength of resonances in boron-like Fe XXII and helium-like Fe XXV. As expected, radiation damping increases with the principal quantum number $`n`$ within a give Rydberg series of resonances. This is because autoionization rates typically decrease as $`n^3`$ while radiative rates remain approximately constant.
Robicheaux et al. (1995) and Gorczyca and Badnell (1996) have studied radiation damping in photoionization cross sections. However, in practical applications of photoionization and and scattering the overall radiation damping effect may be not so great in most cases. This is because while the higher $`n`$ resonances can be highly damped, the lower resonances in the series, which are usually dominant, are not affected significantly, except in highly charged ions (Pradhan and Zhang 1997; Robicheaux 1998).
## 4. Conclusions
A number of methods for atomic structure and electron-ion scattering calculations have been developed, which are able to provide atomic data for spectroscopic applications. These methods vary in complexity and accuracy of the results in such a way that one must often decide between promptly available and highly accurate data to fulfill current spectroscopic needs. In the case of X-ray spectroscopy most of the atomic data currently available have been produced using some of the simplest techniques described here, while improved data such as that from the IRON Project (see review by Pradhan) is slowly been calculated. In some cases, new more elaborate calculations will help in refining the spectral models, while in other cases the new atomic data could vastly change the previous results leading to dramatic changes in the spectral models. Therefore, it is of prime importance to researchers trying to analyze and/or model spectra to have at least some basic understanding of the quality of the data bing used.
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# Enhanced quantized current driven by surface acoustic waves
## ACKNOWLEDGMENTS
We would like to thank T. Weimann for fruitful discussions, H. Marx for providing the 2DEGs and the Deutsche Forschungsgemeinschaft (DFG Ah79/2-1) for financial support.
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# Large 𝑁 Quantum Time Evolution Beyond Leading Order
## I Introduction
The time evolution of quantum systems away from equilibrium is of interest in many applications including, but certainly not limited to, phase transition dynamics, inflationary reheating, and heavy ion collisions. Large $`N`$ expansions have provided a widely used technique for studying equilibrium properties in statistical physics and field theory , and it is natural to apply a similar strategy for studying non-equilibrium problems. The large $`N`$ limit (as typically formulated) is actually a special type of classical limit . Suitable observables behave classically and the quantum dynamics reduces to classical dynamics on an appropriate phase space.
Considerable work has been done examining the dynamics of far from equilibrium states in a variety of applications using leading large-$`N`$ time-evolution . A major virtue of large $`N`$ techniques (compared to alternative wholly uncontrolled approximation schemes) is that one should be able to improve the approximation by systematically including sub-leading effects suppressed by powers of $`1/N`$. For a variety of equilibrium problems (such as critical phenomena), this approach can work quite well .
For initial value problems, in which one would like to choose a non-equilibrium initial state and then examine the subsequent time evolution, traditional formulations of large-$`N`$ expansions using graphical or functional integral techniques are very awkward. A major difficulty with these approaches is that they generate integral equations which are non-local in time when sub-leading $`1/N`$ corrections are retained. For practical (numerical) applications, one would vastly prefer a formulation in which locality in time is always preserved.
In this paper, we describe a formulation of large $`N`$ (or semi-classical) dynamics which leads to a coupled hierarchy of time-local evolution equations for equal time correlation functions. Our approach directly exploits the appropriate group structure underlying the construction of suitable coherent states and the existence of the classical limit . We specifically focus on the time evolution of initial states chosen to equal one of these coherent states. We will give explicit next-to-leading order (NLO), and next-to-next-to-leading order (NNLO), expressions for the required evolution equations. Somewhat related hierarchies of evolution equations have been discussed in several recent papers . Because of our exploitation of the underlying group structure, the formulation we derive is more efficient, in the sense that it requires integration of fewer coupled equations at a given order in $`1/N`$.
A major question which we discuss, but do not fully resolve, is the propagation of errors induced by truncating the exact (infinite) hierarchy at a given order in $`1/N`$. It is known that the $`N\mathrm{}`$ limit is not uniform in time. For example, in typical large $`N`$ field theories the characteristic time scales for scattering or thermalization are known to scale as $`N`$ to some positive power.See, for example, the end of section III of Ref. . For a fixed time interval $`t`$, results obtained by integrating evolution equations truncated at, for example, next-to-leading order, will have only order $`1/N^2`$ errors. For sufficiently large $`N`$, and fixed $`t`$, including successively higher orders in the $`1/N`$ hierarchy will yield more accurate results. But for fixed $`N`$ and some given truncation of the $`1/N`$ hierarchy, it should be expected that the truncation error will grow with increasing time and eventually become order unity. A key question is how this “breakdown” time scales with $`N`$ and the order of the truncation. One might hope that a next-to-leading order approximation would be useful \[that is, have at most $`𝒪(1/N)`$ global errors\] for times of order $`N`$, while a next-to-next-to-leading order scheme would be useful out to times of order $`N^2`$, etc. But it is quite conceivable that errors in an order-$`k`$ truncation will grow with time like $`(t^\alpha /N)^k`$ for some positive $`\alpha `$, which would imply that all truncations break down after a time of order $`N^{1/\alpha }`$. This behavior, which we consider likely, may well depend on the specific theory and choice of initial state. Available numerical work, such as , sheds little light on this issue. We discuss several examples where it is possible to argue that quantum “decoherence” produces exactly this type of limit on the range of validity of large $`N`$ truncations.
The paper is arranged as follows. The general framework which allows us to treat many theories with a classical limit in a uniform fashion is outlined in section II. This material is largely taken from Ref. . Section III describes the particular class of operators we will consider, and examines the structure of their coherent state equal time correlators. Section IV presents the resulting time evolution equations and discusses error propagation. These general results are applied to the examples of point particle quantum mechanics, and a general $`N`$-component vector model, in section V. For point particle quantum mechanics, we argue that the decoherence time generically scales as $`\mathrm{}^{1/2}`$, while for vector models it should scale as $`N^{1/2}`$. A brief concluding discussion follows.
## II Coherence Group and Coherent States
The following slightly abstract framework is applicable to typical large $`N`$ limits (including $`O(N)`$ or $`U(N)`$ invariant vector models, matrix models, and non-Abelian gauge theories), as well as the $`\mathrm{}0`$ limit of ordinary quantum mechanics .
Consider a quantum theory depending on some parameter $`\chi `$ (such as $`\mathrm{}`$ or $`1/N`$). The Hilbert space (which may depend on $`\chi `$) will be denoted $`_\chi `$. The quantum dynamics is governed by a Hamiltonian which we will write as $`(\mathrm{}/\chi )\widehat{H}_\chi `$. This rescaling of the Hamiltonian will prove to be convenient, and makes the Heisenberg equations of motion take the form
$$\frac{d}{dt}\widehat{A}=\frac{i}{\chi }[\widehat{H}_\chi ,\widehat{A}].$$
(1)
The following assumptions are a set of sufficient conditions implying that the $`\chi 0`$ limit is a classical limit.
Assume there is a Lie group $`G`$ (called the coherence group) which, for every value of $`\chi `$, has a unitary representation on $`_\chi `$, $`G_\chi =\{D_\chi (u):uG\}`$. The states generated by applying elements of the coherence group to some (normalized) base state $`|0_\chi _\chi `$,
$$|u_\chi D_\chi (u)|0_\chi ,uG,$$
(2)
are called coherent states. The coherence group acts on these states in a natural way, $`D_\chi (u^{})|u_\chi =|u^{}u_\chi `$.
We assume that the coherence group $`G_\chi `$ acts irreducibly on the corresponding Hilbert space $`_\chi `$. In other words, no operator (except the identity) commutes with all elements of the coherence group. This condition automatically implies that the set of coherent states form an over-complete basis for the Hilbert space $`_\chi `$. It also implies that any operator acting on $`_\chi `$ may be represented as a linear combination of elements of the coherence group.
For any operator $`\widehat{A}`$ acting in $`_\chi `$, we define its symbol $`A_\chi (u)`$ as the set of coherent state expectation values, $`A_\chi (u)=u|\widehat{A}|u_\chi `$, $`uG`$. We assume that the only operator whose symbol vanishes identically is the null operator. Thus, distinct operators have different symbols, which means that any operator can, in principle, be completely reconstructed solely from its diagonal matrix elements in the coherent state basis.
Classical observables will be associated with operators that remain non-singular as $`\chi `$ goes to zero, that is, whose coherent state matrix elements, $`u|\widehat{A}|u^{}_\chi /u|u^{}_\chi `$, do not blow up as $`\chi 0`$ for all $`u,u^{}G`$. Such operators are called classical.
Two coherent states $`|u`$ and $`|u^{}`$ are termed classically equivalent (we will write $`uu^{}`$) if in the $`\chi 0`$ limit, one can not distinguish between them using only classical operators, i.e., $`lim_{\chi 0}A_\chi (u)=lim_{\chi 0}A_\chi (u^{})`$ for all classical operators $`\widehat{A}`$. We assume that the overlap between any two classically inequivalent coherent states decreases exponentially with $`1/\chi `$ in the $`\chi 0`$ limit.
Under these assumptions, one may show that the $`\chi 0`$ limit of this theory truly is a classical limit . The assumptions hold for $`O(N)`$ or $`U(N)`$ invariant vector models, matrix models, and gauge theories . The quantum dynamics reduces to classical dynamics on a phase space $`\mathrm{\Gamma }`$ given by a coadjoint orbit of the coherence group. Formally, points in $`\mathrm{\Gamma }`$ correspond to equivalence classes of coherent states, $`\mathrm{\Gamma }=\{[u]:uG\}`$, with $`[u]=\{u^{}G:uu^{}\}`$. The symplectic structure on the phase space is completely determined by the Lie algebra structure of the coherence group. The classical Hamiltonian is just the $`\chi 0`$ limit of the coherent state expectation of the quantum Hamiltonian,
$$h_{\mathrm{cl}}(u)=\underset{\chi 0}{lim}\widehat{H}_\chi (u).$$
(3)
To have sensible classical dynamics this limit must exist, i.e., $`\widehat{H}_\chi `$ must be a classical operator. (This is why it was convenient the rescale the Hamiltonian by $`\mathrm{}/\chi `$.) The classical action is
$$S_{\mathrm{cl}}[u(t)]=\underset{\chi 0}{lim}𝑑tu(t)|i\chi _t\widehat{H}_\chi |u(t)_\chi .$$
(4)
Both the classical Hamiltonian (3) and the action (4) depend only on the equivalence class of the coherent state $`|u`$,<sup>§</sup><sup>§</sup>§For the action, this is true up to temporal boundary terms which do not affect the dynamics. and thus do define sensible dynamics on the classical phase space.
The preceding discussion is just a formalization of the usual picture of a classical limit. A quantum mechanical wave packet, with a width of order $`\chi ^{1/2}`$, behaves classically in the $`\chi 0`$ limit, and may be associated with a point in the classical phase space. The equations of motion that govern the classical dynamics are just coherent state expectations of the original quantum evolution equations.
## III Coherent State Expectations
As noted earlier, the irreducibility of the coherence group implies that all operators may be (formally) constructed from the generators of the coherence group. Consequently, for characterizing the structure, and time evolution, of any state, one may focus attention on equal-time expectation values of products of coherence group generators.
Let $`g`$ denote the Lie algebra of the coherence group $`G`$. Let $`\{e_i\}`$ be a basis of $`g`$. The commutator of basis elements defines the structure constants, $`[e_i,e_j]=if_{ij}^ke_k`$. The generators $`e_i`$ themselves are not classical operators, but rather are $`1/\chi `$ times classical operators. For convenience, let $`\widehat{x}_i`$ denote the rescaled generator which is a classical operator, $`\widehat{x}_i\chi e_i`$.
Consider the coherent state expectation value of the monomial $`\widehat{x}_{i_1}\widehat{x}_{i_2}\mathrm{}\widehat{x}_{i_k}`$. We would like to find an expansion of this expectation value in powers of $`\chi `$. A convenient representation for our purposes involves subtracted expectations To simplify notation, we will omit the superscript “$`(k)`$” when this can cause no confusion; for example, we will write $`g_{ij}`$ for $`g_{ij}^{(2)}`$, etc. The same remark applies to the connected expectations discussed below.
$`g_{i_1i_2\mathrm{}i_k}^{(k)}`$ $``$ $`(\widehat{x}_{i_1}x_{i_1})\mathrm{}(\widehat{x}_{i_k}x_{i_k}),`$ (5)
where $`\mathrm{}`$ denotes an expectation in some coherent state, and $`x_i\widehat{x}_i`$ are the expectations of the rescaled generators $`\widehat{x}_i`$. Subtracted and un-subtracted expectations are related by
$`\widehat{x}_{i_1}\mathrm{}\widehat{x}_{i_k}`$ $`=`$ $`x_{i_1}\mathrm{}x_{i_k}\times \left\{1+{\displaystyle \underset{(l_1,l_2)}{}}{\displaystyle \frac{g_{l_1l_2}^{(2)}}{x_{l_1}x_{l_2}}}+{\displaystyle \underset{(l_1,l_2,l_3)}{}}{\displaystyle \frac{g_{l_1l_2l_3}^{(3)}}{x_{l_1}x_{l_2}x_{l_3}}}+\mathrm{}+{\displaystyle \frac{g_{i_1\mathrm{}i_k}^{(k)}}{x_{i_1}\mathrm{}x_{i_k}}}\right\},`$ (6)
where the $`n`$-tuples $`(l_1,l_2,\mathrm{},l_n)`$ are ordered subsets of $`\{i_1,\mathrm{},i_k\}`$. (There is no $`g^{(1)}`$ term since $`g_i^{(1)}\widehat{x}_ix_i=0`$.)
Alternatively, one may expand in terms of connected expectations,
$$s_{i_1\mathrm{}i_k}^{(k)}\widehat{x}_{i_1}\mathrm{}\widehat{x}_{i_k}^{\mathrm{conn}}.$$
(7)
The difference, illustrated graphically in Fig. 1, is that expansions in terms of connected expectations involve products of all possible ‘contractions’, while the terms in the expansion in subtracted expectations have only one string of generators ‘contracted’. The difference between subtracted and connected expectations first arises with four generators. Explicitly,
$`\widehat{x}_i\widehat{x}_j\widehat{x}_k\widehat{x}_l`$ $`=`$ $`x_ix_jx_kx_l+x_ix_jg_{kl}^{(2)}+x_ix_kg_{jl}^{(2)}+x_ix_lg_{jk}^{(2)}+x_kx_lg_{ij}^{(2)}+x_jx_lg_{ik}^{(2)}+x_jx_kg_{il}^{(2)}`$ (10)
$`+x_ig_{jkl}^{(3)}+x_jg_{ikl}^{(3)}+x_kg_{ijl}^{(3)}+x_lg_{ijk}^{(3)}+g_{ijkl}^{(4)}`$
$`=`$ $`x_ix_jx_kx_l+x_ix_js_{kl}^{(2)}+x_ix_ks_{jl}^{(2)}+x_ix_ls_{jk}^{(2)}+x_kx_ls_{ij}^{(2)}+x_jx_ls_{ik}^{(2)}+x_jx_ks_{il}^{(2)}`$ (12)
$`+x_is_{jkl}^{(3)}+x_js_{ikl}^{(3)}+x_ks_{ijl}^{(3)}+x_ls_{ijk}^{(3)}+s_{ij}^{(2)}s_{kl}^{(2)}+s_{ik}^{(2)}s_{jl}^{(2)}+s_{il}^{(2)}s_{jk}^{(2)}+s_{ijkl}^{(4)}.`$
The coherent state overlap $`u|u^{}_\chi `$ is the generating functional for expectations of products of generators, since variations of the coherent state $`u^{}`$ can bring down any desired generator of the Lie algebra, $`\delta _i|u^{}=e_i|u^{}`$. The logarithm of this overlap is therefore the generating functional for connected expectations. By assumption, $`\mathrm{ln}u|u^{}_\chi `$ is $`𝒪(1/\chi )`$ as $`\chi 0`$. This immediately implies that the $`k`$-th order connected expectation $`s^{(k)}`$ is $`𝒪(\chi ^{k1})`$. Note also that the commutator of functional derivatives is the functional derivative in the direction of the commutator,
$$\mathrm{}\widehat{x}_i\widehat{x}_j\mathrm{}^{\mathrm{conn}}\mathrm{}\widehat{x}_j\widehat{x}_i\mathrm{}^{\mathrm{conn}}=\mathrm{}[\widehat{x}_i,\widehat{x}_j]\mathrm{}^{\mathrm{conn}}=i\chi f_{ij}^m\mathrm{}\widehat{x}_m\mathrm{}^{\mathrm{conn}},$$
(13)
or $`s_{\mathrm{}ij\mathrm{}}^{(k)}s_{\mathrm{}ji\mathrm{}}^{(k)}=i\chi f_{ij}^ms_\mathrm{}m\mathrm{}^{(k1)}`$.
By considering which connected expectations contribute to $`g^{(k)}`$, one may easily see As $`g^{(k)}=s^{(m_1)}s^{(m_2)}\mathrm{}s^{(m_n)}`$ with $`m_1+\mathrm{}+m_n=k`$ and $`m_i>1`$ for all $`i`$, and $`s^{(m_i)}=𝒪(\chi ^{m_i1})`$, $`g^{(k)}=𝒪(\chi ^{\mathrm{min}\{{\scriptscriptstyle (m_i1)}\}})=𝒪(\chi ^{\mathrm{min}\{kn\}})`$. The largest number of connected diagrams $`n`$ occurs when all $`s^{(m_i)}`$ are $`s^{(2)}`$ (except for one, if $`k`$ is odd, which is $`s^{(3)}`$). that subtracted expectations fall off roughly half as fast as the connected ones, $`g^{(2k)}g^{(2k1)}=𝒪(\chi ^k)`$. Because $`g^{(k+2)}/g^{(k)}=𝒪(\chi )`$, expansion (6) is a power series in $`\chi `$, the parameter measuring how close the system is to being classical. Of course, subtracted expectations may always be rewritten in terms of connected expectations (and vice-versa). Ultimately, equations for connected<sup>\**</sup><sup>\**</sup>\**Or perhaps one-particle irreducible. expectations will be most useful. Nevertheless, using subtracted expectations as an intermediate representation is helpful because of the simple form of expansions in terms of subtracted expectations, as shown by Eq. (6) and Eq. (15) below. For later use, note that
$$g_{\mathrm{}bijc\mathrm{}}^{(k)}g_{\mathrm{}bjic\mathrm{}}^{(k)}=i\chi f_{ij}^m\left(x_mg_{\mathrm{}bc\mathrm{}}^{(k2)}+g_{\mathrm{}bmc\mathrm{}}^{(k1)}\right).$$
(14)
Now consider an operator $`V`$ that can (at least formally) be expanded in powers of generators, $`V=_k_{\{i_1\mathrm{}i_k\}}\alpha _{i_1\mathrm{}i_k}\widehat{x}_{i_1}\mathrm{}\widehat{x}_{i_k}`$, for some set of coefficients $`\{\alpha _{i_1\mathrm{}i_k}\}`$. Operators of this form are well behaved for $`\chi 0`$, and so are good classical operators. Using (6),
$`V`$ $`=`$ $`{\displaystyle \underset{k}{}}{\displaystyle \underset{\{i_1\mathrm{}i_k\}}{}}\alpha _{i_1\mathrm{}i_k}x_{i_1}\mathrm{}x_{i_k}\left\{1+{\displaystyle \underset{(l_1,l_2)}{}}{\displaystyle \frac{g_{l_1l_2}^{(2)}}{x_{l_1}x_{l_2}}}+{\displaystyle \underset{(l_1,l_2,l_3)}{}}{\displaystyle \frac{g_{l_1l_2l_3}^{(3)}}{x_{l_1}x_{l_2}x_{l_3}}}+\mathrm{}+{\displaystyle \frac{g_{i_1\mathrm{}i_k}^{(k)}}{x_{l_1}\mathrm{}x_{l_k}}}\right\}.`$ (15)
This can be packaged in an even more concise form,
$$V=\overline{V}+g_{l_1l_2}^{(2)}V^{(l_1l_2)}+g_{l_1l_2l_3}^{(3)}V^{(l_1l_kl_3)}+g_{l_1l_2l_3l_4}^{(4)}V^{(l_1l_kl_3l_4)}+𝒪(\chi ^3)$$
(16)
where summation on repeated indices is implied, and $`\overline{V}_k_{\{i_1\mathrm{}i_k\}}\alpha _{i_1\mathrm{}i_k}x_{i_1}\mathrm{}x_{i_k}`$ is the number obtained by replacing each generator in $`V`$ by its coherent state expectation. Here we have introduced “ordered derivatives” $`f^{(ij\mathrm{})}\frac{\delta f}{\delta (x_ix_j\mathrm{})}`$ defined by
$`f^{(l)}`$ $`=`$ $`{\displaystyle \frac{\overline{f}}{x_l}},`$ (18)
$`(fg)^{(l_1\mathrm{}l_k)}`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{k}{}}}f^{(l_1\mathrm{}l_i)}g^{(l_{i+1}\mathrm{}l_k)}.`$ (19)
When acting on a string of generators, ordered derivatives produce a sum of products of expectations of the generators which remain after deleting the indicated generators, provided these appear (not necessarily contiguously) somewhere within the string in the order specified by the derivative. For example,<sup>††</sup><sup>††</sup>†† If $`\widehat{x}_{i_1},\mathrm{},\widehat{x}_{i_n}`$ commute, then $`\frac{\delta f}{\delta (x_{i_1}\mathrm{}x_{i_n})}=\frac{1}{n!}\frac{^n\overline{f}}{x_{i_1}\mathrm{}x_{i_n}}`$. In this case, the ordering does not matter, and the $`n!`$ is needed to make up for over-counting. $`\frac{\delta }{\delta (xp)}\widehat{x}^2\widehat{p}\widehat{x}=2x^2`$, and $`\frac{\delta }{\delta (xp)}\widehat{x}\widehat{p}\widehat{x}^2\widehat{p}=x^2p+2xpx+px^2=4x^2p`$.
In the $`\chi 0`$ limit, coherent state expectations of the (rescaled) generators $`\widehat{x}_i`$ turn into coordinates $`x_i`$ on the classical phase space and (classical) operators acting on $`_\chi `$ become functions on phase space, $`V=\overline{V}+O(\chi )`$. For finite $`\chi `$, the successive terms in (16) precisely characterize the corrections to this classical limit.
## IV Time Evolution
Since operators are completely determined by their symbols, to study the time dependence of any observable $`\widehat{A}`$ it is sufficient to take the coherent state expectation value of its Heisenberg equation of motion (1),
$$\frac{d}{dt}\widehat{A}=\frac{i}{\chi }[\widehat{H}_\chi ,\widehat{A}].$$
(20)
In other words, we assume that the initial state is precisely some coherent state $`|u`$, and wish to determine the subsequent time evolution. To do so, we will first find an expansion, in powers of $`\chi `$, for the expectation of the commutator of classical operators.
### A Symbols of Commutators
Consider classical operators $`A`$ and $`B`$ which (as in Section III) may be written as power series in the generators, $`A=\alpha _{i_1\mathrm{}i_m}\widehat{x}_{i_1}\mathrm{}\widehat{x}_{i_m}`$, $`B=\beta _{j_1\mathrm{}j_n}\widehat{x}_{j_1}\mathrm{}\widehat{x}_{j_n}`$. Their product is given by $`AB=\alpha _{i_1\mathrm{}i_m}\beta _{j_1\mathrm{}j_n}\widehat{x}_{i_1}\mathrm{}\widehat{x}_{i_m}\widehat{x}_{j_1}\mathrm{}\widehat{x}_{j_n}`$. Using our previous result (16), we find
$`AB`$ $`=`$ $`\overline{A}\overline{B}+g_{l_1l_2}^{(2)}(AB)^{(l_1l_2)}+g_{l_1l_2l_3}^{(3)}(AB)^{(l_1l_2l_3)}+g_{l_1l_2l_3l_4}^{(4)}(AB)^{(l_1l_2l_3l_4)}+𝒪(\chi ^3)`$ (22)
$`=`$ $`\overline{A}\overline{B}+s_{l_1l_2}^{(2)}(AB)^{(l_1l_2)}+s_{l_1l_2l_3}^{(3)}(AB)^{(l_1l_2l_3)}`$ (24)
$`+\left(s_{l_1l_2}^{(2)}s_{l_3l_4}^{(2)}+s_{l_1l_3}^{(2)}s_{l_2l_4}^{(2)}+s_{l_1l_4}^{(2)}s_{l_2l_3}^{(2)}\right)(AB)^{(l_1l_2l_3l_4)}+𝒪(\chi ^3)`$
where now $`(l_1,l_2,\mathrm{})`$ denote ordered subsets of $`\{i_1,\mathrm{},i_m,j_1,\mathrm{},i_n\}`$. We see from (IV A) that, to leading order, products of classical operators factorize, $`AB=AB+𝒪(\chi )`$.
Using the expansion (IV A), and the reduction formulas for operator derivatives (III), one can evaluate the commutator. A generic term in the result will be
$`g_{l_1l_2\mathrm{}l_k}^{(k)}\left[(AB)^{(l_1l_2\mathrm{}l_k)}(BA)^{(l_1l_2\mathrm{}l_k)}\right]`$ (25)
$`=`$ $`g_{l_1\mathrm{}l_k}^{(k)}\left(A^{(l_1\mathrm{}l_k)}\overline{B}+\mathrm{}+A^{(l_1l_2)}B^{(l_3\mathrm{}l_k)}+A^{(l_1)}B^{(l_2\mathrm{}l_k)}+\overline{A}B^{(l_1\mathrm{}l_k)}\right)`$ (27)
$`g_{l_1\mathrm{}l_k}^{(k)}\left(B^{(l_1\mathrm{}l_k)}\overline{A}+\mathrm{}+B^{(l_1l_2)}A^{(l_3\mathrm{}l_k)}+B^{(l_1)}A^{(l_2\mathrm{}l_k)}+\overline{B}A^{(l_1\mathrm{}l_k)}\right)`$
$`=`$ $`\left(g_{l_1l_2l_3l_4\mathrm{}l_k}^{(k)}g_{l_2l_3l_4\mathrm{}l_kl_1}^{(k)}\right)\left(A^{(l_1)}B^{(l_2l_3l_4\mathrm{}l_k)}B^{(l_1)}A^{(l_2l_3l_4\mathrm{}l_k)}\right)`$ (30)
$`+\left(g_{l_1l_2l_3l_4\mathrm{}l_k}^{(k)}g_{l_3l_4\mathrm{}l_kl_1l_2}^{(k)}\right)\left(A^{(l_1l_2)}B^{(l_3l_4\mathrm{}l_k)}B^{(l_1l_2)}A^{(l_3l_4\mathrm{}l_k)}\right)`$
$`+\left(g_{l_1l_2l_3l_4\mathrm{}l_k}^{(k)}g_{l_4\mathrm{}l_kl_1l_2l_3}^{(k)}\right)\left(A^{(l_1l_2l_3)}B^{(l_4\mathrm{}l_k)}B^{(l_1l_2l_3)}A^{(l_4\mathrm{}l_k)}\right)+\mathrm{}.`$
The last term in the sum (25) is either $`\left(g_{l_1\mathrm{}l_jl_{j+1}\mathrm{}l_{2j}}^{(2j)}g_{l_{j+1}\mathrm{}l_{2j}l_1\mathrm{}l_j}^{(2j)}\right)\left(A^{(l_1\mathrm{}l_j)}B^{(l_{j+1}\mathrm{}l_{2j})}\right)`$ or $`\left(g_{l_1\mathrm{}l_jl_{j+1}\mathrm{}l_{2j+1}}^{(2j+1)}g_{l_{j+1}\mathrm{}l_{2j+1}l_1\mathrm{}l_j}^{(2j+1)}\right)`$$`\left(A^{(l_1\mathrm{}l_j)}B^{(l_{j+1}\mathrm{}l_{2j+1})}B^{(l_1\mathrm{}l_j)}A^{(l_{j+1}\mathrm{}l_{2j+1})}\right),`$ depending on whether $`k`$ is even or odd. Using (14) to reduce the differences $`(g_{\mathrm{}}^{(k)}g_{\mathrm{}}^{(k)})`$ yields the final form for the expectation of the commutator of classical operators. The leading term is precisely the Poisson bracket on the classical phase space, while subsequent terms involve successively higher expectations $`g^{(k)}`$. Displaying subleading $`𝒪(\chi )`$ and $`𝒪(\chi ^2)`$ terms explicitly, one finds
$`{\displaystyle \frac{1}{i\chi }}[A,B]`$ (38)
$`\begin{array}{c}=\left[f_{l_1l_2}^mx_m\right]\left(A^{(l_1)}B^{(l_2)}\right)\end{array}\}𝒪(\chi ^0)`$
$`\begin{array}{c}+\left[f_{l_1l_2}^mg_{ml_3}+f_{l_1l_3}^mg_{l_2m}\right]\left(A^{(l_1)}B^{(l_2l_3)}B^{(l_1)}A^{(l_2l_3)}\right)\\ +\left[x_m\left(f_{l_1l_2}^mg_{l_3l_4}+f_{l_1l_3}^mg_{l_2l_4}+f_{l_1l_4}^mg_{l_2l_3}\right)\right]\left(A^{(l_1)}B^{(l_2l_3l_4)}B^{(l_1)}A^{(l_2l_3l_4)}\right)\\ +\left[x_m\left(f_{l_1l_3}^mg_{l_2l_4}+f_{l_1l_4}^mg_{l_3l_2}+f_{l_2l_3}^mg_{l_1l_4}+f_{l_2l_4}^mg_{l_3l_1}\right)\right]\left(A^{(l_1l_2)}B^{(l_3l_4)}\right)\end{array}\}𝒪(\chi ^1)`$
$`\begin{array}{c}+\left[f_{l_1l_2}^mg_{ml_3l_4}+f_{l_1l_3}^mg_{l_2ml_4}+f_{l_1l_4}^mg_{l_2l_3m}\right]\left(A^{(l_1)}B^{(l_2l_3l_4)}B^{(l_1)}A^{(l_2l_3l_4)}\right)\\ +\left[f_{l_1l_3}^mg_{ml_2l_4}+f_{l_1l_4}^mg_{l_3ml_2}+f_{l_2l_3}^mg_{l_1ml_4}+f_{l_2l_4}^mg_{l_3l_1m}\right]\left(A^{(l_1l_2)}B^{(l_3l_4)}\right)\\ +[x_m(f_{l_1l_2}^mg_{l_3l_4l_5}+f_{l_1l_3}^mg_{l_2l_4l_5}+f_{l_1l_4}^mg_{l_2l_3l_5}+f_{l_1l_5}^mg_{l_2l_3l_4})\\ +f_{l_1l_2}^mg_{ml_3l_4l_5}+f_{l_1l_3}^mg_{l_2ml_4l_5}+f_{l_1l_4}^mg_{l_2l_3ml_5}+f_{l_1l_5}^mg_{l_2l_3l_4m}]\\ \times \left(A^{(l_1)}B^{(l_2l_3l_4l_5)}B^{(l_1)}A^{(l_2l_3l_4l_5)}\right)\\ +[x_m(f_{l_1l_3}^mg_{l_4l_5l_2}+f_{l_1l_4}^mg_{l_3l_5l_2}+f_{l_1l_5}^mg_{l_3l_4l_2}+f_{l_2l_3}^mg_{l_1l_4l_5}+f_{l_2l_4}^mg_{l_1l_3l_5}\\ +f_{l_2l_5}^mg_{l_1l_3l_4})+f^m_{l_1l_3}g_{ml_4l_5l_2}+f^m_{l_1l_4}g_{l_3ml_5l_2}+f^m_{l_1l_5}g_{l_3l_4ml_2}+f^m_{l_2l_3}g_{l_1ml_4l_5}\\ +f_{l_2l_4}^mg_{l_1l_3ml_5}+f_{l_2l_5}^mg_{l_1l_3l_4m}\left]\right(A^{(l_1l_2)}B^{(l_3l_4l_5)}B^{(l_1l_2)}A^{(l_3l_4l_5)})\\ +\left[x_m\left(f_{l_1l_2}^mg_{l_3l_4l_5l_6}+f_{l_1l_3}^mg_{l_2l_4l_5l_6}+f_{l_1l_4}^mg_{l_2l_3l_5l_6}+f_{l_1l_5}^mg_{l_2l_3l_4l_6}+f_{l_1l_6}^mg_{l_2l_3l_4l_5}\right)\right]\\ \times \left(A^{(l_1)}B^{(l_2l_3l_4l_5l_6)}B^{(l_1)}A^{(l_2l_3l_4l_5l_6)}\right)\\ +[x_m(f_{l_1l_3}^mg_{l_4l_5l_6l_2}+f_{l_1l_4}^mg_{l_3l_5l_6l_2}+f_{l_1l_5}^mg_{l_3l_4l_6l_2}+f_{l_1l_6}^mg_{l_3l_4l_5l_2}\\ +f_{l_2l_3}^mg_{l_1l_4l_5l_6}+f_{l_2l_4}^mg_{l_1l_3l_5l_6}+f_{l_2l_5}^mg_{l_1l_3l_4l_6}+f_{l_2l_6}^mg_{l_1l_3l_4l_5})]\\ \times \left(A^{(l_1l_2)}B^{(l_3l_4l_5l_6)}B^{(l_1l_2)}A^{(l_3l_4l_5l_6)}\right)\\ +[x_m(f_{l_1l_4}^mg_{l_5l_6l_2l_3}+f_{l_1l_5}^mg_{l_4l_6l_2l_3}+f_{l_1l_6}^mg_{l_4l_5l_2l_3}+f_{l_2l_4}^mg_{l_1l_5l_6l_3}+f_{l_2l_5}^mg_{l_1l_4l_6l_3}\\ +f_{l_2l_6}^mg_{l_1l_4l_5l_3}+f_{l_3l_4}^mg_{l_1l_2l_5l_6}+f_{l_3l_5}^mg_{l_1l_2l_4l_6}+f_{l_3l_6}^mg_{l_1l_2l_4l_5})\left]\right(A^{(l_1l_2l_3)}B^{(l_4l_5l_6)})\end{array}\}𝒪(\chi ^2)`$
$`\begin{array}{c}+𝒪(\chi ^3).\end{array}`$
### B Equations of Motion
To determine the evolution to order $`𝒪(\chi ^3)`$, we need the time derivatives of $`x_i(t)`$, $`g_{ij}(t)`$, and $`g_{ijk}(t)`$. Take the commutator of products of generators with the Hamiltonian and subtract the disconnected parts to find:
$`{\displaystyle \frac{d}{dt}}x_i`$ $`=`$ $`\left[f_{ij}^ax_a\right]H^{(j)}+\left[f_{ij}^ag_{ak}+f_{ik}^ag_{ja}\right]H^{(jk)}`$ (39)
$`+`$ $`\left[x_a\left(f_{ij}^ag_{kl}+f_{ik}^ag_{jl}+f_{il}^ag_{jk}\right)+f_{ij}^ag_{akl}+f_{ik}^ag_{jal}+f_{il}^ag_{jka}\right]H^{(jkl)}`$ (40)
$`+`$ $`[x_a(f_{ij}^ag_{klm}+f_{ik}^ag_{jlm}+f_{il}^ag_{jkm}+f_{im}^ag_{jkl})`$ (42)
$`+f_{ij}^ag_{aklm}+f_{ik}^ag_{jalm}+f_{il}^ag_{jkam}+f_{im}^ag_{jkla}]H^{(jklm)}`$
$`+`$ $`\left[x_a\left(f_{ij}^ag_{klmn}+f_{ik}^ag_{jlmn}+f_{il}^ag_{jkmn}+f_{im}^ag_{jkln}+f_{in}^ag_{jklm}\right)\right]H^{(jklmn)}`$ (43)
$`+`$ $`𝒪(\chi ^3)`$ (44)
$`{\displaystyle \frac{d}{dt}}g_{ij}`$ $`=`$ $`\left[f_{ik}^ag_{aj}+f_{jk}^ag_{ia}\right]H^{(k)}`$ (45)
$`+`$ $`\left[x_a\left(f_{ik}^ag_{jl}+f_{il}^ag_{kj}+f_{jk}^ag_{il}+f_{jl}^ag_{ki}\right)+f_{ik}^ag_{ajl}+f_{il}^ag_{kaj}+f_{jk}^ag_{ial}+f_{jl}^ag_{kia}\right]H^{(kl)}`$ (46)
$`+`$ $`[x_a(f_{ik}^ag_{lmj}+f_{il}^ag_{kmj}+f_{im}^ag_{klj}+f_{jk}^ag_{ilm}+f_{jl}^ag_{ikm}+f_{jm}^ag_{ikl})`$ (48)
$`+f_{ik}^ag_{almj}+f_{il}^ag_{kamj}+f_{im}^ag_{klaj}+f_{jk}^ag_{ialm}+f_{jl}^ag_{ikam}+f_{jm}^ag_{ikla}]H^{(klm)}`$
$`+`$ $`[x_a(f_{ik}^ag_{lmnj}+f_{il}^ag_{kmnj}+f_{im}^ag_{klnj}+f_{in}^ag_{klmj}`$ (50)
$`+f_{jk}^ag_{ilmn}+f_{jl}^ag_{ikmn}+f_{jm}^ag_{ikln}+f_{jn}^ag_{iklm})]H^{(klmn)}`$
$`+`$ $`𝒪(\chi ^3)`$ (51)
$`{\displaystyle \frac{d}{dt}}g_{ijk}`$ $`=`$ $`\left[f_{il}^ag_{ajk}+f_{jl}^ag_{iak}+f_{kl}^ag_{ija}\right]H^{(l)}`$ (52)
$`+`$ $`\left[x_a\left(f_{il}^ag_{jkm}+f_{jl}^ag_{ikm}+f_{kl}^ag_{ijm}+f_{im}^ag_{ljk}+f_{jm}^ag_{lik}+f_{km}^ag_{lij}\right)\right]H^{(lm)}`$ (53)
$`+`$ $`[f_{il}^a(g_{aj}g_{km}+g_{ak}g_{jm})+f_{jl}^a(g_{ia}g_{km}+g_{im}g_{ak})+f_{kl}^a(g_{ia}g_{jm}+g_{im}g_{aj})`$ (55)
$`+f_{im}^a(g_{lj}g_{ak}+g_{lk}g_{aj})+f_{jm}^a(g_{li}g_{ak}+g_{la}g_{ki})+f_{km}^a(g_{li}g_{ja}+g_{lj}g_{ia})]H^{(lm)}`$
$`+`$ $`[x_a(f_{il}^a(g_{mj}g_{nk}+g_{mk}g_{nj})+f_{im}^a(g_{lj}g_{nk}+g_{lk}g_{nj})+f_{in}^a(g_{lj}g_{mk}+g_{lk}g_{mj})`$ (58)
$`+f_{jl}^a(g_{im}g_{nk}+g_{in}g_{mk})+f_{jm}^a(g_{il}g_{nk}+g_{in}g_{lk})+f_{jn}^a\left(g_{il}g_{mk}+g_{im}g_{lk}\right)`$
$`+f_{kl}^a(g_{im}g_{jn}+g_{in}g_{jm})+f_{km}^a(g_{il}g_{jn}+g_{in}g_{jl})+f_{kn}^a(g_{il}g_{jm}+g_{im}g_{jl}))]H^{(lmn)}`$
$`+`$ $`𝒪(\chi ^3).`$ (59)
Recall that, through third order, there is no difference between the subtracted and connected correlators. Only the disconnected parts of the fourth order correlators appearing in Eq’s. (39) and (45) are needed, since $`g_{ijkl}=g_{ij}g_{kl}+g_{ik}g_{jl}+g_{il}g_{jk}+𝒪(\chi ^3)`$. If equations only accurate to $`𝒪(\chi ^2)`$ are desired, then all terms in Eq’s. (39)–(52) involving third (or higher) order correlators, as well as products of second order correlators, may be dropped.<sup>‡‡</sup><sup>‡‡</sup>‡‡ The resulting next-to-leading order equations are simply $`{\displaystyle \frac{d}{dt}}x_i`$ $`=`$ $`\left(f_{ij}^ax_a\right)H^{(j)}+\left(f_{ij}^ag_{ak}+f_{ik}^ag_{ja}\right)H^{(jk)}+x_a\left(f_{ij}^ag_{kl}+f_{ik}^ag_{jl}+f_{il}^ag_{jk}\right)H^{(jkl)}+𝒪(\chi ^2),`$ (61) and $`{\displaystyle \frac{d}{dt}}g_{ij}`$ $`=`$ $`\left(f_{ik}^ag_{aj}+f_{jk}^ag_{ia}\right)H^{(k)}+x_a\left(f_{ik}^ag_{jl}+f_{il}^ag_{kj}+f_{jk}^ag_{il}+f_{jl}^ag_{ki}\right)H^{(kl)}+𝒪(\chi ^2).`$ (62)
Given these equations of motion for the connected expectations of generators, one can use (16) to describe the dynamics of any classical operator in terms of its symbol. If $`\widehat{V}=\widehat{V}(\{x_i\})`$ is a (time-independent) function of the generators, then its time-dependent expectation value, at next-to-next-to-leading order, is given by
$`\widehat{V}(t)`$ $`=`$ $`\left\{\overline{V}+g_{ij}(t)V^{(ij)}+g_{ijk}(t)V^{(ijk)}+g_{ijkl}(t)V^{(ijkl)}\right\}|_{x_m=x_m(t)}+𝒪(\chi ^3),`$ (63)
where $`x_i(t)`$, $`g_{ij}(t)=s_{ij}(t)`$, and $`g_{ijk}(t)=s_{ijk}(t)`$ are to be obtained by integrating Eq’s. (39)–(52) forward in time, using $`g_{ijkl}(t)=g_{ij}(t)g_{kl}(t)+g_{ik}(t)g_{jl}(t)+g_{il}(t)g_{jk}(t)+𝒪(\chi ^3)`$.
### C Error Accumulation
To any given order in $`\chi `$, we have a system of non-linear, first-order, ordinary differential equations. Initial conditions are imposed by specifying $`x_i(t=0)=u|\widehat{x}_i|u`$ and $`s_{j\mathrm{}l}^{(k)}(t=0)=u|\widehat{x}_j\mathrm{}\widehat{x}_l|u^{\mathrm{conn}}`$, with $`|u`$ some chosen coherent state. Since $`s^{(k)}(t=0)`$ is $`𝒪(\chi ^{k1})`$, and the equations for $`\dot{s}^{(k)}(t)`$ involve only terms of order $`\chi ^{k1}`$ and higher, we still formally have $`s^{(k)}(t)=𝒪(\chi ^{k1})`$ for $`t>0`$. However, as the truncated equations of motion are integrated forward in time, errors accumulate; it is important to understand the rate of growth of this truncation error.
We are dealing with a system of equations which we can write as $`\dot{y}_i=F_i(y)+G_i(y)`$ where $`\{y_i(t)\}`$ are the variables in our problem (that is, the $`x_i`$’s and $`s_{\mathrm{}}^{(k)}`$’s), $`F(y)`$ represents the terms we keep, and $`G(y)`$ stands for everything thrown away by the truncation. Let $`y_0(t)`$ be the solution to the above equation with $`G0`$, and solve perturbatively, $`y(t)=y_0(t)+ϵ(t)`$ with $`ϵ`$ small. Linearizing about $`y(t)=y_0(t)`$, we have
$$\dot{ϵ}=f(t)ϵ+g(t),$$
(64)
where $`f_i^j(t)=F_i(y_0(t))/y_j`$, $`g_i(t)=G_i(y_0(t))`$, and we have dropped $`𝒪(ϵ^2)`$ terms. This linearized system of equations is easy to solve (at least formally). For $`t>0`$,
$$ϵ(t)=\left[𝒯\text{e}^{_0^tf(t^{})𝑑t^{}}\right]ϵ(0)+_0^t\left[𝒯\text{e}^{_t^{}^tf(t^{\prime \prime })𝑑t^{\prime \prime }}\right]g(t^{})𝑑t^{}.$$
(65)
Here, $`𝒯`$ denotes time ordering (with smaller times on the right). If $`f(t)`$ and $`g(t)`$ are globally bounded during the time evolution, $`f(t)\stackrel{~}{f}`$, $`g(t)\stackrel{~}{g}`$, where $`\mathrm{}`$ is some appropriate norm, then a crude estimate of the deviation of the true solution from the approximation is
$$ϵ(t)\text{e}^{\stackrel{~}{f}t}ϵ(0)+\stackrel{~}{g}(\text{e}^{\stackrel{~}{f}t}1)/\stackrel{~}{f}.$$
(66)
Of course for $`t`$ small, errors grow linearly and $`ϵ(t)ϵ(0)(1+\stackrel{~}{f}t)+\stackrel{~}{g}t+𝒪(t^2)`$; with a truncation good to order $`\chi ^k`$ at $`t=0`$, both $`ϵ(0)`$ and $`\stackrel{~}{g}`$ will be $`𝒪(\chi ^{k+1})`$.
In a general treatment, it is hard to do better than the crude bound (66). In dynamical systems with only a few degrees of freedom, there typically are “regular” portions of phase space where perturbations grow only linearly with time . However, it is not at all clear that this is applicable to the truncated quantum dynamics represented by Eq. (64).
In simple examples discussed in the following section, we will find that for times of order $`\chi ^{1/2}`$, the shape of the wavefunction of the evolving state becomes so distorted that the formal hierarchy of correlators, $`s^{(k)}𝒪(\chi ^{k1})`$, upon which the truncation scheme is based, completely breaks down. In terms of the underlying quantum dynamics, if one considers the projection of the initial coherent state wavepacket onto the exact eigenstates of the Hamiltonian, what is happening for sufficiently large time is that the contributions of different eigenstates have decohered to such an extent that the wavepacket has spread beyond recognition. Except for special non-generic cases (such as the harmonic oscillator, where there is no dispersion) one should always expect such decoherence to eventually set in.
## V Examples
We will discuss two examples of theories to which the preceding general results may be applied: the usual semi-classical limit of point particle quantum mechanics, and the large $`N`$ limit of $`O(N)`$ invariant vector models. For brevity of presentation, we will display explicitly only the first corrections to the leading classical approximation, but we emphasize that it is completely straightforward to include yet higher order corrections, such as the $`𝒪(\chi ^2)`$ terms displayed in Eq’s. (39)–(52).
### A $`\mathbf{}\mathbf{}\mathrm{𝟎}`$ Quantum Mechanics
Consider ordinary point particle quantum mechanics, in one dimension for simplicity. The coherence group $`G`$ is the Heisenberg group, generated by $`\{e_i\}=\{\widehat{x}/\mathrm{},\widehat{p}/\mathrm{},\mathrm{𝟏}/\mathrm{}\}`$. The formal parameter that controls how close the theory is to the classical limit is, of course, $`\chi =\mathrm{}`$. The rescaled generators of the coherence group, $`\widehat{x}_i=\mathrm{}e_i`$, include the position $`\widehat{x}`$ and momentum $`\widehat{p}`$ operators whose expectations will serve as classical phase space coordinates. The Heisenberg group, acting on a fixed Gaussian base state, generates conventional coherent states $`\{|p,q\}`$, with wave functions given (up to an overall phase) by
$$x|p,q=(\pi \mathrm{})^{1/4}\mathrm{exp}\left\{\frac{1}{\mathrm{}}\left[ipx\frac{1}{2}(xq)^2\right]\right\}.$$
(67)
We have arbitrarily chosen units such that our Gaussian base state has equal variance in $`\widehat{x}`$ and $`\widehat{p}`$. Consider a Hamiltonian of the typical form $`\widehat{H}=\frac{1}{2}\widehat{p}^2+V(\widehat{x})`$, where, for simplicity, we have set the particle mass to unity. The equations of motion are, of course,
$$\frac{d}{dt}\widehat{x}=\widehat{p},\frac{d}{dt}\widehat{p}=V^{}(\widehat{x}).$$
(68)
We are interested in the time evolution of $`x(t)`$, $`p(t)`$, and the connected correlators $`g_{xx}(t)`$, $`g_{xp}(t)=g_{px}^{}(t)`$, and $`g_{pp}(t)`$, all to order $`\mathrm{}`$. From equations (39) and (45) we find:
$$\begin{array}{cccc}\hfill \dot{x}& =& p\hfill & +𝒪(\mathrm{}^2),\hfill \\ \hfill \dot{p}& =& V^{}\frac{1}{2}V^{\prime \prime \prime }g_{xx}\hfill & +𝒪(\mathrm{}^2),\hfill \\ \hfill \dot{g}_{xx}& =& g_{xp}+g_{px}\hfill & +𝒪(\mathrm{}^2),\hfill \\ \hfill (\dot{g}_{px})^{}=\dot{g}_{xp}& =& g_{pp}V^{\prime \prime }g_{xx}\hfill & +𝒪(\mathrm{}^2),\hfill \\ \hfill \dot{g}_{pp}& =& V^{\prime \prime }\left(g_{xp}+g_{px}\right)\hfill & +𝒪(\mathrm{}^2),\hfill \end{array}$$
(69)
subject to the initial conditions $`x(0)=x_0`$, $`p(0)=p_0`$, $`g_{xx}=g_{pp}=\frac{1}{2}\mathrm{}`$, $`g_{xp}=g_{px}=\frac{1}{2}i\mathrm{}`$. Notice that to this order, $`detg^{(2)}=g_{xx}g_{pp}g_{xp}g_{px}=𝒪(\mathrm{}^3)`$ is a constant of the motion, and Eq’s. (69) are equivalent to a Gaussian variational ansatz (where one approximates the wave packet by a Gaussian with a time-dependent centroid and width). However, if we went to next-to-next-to-leading order in $`\mathrm{}`$ it would become clear that our setup is different. For positive times, higher moments will not be given by simple algebraic expressions in terms of the centroid and variance, and the details of evolution will depend on the shape of the potential.<sup>\**</sup><sup>\**</sup>\** In our Gaussian initial state, all connected correlators higher than second order vanish at time zero, $`s^{(k>2)}(0)0`$. But these moments cannot remain zero unless the potential is harmonic. For example, using Eq. (52) we find that $`\dot{s}_{xpx}=s_{ppx}+s_{xpp}V^{\prime \prime }s_{xxx}V^{\prime \prime \prime }\left(s_{xx}\right)^2+𝒪(\mathrm{}^3)`$, showing explicitly that any nonzero $`V^{\prime \prime \prime }`$ will drive the skewness moments $`s_{ijk}(t)`$ away from zero.
As a trivial warm-up, consider the harmonic oscillator of unit mass and natural frequency $`\mathrm{\Omega }`$: $`\widehat{H}=\frac{1}{2}\widehat{p}^2+\frac{1}{2}\mathrm{\Omega }^2\widehat{x}^2`$. The solutions to (69) are
$`x(t)`$ $`=`$ $`x_0\mathrm{cos}\mathrm{\Omega }t+(p_0/\mathrm{\Omega })\mathrm{sin}\mathrm{\Omega }t,`$ (71)
$`p(t)`$ $`=`$ $`p_0\mathrm{cos}\mathrm{\Omega }t(x_0\mathrm{\Omega })\mathrm{sin}\mathrm{\Omega }t,`$ (72)
$`g_{xx}(t)`$ $`=`$ $`\frac{\mathrm{}}{2}\left[\mathrm{cos}^2\mathrm{\Omega }t+\mathrm{\Omega }^2\mathrm{sin}^2\mathrm{\Omega }t\right],`$ (73)
$`g_{px}^{}(t)=g_{xp}(t)`$ $`=`$ $`\frac{\mathrm{}}{2}\left[i+(\mathrm{\Omega }^1\mathrm{\Omega })\mathrm{cos}\mathrm{\Omega }t\mathrm{sin}\mathrm{\Omega }t\right],`$ (74)
$`g_{pp}(t)`$ $`=`$ $`\frac{\mathrm{}}{2}\left[\mathrm{cos}^2\mathrm{\Omega }t+\mathrm{\Omega }^2\mathrm{sin}^2\mathrm{\Omega }t\right].`$ (75)
Because the potential is quadratic these are exact. Equally simple is an inverted harmonic oscillator. If one takes the Hamiltonian to be $`\widehat{H}=\frac{1}{2}\widehat{p}^2\frac{1}{2}\mathrm{\Omega }^2\widehat{x}^2`$, then the solution of the moment equations (69) becomes
$`x(t)`$ $`=`$ $`x_0\mathrm{cosh}\mathrm{\Omega }t+(p_0/\mathrm{\Omega })\mathrm{sinh}\mathrm{\Omega }t,`$ (77)
$`p(t)`$ $`=`$ $`p_0\mathrm{cosh}\mathrm{\Omega }t+(x_0\mathrm{\Omega })\mathrm{sinh}\mathrm{\Omega }t,`$ (78)
$`g_{xx}(t)`$ $`=`$ $`\frac{\mathrm{}}{2}\left[\mathrm{cosh}^2\mathrm{\Omega }t+\mathrm{\Omega }^2\mathrm{sinh}^2\mathrm{\Omega }t\right],`$ (79)
$`g_{px}^{}(t)=g_{xp}(t)`$ $`=`$ $`\frac{\mathrm{}}{2}\left[i+(\mathrm{\Omega }^1+\mathrm{\Omega })\mathrm{cosh}\mathrm{\Omega }t\mathrm{sinh}\mathrm{\Omega }t\right],`$ (80)
$`g_{pp}(t)`$ $`=`$ $`\frac{\mathrm{}}{2}\left[\mathrm{cosh}^2\mathrm{\Omega }t+\mathrm{\Omega }^2\mathrm{sinh}^2\mathrm{\Omega }t\right].`$ (81)
In both of these examples, the time-evolution of the variances are independent of $`x_0`$ and $`p_0`$. As one would expect, they oscillate (with twice the natural frequency) in the case of the simple harmonic oscillator, and grow (exponentially) for the inverted oscillator.
As a more complicated example, consider the problem of small oscillations in a weakly anharmonic potential,<sup>\*†</sup><sup>\*†</sup>\*† We choose the curvature of the potential at the minimum to equal unity, so that our chosen coherent states have the natural width for the unperturbed potential. This ensures that the resulting dynamics (such as oscillations of $`g^{(2)}`$) are not merely reflecting purely harmonic oscillations. $`V(x)=\frac{1}{2}x^2+\beta x^4`$. The moment equations (69) become
$$\begin{array}{cccc}\hfill \ddot{x}& =& x4\beta x^312\beta xg_{xx}\hfill & +𝒪(\mathrm{}^2),\hfill \\ \hfill \dot{g}_{xx}& =& g_{xp}+g_{px}\hfill & +𝒪(\mathrm{}^2),\hfill \\ \hfill (\dot{g}_{px})^{}=\dot{g}_{xp}& =& g_{pp}g_{xx}12\beta x^2g_{xx}\hfill & +𝒪(\mathrm{}^2),\hfill \\ \hfill \dot{g}_{pp}& =& (1+12\beta x^2)\left(g_{xp}+g_{px}\right)\hfill & +𝒪(\mathrm{}^2).\hfill \end{array}$$
(82)
We will solve these perturbatively; the two small parameters are $`\beta q^2`$ and $`\beta \mathrm{}`$. We will work to first order in $`\beta \mathrm{}`$ (since we have omitted $`𝒪(\mathrm{}^2)`$ terms in the moment equations), and will display explicit results through second order in $`\beta q^2`$. In principle, one could work to any order in $`\beta q^2`$ desired.
In order to keep our error estimates simple, we will treat the time as $`𝒪(1)`$ (in units where the natural frequency is unity). This means we need not worry about the appearance of secular terms — terms which grow as powers of $`t`$ — and may solve Eq’s. (82) strictly perturbatively in the naive fashion. A straightforward calculation, with the initial conditions $`x(0)=q`$, $`p(0)=0`$, $`g_{xx}=g_{pp}=\frac{1}{2}\mathrm{}`$, and $`g_{xp}=g_{px}^{}=\frac{1}{2}i\mathrm{}`$, leads to the solution
$`x(t)`$ $`=`$ $`q\{\mathrm{cos}t+\left(\frac{1}{8}\beta q^2\right)[\mathrm{cos}3t\mathrm{cos}t12t\mathrm{sin}t]`$ (87)
$`+\left(\frac{1}{8}\beta q^2\right)^2\left[\mathrm{cos}5t24\mathrm{cos}3t+23\mathrm{cos}t+96t\mathrm{sin}t36t\mathrm{sin}3t72t^2\mathrm{cos}t\right]`$
$`+(\beta \mathrm{})\left[3t\mathrm{sin}t\right]`$
$`+(\beta \mathrm{})\left(\frac{1}{8}\beta q^2\right)\left[\frac{15}{4}(\mathrm{cos}3t\mathrm{cos}t)18t\mathrm{sin}3t+93t\mathrm{sin}t54t^2\mathrm{cos}t\right]`$
$`+𝒪\left[(\beta q^2)^3\right]+𝒪\left[(\beta \mathrm{})(\beta q^2)^2\right]+𝒪\left[(\beta \mathrm{})^2\right]\},`$
with
$`g_{xx}(t)`$ $`=`$ $`\frac{1}{2}\mathrm{}\left\{13\beta q^2\left[1\mathrm{cos}2t+t\mathrm{sin}2t\right]+𝒪\left[(\beta q^2)^2\right]+𝒪(\beta \mathrm{})\right\},`$ (88)
$`g_{px}^{}(t)=g_{xp}(t)`$ $`=`$ $`\frac{1}{2}\mathrm{}\left\{i\frac{3}{2}\beta q^2\left[3\mathrm{sin}2t+2t\mathrm{cos}2t\right]+𝒪\left[(\beta q^2)^2\right]+𝒪(\beta \mathrm{})\right\},`$ (89)
$`g_{pp}(t)`$ $`=`$ $`\frac{1}{2}\mathrm{}\left\{1+3\beta q^2\left[1\mathrm{cos}2t+t\mathrm{sin}2t\right]+𝒪\left[(\beta q^2)^2\right]+𝒪(\beta \mathrm{})\right\}.`$ (90)
Examining the secular terms in Eq’s. (87) and (88), one sees that terms of order $`\beta ^k`$ are accompanied by at most $`k`$ powers of $`t`$. This is a general result. It implies that our stated condition that the time be $`𝒪(1)`$ is needlessly restrictive. For small $`\beta q^2`$ and $`\beta \mathrm{}`$, the perturbative expansions (87) and (88) are actually valid in the wider domain $`|\beta q^2t|1`$ and $`|\beta \mathrm{}t|1`$, provided a factor of $`t`$ is included with each factor of $`\beta q^2`$ or $`\beta \mathrm{}`$ in the error estimates.
It is instructive to compare this treatment with the result of a perturbative quantum mechanical calculation. Using the brute-force approach of first finding perturbed eigenstates and energy levels, and then evaluating the time-dependent expectation value $`x(t)`$ by projecting the initial coherent state onto individual eigenstates and summing the resultant contributions, a rather tedious calculation using both wavefunctions and energies correct to $`𝒪(\beta )`$ leads to
$`x(t)`$ $`=`$ $`q(e^{(q^2/\mathrm{})\mathrm{sin}^2(3\beta \mathrm{}t/2)}\mathrm{cos}[(1+3\beta \mathrm{})t+{\displaystyle \frac{q^2}{2\mathrm{}}}\mathrm{sin}(3\beta \mathrm{}t)]`$ (97)
$`+(\beta \mathrm{})e^{(q^2/\mathrm{})\mathrm{sin}^2(3\beta \mathrm{}t/2)}\{\mathrm{cos}[(1+3\beta \mathrm{})t+{\displaystyle \frac{q^2}{2\mathrm{}}}\mathrm{sin}(3\beta \mathrm{}t)]({\displaystyle \frac{3}{2}}{\displaystyle \frac{q^2}{\mathrm{}}}{\displaystyle \frac{1}{32}}{\displaystyle \frac{q^4}{\mathrm{}^2}})`$
$`+\mathrm{cos}\left[(1+6\beta \mathrm{})t+{\displaystyle \frac{q^2}{2\mathrm{}}}\mathrm{sin}(3\beta \mathrm{}t)\right]\left({\displaystyle \frac{3}{2}}{\displaystyle \frac{3}{4}}{\displaystyle \frac{q^2}{\mathrm{}}}{\displaystyle \frac{1}{4}}{\displaystyle \frac{q^4}{\mathrm{}^2}}\right)`$
$`+\mathrm{cos}\left[(1+9\beta \mathrm{})t+{\displaystyle \frac{q^2}{2\mathrm{}}}\mathrm{sin}(3\beta \mathrm{}t)\right]\left({\displaystyle \frac{3}{2}}{\displaystyle \frac{q^2}{\mathrm{}}}\right)`$
$`+\mathrm{cos}\left[(1+12\beta \mathrm{})t+{\displaystyle \frac{q^2}{2\mathrm{}}}\mathrm{sin}(3\beta \mathrm{}t)\right]\left({\displaystyle \frac{1}{8}}{\displaystyle \frac{q^2}{\mathrm{}}}+{\displaystyle \frac{1}{4}}{\displaystyle \frac{q^4}{\mathrm{}^2}}\right)`$
$`+\mathrm{cos}[(1+15\beta \mathrm{})t+{\displaystyle \frac{q^2}{2\mathrm{}}}\mathrm{sin}(3\beta \mathrm{}t)]\left({\displaystyle \frac{1}{32}}{\displaystyle \frac{q^4}{\mathrm{}^2}}\right)\}`$
$`+(\beta \mathrm{})e^{(q^2/\mathrm{})\mathrm{sin}^2(9\beta \mathrm{}t/2)}\mathrm{cos}[3(1+6\beta \mathrm{})t+{\displaystyle \frac{q^2}{2\mathrm{}}}\mathrm{sin}(9\beta \mathrm{}t)]\left({\displaystyle \frac{1}{8}}{\displaystyle \frac{q^2}{\mathrm{}}}\right)).`$
This result has $`𝒪(\beta ^2)`$ errors due to the neglect of second (and higher) order corrections in both the eigenstates and energy eigenvalues.
If one restricts $`t`$ to be small compared to both $`1/|\beta \mathrm{}|`$ and $`1/|\beta q^2|`$, then one may expand the result (97) in powers of $`\beta `$. Moreover, in this domain one may easily add in the leading secular $`𝒪[(\beta \mathrm{})^2t]`$ terms omitted from (97), which come from including the $`𝒪(\beta ^2)`$ perturbation to energy levels while using unperturbed wavefunctions.<sup>\*‡</sup><sup>\*‡</sup>\*‡This addition is $`q(t\mathrm{sin}t)\left[\frac{51}{16}(\beta q^2)^2+\frac{153}{8}(\beta q^2)(\beta \mathrm{})+18(\beta \mathrm{})^2\right]`$. If one does not assume that $`\beta \mathrm{}t`$ is small compared to 1, then including the $`𝒪(\beta ^2)`$ energy shift in matrix elements of time-evolution operators unfortunately leads to an analytically intractable infinite sum for $`x(t)`$. One finds
$`x(t)`$ $`=`$ $`q\{\mathrm{cos}t+\left(\frac{1}{8}\beta q^2\right)[\mathrm{cos}3t\mathrm{cos}t12t\mathrm{sin}t]+\left(\frac{1}{8}\beta q^2\right)^2[96t\mathrm{sin}t36t\mathrm{sin}3t72t^2\mathrm{cos}t]`$ (100)
$`+(\beta \mathrm{})\left[3t\mathrm{sin}t\right]+(\beta \mathrm{})\left(\frac{1}{8}\beta q^2\right)\left[18t\mathrm{sin}3t+93t\mathrm{sin}t54t^2\mathrm{cos}t\right]`$
$`+𝒪\left[(\beta q^2+\beta \mathrm{})^2\right]+𝒪\left[(\beta q^2+\beta \mathrm{})^3t^3\right]\}.`$
This result is perfectly consistent with the previous moment-hierarchy result (87), as it must be, except for the non-secular $`𝒪(\beta ^2)`$ terms which are hiding in the first $`𝒪\left[(\beta q^2+\beta \mathrm{})^2\right]`$ error term of (100). If one includes second order perturbations to the eigenstates then these terms also coincide.
In the semi-classical regime, where $`\beta \mathrm{}\beta q^2`$, it is interesting to examine expression (97) when $`\beta \mathrm{}t1`$, making no assumption about the size of $`\beta q^2t`$. In this regime, the first, leading term of Eq. (97) becomes
$$x(t)=qe^{\frac{9}{4}(\beta q^2)(\beta \mathrm{})t^2}\mathrm{cos}\left[(1+3\beta \mathrm{}+\frac{3}{2}\beta q^2)t\right]+𝒪(q\beta \mathrm{}).$$
(101)
In other words, $`x(t)`$ shows damped harmonic behavior, with a shifted $`q`$-dependent frequency, and with an amplitude which decays significantly on the time scale
$$t_d\left[(\beta q^2)(\beta \mathrm{})\right]^{1/2}.$$
(102)
This implies that on this time scale, the initially well localized wavepacket has dispersed so much that its probability distribution is spread out over most of the classically allowed region.Of course, the fact that the amplitude of oscillations in the mean position $`x(t)`$ decays on the decoherence time scale $`t_d`$ cannot mean that the wavepacket has come to rest at the bottom of the potential while remaining a well-localized wavepacket, as this would violate energy conservation. In the semi-classical regime under discussion, the position of the initial wavepacket is significantly displaced from the minimum of the potential, $`q^2\mathrm{}`$, implying that the total energy is large compared to the zero-point energy. Therefore, a negligible mean position at large times necessarily indicates that the wavepacket has spread so much that its probability density, at any late time, is delocalized over the entire classically allowed region, and no longer “sloshes” back-and-forth to any significant extent. Within the classically allowed region, energy conservation implies that the wavefunction must have substantial variations on scales far smaller than the (square root of the) variance in position — which will be comparable to the width of the classically allowed region. Hence, $`t_d`$ should be regarded as a “delocalization” or “decoherence” time. The higher order terms in Eq. (97) all exhibit essentially the same behavior in this regime; each term oscillates with a (slightly different) frequency and has an amplitude which decays on the decoherence time scale $`t_d`$.
Although it will have no bearing on our discussion, it is interesting to note that on yet longer time scales, when $`t`$ is near $`2\pi /(3\beta \mathrm{})`$ or integer multiplies thereof, the exponential factors in Eq. (97) return to near unity, implying that the time-dependent state has “reassembled” itself into a recognizable wavepacket oscillating in the potential. Whether this “reassembly” persists in the exact solution, or is an artifact of our first order perturbative result, is not entirely clear to us. Presumably, this is a reflection of the fact that this is an integrable single degree of freedom system.
The existence of the decoherence time scale (102) has important consequences for the utility of any truncated moment expansion, such as Eq’s. (3952). If the wavepacket has spread to such an extent that it is significantly sampling all of its classically allowed region, while necessarily retaining structure on smaller scales, then the formal hierarchy of connected correlators, $`s^{(k)}\mathrm{}^{k1}`$, will have broken down. Higher order moments will not be small compared to lower order ones. Consequently, the moment expansion presented in the previous section can only be useful for times which are small compared to the decoherence time $`t_d`$.
The $`1/\sqrt{\mathrm{}}`$ dependence of the decoherence time (102) may also be seen in another very simple example. Consider the free evolution of a coherent state in the absence of any potential. As is well known, the width of the wavepacket grows without bound. The evolution equations for $`\widehat{x}`$ and $`\widehat{p}`$ are, of course, trivial, $`\widehat{p}(t)=\widehat{p}(0)`$, and $`\widehat{x}(t)=\widehat{x}(0)+\widehat{p}(0)t`$. Hence, $`\widehat{x}^2(t)=\widehat{x}^2(0)+[\widehat{x}(0)\widehat{p}(0)+\widehat{p}(0)\widehat{x}(0)]t+\widehat{p}(0)^2t^2`$, and so for our initial Gaussian coherent state (with equal variance in $`x`$ and $`p`$),
$`g_{xx}(t)`$ $`=`$ $`g_{xx}(0)+[g_{xp}(0)+g_{px}(0)]t+g_{pp}(0)t^2`$ (103)
$`=`$ $`\frac{1}{2}\mathrm{}(1+t^2).`$ (104)
Here also, we see that for times of order $`\mathrm{}^{1/2}`$ the hierarchy of correlators $`s^{(k)}(t)𝒪(\mathrm{}^k)`$ no longer holds.
We believe this to be a general result. Whenever a semiclassical system exhibits dispersion, the decoherence time is expected to scale as $`\mathrm{}^{1/2}`$, and truncations of the moment hierarchy equations (3952) will only be accurate for times small compared to the decoherence time.
### B Vector Models
Consider an $`O(N)`$ invariant theory whose fundamental degrees of freedom form $`O(N)`$ vectors. For simplicity, we will assume that the degrees of freedom are all bosonic,<sup>\*∥</sup><sup>\*∥</sup>\*∥ Extending this discussion to $`U(N)`$ invariant fermionic models is completely straightforward. and divided into a set of canonical coordinates $`\{\widehat{x}_\alpha ^i\}`$ and corresponding canonical momenta $`\{\widehat{p}_\beta ^j\}`$. Here $`i,j=1,\mathrm{},N`$ are $`O(N)`$ vector indices, while $`\alpha ,\beta =1,\mathrm{},m`$ distinguish different $`O(N)`$ vectors. These basic operators are assumed to satisfy canonical commutation relations, normalized such that $`[\widehat{x}_\alpha ^i,\widehat{p}_\beta ^j]=(i/N)\delta ^{ij}\delta _{\alpha \beta }`$. In other words, we have chosen to scale both coordinates and moments by $`1/\sqrt{N}`$ compared to their textbook form. The small parameter controlling the approach to the classical limit is $`\chi 1/N`$; $`\mathrm{}`$ has been set to unity. The Hamiltonian is assumed to be $`O(N)`$ invariant, and we will completely restrict attention to the $`O(N)`$ invariant sector of the theory. Consequently, the relevant Hilbert space $`_N`$ is the space of all $`O(N)`$ invariant states, and all physical operators can be constructed from the basic bilinears
$`\widehat{A}_{\alpha \beta }`$ $``$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\widehat{x}_\alpha ^i\widehat{x}_\beta ^i,`$ (106)
$`\widehat{B}_{\alpha \beta }`$ $``$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\frac{1}{2}\{\widehat{x}_\alpha ^i\widehat{p}_\beta ^i+\widehat{p}_\beta ^i\widehat{x}_\alpha ^i\},`$ (107)
$`\widehat{C}_{\alpha \beta }`$ $``$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\widehat{p}_\alpha ^i\widehat{p}_\beta ^i.`$ (108)
It will be convenient to regard $`\widehat{x}_\alpha ^i`$ and $`\widehat{p}_\alpha ^i`$ as the components of $`m\times N`$ matrices, so that the basic bilinears (V B) may be assembled into $`m\times m`$ matrices,
$$\widehat{A}=\widehat{x}\widehat{x}^T,\widehat{B}=\widehat{x}\widehat{p}^T\frac{i}{2}\widehat{\mathrm{𝟏}},\text{and}\widehat{C}=\widehat{p}\widehat{p}^T.$$
(109)
Viewed as matrices, $`\widehat{A}`$ and $`\widehat{C}`$ are symmetric, while $`\widehat{B}`$ is non-symmetric. The individual components of $`\widehat{A}`$, $`\widehat{B}`$, and $`\widehat{C}`$ are all Hermitian operators acting on $`_N`$.
We will take the Hamiltonian to have the general form
$$N\widehat{H}_N=N\left[\frac{1}{2}\mathrm{tr}\widehat{C}+V(\widehat{A})\right].$$
(110)
The overall factor of $`N`$ (given our scaling of coordinates and momenta by $`1/\sqrt{N}`$) is exactly what is needed to ensure that the $`N\mathrm{}`$ limit is a classical limit in the framework of section II. The potential energy function $`V(A)`$ may be any chosen scalar-valued function of a symmetric matrix $`A`$. The kinetic energy takes the simple form $`\frac{1}{2}\mathrm{tr}\widehat{C}=\frac{1}{2}_{i,\alpha }(\widehat{p}_\alpha ^i)^2`$ if all degrees of freedom are scaled to have unit mass. Two specific examples in this class of models are:
* A single particle moving in a central potential in $`N`$-dimensions. This is the simplest possible example; the theory has only a single $`O(N)`$ coordinate vector \[i.e., $`m=1`$\]. The Hamiltonian is
$$N\widehat{H}_N=N\left[\frac{1}{2}\stackrel{}{p}\stackrel{}{p}+V(\stackrel{}{x}\stackrel{}{x})\right]=N\left[\frac{1}{2}\widehat{C}+V(\widehat{A})\right],$$
(111)
where $`V(r^2)`$ is now a function of just a one variable.<sup>\***</sup><sup>\***</sup>\*** In terms of coordinates and momenta which have not been rescaled by $`N^{1/2}`$, $`N\widehat{H}_N=\frac{1}{2}\stackrel{}{p}^{\mathrm{\hspace{0.17em}2}}+NV(\frac{1}{N}\stackrel{}{x}^{\mathrm{\hspace{0.17em}2}})`$.
* An $`O(N)`$-invariant $`\varphi ^4`$ field theory. The theory, defined on a spatial lattice, has field operators $`\widehat{\varphi }_s^i`$ and conjugate momenta $`\widehat{\pi }_s^i`$, where $`s`$ labels the sites of some $`d`$-dimensional lattice. The canonical commutation relations (after scaling $`\widehat{\varphi }`$ and $`\widehat{\pi }`$ by $`1/\sqrt{N}`$) are $`[\widehat{\varphi }_s^i,\widehat{\pi }_s^{}^j]=(i/N)\delta ^{ij}\delta _{ss^{}}`$, and the quantum Hamiltonian is
$`N\widehat{H}_N`$ $`=`$ $`N{\displaystyle \underset{s}{}}\left[\frac{1}{2}\widehat{\pi }_s\widehat{\pi }_s+\frac{1}{2}\stackrel{}{}\widehat{\varphi }_s\stackrel{}{}\widehat{\varphi }_s+\frac{1}{2}\mu ^2\widehat{\varphi }_s\widehat{\varphi }_s+\frac{\lambda }{4}\left(\widehat{\varphi }_s\widehat{\varphi }_s\right)^2\right]`$ (113)
$`=`$ $`N{\displaystyle \underset{s}{}}\left\{\frac{1}{2}\widehat{C}_{ss}+\frac{1}{2}\left[(_s^2+\mu ^2)\widehat{A}_{ss^{}}\right]|_{s^{}=s}+\frac{\lambda }{4}\left(\widehat{A}_{ss}\right)^2\right\}.`$ (114)
\[Here $``$ is a lattice forward difference operator, dot products denote the implicit sum over $`O(N)`$ indices, and factors of lattice spacing are suppressed for simplicity.\] The number $`m`$ of $`O(N)`$ vectors \[or the dimension of the matrices $`\widehat{A}`$, $`\widehat{B}`$, and $`\widehat{C}`$\] equals the total number of lattice sites. Ignoring the obvious notational changes ($`x\varphi `$, $`p\pi `$), this theory has precisely the stated form of Eq’s. (109)–(110). The lattice theory may, of course, be viewed as a natural discretization of the formal continuum theory where the field operators $`\widehat{\varphi }^i(x)`$ and $`\widehat{\pi }^i(x)`$ depend on continuous spatial coordinates and
$$N\widehat{H}_N=N(d^dx)\left[\frac{1}{2}\widehat{\pi }(x)\widehat{\pi }(x)+\frac{1}{2}\stackrel{}{}\widehat{\varphi }(x)\stackrel{}{}\widehat{\varphi }(x)+\frac{1}{2}\mu ^2\widehat{\varphi }(x)\widehat{\varphi }(x)+\frac{\lambda }{4}\left(\widehat{\varphi }(x)\widehat{\varphi }(x)\right)^2\right].$$
(115)
Returning to the general discussion, a straightforward calculation shows that the commutators of the basic bilinears are
$`\frac{N}{i}[\widehat{A}_{\alpha \beta },\widehat{A}_{\gamma \delta }]`$ $`=`$ $`\frac{N}{i}[\widehat{C}_{\alpha \beta },\widehat{C}_{\gamma \delta }]=0,`$ (117)
$`\frac{N}{i}[\widehat{A}_{\alpha \beta },\widehat{B}_{\gamma \delta }]`$ $`=`$ $`\widehat{A}_{\alpha \gamma }\delta _{\beta \delta }+\widehat{A}_{\beta \gamma }\delta _{\alpha \delta },`$ (118)
$`\frac{N}{i}[\widehat{B}_{\alpha \beta },\widehat{B}_{\gamma \delta }]`$ $`=`$ $`\widehat{B}_{\gamma \beta }\delta _{\alpha \delta }\widehat{B}_{\alpha \delta }\delta _{\gamma \beta },`$ (119)
$`\frac{N}{i}[\widehat{A}_{\alpha \beta },\widehat{C}_{\gamma \delta }]`$ $`=`$ $`\widehat{B}_{\alpha \gamma }\delta _{\beta \delta }+\widehat{B}_{\beta \gamma }\delta _{\alpha \delta }+\widehat{B}_{\alpha \delta }\delta _{\beta \gamma }+\widehat{B}_{\beta \delta }\delta _{\alpha \gamma },`$ (120)
$`\frac{N}{i}[\widehat{B}_{\alpha \beta },\widehat{C}_{\gamma \delta }]`$ $`=`$ $`\widehat{C}_{\beta \gamma }\delta _{\alpha \delta }+\widehat{C}_{\beta \delta }\delta _{\alpha \gamma }.`$ (121)
In other words, the commutators of $`\widehat{A}`$, $`\widehat{B}`$, and $`\widehat{C}`$ (as well as just $`\widehat{A}`$ and $`\widehat{B}`$) close and these operators generate a Lie algebra.<sup>\*††</sup><sup>\*††</sup>\*†† The Lie algebra structure constants are $`f_{A_{\alpha \beta }B_{\gamma \delta }}^{A_{\mu \nu }}=\frac{1}{2}\left(\delta _{\mu \alpha }\delta _{\nu \gamma }\delta _{\beta \delta }+\delta _{\mu \beta }\delta _{\nu \gamma }\delta _{\alpha \delta }+\delta _{\nu \alpha }\delta _{\mu \gamma }\delta _{\beta \delta }+\delta _{\nu \beta }\delta _{\mu \gamma }\delta _{\alpha \delta }\right)`$, $`f_{B_{\alpha \beta }B_{\gamma \delta }}^{B_{\mu \nu }}=\delta _{\mu \gamma }\delta _{\nu \beta }\delta _{\alpha \delta }\delta _{\mu \alpha }\delta _{\nu \delta }\delta _{\gamma \beta }`$, $`f_{A_{\alpha \beta }C_{\gamma \delta }}^{B_{\mu \nu }}=\delta _{\mu \alpha }\delta _{\nu \gamma }\delta _{\beta \delta }+\delta _{\mu \beta }\delta _{\nu \gamma }\delta _{\alpha \delta }+\delta _{\mu \alpha }\delta _{\nu \delta }\delta _{\beta \gamma }+\delta _{\mu \beta }\delta _{\nu \delta }\delta _{\alpha \gamma }`$, and $`f_{B_{\alpha \beta }C_{\gamma \delta }}^{C_{\mu \nu }}=\frac{1}{2}\left(\delta _{\mu \beta }\delta _{\nu \gamma }\delta _{\alpha \delta }+\delta _{\mu \beta }\delta _{\nu \delta }\delta _{\alpha \gamma }+\delta _{\nu \beta }\delta _{\mu \gamma }\delta _{\alpha \delta }+\delta _{\nu \beta }\delta _{\mu \delta }\delta _{\alpha \gamma }\right)`$, plus those trivially related by antisymmetry; all others vanish. The resulting Lie algebra of operators $`\{\widehat{\mathrm{\Lambda }}(a,b,c)iN_{\alpha \beta }(a_{\alpha \beta }\widehat{A}_{\beta \alpha }+b_{\alpha \beta }\widehat{B}_{\beta \alpha }+c_{\alpha \beta }\widehat{C}_{\beta \alpha })\}`$ is isomorphic to the $`sp(2m)`$ algebra represented by the $`2m`$-dimensional matrices $`\{\lambda (a,b,c)\left[\genfrac{}{}{0pt}{}{b}{a}\genfrac{}{}{0pt}{}{c}{b^T}\right]\}`$, where $`b=b_{\alpha \beta }`$, etc., and $`a`$ and $`c`$ are symmetric. The appropriate coherence group which will create suitable $`O(N)`$ invariant coherent states may be taken to be the group generated by (anti-Hermitian linear combinations of) the operators $`\{\widehat{A}_{\alpha \beta }\}`$ and $`\{\widehat{B}_{\alpha \beta }\}`$. Enlarging the coherence group by including the $`\widehat{C}_{\alpha \beta }`$ operators among the generators is equally acceptable, but unnecessary. The group generated by $`\{\widehat{A}_{\alpha \beta }\}`$ and $`\{\widehat{B}_{\alpha \beta }\}`$ alone satisfies all the conditions for producing an over-complete set of coherent states which behave classically as $`N\mathrm{}`$. Including the $`\widehat{C}_{\alpha \beta }`$ operators among the generators enlarges the coherence group, but has no effect whatsoever on the resulting manifold of coherent states.
Acting on an initial Gaussian base state, the coherence group generates a set of coherent states $`\{|z\}`$, where $`z`$ is a complex symmetric $`m\times m`$ matrix, with positive definite real part, which may be used to uniquely label an individual coherent state. The position space wavefunctions of these coherent states are given by
$$\mathrm{\Psi }_z(x)=det\left[\frac{N}{2\pi }(z+z^{})\right]^{N/4}\mathrm{exp}\left(\frac{1}{2}N\mathrm{tr}x^Tzx\right).$$
(122)
It will be convenient to decompose the matrix $`z`$ into its real and imaginary parts by writing
$$z=\frac{1}{2}a^1i\omega ,$$
(123)
so that $`a=(z+z^{})^1`$ and $`\omega =\frac{i}{2}(zz^{})`$. Both $`a`$ and $`\omega `$ are real symmetric matrices, and $`a`$ is positive definite. Using the fact that $`\widehat{p}_\beta ^i|z=i\widehat{x}_\alpha ^iz_{\alpha \beta }|z`$, a short exercise shows that the coherent state expectation values of the basic bilinears are
$$A(z)=a,B(z)=a\omega ,\text{and}C(z)=\omega a\omega +\frac{1}{4}a^1.$$
(124)
The variances of these operators in the coherent state $`|z`$ are<sup>\*‡‡</sup><sup>\*‡‡</sup>\*‡‡ For aesthetic reasons, we set $`g_{\alpha \beta ,\gamma \delta }^{AB}g_{A_{\alpha \beta }B_{\gamma \delta }}=A_{\alpha \beta }B_{\gamma \delta }A_{\alpha \beta }B_{\gamma \delta }`$, etc.
$`g_{\alpha \beta ,\gamma \delta }^{AA}`$ $`=`$ $`\frac{1}{N}\left[a_{\alpha \gamma }a_{\beta \delta }+a_{\alpha \delta }a_{\beta \gamma }\right],`$ (126)
$`g_{\alpha \beta ,\gamma \delta }^{BB}`$ $`=`$ $`\frac{1}{N}\left[a_{\alpha \gamma }(\frac{1}{4}a^1+\omega a\omega )_{\beta \delta }+(\frac{1}{2}+i\omega a)_{\beta \gamma }(\frac{1}{2}ia\omega )_{\alpha \delta }\right],`$ (127)
$`g_{\alpha \beta ,\gamma \delta }^{CC}`$ $`=`$ $`\frac{1}{N}[(\frac{1}{4}a^1+\omega a\omega )_{\alpha \gamma }(\frac{1}{4}a^1+\omega a\omega )_{\beta \delta }+(\frac{1}{4}a^1+\omega a\omega )_{\alpha \delta }(\frac{1}{4}a^1+\omega a\omega )_{\beta \gamma }],`$ (128)
$`\left(g_{\gamma \delta ,\alpha \beta }^{BA}\right)^{}=g_{\alpha \beta ,\gamma \delta }^{AB}`$ $`=`$ $`\frac{i}{N}\left[a_{\alpha \gamma }(\frac{1}{2}ia\omega )_{\beta \delta }+a_{\beta \gamma }(\frac{1}{2}ia\omega )_{\alpha \delta }\right],`$ (129)
$`\left(g_{\gamma \delta ,\alpha \beta }^{CA}\right)^{}=g_{\alpha \beta ,\gamma \delta }^{AC}`$ $`=`$ $`\frac{1}{N}\left[(\frac{1}{2}ia\omega )_{\alpha \gamma }(\frac{1}{2}ia\omega )_{\beta \delta }+(\frac{1}{2}ia\omega )_{\alpha \delta }(\frac{1}{2}ia\omega )_{\beta \gamma }\right],`$ (130)
$`\left(g_{\gamma \delta ,\alpha \beta }^{CB}\right)^{}=g_{\alpha \beta ,\gamma \delta }^{BC}`$ $`=`$ $`\frac{i}{N}\left[(\frac{1}{2}ia\omega )_{\alpha \gamma }(\frac{1}{4}a^1+\omega a\omega )_{\beta \delta }+(\frac{1}{2}ia\omega )_{\alpha \delta }(\frac{1}{4}a^1+\omega a\omega )_{\beta \gamma }\right].`$ (131)
Given our choice of Hamiltonian (110), the operator equations of motion for the basic bilinears are
$`{\displaystyle \frac{d}{dt}}\widehat{A}_{\alpha \beta }`$ $`=`$ $`\widehat{B}_{\alpha \beta }+\widehat{B}_{\beta \alpha },`$ (133)
$`{\displaystyle \frac{d}{dt}}\widehat{B}_{\alpha \beta }`$ $`=`$ $`\widehat{C}_{\alpha \beta }2\widehat{A}_{\alpha \gamma }\widehat{V}_{\gamma \beta }^{},`$ (134)
$`{\displaystyle \frac{d}{dt}}\widehat{C}_{\alpha \beta }`$ $`=`$ $`2\widehat{B}_{\gamma \alpha }\widehat{V}_{\gamma \beta }^{}2\widehat{V}_{\alpha \gamma }^{}\widehat{B}_{\gamma \beta }.`$ (135)
Here, $`V^{}`$ is shorthand for the variation of $`V(A)`$ with respect to the symmetric matrix $`A`$,
$$V_{\alpha \beta }^{}\frac{\delta V(A)}{\delta A_{(\alpha \beta )}}\frac{1}{2}\left[\frac{\delta V(A)}{\delta A_{\alpha \beta }}+\frac{\delta V(A)}{\delta A_{\beta \alpha }}\right],$$
(136)
and is defined so that $`\delta V(A)=\mathrm{tr}(V^{}\delta A)`$.<sup>†\*</sup><sup>†\*</sup>†\*Note that with this definition, the matrix variation $`V^{}`$ reduces to an ordinary variational derivative in the case of a single vector ($`m=1`$).
Applying the general results (39) and (45) \[actually, only (‡‡IV B) is needed\] to the case at hand, one finds in a straightforward fashion the following equations, valid to next-to-leading order in $`1/N`$, for the time evolution of the expectation values and variances of basic bilinears,
$`{\displaystyle \frac{d}{dt}}A_{\alpha \beta }`$ $`=`$ $`B_{\alpha \beta }+B_{\beta \alpha },`$ (138)
$`{\displaystyle \frac{d}{dt}}B_{\alpha \beta }`$ $`=`$ $`C_{\alpha \beta }2A_{\alpha \eta }V_{\eta \beta }^{}2g_{\alpha \eta ,\mu \nu }^{AA}V_{\eta \beta ,\mu \nu }^{\prime \prime }g_{\mu \nu ,\zeta \xi }^{AA}A_{\alpha \eta }V_{\eta \beta ,\mu \nu ,\zeta \xi }^{\prime \prime \prime }+𝒪(N^2),`$ (139)
$`{\displaystyle \frac{d}{dt}}C_{\alpha \beta }`$ $`=`$ $`2B_{\eta \alpha }V_{\eta \beta }^{}2B_{\eta \beta }V_{\eta \alpha }^{}\left(g_{\eta \alpha ,\mu \nu }^{BA}+g_{\mu \nu ,\eta \alpha }^{AB}\right)V_{\eta \beta ,\mu \nu }^{\prime \prime }\left(g_{\eta \beta ,\mu \nu }^{BA}+g_{\mu \nu ,\eta \beta }^{AB}\right)V_{\eta \alpha ,\mu \nu }^{\prime \prime }`$ (141)
$`g_{\mu \nu ,\zeta \xi }^{AA}\left(B_{\eta \alpha }V_{\eta \beta ,\mu \nu ,\zeta \xi }^{\prime \prime \prime }+B_{\eta \beta }V_{\eta \alpha ,\mu \nu ,\zeta \xi }^{\prime \prime \prime }\right)+𝒪(N^2),`$
together with
$`{\displaystyle \frac{d}{dt}}g_{\alpha \beta ,\gamma \delta }^{AA}`$ $`=`$ $`g_{\alpha \beta ,\gamma \delta }^{BA}+g_{\beta \alpha ,\gamma \delta }^{BA}+g_{\alpha \beta ,\gamma \delta }^{AB}+g_{\alpha \beta ,\delta \gamma }^{AB}+𝒪(N^2),`$ (143)
$`{\displaystyle \frac{d}{dt}}g_{\alpha \beta ,\gamma \delta }^{BB}`$ $`=`$ $`g_{\alpha \beta ,\gamma \delta }^{CB}+g_{\alpha \beta ,\gamma \delta }^{BC}2g_{\alpha \eta ,\gamma \delta }^{AB}V_{\eta \beta }^{}2g_{\alpha \beta ,\gamma \eta }^{BA}V_{\eta \delta }^{}`$ (146)
$`\left(g_{\mu \nu ,\gamma \delta }^{AB}+g_{\gamma \delta ,\mu \nu }^{BA}\right)A_{\alpha \eta }V_{\eta \beta ,\mu \nu }^{\prime \prime }`$
$`\left(g_{\mu \nu ,\alpha \beta }^{AB}+g_{\alpha \beta ,\mu \nu }^{BA}\right)A_{\gamma \eta }V_{\eta \delta ,\mu \nu }^{\prime \prime }+𝒪(N^2),`$
$`{\displaystyle \frac{d}{dt}}g_{\alpha \beta ,\gamma \delta }^{CC}`$ $`=`$ $`2g_{\eta \alpha ,\gamma \delta }^{BC}V_{\eta \beta }^{}2g_{\alpha \beta ,\eta \gamma }^{CB}V_{\eta \delta }^{}2g_{\eta \beta ,\gamma \delta }^{BC}V_{\alpha \eta }^{}2g_{\alpha \beta ,\eta \delta }^{CB}V_{\gamma \eta }^{}`$ (149)
$`\left(g_{\mu \nu ,\alpha \beta }^{AC}+g_{\alpha \beta ,\mu \nu }^{CA}\right)\left(B_{\eta \delta }V_{\gamma \eta ,\mu \nu }^{\prime \prime }+B_{\eta \gamma }V_{\eta \delta ,\mu \nu }^{\prime \prime }\right)`$
$`\left(g_{\mu \nu ,\gamma \delta }^{AC}+g_{\gamma \delta ,\mu \nu }^{CA}\right)\left(B_{\eta \beta }V_{\alpha \eta ,\mu \nu }^{\prime \prime }+B_{\eta \alpha }V_{\eta \beta ,\mu \nu }^{\prime \prime }\right)+𝒪(N^2),`$
$`{\displaystyle \frac{d}{dt}}g_{\alpha \beta ,\gamma \delta }^{AB}`$ $`=`$ $`g_{\alpha \beta ,\gamma \delta }^{BB}+g_{\beta \alpha ,\gamma \delta }^{BB}+g_{\alpha \beta ,\gamma \delta }^{AC}2g_{\alpha \beta ,\gamma \eta }^{AA}V_{\eta \delta }^{}2g_{\alpha \beta ,\mu \nu }^{AA}A_{\gamma \eta }V_{\eta \delta ,\mu \nu }^{\prime \prime }+𝒪(N^2),`$ (150)
$`{\displaystyle \frac{d}{dt}}g_{\alpha \beta ,\gamma \delta }^{AC}`$ $`=`$ $`g_{\alpha \beta ,\gamma \delta }^{BC}+g_{\beta \alpha ,\gamma \delta }^{BC}2g_{\alpha \beta ,\eta \delta }^{AB}V_{\gamma \eta }^{}2g_{\alpha \beta ,\eta \gamma }^{AB}V_{\eta \delta }^{}`$ (152)
$`2g_{\alpha \beta ,\mu \nu }^{AA}\left(B_{\eta \delta }V_{\gamma \eta ,\mu \nu }^{\prime \prime }+B_{\eta \gamma }V_{\eta \delta ,\mu \nu }^{\prime \prime }\right)+𝒪(N^2),`$
$`{\displaystyle \frac{d}{dt}}g_{\alpha \beta ,\gamma \delta }^{BC}`$ $`=`$ $`g_{\alpha \beta ,\gamma \delta }^{CC}2g_{\alpha \eta ,\gamma \delta }^{AC}V_{\eta \beta }^{}2g_{\alpha \beta ,\eta \delta }^{BB}V_{\gamma \eta }^{}2g_{\alpha \beta ,\eta \gamma }^{BB}V_{\eta \delta }^{}`$ (155)
$`\left(g_{\mu \nu ,\gamma \delta }^{AC}+g_{\gamma \delta ,\mu \nu }^{CA}\right)A_{\alpha \eta }V_{\eta \beta ,\mu \nu }^{\prime \prime }`$
$`\left(g_{\mu \nu ,\alpha \beta }^{AB}+g_{\alpha \beta ,\mu \nu }^{BA}\right)\left(B_{\eta \delta }V_{\gamma \eta ,\mu \nu }^{\prime \prime }+B_{\eta \gamma }V_{\eta \delta ,\mu \nu }^{\prime \prime }\right)+𝒪(N^2).`$
Here $`V_{\alpha \beta ,\gamma \delta }^{\prime \prime }\frac{\delta ^2V(A)}{\delta A_{(\alpha \beta )}\delta A_{(\gamma \delta )}}`$, etc. Of course, $`g_{\gamma \delta ,\alpha \beta }^{BA}=\left(g_{\alpha \beta ,\gamma \delta }^{AB}\right)^{}`$ and so on, since the basic bilinears $`\widehat{A}_{\alpha \beta }`$, $`\widehat{B}_{\alpha \beta }`$, and $`\widehat{C}_{\alpha \beta }`$ are all Hermitian.
As they stand, the (truncated) moment equations (V B) and (V B) are highly redundant. This is because the operators $`\widehat{A}_{\alpha \beta }`$, $`\widehat{B}_{\alpha \beta }`$, and $`\widehat{C}_{\alpha \beta }`$ are not independent when acting on the $`O(N)`$ invariant Hilbert space $`_N`$. For many purposes, it is preferable to reduce the evolution equations to a smaller set of independent observables. To see the redundancy, it is convenient first to note that the actions of $`\widehat{B}`$ and $`\widehat{A}`$ on any coherent state $`|z`$ are related,
$$\left(\widehat{B}+\frac{i}{2}\widehat{\mathrm{𝟏}}\right)|z=\widehat{x}\widehat{p}^T|z=\widehat{x}\widehat{x}^T(iz)|z=\widehat{A}(iz)|z.$$
(156)
Hence, the coherent state expectation value of $`\widehat{A}^1\widehat{B}`$ is directly related to that of $`\widehat{A}^1`$,<sup>††</sup><sup>††</sup>††The following discussion assumes that coherent state matrix elements of $`\widehat{A}^1`$ exist, which requires $`N>m+1`$.
$$z|\left(\widehat{A}^1\widehat{B}\right)_{\alpha \beta }|z=iz_{\alpha \beta }\frac{i}{2}z|\left(\widehat{A}^1\right)_{\alpha \beta }|z.$$
(157)
In a similar fashion, the coherent state expectation value of $`\widehat{C}`$ may be expressed as
$`z|\widehat{C}|z`$ $`=`$ $`z|\widehat{p}\widehat{p}^T|z=z|(iz)^{}(\widehat{x}\widehat{x}^T)(iz)|z=z|(\widehat{p}\widehat{x}^T)(\widehat{x}\widehat{x}^T)^1(\widehat{x}\widehat{p}^T)|z`$ (158)
$`=`$ $`z|(\widehat{B}+\frac{i}{2}\widehat{\mathrm{𝟏}})^{}\widehat{A}^1(\widehat{B}+\frac{i}{2}\widehat{\mathrm{𝟏}})|z.`$ (159)
As noted earlier in section II, quantum operators are completely determined by their diagonal expectation values in the over-complete coherent basis. Consequently, the coherent state relations (157) and (159) suffice to infer underlying operator identities. The left-hand side of relation (157) is not manifestly symmetric under interchange of $`\alpha `$ and $`\beta `$, but the right-hand side is symmetric under this interchange. Because Eq. (157) holds for all coherent states $`\{|z\}`$, if one defines
$$\widehat{\mathrm{\Omega }}_{\alpha \beta }\left(\widehat{A}^1\widehat{B}\right)_{\alpha \beta }+i\left(\frac{m+1}{2N}\right)\left(\widehat{A}^1\right)_{\alpha \beta },$$
(160)
then (157) implies that $`\mathrm{\Omega }_{\alpha \beta }=\mathrm{\Omega }_{\beta \alpha }`$, so that $`\mathrm{\Omega }=\mathrm{\Omega }_{\alpha \beta }`$ is a symmetric matrix. Moreover, using the the commutation relations (V B), one may verify that $`\mathrm{\Omega }_{\alpha \beta }`$ is Hermitian. \[Demanding Hermiticity is what determines the coefficient of the second term in (160).\] Similarly, relation (159) implies the operator identity
$$\widehat{C}=\left(\widehat{B}+\frac{i}{2}\widehat{\mathrm{𝟏}}\right)^{}\widehat{A}^1\left(\widehat{B}+\frac{i}{2}\widehat{\mathrm{𝟏}}\right),$$
(161)
showing that the operators $`\{\widehat{C}_{\alpha \beta }\}`$ are not independent of $`\widehat{A}`$ and $`\widehat{B}`$ \[when acting on $`O(N)`$ invariant states\]. Inverting the definition (160) to express $`\widehat{B}`$ in terms of $`\widehat{\mathrm{\Omega }}`$,
$$\widehat{B}=\widehat{A}\widehat{\mathrm{\Omega }}i\left(\frac{m+1}{2N}\right)\widehat{\mathrm{𝟏}},$$
(162)
and using this, plus the Hermiticity of $`\widehat{\mathrm{\Omega }}`$, allows one to rewrite expression (161) for $`\widehat{C}`$ as
$$\widehat{C}=\widehat{\mathrm{\Omega }}\widehat{A}\widehat{\mathrm{\Omega }}+\frac{1}{4}\left(1\frac{m+1}{N}\right)^2\widehat{A}^1.$$
(163)
Hence, within the $`O(N)`$ invariant Hilbert space, instead of working with the basic bilinears $`\widehat{A}`$, $`\widehat{B}`$, and $`\widehat{C}`$ \[totaling $`m(2m+1)`$ distinct operators\], it is sufficient to use only $`\widehat{A}`$ and $`\widehat{\mathrm{\Omega }}`$ \[totaling $`m(m+1)`$ distinct operators\]. These operators are, in fact, canonically conjugate “coordinates” and “momenta”. A short exercise shows that
$`[\widehat{A}_{\alpha \beta },\widehat{A}_{\gamma \delta }]`$ $`=`$ $`[\widehat{\mathrm{\Omega }}_{\alpha \beta },\widehat{\mathrm{\Omega }}_{\gamma \delta }]=0,`$ (165)
$`iN[\widehat{\mathrm{\Omega }}_{\alpha \beta },\widehat{A}_{\gamma \delta }]`$ $`=`$ $`\delta _{\alpha \gamma }\delta _{\beta \delta }+\delta _{\alpha \delta }\delta _{\beta \gamma }.`$ (166)
If the complex symmetric matrix $`z`$ parameterizing coherent states is separated into real and imaginary parts by writing $`z=\frac{1}{2}a^1i\omega `$ \[as in Eq. (123)\], then the coherent state expectations of the canonical operators $`\widehat{A}`$ and $`\widehat{\mathrm{\Omega }}`$ are just $`a`$ and $`\omega `$, respectively,
$$z|\widehat{A}|z=a,z|\widehat{\mathrm{\Omega }}|z=\omega .$$
(167)
\[The first equality was previously noted in Eq. (124).\]
Re-expressing the quantum equations of motion (V B) in terms of the independent canonically conjugate operators gives
$`{\displaystyle \frac{d}{dt}}\widehat{A}`$ $`=`$ $`\widehat{A}\widehat{\mathrm{\Omega }}+\widehat{\mathrm{\Omega }}\widehat{A},`$ (169)
$`{\displaystyle \frac{d}{dt}}\widehat{\mathrm{\Omega }}`$ $`=`$ $`\widehat{\mathrm{\Omega }}^22V_{\mathrm{eff}}^{}(\widehat{A}),`$ (170)
where the “effective” radial potential
$$V_{\mathrm{eff}}(A)V(A)+\frac{1}{8}\left(1\frac{m+1}{N}\right)^2\mathrm{tr}A^1$$
(171)
equals the original potential energy augmented by a “centrifugal potential”.
One may directly evaluate the evolution equations for expectations and variances of the canonically conjugate operators $`\widehat{A}`$ and $`\widehat{\mathrm{\Omega }}`$, or equivalently (and rather tediously) rewrite the previous equations (V B) and (V B) in terms of $`\widehat{A}`$ and $`\widehat{\mathrm{\Omega }}`$. Either way, one finds
$`{\displaystyle \frac{d}{dt}}A_{\alpha \beta }`$ $`=`$ $`(A\mathrm{\Omega }+\mathrm{\Omega }A)_{\alpha \beta }+g_{\alpha \eta ,\eta \beta }^{A\mathrm{\Omega }}+g_{\alpha \eta ,\eta \beta }^{\mathrm{\Omega }A}+𝒪(N^2),`$ (173)
$`{\displaystyle \frac{d}{dt}}\mathrm{\Omega }_{\alpha \beta }`$ $`=`$ $`(\mathrm{\Omega }^2+2V_{\mathrm{eff}}^{})_{\alpha \beta }g_{\alpha \eta ,\eta \beta }^{\mathrm{\Omega }\mathrm{\Omega }}\left(V_{\mathrm{eff}}^{\prime \prime \prime }\right)_{\alpha \beta ,\mu \nu ,\zeta \xi }g_{\mu \nu ,\zeta \xi }^{AA}+𝒪(N^2),`$ (174)
together with
$`{\displaystyle \frac{d}{dt}}g_{\alpha \beta ,\gamma \delta }^{AA}`$ $`=`$ $`g_{\alpha \eta ,\gamma \delta }^{AA}\mathrm{\Omega }_{\eta \beta }+g_{\eta \beta ,\gamma \delta }^{\mathrm{\Omega }A}A_{\alpha \eta }+g_{\alpha \eta ,\gamma \delta }^{\mathrm{\Omega }A}A_{\eta \beta }+g_{\eta \beta ,\gamma \delta }^{AA}\mathrm{\Omega }_{\alpha \eta }`$ (177)
$`+g_{\alpha \beta ,\gamma \eta }^{AA}\mathrm{\Omega }_{\eta \delta }+g_{\alpha \beta ,\eta \delta }^{A\mathrm{\Omega }}A_{\gamma \eta }+g_{\alpha \beta ,\gamma \eta }^{A\mathrm{\Omega }}A_{\eta \delta }+g_{\alpha \beta ,\eta \delta }^{AA}\mathrm{\Omega }_{\gamma \eta }+𝒪(N^2),`$
$`{\displaystyle \frac{d}{dt}}\left(g_{\gamma \delta ,\alpha \beta }^{\mathrm{\Omega }A}\right)^{}={\displaystyle \frac{d}{dt}}g_{\alpha \beta ,\gamma \delta }^{A\mathrm{\Omega }}`$ $`=`$ $`g_{\alpha \eta ,\gamma \delta }^{A\mathrm{\Omega }}\mathrm{\Omega }_{\eta \beta }+g_{\eta \beta ,\gamma \delta }^{\mathrm{\Omega }\mathrm{\Omega }}A_{\alpha \eta }+g_{\alpha \eta ,\gamma \delta }^{\mathrm{\Omega }\mathrm{\Omega }}A_{\eta \beta }+g_{\eta \beta ,\gamma \delta }^{A\mathrm{\Omega }}\mathrm{\Omega }_{\alpha \eta }`$ (179)
$`g_{\alpha \beta ,\eta \delta }^{A\mathrm{\Omega }}\mathrm{\Omega }_{\gamma \eta }g_{\alpha \beta ,\gamma \eta }^{A\mathrm{\Omega }}\mathrm{\Omega }_{\eta \delta }2\left(V_{\mathrm{eff}}^{\prime \prime }\right)_{\gamma \delta ,\mu \nu }g_{\alpha \beta ,\mu \nu }^{AA}+𝒪(N^2),`$
$`{\displaystyle \frac{d}{dt}}g_{\alpha \beta ,\gamma \delta }^{\mathrm{\Omega }\mathrm{\Omega }}`$ $`=`$ $`g_{\alpha \beta ,\eta \delta }^{\mathrm{\Omega }\mathrm{\Omega }}\mathrm{\Omega }_{\gamma \eta }g_{\alpha \beta ,\gamma \eta }^{\mathrm{\Omega }\mathrm{\Omega }}\mathrm{\Omega }_{\eta \delta }2\left(V_{\mathrm{eff}}^{\prime \prime }\right)_{\gamma \delta ,\mu \nu }g_{\alpha \beta ,\mu \nu }^{\mathrm{\Omega }A}`$ (181)
$`g_{\eta \beta ,\gamma \delta }^{\mathrm{\Omega }\mathrm{\Omega }}\mathrm{\Omega }_{\alpha \eta }g_{\alpha \eta ,\gamma \delta }^{\mathrm{\Omega }\mathrm{\Omega }}\mathrm{\Omega }_{\eta \beta }2\left(V_{\mathrm{eff}}^{\prime \prime }\right)_{\alpha \beta ,\mu \nu }g_{\mu \nu ,\gamma \delta }^{A\mathrm{\Omega }}+𝒪(N^2).`$
Initial conditions corresponding to a given coherent state $`|z`$ (with $`z=\frac{1}{2}a^1i\omega `$) are given by $`A(0)=a`$ and $`\mathrm{\Omega }(0)=\omega `$, together with the variances
$$\left(\begin{array}{cc}g_{\alpha \beta ,\gamma \delta }^{AA}& g_{\alpha \beta ,\gamma \delta }^{A\mathrm{\Omega }}\\ g_{\alpha \beta ,\gamma \delta }^{\mathrm{\Omega }A}& g_{\alpha \beta ,\gamma \delta }^{\mathrm{\Omega }\mathrm{\Omega }}\end{array}\right)_{t=0}=\frac{1}{N}\left(\begin{array}{cc}a_{\alpha \delta }a_{\beta \gamma }+a_{\alpha \gamma }a_{\beta \delta }& \frac{i}{2}[\delta _{\alpha \delta }\delta _{\beta \gamma }+\delta _{\alpha \gamma }\delta _{\beta \delta }]\\ \frac{i}{2}[\delta _{\alpha \delta }\delta _{\beta \gamma }+\delta _{\alpha \gamma }\delta _{\beta \delta }]& \frac{1}{4}[a_{\alpha \delta }^1a_{\beta \gamma }^1+a_{\alpha \gamma }^1a_{\beta \delta }^1]\end{array}\right)+𝒪(N^2).$$
(182)
The next-to-leading order evolution equations (V B) and (V B) are directly applicable to any bosonic $`O(N)`$ invariant vector model, such as the $`\varphi ^4`$ theory defined by ($`ii`$)), whose Hamiltonian has the general form (110). The dynamics is encoded in as efficient a form as possible; one has dynamical equations for the $`m(m+1)/2`$ pairs of independent phase space coordinates (V B), and their variances (V B).
In the special case (111) of a single vector (corresponding to a point particle moving in an $`N`$-dimensional spherically symmetric potential) one may drop all the indices and the next-to-leading order evolution equations become
$`{\displaystyle \frac{d}{dt}}A`$ $`=`$ $`A\mathrm{\Omega }+\mathrm{\Omega }A+g_{A\mathrm{\Omega }}+g_{\mathrm{\Omega }A}+𝒪(N^2),`$ (184)
$`{\displaystyle \frac{d}{dt}}\mathrm{\Omega }`$ $`=`$ $`\mathrm{\Omega }^22V_{\mathrm{eff}}^{}g_{\mathrm{\Omega }\mathrm{\Omega }}g_{AA}V_{\mathrm{eff}}^{\prime \prime \prime }+𝒪(N^2),`$ (185)
$`{\displaystyle \frac{d}{dt}}g_{AA}`$ $`=`$ $`4g_{AA}\mathrm{\Omega }+2(g_{\mathrm{\Omega }A}+g_{A\mathrm{\Omega }})A+𝒪(N^2),`$ (186)
$`{\displaystyle \frac{d}{dt}}(g_{\mathrm{\Omega }A})^{}={\displaystyle \frac{d}{dt}}g_{A\mathrm{\Omega }}`$ $`=`$ $`2g_{\mathrm{\Omega }\mathrm{\Omega }}A2g_{AA}V_{\mathrm{eff}}^{\prime \prime }+𝒪(N^2),`$ (187)
$`{\displaystyle \frac{d}{dt}}g_{\mathrm{\Omega }\mathrm{\Omega }}`$ $`=`$ $`4g_{\mathrm{\Omega }\mathrm{\Omega }}\mathrm{\Omega }2(g_{\mathrm{\Omega }A}+g_{A\mathrm{\Omega }})V_{\mathrm{eff}}^{\prime \prime }+𝒪(N^2),`$ (188)
with initial conditions given by $`A(0)=a`$, $`\mathrm{\Omega }(0)=\omega `$, and
$$\left(\begin{array}{cc}g_{AA}& g_{A\mathrm{\Omega }}\\ g_{\mathrm{\Omega }A}& g_{\mathrm{\Omega }\mathrm{\Omega }}\end{array}\right)_{t=0}=\frac{2}{N}\left(\begin{array}{cc}a^2& \frac{i}{2}\\ \frac{i}{2}& \frac{1}{4}a^2\end{array}\right)+𝒪(N^2).$$
(189)
From Eq’s. (V B) and (189) one may again see that to next-to-leading order, the determinant of the variance matrix on the left-hand side of (189) is a constant of the motion, $`detg^{(2)}(t)=𝒪\left(N^3\right)`$. To this order, our method gives exactly same predictions as the Gaussian approximation of . One may, of course, systematically extend the treatment to higher order in $`1/N`$ simply by specializing the next-to-next-to-leading order results in section IV.
The evolution equations (V B) in this single-vector case may be cast in a more transparent form by defining radial position and momentum operators via
$$\widehat{A}=\widehat{r}^2,\widehat{\mathrm{\Omega }}=\frac{1}{2}(\widehat{p}\widehat{r}^1+\widehat{r}^1\widehat{p}),$$
(190)
or equivalently
$$\widehat{r}=\widehat{A}^{1/2},\widehat{p}=\widehat{A}^{1/2}\widehat{\mathrm{\Omega }}\frac{i}{2N}\widehat{A}^{1/2}.$$
(191)
These operators are canonically conjugate,
$$i[\widehat{p},\widehat{r}]=1/N,$$
(192)
and a short exercise rewriting the quantum equations of motion (V B) yields
$`{\displaystyle \frac{d}{dt}}\widehat{r}`$ $`=`$ $`\widehat{p},`$ (194)
$`{\displaystyle \frac{d}{dt}}\widehat{p}`$ $`=`$ $`U_{\mathrm{eff}}^{}(\widehat{r}),`$ (195)
where
$`U_{\mathrm{eff}}(r)`$ $``$ $`V_{\mathrm{eff}}(r^2){\displaystyle \frac{1}{8N^2r^2}}`$ (196)
$`=`$ $`V(r^2)+\frac{1}{8}\left(1\frac{3}{N}\right)\left(1\frac{1}{N}\right)r^2,`$ (197)
and $`U_{\mathrm{eff}}^{}=dU_{\mathrm{eff}}/dr`$. This is a well-known result: $`s`$-wave dynamics in an $`N`$-dimensional central potential is equivalent to one-dimensional quantum dynamics in an effective radial potential $`U_{\mathrm{eff}}`$ containing an additional “centrifugal” potential $`\frac{(N3)(N1)}{8N^2r^2}`$ which is non-vanishing in all dimensions other than 1 and 3 . As seen in the commutation relations (192), the parameter $`1/N`$ plays the role of $`\mathrm{}`$ so that the large $`N`$ limit is precisely equivalent to the semiclassical limit of ordinary one-dimensional quantum mechanics.
The next-to-leading order evolution equations (V B) for the coherent state expectation values and variances of $`A`$ and $`\mathrm{\Omega }`$ may be easily be converted to equivalent next-to-leading order equations for expectations and variances of $`p`$ and $`r`$. One finds,
$`{\displaystyle \frac{d}{dt}}r`$ $`=`$ $`p+𝒪(N^2),`$ (199)
$`{\displaystyle \frac{d}{dt}}p`$ $`=`$ $`U_{\mathrm{eff}}^{}\frac{1}{2}g_{rr}U_{\mathrm{eff}}^{\prime \prime \prime }+𝒪(N^2),`$ (200)
$`{\displaystyle \frac{d}{dt}}g_{rr}`$ $`=`$ $`g_{rp}+g_{pr}+𝒪(N^2),`$ (201)
$`{\displaystyle \frac{d}{dt}}(g_{pr})^{}={\displaystyle \frac{d}{dt}}g_{rp}`$ $`=`$ $`g_{pp}g_{rr}U_{\mathrm{eff}}^{\prime \prime }+𝒪(N^2),`$ (202)
$`{\displaystyle \frac{d}{dt}}g_{pp}`$ $`=`$ $`(g_{rp}+g_{pr})U_{\mathrm{eff}}^{\prime \prime }+𝒪(N^2).`$ (203)
Through next-to-leading order, these evolutions equations are identical to the evolution equations (69) for the usual semiclassical limit.<sup>†‡</sup><sup>†‡</sup>†‡This equivalence persists to all orders, of course, reflecting the exact correspondence between the operator equations of motion (68) and (V B). The initial variances differ, however, due to the differing shapes of the initial wavepackets (67) and (122). For our large $`N`$ coherent states,
$$\left(\begin{array}{cc}g_{rr}& g_{rp}\\ g_{pr}& g_{pp}\end{array}\right)_{t=0}=\frac{1}{2N}\left(\begin{array}{cc}r^2& pr+i\\ pri& p^2+r^2\end{array}\right)+𝒪(N^2),$$
(204)
\[and once again $`detg^{(2)}(t)=𝒪\left(N^3\right)`$\]. The form of this variance matrix (including, for example, the growth in the variance $`g_{rr}`$ with increasing $`r`$) reflects the fact that the underlying $`O(N)`$ invariant coherent state wavefunctions are not constant width one-dimensional Gaussians, but rather $`N`$-dimensional Gaussians centered at the origin with variable width. Hence, the position of the peak in the resulting radial probability distribution is positively correlated with the width of the radial probability distribution about this peak.
For any given choice of the potential, one may integrate the five equations (V B) forward in time and obtain results which are accurate to $`𝒪(N^2)`$ \[for times of order unity\]. For better accuracy, one could extend the treatment to include higher order correlations, as detailed in section IV.
In light of the above exact correspondence between the $`O(N)`$-invariant dynamics of the single-vector model (111), and ordinary one-dimensional quantum dynamics in the the effective radial potential (197) with $`N`$ playing the role of $`\mathrm{}`$, the previous discussion of stability of the truncated moment equations in the semiclassical limit immediately carries over to the large $`N`$ dynamics of the single-vector model. In particular, this means that one should expect to see a decoherence time which scales as $`N^{1/2}`$, beyond which truncations of the moment hierarchy are no longer useful. We have no reason to believe that the scaling of the decoherence time with $`N`$ will be different in more general vector-like large $`N`$ theories, such as the $`\varphi ^4`$ field theory ($`ii`$)), as compared to the single-vector model. Although we have no compelling proof to offer, we expect that a decoherence time of order $`N^{1/2}`$ is a generic feature of vector-like large $`N`$ theories.<sup>†§</sup><sup>†§</sup>†§It is interesting to note that, in contrast to the previous discussion of the semiclassical $`\mathrm{}0`$ limit, examining $`N`$-dimensional free motion in the absence of any potential does not provide an example illustrating breakdown of the moment hierarchy based on $`O(N)`$ invariant coherent states. This is because the growth in the width of a spherically-symmetric Gaussian wavepacket is perfectly represented by a single one of the variable-width $`O(N)`$ invariant coherent states (122), unlike the earlier situation with fixed-width coherent states. Hence $`O(N)`$ invariant free motion is highly non-generic. For $`O(N)`$ invariant free motion (in the general case where $`z`$ is an $`m\times m`$ matrix and $`\widehat{H}_N=\frac{1}{2}\mathrm{tr}\widehat{C}`$), one may show that the exact time evolution maps an initial coherent state $`|z_0`$ into another coherent state $`|z(t)`$ with $`z(t)^1=z_0^1+it\mathbf{\hspace{0.17em}1}`$. The operator equations of motion (V B) may also be integrated exactly and show that $`\widehat{C}(t)=\widehat{C}(0)`$ is a constant of the motion, while $`\widehat{B}=\widehat{B}(0)+2\widehat{C}t`$, and $`\widehat{A}=\widehat{A}(0)+\frac{1}{2}\left[\widehat{B}(0)+\widehat{B}(0)^T\right]t+\widehat{C}t^2`$. This implies, for example, that for large time the variance $`g_{\alpha \beta ,\mu \nu }^{AA}t^4/N`$ and so grows without bound. However, the mean value $`\widehat{A}`$ grows quadratically with $`t`$, and hence rms fluctuations remain bounded and of order $`N^{1/2}`$ for all times.
## VI Conclusions
We have shown that a systematic hierarchy of time-local evolution equations for a minimal set of equal-time correlation functions may be derived in any theory having a classical (or large-$`N`$) limit which fits within the general framework of section II. Truncating this hierarchy at the level of $`k`$’th order moments (i.e., retaining up to $`k`$-point connected correlators) yields results which are accurate up to order $`1/N^k`$.
However, it is clear that the $`t\mathrm{}`$ and $`\mathrm{}0`$ (or $`N\mathrm{}`$) limits are non-uniform. At least in simple one degree of freedom (or single vector) models, we have argued that integrating the truncated moment evolution equations forward in time yields results which, generically, cease to be a good approximation to the true quantum dynamics beyond a decoherence time which scales as $`\mathrm{}^{1/2}`$ (or $`\sqrt{N}`$). The ordering of connected correlators which underlies the truncation of the moment hierarchy is only valid for times small compared to the decoherence time. Going to higher orders in the truncation scheme will not, in general, yield results which remain valid for parametrically longer time intervals.
We expect, but have not demonstrated, that this $`\sqrt{N}`$ scaling of the decoherence time is a general feature of large $`N`$ quantum dynamics. It would obviously be worthwhile to investigate this further, particularly in large $`N`$ models with many vectors. In, for example, an $`O(N)`$ invariant lattice $`\varphi ^4`$ field theory, it would clearly be desirable to understand the dependence of the decoherence time on the energy of the initial state and the lattice volume. If the $`\sqrt{N}`$ scaling of the decoherence time is generically true this would, for example, imply that one cannot use truncated hierarchies of large $`N`$ evolution equations (at least of the form considered here) to study the non-equilibrium dynamics of thermalization or hydrodynamic transport, as the relevant time scales for these processes scale like $`N`$ in the large $`N`$ limit . We hope that future work will shed light on these issues.
## Acknowledgment
One of the authors (LGY) wishes to thank Fred Cooper for asking the question which stimulated this work. A. Morozov was involved in initial portions of this investigation; his efforts are gratefully acknowledged.
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# Invited Paper presented at 16th International Conference on TRI–PP–00–27 Few-Body Problems in Physics, Taipei, March 6-10 May 2000 Symmetries and Symmetry Breaking 11footnote 1Work supported in part by the Natural Sciences and Engineering Council of Canada.
## 1 NUCLEON-NUCLEON PARITY VIOLATION
At low- and intermediate-energies, the parity violating weak $`N`$-$`N`$ interaction can be described in terms of a meson exchange model involving a strong interaction vertex and a weak interaction vertex (assuming one boson exchanges). The strong interaction vertex is well understood; it is represented by the conventional meson-exchange parameterization of the $`N`$-$`N`$ interaction. The weak interaction vertex is calculated from the Weinberg-Salam model assuming that the $`W`$\- and $`Z`$-bosons are exchanged between the intermediate mesons ($`\pi ,\rho `$, and $`\omega `$) and constituent quarks of the nucleon. The parity violating interaction can then be described in terms of seven weak meson-nucleon coupling constants. The six weak meson-nucleon coupling constants ($`f_\pi ^1,h_\rho ^0,h_\rho ^1,h_\rho ^2,h_\omega ^0,h_\omega ^1`$, with the superscripts indicating isospin changes and the subscripts the exchanged meson) have been calculated by Desplanques, Donoghue, and Holstein (DDH) , synthesizing the quark model and SU(6) and treating strong interaction effects in renormalization group theory. The seventh weak meson-nucleon coupling constant $`h_\rho ^1`$ is estimated to be smaller and is usually deleted from further consideration. DDH tabulated ‘best guess values’ and ‘reasonable ranges’ for the six weak meson-nucleon coupling constants. Following the seminal paper by DDH, various other groups have calculated the weak meson-nucleon coupling constants, but the considerable ranges of uncertainty attached to these remain (see Ref. 2). The parity violating $`\pi \mathrm{\Delta }N`$ vertex plays an increasingly important role in elastic and inelastic proton-proton scattering above the pion-production threshold. It is also apparent that the theoretical situation regarding $`f_\pi ^1`$ is not settled by any means. For a recent review see also Haeberli and Holstein .
A complete determination of the six weak meson-nucleon coupling constants demands a minimum of six experimental, linearly independent combinations of the weak meson-nucleon coupling constants. Of to date there do not exist enough experimental constraints of the required precision. This situation can only be remedied by performing a set of judiciously chosen, precise parity violation experiments.
Precise measurements of the $`p`$-$`p`$ parity violating longitudinal analyzing power have been made at 13.6 MeV $`[A_z=(0.93\pm 0.20\pm 0.05)\times 10^7]`$ at the University of Bonn and at 45 MeV $`[A_z=(1.53\pm 0.23)\times 10^7]`$ at SIN (now PSI) . Here $`A_z`$ is defined as $`A_z=(\sigma ^+\sigma ^{})/(\sigma ^++\sigma ^{}),\mathrm{with}\sigma ^+\mathrm{and}\sigma ^{}`$ representing the scattering cross sections for polarized incident protons of positive and negative helicity, respectively, integrated over the range of angles determined by the acceptance of the experimental apparatus in question. From the SIN measurement at 45 MeV and the $`\sqrt{E}`$ dependence of $`A_z`$ at lower energies, one can extrapolate that at 13.6 MeV $`A_z=(0.86\pm 0.13)\times 10^7`$. Consequently, the two precise, low energy measurements are in excellent agreement. Both results allow determining a combination of effective $`\rho `$ and $`\omega `$ weak meson-nucleon coupling constants: $`A_z`$ = 0.153h$`{}_{\rho }{}^{}{}_{}{}^{pp}`$ \+ 0.113h$`{}_{\omega }{}^{}{}_{}{}^{pp}`$, with h$`{}_{\rho }{}^{pp}=h_\rho ^0+h_\rho ^1+h_\rho ^2/\sqrt{6}`$ and h$`{}_{\omega }{}^{pp}=h_\omega ^0+h_\omega ^1`$. One should note that a measurement of $`A_z`$ in $`p`$-$`p`$ scattering is sensitive only to the short range part of the parity violating interaction (parity violating $`\pi ^0`$ exchange would be also CP violating and is therefore suppressed by a factor of about $`2\times 10^3`$).
A partial wave decomposition allows $`A_z`$ to be written as a sum of the various transition amplitudes ($`{}_{}{}^{1}S_{0}^{}{}_{}{}^{3}P_{0}^{}`$), ($`{}_{}{}^{3}P_{2}^{}`$ \- $`{}_{}{}^{1}D_{2}^{}`$), ($`{}_{}{}^{1}D_{2}^{}`$ \- $`{}_{}{}^{3}F_{2}^{}`$) , ($`{}_{}{}^{3}F_{4}^{}`$ \- $`{}_{}{}^{1}G_{4}^{}`$), etc. For incident energies below 100 MeV essentially only the first transition amplitude ($`{}_{}{}^{1}S_{0}^{}`$ \- $`{}_{}{}^{3}P_{0}^{}`$) is contributing, but for higher energies the second parity violating transition amplitude ($`{}_{}{}^{3}P_{2}^{}`$ \- $`{}_{}{}^{1}D_{2}^{}`$) is of increasing importance. The contribution of the next higher order (third) transition amplitude ($`{}_{}{}^{1}D_{2}^{}`$ \- $`{}_{}{}^{3}F_{2}^{}`$) is negligibly small. The TRIUMF 221.3 MeV $`p`$-$`p`$ parity violation experiment is unique in that it selected an energy where the ($`{}_{}{}^{1}S_{0}^{}`$ \- $`{}_{}{}^{3}P_{0}^{}`$) transition amplitude contribution integrates to zero, taking into account the angular acceptance of the apparatus . This is a reflection of the fact that the $`{}_{}{}^{1}S_{0}^{}`$ and $`{}_{}{}^{3}P_{0}^{}`$ phase shifts change sign near 230 MeV. Simonius has shown that the ($`{}_{}{}^{3}P_{2}^{}`$ \- $`{}_{}{}^{1}D_{2}^{})`$ transition amplitude depends only weakly on $`\omega `$ exchange (to an extent depending on the strong vector meson-nucleon coupling constants). Consequently, the TRIUMF 221.3 MeV experiment presents a determination of $`h_\rho ^{pp}`$. In the TRIUMF experiment a 200 nA proton beam with a polarization of 0.80 is incident on a 0.40 m long LH<sub>2</sub> target, after extraction from the optically pumped polarized ion source (OPPIS), passing a Wien filter in the injection beam line, acceleration through the cyclotron to an energy of 221.3 MeV, and multiturn extraction by a stripping foil. A combination of solenoid-dipole-solenoid-dipole magnets on the external beam line provides a longitudinally polarized beam with either positive or negative helicity or vice versa. $`A_z`$ follows from the helicity dependence of the $`p`$-$`p`$ total cross section as determined in precise measurements of the normalized transmission asymmetry through the 0.40 long LH<sub>2</sub> target: $`A_z=(1/P)(T/S)(T^+T^{})/(T^++T^{})`$, where $`P`$ is the incident beam longitudinal polarization, $`T=1S`$ is the average transmission through the target, and the + and - signs indicate the helicity. Very strict constraints are imposed on the incident longitudinally polarized beam in terms of intensity, transverse horizontal (x) and vertical (y) beam position and direction, beam width ($`\sigma _x`$ and $`\sigma _y`$), longitudinal polarization ($`P_z`$), transverse polarization ($`P_x`$ and $`P_y`$), first moments of the transverse polarization ($`<xP_y>`$ and $`<yP_x>`$), and energy, together with deviations of the transmission measuring apparatus from cylindrical symmetry. Helicity correlated modulations in the beam parameters originate at OPPIS, but can be amplified by the beam transport through the injection beam line, the cyclotron accelerator, and the extraction beam line. Residual systematic errors arising from the imperfections of the incident beam and of the response of the transmission measuring apparatus, are individually not to exceed one-tenth of the expected value of $`A_z`$ (or $`6\times 10^9`$). Particular troublesome are the first moments of residual transverse polarization, as well as helicity correlated changes in energy. The approach which has been followed is to further measure the sensitivity or response to residual imperfections, to monitor these imperfections during data taking, and to make corrections where necessary. Random changes in the incident beam parameters cause a dilution of the effect to be measured and therefore necessitate longer data taking times. Details about the experimental arrangements and procedures can be found in Ref. 2, while an account of the data analysis will be reported elsewhere. Fig. 1 presents the three lower energy results in a comparison with the meson exchange model theoretical predictions of Driscoll and Miller and Iqbal and Niskanen , the chiral soliton model prediction of Driscoll and Meissner , and the quark model prediction of Grach and Shmatikov . Note that the first two predictions have used the DDH weak meson-nucleon coupling constants. For the TRIUMF 221.3 MeV experiment one can derive that $`A_z=0.0296h_\rho ^{pp}`$. The three lower energy measurements establish the (expected) energy behaviour and allow a delineation of $`h_\rho ^{pp}`$ and $`h_\omega ^{pp}`$.
Fig. 1 Theoretical predictions of Driscoll and Miller , Iqbal and Niskanen , Driscoll and Meissner , and Grach and Shmatikov compared to the low-energy $`p`$-$`p`$ parity violating longitudinal analyzing powers $`A_z`$.
There exist two further higher energy parity violation experiments. The first one is a $`p`$-$`p`$ parity violation measurement at 800 MeV with $`A_z=(2.4\pm 1.1)\times 10^7`$ at LANL. Its interpretation in terms of the effective $`\rho `$ and $`\omega `$ weak meson-nucleon coupling constants is more difficult due to the presence of a large inelasticity (pion production). The second one is a $`pN`$ parity violation measurement at 5.13 GeV on a water target with $`A_z=(26.4\pm 6.0\pm 3.6)\times 10^7`$ at ANL with the ZGS. This result is an order of magnitude larger than what is expected based upon using simple scaling arguments. New $`p`$-$`p`$ parity violation experiments are being planned at TRIUMF possibly at 450 MeV and with COSY at the Forschungszentrum Jülich near 2 GeV as a storage ring experiment.
Other $`N`$-$`N`$ parity violation measurements have dealt with the circular polarization $`P_\gamma `$ of the $`\gamma `$-rays in $`n`$-$`p`$ capture and with the asymmetry $`A_\gamma `$ in $`n`$-$`p`$ capture with polarized cold neutrons, as well as the inverse reaction, deuteron photodisintegration with circularly polarized $`\gamma `$-rays. However, these experiments were not of enough statistical precision to have an impact on the determination of $`f_\pi ^1`$. A new measurement of $`A_\gamma `$ is being prepared at LANSCE, aiming at a ten-fold improvement in accuracy (to a precision of $`\pm 0.5\times 10^8`$, which will determine $`f_\pi ^1`$ to $`\pm 0.4\times 10^7)`$. In the experiment, neutrons from the pulsed spallation source are moderated by a LH<sub>2</sub> moderator, and their energy determined by time-of-flight. The cold neutrons are polarized by transmission through polarized <sup>3</sup>He gas; the neutron spin direction can be subsequently reversed by a rf resonance flipper. The neutrons are then guided to a liquid para-hydrogen target which is surrounded by an array of $`\gamma `$-ray detectors. Similarly, a new measurement of the asymmetry in photodisintegration of the deuteron with circularly polarized photons, obtained when a 8 MeV polarized electron beam from the CEBAF injector is incident on a gold Bremsstrahlung target, is being developed at Jefferson Laboratory. It is planned to measure $`f_\pi ^1`$ to an accuracy better than 30% of the current theoretical ‘best value’ of $`f_\pi ^1`$. The latter two experiments are designed to resolve the experimental discrepancy between the value of $`f_\pi ^1`$ as deduced from measurements of the circularly polarized 1.081 MeV $`\gamma `$-rays from the well known parity-mixed doublet in <sup>18</sup>F and as deduced from the measurements of the anapole moment
Fig. 2 Plot of the constraints on the isoscalar and isovector ($`f_\pi ^1`$) weak meson-nucleon coupling constants.
of <sup>133</sup>Cs. A plot of the constraints on the isoscalar and isovector weak meson-nucleon coupling constants, obtained from the more precise low-energy parity violation data, is shown in Fig. 2.
## 2 PARITY-EVEN/TIME-REVERSAL-ODD INTERACTION
Time-reversal-invariance non-conservation has for the first time been unequivocally demonstrated in a direct measurement in the CPLEAR experiment. The experiment measured the difference in the transition probabilities $`P(\overline{K}^0K^0)`$ and $`P(K^0\overline{K}^0)`$. Assuming CPT conservation but allowing for a possible breaking of the $`\mathrm{\Delta }S=\mathrm{\Delta }Q`$ rule, the result obtained for $`A_T=[R(\overline{K}^0K^0)R(K^0\overline{K}^0)]/[R(\overline{K}^0K^0)+R(K^0\overline{K}^0)]=[6.6\pm 1.3(\mathrm{stat}.)\pm 1.0(\mathrm{syst}.)]\times 10^3`$ is in good agreement with the measure of CP violation in neutral kaon decay. A more recently reported result is a large asymmetry in the distribution of $`K_L\pi ^+\pi ^{}e^+e^{}`$ events in the CP-odd/T-odd angle $`\varphi `$ between the decay planes of the $`\pi ^+\pi ^{}`$ and $`e^+e^{}`$ pairs in the $`K_L`$ center of mass system. The overall asymmetry found was $`[13.6\pm 2.5(\mathrm{stat}.)\pm 1.2(\mathrm{syst}.)`$\]%. This raises the question about time-reversal-invariance non-conservation in systems other than the kaon system.
Tests of time-reversal-invariance can be distinguished as belonging to two classes: the first one deals with P-odd/T-odd interactions, while the second one deals with P-even/T-odd interactions (assuming CPT conservation this implies C-conjugation non-conservation). But it is to be noted that constraints on these two classes of interactions are not independent since the effects due to P-odd/T-odd interactions may also be produced by P-even/T-odd interactions in conjunction with Standard Model parity violating radiative corrections. The latter can occur at the $`10^7`$ level and may present a limit on the constraint of a P-even/T-odd interaction derived from experiment. Limits on a P-odd/T-odd interaction follow from measurements of the electric dipole moment (edm) of the neutron (which currently stands at $`<6\times 10^{26}`$ e.cm \[95% C.L.\]). This provides a limit on a P-odd/T-odd pion-nucleon coupling constant which is less than 10<sup>-4</sup> times the weak interaction strength. Measurements of <sup>129</sup>Xe and <sup>199</sup>Hg edm’s ($`<8\times 10^{28}`$ e.cm \[95% C.L.\]) give similar constraints. \[see Ref. 20\]
Experimental limits on a P-even/T-odd interaction are much less stringent. Following the conventional approach of describing the $`N`$-$`N`$ interaction in terms of meson exchanges, it can be shown that only charged rho-meson exchange and a<sub>1</sub>-meson exchange can lead to a P-even/T-odd interaction. The better constraints stem first from measurements of the edm of the neutron and second from measurements of charge symmetry breaking (CSB) in $`n`$-$`p`$ elastic scattering. Haxton, Hoering, and Ramsay-Musolf have deduced constraints on a P-even/T-odd interaction from nucleon, nuclear, and atomic edm’s with the better constraint coming from the measurement of the edm of the neutron. In terms of a ratio to the strong rho-meson nucleon coupling constant, they deduced for the P-even/T-odd rho-meson nucleon coupling: $`|\overline{g}_\rho |`$ $`<0.53\times 10^3\times |f_\pi ^{\mathrm{DDH}}/f_\pi ^{\mathrm{meas}.}|`$. But as indicated above there exists great uncertainty about the value of $`f_\pi ^1`$; the ratio of the theoretical to experimental value of $`f_\pi ^1`$ may be as large as 15! However, constraints derived from one-loop contributions to the edm of the neutron exceed the two-loop limits by more than an order of magnitude and are much more stringent. It is to be noted that a translation in terms of coupling strengths in the hadronic sector still needs to be made.
It is very difficult to accommodate a P-even/T-odd interaction in the Standard Model. It requires C-conjugation non-conservation, which cannot be introduced at the first generation quark level. It can neither be introduced into the gluon self-interaction. Consequently, one needs to consider C-conjugation non-conservation between quarks of different generations and/or between interacting fields.
Charge symmetry breaking in $`n`$-$`p`$ elastic scattering manifests itself as a non-zero difference of the neutron ($`A_n`$) and proton ($`A_p`$) analyzing powers, $`\mathrm{\Delta }A=A_nA_p=2\times [Re(b^{}f)+Im(c^{}h)]/\sigma _0`$. Here the complex amplitude $`f`$ is charge symmetry breaking, while the complex amplitude $`h`$ is both charge symmetry breaking and time-reversal-invariance non-conserving. The complex amplitudes $`b`$ and $`c`$ belong to the usual five $`n`$-$`p`$ scattering amplitudes and $`\sigma _0`$ is the unpolarized differential cross section. The three precision experiments performed (at TRIUMF at 477 MeV and at 347 MeV , and at IUCF at 183 MeV ) have unambiguously shown that charge symmetry is broken and that the results for $`\mathrm{\Delta }A`$ at the zero-crossing of the average analyzing are very well reproduced by meson exchange model calculations. (see Fig. 3) A P-even/T-odd interaction introduces a term in the scattering amplitude which is simultaneously charge symmetry breaking (the complex amplitude $`h`$ in the above expression). Thus, Simonius deduced an upper limit on a P-even/T-odd interaction from a comparison of the three experimental results with the theoretical predictions. The upper limit so derived is $`|\overline{g}_\rho |<6.6\times 10^3`$ \[95% C.L.\]. This result is therefore comparable to the upper limit deduced from the edm of the neutron, taking the current experimental value of $`f_\pi ^1`$ extracted from <sup>18</sup>F, and is considerably lower than the limits inferred from direct tests of a P-even/T-odd interaction. Even though it is inconceivable in the Standard Model to account for a P-even/T-odd interaction, there is a need to clarify the experimental situation by providing a better experimental result.
Fig. 3 Experimental results of $`\mathrm{\Delta }A`$ at the zero-crossing at incident neutron energies of 183, 347, 477 MeV compared with theoretical predictions of Iqbal and Niskanen, and Holzenkamp, Holinde, and Thomas. The inner error bars present the statistical uncertainties; the outer error bars have the systematic uncertainties included (added in quadrature). For details see Ref. 25.
A better experimental constraint may be provided by an improved upper limit on the electric dipole moment of the neutron. Indeed a new measurement with a sensitivity of $`4\times 10^{28}`$ e.cm has been proposed at LANSCE. This would constitute a more than two orders of magnitude improvement over the present upper limit. Performing an improved $`n`$-$`p`$ elastic scattering CSB experiment also appears to be an attractive possibility. One can calculate with a great deal of accuracy the contributions to CSB due to one-photon exchange and due to the $`n`$-$`p`$ mass difference affecting one-pion and rho-meson exchange. Furthermore, one can select an energy where the $`\rho ^0\omega `$ meson mixing contribution changes sign at the same angle where the average analyzing power changes sign and therefore does not contribute to $`\mathrm{\Delta }A`$. This occurs at an incident neutron energy of 320 MeV and is caused by the particular interplay of the $`n`$-$`p`$ phase shifts and the form of the spin/isospin operator connected with the $`\rho ^0\omega `$ mixing term. Also the one-photon exchange term at 320 MeV changes sign at about the same angle as the average of the analyzing powers. The contribution due to two-pion exchange with an intermediate $`\mathrm{\Delta }`$ is expected to be less than one tenth of the overall $`\mathrm{\Delta }A`$, essentially determining an upper limit on the theoretical uncertainty. (see Fig. 4) It has been shown that simultaneous $`\gamma \pi `$ exchanges can only contribute to $`\mathrm{\Delta }A`$ through second order processes and can therefore be neglected. At 320 MeV the effects of inelasticity (pion production) are negligibly small. It appears therefore well within reach to reduce the theoretical uncertainty in the comparison of theory with experiment. Both the statistical and systematic errors, obtained in the 347 MeV TRIUMF experiment, can be considerably improved upon (by a factor three to four). With the developments of optically pumped polarized ion sources which have taken place in the intervening years it will be possible to obtain up to 50 $`\mu A`$ of 342 MeV 80% polarized proton beam on the neutron production target (a factor of 50 increase in neutron beam
Fig. 4 Angular distributions of the different contributions to $`\mathrm{\Delta }A`$ at an incident neutron energy of 320 MeV. (Ref. 29) Note that the $`\rho ^0\omega `$ mixing contribution passes through zero at the same angle as the average of $`A_n`$ and $`A_p`$ (vertical bars). The figure on the right gives the total $`\mathrm{\Delta }A`$ angular distribution.
intensity at 320 MeV over the previous 347 MeV $`n`$-$`p`$ CSB experiment). In addition various systematic error reducing improvements can be introduced. Such an experiment would constitute a measurement of CSB in $`n`$-$`p`$ elastic scattering of unprecedented precision of great value of its own and would simultaneously provide a greatly improved upper limit on a P-even/T-odd interaction.
## 3 ELECTRON-PROTON PARITY VIOLATION
The structure of the nucleon at low energies in terms of the quark and gluon degrees of freedom is not well understood. Of particular interest are the two proton ground state matrix elements which are sensitive to point-like “strange” quarks and hence to the quark-antiquark sea in the proton. The two matrix elements of interest are the elastic scattering vector weak neutral current ‘charge’ and ‘magnetic’ form factors, $`G_E^Z`$ and $`G_M^Z`$, respectively. These form factors can be deduced from parity violating electron-proton elastic scattering measurements. Assuming charge symmetry, i.e., the proton and neutron differ only by the interchange of the “up” and “down” quarks, one can determine the “up”, “down”, and ”strange” quark contributions to the ‘charge’ and ‘magnetic’ form factors of the nucleon. These contributions would result from taking the appropriate linear combinations of the weak neutral form factors and their electromagnetic counterparts.
Determinations of both the ‘charge’ and ‘magnetic’ “strange” quark form factors, $`G_E^s`$ and $`G_M^s`$, would constitute the first direct information on the quark sea in low energy observables. Electron-proton parity violation experiments are to determine these contributions to the proton form factors at the few percent level. High energy experiments suggest that the “strange” quarks carry about half as much momentum as the “up” and “down” quarks in the sea. The matrix elements, $`G_E^Z`$ and $`G_M^Z`$, are also of relevance to discussions of the Ellis-Jaffe sum rule and the $`\pi N`$ sigma term; there is uncertainty in both of these about the “strange” quark contributions. The quantity to be measured is the longitudinal analyzing power $`A_z`$, which is defined completely analogous to the $`p`$-$`p`$ one. Making pairs of measurements at forward and backward angles will allow the separation of $`G_E^Z`$ and $`G_M^Z`$. Predicted analyzing powers are in the $`10^6`$ to $`10^5`$ range.
Various electron-proton parity violation experiments have been performed or are being prepared for the near future. The SAMPLE experiment at the MIT-Bates Linear Accelerator detected backward scattered electrons in large air Čerenkov detectors in 200 MeV elastic e-p and quasielastic e-d scattering. The Čerenkov detectors subtended the laboratory angular range from 130 to 170, corresponding to a four momentum transfer $`Q^2`$ of about 0.1 (GeV/c)<sup>2</sup>. The value obtained in $`e`$-$`p`$ scattering of $`A_z`$ = (-4.92 $`\pm `$ 0.61 $`\pm `$ 0.73) $`\times `$ 10<sup>-6</sup> results in $`G_M^s=0.45G_A^Z+0.20\pm 0.17(\mathrm{stat}.)\pm 0.21(\mathrm{syst}.)`$. Taking theoretical estimates for the axial form factor $`G_A^Z`$ leads to a substantially positive $`G_M^s`$, however the preliminary value obtained in $`e`$-$`d`$ scattering does not corroborate the estimated value for $`G_A^Z`$ and consequently the isoscalar and isovector axial radiative corrections, $`R_A^0`$ and $`R_A^1`$, used in obtaining the estimate for $`G_A^Z`$ may be in error. Note that the isovector radiative correction has a connection to hadronic parity violation. The HAPPEX experiment at Jefferson Laboratory detected forward scattered electrons in the two Hall-A HRS spectrometers placed left and right of the incident beam at 12.5, corresponding to a $`Q^2`$ of 0.47 (GeV/c)<sup>2</sup> in 3.335 GeV elastic $`e`$-$`p`$ scattering. The latter of the two data taking runs used strained GaAs crystals to give an electron beam with about 70% polarization. The experiment measured the combination $`G_E^s+0.39G_M^s`$. The result from the first data taking run is $`A_z=(14.5\pm 2.0\pm 1.1)\times 10^6`$, which gives $`G_E^s+0.39G_M^s=0.023\pm 0.034(\mathrm{stat}.)\pm 0.022(\mathrm{syst}.)\pm 0.026G_E^n`$. Taking current information on $`G_E^n`$ this is essentially a null-result. A new round of experiments to measure $`G_E^n`$ at Jefferson Laboratory will remove the remaining uncertainty. The preliminary result of the second data taking run is $`A_z=(14.6\pm 1.1\pm 0.6)\times 10^6`$ in excellent agreement with the result of the first data taking run.
The A4 experiment at MAMI will detect forward scattered electrons in 855 MeV parity violating elastic $`e`$-$`p`$ scattering at 35 in a cylindrical calorimeter made of 1022 PbF<sub>2</sub> crystals. The experiment will determine a linear combination of $`G_E^s+0.22G_M^S`$ at a $`Q^2`$ value of 0.23 (GeV/c)<sup>2</sup>. The $`G^0`$ experiment at Jefferson Laboratory is the most comprehensive effort to date to measure both $`G_E^s`$ and $`G_M^s`$ over the range of $`Q^20.11.0`$ (GeV/c)<sup>2</sup>. Forward and backward angle parity violating elastic $`e`$-$`p`$ and quasielastic $`e`$-$`d`$ will be measured. In the forward mode of operation, protons scattered in a polar angular range of $`\pm 10^{}`$ around 70 will be detected in eightfold symmetry around the incident beam axis. In the backward mode of operation electrons scattered around 108 will be detected. Parity violating quasielastic $`e`$-$`d`$ scattering is again required to determine the axial form factor $`G_A^Z`$ contribution. A custom designed superconducting toroidal spectrometer with eightfold symmetry is being constructed. The scattered particles are detected in segmented scintillator arrays in the focal plane of the spectrometer. Commissioning running is scheduled for late 2001. The errors anticipated to be obtained for $`G_E^s`$ and $`G_M^s`$ are shown in Fig. 5.
Fig. 5 Anticipated errors on $`G_E^s`$ and $`G_M^s`$ as function of $`Q^2`$ from the $`G^0`$ experiment. Note that the proton electric and magnetic form factors are divided by a factor 10. The SAMPLE experiment result $`G_M^s=0.61\pm 0.17\pm 0.21`$ n.m. at $`Q^2`$ = 0.1 (GeV/c)<sup>2</sup> is also shown.
## 4 SUMMARY
Fundamental symmetries are tested with unheard precision leading to insight in the underlying structure and interactions. Subtle effects are observed in which the hadronic weak interaction plays a prominent role.
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# On Embedding Circle-Bundles in Four-Manifolds
## 1. Introduction
Let $`Y(n,g)`$ denote the circle bundle over a genus $`g`$ surface with Euler number $`n`$. Our main result in this paper is the following:
###### Theorem 1.1.
If $`X`$ is a complex surface of general type, and $`Y(n,g)`$ is circle-bundle over a Riemann surface of genus $`g`$, whose Euler number $`n`$ satisfies $`|n|2g1`$, then $`X`$ admits no splittings along an embedded copy of $`Y=Y(n,g)`$ of the form $`X=X_1\mathrm{\#}_YX_2`$ with $`b_2^+(X_1),b_2^+(X_2)>0`$.
In the above theorem, the quantity $`b_2^+(Z)`$ of a four-manifold $`Z`$ with boundary denotes the maximal dimension of a positive-definite subspace for the intersection form on $`H^2(Z,Z;)`$. It is suggestive to compare the hypothesis that $`|n|2g1`$ with the “adjunction inequality” for surfaces of non-negative square (see or ). Indeed, the hypotheses of Theorem 1.1 are sharp: if we allow $`b_2^+(X_2)=0`$ or $`|n|2g2`$, there are many examples of such splittings, obtained by blowing up smoothly embedded complex curves $`C`$ in $`X`$, and splitting $`X`$ along the boundary of the tubular neighborhood of $`C`$.
Moreover, the situation for elliptic surfaces is quite different, as we see below:
###### Theorem 1.2.
* Every simply-connected elliptic surface with $`b_2^+(X)>3`$ admits a splitting along $`Y(1,1)`$ with $`b_2^+(X_i)>0`$.
* For each $`n>0`$, there is a simply-connected elliptic surface $`X`$ which admits a splitting along $`Y(n,1)`$ with $`b_2^+(X_i)>0`$.
Proof. A splitting of Type (1) is given as follows. Note that $`Y(1,1)`$ is the mapping cylinder of a (single) Dehn twist on the torus. Thus, if we begin with the rational elliptic surface
$$\pi :E(1)=^2\mathrm{\#}9\overline{}^2^1,$$
let $`x^1`$ denote the image of a fishtail fiber, and let $`\mathrm{\Delta }`$ be a disk around $`x`$ containing no other singular points for $`\pi `$, then $`\pi ^1(\mathrm{\Delta })=Y(1,1)`$ splits $`E(1)`$ into a pair of elliptic fibrations $`Z_1`$ and $`Z_2`$ over disks. Thus we can realize $`E(3)`$ as a union of fiber sums $`E(1)\mathrm{\#}Z_1`$ and $`Z_2\mathrm{\#}E(1)`$ joined along $`Y(1,1)`$, where $`E(3)`$ is the fiber sum of three copies of $`E(1)`$. Neither side is negative-definite: both sides contain a torus of square zero and a sphere (constructed from vanishing cycles) which meets this torus in a single, positive point. Since every simply-connected elliptic surface with $`b_2^+3`$ can be obtained from $`E(3)`$ by fiber sums with $`E(1)`$, logarithmic transformations, and blow-ups (see , or ) the result follows.
A splitting of Type (2) is realized by finding an elliptic surface $`Z`$ over $`^1`$ which contains $`n`$ singular values for the elliptic fibration whose holonomy is a Dehn twist along a given curve in the fiber. In fact, it is a theorem of Moishezon (see , also Theorem 3.6 in Chapter 2 of ) that if $`Z`$ is a nodal elliptic surface without multiple fibers and $`2m`$ singular fibers, then we can think of the monodromy representation around $`m`$ of the singular fibers, of which we select $`n`$, as being a ($`+1`$) Dehn twist around a fixed non-separating curve in the fiber, and the monodromy around the remaining $`m`$ as being a Dehn twist around another curve. Let $`\mathrm{\Delta }`$ be a disk in $`^1`$ which contains only the $`n`$ distinguished singular points and no others. Now, it is easy to see that $`\pi ^1(\mathrm{\Delta })=Y(n,1)`$, which separates the elliptic surface. Forming fiber sums with rational elliptic surfaces on both sides as before, we get a decomposition of the elliptic surface $`E(1)\mathrm{\#}Z\mathrm{\#}E(1)`$ along $`Y(n,1)`$ into two pieces with $`b_2^+>0`$.
###### Remark 1.3.
Note that the hypothesis that $`b_2^+(Z)>3`$ above is necessary: the elliptic surface $`E(2)`$ admits no decomposition along $`Y(n,g)`$ with $`g2n1`$ and $`b_2^+(X_i)>0`$. This follows from the fact that $`E(2)`$ has a single basic class, together with the vanishing result, Theorem 2.1, from Section 2.
###### Remark 1.4.
Using the above decomposition (Type 2) as a building block, it is possible to construct symplectic four-manifolds $`Z`$ which decompose along $`Y(n,g)`$ with $`n`$ and $`g`$ arbitrarily large, such that both sides have positive $`b_2^+`$. For example, one can start with an elliptic surface $`X`$ decomposed along $`Y(n,1)`$ in the manner of Theorem 1.2, and find a symplectic torus $`TX`$ (which is symplectic for a form arbitrarily close to a Kähler form for $`X`$) which meets $`Y(n,1)`$ in a fiber circle for the Seifert fibration of $`Y(n,1)`$, and has square zero. Forming the fiber sum of $`X`$ with, say, $`T^2\times \mathrm{\Sigma }_{g1}`$ (by gluing $`TX`$ to $`T^2\times p`$), we obtain $`Z`$ as claimed.
Theorem 1.1 follows from a “vanishing theorem,” Theorem 2.1, according to which a certain sum of Seiberg-Witten invariants for $`X`$ vanishes whenever $`X`$ splits into two pieces with $`b_2^+(X_i)>0`$ along $`Y(n,g)`$, when $`|n|2g1`$. This is a more refined vanishing statement than the usual vanishing theorem over $`S^3`$: in particular there are manifolds with non-trivial Seiberg-Witten invariants which admit such splittings, as is illustrated by Theorem 1.2. The vanishing theorem is proved by looking at the ends of the moduli spaces of flows to the reducibles: this is also the philosophy adopted by Austin and Braam in , see also . In the case where $`g=1`$, it is interesting to compare the vanishing theorem with a certain vanishing theorem for Donaldson polynomials proved by Morgan, Mrowka, and Ruberman (Theorem 16.0.1 of ).
We will give the proof of Theorem 1.1 in Section 2, after stating and proving the more general vanishing result on which it is based.
## 2. The Vanishing Theorem
To state the vanishing theorem which implies Theorem 1.1, we must introduce some notation. We think of the Seiberg-Witten invariant of a smooth, oriented, closed four-manifold $`X`$ (with a “homology orientation” – an orientation on $`(H^0H^1H^+)(X;)`$) and $`\mathrm{Spin}^c`$ structure $`𝔰`$ as a homogenous polynomial map
$$SW_{(X,𝔰)}:𝔸(X)$$
of degree
$$d(𝔰)=\frac{c_1(𝔰)^2(2\chi +3\sigma )}{4}$$
on the algebra
$$𝔸(X)=[U]_{}\mathrm{\Lambda }^{}H_1(X;),$$
where $`U`$ is a two-dimensional generator, and $`\mathrm{\Lambda }^{}H_1(X;)`$ is the exterior algebra on the first homology of $`X`$ (graded in the obvious manner). This algebra maps surjectively to the cohomology ring of the irreducible configuration space $`^{}(X,𝔰)`$ of pairs $`[A,\mathrm{\Phi }]`$ of $`\mathrm{Spin}^c`$ connections $`A`$ and somewhere non-vanishing spinors $`\mathrm{\Phi }`$ modulo gauge. (We denote the full configuration space of pairs modulo gauge, i.e. where $`\mathrm{\Phi }`$ is allowed to vanish, by $`(X,𝔰)`$.) As usual, the Seiberg-Witten invariant is obtained by cohomological pairings of these cohomology classes with the fundamental cycle of the moduli space $`(X,𝔰)`$ of solutions to the Seiberg-Witten equations, which is naturally induced from the homology orientation.
As in Section 1, let $`Y=Y(n,g)`$ be the circle bundle over a Riemann surface with Euler number $`n`$ over a surface $`\mathrm{\Sigma }`$ of genus $`g>0`$. Throughout this section, we assume that
$$|n|2g1.$$
Recall that $`H^2(Y;)^{2g}(/n)`$, where the $`/n`$ factor is generated by multiples of the pull-back $`\pi ^{}`$ of the orientation class of $`\mathrm{\Sigma }`$. Indeed, there is a canonical $`\mathrm{Spin}^c`$ structure $`𝔱_0`$ over $`Y`$ associated to the two-plane field orthogonal to the circle directions. Thus, forming the tensor product with $`𝔱_0`$ gives a canonical identification
$$\mathrm{Spin}^c(Y)H^2(Y;).$$
In particular, there are $`n`$ distinguished $`\mathrm{Spin}^c`$ structures $`𝔱_e`$ over $`Y`$, indexed by $`e/n`$ (thought of as a subset of $`H^2(Y;)`$).
###### Theorem 2.1.
Let $`X`$ be a smooth, closed, oriented four-manifold which splits along an embedded copy of $`Y=Y(n,g)`$ with $`|n|>2g1`$, so that $`X=X_1\mathrm{\#}_YX_2`$ with $`b_2^+(X_i)>0`$ for $`i=1,2`$. Fix a $`\mathrm{Spin}^c`$ structure $`𝔰`$ on $`X`$, and let $`𝔰|_Y=𝔱`$. If $`𝔱`$ is not one of the $`n`$ distinguished $`\mathrm{Spin}^c`$ structures on $`Y`$, then $`SW_{(X,𝔰)}0`$. Similarly, if $`𝔱=𝔱_e`$ for $`2g2<e<n`$, then $`SW_{(X,𝔰)}0`$. Otherwise, if $`𝔱=𝔱_e`$ for $`i=0,\mathrm{},2g2`$, we have that
(1)
$$\underset{\{𝔰^{}|𝔰^{}|_{X_1}=𝔰|_{X_1},𝔰^{}|_{X_2}=𝔰|_{X_2}\}}{}SW_{(X,𝔰^{})}0.$$
Note that the inclusion $`YX`$ gives rise to a coboundary map $`\delta :H^1(Y;)H^2(Y;)`$, whose image we denote by $`\delta H^1(Y;)`$. Another way of stating Equation (1) is:
$$\underset{\eta \delta H^1(Y;)}{}SW_{(X,𝔰+\eta )}0.$$
The above theorem is proved by considering the ends of certain moduli spaces over cylindrical-end manifolds. In general, these ends are described in terms of the moduli spaces of the boundary $`Y`$, and the moduli spaces of solutions on the cylinder $`\times Y`$ (using a product metric and perturbation).
Specifically, let $`Y`$ be a three-manifold, and let $`𝒩_Y(𝔱)`$ denote the moduli space of solutions to the three-dimensional Seiberg-Witten equations over $`Y`$ in the $`\mathrm{Spin}^c`$ structure $`𝔱`$. Given a pair of components $`C_1`$, $`C_2`$ in $`𝒩_Y(𝔱)`$, let $`(C_1,C_2)`$ denote the moduli space of solutions $`[A,\mathrm{\Phi }]`$ to the Seiberg-Witten equations on $`\times Y`$ for which
$`\underset{t\mathrm{}}{lim}[A,\mathrm{\Phi }]|_{\{t\}\times Y}C_1`$ and $`\underset{t\mathrm{}}{lim}[A,\mathrm{\Phi }]|_{\{t\}\times Y}C_2`$
The theory of can be adapted to give the moduli space $`(C_1,C_2)`$ a Fredholm deformation theory, and a pair of continuous “boundary value maps” for $`i=1,2`$
$$\rho _{_{C_i}}:(C_1,C_2)C_i$$
characterized by
$`\rho _{_{C_1}}[A,\mathrm{\Phi }]=\underset{t\mathrm{}}{lim}[A,\mathrm{\Phi }]|_{\{t\}\times Y}`$ and $`\rho _{_{C_2}}[A,\mathrm{\Phi }]=\underset{t+\mathrm{}}{lim}[A,\mathrm{\Phi }]|_{\{t\}\times Y}.`$
The moduli space $`(C_1,C_2)`$ admits a translation action by $``$, and we let $`\widehat{}(C_1,C_2)`$ denote the quotient of this space by this action. The boundary value maps are $``$-invariant, and hence induce boundary value maps on the quotient
$$\rho _{_{C_i}},:\widehat{}(C_1,C_2)C_i.$$
As in (by analogy with the cases considered by Floer, see for instance ), the solutions to the three-dimensional Seiberg-Witten equations are the critical points for a “Chern-Simons-Dirac” functional $`\mathrm{CSD}`$ defined on the configuration space $`(Y,𝔱)`$. The Seiberg-Witten equations on $`\times Y`$ can be naturally identified with upward gradient flowlines for this functional. (Strictly speaking, the functional $`\mathrm{CSD}`$ is real-valued only when the first Chern class $`c_1(𝔱)`$ is torsion; otherwise it is circle-valued.)
Solutions in $`𝒩(Y,𝔱)`$ whose spinor vanishes identically correspond to flat connections on the determinant line bundle for $`𝔱`$. By analogy with the Donaldson-Floer theory, these solutions are usually called reducibles, and those with somewhere non-vanishing spinor are called irreducibles.
In the case where $`Y`$ is a non-trivial circle bundle over a Riemann surface these moduli spaces were studied in , see also (where $`Y`$ is endowed with a circle-invariant metric and the Seiberg-Witten equations over it are suitably perturbed). Specifically, there is the following result:
###### Theorem 2.2.
Let $`Y`$ be a circle bundle over a Riemann surface with genus $`g>0`$, and Euler number $`|n|>2g2`$ (oriented as circle bundle with negative Euler number). Then, the moduli space $`𝒩_Y(𝔱)`$ is empty unless $`𝔱`$ corresponds to a torsion class in $`H^2(Y;)`$. Suppose that $`𝔱=𝔱_e`$ for $`e/nH^2(Y;)`$. Then
* If $`0eg1`$ then $`𝒩_Y(𝔱)`$ contains two components, a reducible one $`𝒥`$, identified with the Jacobian torus $`H^1(\mathrm{\Sigma };S^1)`$, and a smooth irreducible component $`C`$ diffeomorphic to $`\mathrm{Sym}^e(\mathrm{\Sigma })`$. Both of these components are non-degenerate in the sense of Morse-Bott. There is an inequality $`\mathrm{CSD}(𝒥)>\mathrm{CSD}(C)`$, so the space $`\widehat{}(𝒥,C)`$ is empty. The space $`\widehat{}(C,𝒥)`$ is smooth of expected dimension $`2e`$; indeed it is diffeomorphic to $`\mathrm{Sym}^e(\mathrm{\Sigma })`$ under the boundary value map
$$\rho __C:\widehat{}(C,𝒥)C\mathrm{Sym}^e(\mathrm{\Sigma }).$$
* If $`g1<e2g2`$, the Seiberg-Witten moduli spaces over both $`Y`$ and $`\times Y`$ in this $`\mathrm{Spin}^c`$ structure are naturally identified with the corresponding moduli spaces in the $`\mathrm{Spin}^c`$ structure $`2g2e`$, which we just described.
* For all other $`e`$, $`𝒩_Y(𝔱)`$ contains only reducibles. Furthermore, it is smoothly identified with the Jacobian torus $`𝒥`$.
When $`eg1`$, the above theorem is a special case of Theorems 1 and 2 of (see especially Corollary 1.5 of ). When $`e=g1`$, the case considered in that paper is not “generic”. In fact, there is a natural perturbation (by some small multiple of the connection $`1`$-form of the Seifert fibration), which achieves the genericity stated above. This perturbation was used in to prove strong “adjunction inequalities” for manifolds which are not of simple type, and the above theorem in the case where $`e=g1`$ is precisely Theorem 8.1 of . Note that the hypothesis $`n>2g2`$ is required to separate the irreducible manifolds into distinct $`\mathrm{Spin}^c`$ structures. Note also that if the orientation on $`Y`$ is reversed, the flow-lines reverse direction.
The proof of Theorem 2.1 is obtained by considering the ends of the moduli spaces $`(X_1,𝔰_1,𝒥)`$ of Seiberg-Witten monopoles over the cylindrical-end manifold
$$X_1^+=X_1_{X_1=\{0\}\times Y}[0,\mathrm{})\times Y$$
in the $`\mathrm{Spin}^c`$ structure $`𝔰_1`$, whose boundary values are reducible. We will assume, as in that theorem, that $`b_2^+(X_1)>0`$. In general, moduli spaces of finite energy solutions to the Seiberg-Witten equations on a manifold with cylindrical ends are not compact. (The “finite energy condition” in this context is equivalent to the hypothesis that the pair $`[A,\mathrm{\Phi }]`$ has a well-defined boundary value.) They do, however, have “broken flowline” compactifications (see and ). In particular, if $`C`$ is a component of $`𝒩(Y,𝔰_1|_Y)`$, then for generic perturbations, the moduli space $`(X_1,𝔰_1,C)`$ is a smooth manifold with finitely many ends indexed by components $`C_1,\mathrm{},C_n`$ in the moduli space $`𝒩(Y,𝔰_1|_Y)`$, with $`\mathrm{CSD}(C_1)<\mathrm{CSD}(C_2)<\mathrm{}<\mathrm{CSD}(C_n)<\mathrm{CSD}(C)`$. When all the $`C_i`$ are non-degenerate in the sense of Morse-Bott, and consist of irreducibles, a neighborhood of the corresponding end is diffeomorphic to the fibered product
$$(X_1,𝔰_1,C_1)\times _{C_1}(C_1,C_2)\times _{C_2}\mathrm{}\times _{C_n}(C_n,C),$$
under a certain gluing map (provided that this space is a manifold – i.e. provided that the various boundary value maps are transverse).
In particular, suppose $`X_1`$ is an oriented four-manifold with boundary, whose boundary $`X_1`$ is identified with $`Y=Y(n,g)`$ with the orientation described in Theorem 2.2. Then, it follows from that theorem that if $`𝔰_1|_Y=𝔱_e`$ for $`0e2g2`$, then $`(X_1,𝔰_1,C)`$ is compact (since there are no “intermediate” critical manifolds to be added), and $`(X_1,𝔰_1,𝒥)`$ has a single end whose neighborhood is diffeomorphic to
$$(X_1,𝔰_1,C)\times (0,\mathrm{}).$$
(We use here the fact that the restriction map $`\rho __C:\widehat{}(C,𝒥)C`$ is a diffeomorphism.) This gluing map
$$\gamma :(X_1,𝔰_1,C)\times (0,\mathrm{})(X_1,𝒥)$$
is compatible with restriction to compact subsets of $`X_1^+`$; e.g. if we consider the compact subset $`X_1X_1^+`$, then
$$\underset{T\mathrm{}}{lim}\gamma ([A,\mathrm{\Phi }],T)|_{X_1}=[A,\mathrm{\Phi }]|_{X_1}.$$
We make use of the end of $`(X_1,𝔰_1,𝒥)`$ in the following proposition. Recall that the moduli space $`(X_1,𝔰_1,C)`$ is a smooth, compact submanifold of the irreducible configuration space of $`X_1^+`$. It has a canonical top-dimensional homology class, denoted $`[(X_1,𝔰_1,C)]`$, induced from the “homology orientation” of $`X_1`$. It inherits cohomology classes by pulling back via the boundary value map
$$\rho __C:(X_1,𝔰_1,C)$$
and from the natural map
$$i_{X_1}:(X_1,𝔰)^{}(X_1,𝔰)$$
given by restricting the pair $`[A,\mathrm{\Phi }]`$ to the compact subset $`X_1X_1^+`$ (this restriction is irreducible from the unique continuation theorem for the Dirac operator). The pairings with these classes can be thought of as a “relative Seiberg-Witten” invariant.
###### Proposition 2.3.
Suppose $`b_2^+(X_1)>0`$. Given any cohomology classes $`aH^{}(^{}(X_1,𝔰_1))`$ and $`bH^{}(C)`$, the homology class $`[(X_1,𝔰_1,C)]`$ pairs trivially with the class $`i_{X_1}^{}(a)\rho __C^{}(b)`$.
Proof.
First, we reduce to the case where $`b`$ is absent (i.e. zero-dimensional). This is done in two steps, first establishing an inclusion
(2)
$$(i_Y\rho __C)^{}H^{}(^{}(Y,𝔱))i_{X_1}^{}(^{}(X_1,𝔰_1)),$$
where both are thought of as subsets of $`H^{}((X_1,𝔰_1,C))`$, and then seeing that the map
$$i_Y^{}:H^{}(^{}(Y,𝔱))H^{}(C)$$
is surjective.
To see Inclusion (2) we describe the geometric representatives for the generators of the cohomology ring
$$H^{}(^{}(Y,𝔱))[U]_{}\mathrm{\Lambda }^{}(H_1(Y;)).$$
Given a point $`yY`$ and a line $`\mathrm{\Lambda }_yW_y`$ in the fiber of the spinor bundle over $`y`$, the class $`U`$ is Poincaré dual to the locus $`V_{(y,\mathrm{\Lambda }_y)}`$ of pairs $`[B,\mathrm{\Psi }]`$ with $`\mathrm{\Psi }_y\mathrm{\Lambda }_y`$. Moreover, given a curve $`\gamma :S^1X`$, the corresponding one-dimensional cohomology class determined by the homotopy type of the map
$$h_\gamma :^{}(Y,𝔱)S^1$$
given by measuring the holonomy of $`B`$ (relative to some fixed reference connection $`B_0`$) around $`\gamma `$; i.e. it is Poincaré dual to the preimage $`V_\gamma `$ of a regular value of $`h_\gamma `$. This cohomology class is denoted $`\mu [\gamma ]H^1(^{}(Y,𝔱))`$. (Note that geometric representatives cohomology classes in the configuration spaces of four-manifolds are constructed in an analogous manner.)
Now, fix a curve $`\gamma Y`$ and consider the one-parameter family of maps
$$h_t:(X_1,𝔰_1,C)S^1$$
indexed by $`t(0,1]`$ defined by measuring the holonomy of $`A`$ around the curve $`\{1/t\}\times \gamma X_1^+`$. Since the configurations in $`[A,\mathrm{\Phi }](X_1,𝔰_1,C)`$ converge exponentially to a stationary solution (see ), $`h_t`$ extends continuously to $`t=0`$. Now, $`h_0`$ represents $`(i_Y\rho __C)^{}`$ of the one-dimensional class $`\mu [\gamma ]H^1(^{}(Y,𝔱))`$, while $`h_1`$ represents the restriction (to the moduli space) of the one-dimensional class $`\mu [\gamma _1]H^{}(^{}(X_1,𝔰_1))`$, where $`\mu _1=\{1\}\times YX_1`$. A similar discusion applies to the two-dimensional class to show that $`\rho __C^{}U=i_{X_1}^{}U`$ (now we use the connection $`A`$ to identify the fiber “at infinity” with the fiber at some point inside $`X`$). This completes the verification of Inclusion (2).
Surjectivity of
$$i_Y^{}:H^{}(^{}(Y,𝔱);)𝔸(Y)H^{}(C;)$$
follows from classical properties of the cohomology of symmetric products $`\mathrm{Sym}^e(\mathrm{\Sigma })`$ (see ), according to which the cohomology ring is generated by “symmetrizations” of the cohomology of $`\mathrm{\Sigma }`$. It is then a straightforward verification (which is spelled out in Proposition 6.10 of ) using the geometric interpretations of the cohomology classes given above to see that that $`i_Y^{}\mu (\gamma )`$ corresponds to the symmetrization of $`\pi (\gamma )`$, while $`i_Y^{}U`$ corresponds to the symmetrization of the point $`\pi (x)`$ on $`\mathrm{\Sigma }`$ (where we think of $`U`$ as the Poincaré dual of $`V_{(x,\mathrm{\Lambda })}`$ for some choice of line $`\mathrm{\Lambda }W_x`$).
Thus, it remains to prove that $`[(X_1,𝔰_1,C)]`$ pairs trivially with classes in
$$i_{X_1}^{}H^{}(^{}(X_1,𝔰_1)).$$
We can think of the cohomological pairing $`\rho __C^{}(b),[(X_1,𝔰_1,C)]`$ as counting the (signed) number of points to the Seiberg-Witten equations, which satisfy constraints in the compact subset $`X_1X_1^+`$; i.e. if $`b=U^d[\mu _1]\mathrm{}[\mu _{\mathrm{}}]`$, where $`\mu _i`$ are curves in $`X_1`$, and $`x_1,\mathrm{},x_d`$ are generic points in $`X_1`$, and $`\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_d`$ are generic lines $`\mathrm{\Lambda }_iW^+|_{\{x_i\}}`$, then we have gemetric representatives $`V_{(x_i,\mathrm{\Lambda }_i)}`$ and $`V_{\mu _i}`$ for these cohomology classes, so that
$$[(X_1,𝔰_1)],b=\mathrm{\#}(X_1)V,$$
where $`V=V_{(x_1,\mathrm{\Lambda }_1)}\mathrm{}V_{(x_d,\mathrm{\Lambda }_d)}V_{\mu _1}\mathrm{}V_\mu _{\mathrm{}}`$. In fact, if we consider the solutions $`(X_1,𝔰_1,𝒥)`$ which satisfy these same constraints, then we get a manifold of dimension one with ends corresponding to $`(X_1,𝔰_1,C)V`$. Thus, counting boundary points with sign, we see that
$$\mathrm{\#}(X_1,𝔰_1,C)V=0.$$
Proof of Theorem 2.1. In the splitting $`X=X_1_YX_2`$, we can number the sides so that the boundary of $`X_1`$ is $`Y`$ oriented as in Theorem 2.2. Let $`𝔰`$ be a $`\mathrm{Spin}^c`$ structure over $`X`$, and let $`𝔰_i\mathrm{Spin}^c(X_i)`$ denote the restriction $`𝔰_i=𝔰|_{X_i}`$ for $`i=1,2`$, and let $`𝔱\mathrm{Spin}^c(Y)`$ denote the restriction $`𝔱=𝔰|_Y`$. Let $`X(T)`$ denote the Riemannian structure on $`X`$ obtained by inserting a cylinder $`[T,T]\times Y`$ between $`X_1`$ and $`X_2`$ (but keeping the metrics on these two pieces to be fixed, and product-like near the boundary). If the Seiberg-Witten invariants for a $`\mathrm{Spin}^c`$ structure $`𝔰`$ over $`X`$ is non-trivial, for any unbounded, increasing sequence of real numbers $`\{T_i\}`$, there must be a sequence of Seiberg-Witten monopoles $`[A_i,\mathrm{\Phi }_i](X(T_i),𝔰)`$. The uniform bound in the energy, and local compactness (see ) allows one to find a sequence $`\{t_i\}`$ with $`t_iT_i`$, so that, after passing to a subsequence if necessary, $`[A_i,\mathrm{\Phi }_i]|_{\{t_i\}\times Y}`$ converges to a stationary solution; i.e. it converges to a point in $`𝒩(Y,𝔱)`$. Thus, from Theorem 2.2, it follows that the Seiberg-Witten invariant for $`𝔰`$ vanishes unless $`𝔱`$ is one of the $`n`$ distinguished $`\mathrm{Spin}^c`$ structures $`𝔱_e`$ over $`Y`$.
Suppose that $`𝔱=𝔱_e`$, for $`e=0,\mathrm{},2g`$. Note the excision principle for the index gives that
$$d(𝔰)=dim(X_1,𝔰_1,C)+dim(X_2,𝔰_2,C)dim(C)$$
for any $`\mathrm{Spin}^c`$ structure $`𝔰\mathrm{Spin}^c(X)`$ with $`𝔰|_{X_i}=𝔰_i`$ for $`i=1,2`$ (and generic, compactly supported perturbations of the equations over $`X_i`$). We fix an integer $`\mathrm{}0`$ and homology classes $`a_1,\mathrm{},a_mH_1(X_1;)`$, $`b_1,\mathrm{},b_nH_1(X_2;)`$ with
$$2\mathrm{}+m+n=d(𝔰).$$
Let $`𝔰_1\mathrm{\#}𝔰_2\mathrm{Spin}^c(X)`$ denote the subset of $`\mathrm{Spin}^c`$ structures on $`X`$:
$$𝔰_1\mathrm{\#}𝔰_2=\{𝔰^{}\mathrm{Spin}^c(X)|𝔰^{}|_{X_1}=𝔰_1,𝔰^{}|_{X_2}=𝔰_2\},$$
and let $`(X(T),𝔰_1\mathrm{\#}𝔰_2)`$ denote the union
$$(X(T),𝔰_1\mathrm{\#}𝔰_2)=\underset{𝔰^{}𝔰_1\mathrm{\#}𝔰_2}{}(X,𝔰^{}).$$
Clearly, we have that
$$\mathrm{\#}(X(T),𝔰_1\mathrm{\#}𝔰_2)V_1V_2=\underset{𝔰^{}𝔰_1\mathrm{\#}𝔰_2}{}SW_{X,𝔰^{}}(U^{\mathrm{}}[a_1]\mathrm{}[a_m][b_1]\mathrm{}[b_n])$$
where $`V_i`$ are the intersection of the constraints from the $`X_i`$ side; e.g.
$$V_1=V_{(x_1,\mathrm{\Lambda }_1)}\mathrm{}V_{(x_{\mathrm{}},\mathrm{\Lambda }_m)}V_{a_1}\mathrm{}V_{a_m}$$
and
$$V_2=V_{b_1}\mathrm{}V_{b_n}.$$
Thus, our aim is to prove that the total signed number of points in the cut-down moduli space $`(X(T),𝔰_1\mathrm{\#}𝔰_2)V_1V_2`$ is zero.
Given pre-compact sets $`K_i(X_i,𝔰_i,C)`$ for $`i=1,2`$, there are gluing maps defined for all sufficiently large $`T`$,
$$\gamma _{C;T}:K_1\mathrm{\#}_CK_2(X(T),𝔰_1\mathrm{\#}𝔰_2),$$
where the domain is the fibered product of $`K_1`$ and $`K_2`$ over $`\rho _1`$ and $`\rho _2`$, i.e. the set of $`[A_1,\mathrm{\Phi }_1]K_1,[A_2,\mathrm{\Phi }_2]K_2`$ with
$$\rho _1([A_1,\mathrm{\Phi }_1])=\rho _2([A_2,\mathrm{\Phi }_2]),$$
and the range consists of all configurations $`[A,\mathrm{\Phi }]`$ which are whose restrictions to $`X_1`$ and $`X_2`$ are sufficiently close to restrictions (to $`X_1`$ and $`X_2`$) of configurations $`[A_1,\mathrm{\Phi }]\times [A_2,\mathrm{\Phi }_2]`$ in the fibered product.
We claim that for all sufficiently large $`T`$, the cut-down moduli space lies in the range of this map. Specifically, if we had a sequence $`[A_i,\mathrm{\Phi }_i](X(T_i),𝔰_1\mathrm{\#}𝔰_2)`$ for an increasing, unbounded sequence $`\{T_i\}_{i=1}^{\mathrm{}}`$ of real numbers, the sequence converges $`C^{\mathrm{}}`$ locally to give a pair of Seiberg-Witten monopoles monopoles $`[A_1,\mathrm{\Phi }_1](X_1,𝔰_1)`$ and $`[A_2,\mathrm{\Phi }_2](X_2,𝔰_2)`$. These monopoles have finite energy (since the total variation of $`\mathrm{CSD}`$ is bounded in the limit), so they have boundary values, which must lie in either $`C`$ or $`𝒥`$. We exclude all but one of the four cases as follows:
* The case where $`\rho __1[A_1,\mathrm{\Phi }_1]𝒥`$ and $`\rho __2[A_2,\mathrm{\Phi }_2]C`$ is excluded since $`\mathrm{CSD}(C)<\mathrm{CSD}(𝒥)`$.
* The case where $`\rho __1[A_1,\mathrm{\Phi }_1]𝒥`$ and $`\rho __2[A_2,\mathrm{\Phi }_2]𝒥`$ is excluded by a dimension count. Specifically, we must have that
$$\rho __1[A_1,\mathrm{\Phi }_1]=\rho __2[A_2,\mathrm{\Phi }_2]$$
and $`[A_1,\mathrm{\Phi }_1](X_1,𝒥)V_1`$ and $`[A_2,\mathrm{\Phi }_2](X_2,𝒥)V_2`$, i.e. the pair $`[A_1,\mathrm{\Phi }_1]\times [A_2,\mathrm{\Phi }_2]`$ lies in the fibered product $`(X_1,𝔰_1,𝒥)\times _𝒥(X_2,𝔰_2,𝒥)`$, a space whose dimension is one less than the expected dimension $`d(𝔰)`$ of the moduli space. Thus, for generic representatives $`V_1`$ and $`V_2`$, this intersection is empty.
* The case where $`\rho __1[A_1,\mathrm{\Phi }_1]C`$ and $`\rho __2[A_2,\mathrm{\Phi }_2]𝒥`$ is excluded by a similar dimension count. We have that $`[A_1,\mathrm{\Phi }_1](X_1,𝔰_1,C)V_1`$, $`[A_2,\mathrm{\Phi }_2](X_2,𝔰_2,𝒥)V_2`$, and $`\rho __1[A_1,\mathrm{\Phi }_1]`$ is connected to $`\rho __2[A_2,\mathrm{\Phi }_2]`$ by a (uniquely determined) flow in $`\widehat{}(C,𝒥)`$. This set has expected dimension $`2`$.
The remaining case is that $`[A_1,\mathrm{\Phi }_1](X_1,𝔰_1,C)`$, and $`[A_2,\mathrm{\Phi }_2](X_2,𝔰_2,C)`$, with
$$\rho __1[A_1,\mathrm{\Phi }_1]=\rho __2[A_2,\mathrm{\Phi }_2].$$
In particular, $`[A_1,\mathrm{\Phi }_1]`$ lies in the compact set $`(X_1,𝔰_1,C)`$, while $`[A_2,\mathrm{\Phi }_2]`$ lies in the set $`\rho __2^1(\rho __1(X_1,𝔰_1,C)V_1)V_2`$ which is also compact (according to the dimension count used to exclude Case (C-3) above). Thus, for all sufficiently large $`T`$, the cut-down moduli space lies in the image of the gluing map $`\gamma _{C;T}`$.
On compact subsets of $`X(T)`$ away from the “neck”, gluing is a $`C^{\mathrm{}}`$ small perturbation, which goes to zero as the neck-length is increased; in particular, for $`i=1,2`$,
$$\underset{T\mathrm{}}{lim}\gamma _{C;T}([A_1,\mathrm{\Phi }_1]\mathrm{\#}[A_2,\mathrm{\Phi }_2])|_{X_i}=[A_i,\mathrm{\Phi }_i]|_{X_i}.$$
It follows from this that
(3)
$$\mathrm{\#}(X(T),𝔰_1\mathrm{\#}𝔰_2)V_1V_2=\mathrm{\#}\left(\left((X_1,𝔰_1,C)V_1\right)\times _C\left((X_2,𝔰_2,C)V_2\right)\right).$$
The latter quantity can be thought of as cohomological pairing in $`(X_1,𝔰_1,C)`$ as follows.
Fix an oriented, $`v`$-dimensional submanifold $`V(X_2,𝔰_2,C)`$, and consider the function which assigns to each smooth map
$$f:ZC$$
(where $`Z`$ is a smooth, oriented, compact manifold whose dimensionl equals the codimension of $`V`$) the number of points in the fibered product $`\mathrm{\#}(Z\times _CV)`$ (counting with sign, after arranging $`f`$ to be transverse to $`V`$). This is the pairing of the fundamental cycle of $`Z`$ with an induced cohomology class in $`H^{d_2v}(C,)`$. Indeed, this class can be thought of as the “push-forward” of the Poincaré dual to $`V`$, under a map
$$(\rho __2)_{}:H^i((X_2,C);)H^{i+dim(C)d_2}(C;).$$
Thus, the count in Equation (3) can be thought of as the pairing
$$\mathrm{\#}(X(T),𝔰_1\mathrm{\#}𝔰_2)V_1V_2=[(X_1,𝔰_1,C)],\mathrm{PD}(V_1)\rho __1^{}(\rho __2)_{}\mathrm{PD}(V_2).$$
This pairing vanishes, according to Proposition 2.3. This completes Theorem 2.1 in the case when $`𝔱=𝔱_e`$ for $`e=0,\mathrm{},2g2`$.
In the case when $`𝔱=𝔱_e`$ for $`2g1<e<n`$, the vanishing of the Seiberg-Witten invariant for any $`𝔰`$ structure with $`𝔰|_Y=𝔱`$ is guaranteed by the same dimension count which we used to exclude Case (C-2) above. ∎
The proof of Theorem 1.1 follows from an application of Theorem 2.1, together with the known properties of Seiberg-Witten invariants for complex surfaces of general type (see for instance or , and also ), according to which a minimal surface of general type has only two “basic classes” ($`\mathrm{Spin}^c`$ structures for which the Seiberg-Witten invariant is non-zero), the “canonical” $`\mathrm{Spin}^c`$ structure $`𝔰_0`$ (whose first Chern class is given by $`c_1(𝔰_0)=K_X`$, where $`K_XH^2(X;)`$ is the first Chern class of the complex cotangent bundle of $`X`$), and its conjugate. Moreover, the basic classes of the $`n`$-fold blow-up $`\widehat{X}=X\mathrm{\#}n\overline{}^2`$ are those $`\mathrm{Spin}^c`$ structures $`𝔰`$ whose restriction away from the exceptional spheres agrees with $`𝔰_0`$ or its conjugate, and whose first Chern class evaluated on each exceptional sphere $`E_i`$ satisfies
$$c_1(𝔰),[E_i]=\pm 1.$$
In fact, since $`K_X^2>0`$ for a minimal surface of general type, the basic classes are in one-to-one correspondence with their first Chern classes. In view of this fact, throughout the following proof, we label the basic classes of $`\widehat{X}`$ by their first Chern classes.
Proof of Theorem 1.1. The subgroup $`\delta H^1(Y;)`$ partitions $`\mathrm{Spin}^c(X)`$ into orbits, and Theorem 2.1 states that if $`X`$ could be decomposed, then the sum of invariants under each orbit vanishes. Note moreover that if $`Y`$ separates $`X`$, then the intersection form restricted to the subgroup $`\delta H^1(Y;)`$ is trivial: this is true because we can represent cohomology classes $`[\omega ],[\eta ]\delta H^1(Y;)`$ by differential form representatives $`\omega `$ and $`\eta `$, with $`\omega |_{X_1}0`$ and $`\eta |_{X_2}0`$, so that the representative for $`[\omega ][\eta ]`$, $`\omega \eta `$, vanishes identically. It follows from this that in each orbit, there can exist at most two basic classes, for if we had two basic classes which had the same coefficient in $`K_X`$, then their difference would have negative square. Now, suppose that $`K_XE_1\mathrm{}E_n`$ had another basic class in its orbit. We know that the other basic class would be of the form
$$K_X+E_1+\mathrm{}+E_aE_{a+1}\mathrm{}E_n,$$
after renumbering the exceptional curves if necessary. The difference $`\mathrm{\Delta }`$ is $`2(K_XE_1\mathrm{}E_a)`$, which must have square zero, which forces $`a>0`$ (recall that $`K_X^2>0`$ for a minimal surface of general type). Now, consider the basic class $`K_X+E_1+\mathrm{}+E_n`$. It, too, can have at most one other basic class in its orbit, and the difference has the form
$$\mathrm{\Delta }^{}=2(K_X+ϵ_1E_1+\mathrm{}ϵ_nE_n),$$
where we know that $`ϵ_1,\mathrm{},ϵ_n0`$, in particular $`\mathrm{\Delta }\mathrm{\Delta }^{}`$ is a non-zero class, which is easily seen to have negative square. But this contradicts the fact that $`\mathrm{\Delta }\mathrm{\Delta }^{}\delta H^1(Y;)`$. Thus, it follows that either the basic class $`K_XE_1\mathrm{}E_n`$ or $`K_X+E_1+\mathrm{}+E_aE_{a+1}\mathrm{}E_n`$ is alone in its $`\delta H^1(Y;)`$ orbit. But this contradicts the conclusion of Theorem 2.1. ∎
### 2.1. Final Remarks
It is suggestive to compare the formal framework adopted here with that of equivariant Morse theory. Specifically, the “Chern-Simons-Dirac” operator on $`Y`$ in the set-up of Theorem 2.2 has precisely two critical manifolds, a manifold of reducibles $`𝒥`$ (consisting of configurations whose stabilizer in the gauge group is a circle), and a manifold of irreducibles $`C`$ (consisting of configurations whose stabilizers are trivial). From the point of view of equivariant cohomology, then, there should be an “equivariant Floer homology”, and an analogue of the Bott spectral sequence, whose $`E_2`$ term consists of the homology of the irreducible critical point set $`H_{}(\mathrm{Sym}^e(\mathrm{\Sigma });)`$, and the $`S^1`$-equivariant homology of the reducible manifold, which is given by
$$H^{}(^{\mathrm{}}\times 𝒥;)[U]_{}\mathrm{\Lambda }^{}H_1(Y;).$$
From this point of view, Proposition 2.3, upon which the vanishing theorem rests, can be seen then as the calculation of the differential in this spectral sequence.
The equivariant point of view has been stressed by a number of researchers in the field, including (especially in the context of gluing along rational homology three-spheres) , , .
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# Cosmology with Curvature-Saturated Gravitational Lagrangian 𝑅/√{1+𝑙⁴𝑅²}
## I Introduction
According to an old idea by Sakharov , the gravitational properties of spacetime are caused by the bending stiffness of all quantum fields in a spacetime of scalar curvature $`R`$. This idea of induced gravity has inspired many subsequent theories of gravitation, from Adler’s proposal to consider Einstein gravity as a symmetry breaking effect in quantum field theory to the modern induced gravity derived from string fluctuations . Whatever the precise mechanism, any induced gravity will lead to a Lagrangian which is bounded at large $`R`$, and may also go to zero. The latter case would be analogous to the elastic stiffness of solids, which is constant for small distortions, but vanishes after the solid cracks.
In this paper we investigate the physical consequences of a simple Lagrangian which goes to a constant at large $`R`$, thus interpolating between the Einstein-Hilbert Lagrangian for small $`R`$ and a pure cosmological constant for large $`R`$. This Lagrangian will be referred to as curvature-saturated and reads
$$_{\mathrm{cs}}=\frac{1}{16\pi G}\frac{R}{\sqrt{1+l^4R^2}}.$$
(1)
The length parameter $`l`$ may range from an order of the Planck length $`l_\mathrm{P}`$ or a few orders of magnitude larger than $`l_\mathrm{P}`$. Applying standard methods and those of Refs. , we shall derive the cosmological consequences of the saturation and compare our ansatz with others.
One of the motivations for a renewed interest in a more detailed consideration of cosmology with non-linear curvature terms comes from M-theory, see Ref. “Brane new world”. In a conformal anomaly is considered, which turns out to have analogous consequences as Starobinsky’s anomaly-driven inflation with $`R`$– and $`R^2`$-terms, see e.g. Refs. for the older results. Ref. contains the latest results concerning the effective $`\mathrm{\Lambda }`$-term in such models.
Our own direct motivation to tackle the model discussed below was as follows: We tried to make the analogy proposed in more closer than done by others; the analogy with solid state physics is this one: For small forces, the resistance to bending is proportional to this force, but after a certain threshold – defined by cracking the solid – the resistance vanishes.
A similar line of reasoning was deduced in Ref. : There the finite-size effects from the closed Friedmann universe to the quantum states of fields have been calculated. Instead of continuous distribution of the energy levels of the quantum fields, one has a discrete spectrum. Qualitatively, the result is: If the radius $`a`$ of the spatial part of spacetime shrinks close to zero, which is almost the same as very large $`R`$, then the spacings between the energy levels become larger and larger, and after a certain threshold, all fields will be in the ground state. This behaviour shall be represented by an effective action. The concrete form of the corresponding effective Lagrangian is not yet fully determined (that shall be the topic of later work), but preliminarily we found out that the behaviour for large $`R`$ will quite probably be of a Lagrangian bounded by a special effective $`\mathrm{\Lambda }`$; so we have chosen one of the easiest analytic functions possessing this large-$`R`$ behaviour together with the correct weak-field shape.
The paper is organized as follows: In Sec. II we calculate the consequences of the effective Lagrangian $`_{\mathrm{cs}}`$.
In Sec. III we investigate the consequences of the $`R`$-dependence of the effective gravitational constant defined by
$$\frac{1}{16\pi G_{\mathrm{eff}}}\frac{d}{dR},$$
(2)
which is
$$G_{\mathrm{eff}}=G\sqrt{1+l^4R^2}^3$$
(3)
for $`=_{\mathrm{cs}}`$ and tends to infinity as $`R\pm \mathrm{}`$.
Then we apply two different conformal transformations to $`_{\mathrm{cs}}`$. One of them, presented in Sec. IV, makes $`_{\mathrm{cs}}`$ asymptotically equivalent to the Gurovich-ansatz ,
$$=\frac{R}{16\pi G}+c_1|R|^{4/3}.$$
(4)
The other transformation, by the Bicknell theorem given in Sec. V, establishes a conformal relation to Einstein’s theory, with a minimally coupled scalar field. In the literature, see and the references cited there, only the second of these conformal transformations has so far been used. The physical consequences of these three theories are, of course, quite different since the metrics are not related to each other by coordinate transformations.
Our approach differs fundamentally from that derived from the limiting curvature hypothesis (LCH) in Refs. , where the gravitational Lagrangian reads
$$=R+\frac{\mathrm{\Lambda }}{2}\left(\sqrt{1R^2/\mathrm{\Lambda }^2}1\right)$$
(5)
whose derivative with respect to $`R`$ diverges for $`R\mathrm{\Lambda }`$. This divergence was supposed to prevent a curvature singularity, a purpose not completely reached by the model presented in the first of Refs. because other curvature invariants may still diverge. (Let us note for completeness: In the second of Refs. , a more detailed version of the LCH is presented which covers also the bounding of the other curvature invariants; it is restricted to isotropic cosmological models. For more general space-times one faces the problem that sometimes a curvature singularity exists, but all polynomial curvature invariants remain bounded there.)
In contrast to Eq. (5), our model favors high curvature values.
It turns out that the use of synchronized or conformal time is not optimal for our problem. We therefore use a new time coordinate which we call curvature time for the spatially flat Friedmann model. The general properties of this coordinate choice are described in Sec. VI.
In Sec. VII we study the consequences of curvature-saturation for some cosmological models using the coordinates of Sec. VI. In Sec. VIII, finally, we summarize our results and compare with the related papers to .
## II Field Equations of Curvature-Saturated Gravity
The curvature-saturated Lagrangian (1) interpolates between the Einstein-Hilbert Lagrangian
$$_{\mathrm{EH}}=\frac{R}{16\pi G},$$
(6)
which is experimentally confirmed at weak fields, and a pure cosmological constant at strong fields
$$=\frac{\pm 1}{16\pi Gl^2}.$$
(7)
The $`R`$ dependence is plotted in Fig. 1.
The usual gravitational constant is obtained from the derivative of the Einstein-Hilbert Lagrangian:
$$\frac{1}{16\pi G}=\frac{d_{\mathrm{EH}}}{dR}.$$
(8)
From our curvature-saturated Lagrangian (1) we obtain, with this derivative, the effective gravitational constant (3). The definition (8) is motivated as follows: If one considers the Newtonian limit for a general Lagrangian $`(R)`$ which may contain a nonvanishing cosmological constant, the potential between two point masses contains a Newtonian $`1/r`$-part plus a Yukawa-like part $`\mathrm{exp}(r/r_Y)`$ stemming from the nonlinearities of the Lagrangian; the details are given in the Appendix. At distances much larger than $`r_Y`$, but much smaller than $`1/\sqrt{R}`$, only the $`1/r`$-term survives, and the coupling strength of the $`1/r`$-term is given by the effective gravitational constant $`G_{\mathrm{eff}}`$. For a recent version to deduce such weak-field expressions, see Ref. .
For a general Lagrangian $`(R)`$ such as (1), the calculation of the field equation is somewhat tedious, since the Palatini formalism which simplifies the calculation in Einstein’s theory is no longer applicable. Recall that in this, metric and the affine connection are varied independently, the latter being identified with the Christoffel symbol only at the end.
Here the following indirect procedure leads rather efficiently to the correct field equations. Let
$$^{}\frac{d}{dR},^{\prime \prime }\frac{d^2}{dR^2},$$
(9)
and form the covariant energy-momentum tensor of the gravitational field which is given by the variational derivative of $``$ with respect to the metric $`g_{ab}`$:
$$\mathrm{\Theta }_{ab}\frac{2}{\sqrt{g}}\frac{\delta \sqrt{g}}{\delta g_{ab}},$$
(10)
where $`g`$ denotes the determinant of $`g_{ab}`$. For dimensional reasons, $`\mathrm{\Theta }_{ab}`$ has the following structure
$$\mathrm{\Theta }_{ab}=\alpha ^{}R_{ab}+\beta ^{}Rg_{ab}+\gamma g_{ab}+\delta \text{ }\text{ }\text{ }^{}g_{ab}+ϵ_{;ab}^{}$$
(11)
with the 5 real constants $`\alpha \mathrm{}ϵ`$. These constants can be uniquely determined up to one overall constant factor by the covariant conservation law
$$\mathrm{\Theta }_{;b}^{ab}=0.$$
(12)
The overall factor is fixed by the Einstein limit $`l0`$ of the theory, where $`\mathrm{\Theta }_{ab}=(R_{ab}\frac{1}{2}Rg_{ab})/8\pi G`$. In this way we derive the following form of the covariantly conserved energy-momentum tensor of the gravitational field
$$\mathrm{\Theta }_{ab}=\frac{1}{16\pi G}\left(2^{}R_{ab}g_{ab}+2\text{ }\text{ }\text{ }^{}g_{ab}2_{;ab}^{}\right).$$
(13)
The calculation is straightforward, if one is careful to distinguish between $`(\text{ }\text{ }\text{ }^{})_{;a}`$ and $`\text{ }\text{ }\text{ }(_{;a}^{})`$, which differ by a multiple of the curvature scalar.
Inserting our curvature-saturated Lagrangian (1) into (9) and omitting the subscript, we have
$$=\frac{R}{2}\left(1+l^4R^2\right)^{1/2},^{}=\frac{d}{dR}=\frac{1}{2}\left(1+l^4R^2\right)^{3/2},$$
(14)
and find from (13)
$$\mathrm{\Theta }_{ab}=\frac{1}{8\pi G}\left\{\frac{R_{ab}}{\left(1+l^4R^2\right)^{3/2}}\frac{Rg_{ab}}{2\left(1+l^4R^2\right)^{1/2}}+g_{ab}\text{ }\text{ }\text{ }\left[\frac{1}{\left(1+l^4R^2\right)^{3/2}}\right]\left[\frac{1}{\left(1+l^4R^2\right)^{3/2}}\right]_{;ab}\right\}.$$
(15)
Setting $`l=0`$ reduces this to $`1/16\pi G`$ times the Einstein tensor. The trace of (15) is
$$\mathrm{\Theta }_a{}_{}{}^{a}=\frac{1}{8\pi G}\{\frac{R+2l^4R^3}{\left(1+l^4R^2\right)^{3/2}}3\text{ }\text{ }\text{ }\left[\frac{1}{\left(1+l^4R^2\right)^{3/2}}\right]\}.$$
(16)
According to Einstein’s equation, $`\mathrm{\Theta }_{ab}`$ has to be equal to the energy momentum tensor of the matter $`T_{ab}`$, i.e., $`T_{ab}=\mathrm{\Theta }_{ab}`$. Equation (16) implies that in the vacuum, the only constant curvature scalar is $`R=0`$, such that this model does not possess a de Sitter solution. Further, we can see from Eq. (15), that a curvature singularity does not necessarily imply a divergence of energy-momentum, but may be compensated by the infinity of $`G_{\mathrm{eff}}`$.
## III Effective gravitational constant and weak-field behavior
Let us compare the effective gravitational constant $`G_{\mathrm{eff}}`$ of our curvature-saturated model with those of other models discussed in the literature. From (3) we see that $`G_{\mathrm{eff}}`$ has the weak-field expansion
$$G_{\mathrm{eff}}=G\left(1+\frac{3}{2}l^4R^2+\mathrm{}\right),$$
(17)
and the strong-field expansion
$$G_{\mathrm{eff}}=Gl^6|R|^3\left(1+\frac{3}{2l^4R^2}+\mathrm{}\right).$$
(18)
The full $`R`$-behavior is plotted in Fig. 2.
The weak-field expansion of $`_{\mathrm{cs}}`$ is given by
$$_{\mathrm{cs}}=\frac{R}{16\pi G\sqrt{1+l^4R^2}}=\frac{R}{16\pi G}+\underset{k=1}{\overset{\mathrm{}}{}}b_kR^{2k+1}$$
(19)
with real coefficients $`b_k`$, where $`b_1=l^4/32\pi G`$.
As one can see, the quadratic term is absent, so that the linearized field equation coincides with the linearized Einstein equation. Thus we encounter neither ghosts nor tachyons; for details see Appendix B.
There is, however, a price to pay for it. The theory has lost linearization stability of the solutions. This latter property has the following consequences: If one performs a weak-field expansion
$$g_{ij}=\eta _{ij}+\underset{m=1}{\overset{\mathrm{}}{}}ϵ^mg_{ij}^{(m)}$$
(20)
around flat spacetime to solve the field equation, one has to use the terms up to the order $`m=2`$ to get the complete weak-field part of the set of solutions. With this peculiarity, we obtain a well-posed Cauchy problem for the gravity theory following from the Lagrangian $`_{\mathrm{cs}}`$.
Let us now compare our theory with others available in the literature. Let
$$_{\alpha ,n}(R)=\frac{R}{16\pi G}+\alpha R^n$$
(21)
with some number $`n>1`$ and constant $`\alpha 0`$. In analogy with Eq. (2) we calculate the effective gravitational constant from
$$\frac{1}{16\pi G_{\mathrm{eff}}}=\frac{d_{\alpha ,n}}{dR}=\frac{1}{16\pi G}+\alpha nR^{n1}$$
(22)
such that
$$G_{\mathrm{eff}}=\frac{G}{1+16\pi \alpha nR^{n1}G},$$
(23)
i.e., $`G_{\mathrm{eff}}0`$ as $`R\pm \mathrm{}`$. For $`n=2`$, more exactly: for all even natural numbers $`n`$, we meet an additional peculiarity that $`G_{\mathrm{eff}}`$ can diverge for finite values of $`R`$ already. Such values of $`R=R_{\mathrm{crit}}`$ are called critical . For $`n=2`$ we get
$$R_{\mathrm{crit}}=\frac{1}{32\alpha \pi G},$$
(24)
and this is the region where $`G_{\mathrm{eff}}`$ changes its sign, as shown in Figures 3 and 4. At critical values of the curvature scalar, the Cauchy problem fails to be a well-posed one.
## IV Conformal duality
In Ref. , a duality transformation relating between different types of nonlinear Lagrangians has been found. In the present notation it implies the following relation. Let
$$\widehat{g}_{ab}=^2g_{ab}$$
(25)
be the conformally transformed metric with $`^{}0`$, which is fulfilled by our Lagrangian (1). Then the conformally transformed curvature scalar equals
$$\widehat{R}=\frac{3R}{^2}\frac{4}{^3},$$
(26)
and the associated Lagrangian is
$$\widehat{}=\frac{2R}{^3}\frac{3}{^4}.$$
(27)
We easily verify that $`\widehat{}^{}^{}=1`$. Then one can prove that $`g_{ab}`$ solves the vacuum field equation following from $`(R)`$ if and only if $`\widehat{g}_{ab}`$ of Eq. (25) solves the corresponding equation for $`\widehat{}(\widehat{R})`$ of Eq. (27).
Example: For $`=R^{k+1}`$ we find, up to an inessential constant factor, $`\widehat{}=\widehat{R}^{\widehat{k}+1}`$ with $`\widehat{k}=1/(21/k)`$, such that for a purely quadratic theory with $`=R^2`$, also $`\widehat{}=\widehat{R}^2`$. For our curvature-saturated model $``$ const. we should expect a behavior with $`k1`$, i.e., $`\widehat{k}1/3`$, this leads to $`\widehat{}\widehat{R}^{4/3}`$, which is the Gurovich-model , cf. Eq. (4).
Let us study this in more detail. To simplify the expressions we use, in this subsection only, reduced units with $`16\pi G=1`$ to best exhibit the fixed point $`l=0`$ of this transformation making it an identity transformation if applied to Einstein’s theory where $`k=1`$. In the present units, Eqs. (14) have to be multiplied by 2 and become
$$=R\left(1+l^4R^2\right)^{1/2},^{}=\frac{d}{dR}=\left(1+l^4R^2\right)^{3/2}.$$
(28)
Inserting these into (25)–(27), we obtain
$$\widehat{g}_{ab}=\frac{g_{ab}}{(1+l^4R^2)^{\mathrm{\hspace{0.17em}3}}}$$
(29)
and
$$\widehat{R}=R(1+l^4R^2)^{\mathrm{\hspace{0.17em}3}}(14l^4R^2).$$
(30)
For small $`R`$ we have
$$\widehat{R}=R(1l^4R^2+\mathrm{}),$$
(31)
and for large $`|R|`$
$$\widehat{R}=4l^{16}R^9\left(1+\frac{11}{4l^4R^2}+\mathrm{}\right).$$
(32)
The inverse function $`R(\widehat{R})`$ of (30) is not expressible in closed form, but its small- and large-curvature expansion can be calculated from (31) and (32)
$$R=\widehat{R}(1+l^4\widehat{R}^2+\mathrm{}),R=\left(\frac{\widehat{R}}{4l^{16}}\right)^{\mathrm{\hspace{0.17em}1}/9}\left[1\frac{11}{36l^4}\left(\frac{4l^{16}}{\widehat{R}}\right)^{\mathrm{\hspace{0.17em}2}/9}+\mathrm{}\right]$$
(33)
¿From Eq. (27) we see that
$$\widehat{}=R(1+l^4R^2)^{\mathrm{\hspace{0.17em}9}/2}(13l^4R^2)$$
(34)
where $`R(\widehat{R})`$ has to be inserted. For large $`R`$ we use the right-hand equation in (33) and obtain the limiting behavior
$$\widehat{}=3l^{22}\left(\frac{\widehat{R}}{4l^{16}}\right)^{\mathrm{\hspace{0.17em}4}/3}\left[1\frac{51}{6l^4}\left(\frac{4l^{16}}{\widehat{R}}\right)^{\mathrm{\hspace{0.17em}2}/9}+\mathrm{}\right].$$
(35)
## V Bicknell’s theorem
Bicknell’s theore , in the form described in Ref. , relates Lagrangians of the type (14) to Einstein’s theory coupled minimally to a scalar field $`\varphi `$ with a certain interaction potential $`\stackrel{~}{V}(\varphi )`$. This Lagrangian is given by
$$_{\mathrm{EH}}+\frac{1}{2}\varphi _{,i}\varphi ^{,i}\stackrel{~}{V}(\varphi ).$$
(36)
The relation of $`\stackrel{~}{V}(\varphi )`$ with $`(R)`$ is expressed most simply by defining a field with a different normalization $`\psi =\sqrt{2/3}\varphi `$, in terms of which the potential $`\stackrel{~}{V}(\varphi )=V(\psi )`$ reads
$$V(\psi )=(R)e^{2\psi }\frac{R}{2}e^\psi ,$$
(37)
with $`R`$ being the inverse function of
$$\psi =\mathrm{ln}[2^{}(R)].$$
(38)
The metric in the transformed Lagrangian (36) is
$$\stackrel{~}{g}_{ab}=e^\psi g_{ab}.$$
(39)
For our particular Lagrangian (14) we have from (38):
$$\psi =\frac{3}{2}\mathrm{ln}(1+l^4R^2).$$
(40)
Now we restrict our attention to the range $`R>0`$ where $`\psi <0`$; the other sign can be treated analogously. Then (40) is inverted to
$$R=\frac{1}{l^2}\sqrt{e^{2\psi /3}1},$$
(41)
such that (37) becomes
$$V(\psi )=\frac{1}{2l^2}(e^{5\psi /3}e^\psi )\sqrt{e^{2\psi /3}1}.$$
(42)
In the range under consideration, this is a positive and monotonously increasing function of $`\psi `$ (see Fig. 5), with the large-$`\mathrm{\Phi }`$ behavior
$$V=\frac{1}{2l^2}e^{2\psi }.$$
(43)
This is the typical exponential potential for power-law inflation. As mentioned at the end of Section II, no exact de Sitter inflation exists. For $`\psi 0`$, also $`V(\psi )0`$ like $`4\sqrt{2/3}\psi ^{3/2}`$.
If $`V(\psi )`$ has a quadratic minimum at some $`\psi _0`$ with positive value $`V_0=V(\psi _0)`$, then there exists a stable de Sitter inflationary phase. As a pleasant feature, the potential $`V(\psi )`$ has no maximum which have given rise to tachyons.
From Eq. (40) one can see that for weak fields, $`\psi R^2`$, whereas a $`R+R^2`$-theory has $`\psi R`$. In other words: In our model it is a better approximation to assume the conformal factor $`e^\psi `$ to be approximately constant for weak fields then in $`R+R^2`$-theories, since at the level keeping only terms linear in $`R`$ the two metrics $`g_{ab}`$ and $`\stackrel{~}{g}_{ab}`$ in (39) coincide.
## VI Friedmann models in curvature time
The expanding spatially flat Friedmann model may be parametrized with the help of curvature time $`a>0`$ as follows:
$`ds^2=a^2\left[{\displaystyle \frac{da^2}{B(a)}}dx^2dy^2dz^2\right],`$ (44)
where $`B(a)`$ is an arbitrary positive function determining $`R`$ as
$`R={\displaystyle \frac{3}{a^3}}{\displaystyle \frac{dB}{da}},`$ (45)
depending only on the first derivative of $`B(a)`$. This is a special feature of (44) since, in general, the curvature scalar depends on the second derivative of the metric components. Note also the linear dependence of $`R`$ on $`B^{}dB/da`$, in contrast to the usual nonlinear dependence of the curvature scalar on the first derivative of the metric coefficients.
Let us recall some facts on Friedmann models in curvature time and exhibit the corresponding transformation to synchronized time.
### A From curvature time to synchronized time
The spatially flat Friedmann model in synchronized time has the metric
$`ds^2=dt^2a^2(t)(dx^2+dy^2+dz^2).`$ (46)
Metric (44) goes over to metric (46) via
$$dt=\frac{ada}{\sqrt{B(a)}},$$
(47)
such that
$`t=t(a)={\displaystyle \frac{ada}{\sqrt{B(a)}}}.`$ (48)
The inverse function $`a(t)`$ provides us with the desired transformation.
### B From synchronized time to curvature time
Consider $`a(t)`$ in an expanding model with
$$\dot{a}\frac{da}{dt}>0.$$
(49)
Then we can invert $`a(t)`$ to $`t(a)`$, and have
$`B(a)=a^2[\dot{a}(t(a)]^2.`$ (50)
¿From this relation we understand why $`R`$ depends on the first derivative of $`B`$ only: $`B`$ itself contains a derivative of $`a`$, and $`R`$ is known to contain up to second order derivatives of $`a(t)`$.
### C Examples
Let $`a(t)=t^n`$, i.e. $`t=a^{1/n}`$, $`\dot{a}(t)=nt^{n1}`$, $`\dot{a}(t(a))=na^{11/n}`$. Then Eq. (50) yields
$`B(a)=n^2a^{42/n}.`$ (51)
Let further $`a(t)=e^{Ht}`$, $`H=\text{const.}>0`$, $`\dot{a}=Ha`$. Then
$`B(a)=H^2a^4.`$ (52)
Obviously, Eq. (52) is a limiting form of Eq. (51) for $`n\mathrm{}`$. Equation (44) with $`B(a)`$ from (52) represents a vacuum solution of Einstein’s theory with $`\mathrm{\Lambda }`$-term where $`\mathrm{\Lambda }=3H^2`$, namely the de Sitter spacetime.
Let us also give some examples for the direct use the curvature time:
* From Eq. (45) we see that $`R=0`$ implies $`B`$ const., corresponding to $`n=\frac{1}{2}`$ in Eq. (51), i.e., $`a=t^{1/2}`$ in synchronized time. This is the usual Friedmann radiation model.
* Also from Eq. (45), a constant $`R`$ $`0`$ implies $`B=C_1+C_2a^4`$ with constants $`C_1`$ and $`C_2`$$`C_20`$.
For $`C_1=0,C_2=H^2`$, this represents the de Sitter spacetime Eq. (52).
* The dust-model in synchronized coordinates is given by $`a=t^{2/3}`$, i.e., with Eq. (51) we get
$`B(a)={\displaystyle \frac{4}{9}}a,`$ (53)
such that $`B^{}=`$ const. Together with Eq. (45), this leads to
$$Ra^3=\text{const,}$$
(54)
ensuring mass conservation, because $`R`$ is proportional to the mass density, and the pressure is negligible for dust.
### D The variational derivative
For the metric (44) we have
$`\sqrt{g}\sqrt{detg_{ij}}={\displaystyle \frac{a^4}{\sqrt{B}}}.`$ (55)
The Lagrangian for Einstein’s theory with $`\mathrm{\Lambda }`$-term reads
$`=(R+2\mathrm{\Lambda })\sqrt{g}.`$ (56)
With (45) and (50) we get from (56)
$$=\left(2\mathrm{\Lambda }\frac{3B^{}}{a^3}\right)a^4B^{1/2}.$$
(57)
The vanishing of the variational derivative
$$\frac{\delta }{\delta B}\frac{}{B}\left(\frac{}{B^{}}\right)^{}=0$$
(58)
gives $`B=H^2a^4`$ with $`\mathrm{\Lambda }=3H^2`$, i.e., the usual de Sitter spacetime. No integration is necessary, since the derivative of $`B`$ cancels. Intermediate expressions are
$$\frac{}{B}=\left(2\mathrm{\Lambda }\frac{3B^{}}{a^3}\right)a^4\left(\frac{1}{2}\right)B^{3/2},$$
(59)
$$\frac{}{B^{}}=3aB^{1/2},\left(\frac{}{B^{}}\right)^{}=3B^{1/2}+\frac{3}{2}aB^{}B^{3/2}.$$
(60)
### E Remaining coordinate-freedom
Translations in $`t`$ do not change the form of the metric (46). This freedom is related to the fact that the integration constant in the integral (48) remains undetermined; this coordinate freedom has no analog in the metric in curvature time Eq. (44).
The metric (44) has the following property: It remains unchanged under multiplication of $`a^4`$ and $`B`$ by the same positive constant. Such a constant factor appears if we multiply the spatial coordinates by a constant factor. In synchronized coordinates this property means that not $`a`$ itself, but only the Hubble parameter
$`H(t):={\displaystyle \frac{\dot{a}}{a}}`$ (61)
has an invariant meaning. By the same token, not $`B(a)`$ itself, but only $`B(a)/a^4`$ has an invariant meaning. In fact, from Eq. (50) we see that
$`{\displaystyle \frac{B}{a^4}}=H^2.`$ (62)
## VII Cosmological solutions
Here we recall some formulas of Ref. , and present some new results for the curvature-saturated Lagrangian.
### A Solutions for Lagrangian $`R^m`$
For the Lagrangian $`=R^m`$, we obtain the following exact solutions for a closed Friedmann universe:
$$ds^2=dt^2\frac{t^2}{2m^22m1}d\sigma _{(+)}^2,$$
(63)
where $`d\sigma _{(+)}^2`$ is the metric of the unit 3-sphere.
Analogously, for the open model
$$ds^2=dt^2\frac{t^2}{2m2m^2+1}d\sigma _{(+)}^2.$$
(64)
Of course, both expressions are valid for positive denominators only.
For the spatially flat Friedmann model, it proves useful to employ the cosmic scale factor $`a`$ itself as a time-like coordinate.
$$ds^2=a^2\left[Q^2(a)da^2dx^2dy^2dz^2\right].$$
(65)
This coordinate is meaningful as long as the Hubble parameter is different from zero, so that we cover only time intervals where the universe is either expanding or contracting. Possibly existing maxima or minima of the cosmic scale factor as seen in synchronized time can, however, been dealt by a suitable limiting process and patching. The curvature scalar reads now
$$R=\frac{6}{a^3Q^3}\frac{dQ}{da},$$
(66)
and to reduce the order of the field equation it proves useful to define
$$P(a)=\frac{d\mathrm{ln}Q}{da}.$$
(67)
Then the field equation is fulfilled if
$$0=m(m1)\frac{dP}{da}+(m1)(12m)P^2+m(43m)\frac{P}{a}.$$
(68)
Therefore, the spatially flat Friedmann models can be solved in closed form, but not always in synchronized coordinates.
### B Solutions for Lagrangian $`_{\mathrm{cs}}`$
In the context of our curvature-saturated model, we shall restrict ourselves to the expanding spatially flat Friedmann model. The field equation written in synchronized or conformal time—the two most often used time coordinates used for this purpose—have the disadvantage that the number of terms is quite large, and that even in the simplest case $`=\frac{1}{2}R^2`$ we cannot give closed-form solutions, apart from the trivial solutions $`R0`$ having the same geometry as the radiation universe ($`a=\sqrt{t}`$ in synchronized time $`t`$) and the de Sitter universe ($`a=e^t`$ in synchronized time $`t`$). So, we prefer to work in the less popular coordinates (65). In principle, the field equation should be of fourth order, but we shall reduce it to second order.
To find the field equation for a spatially flat Friedmann model with our Lagrangian, it is useful to consider first a general nonlinear Lagrangian and specialize to $`_{\mathrm{cs}}`$ afterwards. To simplify (66), we define instead of $`Q(a)`$ the function $`B(a)=Q(a)^2>0`$ as a new dependent function. Then (65) reads
$$ds^2=a^2\left[\frac{da^2}{B(a)}dx^2dy^2dz^2\right]$$
(69)
and (66) goes over to
$$R=\frac{3}{a^3}\frac{dB}{da}.$$
(70)
Thus, $`B`$ itself does not appear explicitly, and only first, and not second derivatives are present. The geometric origin of this property is the same as in Schwarzschild coordinates—one integration constant is lost in the definition of the coordinates, and this makes curvature depend only on the first derivative of the metric.
From the 10 vacuum field equations (15) only the 00-component is essential; it is the constraint equation, therefore it has one order less than the full field equation, but if the constraint is fulfilled always, then all other components are fulfilled, too.<sup>*</sup><sup>*</sup>*This behavior is known already from the Friedmann equation in General Relativity: Energy density is proportional to the square of the Hubble parameter which contains only a first derivative. Together with Eq. (70) we should now expect that the fourth order field equation (15) can be reduced to one single second order equation for $`B(a)`$, where hopefully, $`B`$ itself no more appears.
The equation $`\mathrm{\Theta }_{00}=0`$ is via (9) and (13) equivalent to
$$0=3^{}\left(2Ba\frac{dB}{da}\right)a^418aB^{\prime \prime }\frac{d}{da}\left(\frac{1}{a^3}\frac{dB}{da}\right),$$
(71)
which is much simpler than the analogous equation in synchronous time, as observed here for the first time.
Before we insert our Lagrangian $`_{\mathrm{cs}}`$ into (71), let us cross check its validity by solving known problems: If $`^{\prime \prime }`$ vanishes identically, then $`^{}`$ is a constant, and we return to Einstein’s theory. The case $`B`$ const. gives the radiation universe, while $`B=a^4`$ is the exact de Sitter solution. For the Lagrangian $`=\frac{1}{2}R^2`$ with $`^{}=R`$ and $`^{\prime \prime }=1`$, and Eq. (71) reduces to
$$0=a\dot{B}^24aB\ddot{B}+8B\dot{B},$$
(72)
where a dot denotes differentiation with respect to $`a`$. Again, $`B=a^4`$ is the exact de Sitter solution. Defining $`\beta =\mathrm{ln}B`$ and $`z=a\dot{\beta }`$, Eq. (72) goes over in
$$4a\dot{z}=3z(4z).$$
(73)
With $`\alpha =\mathrm{ln}a`$ we arrive at
$$4\frac{dz}{d\alpha }=3z(4z),$$
(74)
which can be solved in closed form. Qualitatively it is clear that $`z=4`$, i.e., the de Sitter solution, represents an attractor. Solving Eq. (74) we obtain in the region $`0<z<4`$:
$$z=2+2\mathrm{tanh}\left(\frac{3}{2}\alpha \right),$$
(75)
showing explicitly that $`z4`$ for $`\alpha \mathrm{}`$. The metric can be calculated from
$$\dot{\beta }=\frac{2}{a}\left(1+\frac{a^31}{a^3+1}\right),$$
(76)
using the identity
$$\mathrm{tanh}\mathrm{ln}x=\frac{x^21}{x^2+1}.$$
(77)
After these preparations we are ready to deal with our Lagrangian $`_{\mathrm{cs}}`$. We insert $``$ and $`^{}`$ from Eq. (28), and
$$^{\prime \prime }=3l^4R(1+l^4R^2)^{5/2}$$
(78)
into Eq. (71) and obtain, after setting $`l=1`$, the simple expression
$$54a^9B\dot{B}\frac{d}{da}(a^3\dot{B})=a^5(a^6+9\dot{B}^2)(2Ba\dot{B})+\dot{B}(a^6+9\dot{B}^2)^2.$$
(79)
In these coordinates, the flat Minkowski spacetime does not exist, and the radiation universe $`R=0`$ is not a solution. This is why $`B=`$ const. yields no solution to Eq. (79). Also, as was known from the beginning: the de Sitter spacetime $`B=a^4`$ is not an exact solution here. However, in the nearby-region where the Lagrangian is well approximated by a quadratic function in $`R`$ with a nonvanishing linear term, the behavior of the solutions is quite similar to that of $`R+R^2`$-models, where no exact de Sitter solution exists, but a quasi de Sitter solution represents a transient attractor with sufficient long duration to solve the known cosmological problems. These calculations have been presented at different places, most explicitly in Ref. . After this phase, the universe goes to the weak-field behavior, where our model behaves as usual.
The main departure of our model from the usual one is in the region of large curvature scalar, where $`|\dot{B}|`$ is large compared to $`a^3`$. To find out the behavior of the solutions in this limit, we compare the leading terms in Eq. (79) and see that $`\ddot{B}`$ is proportional to $`\dot{B}^4`$, where the coefficient of proportionality is positive and slowly varying. Thus, we find approximately $`B(a)a^{2/3}`$ for small $`a`$. This implies the existence of a big-bang singularity, but with a different behavior: From Eq. (69) we obtain
$$ds^2=a^2\left[\frac{da^2}{a^{2/3}}dx^2dy^2dz^2\right],$$
(80)
which corresponds in synchronized time to the behavior
$$ds^2=dt^2t^{6/5}(dx^2+dy^2+dz^2),$$
(81)
this being a good approximation to the exact metric for small $`t`$, differing from the usual big-bang behavior in almost all other models. Further details of our model will be presented elsewhere.
### C The cosmological singularity
Here we present the argument with the singularity behaviour mentioned at the end of section II: In our model, differently from Einstein’s theory, the divergence of the curvature does not necessarily imply the divergence of any part of the energy-momentum-tensor. Let us concentrate on the trace. The r.h.s. of Eq. (16) reads
$$\frac{1}{8\pi G}\left\{\frac{R+2l^4R^3}{\left(1+l^4R^2\right)^{3/2}}3\text{ }\text{ }\text{ }\left[\frac{1}{\left(1+l^4R^2\right)^{3/2}}\right]\right\}$$
and this expression must be equal to the trace $`T`$ of the energy-momentum tensor. In Einstein’s theory, $`R\pm \mathrm{}`$ necessarily implies $`T\pm \mathrm{}`$, whereas here, $`T`$ may remain finite even if $`R\mathrm{}`$.
Detailed numerical calculations would support this qualitative picture, however, we postpone such calculations until we have a more strictly physically motivated form of the Lagrangian.
## VIII Discussion
We have argued that the gravitational action $`𝒜`$ has a decreasing dependence on $`R`$ for increasing $`|R|`$. Such a behavior is expected from the spacetime stiffness caused by the vacuum fluctuations of all quantum fields in the universe.
Our model does not have the tachyonic disease of $`R+R^2`$ models studies by Stelle and others .
Since our model has an action which interpolates between Einstein’s action and a pure cosmological term, it promises to have interesting observable consequences which may explain some of the experimental cosmological data.
The heat-kernel expansion of the effective action in a curved background is closely related to the Seeley-Gilkey coefficients , and for higher loop expansion also higher powers of curvature appear: To get the $`n`$-loop approximation one has to add terms until $`R^{n+1}`$, a behavior which also happens in the string effective action . So, if one cuts this procedure at a certain value of $`n`$, one gets always as leading term for high curvature values a term like $`R^{n+1}`$. However, the $`n`$-loop approximations need not converge to the correct result if one simply takes $`n\mathrm{}`$ in the $`n`$-loop-result. In fact, what we have used in the present paper is such an example:
$$_{\mathrm{cs}}=\frac{R}{16\pi G\sqrt{1+l^4R^2}}=\frac{R}{16\pi G}+\underset{k=1}{\overset{\mathrm{}}{}}b_kR^{2k+1}$$
(82)
with some real constants $`b_k`$, where
$$b_1=\frac{l^4}{32\pi G}$$
(83)
but the Taylor expansion on the right hand side diverges for $`R>l^2`$. So, the Taylor expansion is useful for small $`R`$-values only, and for large values $`R`$ we need a correct analytical continuation.
Prigogine et al. have proposed in Eq. (18) of Ref. a model where the effective gravitational constant depends on the Hubble parameter of a Friedmann model. Though this ansatz depends on the special 3+1-decomposition of spacetime, it shares some similarities with the model discussed here. More recent developments how to find a well-founded gravitational action from considering quantum effects can been found in and .
Quite recently, see for instance , accelerated expansion models of the universe have been discussed and compared with new observations. We postpone the comparison of our model with these observations to later work.
## Acknowledgment
H.-J. S. gratefully acknowledges financial support from DFG and from the HSP III-program. We thank V. Gurovich and the colleagues of the Free University Berlin, where this work has been done, especially M. Bachmann and A. Pelster, for valuable comments.
## A Newtonian limit in a nonflat background
The Newtonian limit of a theory of gravity is defined as follows: It is the weak-field slow-motion limit for fields whose energy-momentum tensor is dominated by its zero-zero component in comoving time. Usually, the limit is formed in a flat background, and sometimes, this is assumed to be a necessary assumption. This is, however, not true, and we show here briefly how to calculate the Newtonian limit in a nonflat background, Moreover, our approach is different from what is usually called Newtonian cosmology. To have a concrete example, we take the background as a de Sitter spacetime.
The slow-motion assumption allows us to work with static spacetime and the matter, assuming the energy-momentum tensor to be
$$T_{ij}=\rho \delta _i^0\delta _j^0,$$
(A1)
where $`\rho `$ is the energy density, and time is assumed to be synchronized. The de Sitter spacetime in its static form can be given as
$`ds^2=(1kr^2)dt^2+{\displaystyle \frac{dr^2}{1kr^2}}+r^2d\mathrm{\Omega }^2,`$ (A2)
where $`x^0=t,x^1=r,x^2=\chi ,x^3=\theta `$ and $`d\mathrm{\Omega }^2=d\chi ^2+\mathrm{sin}^2\chi d\theta ^2`$ is the metric of the 2–sphere. In this Appendix, we have changed the signature of the metric from $`(+)`$, which is usual in cosmology, to $`(+++)`$, which leads to the standard definition of the Laplacian.
The parameter $`k`$ characterizes the following physical situations: For $`k=0`$, we have the usual flat background. By setting $`k=0`$ we can therefore compare the results with the well-known ones. The case $`k>0`$ corresponds to a positive cosmological constant $`\mathrm{\Lambda }`$. In the calculations, we must observe that the time coordinate $`t`$ fails to be a synchronized for $`k0`$, but it is obvious from the context how to obtain the synchronized time from it.
In the coordinates (A2), there is a horizon at $`r=r_0\frac{1}{\sqrt{k}}`$. So, our approach makes sense in the interval $`0<r<r_0`$. However, $`r_0`$ shall be quite large in comparison with the system under consideration, so that we do not meet a problem here.
Now, the following ansatz seems appropriate:
$`ds^2=(1kr^2)(12\phi )dt^2+\left({\displaystyle \frac{dr^2}{1kr^2}}+r^2d\mathrm{\Omega }^2\right)(1+2\psi ),`$ (A3)
where $`\phi `$ and $`\psi `$ depend on the spatial coordinates only. The weak-field assumption allows us to make linearization with respect to $`\phi `$ and $`\psi `$. An extended matter configuration can be obtained by superposition of point particles, so we only need to solve the problem for a $`\delta `$-source at $`r=0`$. This one is spherically symmetric, so we may assume $`\phi =\phi (r)`$ and $`\psi =\psi (r)`$ in Eq. (A3). For the metric components we get:
$$g_{00}=(1kr^2)(12\phi ),g_{11}=\frac{1+2\psi }{1kr^2},g_{22}=r^2(1+2\psi ),g_{33}=g_{22}\mathrm{sin}^2\chi .$$
(A4)
The inverted components are up to linear order in $`\phi `$ and $`\psi `$:
$$g^{00}=\frac{1+2\phi }{1kr^2},g^{11}=(1kr^2)(12\psi ),g^{22}=\frac{12\psi }{r^2},g^{33}=g^{22}\mathrm{sin}^2\chi ,$$
(A5)
which gives the Christoffel symbols
$`\mathrm{\Gamma }_{01}^0`$ $`=`$ $`\phi ^{}{\displaystyle \frac{kr}{1kr^2}},`$ (A6)
$`\mathrm{\Gamma }_{00}^1`$ $`=`$ $`(1kr^2)\left[kr+2kr(\phi +\psi )\phi ^{}(1kr^2)\right],`$ (A7)
$`\mathrm{\Gamma }_{11}^1`$ $`=`$ $`\psi ^{}+{\displaystyle \frac{kr}{1kr^2}},`$ (A8)
$`\mathrm{\Gamma }_{12}^2`$ $`=`$ $`\mathrm{\Gamma }_{13}^3=\psi ^{}+{\displaystyle \frac{1}{r}},`$ (A9)
$`\mathrm{\Gamma }_{22}^1`$ $`=`$ $`r(1kr^2)\psi ^{}r^2(1kr^2),`$ (A10)
$`\mathrm{\Gamma }_{33}^1`$ $`=`$ $`\mathrm{sin}^2\chi \mathrm{\Gamma }_{22}^1,`$ (A11)
$`\mathrm{\Gamma }_{32}^3`$ $`=`$ $`\mathrm{cot}\chi ,`$ (A12)
$`\mathrm{\Gamma }_{33}^2`$ $`=`$ $`\mathrm{sin}\chi \mathrm{cos}\chi ,`$ (A13)
and the Ricci tensor reads
$`R_{00}`$ $`=`$ $`3k(1kr^2)\phi ^{\prime \prime }(1kr^2)^2{\displaystyle \frac{2\phi ^{}}{r}}(1kr^2)+6k(\phi +\psi )(1kr^2)+kr(1kr^2)(5\phi ^{}\psi ^{}),`$ (A14)
$`R_{11}`$ $`=`$ $`2\psi ^{\prime \prime }+\phi ^{\prime \prime }{\displaystyle \frac{2}{r}}\psi ^{}+{\displaystyle \frac{3k}{1kr^2}}+{\displaystyle \frac{kr}{1kr^2}}(\psi ^{}3\phi ^{}),`$ (A15)
$`R_{22}`$ $`=`$ $`3kr^2\psi ^{\prime \prime }r^2(1kr^2)\psi ^{}(2r4kr^3)+(\phi ^{}\psi ^{})(rkr^3),`$ (A16)
$`R_{33}`$ $`=`$ $`R_{22}\mathrm{sin}^2\chi .`$ (A17)
Before we discuss these equations, we consider two obvious limits:
For $`k=0`$, we see that $`R_{00}=\phi ^{\prime \prime }2\phi ^{}/r=\mathrm{\Delta }\phi `$, leading to the usual Newtonian limit $`\mathrm{\Delta }\phi =4\pi G\rho `$.
For $`\phi =\psi =0`$ we get for the Ricci tensor:
$$R_0^0=R_1^1=R_2^2=R_3^3=3k,$$
(A18)
and thus the de Sitter spacetime with $`R=12k`$ for $`k>0`$.
Returning to the general case we have
$$\frac{R}{2}=6k12k\psi +(\phi ^{\prime \prime }2\psi ^{\prime \prime })(1kr^2)+\frac{2}{r}\phi ^{}5kr\phi ^{}\frac{4}{r}\psi ^{}+7kr\psi ^{}$$
(A19)
and then
$$R_0^0\frac{R}{2}=3k+6k\psi +2\psi ^{\prime \prime }(1kr^2)6kr\psi ^{}+\frac{4}{r}\psi ^{}.$$
(A20)
The other components have a similar structure and can be calculated easily from the above equations. The first term of the r.h.s., $`3k`$, will be compensated by the $`\mathrm{\Lambda }`$-term. The usual gauging to $`\psi 0`$ and $`\phi 0`$ as $`r\mathrm{}`$ is no more possible because for $`r>r_0`$ our approximation is no more valid. As an alternative gauge we add such constant values to $`\psi `$ and $`\phi `$ that they are approximately zero in the region under consideration. So we may disregard the term $`6k\psi `$. All remaining terms with $`k`$ can be obtained from those without $`k`$ by multiplying with factors of the type $`1+ϵ`$ where $`ϵkr^2`$, $`k=1/r_0^2`$, with $`r_0`$ being of the order of magnitude of the world radius. In a first approximation, this gives only a small correction to the gravitational constant. In a second approximation, there are deviations from the $`1/r`$-behavior.
An analogous discussion for the Lagrangian $`R+l^2R^2`$ tells us that in a range where $`lrr_0`$, the potential behaves like $`(1c_1e^{r/l})/r`$, as in flat space.
## B The absence of ghosts and tachyons
Here we show in more details what has been stated after Eq. (19). In the conformally transformed picture with a scalar field, the absence of tachyons (i.e., particles with wrong sign in front of the potential term) becomes clear from the form of the potential. For checking ghosts (i.e., particles with wrong sign in front of the kinetic term) we have to go a little more into the details: In Stelle the particle content of fourth order gravity with terms up to quadratic order has been determined, and the existence/absence of ghosts and tachyons has been given in dependence on the free constants of the theory. In the first of Refs. , the analogous calculation as in has been done for a term $`R^3`$ added to the Einstein–Hilbert-Lagragian. Let us give here the argument for general $`n3`$: If $`R^n`$ is in $``$, then the term $`R^{n1}`$ and its derivatives are in the corresponding expression after variational derivative with respect to the metric. In the result, all terms represent products of at least $`n1`$ small quantities; because of $`n3`$ these are always at least two factors; thus, they all vanish in the linearization about the Minkowski space–time.
Now, one might be tempted to require the analogous linearization properties for a Friedmann–Robertson–Walker background. However, linearization around other than flat space–times is not at all a trivial task, see , even for Einstein’s theory: For the closed Friedmann model, Einstein’s theory is linearization unstable, for spatially flat models it is stable, and for the open Friedmann model the result is – contrary to other claims in the older literature – not yet known. We face the further problem that linearization around the de Sitter space-time is complicated to determine, because the same geometry can be locally represented as a spatially flat as well as a closed Friedmann model. So, we leave the question of linearization stability with non-flat background of our model unanswered.
Another type of reasoning was given quite recently: In the possibility has been discussed that the contributions to the Lagrangian coming of gravitons on the one hand and of gravitinos on the other may cancel each other to avoid the ghost problem.
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# The interference of a nonclassical light pulse with a coherent one and the sub-Poissonian statistics formation.
## I Introduction
The formation of ultrashort light pulses (USPs) with suppressed photon fluctuations remains in the focus of considerable attention. The formation and application of USPs in a nonclassical state make possible to combine in experiments a high time resolution with a low level of fluctuations. There are some methods to obtain the USPs in the nonclassical state.
Parametric amplification is a technique that is most extensively used for the production of the USPs in the nonclassical state. In the case of degenerate three-frequency parametric amplification, quadrature-squeezed light is produced. However, this light is found to have super-Poissonian photon statistics directly at the output of an amplifier, and needs interferometers to transform it to one with sub-Poissonian statistics.
One of the most interesting methods for production of USPs in the nonclassical state is the self-phase modulation (SPM) in a nonlinear inertial medium . In this case the USP in the quadrature-squeezing state can be formed with conservation of the photon statistics . The SPM itself is not accompanied by a change in photon statistics. With the aid of nonlinear optical devices in the presence of the SPM one can obtain light with sub-Poissonian photon statistics (see ).
For the first time in a simple method for the production of the USPs with suppressed photon number fluctuations was considered, which is based on the SPM of a USP in a nonlinear inertial medium and subsequent transmission of a pulse through a dispersive optical element. The accurate calculation of this process was conducted making use of technique developed in .
In the present paper, the interference of the SPM-USP with the coherent USP is analysed. The interference is realized using the optical beam-splitter and the process under consideration is analysed in the framework of the quantum theory of SPM of USPs developed in .
The aim of the present article is to analyse the spectra of quantum fluctuations of quadrature components and the spectra of quantum fluctuations of the photon number of the pulses at the outputs of an optical beam-splitter. An extended Mandel parameter will be introduced in order to analyse the photon statistics at all frequencies. The modulation of the total photon number will be analysed too.
## II Spectra of quantum fluctuations of quadratures of SPM-USPs
For the first time in the consequently quantum theory of self-action of USPs in nonlinear inertial medium based on the algebra of time-dependent Bose-operators has been developed. When analysed from the quantum point of view, SPM of USPs is described by the expression
$$\widehat{A}(t,z)=e^{\widehat{O}(t)}\widehat{A}_0(t),$$
(1)
where $`\widehat{A}(t,z)`$ is the annihilation Bose-operator in the given cross section $`z`$ of the nonlinear inertial medium, $`\widehat{O}(t)=i\gamma q[\widehat{n}_0(t)]`$ and $`\gamma q[\widehat{n}_0(t)]`$ is nonlinear phase incursion. The permutation relation takes place:
$$\widehat{A}_0(t_1)e^{\widehat{O}(t_2)}=e^{\widehat{O}(t_2)+G(t_2t_1)}\widehat{A}_0(t_1),$$
(2)
and the theorem of the normal ordering is valid
$$e^{\widehat{O}(t)}=\widehat{𝐍}\mathrm{exp}\left\{_{\mathrm{}}^{\mathrm{}}\left[e^{(\theta )}1\right]\widehat{n}_0(t\theta \tau _r)𝑑\theta \right\},$$
(3)
where $`G(t_2t_1)=i\gamma h(t_2t_1)`$, $`h(t)=H(|t|)`$, $`(\theta )=i\gamma \stackrel{~}{h}(\theta )`$, $`\stackrel{~}{h}(\theta )=\tau _rh(\theta \tau _r)`$, $`\tau _r`$ is the relaxation time of the nonlinear medium and $`\widehat{𝐍}`$ is the operator of normal ordering. Here $`H(t)`$ represents the nonlinear response function of the nonlinear medium ($`H(t)>0`$ for $`t>0`$ and $`H(t)=0`$ for $`t<0`$). Eqs.(1-3) are written in the moving frame $`t=t^{}z/u`$, $`z=z^{}`$, where $`t^{}`$ is the time measured in co-moving frame and $`u`$ is the speed of a pulse in the nonlinear inertial medium.
The quadrature components of the light pulse are defined by the expressions
$`\widehat{X}(t,z)`$ $`=`$ $`[\widehat{A}(t,z)+\widehat{A}^+(t,z)]/2,`$ (5)
$`\widehat{Y}(t,z)`$ $`=`$ $`[\widehat{A}(t,z)\widehat{A}^+(t,z)]/2i,`$ (6)
and for average values of quadratures we have
$`\widehat{X}(t,z)`$ $`=`$ $`|\alpha _0(t)|e^{\mu (t)}\mathrm{cos}\mathrm{\Phi }(t),`$ (8)
$`\widehat{Y}(t,z)`$ $`=`$ $`|\alpha _0(t)|e^{\mu (t)}\mathrm{sin}\mathrm{\Phi }(t).`$ (9)
In (8-9) we denoted $`\mathrm{\Phi }(t)=\psi (t)+\phi (t)`$, $`\psi (t)=2\gamma |\alpha _0(t)|^2=2\gamma \overline{n}_0(t)`$ and $`\mu (t)=\gamma ^2\overline{n}_0(t)=\gamma \psi (t)/2`$. The correlation functions of the quadrature components must be introduced as
$`R_X(t,t+\tau )={\displaystyle \frac{1}{2}}`$ $`[`$ $`\widehat{X}(t,z)\widehat{X}(t+\tau ,z)+\widehat{X}(t+\tau ,z)\widehat{X}(t,z)`$ (11)
$``$ $`2\widehat{X}(t,z)\widehat{X}(t+\tau ,z)],`$ (12)
$`R_Y(t,t+\tau )={\displaystyle \frac{1}{2}}`$ $`[`$ $`\widehat{Y}(t,z)\widehat{Y}(t+\tau ,z)+\widehat{Y}(t+\tau ,z)\widehat{Y}(t,z)`$ (13)
$``$ $`2\widehat{Y}(t,z)\widehat{Y}(t+\tau ,z)],`$ (14)
where $`\tau =t_1t`$ ($`t_1`$ is an arbitrary moment of time, $`t_1>t`$). Using algebra of time-dependent Bose-operators (see (2-3)) for the correlation functions we find
$`R_X(t,t+\tau )={\displaystyle \frac{1}{4}}\{\delta (\tau )`$ $``$ $`\psi (t)h(\tau )\mathrm{sin}2\mathrm{\Phi }(t)`$ (16)
$`+`$ $`\psi ^2(t)g(\tau )\mathrm{sin}^2\mathrm{\Phi }(t)\},`$ (17)
$`R_Y(t,t+\tau )={\displaystyle \frac{1}{4}}\{\delta (\tau )`$ $`+`$ $`\psi (t)h(\tau )\mathrm{sin}2\mathrm{\Phi }(t)`$ (18)
$`+`$ $`\psi ^2(t)g(\tau )\mathrm{cos}^2\mathrm{\Phi }(t)\},`$ (19)
where we denoted $`h(\tau )=\tau _r^1\mathrm{exp}\{|\tau |/\tau _r\}`$, $`g(\tau )=\tau _r^1(1+|\tau |/\tau _r)\mathrm{exp}\{|\tau |/\tau _r\}`$ and considered that nonlinear inertial medium is of a Kerr type. To get the expressions (17-19), the approximations $`\tau _r\tau _p`$ , $`\gamma 1`$ have been used ($`\tau _p`$ is the pulse duration). The spectra of the fluctuations of quadrature components are
$`S_X(\mathrm{\Omega },t)={\displaystyle \frac{1}{4}}[1`$ $``$ $`2\psi (t)L(\mathrm{\Omega })\mathrm{sin}2\mathrm{\Phi }(t)`$ (22)
$`+4\psi ^2(t)L^2(\mathrm{\Omega })\mathrm{sin}^2\mathrm{\Phi }(t)],`$
$`S_Y(\mathrm{\Omega },t)={\displaystyle \frac{1}{4}}[1`$ $`+`$ $`2\psi (t)L(\mathrm{\Omega })\mathrm{sin}2\mathrm{\Phi }(t)`$ (24)
$`+4\psi ^2(t)L^2(\mathrm{\Omega })\mathrm{cos}^2\mathrm{\Phi }(t)],`$
where we denoted $`\mathrm{\Omega }=\omega \tau _r`$ and $`L(\mathrm{\Omega })=1/[1+\mathrm{\Omega }^2]`$. At initial phase chosen optimal at the frequency $`\mathrm{\Omega }_0`$ (see )
$$\phi _0(t)=\frac{1}{2}\mathrm{arctan}\left\{\frac{1}{\psi (t)L(\mathrm{\Omega }_0)}\right\}\psi (t)$$
(25)
the spectra (22-24) have the forms
$`S_X(\mathrm{\Omega }_0,t)`$ $`=`$ $`{\displaystyle \frac{1}{4}}[\sqrt{1+\psi ^2(t)L^2(\mathrm{\Omega }_0)}\psi (t)L(\mathrm{\Omega }_0)]^2,`$ (27)
$`S_Y(\mathrm{\Omega }_0,t)`$ $`=`$ $`{\displaystyle \frac{1}{4}}[\sqrt{1+\psi ^2(t)L^2(\mathrm{\Omega }_0)}+\psi (t)L(\mathrm{\Omega }_0)]^2.`$ (28)
At any frequency $`\mathrm{\Omega }`$ the spectra (22-24) are given by:
$`S_X(\mathrm{\Omega },t)`$ $`=`$ $`S_X(\mathrm{\Omega }_0,t)+{\displaystyle \frac{1}{2}}\psi (t)[L(\mathrm{\Omega })L(\mathrm{\Omega }_0)]`$ (32)
$`\times \{[L(\mathrm{\Omega })+L(\mathrm{\Omega }_0)]\psi (t)[1+(L(\mathrm{\Omega })+L(\mathrm{\Omega }_0))`$
$`\times L(\mathrm{\Omega }_0)\psi ^2(t)][1+\psi ^2(t)L^2(\mathrm{\Omega })]^{1/2}\},`$
$`S_Y(\mathrm{\Omega },t)`$ $`=`$ $`S_Y(\mathrm{\Omega }_0,t)+{\displaystyle \frac{1}{2}}\psi (t)[L(\mathrm{\Omega })L(\mathrm{\Omega }_0)]`$ (35)
$`\times \{[L(\mathrm{\Omega })+L(\mathrm{\Omega }_0)]\psi (t)+[1+(L(\mathrm{\Omega })+L(\mathrm{\Omega }_0))`$
$`\times L(\mathrm{\Omega }_0)\psi ^2(t)][1+\psi ^2(t)L^2(\mathrm{\Omega })]^{1/2}\}.`$
The spectra of the quantum fluctuations of squeezed $`X`$-quadrature at the initial phase chosen optimal at the reduced frequency $`\mathrm{\Omega }_0=0`$ are displayed in Fig.1.
From Fig.1 one can see that the squeezing of quantum fluctuations is greatest at the frequency for which the phase of the initial pulse was chosen optimal. In addition, in have been shown that the spectra of quantum fluctuations of quadrature components can be controlled by the choice of the phase of the initial light pulse and that, at the initial phase chosen optimal at the frequency $`\mathrm{\Omega }_0=1`$, the squeezing of quadrature component $`X`$ is maximum at frequencies $`\omega 1/\tau _r`$ for $`\psi (t)>1`$.
The choice of the optimal phase as (25) means that in the quadrature’s space $`XY`$ the big axis of the eclipse of squeezing is parallel with the $`Y`$ axis and consequently, the squeezing of the $`X`$ quadrature is maximum (see (32)). The transformation of the uncertainty region of quadrature fluctuations, from the circle for the initial coherent light pulse to the ellipse for the SPM-USP, always take place. The choice of the phase as (25) means the orientation of the system of quadrature squeezing observation $`XY`$ as the quantum fluctuations of $`X`$ quadrature are maximum suppressed.
## III The interference of SPM-USP with a coherent one
In what follows we are interested in analysing of the interference between the SPM-USP and coherent one. As a theoretical model we use the model of the symmetric beam-splitter (double reflection angle of the incident SPM-USP is $`\pi /2`$). The interference scheme is presented in Fig.2. The input-output operator transformations at the output # $`1`$ are given by
$`\widehat{B}_1(t)`$ $`=`$ $`i\sqrt{R}\widehat{A}_1(t,l)+\sqrt{T}\widehat{A}_2(t),`$ (37)
$`\widehat{B}_1^+(t)`$ $`=`$ $`i\sqrt{R}\widehat{A}_1^+(t,l)+\sqrt{T}\widehat{A}_2(t),`$ (38)
and at the output # 2 they are
$`\widehat{B}_2(t)`$ $`=`$ $`\sqrt{T}\widehat{A}_1(t,l)+i\sqrt{R}\widehat{A}_2(t),`$ (40)
$`\widehat{B}_2^+(t)`$ $`=`$ $`\sqrt{T}\widehat{A}_1^+(t,l)i\sqrt{R}\widehat{A}_2^+(t),`$ (41)
where $`l`$ is the length of the nonlinear inertial medium and $`R`$ is the coefficient of reflection, $`R+T=1`$.
### A The spectra of quantum fluctuations of quadratures at the beam-splitter output # $`1`$
We define the quadrature components at the output # $`1`$ of the beam splitter as (see (5-6))
$`\widehat{𝒳}_1(t)`$ $`=`$ $`[\widehat{B}_1(t)+\widehat{B}_1^+(t)]/2,`$ (43)
$`\widehat{𝒴}_1(t)`$ $`=`$ $`[\widehat{B}_1(t)\widehat{B}_1^+(t)]/2i.`$ (44)
For average values of the quadrature components (43-44) we find
$`\widehat{𝒳}_1(t)`$ $`=`$ $`|\alpha _0(t)|e^{\mu (t)}[\sqrt{T}\mathrm{cos}\phi _2(t)\sqrt{R}\mathrm{sin}\mathrm{\Phi }_1(t)],`$ (46)
$`\widehat{𝒴}_1(t)`$ $`=`$ $`|\alpha _0(t)|e^{\mu (t)}[\sqrt{T}\mathrm{sin}\phi _2(t)+\sqrt{R}\mathrm{cos}\mathrm{\Phi }_1(t)],`$ (47)
where $`\mathrm{\Phi }_1(t)=\psi (t)+\phi _1(t)`$, $`\psi (t)=2\gamma \overline{n}_1(t)`$, $`\mu (t)=\gamma \psi (t)/2`$ and $`\phi _1(t)`$ is the initial phase of the pulse in nonlinear section. For the correlation functions of the quadrature components we have
$`R_{𝒳_1}(t,t+\tau )={\displaystyle \frac{1}{4}}\{\delta (\tau )`$ $`+`$ $`R\psi (t)h(\tau )\mathrm{sin}2\mathrm{\Phi }_1(t)`$ (49)
$`+`$ $`R\psi ^2(t)g(\tau )\mathrm{cos}^2\mathrm{\Phi }_1(t)\},`$ (50)
$`R_{𝒴_1}(t,t+\tau )={\displaystyle \frac{1}{4}}\{\delta (\tau )`$ $``$ $`R\psi (t)h(\tau )\mathrm{sin}2\mathrm{\Phi }_1(t)`$ (51)
$`+`$ $`R\psi ^2(t)g(\tau )\mathrm{cos}^2\mathrm{\Phi }_1(t)\}.`$ (52)
For the spectra of quadrature fluctuations we find
$`S_{𝒳_1}(\mathrm{\Omega },t)={\displaystyle \frac{1}{4}}[1`$ $`+`$ $`2R\psi (t)L(\mathrm{\Omega })\mathrm{sin}2\mathrm{\Phi }_1(t)`$ (55)
$`+4R\psi ^2(t)L^2(\mathrm{\Omega })\mathrm{cos}^2\mathrm{\Phi }_1(t)],`$
$`S_{𝒴_1}(\mathrm{\Omega },t)={\displaystyle \frac{1}{4}}[1`$ $``$ $`2R\psi (t)L(\mathrm{\Omega })\mathrm{sin}2\mathrm{\Phi }_1(t)`$ (57)
$`+4R\psi ^2(t)L^2(\mathrm{\Omega })\mathrm{sin}^2\mathrm{\Phi }_1(t)].`$
At optimal phase (25) the spectra (55-57) take the forms
$`S_{𝒳_1}(\mathrm{\Omega }_0,t)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\{[\sqrt{1+\psi ^2(t)L^2(\mathrm{\Omega }_0)}+R\psi (t)L(\mathrm{\Omega }_0)]^2`$ (60)
$`(2R1)^2\psi ^2(t)L^2(\mathrm{\Omega }_0)\},`$
$`S_{𝒴_1}(\mathrm{\Omega }_0,t)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\{[\sqrt{1+\psi ^2(t)L^2(\mathrm{\Omega }_0)}R\psi (t)L(\mathrm{\Omega }_0)]^2`$ (62)
$`(2R1)^2\psi ^2(t)L^2(\mathrm{\Omega }_0)\}.`$
At any frequency $`\mathrm{\Omega }`$ we have
$`S_{𝒳_1}(\mathrm{\Omega },t)`$ $`=`$ $`S_{𝒳_1}(\mathrm{\Omega }_0,t)+{\displaystyle \frac{1}{2}}R\psi (t)[L(\mathrm{\Omega })L(\mathrm{\Omega }_0)]`$ (66)
$`\times \{[L(\mathrm{\Omega })+L(\mathrm{\Omega }_0)]\psi (t)+[1+(L(\mathrm{\Omega })+L(\mathrm{\Omega }_0))`$
$`\times L(\mathrm{\Omega }_0)\psi ^2(t)][1+\psi ^2(t)L^2(\mathrm{\Omega }_0)]^1\},`$
$`S_{𝒴_1}(\mathrm{\Omega },t)`$ $`=`$ $`S_{𝒴_1}(\mathrm{\Omega }_0,t)+{\displaystyle \frac{1}{2}}R\psi (t)[L(\mathrm{\Omega })L(\mathrm{\Omega }_0)]`$ (69)
$`\times \{[L(\mathrm{\Omega })+L(\mathrm{\Omega }_0)]\psi (t)[1+(L(\mathrm{\Omega })+L(\mathrm{\Omega }_0))`$
$`\times L(\mathrm{\Omega }_0)\psi ^2(t)][1+\psi ^2(t)L^2(\mathrm{\Omega }_0)]^1\}.`$
As one can see from (66-69) the interference of the SPM-USP with the coherent USP does not influence the quadrature squeezing of the SPM-USP. The choice of the coefficient of reflection $`R`$ of the beam-splitter allows us to receive the spectra of squeezing (see (66-69)) of interest to us. One can conclude that, the lost of photons from SPM-USP as a result of division in the beam-splitter is follows by the reduction of squeezing. In fact this lost is compensated by the coherent USP. The spectra of quantum fluctuations of $`𝒴_1`$ quadrature at time $`t=0`$ as a function of maximum nonlinear phase $`\psi (0)=\psi _0`$ and the reduced frequency $`\mathrm{\Omega }`$ in the case of the $`50\%`$ beam-splitter ($`R=T=1/2`$) is displayed in Fig.3. On Fig.3 one can see that, in this case the squeezing is reduced at one half.
Also, it is important to remark that, as a result of symmetric reflection the squeezing from the $`X`$ quadrature is moved in the $`Y`$ quadrature. This means that the ellipse of squeezing is moved in the $`XY`$ frame with an angle equal to double reflection angle. From theoretical point of view this is expressed by the presence of the index $`i=\mathrm{exp}(i\pi /2)`$ in the (37-38) and (40-41). The double angle of reflection can be interpreted as a geometrical phase. The choice of the position of the beam-splitter relative to direction of propagation of the SPM-USP allows us to control the position of the ellipse of squeezing in the $`XY`$ frame.
### B The spectra of quantum fluctuations of quadratures at the beam-splitter output # $`2`$
At the output # $`2`$ the quadrature components are defined as
$`\widehat{𝒳}_2(t)`$ $`=`$ $`[\widehat{B}_2(t)+\widehat{B}_2^+(t)]/2,`$ (71)
$`\widehat{𝒴}_2(t)`$ $`=`$ $`[\widehat{B}_2(t)\widehat{B}_2^+(t)]/2i.`$ (72)
For the average values of quadrature components we find
$`\widehat{𝒳}_2(t)`$ $`=`$ $`|\alpha _0(t)|e^{\mu (t)}[\sqrt{T}\mathrm{cos}\mathrm{\Phi }_1(t)\sqrt{R}\mathrm{sin}\phi _2(t)],`$ (74)
$`\widehat{𝒴}_2(t)`$ $`=`$ $`|\alpha _0(t)|e^{\mu (t)}[\sqrt{T}\mathrm{sin}\mathrm{\Phi }_1(t)+\sqrt{R}\mathrm{cos}\phi _2(t)].`$ (75)
Leaving out the preliminary accounts for the correlation functions of the quadrature components we get
$`R_{𝒳_2}(t,t+\tau )={\displaystyle \frac{1}{4}}\{\delta (\tau )`$ $``$ $`T\psi (t)h(\tau )\mathrm{sin}2\mathrm{\Phi }_1(t)`$ (77)
$`+`$ $`T\psi ^2(t)g(\tau )\mathrm{cos}^2\mathrm{\Phi }_1(t)\},`$ (78)
$`R_{𝒴_2}(t,t+\tau )={\displaystyle \frac{1}{4}}\{\delta (\tau )`$ $`+`$ $`T\psi (t)h(\tau )\mathrm{sin}2\mathrm{\Phi }_1(t)`$ (79)
$`+`$ $`T\psi ^2(t)g(\tau )\mathrm{cos}^2\mathrm{\Phi }_1(t)\}.`$ (80)
In consequence, for the spectra of quantum fluctuations of quadratures we have
$`S_{𝒳_2}(\mathrm{\Omega },t)={\displaystyle \frac{1}{4}}[1`$ $``$ $`2T\psi (t)L(\mathrm{\Omega })\mathrm{sin}2\mathrm{\Phi }_1(t)`$ (83)
$`+4T\psi ^2(t)L^2(\mathrm{\Omega })\mathrm{sin}^2\mathrm{\Phi }_1(t)],`$
$`S_{𝒴_2}(\mathrm{\Omega },t)={\displaystyle \frac{1}{4}}[1`$ $`+`$ $`2T\psi (t)L(\mathrm{\Omega })\mathrm{sin}2\mathrm{\Phi }_1(t)`$ (85)
$`+4T\psi ^2(t)L^2(\mathrm{\Omega })\mathrm{cos}^2\mathrm{\Phi }_1(t)].`$
The spectra (83-85) at optimal phase (25) are:
$`S_{𝒳_2}(\mathrm{\Omega }_0,t)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\{[\sqrt{1+\psi ^2(t)L^2(\mathrm{\Omega }_0)}T\psi (t)L(\mathrm{\Omega }_0)]^2`$ (88)
$`(2T1)^2\psi ^2(t)L^2(\mathrm{\Omega }_0)\},`$
$`S_{𝒴_2}(\mathrm{\Omega }_0,t)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\{[\sqrt{1+\psi ^2(t)L^2(\mathrm{\Omega }_0)}+T\psi (t)L(\mathrm{\Omega }_0)]^2`$ (90)
$`(2T1)^2\psi ^2(t)L^2(\mathrm{\Omega }_0)\}.`$
At any frequency $`\mathrm{\Omega }`$ the spectra (83-85) take the forms
$`S_{𝒳_2}(\mathrm{\Omega },t)`$ $`=`$ $`S_{𝒳_2}(\mathrm{\Omega }_0,t)+{\displaystyle \frac{1}{2}}T\psi (t)[L(\mathrm{\Omega })L(\mathrm{\Omega }_0)]`$ (94)
$`\times \{[L(\mathrm{\Omega })+L(\mathrm{\Omega }_0)]\psi (t)[1+(L(\mathrm{\Omega })+L(\mathrm{\Omega }_0))`$
$`\times L(\mathrm{\Omega }_0)\psi ^2(t)][1+\psi ^2(t)L^2(\mathrm{\Omega }_0)]^1\},`$
$`S_{𝒴_2}(\mathrm{\Omega },t)`$ $`=`$ $`S_{𝒴_2}(\mathrm{\Omega }_0,t)+{\displaystyle \frac{1}{2}}T\psi (t)[L(\mathrm{\Omega })L(\mathrm{\Omega }_0)]`$ (97)
$`\times \{[L(\mathrm{\Omega })+L(\mathrm{\Omega }_0)]\psi (t)+[1+(L(\mathrm{\Omega })+L(\mathrm{\Omega }_0))`$
$`\times L(\mathrm{\Omega }_0)\psi ^2(t)][1+\psi ^2(t)L^2(\mathrm{\Omega }_0)]^1\}.`$
As already have been remarked in the previous analyse, the spectra of quadrature fluctuations can be controlled by the choice of the coefficients of the beam-splitter (see (66-69)). It is interesting to remark that, in case of the symmetric $`50\%`$ beam-splitter, the squeezing of quadrature fluctuations at output # $`2`$ is present only in the $`𝒳_2`$ quadrature and it is reduced at one half (see (94-97)). Since the geometrical phase is equal to $`0`$ in the analysed case (the refracted part of the SPM-USP does not change the direction of propagation in comparison with the initial SPM-USP) the ellipse of squeezing does not change this initial position in the $`XY`$ frame and has the big axis parallel to $`Y`$ axis.
## IV The spectra of fluctuations of photon number and the statistics
We introduce the correlation function of photon number at the output number $`j`$ ($`j=\overline{1,2}`$) in the following symmetric form
$`R_{N,j}(t_1,t_2)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[\widehat{N}_j(t_1)\widehat{N}_j(t_2)+\widehat{N}_j(t_2)\widehat{N}_j(t_1)`$ (99)
$`2\widehat{N}_j(t_1)\widehat{N}_j(t_2)].`$
To simplify the accounts we consider that the initial pulses have the identical average values of the photon number ($`\overline{n}_1(t)=\overline{n}_2(t)=\overline{n}_0(t)`$). Using the algebra of time-dependent Bose-operators (2-3) for the correlation functions we get
$`R_{N,1}(t,t+\tau )`$ $`=`$ $`\overline{n}_0(t)[\delta (\tau )2R\sqrt{RT}\psi (t)h(\tau )\mathrm{cos}\stackrel{~}{\mathrm{\Phi }}(t)`$ (103)
$`+RT\psi (t)h(\tau )\mathrm{sin}2\stackrel{~}{\mathrm{\Phi }}(t)`$
$`+RT\psi ^2(t)g(\tau )\mathrm{cos}^2\stackrel{~}{\mathrm{\Phi }}(t)],`$
$`R_{N,2}(t,t+\tau )`$ $`=`$ $`\overline{n}_0(t)[\delta (\tau )+2T\sqrt{RT}\psi (t)h(\tau )\mathrm{cos}\stackrel{~}{\mathrm{\Phi }}(t)`$ (106)
$`+RT\psi (t)h(\tau )\mathrm{sin}2\stackrel{~}{\mathrm{\Phi }}(t)`$
$`+RT\psi ^2(t)g(\tau )\mathrm{cos}^2\stackrel{~}{\mathrm{\Phi }}(t)].`$
where in (103-106) is denoted $`\tau =t_1t`$ (see (12-14)), $`\stackrel{~}{\mathrm{\Phi }}(t)=\psi (t)+\mathrm{\Delta }\phi (t)`$, $`\mathrm{\Delta }\phi (t)=\phi _1(t)\phi _2(t)`$. In the case of the symmetrical $`50\%`$ beam-splitter, the spectra of quantum fluctuations of the photon number at the measured time $`𝒯`$ have the forms
$`S_{𝒯,1}(\mathrm{\Omega },t)`$ $`=`$ $`\overline{𝒩}+\overline{n}_0(t)\psi (t)L(\mathrm{\Omega })\left({\displaystyle \frac{𝒯}{\tau _p}}\right)[\mathrm{cos}\stackrel{~}{\mathrm{\Phi }}(t)`$ (109)
$`+{\displaystyle \frac{1}{2}}\mathrm{sin}2\stackrel{~}{\mathrm{\Phi }}(t)+\psi (t)L(\mathrm{\Omega })\mathrm{cos}^2\stackrel{~}{\mathrm{\Phi }}(t)],`$
$`S_{𝒯,2}(\mathrm{\Omega },t)`$ $`=`$ $`\overline{𝒩}+\overline{n}_0(t)\psi (t)L(\mathrm{\Omega })\left({\displaystyle \frac{𝒯}{\tau _p}}\right)[\mathrm{cos}\stackrel{~}{\mathrm{\Phi }}(t)`$ (111)
$`+{\displaystyle \frac{1}{2}}\mathrm{sin}2\stackrel{~}{\mathrm{\Phi }}(t)+\psi (t)L(\mathrm{\Omega })\mathrm{cos}^2\stackrel{~}{\mathrm{\Phi }}(t)],`$
where $`\overline{𝒩}=_{\mathrm{}}^{\mathrm{}}\overline{n}_0(t)𝑑t`$ is the total photon number of the initial coherent USP. An extended Mandel parameter at all frequencies describing the photon statistics at measurement time $`𝒯`$ can be introduced as
$`Q_{𝒯,1}(\mathrm{\Omega },t)`$ $`=`$ $`\epsilon _{𝒯,1}(\mathrm{\Omega },t)/\overline{n}_0(t)`$ (113)
$`=`$ $`\psi (t)L(\mathrm{\Omega })\left({\displaystyle \frac{𝒯}{\tau _p}}\right)[\mathrm{cos}\stackrel{~}{\mathrm{\Phi }}(t)+{\displaystyle \frac{1}{2}}\mathrm{sin}2\stackrel{~}{\mathrm{\Phi }}(t)`$ (115)
$`+\psi (t)L(\mathrm{\Omega })\mathrm{cos}^2\stackrel{~}{\mathrm{\Phi }}(t)],`$
$`Q_{𝒯,2}(\mathrm{\Omega },t)`$ $`=`$ $`\epsilon _{𝒯,2}(\mathrm{\Omega },t)/\overline{n}_0(t)`$ (116)
$`=`$ $`\psi (t)L(\mathrm{\Omega })\left({\displaystyle \frac{𝒯}{\tau _p}}\right)[\mathrm{cos}\stackrel{~}{\mathrm{\Phi }}(t)+{\displaystyle \frac{1}{2}}\mathrm{sin}2\stackrel{~}{\mathrm{\Phi }}(t)`$ (118)
$`+\psi (t)L(\mathrm{\Omega })\mathrm{cos}^2\stackrel{~}{\mathrm{\Phi }}(t)],`$
where in (115-118) we denoted $`\epsilon _{𝒯,j}(\mathrm{\Omega },t)=S_{𝒯,j}(\mathrm{\Omega },t)\overline{𝒩}`$.
The statistics of photon number can be controlled by the choice of the nonlinear phase addition $`\psi (t)`$ and linear phase shift $`\mathrm{}\phi (t)`$. The dependence of the extended Mandel parameter $`Q_{𝒯,1}(\mathrm{\Omega },t)`$ at time $`t=0`$, $`𝒯/\tau _p=1/10`$ and $`\mathrm{}\phi (t)=\pi /2`$ on $`\psi _0`$ and $`\mathrm{\Omega }`$ is displayed in Fig.4.
As follows from Fig.4, for different values of the nonlinear phase addition, the statistics of photon number can be as super-Poissonian as sub-Poissonian. The maximum increase or decrease of photon number fluctuations take place around the central frequency $`\mathrm{\Omega }=0`$. From Fig.4 it follows that in case $`\mathrm{}\phi (t)=\pi /2`$ the first suppression of the photon number fluctuations takes place around $`\psi (t)=3.5`$ and it can be observed at all frequencies from $`0`$ to $`\omega 2/\tau _r`$.
It is important to mention that, if the radiation is monochromatic, the formulas (115-118) remain valid. In this case, there are no second and third terms in (115-118) and the reduced frequency $`\mathrm{\Omega }=0`$, as since all the photons have the same frequency. The presence of these terms is connected with the correlation between modes of SPM-USP (see ) and in consequence, the quantum fluctuations of the photon number for the pulse field will be “smoothed-out” at the frequency $`\mathrm{\Omega }=0`$ in comparison with the monochromatic radiation.
## V The modulation of the total photon number
Let us introduce the photon number operator at the measurement time $`𝒯`$
$$\widehat{𝒩}_{𝒯,j}(t)=_{t𝒯/2}^{t+𝒯/2}\widehat{N}_j(t^{^{}})𝑑t^{^{}},$$
(119)
where
$$\widehat{N}_j(t)=\widehat{B}_j^+(t)\widehat{B}_j(t).$$
(120)
To simplify the accounts, in following we consider that $`\overline{n}_1(t)=\overline{n}_2(t)=\overline{n}_0(t)`$ and then we get
$`\widehat{𝒩}_{𝒯,1}(t)`$ $`=`$ $`𝒯\overline{n}_0(t)\left\{12\sqrt{RT}\mathrm{sin}\left[\psi (t)+\mathrm{}\phi (t)\right]\right\},`$ (122)
$`\widehat{𝒩}_{𝒯,2}(t)`$ $`=`$ $`𝒯\overline{n}_0(t)\left\{1+2\sqrt{RT}\mathrm{sin}\left[\psi (t)+\mathrm{}\phi (t)\right]\right\}.`$ (123)
where $`\overline{n}_0(t)=\overline{n}_0\rho (t)`$ and $`\rho (t)`$ is the envelope of the pulse. As can be seen from (122-123), for different values of the linear phase shift $`\mathrm{}\phi (t)`$ of the initial coherent light pulses, all photons can be distributed only in one of the outputs of the beam-splitter. This property is specific for the quantum interference only and has no classical analogue.
Let us introduce the total number operator at the beam-splitter output # $`1`$ as
$$\widehat{𝒩}_1=_{\mathrm{}}^{\mathrm{}}\widehat{N}_1(t)𝑑t,$$
(124)
and for its average value we obtain
$$\overline{𝒩}_1=\overline{𝒩}\left[12\sqrt{RT}\frac{\mathrm{sin}(\psi _0/4)}{\psi _0/4}\mathrm{sin}\left(\frac{\psi _0}{4}+\mathrm{}\phi \right)\right].$$
(125)
To get (125) we considered that the initial pulses are quite Gaussian, their envelopes having the form $`\rho (t)\mathrm{exp}\left\{t^2/2\tau _p^2\right\}/\sqrt{2}`$, and that $`\mathrm{}\phi (t)=\mathrm{}\phi `$ is a constant.
The dependence of $`\overline{𝒩}_1/\overline{𝒩}`$ at maximum nonlinear phase addition $`\psi _0`$ for 50% beam-splitter and $`\mathrm{}\phi =0`$ is displayed in Fig.5.
As follows from (125) the SPM process leads to the additional modulation of the total photon number. In consequence, the choice of the maximum nonlinear phase addition $`\psi _0`$ and the initial phase shift $`\mathrm{}\phi `$ allows us to control the total photon number at the outputs of the beam-splitter.
## VI Discussions and conclusions
The analyse of the interference process of the SPM-USP with coherent USP, based on the algebra of time-dependent Bose-operators, allows us to understand in which way the interference leads to the sub- and super-Poissonian photon statistics formation. Analysing the spectra of quantum fluctuations of quadrature components was concluded that, the position of the optical beam-splitter plays an important role in the quadrature squeezing observation. The initial phase, chosen optimal for a determined frequency plays a role of the phase of reference. Its choice means that, the $`XY`$ frame is placed so as the observed squeezing of quadrature $`X`$ is maximum. As a result of reflection of the SPM-USP at the beam-splitter, the ellipse of squeezing will move itself in the $`XY`$ frame with an angle equal to the geometrical phase. The geometrical phase can take the values for which the squeezing cannot be observed in $`X`$ or $`Y`$ quadratures. Let us take into consideration the case in which we get the initial phase of the light pulse in the nonlinear section arbitrary. In this case the ellipse of squeezing is located arbitrary in the $`XY`$ frame and it is possible do not observe the squeezing of quadrature fluctuations in the measurements. In case we chose the initial phase optimal, we locate the ellipse of squeezing in the $`XY`$ frame so as the small axis of the ellipse will lie along the $`X`$ axis. In consequence, the observed squeezing of $`X`$ quadrature fluctuations is maximum. If we do not know the position of the ellipse in the $`XY`$ frame (taking the initial phase arbitrary), then rotating the beam-splitter we can orientate the $`XY`$ frame so as the ellipse of squeezing can be displayed with small axis along $`X`$ axis and the squeezing of the $`X`$ quadrature fluctuations can be observed.
Up to now, it is considered that, after reflection from the beam-splitter, the squeezing of the quadrature fluctuations can be affected and complete disappeared as a result of the vacuum fluctuations participation at the other input of the beam-splitter. This interpretation is not quite correct. It is important to mention that, the presence of vacuum fluctuations was already taken into account when the SPM process in the nonlinear inertial medium was analysed from quantum point of view .
The presented in the recent work analyse gives the indication about the use of the nonclassical light pulses in gravitational wave detection. There are two points of interest. One is represented by the choice of the initial phase of a pulse at the input in nonlinear inertial medium which will determine the position of the ellipse of squeezing of quadrature fluctuations in the $`XY`$ frame. This optimal phase must be interpreted as the phase of reference. Other point is represented by the reflection of the SPM-USP from the surfaces, as since the addition of the geometrical phase will rotate the ellipse in the $`XY`$ frame consequently. These points must be taken into account when the experiments for gravitational wave detection are implemented using the USP in nonclassical state.
The extended Mandel parameter is introduced and the photon statistics is scanned at all frequencies. It is shown that, the statistics of the photon number can be controlled by choice of the nonlinear phase addition and the linear phase shift of the initial coherent light pulses. The choice of the linear phase shift represents an effective method of control of the sub- or super-Poissonian statistics formation and of the modulation of the total photon number.
It is interesting to mention that, the analyse of the interference process between two SPM-USPs can leads to the determination of the evolution of the linear parameters of the ellipse of squeezing as a functions of the nonlinear phase addition. This analyse will be the subject of another publication.
The author is grateful to S. Codoban (JINR, Dubna) for useful discussions and rendered help.
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# On the Nucleon Stability and Σ_𝑁 Term Puzzles
## 1 Introduction
Two puzzles for a nucleon, which is believed to connect to its chiral properties, remain to be explained theoretically. The first one concerns with the mechanism for the stability of a nucleon; the second one is the nucleon $`\mathrm{\Sigma }_N`$ term discrepancy problem, which is made more severe by recent analysis of the pion–nuclon scattering data.
1) Despite the success of the bag and/or non-topological solitonic type of models for a nucleon in the case of a single nucleon, a serious potential problem for these models exists when one tries to describe many nucleon systems, like a nucleus, base on these models. On the one hand, there are empirical evidences which imply that a nucleon inside of a nucleus keeps most of its identity. On the other hand, the theoretical models allow the solution for a nucleon inside a nucleus to change significantly compared to the free space solution to favor a bag or a soliton of multi-quarks rather than the phenomenologically successful picture of a nucleus that can be described using the nucleon and meson degrees of freedom. Albeit arguments for a kind of equivalence between these two pictures were put forward in literature, the quantitative justification of this picture in realistic situations is still lacking. It implies that perhaps certain elements of an essential physics have not been understood so far. It is expected that a study of this possibility using the local theory recently developed could lead to a way out of the dilemma .
2) Due to the approximate chiral symmetry in strong interaction, which is only explicitly broken by a tiny current mass for the light quarks, the low energy dynamics involving pion is very accurately constrained by the chiral symmetry of QCD. Indeed the Goldberger–Treiman relation, the Adler–Weisberger relation and other chiral relations are confirmed by data even with an extension of the PCAC relation , which includes the possibility that there is a vector type virtual color superconducting phase. Such an expected success fails to reappear in the next leading order in the current quark mass. The resulting nucleon $`\sigma _N`$ term, which is a measure of the scalar charge of the nucleon multiplied by the current quark mass, seems to be poorly connected to the on shell pion nucleon compton scattering data. The analysis of the most current data generates even greater values for $`\mathrm{\Sigma }_N`$ on the Cheng–Dashen point than the old one –making the descrepancy larger ($`100\%`$ effect). The lattice QCD evaluation of the same quantity also produces a larger value. Under the conventional chiral symmetry breaking picture, this is indeed very puzzling: why the leading order relations are so accurately confirmed by data and at the same time the next leading order one fails $`100\%`$, given the fact that the expected error should be at most around $`m_\pi /M_A10\%`$ with $`m_\pi `$ the mass of a pion and $`M_A1`$ GeV. Assuming a large strangeness content for a nucleon may not be consistent with other observables for the nucleon . The way to resolve the nucleon $`\mathrm{\Sigma }_N`$ problem without certain new insight about the structure of the nucleon is a challenge if the chiral symmetry is regarded to be only slightly broken at low energy, as it manifests itself in the leading order and in many other observables.
The reason for the leading order results not to be affected by whether or not there is a metastable color superconducting phase is not hard to understand. The quantities that are sensitive to the possible metastable color superconducting phases are evaluated on the pion mass shell, resulting in a pion dominance. The quantity, namely $`g_A`$, which is not evaluated on the pion mass shell is however a quantity that is insensitive to the low energy chiral dynamics. In fact the scale that determines the change of $`g_A`$ with the momentum transfer is controlled by a scale of order $`M_A1`$ GeV . The nucleon $`\mathrm{\Sigma }_N`$ term on the other hand does not enjoy such an insensitivity.
It is therefore interesting to study whether or not these difficulties of the conventional picture for the nucleon could provide a useful source of information for a semi-quantitative discussion of various basic properties of a nucleon in the light of the local theory developed recently.
There are several new features in the local theory. First of all, it is based on an 8–component “real” representation for fermions, which is not equivalent to the currently adopted frame–work for relativistic processes at finite density . Some of the advantages of this representation are discussed in Ref. . Second, the statistical gauge invariance, dark component for local observables and the statistical blocking effects can be studied in the new frame-work . It would be interesting to study what is its implication on a possible simultaneous solution of the aforementioned problems. It serves as a further consistency check of the mechanisms proposed here and in our earlier publications.
Qualitative discussion of such a solution to the nucleon stability problem is proposed recently , in which it is argued that the statistical blocking effects of the strong interaction vacuum state can be the underlying physical mechanism for the stability of a nucleon inside of a nucleus and nuclear matter. It remains to be investigated quantitatively as to whether or not such a mechanism is right in the sense that 1) it respects constraints of the chiral Ward–Takahashi identities and 2) it could reproduce at least two of the most important physical properties of the nucleon, namely, its size and its mass with reasonable model assumptions and parameters. This work aims at a semi-quantitative study of this question. A model picture for the nucleon has to be established and the physics and order of magnitude estimate of the possible quantitative properties of the model could then be derived.
The nucleon $`\mathrm{\Sigma }_N`$ term problem is also not a qualitative problem in the local theory if one assumes that there is at least one metastable color superconducting phase . New developments are made both in the theoretical and phenomenological directions , which indicate that such a possibility can be oberved and further more is favored by the evidences considered. However these studies contains very little specific properties of the nucleon structure. It seems that a more quantitative study of the nucleon $`\mathrm{\Sigma }_N`$ problem is the natural next step.
The paper is organized in the following way. Section 2 contains an introduction of the chiral models for low and intermediate energy strong interaction, a discussion of the vacuum phenomenology in terms of wave function renormalized quasi-particles of the spontaneous chiral symmetry breaking, an introduction of a model for a nucleon and a determination of its size under such a scenario when the value of $`\mathrm{\Sigma }_N`$ is kept constant. A semi-quantitative mechanism for the stability of a nucleon inside of a nucleus and/or nuclear matter is suggested and discussed in Section 3. It is based upon the local theory in which the fermion fields are represented by an 8-component “real” spinor. The possible role played by the metastable color superconducting phases for the strong interaction vacuum state is discussed in section 4. The discrepancy of the nucleon $`\mathrm{\Sigma }_N`$ term extracted from different sources is resolved by assuming the existence of a virtual color superconducting phase for the strong interaction vacuum state, in which, a virtual color superconducting component for a nucleon lives. The difference in energy density between the true chiral symmetry breaking ground state and the virtual color superconducting phase is estimated based on available data. The section 5 contains further discussions. A summary is given in the last section.
## 2 The Global $`\mathrm{\Sigma }_N`$ Term and a Model for Nucleons
The standard model for strong interaction is QCD from which relevant non-perturbative information about the vacuum and hadron structure interested in this study is not easy to extract. Although lattice simulation could provide a first principle investigation of the problem, it is still not satisfactory. Full QCD is also too complicated for the present purpose. Simplified models constructed based on the slightly (explicit) broken chiral symmetry of the QCD Lagrangian in the fermionic sectors are frequently used, which are regarded as corresponding to the effective Lagrangian of QCD when the gluonic degrees of freedom are functionally integrated out. This effective Lagrangian is further simplified by assuming that only 4-point contact 4-point quark–quark interaction terms are dominant. The well known ’t Hooft interaction due to instanton is one of them. This kind of models are quite successful in describing a large collection of the hadronic phenomena (see, e.g., Refs. and the references therein). It indicates that these models contain elements of the truth about the physics of strong interaction.
In the present section, attempts are made in the following to justify the picture that nucleons are made up of loosely bound constituent quarks in a chiral bag. The mass and wave function renormalization of these dressed quarks are generated by the spontaneous breaking of the chiral symmetry, which is governed by an 4–point quark–quark interaction. Albeit the long range confinement effects are not considered in these models, they are expected to be reasonable ones for the ground state chiral properties of hadrons, in which quarks are not separated far apart. The underlying QCD plays very little direct role as far as the chiral properties are concerned once these quantities are renormalized using relevant data from observations. Other quantities of interest here that detailed QCD dynamics may play more important role like the confinement, the wave function renormalization of the quark fields, the bag constant and the energy density of the possible virtual color superconducting phase are, instead of predicted by the models, fitted to experimental data. If possible, a comparison with available lattice QCD data is going to be carried out to test the consistency of the model assumptions.
### 2.1 The strong interaction models
The 4–fermion interaction models for the light quark system involving up and down quarks have the generic form
$``$ $`=`$ $`\overline{\psi }\left(i\text{/}m_0\right)\psi +_{int}`$ (1)
with the 4–fermion interaction terms classified into two categories in the quark–antiquark channel
$`_{int}`$ $`=`$ $`\text{ }\text{}\}\begin{array}{c}color\\ singlet\end{array}+\text{ }\text{}\}\begin{array}{c}color\\ octet\end{array}+Fierzterm`$ (6)
$`=`$ $`_{int}^{(0)}+_{int}^{(8)}.`$ (7)
$`_{int}^{(0)}`$ generates quark–antiquark scattering in color singlet channel and $`_{int}^{(8)}`$ generates quark–antiquark scattering in color octet channel. For $`_{int}^{(0)}`$, the well known two flavor chiral symmetric Nambu Jona–Lasinio (NJL) interaction can be chosen, namely
$`_{int}^{(0)}`$ $`=`$ $`G_0\left[(\overline{\psi }\psi )^2+(\overline{\psi }i\gamma ^5𝝉\psi )^2\right].`$ (8)
The color octet $`^{(8)}`$ is written in the quark–quark (antiquark–antiquark) channel form for our purposes, namely
$`_{int}^{(8)}`$ $`=`$ $`\text{}\}\begin{array}{c}color\\ triplet\end{array}+\text{}\}\begin{array}{c}color\\ sextet\end{array}+(q\overline{q})`$ (13)
$`=`$ $`_{int}^{(3)}+_{int}^{(6)}.`$ (14)
The color sextet term is repulsive in the one gluon exchange case. Due to the non-existence of colored baryon containing three quarks in nature, it is assumed to be generally true. So we restrict ourselves to the attractive color triplet two quark interaction terms. The attractive color triplet quark bilinear terms can be classified according to their transformation properties under Lorentz and chiral $`SU(2)_L\times SU(2)_R`$ groups . In general, if only terms without derivative in fermion fields are considered, $`_{int}^{(3)}`$ has the following form
$`_{int}^{(3)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{r}{}}G_r{\displaystyle \underset{ab}{}}C_{ab}^r(\overline{\psi }\mathrm{\Gamma }_a^r\stackrel{~}{\overline{\psi }})(\stackrel{~}{\psi }\mathrm{\Gamma }_b^r\psi ),`$ (15)
with $`\mathrm{\Gamma }_a^r`$, $`\mathrm{\Gamma }_b^r`$ matrices in Dirac, flavor and color spaces generating representation “$`r`$”. The tilded fermion field operators $`\stackrel{~}{\psi }`$ and $`\stackrel{~}{\overline{\psi }}`$ in the 4-component representation are defined as
$`\stackrel{~}{\psi }`$ $`=`$ $`\psi ^T(i\tau _2)C^1,`$ (16)
$`\stackrel{~}{\overline{\psi }}`$ $`=`$ $`Ci\tau _2\overline{\psi }^T.`$ (17)
In the 8-component representation for the fermion field in the local theory, these two identities are implemented as a pair of constraints.
Operator $`\stackrel{~}{\psi }\mathrm{\Gamma }_b^r\psi `$ belongs to an irreducible representation “r” of chiral, Lorentz and color groups and $`\overline{\psi }\mathrm{\Gamma }_a^r\stackrel{~}{\overline{\psi }}`$ belongs to the conjugate representation. Coefficients $`C_{ab}^r`$ render the summation $`_{ab}\mathrm{}`$ invariant under Lorentz, chiral $`SU(2)_L\times SU(2)_R`$ and color $`SU(3)_c`$ groups. $`\{G_r\}`$ is a set of independent 4–fermion coupling constants characterizing the color triplet quark–quark interactions.
Since the QCD Lagrangian at low energy or long distance is presently unknown, the coupling constants $`G_r`$ are not known. Albeit instanton induced ’t Hooft interaction is frequently used in literature, which can reduce the independent coupling constant to only one, it will not be assumed here from the start. Rather, by studying the virtual phases of the strong interaction vacuum states using experimental observables, some of the important $`G_r`$, which determines the neighborhood of the strong interaction vacuum state, can at least be inferred. The question of whether or not instanton induced interaction dominates is left to be determined by observations.
For that purpose, two models that have different virtual superconducting phases for the vacuum state are studied. The half bosonized version of them with the quark field represented in an 8-component real representation are:
#### 2.1.1 A model for scalar color superconductivity
$`_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}\overline{\mathrm{\Psi }}\left[i\text{/}\sigma i𝝅𝝉\gamma ^5O_3\gamma ^5𝒜_c\chi ^cO_{(+)}\gamma ^5𝒜^c\overline{\chi }_cO_{()}\right]\mathrm{\Psi }`$ (18)
$`{\displaystyle \frac{1}{4G_0}}(\sigma ^2+𝝅^2)+{\displaystyle \frac{1}{2G_{3s}}}\overline{\chi }_c\chi ^c,`$
where $`\sigma `$, $`𝝅`$, $`\overline{\chi }_c`$ and $`\chi ^c`$ are auxiliary fields, $`(\chi ^c)^{}=\overline{\chi }_c`$, $`G_0`$ and $`G_{3s}`$ are coupling constants of the model. The matrices $`O_i`$ (i=1,2,3) are Pauli matrices acting on the upper and lower 4-components of $`\mathrm{\Psi }`$ and $`𝒜_c`$ (c=R,B,G) are a set of three $`2\times 2`$ antisymmetric matrices acting on the color indices of $`\mathrm{\Psi }`$.
This model Lagrangian has a color superconducting phase when the coupling constant $`G_{3s}`$ is sufficiently large. When the phase in which the normal chiral symmetry breaking is the true ground state of the vacuum, the corresponding color superconducting phase is a metastable virtual phase in which the order parameter $`\sigma `$ for the normal chiral symmetry breaking phase vanishes .
#### 2.1.2 A model for vector color superconductivity
$`_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\overline{\mathrm{\Psi }}[i\text{/}\sigma i𝝅𝝉\gamma ^5O_3+O_{(+)}(\varphi _\mu ^c\gamma ^\mu \gamma ^5𝒜_c+𝜹_\mu ^c𝝉\gamma ^\mu 𝒜_c)`$ (19)
$`O_{()}(\overline{\varphi }_{\mu c}\gamma ^\mu \gamma ^5𝒜^c+\overline{𝜹}_{\mu c}𝝉\gamma ^\mu 𝒜^c)]\mathrm{\Psi }{\displaystyle \frac{1}{4G_0}}(\sigma ^2+𝝅^2)`$
$`{\displaystyle \frac{1}{2G_{3v}}}(\overline{\varphi }_{\mu c}\varphi ^{\mu c}+\overline{𝜹}_{\mu c}𝜹^{\mu c}),`$
where $`\overline{\varphi }_{\mu c}`$, $`\varphi _\mu ^c`$ with $`(\varphi _\mu ^{})_c=\overline{\varphi }_{\mu c}`$, $`\overline{𝜹}_{\mu c}`$, $`𝜹_\mu ^c`$ with $`(𝜹_\mu ^{})_c=\overline{𝜹}_{\mu c}`$ are auxiliary fields introduced.
This model Lagrangian has a vector color superconducting phase induced by vector diquark condensation when the coupling constant $`G_{3v}`$ is sufficiently large. The color superconducting phase also breaks the chiral symmetry. When the normal chiral symmetry breaking phase is the true ground state of the vacuum, the corresponding color superconducting phase is a metastable virtual phase in which the order parameter $`\sigma `$ for the normal chiral symmetry breaking phase vanishes .
#### 2.1.3 The normal chiral symmetry breaking phase
The gap equation for the particles in the normal chiral symmetry breaking phase for both of these two models can be derived from a stability equation for the quark mass in the Hartree–Fock approximation. The result is
$$\frac{\sigma ^2}{\mathrm{\Lambda }^2}\text{ln}\left(1+\frac{\mathrm{\Lambda }^2}{\sigma ^2}\right)=\left(1\frac{\pi }{12\alpha _0}\right)$$
(20)
with
$$\alpha _0=\frac{\stackrel{~}{G}_0\mathrm{\Lambda }^2}{4\pi },$$
(21)
where $`\stackrel{~}{G}_0`$ is a linear combination of all the 4–point coupling constants of theory .
Since the vacuum $`\sigma _{vac}350`$ MeV, Eq. 20 implies $`\alpha _00.38`$ when $`\mathrm{\Lambda }`$ takes the value 0.9 GeV. The vacuum expectation value $`<0|\overline{\psi }\psi |0>`$ that can be extracted from Gell-Mann, Oakes and Renner (GOR) relation
$$f_\pi ^2m_\pi ^2=m_0<0|\overline{\psi }\psi |0>$$
(22)
is related to $`\sigma _{vac}`$ by
$$\sigma _{vac}=2\stackrel{~}{G}_0<0|\overline{\psi }\psi |0>,$$
(23)
in the mean field or Hartree–Fock approximation. Here $`m_0=(m_u+m_d)/2`$, $`m_u`$ and $`m_d`$ is the mass of the current up and down quark respectively.
For a realistic description of the physics of the strong interaction, the Hartree–Fock approximation may not be sufficient. One of the most important corrections is to sum over all the self energy diagrams for the auxiliary field $`\sigma ^{}=\sigma \sigma _{vac}`$, which is shown in Fig. 1.
The result is a modification of the tree level propagator for $`\sigma ^{}`$, which is
$$D_\sigma (p)=2iG_0,$$
(24)
to
$$D_\sigma (p)=\frac{2iG_0}{1J_n(p)}$$
(25)
with $`J_nJ_n(0)>0`$. It causes the anti-screening of the scalar charge of a dressed quark.
The gap equation is however robust against this kind of corrections since the same $`1/(1J_n)`$ factor enters both of the tadpole terms and the Hartree part of the loop terms of the stability equation . So the gap equation in the Hartree approximation with the dressed $`\sigma ^{}`$ propagator is of the following form
$$i\delta \mathrm{\Sigma }=\frac{1}{1J_n}\left(\text{}+\text{}\right)=0,$$
(26)
where the larger black dot represents the tadpole due to a shift of the auxiliary field $`\sigma `$. The resummation over one-loop self-energy corrections to the $`\sigma ^{}`$ field propagator does not modify of the gap equation derived by using tree level propagator for the $`\sigma ^{}`$ field. The Fock terms for the stability equation can leads to a modification against the naive one loop equations due to the fact that the dressed propagator for $`\sigma ^{}`$ inside the loop depends on momentum and the Fock terms due to exchange of other auxiliary fields <sup>1</sup><sup>1</sup>1The physical Goldstone pion contributions should be treated using a derivative coupling to the quarks to implement the chiral symmetry rather than using the pseudo-scalar coupling of the tree level $`𝝅`$ field . This could leads to a significant reduction of the pion loop effects, e.g., at finite temperature . does not renormalizes in the same way as $`\sigma ^{}`$. This is especially true for the auxiliary fields corresponding to the $`^{(8)}`$ term. The difference should be ignored in the following discussion based on the large $`N_c`$ argumentation which suppresses the Folk terms but should be studied in more detailed researches.
The Hartree–Fock relationship given in Eq. 23, on the other hand, has to be modified to
$$\sigma _{vac}=\frac{2\stackrel{~}{G}_0}{1J_n}<0|\overline{\psi }\psi |0>.$$
(27)
This is because the effective 4-fermion interaction strength $`\stackrel{~}{G}_0^{eff}=\stackrel{~}{G}_0/(1J_n)`$ at the level of approximation considered. Since the resummation over the one-loop self-energy terms for the $`\sigma ^{}`$ field does not shift the pole position of the quark propagator, its effects on the quarks are represented by a wave function renormalization with $`Z_\psi =1J_n`$, leading to a reduction of the vacuum expectation value of the $`\overline{\psi }\psi `$ operator against its mean field value.
### 2.2 The physical vacuum phenomenology and the global $`\mathrm{\Sigma }_N`$ term
#### 2.2.1 The quasiparticle picture for the chiral condensate
The physics of the global observables are dominated by the quasiparticles of the true ground state of the vacuum . Its physical picture, which could be useful for the discussions in the following subsection, can be derived from the gap equation for the dynamical mass of the quarks. The generic form of the one loop gap equation is of the following form
$$\sigma _{vac}=i8n_fn_c\stackrel{~}{G}_0\frac{d^4p}{(2\pi )^4}\frac{\sigma _{vac}}{p^2\sigma _{vac}^2+iϵ},$$
(28)
where $`n_f=2`$ and $`n_c=3`$ are the number of flavor and color respectively. Here, the current mass $`m_0`$ of the quarks is assumed to be zero for simplicity. If treated properly, the solution of this gap equation is beyond the Hartree–Fock approximation .
In order to project out the contributions of the quasiparticles of the true ground state of the vacuum, the 4–momentum integration can be carried out by doing the $`p^0`$ integration first . The result is
$$\sigma _{vac}=4n_fn_c\stackrel{~}{G}_0^{\mathrm{\Lambda }_3}\frac{d^3p}{(2\pi )^3}\frac{\sigma _{vac}}{E_p}$$
(29)
with $`E_p=\sqrt{p^2+\sigma ^2}`$ the energy of the quasiparticle and $`\mathrm{\Lambda }_3`$ the formal cut off in the 3-momentum, which is not going to be used for the numerical evaluations. From Eq. 27, the scalar density of the vacuum state in the one loop approximation is
$$<0|\overline{\psi }\psi |0>=2n_fn_c^{\mathrm{\Lambda }_3}\frac{d^3p}{(2\pi )^3}\left[(1J_n)\frac{\sigma _{vac}}{E_p}\right]=2n_fn_c^{\mathrm{\Lambda }_3}\frac{d^3p}{(2\pi )^3}(1J_n)\overline{u}(p)u(p)$$
(30)
with $`u(p)`$ a solution of the Dirac equation for a fermion with mass $`\sigma `$ and energy $`E_p`$. It can be seen that the non-perturbative<sup>2</sup><sup>2</sup>2 The perturbative wave function renormalization around the $`\sigma _{vac}`$ is absent since the radiative corrections to the self-energy of the quarks is self-consistently adjusted to zero in the auxiliary field approach adopted here. factor $`Z_\psi =1J_n`$ serves as a wave function renormalization factor for the fermion field $`\psi `$.
Eq. 29 implies that the scalar charge density of the vacuum state is obtained by summing over the scalar charge
$$\underset{q_\mu 0}{lim}j_s^{()}(p+q,p)=\overline{u}(p)u(p)$$
(31)
of contributing negative energy fermions in unit space volume renormalized by factor $`1J_n`$.
The physical value for $`<0|\overline{\psi }\psi |0>`$ can be computed from model independent GOR relation, which is
$$<0|\overline{\psi }\psi |0>=0.024\text{GeV}^3$$
(32)
when $`m_0=(m_u+m_d)/2=7`$ MeV is assumed.
The value of $`\sigma _{vac}`$ is generally taken to be around $`0.35`$ GeV, which corresponds to the valence quark mass. Assuming $`\mathrm{\Lambda }=0.9`$ GeV, then Eq. 32 requires
$$J_n=0.21$$
(33)
which is a positive number as expected. This number is consistent with the lattice result $`Z_\psi =0.70.8`$ of Ref. .
#### 2.2.2 Chiral symmetry and pion decay constant
One of the model independent consequences of the chiral symmetry of the QCD Lagrangian is derived from the chiral Ward–Takahashi identity
$$q^\mu 𝚪_\mu ^5=\frac{1}{2}m_0D_\pi \gamma ^5𝝉\frac{1}{4}\{\mathrm{\Sigma },\gamma ^5𝝉\}$$
(34)
with $`D_\pi `$ a scalar function and $`𝚪_\mu ^5`$ the proper axial-vector vertex for quarks that is dominated by the one pion pole at low $`q^2`$ and $`\mathrm{\Sigma }`$ their self-energy. $`𝝉`$ is the set of Pauli matrices in the isospin space. It leads to the GOR relation as the first order correction to the chiral symmetric results for the vacuum case, and, in certain approximate scheme like the one loop one, an evaluation of the decay constant for the chiral Goldstone boson , namely,
$$f_\pi ^2=\frac{n_fn_c\sigma _{vac}^2}{8\pi ^2}\left(\mathrm{ln}\frac{\sigma _{vac}^2+\mathrm{\Lambda }^2}{\sigma _{vac}^2}\frac{\mathrm{\Lambda }^2}{\sigma _{vac}^2+\mathrm{\Lambda }^2}\right).$$
(35)
It is one of the relationships used to fix the model parameters in the one loop approximation (see, e.g., ).
The picture developed here goes beyond the one loop approximation by formally summing over the self-energy terms for the auxiliary $`\sigma ^{}`$ fields, which results in a wave function renormalization factor $`1J_n`$ for the quasi-particles. Before the spontaneous chiral symmetry breaking taking place, the chiral Ward–Takahashi identity requires that there should also be a $`1/(1J_n)`$ renormalization factor for the axial-vector vertex as it is shown in the first graph of Fig. 2. After the chiral symmetry is spontaneously broken down, a massless Goldstone pion pole appears in the axial-vector vertex which is shown in the second graph of Fig. 2. It is reasonable to assume that the same renormalization for the pion–quark coupling vertex. The self-energy $`\mathrm{\Sigma }`$ that enters Eq. 34 is related to the quark mass $`\sigma _{vac}`$ as
$`\mathrm{\Sigma }`$ $``$ $`{\displaystyle \frac{\sigma _{vac}}{1J_n}}.`$ (36)
The various renormalization factors are displayed in Fig. 2. for the axial-vector vertex function with the pion pole separated out.
The renormalized one loop result for the pion decay constant is then found to be
$$f_\pi ^2=\frac{n_fn_c\sigma _{vac}^2}{8\pi ^2}\left(1J_n\right)\left(\mathrm{ln}\frac{\sigma _{vac}^2+\mathrm{\Lambda }^2}{\sigma _{vac}^2}\frac{\mathrm{\Lambda }^2}{\sigma _{vac}^2+\mathrm{\Lambda }^2}\right).$$
(37)
The values $`\mathrm{\Lambda }=0.9`$ GeV, $`f_\pi =93`$ MeV, $`\sigma _{vac}350`$ MeV and $`J_n=0.21`$ adopted above are consistent with this equation and thus with the chiral Ward–Takahashi identity 34.
Thus, the approximation approach taken here, which improves the canonical quasi-particle picture by introducing a non-perturbative wave function renormalization factor $`1J_n`$, is quit successful in the vacuum sector.
#### 2.2.3 Model independent information about the scalar charge (density) inside a nucleon
In the chiral quark model adopted here, the global nucleon $`\mathrm{\Sigma }_N`$ term, which can be extracted from the low energy pion–nucleon Compton scattering, can be evaluated in the following way
$$\mathrm{\Sigma }_N(0)=<N|:\overline{\psi }\psi :|N>=m_0\left[3(1J_n)j_s^{(+)}+\mathrm{\Delta }v\delta <0|\overline{\psi }\psi |0>\right].$$
(38)
The scalar charge for the valence quarks $`j_s^{(+)}1`$, $`\mathrm{\Delta }v`$ is the volume taken by a nucleon and $`\delta <0|\overline{\psi }\psi |0>`$ is the change of the scalar charge density of the quarks on the negative energy states inside of the nucleon.
The global $`\delta <0|\overline{\psi }\psi |0>`$ term interested here is allowed to differ from zero. It is the sum of all the renormalized quark scalar charge density of all the contributing quarks orbits with a distortion of $`\sigma _{vac}`$ inside of a nucleon minus those without a distortion of $`\sigma _{vac}`$. A non-vanishing value for this quantity thus implies a bag picture for a nucleon with the scalar field $`\sigma `$ acting as a medium which is distorted to form the non-topological “bag”.
Assuming that $`m_0=7`$ MeV, and that $`\mathrm{\Sigma }_N(0)5575`$ MeV after extracting from the Cheng–Dashen point to the $`t=0`$ point, the vacuum polarization effects are
$$\mathrm{\Delta }v\delta <0|\overline{\psi }\psi |0>=5.58.4$$
(39)
which is relatively large. As it is discussed in the following, it leads to a non-topological bag picture for a nucleon.
Before continuing, it should be pointed out that some of the studies of the nucleon $`\mathrm{\Sigma }_N`$ problem in the constituent quark model (see, e.g. Ref. and many others) are based on the following equation
$$\mathrm{\Sigma }_N(0)=<N|:\overline{\psi }\psi :|N>=m_0\frac{3j_s^{(+)}}{1J_n}.$$
(40)
It can accommodate the old value $`\mathrm{\Sigma }_N(0)45`$ MeV . Such an equation is obtained by observing that $`\mathrm{\Sigma }_N(0)`$ is a sum of the static scalar charge of quarks in an additive quark model. The effect of interaction can be obtained from the scalar vertex function by summing over ladder diagrams of the form given by Figs. 1 and 3, which gives us the factor $`1/(1J_n)`$ in the above equation. The wave function renormalization of the two amputated quark legs in such an approach is not taken into account at the same level of approximation as the vertex function.
The inclusion of the effects of the wave function renormalization is required from unitarity point of view. The interaction of the of the self-energy graphs in Fig. 3 generates poles in the s-channel, which correspond to the physical mesons , these new poles should take some of the strength of the quasi-quark due to unitarity constraint. Without $`1J_n`$ term, the so called “double or over counting” of the degrees of freedom is going to happen.
Therefore Eq. 40 is not consistent since it is well known from the study of the gauge invariance in perturbative quantum electrodynamics that it is important to consider both the vertex function renormalization and the wave function renormalization in a given order in order to maintain the Ward–Takahashi identity. To sharpen the problem, it should be noticed that Eq. 30 should be
$$<0|\overline{\psi }\psi |0>=2n_fn_c^{\mathrm{\Lambda }_3}\frac{d^3p}{(2\pi )^3}\frac{\overline{u}(p)u(p)}{1J_n}$$
(41)
in such an approximation. It would lead to a significant inconsistency between Eq. 32 and the phenomenological values $`\sigma _{vac}=0.35`$ GeV, $`\alpha _0=0.38`$, $`\mathrm{\Lambda }=0.9`$ GeV, the value of the nucleon $`\mathrm{\Sigma }_N(0)`$ value and the requirement that $`J_n=0.20.3`$ from lattice simulation of QCD .
Under Eq. 38, it is inevitable to include a non-vanishing value for $`\delta <0|\overline{\psi }\psi |0>`$, as given by Eq. 39. It implies a non-topological bag picture for the nucleon.
The radius of the bag can be determined in the local theory by requiring that a nucleon is stable in a nucleus, which is revealed in experimental observations. In the aforementioned local theory, the energy of the lowest energy level in the bag must be at
$$E=ϵ_0\mu =0.$$
(42)
This is explained in the next section. Here $`\mu `$ is the time component of the statistical gauge field that is related to the fermion number density through
$$\rho =\frac{2}{\pi ^2}\mu ^3$$
(43)
when the value of the statistical blocking parameter $`\epsilon _{vac}`$ is much smaller than $`\mu `$ which is non-vanishing inside of a nucleon. Assuming that $`\mu `$ is constant inside the bag and zero outside, its value inside of the bag is related to the bag radius
$$\frac{2}{\pi ^2}\mathrm{\Delta }v\mu ^3=\frac{8}{3\pi }(R\mu )^3=3$$
(44)
since the fermion number of a nucleon is 3. So
$$\mu =\frac{3}{2R}\left(\frac{\pi }{3}\right)^{1/3}.$$
(45)
The energy $`ϵ_0`$ can be estimated following the MIT bag model calculation
$$ϵ_0=\frac{c_0}{R}+\sigma _{in}$$
(46)
with $`c_0=2.04`$, $`\sigma _{in}`$ the $`\sigma `$ field inside of the bag, which is assumed zero in the authentic MIT bag model. Here, however,
$$\sigma _{in}=\frac{2\stackrel{~}{G}_0}{1J_n}\left(<0|\overline{\psi }\psi |0>+\delta <0|\overline{\psi }\psi |0>\right).$$
(47)
Therefore $`ϵ_0\mu =0`$ implies
$$\left[c_0\frac{3}{2}\left(\frac{\pi }{3}\right)^{1/3}\right]\frac{1J_n}{R}\frac{6\alpha _0}{\mathrm{\Lambda }^2}\frac{D}{R^3}=\frac{8\pi \alpha _0}{\mathrm{\Lambda }^2}<0|\overline{\psi }\psi |0>,$$
(48)
where $`D=\mathrm{\Delta }v\delta <0|\overline{\psi }\psi |0>=5.58.4`$ (see Eq. 39). Substituting the values $`\alpha _0=0.38`$ and $`\mathrm{\Lambda }=0.9`$ GeV, we have
$$R=0.670.78\text{fm}$$
(49)
for the core part of a nucleon. This is a quite reasonable range of numbers.
To get a rough picture for the distribution of $`\sigma `$ field around a nucleon, let us estimate $`\sigma _{in}`$ inside of a nucleon using the value obtained above
$$\sigma _{in}=153131\text{MeV}$$
(50)
Therefore we have a bag type of solution for a nucleon in this chiral model with its basic character determined by the measured global $`\mathrm{\Sigma }_N`$ term value. The value of the $`\sigma `$ field outside of the bag is of order 350 MeV and the value for $`\sigma `$ inside of the bag is of order $`153131`$ MeV.
The distribution of $`\sigma `$ around a nucleon is schematically plotted in Fig. 5 in which the three black dots represent the energy level for the three valence quarks.
## 3 The Stability of a Nucleon
The stability of a nucleon inside a nucleus is an important enough problem to the considered, since, as far as it is known now, that the properties of a nucleus can be described well in terms of the nucleon degrees of freedom. On the other hand, the current theory tends to imply a change of the properties of a nucleon inside a nucleon. For example, the bag constant of a nucleon should be reduced inside of a heavy nucleus, which would give rise to an increase of the size of the nucleon. An overlapping configuration is also favored due the reduction in the surface area between different phases, which is not usually considered. There is little room for nucleons inside a heavy nucleus to increase their size before touch each other. There does not seems to be a reasonable scheme or mechanism to reconcile these two contradictory pictures at the present.
The stability problem of a nucleon inside of a nucleus is discussed qualitatively in Ref. under the local theory for the relativistic finite density problems. The essential mechanism is derived from the statistical blocking effects in the normal chiral symmetry breaking phase of the strong interaction vacuum state that is revealed in the new theoretical frame-work.
Before getting into the quantitative discussions, an important qualitative property of the statistical blocking parameter $`\epsilon ^\mu `$ should be mentioned. It is related to the possibility of a distortion of the $`\epsilon ^\mu `$ around the nucleon, which seems to happen at a first look since the statistical gauge field $`\mu ^\alpha `$ has distinct different values inside and outside of a nucleon. Such a distortion does not happen for $`\epsilon ^\mu `$. Because the statistical gauge field $`\mu ^\alpha `$ is coupled to a conserved current. A change of $`\mu ^\alpha `$ in a space-time region would cause a change of the baryon number density in the same region. Such a change can not be balanced locally since baryon number can not be created, it has to be transported inside or outside of the region from other regions far apart which cause corresponding changes of the baryon number in these other regions. In the end, the statistical gauge field $`\mu ^\alpha `$ correlates with each other at long distances<sup>3</sup><sup>3</sup>3Such a long distance correlation is destroyed by the statistical blocking effects in the normal chiral symmetry breaking phase or the $`\alpha `$-phase of the strong interaction vacuum .. The statistical blocking parameter does not couple to a conserved number; it can be seen as effectively couples to a current that is the difference of the fermion’s baryon number and that of the anti-fermion’s baryon number. This number can be created or destroyed locally (in space-time) in an extremely efficient way since a generation or annihilation of one fermion–anti-fermion pair inside of the vacuum would change two unit of the charge of this later effective current. This means that there is little correlation between $`\epsilon ^\mu `$ at different space-time points. Therefore, $`\epsilon ^\mu `$ can be viewed as corresponding to a very massive excitation of the system, which will not response to the presence of a nucleon once its vacuum expectation value is established.
### 3.1 The vacuum statistical blocking parameter
The absolute minima of the effective potential in the normal chiral symmetry breaking phase are located at non-vanishing values of $`\epsilon `$ and vanishing $`\mu ^\alpha `$. The effective potential is of the following form
$$V_{\text{eff}}(\sigma ,\epsilon )=i2n_fn_c_𝒞\frac{d^4p}{(2\pi )^4}\mathrm{ln}\left(1\frac{\sigma ^2}{p^2}\right)+\frac{1}{4\stackrel{~}{G}_0}\sigma ^2+\frac{n_fn_c}{2\pi ^2}\epsilon ^4,$$
(51)
where the dependency on the statistical gauge field $`\mu ^\alpha `$ is suppressed and the integration contour $`𝒞`$ lies beneath the real $`p^0`$ axis for $`\mathrm{}<p^0<\epsilon `$ and $`0<p^0<\epsilon `$ and above the real $`p^0`$ axis when $`\epsilon <p^0<0`$ and $`\epsilon <p^0<\mathrm{}`$ . The minima for $`V_{\text{eff}}(\sigma ,\epsilon )`$ are not located at zero $`\epsilon `$ for any non-vanishing $`\sigma `$.
For the case interested here, $`\sigma =\sigma _{vac}=350`$ MeV. The value of the corresponding statistical blocking parameter $`\epsilon _{vac}`$ is found to be
$$\epsilon _{vac}40\text{MeV}$$
(52)
after minimizing of the effective potential.
### 3.2 The energy of the valence quarks
In the 8–component “real” theory for fermions, there are two branches of single quark orbits with energy given by
$$E_\pm =\pm \left(ϵ_0\mu \right)$$
(53)
that are relevant to the problem. Here $`ϵ_0`$ is given by Eq. 46 and $`\mu `$ is given by 45. For each specific value of $`R`$ value, only one of these two states corresponds to particle excitation and the hole of the other state correspond to mirror antiparticles . The later excitation is absent in the conventional 4-component theory for fermions. The other states with
$$E_\pm ^{}=\pm \left(ϵ_0+\mu \right)$$
(54)
lie in the continue spectra. They are not directly relevant to the discussion here.
The allowed state in the presence of the statistical blocking is shown in Fig. 6 for $`\mathrm{\Sigma }_N(0)=62`$ MeV. The spectrum of the particle excitation is draw with a solid line there. There are two discontinuities for the particle energy level for the valence quarks, with a gap of order of $`80`$ MeV. The rest of the curves draw in dotted lines are for (mirror) anti-particles.
It is likely that the crossing of $`E_+`$ and $`E_{}`$ at degenerate point of about $`R0.72`$ fm is separated by a gap in more quantitative studies. Such a small complication is not going to be discussed further here. The important feature here is the presence of a gap for the lowest energy of the quark orbits for bag radius $`R0.70`$ and $`R0.74`$ fm. These two gaps provide a stabilizing force for a nucleon.
### 3.3 The “quantization” of the size of a nucleon
The mass of a nucleon consists of two parts. The first one is the energy of the valence quarks discussed above. The second is the energy needed to establish the soliton configuration in which the value of $`\sigma `$ is different from $`\sigma _{vac}`$ and there is a finite statistical gauge field $`\mu ^0`$. Combining these contributions, the “classical value” of the mass is given by
$$M_N=3E_0(R)+\frac{4\pi }{3}BR^3+\frac{3}{2\pi ^2}\mathrm{\Delta }v\mu ^4\frac{1.5}{R},$$
(55)
where $`R`$ is treated as classical quantity. Here $`E_0(R)`$ is the allowed energy for valence quarks, $`B`$ is the bag constant similar to the MIT bag model and the third term is the corresponding term in the local theory . The last term is a term from the free fermion orbits corresponding to the valence quarks that has to be subtracted according to the prescriptions of the local theory . There should be a surface term that is proportional to $`R^2`$ in principle. I shall ignore such a term in the following discussions.
The value of $`B`$ can in principle be obtained from $`V_{\text{eff}}`$ too. But it will not be fixed this way here since the model for evaluating $`V_{\text{eff}}`$ does not include many important elements such as the gluon degrees of freedom, which may be important especially when $`\sigma \sigma _{vac}`$. It will be treated as free parameters here. If the value of $`\mathrm{\Sigma }_N(0)`$ is chosen to be
$$\mathrm{\Sigma }_N(0)=62\text{MeV}$$
(56)
within the experimental range, then, it is found that $`M_N1`$ GeV if
$$B=(219\text{MeV})^4=2.3\times 10^3\text{GeV}^4.$$
(57)
The change of $`M_N`$ with $`R`$ is shown in Fig. 7. It can be seen that for the value of $`\mathrm{\Sigma }_N(0)`$ chosen, the radius of a nucleon is around $`0.670.78`$ fm. Albeit there are two stable positions for $`M_N`$ given in Eq. 55, a nucleon has a unique radius due to the fact that so far $`R`$ is treated as a classical variable. Since $`R`$ is a collective coordinate for a finite system of confined quarks (including the sea quarks), it has to be quantized, despite its fluctuation is expected to be much reduced <sup>4</sup><sup>4</sup>4If a collective coordinate, like $`R`$ considered here, connects to the common motion of $`N`$ particles, the magnitude of its quantum spreading is reduced by a factor of $`1/\sqrt{N}`$ compared to the spreading of each of the single particles. , it is non-vanishing due to the finiteness of the system. If such a quantization of $`R`$ is considered, the tunneling between these two stable position leads to a unique stable value for $`R`$, which is in the middle of the trapping well.
Also shown in Fig. 7 are $`M_N`$ for different values of $`B`$.
It can be seen that the location of the gap for $`M_N`$ is not changed as $`B`$ is varied despite the significant change of its value. This is an important property since one of the medium effects inside of a nucleus is represented by the reduction of $`B`$ in the medium, which causes the increase of the size of the nucleon in the conventional picture. The statistical blocking effects in the local finite density theory prevent such a change from happening.
In the authentic MIT model, the dependence of $`M_N`$ on the bag radius $`R`$ is much flatter, which means that a nucleon is much softer than the one studied here. For example, the size of the nucleon changes rapidly from 0.75 fm to 1.1 fm as the bag constant $`B^{1/4}`$ is reduced from $`219`$ MeV to $`157`$ MeV, which is already too large.
Due to the presence of a gap of order $`240`$ MeV for $`M_N`$, which is three times of the gap for a single valence quark, a nucleon should appear to be quite rigid against any form of external forces to change its radial size.
## 4 The Local $`\sigma _N`$ Term and the Virtual Component of a Nucleon
After establishing a possible stable chiral model for a nucleon in a semi-quantitative way base on the phenomenological information about a nucleon, we now turn to the local properties of the nucleon concerning its possible color superconducting companion that lives in the possible virtual color superconducting phase of the strong interaction vacuum state. Such a possibility seems to be supported by various empirical evidences . The results of these works, however, depend only on general model independent properties of the virtual color superconducting phase, like the spontaneous partial breaking of the electromagnetic local gauge symmetry and the spontaneous breaking of chiral symmetry . Detailed structure of this possible superconducting component of a nucleon is not yet discussed.
A few properties of the color superconducting component of a nucleon that is likely to be true can nevertheless be listed. First, the radius of such a component should be comparable to the normal component discussed above. Second, there are also three valence quarks, which is the quasi-particles of the color superconducting phase, in such a component. This is because a nucleon is color neutral. Third, the energy density of the observable color superconducting phase(s) should not be much larger than the true ground state of the vacuum. And fourth, the normal and the virtual component of a nucleon should appear in pair in low energy processes. Keeping this picture in mind, we shall turn to the question of the static scalar charge of a nucleon.
The nucleon $`\mathrm{\Sigma }_N`$ problem is different from the leading order relations because both side of the isoscalar part of the corresponding chiral Ward–Takahashi identity for the nucleon $`\mathrm{\Sigma }_N`$ problem are sensitive to the low energy chiral dynamics. The value for $`\mathrm{\Sigma }_N`$ on the Cheng–Dashen point extracted from pion–nucleon scattering data depends only on the pion nucleon coupling and, at low energy interested here, it is a global observable. The $`\sigma _N`$ obtained from baryon spectrum contains the information of the possible metastable color superconducting phase if it indeed exists. This is because it is a local operator in space-time defined on an equal time hypersurface within the spatial region occupied by the nucleon. Therefore, according to the local theory developed for such kind of situations, it contains the dark component that is not seen in the global low energy pion–nucleon scattering observables. It is therefore interesting to see to what extent one can learn, from the difference between the nucleon $`\mathrm{\Sigma }_N`$ term extracted from the pion–nucleon scattering data and the corresponding term extracted from the baryon mass spectra, about the structure of the companion virtual superconducting component of a nucleon that lives in the possible metastable color superconducting phase of the strong interaction vacuum state and also about the difference between the energy density of the true chiral symmetry breaking phase and the virtual color superconducting phase of the strong interaction vacuum.
### 4.1 The static scalar charge of a nucleon as a local observable
The scalar charge of the nucleon that is extracted from the baryon spectrum is the following expectation value of the normal ordered scalar operator
$$\sigma _N=m_0d^3x<N|:\overline{\psi }\psi :(𝒙,t=0)|N>.$$
(58)
The normal ordering is relative the vacuum state in which the chiral symmetry is spontaneously broken and it is measured at the equal-time hypersurface of the rest frame of the nucleon.
Although the detailed information about the right hand side (r.h.s) of Eq. 58 can not be easily extracted given the QCD dynamics, some general properties of it can nevertheless be given. First, the matrix element
$$<N|:\overline{\psi }\psi :(𝒙,t)|N>$$
is non-vanishing only around the nucleon in its rest frame due to the normal ordering. Therefore, the presence of a nucleon at the time $`t=0`$ as specified in Eq. 58 can be viewed as a local perturbation (measurement) of the vacuum state with spatial dimension of order of the size of a nucleon and with a infinite resolution in time. According to the local theory developed in Refs. , the quantity $`\sigma _N`$, which is a measure of certain component of the energy density inside a nucleon, contains a dark component.
The dark component can be classified into two kinds. The first kind of the dark component is from the deviation of the local observable from the one that contains only the contribution of the quasi–particles of the vacuum state. They are always present for local observables. For an semi-quantitative study given here, such a component is relatively hard to evaluate and is expected to be small. It is only included formally here. The second kind, which is perhaps more interesting, is due to the possibility that there is a close by metastable phase for the vacuum state. In this later case, the quasi-particles from the metastable phase also contribute to the local observables with a suppression factor $``$ determined by the resolution of the observation $`\mathrm{\Delta }\omega `$ and the difference $`\mathrm{\Delta }ϵ`$ between the energy density of the true phase and the virtual phase(s) of the vacuum state.
The suppression factor can be determined according to the prescription given in Refs. . First, the quantum correlation length for the observable interested is determined by the scale $`1/M_A1`$ $`\text{GeV}^10.2`$ fm, it is smaller than the size of the nucleon but is apparently larger than the inverse of the temporal resolution, which is zero. Therefore we have
$$=e^{\mathrm{\Delta }ϵ\mathrm{\Delta }v/M_A}$$
(59)
for all one particle irreducible graphs constructed from the quasiparticles of the virtual phase . Here $`\mathrm{\Delta }v`$ is of the order of the volume taken by a single nucleon in space. Suppose that there is a virtual color superconducting phase, then Eq. 58 can be further specified as
$$\sigma _N=\sigma _n+\sigma _s$$
(60)
where $`\sigma _n`$ is the contributions from the quasi-particles (together with their dark component of the first kind discussed above) of the ground state of the vacuum in which the chiral symmetry is spontaneously broken and $`\sigma _s`$ is the contributions from the quasi-particles of the metastable virtual phase, which is assumed to be color superconducting here, including all the dark component of the first kind.
### 4.2 The local $`\sigma _N`$ term
For the local observables, the two type of dark components discussed in section 4.1 has to be included. The dark components of the first kind are included formally by doing the loop 4-momentum integration in the Euclidean space which is cut off covariantly. The dark component of the second kind, namely, the one coming from the quasiparticles of the possible metastable phases of the vacuum, can be evaluated using standard field theoretical method.
The local scalar charge Eq. 58, can be written as
$$\sigma _N=\mathrm{\Sigma }_N+\delta \sigma _N^{(val)}+\delta \sigma _N^{(sea)}$$
(61)
where $`\mathrm{\Sigma }_N`$ is the global value discussed above, $`\delta \sigma _N^{(val)}`$ is the contribution of the quasiparticles of the metastable phase to the valence quarks and $`\delta \sigma _N^{(sea)}`$ is the contribution of the quasiparticles of the metastable phase to the sea quarks.
The type of the corrections to the scalar charge of a quark are be diagrammatically expressed in Fig. 8. Since the valence quarks of the color superconducting component of a nucleon does not contain a scalar charge in the chiral limit. They do not contribute to the quantities studied here. Therefore only the Feynman graphs containing external quark lines living in the normal phase of the vacuum state are draw in Fig. 8.
After the summation of the series of the ladder diagrams and including the wave function renormalization for the quark legs, the correction to the scalar charge of a quark is found to be
$$\delta j_s=J_s\overline{u}(p)u(p)$$
(62)
with the one loop diagram contribution denoted by $`J_s`$. The value of $`J_s`$ depends on the nature of the color superconducting phase.
#### 4.2.1 The possible virtual scalar color superconducting phase
The propagators for the quarks in the scalar color superconducting phase can be found after fixing the complex phase and the direction in the color space of $`\chi ^c`$. If $`\chi ^c`$ is chosen to be along, say, the “red” direction and to be real, then
$$S_{F1}(p)=\frac{i}{\text{/}p\gamma ^5𝒜\chi O_1}$$
(63)
for “blue” and “green” quarks and
$$S_{F2}(p)=\frac{i}{\text{/}p}$$
(64)
for “red” quark.
The bubble diagram in Fig. 8 is found to be
$`J_s`$ $`=`$ $`i\stackrel{~}{G}_0\text{Tr}{\displaystyle \frac{d^4p}{(2\pi )^4}\left[S_{F1}(p)S_{F1}(p)+S_{F2}(p)S_{F2}(p)\right]}`$ (65)
$`=`$ $`16\stackrel{~}{G}_0{\displaystyle \frac{d^4p_E}{(2\pi )^4}\left(\frac{2}{p_E^2+\chi ^2}+\frac{1}{p_E^2}\right)}`$
with $`p_E^\mu `$ the Euclidean 4-momentum. The result is
$$J_s=\frac{12\alpha _0}{\pi }\left[1\frac{2}{3}\frac{\chi ^2}{\mathrm{\Lambda }^2}\mathrm{ln}\left(1+\frac{\mathrm{\Lambda }^2}{\chi ^2}\right)\right]=\frac{1\frac{2}{3}\frac{\chi ^2}{\mathrm{\Lambda }^2}\mathrm{ln}\left(1+\frac{\mathrm{\Lambda }^2}{\chi ^2}\right)}{1\frac{\sigma ^2}{\mathrm{\Lambda }^2}\mathrm{ln}\left(1+\frac{\mathrm{\Lambda }^2}{\sigma ^2}\right)}.$$
(66)
Here, use has been made of the gap equation for $`\sigma `$, namely Eq. 20. Since $`\sigma 350`$ MeV and assuming $`\chi 200`$ MeV, then $`J_s=1.3\times e^{\mathrm{\Delta }ϵ\mathrm{\Delta }v/M_A}`$.
The contribution of the quasiparticles of the virtual phase to the value of the local $`\sigma _N(0)`$ can be obtained from Eq. 38 through a replacement of the global wave function renormalization factor $`1J_n`$ by local one, namely $`1J_n1J_nJ_s`$. Since the wave function renormalization factor for every contributing quark orbits, including the sea quarks, is modified in such a way, we can write
$$\sigma _N(0)=\mathrm{\Sigma }_N(0)\frac{J_s}{1J_n}\mathrm{\Sigma }_N(0).$$
(67)
Therefore we have
$$J_s=\frac{(1J_n)\mathrm{\Delta }\mathrm{\Sigma }_N}{\mathrm{\Sigma }_N(0)}.$$
(68)
For the value of $`\chi `$ taken, we have
$$\mathrm{\Delta }ϵ=\frac{3}{4\pi }\frac{M_A}{R^3}\mathrm{ln}\frac{1.3\times \mathrm{\Sigma }_N(0)}{(1J_n)\mathrm{\Delta }\mathrm{\Sigma }_N(0)}.$$
(69)
Using the estimates $`\mathrm{\Sigma }_N(2\mu ^2)\mathrm{\Sigma }_N(0)=1015`$ MeV from the chiral perturbation theory and dispersion relation calculation , $`\mathrm{\Sigma }_N(0)=5575`$ MeV. If no strangeness is assumed for a nucleon and $`M_A=\mathrm{\Lambda }=0.9`$ GeV, then
$$\mathrm{\Delta }ϵ=789409\text{MeV/fm}\text{3}.$$
(70)
This is a reasonable number.
#### 4.2.2 The possible virtual vector color superconducting phase
The propagator represented by the dashed lines in Fig. 8 for the vector color superconducting phase are also found after fixing the color and complex phase for the order parameter $`\varphi _\mu ^c`$. One of the choices is that for “blue” and “green” quarks
$`S_{F1}(p)`$ $`=`$ $`\left(1iO_2{\displaystyle \frac{\text{/}p}{p^2}}\text{/}\varphi \gamma ^5𝒜\right)F(p),`$ (71)
$`F(p)`$ $`=`$ $`i{\displaystyle \frac{(p^2\varphi ^2)\text{/}p2p\varphi \text{/}\varphi }{(p^2\varphi ^2)^24(p\varphi )^2}}`$ (72)
and
$$S_{F2}(p)=\frac{i}{\text{/}p}$$
(73)
for “red” quark.
The bubble diagram in Fig. 8 is found to be
$`J_s`$ $`=`$ $`i\stackrel{~}{G}_0\text{Tr}{\displaystyle \frac{d^4p}{(2\pi )^4}\left[S_{F1}(p)S_{F1}(p)+S_{F2}(p)S_{F2}(p)\right]}`$ (74)
$`=`$ $`16\stackrel{~}{G}_0{\displaystyle \frac{d^4p_E}{(2\pi )^4}\left(\frac{2(p_E^2+\varphi ^2)}{(p_E^2+\varphi ^2)4(p_E\varphi )^2}+\frac{1}{p_E^2}\right)}`$
with $`p_E^\mu `$ the Euclidean 4-momentum. The result is quite simple
$$J_s=\frac{12\alpha _0}{\pi }\left(1\frac{1}{3}\frac{\varphi ^2}{\mathrm{\Lambda }^2}\right)=\frac{1\frac{1}{3}\frac{\varphi ^2}{\mathrm{\Lambda }^2}}{1\frac{\sigma ^2}{\mathrm{\Lambda }^2}\mathrm{ln}\left(1+\frac{\mathrm{\Lambda }^2}{\sigma ^2}\right)}.$$
(75)
For typical values of $`\sigma _{vac}=350`$ MeV and $`\sqrt{\varphi ^2}=200`$ MeV, $`J_s=1.4\times e^{\mathrm{\Delta }ϵ\mathrm{\Delta }v/M_A}`$ and Eq. 69 is replaced by
$$\mathrm{\Delta }ϵ=\frac{3}{4\pi }\frac{M_A}{R^3}\mathrm{ln}\frac{1.4\times \mathrm{\Sigma }_N(0)}{(1J_n)\mathrm{\Delta }\mathrm{\Sigma }_N(0)},$$
(76)
which is not much different from Eq. 69. The values for $`\mathrm{\Delta }ϵ`$ is of the same order as Eq. 70, namely
$$\mathrm{\Delta }ϵ=853450\text{MeV/fm}\text{3}.$$
(77)
This is a reasonable number.
This range of $`\mathrm{\Delta }ϵ`$ is quite stable against change of the order parameters for the color superconducting phase. For example, even in the limit $`\sqrt{\varphi ^2}0`$,
$$\mathrm{\Delta }ϵ=865457\text{MeV/fm}\text{3},$$
(78)
which changes very little. The same is true for the scalar color superconducting phase.
## 5 Discussions
### 5.1 The inputs and outputs
The input parameters are listed in Table 1.
The quantitative outputs of our investigation based on the local theory are given in Table 2. It can be seen that using 6 input parameters/equations, more than double physical quantities/equations can be obtained/satisfied. The gain is obvious. In addition, a new mechanism for the stability of a nucleon inside a nucleus or nuclear matter is found and semi-quantitative pictures for the nucleon and the vacuum state of the strong interaction is also established.
The presence of at least one metastable virtual color superconducting phase for the strong interaction vacuum is assumed instead of derived in this paper. Unlike in other work, like Refs. , the existence of the metastable color superconducting phase for the strong interaction vacuum state is a sufficient condition for the solution of the nucleon $`\mathrm{\Sigma }_N`$ term puzzle. It is not a necessary one despite the fact that such an assumption is rather natural for the solution of the puzzle.
### 5.2 The nucleon as a non-topological soliton
It is demonstrated that when the quark wave function renormalization is taken into account, the experimentally observed global nucleon $`\mathrm{\Sigma }_N`$ term implies, in general, that a modification of the quasi–particle picture for the hadrons seems to be inevitable. This kind of picture for a nucleon is somewhat different from the authentic constituent quark model but are in qualitative agreement with the bag type of models for a nucleon introduced on more phenomenological basis where the nature of the scalar $`\sigma `$ field, which is claimed to be related to gluon contributions, is actually not clear. Here the scalar field is the order parameter for the chiral symmetry, which contains contributions from gluon as well as quarks. The sizes of the non-topological soliton are found to be around $`0.7`$ fm.
Given the fact that $`\sigma _{vac}`$, the bag constant $`B`$ and the renormalization factor $`J_n`$ are related to the complicated underlying QCD dynamics, they are not computed using the 4-fermion interaction models adopted here, which serve the purpose of providing the general structure of the solution to the problems and of discussing the chiral symmetry issues. Rather, these parameters are determined by fitting empirical data. This is done despite the fact that model computations may not leads to a set of values significantly deviate from the fitted ones, such a success may not be regarded as an indication that QCD formulated in the local theory can be replaced by the models but only that these models are effective ones.
### 5.3 The solution of the two puzzles
After the basic picture of a nucleon is established, the role played by the statistical blocking effects of the chiral symmetry breaking phase in maintaining the stability of a nucleon inside a large nucleus and in a nuclear matter is discussed. It is shown that two of the essential elements of the local theory, namely, the statistical blocking effects and the inequivalent 8-component “real” representation for the fermion field are both needed for the mechanism. The bag constant $`B`$ is determined in the frame-work of the local theory by producing the right mass for the nucleon.
The discrepancy between the nucleon $`\mathrm{\Sigma }_N`$ term from low energy pion–nucleon scattering data and the one extracted from baryon spectra is resolved by realizing that the former one is a global observable while the later one is a local one, which according to the local theory, contains a dark component. The dark component is further attributed to a possible virtual color superconducting phase that seems to manifest in the high energy electromagnetic interactions involving a nucleon and some other ones . The difference in energy density of the metastable phase and the stable phase of the strong interaction vacuum state is estimated.
Thus, the new features of the local theory and the assumption that there is a virtual color superconducting phase for the strong interaction vacuum state provides a set of conditions for a simultaneous solution of the two puzzles.
### 5.4 The diquark model for a nucleon?
The diquark models for a nucleon has a long history. It has certain advantages on the phenomenological ground. There are theoretical difficulties for these models. It is criticized in a similar, albeit more detailed in certain aspects, approach to the quark–quark interaction, which does not consider the required wave function renormalization . We belief that their argument is sound at least for this specific problem.
But under the current scenario, the result of Ref. is not the end of the diquark model for a nucleon since its argument applies only to the normal component of the nucleon, which lives in the true ground state of the vacuum in which the chiral symmetry is spontaneously broken. In the case of the possible virtual color superconducting component for a nucleon, which lives in the possible virtual color superconducting phase of the vacuum, an entirely different set of channels, like in the $`\chi `$ or $`\varphi ^\mu `$, $`𝜹^\mu `$ etc. fields must be iterated. It is not known at the present whether or not a quark–quark clustering for the “constituent” quarks in the color superconducting component of a nucleon is favored or not. It is one of the questions to be understood in the future.
### 5.5 Speculations
Given these figures, it is tempting to speculate that perhaps the Ropper resonance at $`1.44`$ GeV is an excited $`N^{}`$ resonance produced in such a way that its normal component is striped away by some means so that it is dominated by the color superconducting companion of a nucleon? It would be very interesting to understand the structure of the color superconducting virtual companion of the nucleon and more importantly the virtual color superconducting phase of the strong interaction vacuum state if they are indeed there.
## 6 Summary
A mechanism for the stability of a nucleon and a resolution of the nucleon $`\mathrm{\Sigma }_N`$ term problem are studied based on an implementation of the chiral Ward–Takahashi identity of QCD at a level that goes beyond the mean field approximation by formally summing over ladder diagrams for $`\sigma ^{}`$ field. The local theory for the problem is adopted as the theoretical frame-work for a coherent discussion and solution of these two seemingly unrelated puzzles.
In the future, self-consistent numerical study of the non-topological soliton model proposed here is an interesting topic to be studied . Theoretical and experimental means should be searched and/or developed to study the structure of the possible superconducting component of the nucleon in more details. It is interesting not only to the nucleon structure itself, but also because it carries information about the possible virtual color superconducting phase of the vacuum state, which is one of the most important objects to be understood in the context of contemporary discoveries in cosmology.
## Acknowledgement
This work is supported by the National Natural Science Foundation of China under contract 19875009. Part of the present work was carried out while the author was visiting the CSSM of the University of Adelaide, Adelaide, Australia.
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# Detection of gravitational waves with quantum encryption technology
## Abstract
We propose a new technique for detecting gravitational waves using Quantum Entangled STate (QUEST) technology. Gravitational waves reduce the non-locality of correlated quanta controlled by Bell’s inequalities, distorting quantum encryption key statistics away from a pure white noise. Gravitational waves therefore act as shadow eavesdroppers. The resulting colour distortions can, at least in principle, be separated from noise and can differentiate both deterministic and stochastic sources.
http://sun1.sms.port.ac.uk/cosmos/users/bruce/mousetrap/
PACS numbers: 04.30.-w, 03.67Dd, 95.55Ym, 03.65Bz
Quantum cryptography provides a stunning application of Einstein-Podolsky-Rosen (EPR) correlations and Bell’s inequalities . Not only does it promise perfectly secure key distribution but we argue that it may also allow the detection of the shadowy traces of gravitational waves whose existence is the most important outstanding test of Einstein’s General Relativity and the subject of massive current and next-generation experiments .
A general quantum encryptographic scheme consists of key generation using entangled quantum states, by two parties (Alice and Bob) interested in communicating securely. An attack may be made by an eavesdropper (Eve) who secretly attempts to determine the key as the pairs of entangled EPR quanta travel to Alice and Bob. By performing a sequence of measurements on these entangled pairs of photons, Alice and Bob determine the key they will use to encrypt their message. The vital advantage quantum mechanics provides lies in the impossibility that Eve can intercept the secret key, made up of individual quanta, without giving away her presence to Alice and Bob, since such interception unavoidably alters the entanglement of the EPR pairs, as measured by violations of Bell’s inequalities.
Variants of the standard BB84 protocol based on the transmission of single pairs of EPR photons have been used recently in practical quantum key distribution over optic fiber networks up to $`48`$ km in length . Similar experiments have illustrated the feasibility of quantum encryption in practical situations and the field is now sufficiently mature to be a tool in fundamental research beyond the foundations of quantum mechanics.
In this letter we propose a simple gedanken experiment to detect the effects of gravitational waves through the distortions they cause in the statistics of the quantum keys determined by Alice and Bob. The use of quantum encryption technology may be implemented in at least two ways: one based on randomly swapping polarisers, the second based on laser interferometry.
Consider the Ekert protocol in which Alice and Bob are equipped with randomly swapping polarizers. Entangled pairs of photons are emitted in the singlet state
$$|\mathrm{\Psi }=\frac{1}{\sqrt{2}}[|H_A|V_B+|V_A|H_B],$$
(1)
where the photon $`A`$ is sent to Alice, and the photon $`B`$ to Bob. $`H`$ and $`V`$ denote the horizontal and vertical polarizations, prepared by a laser coupled to a parametric down-conversion device. The arrival time of the photons at the polarizers is synchronized with their random swapping.
If a polarizer happens to be correctly oriented, the incident photon is detected, and a “1” is recorded. Otherwise a “0”. Repetition generates two equal length binary strings $`𝐀`$ and $`𝐁`$, corresponding to the measurements of Alice’s and Bob’s detectors (see Fig. (1)).
Alice and Bob then publicly announce the orientations of their polarizers corresponding to each element in $`𝐀`$ and $`𝐁`$. They then eliminate the elements of $`𝐀`$ and $`𝐁`$ corresponding to non-coincident orientations of the two polarizers. The string entries of the remaining subsets of $`𝐀`$ and $`𝐁`$ form the two quantum keys, $`K_A`$ and $`K_B`$. In the absence of gravitational waves and noise the two keys coincide, $`K_A=K_B`$, since the photon pairs were perfectly entangled.
With the keys determined Alice proceeds with the transmission of a message encrypted with her key, using e.g. logical AND or XOR, to Bob, who decodes it with his key, $`K_B`$ However, this is not of interest to us. Instead, cross-correlation of the keys $`K_A`$ and $`K_B`$ allows, in principle, the detection of gravitational waves.
This detection proceeds thanks to a fundamental property of quantum cryptography: the keys derived from an ideal experiment are Markovian, pure white noise random strings of “0”s and “1”s. The presence of a gravitational wave colours the cross-correlation statistics so they are no longer white (see Figures 1 and 2).
A gravitational wave introduces a discoloration by changing the arrival time of the photons at Alice and Bob, by altering the detectors’ local time and the path length travelled by the photon. This implies that, in the key strings $`K_A`$ and $`K_B`$, the probability of a “1” (a detection) is no longer equal to the probability of a “0” (a non-detection). In addition the two strings will no longer coincide element by element: $`K_AK_B`$. In order to analyse this effect one may construct the cross-correlation matrix between $`K_A`$ and $`K_B`$ (see Fig. 2) and search for off-diagonal power.
Alternatively it is convenient to consider the string $`𝐊K_AK_B`$ formed using an appropriate operator $``$, such as logical AND. We then define the accumulated fluctuation $`\xi (t)`$ as the absolute value of the difference, for a given temporal length $`t`$, of the number of non-detections, $`N_{[0]}`$, and detections, $`N_{[1]}`$, in $`𝐊`$, viz. $`\xi (t)|N_{[1]}N_{[0]}|`$.
$`\xi (t)`$ then obeys the stochastic differential equation
$$\frac{d\xi }{dt}=\frac{\mathrm{\Gamma }_{ph}}{2}(w(t)+h(t)),$$
(2)
where $`\mathrm{\Gamma }_{ph}`$ is the photons pair rate, $`w(t)`$ is a stochastic process which describes the intrinsic noise of the system and $`h(t)`$ is the strain produced by the gravitational wave. The key feature of this equation is that, since the intrinsic noise $`w(t)`$ is due only to the effective randomness of the polarizer, it has the ideal statistical properties
$$w(t)=0,w(t)w(t^{})=D\delta (tt^{}).$$
(3)
In the idealized case where complex and experiment-specific noise sources (such as thermal and seismic fluctuations) are neglected, the intrinsic fluctuations in the time series $`\xi `$ extracted from the polarizers are described by a frequency-independent random process characterized by the noise spectral density $`D`$ \- the noise-induced mean square fluctuations per unit frequency.
Since QUEST detectors use single photon pairs there is no shot noise and if $`h(t)=0`$, the accumulated fluctuation is a random process with zero mean and a linearly increasing variance
$$(\text{Var}\xi )^2=\xi (t)\xi (t)=\frac{\mathrm{\Gamma }_{ph}^2}{4}Dt.$$
(4)
A gravitational wave affects the statistical properties of the accumulated fluctuations. For the sake of completeness, we distinguish between (i) a deterministic gravitational wave (such as that originating from a binary system of compact objects such as black holes or neutron stars) and (ii) a stochastic background of gravitational waves (arising from inflation or other cosmological sources) . If a deterministic gravitational wave is present the accumulated fluctuation is no longer described by a zero-mean random variable. Rather its mean value is
$$\xi (t)=\frac{\mathrm{\Gamma }_{ph}}{2}_0^t𝑑\tau h(\tau ).$$
(5)
If instead a stochastic background of gravitational waves is present, the accumulated fluctuations are still described by a zero mean variable, but the gravitational wave affects the variance of $`\xi `$, which for a stationary gravitational wave background is
$$(\text{Var}\xi )^2=\frac{\mathrm{\Gamma }_{ph}^2}{4}Dt\left[1+\frac{1}{D}_0^t𝑑\tau H(\tau )\right]$$
(6)
where $`H(\tau )=h(t)h(0)`$.
In order to assess the sensitivity of our gedanken experiment to a gravitational wave we need to estimate the noise background induced by the effective randomness of the polarizer. Consider the data set collected by one observer divided into sub-sets of $`N`$ points. Since each point corresponds to a photon, $`N=(\mathrm{\Gamma }_{ph}/2)\tau _N`$, where $`\tau _N`$ is the temporal length of the sub-set. For each data sub-set, the background noise due to the polarizer is then
$`S_N=D/N210^{43}\text{Hz}^1\times \left({\displaystyle \frac{10^8\text{s}^1}{\mathrm{\Gamma }_{ph}}}\right)\left({\displaystyle \frac{\mathrm{\Theta }_{sw}}{10^{10}}}\right)^2\left({\displaystyle \frac{\tau _{coh}}{10^{12}\text{s}}}\right)\left({\displaystyle \frac{10^3\text{s}}{\tau _N}}\right),`$ (7)
where $`\mathrm{\Theta }_{sw}`$ is the percentage error in the swapping of the random polarizers and $`\tau _{coh}`$ is the photon coherence time. At a given frequency $`f`$ and for a data sub-set of about 20 minutes, the characteristic amplitude of the noise induced by the polarization swapper $`h_{rms}=(2fS_N)^{1/2}`$ is then $`6.3\times 10^{22}(f/1\text{Hz})^{1/2}`$. For comparison, the expected characteristic noise amplitude for the LIGO interferometer around $`200Hz`$ is $`2\times 10^{22}`$. The noise associated with the polarizers is therefore minimized at frequencies lower than the typical frequencies where large scale earth-based interferometers reach their best sensitivity.
It is also important to stress that, unlike large scale interferometers, this device effectively operates when the gravitational wavelength $`\lambda _{GW}`$ is less than the distance $`d_{AB}`$ between the two receivers $`A,B`$. In particular, the probability of “detection” and “non-detection” become equal when $`\lambda _{GW}2d_{AB}`$ and therefore there is a low frequency cut-off around $`fc/2d_{AB}`$. The precise response of such an experiment will depend on seismic and thermal noise which in turn depend on the exact experimental set-up. Since we are interested in general issues we do not address this important issue here.
An alternative to the swapping polarizer set-up is to exploit a quantum cryptographic protocol based on the continuous detection of photons in interferometers, relying on energy-time correlations rather than on polarization correlations . We shall discuss in detail this implementation and noise-related issues in future work.
To clarify our proposal consider the effects of gravitational waves on the famous Bell inequalities describing quantum non-locality. Both the polarizer and interferometer implementations can be unified in the following formalism. Let $`R_{ij}(\delta _A,\delta _B),i,j=0,1`$ be the number of time-correlated events detected by Alice (A) and Bob (B) as a function of instrument parameters $`\delta _i`$, which will be polarization orientations or phase shifts in the case of an interferometer setup. The normalized correlation coefficient of the measurements made by the detectors $`A`$ and $`B`$ is then
$`E(\delta _A,\delta _B)={\displaystyle \frac{R_{00}(\delta _A,\delta _B)R_{01}(\delta _A,\delta _B)R_{10}(\delta _A,\delta _B)+R_{11}(\delta _A,\delta _B)}{R_{00}(\delta _A,\delta _B)+R_{01}(\delta _A,\delta _B)+R_{10}(\delta _A,\delta _B)+R_{11}(\delta _A,\delta _B)}}.`$
Following , one may then define the composite operator
$`S|E(\delta _A^{},\delta _B^{\prime \prime })+E(\delta _A^{},\delta _B^{\prime \prime })+E(\delta ^{\prime \prime }{}_{A}{}^{},\delta _B^{})E(\delta ^{\prime \prime }{}_{A}{}^{},\delta ^{\prime \prime }{}_{B}{}^{})|`$ (8)
where $`\delta _i^{}`$ and $`\delta _i^{\prime \prime }`$ represent specific values of the parameters $`\delta _i`$. The quantity $`S`$ allows one to test the degree of violation of Bell’s inequalities; in particular, the value $`S=2\sqrt{2}`$ is achieved for maximal entanglement of the states. Gravitational waves reduce $`S`$, as will a general eavesdropper.
We have outlined how quantum encryption technology may be exploited to yield a potentially sensitive detector of weak cosmic gravitational waves. These QUantum Entangled STate (QUEST) detectors are complementary to current and planned interferometric detectors such as LIGO and LISA. While QUEST detectors are not affected by shot noise, detailed understanding of all other relevant noise sources is lacking and depends on exact details of the detector set-up. Hence whether QUEST detectors can reach the sensitivity of the interferometric detectors by exploiting the non-locality fundamental to quantum mechanics is still unknown. What is certain is that gravitational waves will act as shadow eavesdroppers, reducing the degree of entanglement between quantum states controlled by Bell’s inequalities, which is precisely how they would be detected.
###### Acknowledgements.
We thank Marco Bruni, David Kaiser, Philippos Papodopoulos and Alberto Vecchio for insightful comments on the manuscript.
Captions
Figure 1
Schematic illustration of the proposed experiment: a source of EPR photons repeatedly sends a single pair of polarization-entangled photons to Alice and Bob who are equipped with randomly swapping polarizers. A gravitational wave, by affecting the path length and local proper time of the observer (here Bob) will reduce the probability of detection of the photon causing distortions of the detection statistics used to build the quantum keys (see Figure 2).
Figure 2
The averaged cross-correlation matrix of sample $`50`$-element keys $`K_A`$ and $`K_B`$. Inset: The idealized white-noise case (without gravitational waves). The diagonal dominates in the large key length limit where the cross-correlation is simply $`\delta _{ij}`$. The main figure schematically shows the effects of a deterministic gravitational wave which induces off-diagonal power representing non-white correlations.
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# On the Reidemeister torsion of rational homology spheres
## Introduction
In the paper V. Turaev has proved a certain identity involving the Reidemeister torsion of a rational homology sphere. In this very short note we will suitably interpret this identity as a second order finite difference equation satisfied by the torsion which will allow us to prove a general structure result for the $`mod`$ reduction of the torsion. More precisely we prove that the mod $`𝐙`$ reduction of the torsion is completely determined by three data.
$``$ a certain canonical $`spin^c`$ structure,
$``$ the linking form of the rational homology sphere and
$``$ a constant $`c/`$.
As a consequence, the constant $`c`$ is a $`/`$-valued invariant of the rational homology sphere. Experimentations with lens spaces suggest this invariant is as powerful as the torsion itself.
## 1 The Reidemeister torsion
We review briefly a few basic facts about the Reidemeister torsion a rational homology $`3`$-sphere. For more details and examples we refer to .
Suppose $`M`$ is a rational homology sphere. We set $`H:=H_1(M,)`$ and use the multiplicative notation to denote the group operation on $`H`$. Denote $`Spin^c(M)`$ the $`H`$-torsor of isomorphism classes of $`spin^c`$ structure on $`M`$. We denote by $`𝔉`$ the space of functions
$$\varphi :H.$$
The group $`H`$ acts on $`𝔉_M`$ by
$$H\times 𝔉(g,\varphi )g\varphi $$
where
$$(g\varphi )(h)=\varphi (hg).$$
We denote by $`_H`$ the augmentation map
$$𝔉_M,_H\varphi :=\underset{hH}{}\varphi (h).$$
According to Reidemeister torsion is a $`H`$-equivariant map
$$\tau :Spin^c(M)𝔉_M,Spin^c(M)\sigma \tau _\sigma 𝔉_M$$
such that
$$_H\tau _\sigma =0$$
Denote by $`\mathrm{𝐥𝐤}_M`$ the linking form of $`M`$,
$$\mathrm{𝐥𝐤}_M:H\times H/.$$
V. Turaev has proved in that $`\tau _\sigma `$ satisfies the identity
$$\tau _\sigma (g_1g_2h)\tau _\sigma (g_1h)\tau _\sigma (g_2h)+\tau _\sigma (h)=\mathrm{𝐥𝐤}_M(g_1,g_2)mod$$
(1.1)
$`g_1,g_2,hH,\sigma Spin^c(M)`$.
## 2 A second order “differential equation”
The identity (1.1) admits a more suggestive interpretation. To describe it we need a few more notation.
Denote by $`𝒮`$ the space of functions $`H/`$. Each $`gH`$ defines a first order differential operator
$$\mathrm{\Delta }_g:𝒮𝒮,(\mathrm{\Delta }_gu)(h):=u(gh)u(h),u𝒮,hH.$$
If $`\mathrm{\Xi }=\mathrm{\Xi }_\sigma `$ denotes the mod $``$ reduction of $`\tau _\sigma `$ then we can rewrite (1.1) as
$$(\mathrm{\Delta }_{g_1}\mathrm{\Delta }_{g_2}\mathrm{\Xi })(h)=\mathrm{𝐥𝐤}_M(g_1,g_2)$$
(2.1)
We will prove uniqueness and existence results for this equation. We begin with the (almost) uniqueness part.
###### Lemma 2.1.
The second order linear differential equation (2.1) determines $`\mathrm{\Xi }`$ up to an “affine” function.
Proof Suppose $`\mathrm{\Xi }_1`$, $`\mathrm{\Xi }_2`$ are two solutions of the above equation. Set $`\mathrm{\Psi }:=\mathrm{\Xi }_1\mathrm{\Xi }_2`$. $`\mathrm{\Psi }`$ satisfies the equation
$$\mathrm{\Delta }_{g_1}\mathrm{\Delta }_{g_2}\mathrm{\Psi }=0.$$
Now observe that any function $`F𝒮`$ satisfying the second order equation
$$\mathrm{\Delta }_u\mathrm{\Delta }_vF=0,u,vH$$
is affine, i.e. it has the form
$$F=c+\lambda $$
where $`c/`$ is a constant and $`\lambda :H/`$ is a character. Indeed, the condition
$$\mathrm{\Delta }_u(\mathrm{\Delta }_vF)=0,u$$
implies $`\mathrm{\Delta }_vF`$ is a constant depending on $`v`$, $`c(v)`$. Thus
$$F(vh)F(h)=c(v),h.$$
The function $`G=FF(1)`$ satisfies the same differential equation
$$G(vh)G(h)=c(v)$$
and the additional condition $`G(1)=0`$. If we set $`h=1`$ in the above equation we deduce
$$G(v)=c(v).$$
Hence
$$G(vh)=G(h)+G(v),v,h$$
so that $`G`$ is a character and $`F=F(1)+G`$. Thus, the differential equation (2.1) determines $`\mathrm{\Xi }`$ up to a constant and a character. $`\mathrm{}`$
###### Lemma 2.2.
Suppose $`b:H\times H/`$ is a nonsigular, symmetric bilinear form on $`H`$. Then there exists a quadratic form $`q:H/`$ such that
$$\mathrm{\Delta }q=b$$
where
$$(\mathrm{\Delta }q)(uv):=q(uv)q(u)q(v).$$
Proof<sup>1</sup><sup>1</sup>1We are indebted to Andrew Ranicki for suggesting this approach. Let us briefly recall the terminology in this lemma. $`b`$ is nonsingular if the induced map $`GG^{\mathrm{}}`$ is an isomorphism. A quadratic map form is a function $`q:H/`$ such that
$$q(0)=0,q(u^k)=k^2q(u),uH,k$$
and $`\mathrm{\Delta }q`$ is a bilinear form.
Suppose $`b`$ is a nonsingular, symmetric , bilinear form $`H\times H/`$. Then, according to \[4, §7\], $`b`$ admits a resolution. This is a nondegenerate, symmetric, bilinear form
$$B:\mathrm{\Lambda }\times \mathrm{\Lambda }$$
on a free abelian group $`\mathrm{\Lambda }`$ such that, the induced monomorphism $`J_B:\mathrm{\Lambda }\mathrm{\Lambda }^{}`$ is a resolution of $`H`$
$$0\mathrm{\Lambda }\stackrel{J_B}{}\mathrm{\Lambda }^{}\stackrel{\pi }{}H0$$
and $`b`$ coincides with the induced bilinear form on $`\mathrm{\Lambda }^{}/(J_B\mathrm{\Lambda })`$ ($`n:=\mathrm{\#}H`$)
$$b(\pi (u),\pi (v))=\frac{1}{n^2}B(J_B^1(nu),J_B^1(nv))mod,u,v\mathrm{\Lambda }^{}.$$
Now set
$$q(\pi (u))=\frac{1}{2n^2}B(J_B^1(nu),J_B^1(nu))mod$$
It is clear that this quantity is well defined i.e.
$$\frac{1}{2n^2}B(J_B^1(nu),J_B^1(nu))=\frac{1}{2n^2}B(J_B^1(nv),J_B^1(nv))mod$$
if $`v=u+J_B\lambda ,\lambda \mathrm{\Lambda }`$. Clearly
$$\mathrm{\Delta }q=b.\mathrm{}$$
We deduce that there exists a constant $`c`$, a character $`\lambda :H/`$ and a quadratic form $`q`$ such that
$$\mathrm{\Xi }(h)=\mathrm{\Xi }_\sigma (h)=c+\lambda (h)+q(h),\mathrm{\Delta }q=\mathrm{𝐥𝐤}_M.$$
In the above discussion the choice of the $`spin^c`$ structure $`\sigma `$ is tantamount to a choice of an origin of $`H`$ which allowed us to identify the torsion of $`M`$ as a function $`H`$. Once we make such a non-canonical choice, we have to replace $`\mathrm{\Xi }`$ with the family of translates
$$\{\mathrm{\Xi }_g():=\mathrm{\Xi }(g);gH\}$$
In particular
$$\mathrm{\Xi }_g(h):=\mathrm{\Xi }(gh)=c+\lambda (gh)+q(gh)=\underset{c(g)}{\underset{}{\left(c+\lambda (g)+q(g)\right)}}+\underset{\lambda _g(h)}{\underset{}{\left(\lambda (h)+(\mathrm{\Delta }q)(g,h)\right)}}+q(h)$$
where $`\lambda _g()=\lambda ()+\mathrm{𝐥𝐤}_M(g,)`$. Since the linking from is nondegenerate we can find an unique $`g`$ such that $`\lambda _g=0`$.
We have proved the following result.
###### Proposition 2.3.
Suppose $`M`$ is a rational homology sphere. Then there exists an unique $`spin^c`$-structure $`\sigma `$ on $`M`$ so that, with respect to this choice the mod $``$ reduction of $`\tau _{M,\sigma }`$
$$\mathrm{\Xi }(h):=\tau _\sigma (h)mod$$
has the form
$$\mathrm{\Xi }(h)=c+q(h)$$
where $`c/`$ is a constant while $`q(u)`$ is the unique quadratic form such that
$$\mathrm{\Delta }q=\mathrm{𝐥𝐤}_M.$$
In particular,
$$\mathrm{\Xi }(h)=\mathrm{\Xi }(h^1)mod,$$
and the constant $`c/`$ is a topological invariant of $`M`$.
## 3 Examples
We want to show on some simple examples that the invariant $`c`$ is nontrivial.
(a) Suppose $`M=L(8,3)`$. Then its torsion is (see )
$$T_{8,\mathrm{\hspace{0.17em}3}}\frac{9}{32}x^7\frac{3}{32}x^6\frac{9}{32}x^5+\frac{5}{32}x^4+\frac{7}{32}x^3\frac{3}{32}x^2+\frac{7}{32}x+\frac{5}{32}$$
where $`x^8=1`$ is a generator of $`^8`$. Then
$$q(x^n)=\frac{3k^2n^2}{16}$$
The set of possible values $`\frac{3m^2}{16}`$ mod $``$ is
$$A:=\{0,\frac{3}{16},\frac{4}{16},\frac{5}{16}\}$$
The set possible values of $`\mathrm{\Xi }(h)`$ is
$$B:=\{\frac{9}{32},\frac{3}{32},\frac{5}{32},\frac{7}{32}\}.$$
We need to find a constant $`c/`$ such that
$$Bc=A.$$
Equivalently, we need to figure out orderings $`\{a_1,a_2,a_3,a_4\}`$ and $`\{b_1,b_2,b_3,b_4\}`$ of $`A`$ and $`B`$ such that $`b_ia_i`$ mod $``$ is a constant independent of $`i`$. A little trial and error shows that
$$\stackrel{}{A}=(0,\frac{3}{16},\frac{4}{16},\frac{5}{16}),\stackrel{}{B}=(\frac{3}{32},\frac{9}{32},\frac{5}{32},\frac{7}{32})$$
and the constant is $`c=\frac{3}{32}`$. This is the coefficient of $`x^2`$. We deduce that (modulo $``$)
$$F:=T_{8,3}(x)+\frac{3}{32}\frac{3}{16}x^70x^6\frac{3}{16}x^5+\frac{1}{4}x^4+\frac{1}{4}x^30x^2+\frac{1}{4}x+\frac{1}{4}$$
The translation of $`F`$ by $`x^2`$ is
$$x^2(T_{8,3}+\frac{3}{32})=\frac{1}{4}x^7+\frac{1}{4}x^6\frac{3}{16}x^5\frac{3}{16}x^3+\frac{1}{4}x^2+\frac{1}{4}x.$$
(b) Suppose $`M=L(7,2)`$. Then, its torsion is (see )
$$T_{7,\mathrm{\hspace{0.17em}2}}\frac{2}{7}x^6+\frac{1}{7}x^5+\frac{2}{7}x^3+\frac{1}{7}x\frac{2}{7}$$
where $`x^7=1`$ is a generator of $`_7`$. We see that in this form $`T_{7,2}`$ is symmetric, i.e. the coefficient of $`x^k`$ is equal to the coefficient of $`x^{6k}`$. The constant $`c`$ in this case must be the coefficient of the middle monomial $`x^3`$, which is $`\frac{2}{7}`$.
(c) Suppose $`M=L(7,1)`$. Then
$$T_{7,\mathrm{\hspace{0.17em}1}}\frac{2}{7}x^6+\frac{1}{7}x^5\frac{1}{7}x^4\frac{4}{7}x^3\frac{1}{7}x^2+\frac{1}{7}x+\frac{2}{7}$$
This is again a symmetric polynomial and the coefficient of the middle monomial is $`4/7`$. We see that this invariant distinguishes the lens spaces $`L(7,1)`$, $`L(7,2)`$.
(d) For $`M=L(9,2)`$ we have
$$T_{9,\mathrm{\hspace{0.17em}2}}\frac{10}{27}x^8+\frac{2}{27}x^7\frac{1}{27}x^6+\frac{8}{27}x^5+\frac{2}{27}x^4+\frac{8}{27}x^3\frac{1}{27}x^2+\frac{2}{27}x\frac{10}{27}$$
Again, this is a symmetric function, i.e the coefficient of $`x^k`$ is equal to the coefficient of $`x^{8k}`$, $`x^9=1`$. The constant is the coefficient of $`x^5`$, which is $`2/27`$. We deduce that, mod $``$, we have
$$T_{9,2}=\frac{2}{3}x^8\frac{2}{9}x^7\frac{1}{3}x^6\frac{2}{9}x^7$$
(e) Finally when $`M=L(9,7)`$ we have
$$T_{9,\mathrm{\hspace{0.17em}7}}\frac{8}{27}x^8\frac{2}{27}x^7+\frac{10}{27}x^6+\frac{1}{27}x^5\frac{2}{27}x^4+\frac{1}{27}x^3+\frac{10}{27}x^2\frac{2}{27}x\frac{8}{27}$$
the polynomial is again symmetric so that the constant $`c`$ is the coefficient of $`x^4`$ which is $`2/7`$.
It would be very interesting to know whether the invariant $`c`$ satisfies any surgery properties. This is not a trivial issue because we cannot relate the potential surgery properties of $`c`$ to the surgery properties of the torsion. In the case of torsion the surgery formula involve finite difference operators which kill the constants so $`c`$ will not appear in any of them.
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# On the Internal Structure of Relativistic Jets11footnote 1Astronomy Letters, 2000, 26, (4), 208–218, Translated from Russian by V. Astakhov
## 1 Introduction
The formation mechanism of jets is a key issue in the study of the magnetospheric structure of compact astrophysical objects. Indeed, jets are observed in most compact sources, ranging from active galactic nuclei (AGNs), quasars, and radio galaxies to accreting neutron stars, solar-mass black holes (SS 433, X-ray novae) , and young stellar objects . Moreover, jets have also been recently discovered in young radio pulsars . At the same time, in most studies devoted to the magnetohydrodynamic (MHD) model of such objects , in which the formation of jets is coupled with the attraction of longitudinal currents flowing in the magnetosphere, the attention was focused on intrinsic collimation in the sense that the effect of the external medium was assumed to be marginal. However, such a situation is possible only for a nonzero total current $`I`$ flowing within the jet , so the question of its closure in the outer parts of the magnetosphere arises. On the other hand, the longitudinal current is often constrained by the regularity condition at the fast magnetosonic surface, which by no means always leads to the sufficiently large longitudinal currents required for collimation . In other words, with the exception of the force-free case , as yet no working model of a jet in which, on the one hand, the total electric current would be zero and, on the other hand, the total magnetic flux $`\mathrm{\Psi }_0`$ in the jet would be finite, has been constructed. However, the force-free approximation, in which, by definition, the particle energy density is disregarded, does not allow the fraction of energy transferred by the outflowing plasma to be determined.
At the same time, the question of collimation cannot be solved in isolation from the external conditions (see, e.g., ). In particular, this is clear even from a popular example of the magnetosphere of a compact object with a monopole magnetic field, because for any arbitrarily weak external regular magnetic field, the monopole solution (for which the magnetic field falls off as $`r^2`$) cannot be extended to infinity. Moreover, as is well known from an example of moving cosmic bodies, such as Jupiter’s moons or artificial Earth satellites , as well as radio pulsars , the external magnetic field can serve as an effective transfer link, which occasionally determines the general energy losses of the system. For this reason, constructing a consistent magnetospheric model for compact objects immersed in an external magnetic field is, in our view, of undeniable interest, especially since, as was noted above, such a jet model was previously constructed in the force-free approximation .
Undoubtedly, the existence of an external regular magnetic field in the vicinity of compact objects is largely open to question. The regular magnetic field in our Galaxy, i.e., the field that is constant on scales comparable to the sizes of our Galaxy, is known to be
$$B_{\mathrm{ext}}10^6\mathrm{G},$$
(1)
and essentially matches the random magnetic-field component, which varies even on scales of several parsecs . However, if the collimation is assumed to be actually produced by an external magnetic field, it becomes possible to estimate the jet radius. Indeed, assuming the magnetic field in the jet to be similar to the external magnetic field (1), we obtain from the condition for the conservation of magnetic flux
$$r_\mathrm{j}R\left(\frac{B_{\mathrm{in}}}{B_{\mathrm{ext}}}\right)^{1/2},$$
(2)
where $`R`$ and $`B_{\mathrm{in}}`$ are the radius and magnetic field of the compact object, respectively. For example, for AGNs ($`B_{\mathrm{in}}10^4\mathrm{G}`$, $`R10^{13}\mathrm{cm}`$), we have
$$r_\mathrm{j}1\mathrm{pc},$$
(3)
which corresponds to the observed jet radii . One might expect such a picture to be also preserved for an external medium with pressure $`PB^2/8\pi `$; therefore, it seems of interest to consider the internal structure of a one-dimensional jet immersed in an external uniform magnetic field. However, a discussion of the more realistic case of a medium with pressure is beyond the scope of this study. Nor do we discuss the collimation itself but only consider the internal structure of observed one-dimensional jets. This issue has become particularly urgent because of the new possibilities offered by space radio interferometry, which enables the internal structure of such jets to be resolved. The effect of an external medium on the internal structure of relativistic jets in the MHD model discussed here was previously studied only by Appl and Camenzind . They considered only a special case with a constant angular velocity of the plasma, in which the solution with a zero total electric field flowing inside the jet could not be constructed. As we show below, it is for the case of an angular velocity decreasing toward the jet periphery (which, incidentally, is typical of all models with a magnetic field passing through the accretion disk) that the solution with finite magnetic flux $`\mathrm{\Psi }_0`$ and zero total current $`I(\mathrm{\Psi }_0)=0`$ can be constructed.
On the other hand, many authors obtained a universal solution with a central core for a cylindrical jet:
$$B_z=\frac{B_0}{1+\varpi ^2/\varpi _c^2},$$
(4)
where, in the relativistic case,
$$\varpi _c=\frac{c\gamma }{\mathrm{\Omega }}$$
(5)
is the size of the central core, $`\mathrm{\Omega }`$ is the angular velocity of the compact object, and $`\gamma `$ is the characteristic Lorentz factor of the outflowing plasma. As can be easily seen, such a solution results in a rapid falloff of the poloidal field $`B_z\varpi ^2`$ far from the rotation axis $`\varpi \varpi _c`$. However, this solution is in conflict with the force-free approximation, in which the poloidal magnetic field remains essentially constant . Indeed, when the energy density of the electromagnetic field exceeds appreciably the plasma energy density (and it is this case that was considered), it would be natural to assume that the internal jet structure must be similar to the force-free one.
The examples given above show that a more detailed study with allowance for all possible solutions is required even in the simplest case of a one-dimensional cylindrical jet considered in terms of ideal magnetohydrodynamics. Our study aims at a consistent investigation of this issue.
Several features that we use when studying the structure of relativistic jets typical of AGNs and radio pulsars should be immediately noted. First of all, the jet radius $`r_\mathrm{j}`$ in all real cases proves to be considerably larger than the light-cylinder radius $`R_\mathrm{L}=c/\mathrm{\Omega }`$. This implies that, when the internal structure of jets is investigated, the corresponding equations must be written in complete relativistic form. On the other hand, the gravitational forces can be disregarded in them far from the compact object. Finally, for simplicity, we consider below a cold plasma, which is justifiable because the thermal processes in the magnetospheres of radio pulsars play no crucial role. As for the jets from AGNs, this approximation is applicable here in those magnetospheric regions in which the plasma density is low. In any case, this is true for the field lines passing through the surface of a black hole.
The one-dimensional solutions describing collimated jets are obtained in Sect. 2; analytic and numerical solutions for the basic physical quantities characterizing the structure and physics of jets are given in Sect. 3. The problem is solved in straightforward statement; i.e., all jets characteristics are determined by a set of parameters in the compact source and, most importantly, by the physical conditions in the external medium. As a result, we have found the conditions under which most of the energy in actual relativistic jets must be transferred by the electromagnetic field, while a region with subsonic flow exists in the central jet regions. We also show that the solution with a central core (4) and (5) cannot be realized in an external magnetic field. Finally, some astrophysical implications of the theory developed for one-dimensional jets are discussed in Sect. 4.
## 2 Basic equations
Let us consider the structure of a one-dimensional jet where all quantities depend only on radius $`\varpi `$; in what follows, the temperature of the matter is assumed to be zero, and $`c=1`$. As in the general axisymmetric case, it is convenient to describe the magnetic-field structure in terms of magnetic-flux function $`\mathrm{\Psi }(\varpi )`$, which is related to the longitudinal magnetic field by
$$B_z(\varpi )=\frac{1}{2\pi \varpi }\frac{d\mathrm{\Psi }}{d\varpi }.$$
(6)
Accordingly, it is convenient to write the toroidal magnetic field, the electric field, and the 4-velocity vector of the matter as
$$B_\phi (\varpi )=\frac{2I}{\varpi },$$
(7)
$$𝐄=\frac{\mathrm{\Omega }_\mathrm{F}}{2\pi }\frac{d\mathrm{\Psi }}{d\varpi }𝐞_\varpi ,$$
(8)
$$𝐮=\frac{\eta }{n}𝐁+\gamma \mathrm{\Omega }_\mathrm{F}\varpi 𝐞_\phi ,$$
(9)
where $`I(\varpi _0)`$ is the total current within $`\varpi <\varpi _0`$. In the case of a cold plasma, at the cylindrical magnetic surfaces $`\mathrm{\Psi }=\mathrm{const}`$, four ”integrals of motion” can be introduced, which should be considered precisely as functions of magnetic flux $`\mathrm{\Psi }`$ in the most general statement. These are primarily $`\mathrm{\Omega }_\mathrm{F}(\mathrm{\Psi })`$ and $`\eta (\mathrm{\Psi })`$ in the definitions (8) and (9), as well as the $`z`$-component of angular momentum $`L(\mathrm{\Psi })=I/2\pi +\mu \eta \varpi u_\varphi `$ and the energy flux $`E(\mathrm{\Psi })=\mathrm{\Omega }_\mathrm{F}I/2\pi +\gamma \mu \eta `$. Here, $`\mu `$ is the relativistic specific enthalpy, which is equal to the mass of particles for a cold plasma. The specific form of the integrals of motion must be determined from boundary conditions in the compact source and from critical conditions at the singular surfaces.
As a result, the equilibrium equation for magnetic surfaces far from gravitating bodies (Grad–Shafranov’s equation) can be written as (see, e.g., )
$`{\displaystyle \frac{1}{\varpi }}{\displaystyle \frac{d}{d\varpi }}\left({\displaystyle \frac{A}{\varpi }}{\displaystyle \frac{d\mathrm{\Psi }}{d\varpi }}\right)+\mathrm{\Omega }_\mathrm{F}(\mathrm{\Psi })^2{\displaystyle \frac{d\mathrm{\Omega }_\mathrm{F}}{d\mathrm{\Psi }}}+{\displaystyle \frac{64\pi ^4}{\varpi ^2}}{\displaystyle \frac{1}{2M^2}}{\displaystyle \frac{d}{d\mathrm{\Psi }}}\left({\displaystyle \frac{G}{A}}\right){\displaystyle \frac{32\pi ^4}{M^2}}{\displaystyle \frac{d(\mu ^2\eta ^2)}{d\mathrm{\Psi }}}=0,`$ (10)
where
$$G=\varpi ^2e^2+M^2L^2M^2\varpi ^2E^2,$$
$$A=1\mathrm{\Omega }_\mathrm{F}^2\varpi ^2M^2,$$
$$e(\mathrm{\Psi })=E(\mathrm{\Psi })\mathrm{\Omega }_\mathrm{F}(\mathrm{\Psi })L(\mathrm{\Psi }).$$
Here, $`M^2=𝐮_p^2/𝐮_a^2`$ is the square of the Mach number with respect to the Alfven velocity $`u_a=B_z/(4\pi n\mu )^{1/2}`$, and the derivative $`d/d\mathrm{\Psi }`$ acts only on the integrals of motion. The remaining jet parameters are given by the well-known algebraic relations (see, e.g., ):
$`{\displaystyle \frac{I}{2\pi }}`$ $`=`$ $`{\displaystyle \frac{L\mathrm{\Omega }_\mathrm{F}\varpi ^2E}{1\mathrm{\Omega }_\mathrm{F}^2\varpi ^2M^2}},`$ (11)
$`\gamma `$ $`=`$ $`{\displaystyle \frac{1}{\mu \eta }}{\displaystyle \frac{(E\mathrm{\Omega }_\mathrm{F}L)M^2E}{1\mathrm{\Omega }_\mathrm{F}^2\varpi ^2M^2}},`$ (12)
$`u_\phi `$ $`=`$ $`{\displaystyle \frac{1}{\varpi \mu \eta }}{\displaystyle \frac{(E\mathrm{\Omega }_\mathrm{F}L)\mathrm{\Omega }_\mathrm{F}\varpi ^2LM^2}{1\mathrm{\Omega }_\mathrm{F}^2\varpi ^2M^2}}.`$ (13)
Equation (10) contains four integrals of motion; this equation has no singularity at the fast magnetosonic surface, because it depends only on coordinate $`\varpi `$. As for the Alfven surface, $`A=0`$, the problem of the boundary conditions generally requires a further study beyond the theory of ideal magnetohydrodynamics. At the same time, for the fairly large currents $`II_{\mathrm{GJ}}`$ considered here, a solution continuous across this surface can always be constructed by a small change in the integrals of motion near the Alfven surface, $`\mathrm{\Psi }\mathrm{\Psi }_\mathrm{A}`$. Consequently, equation (10) requires six boundary conditions. These boundary conditions primarily include the external uniform magnetic field
$$B_z(r_\mathrm{j})=B_{\mathrm{ext}},$$
(14)
and the regularity condition at the magnetic axis $`\varpi 0`$
$$\mathrm{\Psi }(\varpi )C\varpi ^2.$$
(15)
In addition, all four integrals $`\mathrm{\Omega }_\mathrm{F}`$, $`E`$, $`L`$, and $`\eta `$ must be specified. As for the remaining quantities characterizing the flow, such as the jet radius $`r_\mathrm{j}`$ and the outflowing plasma energy, they must be determined as a solution of the problem formulated above. Similarly, the solution of the problem must also give an answer to the question of whether the flow in the jet is supersonic.
Let us now consider the determination of the integrals of motion in more detail. It would be natural to assume that, at the jet boundary where there is no longitudinal motion of the matter, all four integrals of motion become zero
$`\mathrm{\Omega }_\mathrm{F}(\mathrm{\Psi }_0)=0,E(\mathrm{\Psi }_0)=0,L(\mathrm{\Psi }_0)=0,\eta (\mathrm{\Psi }_0)=0.`$ (16)
Here, $`\mathrm{\Psi }_0`$ is the finite total magnetic flux concentrated in the jet. This case corresponds to the absence of tangential discontinuities at the jet boundary; according to (11), the total electric current within the jet is automatically equal to zero.
We use the integrals of motion $`\mathrm{\Omega }_\mathrm{F}(\mathrm{\Psi })`$, $`L(\mathrm{\Psi })=I(\mathrm{\Psi })/2\pi `$, and $`E(\mathrm{\Psi })=\mathrm{\Omega }_\mathrm{F}(\mathrm{\Psi })L(\mathrm{\Psi })`$ derived by Beskin et al. for the force-free magnetosphere of a black hole, which satisfy the conditions (16) and, consequently, can be directly used to study the jet structure. The only but very important change here is the fact that, for a finite magnetization parameter $`\sigma `$ ,
$$\sigma =\frac{\mathrm{\Omega }_\mathrm{F}^2(0)}{8\pi ^2}\frac{\mathrm{\Psi }_0}{\mu \eta (0)},$$
(17)
which tends to infinity in the force-free approximation, the particle contribution must be added to the energy integral $`E(\mathrm{\Psi })`$, because the energy flux of the electromagnetic field near the rotation axis must inevitably vanish. As a result, we have with accuracy up to $`\sigma ^11`$
$$\mathrm{\Omega }_\mathrm{F}(\mathrm{\Psi })=\frac{2\sqrt{1\mathrm{\Psi }/\mathrm{\Psi }_0}}{1+\sqrt{1\mathrm{\Psi }/\mathrm{\Psi }_0}}\mathrm{\Omega }_\mathrm{F}(0),$$
(18)
$$I(\mathrm{\Psi })=\frac{1}{2\pi ^2}\frac{\sqrt{1\mathrm{\Psi }/\mathrm{\Psi }_0}}{1+\sqrt{1\mathrm{\Psi }/\mathrm{\Psi }_0}}\mathrm{\Omega }_\mathrm{F}(0)\mathrm{\Psi },$$
(19)
$$E(\mathrm{\Psi })=\gamma _{\mathrm{in}}\mu \eta +\mathrm{\Omega }_\mathrm{F}(\mathrm{\Psi })L(\mathrm{\Psi }).$$
(20)
Below, we assume, for simplicity, that
$$\gamma _{\mathrm{in}}=\mathrm{const}.$$
(21)
We emphasize that $`\gamma _{\mathrm{in}}`$ in expression (20) has the meaning of the injection Lorentz factor in the region of the compact object and it is not equal to the Lorentz factor of the jet particles.
Thus, we see from the formula for the energy flux $`E(\mathrm{\Psi })`$ that the contribution by the electromagnetic field becomes dominant only at $`\mathrm{\Psi }>\mathrm{\Psi }_{\mathrm{in}}`$, where
$$\mathrm{\Psi }_{\mathrm{in}}=\frac{\gamma _{\mathrm{in}}}{\sigma }\mathrm{\Psi }_0.$$
(22)
At low values of $`\mathrm{\Psi }`$, most of the energy is transferred by the relativistic particles; as directly follows from relation (20), their Lorentz factor is constant and equal to their initial value $`\gamma _{\mathrm{in}}`$. As for the integral $`\eta (\mathrm{\Psi })`$, the particle-to-magnetic flux ratio, we chose it in the form
$$\eta (\mathrm{\Psi })=\eta _0(1\mathrm{\Psi }/\mathrm{\Psi }_0),$$
(23)
which satisfies the condition (16).
We emphasize that the very possibility of using the integrals of motion obtained by analyzing the inner magnetospheric regions, is not trivial. Indeed, the flow outside the fast magnetosonic surface is completely determined by four boundary conditions at the surface of a rotating body. At the same time, a one-dimensional flow can be produced by the interaction with the external medium, which gives rise (see, e.g., ) to perturbations or shock waves propagating from ”acute angles” and other irregularities. Therefore, in regions where the conditions for the validity of ideal magnetohydrodynamics are violated, a significant redistribution of energy $`E`$ and angular momentum $`L`$ is possible (e.g. a part of them can be lost via radiation). Nevertheless, we assume here, for simplicity, that the integrals of motion $`E(\mathrm{\Psi })`$ and $`L(\mathrm{\Psi })`$, functions of flux $`\mathrm{\Psi }`$, remain exactly the same as those in the inner magnetospheric regions.
In the one-dimensional case we consider, it is convenient to reduce the second-order equation (10) to a set of two first-order equations for $`\mathrm{\Psi }(\varpi )`$ and $`M^2(\varpi )`$. Multiplying equation (10) by $`2A(d\mathrm{\Psi }/d\varpi )`$, we obtain
$`{\displaystyle \frac{d}{d\varpi }}\left[{\displaystyle \frac{A^2}{\varpi ^2}}\left({\displaystyle \frac{d\mathrm{\Psi }}{d\varpi }}\right)^2\right]+A\left({\displaystyle \frac{d\mathrm{\Psi }}{d\varpi }}\right)^2{\displaystyle \frac{d^{}\mathrm{\Omega }_\mathrm{F}^2}{d\varpi }}+{\displaystyle \frac{64\pi ^4A}{\varpi ^2M^2}}{\displaystyle \frac{d^{}}{d\varpi }}\left({\displaystyle \frac{G}{A}}\right){\displaystyle \frac{64\pi ^4A}{M^2}}{\displaystyle \frac{d}{d\varpi }}(\mu ^2\eta ^2)=0,`$ (24)
with the derivative $`d^{}/d\varpi `$ acting only on the integrals of motion. Finally, we use ”Bernoulli’s relativistic equation” $`\gamma ^2𝐮^2=1`$, which, given the definitions of the integrals of motion $`E(\mathrm{\Psi })`$ and $`L(\mathrm{\Psi })`$, can be written as
$`A^2\left({\displaystyle \frac{dy}{dx}}\right)^2={\displaystyle \frac{e^2}{\mu ^2\eta ^2}}{\displaystyle \frac{x^2(AM^2)}{M^4}}+{\displaystyle \frac{x^2E^2}{\mu ^2\eta ^2}}{\displaystyle \frac{\mathrm{\Omega }_\mathrm{F}^2(0)L^2}{\mu ^2\eta ^2}}{\displaystyle \frac{x^2A^2}{M^4}},`$ (25)
where we introduced the dimensionless variables
$`x`$ $`=`$ $`\mathrm{\Omega }_\mathrm{F}(0)\varpi =\mathrm{\Omega }\varpi ,`$ (26)
$`y`$ $`=`$ $`\sigma \mathrm{\Psi }/\mathrm{\Psi }_0.`$ (27)
As a result, substituting the right-hand part of (25) into the first term of (24) and performing differentiation, we obtain the first first-order differential equation
$`\left[{\displaystyle \frac{e^2}{\mu ^2\eta ^2}}+{\displaystyle \frac{\mathrm{\Omega }_\mathrm{F}^2}{\mathrm{\Omega }_\mathrm{F}^2(0)}}x^21\right]{\displaystyle \frac{dM^2}{dx}}={\displaystyle \frac{M^6}{x^3A}}{\displaystyle \frac{\mathrm{\Omega }_\mathrm{F}^2(0)L^2}{\mu ^2\eta ^2}}{\displaystyle \frac{xM^2}{A}}{\displaystyle \frac{\mathrm{\Omega }_\mathrm{F}^2}{\mathrm{\Omega }_\mathrm{F}^2(0)}}\left({\displaystyle \frac{e^2}{\mu ^2\eta ^2}}2A\right)`$
$`+{\displaystyle \frac{M^2}{2}}{\displaystyle \frac{dy}{dx}}\left[{\displaystyle \frac{1}{\mu ^2\eta ^2}}{\displaystyle \frac{de^2}{dy}}+{\displaystyle \frac{x^2}{\mathrm{\Omega }_\mathrm{F}^2(0)}}{\displaystyle \frac{d\mathrm{\Omega }_\mathrm{F}^2}{dy}}2\left(1{\displaystyle \frac{\mathrm{\Omega }_\mathrm{F}^2}{\mathrm{\Omega }_\mathrm{F}^2(0)}}x^2\right){\displaystyle \frac{1}{\eta }}{\displaystyle \frac{d\eta }{dy}}\right].`$ (28)
The second first-order differential equation is Bernoulli’s equation (25), which should now be considered as an equation for the derivative $`dy/dx`$. The set of equations (25) and (28) allows a general solution to be constructed for a one-dimensional jet immersed in an external magnetic field.
We emphasize one important advantage of the set of first-order equations (25) and (28) over the initial second-order equation (10). The point is that the relativistic equation (10), which is basically the force balance equation, contains the electromagnetic force
$$𝐅_{\mathrm{em}}=\rho _\mathrm{e}𝐄+𝐣\times 𝐁,$$
(29)
in which the electric and magnetic contributions virtually cancel each other far out from the rotation axis $`\varpi R_\mathrm{L}`$. Using Bernoulli’s equation (25), we can derive
$$\frac{|\rho _\mathrm{e}𝐄+𝐣\times 𝐁|}{|𝐣\times 𝐁|}\frac{1}{\gamma ^2}.$$
(30)
When analyzing (10), we therefore must retain all higher order terms $`\gamma ^2`$, while the zero-order quantities $`\rho _e𝐄`$ and $`𝐣\times 𝐁`$ in (28) are analytically removed using Bernoulli’s equation, so all terms of this equation are of the same order.
Finally, it is also important that the exact equation (28) has no singularity near the rotation axis. In other words, its solution contains no $`\delta `$-shaped current $`I\delta (\varpi )`$ flowing along the jet axis; several authors pointed out to the necessity of it .
## 3 Exact solutions and numerical results
Let us consider the basic properties of the set of equations (25) and (28). As can be easily verified, in the relativistic case under consideration, we may assume $`\gamma =u_z`$ with high accuracy. Far from the rotation axis, $`x\gamma _{\mathrm{in}}`$ ($`\varpi \gamma _{\mathrm{in}}R_\mathrm{L}`$), equation (25) can be rewritten in the limit $`M^2x^2`$ as
$$\frac{d\mathrm{\Psi }}{d\varpi }=\frac{8\pi ^2E(\mathrm{\Psi })}{\varpi \mathrm{\Omega }_\mathrm{F}^2(\mathrm{\Psi })}$$
(31)
or, equivalently,
$$B_z(\varpi )=\frac{4\pi E(\mathrm{\Psi })}{\varpi ^2\mathrm{\Omega }_\mathrm{F}^2(\mathrm{\Psi })}.$$
(32)
As we see, equation (31) does not contain $`M^2`$ at all and can therefore be integrated independently. This must be the case, because equation (31) must coincide with the asymptotics of the force-free equation, which can be derived from (25) by going to the limit $`M^20`$. Assuming now that $`B_z(r_\mathrm{j})=B_{\mathrm{ext}}`$ in (32), we obtain, in particular, for the jet radius
$$r_\mathrm{j}^2=\underset{\mathrm{\Psi }\mathrm{\Psi }_0}{lim}\frac{4\pi E(\mathrm{\Psi })}{\mathrm{\Omega }_\mathrm{F}^2(\mathrm{\Psi })B_{\mathrm{ext}}}.$$
(33)
Consequently, the jet radius is determined by the limit of the $`E(\mathrm{\Psi })/\mathrm{\Omega }_\mathrm{F}^2(\mathrm{\Psi })`$ ratio as $`\mathrm{\Psi }\mathrm{\Psi }_0`$. In particular, for $`E(\mathrm{\Psi })`$ and $`\mathrm{\Omega }_\mathrm{F}(\mathrm{\Psi })`$ given by (18)–(20), we have
$$\underset{\mathrm{\Psi }\mathrm{\Psi }_0}{lim}\frac{E(\mathrm{\Psi })}{\mathrm{\Omega }_\mathrm{F}^2(\mathrm{\Psi })}=\frac{1}{4\pi ^2}\mathrm{\Psi }_0,$$
(34)
so the limit (33) does actually exist. As a result, we obtain
$$r_\mathrm{j}=\sqrt{\frac{\mathrm{\Psi }_0}{\pi B_{\mathrm{ext}}}},$$
(35)
which essentially coincides with estimate (2). This is no surprise, because we show below that equations (25) and (28) for the integrals of motion (18), (20), and (23) have a constant magnetic field as their solution over a wide range of $`\varpi `$.
Let us now consider in more details the behavior of the solution in the inner jet region, where $`\mathrm{\Psi }\mathrm{\Psi }_0`$ and, hence, the integrals of motion can be approximately written as
$`L(\mathrm{\Psi })`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }_\mathrm{F}}{4\pi ^2}}\mathrm{\Psi },`$ (36)
$`\mathrm{\Omega }_\mathrm{F}(\mathrm{\Psi })`$ $`=`$ $`\mathrm{\Omega }=\mathrm{const},`$ (37)
$`\eta (\mathrm{\Psi })`$ $`=`$ $`\eta _0=\mathrm{const},`$ (38)
with $`E(\mathrm{\Psi })=\gamma _{\mathrm{in}}\mu \eta _0+\mathrm{\Omega }_\mathrm{F}L`$ and $`e(\mathrm{\Psi })=\gamma _{\mathrm{in}}\mu \eta _0=\mathrm{const}`$. As a result, $`\mathrm{\Omega }_\mathrm{F}L/\mu \eta _0=2y`$, and we can rewrite equations (25) and (28) as
$`(1x^2M^2)^2\left({\displaystyle \frac{dy}{dx}}\right)^2={\displaystyle \frac{\gamma _{\mathrm{in}}^2x^2}{M^4}}(1x^22M^2)+x^2(\gamma _{\mathrm{in}}+2y)^24y^2{\displaystyle \frac{x^2}{M^4}}(1x^2M^2)^2,`$ (39)
$$(\gamma _{\mathrm{in}}^2+x^21)\frac{dM^2}{dx}=2xM^2\frac{\gamma _{\mathrm{in}}^2xM^2}{(1x^2M^2)}+\frac{4y^2M^6}{x^3(1x^2M^2)}.$$
(40)
Equations (39) and (40) describing the internal jet structure can be solved analytically. It can be verified by direct substitution that we have the following asymptotics for $`x\gamma _{\mathrm{in}}`$:
$$M^2(x)=M_0^2=\mathrm{const},$$
(41)
$$y(x)=\frac{\gamma _{\mathrm{in}}}{2M_0^2}x^2,$$
(42)
which correspond to a constant magnetic field
$$B_z=B_z(0)=\frac{4\pi \gamma _{\mathrm{in}}\mu \eta _0}{M_0^2}=\frac{\gamma _{\mathrm{in}}}{\sigma M_0^2}B(R_\mathrm{L})=\mathrm{const},$$
(43)
where $`B(R_\mathrm{L})=\mathrm{\Psi }_0/R_\mathrm{L}^2`$. Here, we assume that $`\gamma _{\mathrm{in}}1`$, which is the typical for jets from AGNs and radio pulsars.
The solution of (39) and (40) for $`x\gamma _{\mathrm{in}}`$, i.e., at $`\varpi \gamma _{\mathrm{in}}R_\mathrm{L}`$, depends on the relationship between $`\gamma _{\mathrm{in}}`$ and $`M_0=M(0)`$. For example, at $`M_0^2>\gamma _{\mathrm{in}}^2`$, when, according to (43), the axial magnetic field is fairly weak, the total magnetic flux within $`\varpi <\gamma _{\mathrm{in}}R_\mathrm{L}`$
$$\mathrm{\Psi }(\gamma _{\mathrm{in}}R_\mathrm{L})\pi \gamma _{\mathrm{in}}^2R_\mathrm{L}^2B_z(0)$$
(44)
can be written as
$$\mathrm{\Psi }(\gamma _{\mathrm{in}}R_\mathrm{L})\frac{\gamma _{\mathrm{in}}^2}{M_0^2}\mathrm{\Psi }_{\mathrm{in}},$$
(45)
where the flux $`\mathrm{\Psi }_{\mathrm{in}}`$ is given by (22). We see that, if the condition $`M_0^2>\gamma _{\mathrm{in}}^2`$ is satisfied, then the total magnetic flux within $`\varpi <\gamma _{\mathrm{in}}R_\mathrm{L}`$ is lower than $`\mathrm{\Psi }_{\mathrm{in}}`$; so, outside this region, the particles also make the main contribution to $`E(\mathrm{\Psi })`$ as before, while the contribution of the electromagnetic field may be neglected. As a result, at $`\varpi \gamma _{\mathrm{in}}R_\mathrm{L}`$, the solution of (39) and (40) has a quadratic rise of $`M^2`$ and a power-law falloff of the magnetic field :
$$M^2(x)=M_0^2\frac{x^2}{\gamma _{\mathrm{in}}^2}x^2,$$
(46)
$$B_z(x)=B_z(0)\frac{\gamma _{\mathrm{in}}^2}{x^2}.$$
(47)
Consequently, the magnetic flux increases very slowly (logarithmically) with the distance from the rotation axis:
$$\mathrm{\Psi }(x)\mathrm{ln}(x/\gamma _{\mathrm{in}}).$$
(48)
Such a behavior of the magnetic field, in turn, shows that the transition flux $`\mathrm{\Psi }=\mathrm{\Psi }_{\mathrm{in}}`$ is reached exponentially far from the rotation axis, in conflict with the estimate (2) corresponding to the assumption of jet collimation.
We may thus conclude that an external constant magnetic field limits the Mach number at the rotation axis above
$$M^2(0)<M_{\mathrm{max}}^2=\gamma _{\mathrm{in}}^2.$$
(49)
Accordingly, as follows from (43), the magnetic field at the rotation axis cannot be weaker than
$$B_{\mathrm{min}}=\frac{1}{\sigma \gamma _{\mathrm{in}}}B(R_\mathrm{L}).$$
(50)
If, however, the Mach number at the rotation axis does not exceed $`\gamma _{\mathrm{in}}`$ (i.e. if $`M_0^2<\gamma _{\mathrm{in}}`$), then, as for the similar asymptotics $`\varpi \gamma _{\mathrm{in}}R_\mathrm{L}`$, the solution of (39) and (40) for $`\gamma _{\mathrm{in}}R_\mathrm{L}\varpi r_\mathrm{j}`$ gives a constant magnetic field (43), which corresponds to the solution
$$y(x)=\frac{\gamma _{\mathrm{in}}}{2M_0^2}x^2.$$
(51)
At the same time, in this case, we have only a linear increase in the square of the Mach number
$$M^2(x)=M_0^2\frac{x}{\gamma _{\mathrm{in}}}x^2.$$
(52)
Then, according to (27) and (51), the jet radius can be written as
$$r_\mathrm{j}=\sqrt{\frac{\sigma M_0^2}{\gamma _{\mathrm{in}}}}R_\mathrm{L},$$
(53)
which is equivalent to (35) \[and in agreement with (2)\]. Moreover, as can be easily verified, the constant magnetic field $`B_z=B(0)`$ for the invariants (18)–(20) proves to be an exact solution of (31) in the entire jet up to the jet boundary, $`\varpi =r_\mathrm{j}`$. Here, we may therefore assume $`B(0)=B_{\mathrm{ext}}`$. Consequently, according to (43), we obtain
$$M_0^2=\frac{\gamma _{\mathrm{in}}}{\sigma }\frac{B(R_\mathrm{L})}{B_{\mathrm{ext}}}.$$
(54)
Using relation (54), we can also express all the remaining jet parameters in terms of the external magnetic field.
Note that the absence of a declining solution $`B_z\varpi ^2`$ \[see (4)\] is associated with the first term in the right-hand part of (28) proportional to $`L^2`$. This term, which changes appreciably the behavior of the solution, appears to be missed previously. As it was already emphasized above, this is not surprising because the corresponding term in the second-order equation (10) is of high order and small. On the other hand, far from the rotation axis $`\varpi \gamma _{\mathrm{in}}`$, equation (28) can be rewritten as
$$\frac{d}{dx}\left[\frac{\mu \eta \mathrm{\Omega }_\mathrm{F}x^2}{M^2}\right]+\frac{\mathrm{\Omega }_\mathrm{F}^4(0)M^2}{\mathrm{\Omega }_\mathrm{F}\mu \eta x^3(x^2+M^2)}L^2=0,$$
(55)
in which both terms are of the same order. Neglecting the term proportional to $`L^2`$, we arrive at the solution (46), $`M^2x^2`$, for $`\mathrm{\Omega }_\mathrm{F}=\mathrm{const}`$ and $`\eta =\mathrm{const}`$. The conservation of function
$$H=\frac{\mu \eta \mathrm{\Omega }_\mathrm{F}x^2}{M^2}$$
(56)
was first found by Heyvaerts and Norman for conical solutions, when all quantities depend only on spherical coordinate $`\theta `$, but has also been repeatedly discussed when analyzing cylindrical flows. However, as we see, $`H`$ is generally not conserved in the cylindrical geometry for relativistic jets. To be more precise, the second term in (55) turns out to be significant for all models with a nearly constant density of the longitudinal electric current in the central jet region, where the invariant $`L(\mathrm{\Psi })`$ linearly increases with magnetic flux $`\mathrm{\Psi }`$ if $`\mathrm{\Psi }\mathrm{\Psi }_0`$.
Thus, we conclude that the solution with a central core (4) cannot be realized in the presence of an external medium with a finite regular magnetic field. This conclusion appears to be also valid in the presence of a medium with finite pressure $`P`$. Indeed, since the magnetic flux (48) increases very slowly (logarithmically), the solution (4) yields an exponentially large jet radius $`r_\mathrm{j}R_\mathrm{L}\mathrm{exp}(\mathrm{\Psi }_0/\mathrm{\Psi }_{\mathrm{in}})`$. Accordingly, the magnetic energy density must also be low at $`\varpi r_\mathrm{j}`$. However, this configuration cannot exist in the presence of an external medium with finite pressure $`P`$, irrespective of whether it is produced by a magnetic field or by a plasma. We may therefore conclude that the solutions with a central core can be realized only for a special choice of the integral $`L(\mathrm{\Psi })`$, which increases only slightly with the magnetic flux, and only in the absence of an external medium. For the most natural (from our point of view) models with a constant current density in the central jet regions, the solution with a central core cannot be realized even in the absence of an external medium.
In order to derive now the energy distribution in the jet and the particle Lorentz factor, it is convenient to introduce the quantity
$$g(x)=\frac{M^2}{x^2}.$$
(57)
Since at large distances $`\varpi \gamma _{\mathrm{in}}R_\mathrm{L}`$, according to (12), we have
$$\frac{\gamma \mu \eta }{E}=\frac{g}{g+1},$$
(58)
$`g`$ is simply the ratio of the energy flux transferred by particles $`W_{\mathrm{part}}`$ to the energy flux of the electromagnetic field. As a result, from relation (52) for $`x\gamma _{\mathrm{in}}`$, we obtain
$$\frac{W_{\mathrm{part}}}{W_{\mathrm{tot}}}\frac{M_0^2}{\gamma _{\mathrm{in}}}x^11,$$
(59)
Accordingly, from (57) and (58) for the particle Lorentz factor at $`x\gamma _{\mathrm{in}}`$, we derive
$$\gamma (x)=x.$$
(60)
Finally, expression (13) for the $`\phi `$-component of the 4-velocity vector $`u_\phi `$ yields the following toroidal velocity $`v_\phi =u_\phi /\gamma `$ at $`x\gamma _{\mathrm{in}}`$:
$$v_\phi (x)=\frac{1}{x}.$$
(61)
We see that the particle energy approaches the universal asymptotic limit (60) at $`\varpi \gamma _{\mathrm{in}}R_\mathrm{L}`$. Naturally, such a simple asymptotics can also be derived from simpler considerations. Indeed, using the “frozen-in” equation $`𝐄+𝐯\times 𝐁=0`$, we obtain for the drift velocity
$$U_{\mathrm{dr}}^2=\frac{|𝐄|^2}{|𝐁|^2}=\left(\frac{B_\phi ^2}{|𝐄|^2}+\frac{B_z^2}{|𝐄|^2}\right)^1.$$
(62)
In our case, however, according to (6) and (7), we have
$$B_\phi ^2|𝐄|^2,$$
(63)
$$|𝐄|^2x^2B_z^2.$$
(64)
As a result, relations (62)–(64) immediately lead to the exact asymptotics (60). It thus follows that, for example, for electron–positron jets from AGNs, the jet particle energy is typically
$$Em_\mathrm{e}c^2\frac{r_j}{R_L}(10^410^5)\mathrm{Mev}.$$
(65)
On the other hand, according to (57), we reach a very important conclusion that, far from the light cylinder, $`\varpi \gamma _{\mathrm{in}}R_\mathrm{L}`$,
$$g1.$$
(66)
Consequently, according to (58), the particle contribution to the general energy flux balance proves to be minor. For example, at $`B_{\mathrm{ext}}B_{\mathrm{min}}`$ for $`rr_\mathrm{j}`$, we have
$$\frac{W_{\mathrm{part}}}{W_{\mathrm{tot}}}\sqrt{\frac{\gamma _{\mathrm{in}}}{\sigma }}\text{ .}$$
(67)
While in the general case, we obtain
$$\frac{W_{\mathrm{part}}}{W_{\mathrm{tot}}}\frac{1}{\sigma }\left[\frac{B(R_\mathrm{L})}{B_{\mathrm{ext}}}\right]^{1/2}.$$
(68)
Thus, we reach a fairly nontrivial conclusion that the fraction of energy transferred by particles in a one-dimensional jet must be determined by the parameters of the external medium.
Let us now discuss the results of exact calculations obtained by numerical integration of equations (25) and (28) with the integrals of motion (18)–(20) and (23). In figures 1(a) and 1(b), the Mach number and the energy flux concentrated in particles $`\gamma \mu \eta `$ are plotted against $`x=\mathrm{\Omega }_\mathrm{F}\varpi `$ for $`M_0^2=16`$, $`\gamma _{\mathrm{in}}=8`$, and $`\sigma =1000`$. The dashed lines indicate the behavior of these quantities that follows from the analytical asymptotics (52) and (60). As we see, at sufficiently small $`x`$ when the integrals of motion (18)–(20) and (23) are similar to (36)–(38), the analytical asymptotics match the exact numerical results. On the other hand, as expected, $`\gamma \mu \eta `$ and $`M^2`$ are zero at $`\mathrm{\Psi }=\mathrm{\Psi }_0`$, i.e., at the jet edge. Figure 1(c) shows the dependence of the poloidal field, $`x^1(dy/dx)`$, for the inner parts of the jet, $`\mathrm{\Psi }<\mathrm{\Psi }_0`$. We see from this figure that the magnetic field is nearly constant at $`x>\gamma _{\mathrm{in}}`$, in agreement with the analytic estimates (43) and (51). Of course, in general, the structure of the poloidal magnetic field is determined by a specific choice of the integrals $`E(\mathrm{\Psi })`$ and $`L(\mathrm{\Psi })`$.
In conclusion, it is of interest to compare the energy of jet particle with the limiting energy acquired by the particles as they outflow from the magnetosphere of a compact object with a monopole magnetic field. According to calculations by Beskin et al. , the particle Lorentz factor outside the fast magnetosonic surface $`r>\sigma ^{1/3}R_\mathrm{L}`$ in the absence of an external medium can be written as
$`\gamma (y)`$ $`=`$ $`y^{1/3},y>\gamma _{\mathrm{in}}^3,`$ (69)
$`\gamma (y)`$ $`=`$ $`\gamma _{\mathrm{in}},y<\gamma _{\mathrm{in}}^3,`$ (70)
where $`y`$ is given by (27). On the other hand, relations (51) and (60) for the jet yield
$`\gamma (y)`$ $`=`$ $`\left({\displaystyle \frac{M_0^2}{\gamma _{\mathrm{in}}}}\right)^{1/2}y^{1/2},y>{\displaystyle \frac{\gamma _{\mathrm{in}}^3}{M_0^2}},`$ (71)
$`\gamma (y)`$ $`=`$ $`\gamma _{\mathrm{in}},y<{\displaystyle \frac{\gamma _{\mathrm{in}}^3}{M_0^2}}.`$ (72)
As shown in figure 2(a), for $`M^2>1`$, i.e., for $`B_{\mathrm{ext}}<B_{\mathrm{cr}}`$, where
$$B_{\mathrm{cr}}=\frac{\gamma _{\mathrm{in}}}{\sigma }B(R_\mathrm{L}),$$
(73)
the Lorentz factor of the jet particles (71) is always larger than the Lorentz factor acquired by the particles as they outflow from a magnetosphere with a monopole magnetic field, but, of course, is always smaller than the critical Lorentz factor
$$\gamma (y)=y,$$
(74)
which corresponds to the complete transformation of the electromagnetic energy into the particle energy. This implies that, at $`B_{\mathrm{ext}}<B_{\mathrm{cr}}`$, the particles must be additionally accelerated during the collimation coupled with the interaction of the outflowing plasma with the external medium. If, alternatively, $`B_{\mathrm{ext}}>B_{\mathrm{cr}}`$, then, in the inner jet regions, at
$$\varpi <\frac{\gamma _{\mathrm{in}}}{M_0^2}R_\mathrm{L},$$
(75)
the particle energy on a given field line turns out to be even lower than that for a monopole magnetic field, as shown in figure 2(b).
The latter result can be easily explained. Indeed, for the integrals of motion (36)–(38) we consider, the factor $`D`$, whose zero value determines the location of the fast MHD surface (see for more detail), can be rewritten in the case of a cold plasma as
$`M^2DA+{\displaystyle \frac{B_\phi ^2}{B_\mathrm{p}^2}}=A+{\displaystyle \frac{4x^2y^2M^4}{4y^2M^42x^2M^2x^4}}.`$ (76)
It is easy to show that expression (76) for $`y`$ and $`M^2`$ given by (51) and (52) is negative at
$$M^21,$$
(77)
i.e., at $`x`$ corresponding to (75). Consequently, we may reach another important conclusion that, for sufficiently strong external magnetic fields $`B_{\mathrm{ext}}>B_{\mathrm{cr}}`$ (73) when $`M^2<1`$, a region with a subsonic flow inevitably emerges in the inner jet regions $`\varpi <r_s`$, where
$$r_\mathrm{s}\sigma \left[\frac{B_{\mathrm{ext}}}{B(R_\mathrm{L})}\right]^{3/2}R_\mathrm{L}.$$
(78)
At the same time, a region with the subsonic flow can be produced far from the compact object either by a shock wave or by a strong distortion of the magnetic field within the fast magnetosonic surface located in the vicinity of the compact object. In both cases, the magnetic-field perturbation causes the particle energy to decrease.
## 4 Discussion
Thus, we conclude that the exact equilibrium equations (25) and (28) do actually allow a construction of a self-consistent model for a jet immersed in an external uniform magnetic field. The advantage of these equations over equation (10) results from the fact that all terms in (28) are of the same order. In this case, the uniformity of the poloidal magnetic field within the jet (43) results from the choice of integrals (36)–(38). In general, the poloidal magnetic field depends on the specific form of the integrals. A full analysis of the possible solutions is beyond the scope of this paper.
We have shown that the fraction of energy transferred by particles $`W_{\mathrm{part}}/W_{\mathrm{tot}}`$ must be largely determined by the parameters of the external medium. In case $`\sigma >>\sigma _{\mathrm{cr}}`$, where
$$\sigma _{\mathrm{cr}}=\left[\frac{B(R_\mathrm{L})}{B_{\mathrm{ext}}}\right]^{1/2},$$
(79)
the energy transferred by particles is only a small fraction of the energy flux $`W_{\mathrm{em}}`$ transferred by the electromagnetic field. Consequently, the jet is strongly magnetized ($`W_{\mathrm{part}}W_{\mathrm{tot}}`$) only at sufficiently large $`\sigma `$. If, however, the magnetization parameter does not exceed $`\sigma _{\mathrm{cr}}`$, then, in this case, an appreciable part of the energy in the jet is transferred by particles. This, in turn, implies that a considerable part of the energy must be transferred from the electromagnetic field to the plasma particles during jet collimation. It is interesting that $`\sigma _{\mathrm{cr}}`$ turns out to be approximately the same both for AGNs and for fast radio pulsars:
$$\sigma _{\mathrm{cr}}10^510^6.$$
(80)
We have shown that the central part of the jet must be subsonic for sufficiently strong external magnetic fields. Thus, the theory gives direct predictions whose validity can be verified by observations. It should also be noted that the results obtained above are applicable both to electron–positron and to electron–proton jets. However, a direct evidence that the jets in AGNs are actually electron–positron ones has recently appeared .
In our view, an important result is that, if the external regular magnetic field is taken into account, the MHD equations allow a self-consistent model to be constructed for a jet with a zero total longitudinal electric current, $`I(\mathrm{\Psi }_0)=0`$. In this case, a uniform magnetic field that matches the external magnetic field can also be a solution for the inner jet regions. As was already emphasized above, the radii of the jets from AGNs can thus be also explained in a natural way. In addition, since only a small fraction of the electromagnetic-field energy is transformed into the particle energy, the energy transfer from the compact object in the region of energy release can be explained as well. At the same time, extending the MHD solution to the jet region requires very high particle energies ($`10^4\mathrm{MeV}`$), which have not been recorded yet. However, a consistent discussion of the outflowing-plasma energy requires a proper allowance for the particle interaction with the surrounding medium (for example, with background radiation), which may cause a significant change in particle energy.
As for the quantitative predictions about the real physical parameters of jets, they, as we showed above, essentially depend only on the following three quantities: the magnetization parameter $`\sigma `$ (17), the Lorentz factor $`\gamma _{\mathrm{in}}`$ in the generation region, and the external magnetic field $`B_{\mathrm{ext}}`$. In this case, the main uncertainty for electron–positron jets from AGNs ($`B_{\mathrm{in}}10^4\mathrm{G}`$, $`R_{\mathrm{in}}10^{14}\mathrm{cm}`$) is the value of the magnetization parameter. Indeed, this quantity depends on the efficiency of pair production in the magnetosphere of a black hole, which, in turn, is determined by the density of hard gamma-ray photons. As a result, if the density of hard gamma-ray photons with energies $`E_\gamma >1\mathrm{MeV}`$ near the black hole is high enough, then the particles will be produced by direct collisions of photons $`\gamma +\gamma e^++e^{}`$ . This causes an abrupt increase in the multiplicity parameter $`\lambda =n/n_{\mathrm{GJ}}10^{10}`$$`10^{12}`$, where $`n_{\mathrm{GL}}\mathrm{\Omega }B/2\pi c`$ is the characteristic particle density required to shield the longitudinal electric field. Using the well-known estimate (see, e.g., )
$$\sigma \frac{\mathrm{\Omega }eB_{\mathrm{in}}R_{\mathrm{in}}^2}{\lambda m_ec^3},$$
(81)
we obtain
$$\sigma 1010^3.\gamma _{\mathrm{in}}310.$$
(82)
On the other hand, for low densities of gamma-ray photons when an electron–positron plasma can be produced only in regions with a nonzero longitudinal electric field, which are equivalent to the outer gaps in the magnetospheres of radio pulsars , the multiplicity of the particle production is fairly small: $`\lambda 10`$$`100`$. In this case, we obtain
$$\sigma 10^{11}10^{13},\gamma _{\mathrm{in}}10.$$
(83)
Finally, for fast Crab- or Vela-type radio pulsars ($`B_{\mathrm{in}}10^{13}\mathrm{G}`$, polar-cap radius $`R_{\mathrm{in}}10^5\mathrm{cm}`$, and $`\lambda =n/n_{\mathrm{GJ}}10^4`$) in which jets are observed, we have
$$\sigma 10^610^7,\gamma _{\mathrm{in}}10^210^3.$$
(84)
Relations (82) and (83) show that the properties of the jets from AGNs considerably depend on the magnetization parameter $`\sigma `$. For example, according to (79), the jet particles for sources with large $`\sigma 10^{12}`$ transfer only a small fraction of the energy compared to the electromagnetic flux, so the flow within the jet differs only slightly from the force-free flow. In addition, in this case the external magnetic field $`B_{\mathrm{ext}}10^6\mathrm{G}`$ exceeds the critical magnetic field $`B_{\mathrm{cr}}10^7\mathrm{G}`$. According to (43) and (73), this implies that a subsonic region must exist in the inner regions of such jets. On the other hand, a substantial part of the energy in sources with magnetization parameter $`\sigma 100`$ during jet collimation must be transferred by plasma particles, and no subsonic region is formed near the rotation axis. As for the fast radio pulsars, the condition $`B_{\mathrm{ext}}B_{\mathrm{cr}}`$ is satisfied for them, so a subsonic region in the central parts of pulsar jets is not achieved either. On the other hand, the estimate (79) shows that an appreciable part of the jet total energy must be coupled with particles.
We emphasize that, since the jet radius (53) for AGNs always exceeds the light-cylinder radius by several orders of magnitude,
$$r_\mathrm{j}(10^410^5)R_\mathrm{L},$$
(85)
the toroidal magnetic field $`B_\varphi `$ within the jet must exceed the poloidal magnetic field $`B_p`$ in the same proportion,
$$B_\phi \frac{r_\mathrm{j}}{R_\mathrm{L}}B_\mathrm{p}(10^410^5)B_{\mathrm{ext}}.$$
(86)
Consequently, detection of such a strong toroidal component would be a crucial argument for the picture discussed here. Unfortunately, determination of the actual physical conditions in jets currently involves considerable difficulties. Nevertheless, not only data on the direct detection of such a structure but also evidence for the existence of magnetic fields $`B0.1\mathrm{G}`$, closely matching the estimate (86, have recently appeared.
## Acknowledgments
We wish to thank A.V. Gurevich for interest in this study, a useful discussion, and support. We also wish to thank S.V. Bogovalov, L.I. Gurvits, and S.A. Lamzin for fruitful discussions. This study was supported in part by the INTAS grant no. 96–154 and the Russian Foundation for Basic Research (project number 99-02-17184). L. Malyshkin is also grateful to the International Science Foundation.
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# How weak is weak extent?
Abstract.We show that the extent of a Tychonoff space of countable weak extent can be arbitrary big. The extent of $`X`$ is $`e(X)=\mathrm{sup}\{|F|:FX`$ is closed and discrete$`\}`$ while $`we(X)=\mathrm{min}\{\tau :\text{ for every open cover }𝒰\text{ of }X\text{ there is }AX\text{ such that }|A|\tau \text{ and }St(A,𝒰)=X\}`$ is the weak extent of $`X`$ (also called the star-Lindelöf number of $`X`$). Also we show that the extent of a normal space with countable weak extent is not greater than $`𝐜`$.
Keywords: extent, weak extent, star-Lindelöf number, linked-Lindelöf number, normal space.
AMS Subject Classification: 54A25, 54D20
Recall that the extent of a topological space $`X`$ is the cardinal $`e(X)=\mathrm{sup}\{|F|:FX`$ is closed and discrete$`\}`$. The weak extent of $`X`$ is the cardinal $`we(X)=\mathrm{min}\{\tau :`$ for every open cover $`𝒰`$ of $`X`$ there is $`AX`$ such that $`|A|\tau `$ and $`St(A,𝒰)=X\}`$ . The reason for this name is that for any $`X\mathrm{T}_1`$, $`we(X)e(X)`$; indeed, supposing $`we(X)>\kappa `$, there is an open cover such that for every $`AX`$ with $`|A|\kappa `$ one has $`St(A,𝒰)X`$; then one can inductively choose points $`x_\alpha `$, $`\alpha <\kappa `$, so that $`x_\alpha St(\{x_\beta :\beta <\alpha \},𝒰)`$ for each $`\alpha `$; once the points have been choosen the set $`\{x_\alpha :\alpha <\kappa \}`$ is closed, discrete and of cardinality $`\kappa `$, so $`e(X)\kappa `$. Note also that $`we(X)d(X)`$ obviously holds for every $`X`$. Some cardinal inequalities involving extent can be improved by replacing extent by weak extent. Thus for $`XT_1`$, $`|K(X)|we(X)^{psw(X)}`$ . A natural question was stated in , : how big can the difference between the extent and the weak extent of a $`\mathrm{T}_i`$ space be? First, we give the answer for the Tychonoff case.
###### Theorem 1
For every cardinal $`\tau `$ there is a Tychonoff space $`X`$ such that $`e(X)\tau `$ and $`we(X)=\omega `$.
Before the paper the cardinal function $`we(X)`$ was called the star-Lindelöf number , , . In particular, a space $`X`$ such that $`we(X)=\omega `$ is called star-Lindelöf or Lindelöf, see e.g. , , .
Note that $`e(X)2^{we(X)\chi (X)}`$ for every regular space $`X`$ .
To prove Theorem 1, we use a set-theoretic fact in Theorem 2 below.. Let $`S`$ be a set and $`\lambda `$ a cardinal. A set mapping of order $`\lambda `$ on $`S`$ is a mapping that assigns to each $`sS`$ a subset $`f(s)S`$ so that $`|f(s)|<\lambda `$ and $`sf(s)`$. A subset $`TS`$ is called $`f`$-free if $`f(t)T=\mathrm{}`$ for every $`tT`$. Answering a question of Erdös, Fodor proved in 1952 (, see also , Theorem 3.1.5) a general theorem a partial case of which is the following
###### Theorem 2
(Fodor) Let $`S`$ be a set of cardinality $`\tau `$ and let $`f`$ be a set mapping on $`S`$ of order $`\omega `$. Then there is a countable family $``$ of $`f`$-free subsets of $`S`$ such that $`=S`$.
Proof of Theorem 1: Let $`\tau `$ be an infinite cardinal. For each $`\alpha <\tau `$, $`z_\alpha `$ denotes the point in $`D^\tau `$ with only the $`\alpha `$-th coordinate equal to $`1`$. Put $`Z=\{z_\alpha :\alpha <\tau \}`$. Then $`Z`$ is a discrete subspace of $`D^\tau `$. Further, let $`\kappa `$ be a cardinal such that $`\mathrm{cf}(\kappa )>\tau `$. Put
$$X=\left(D^\tau \times (k+1)\right)\left(\left(D^\tau Z\right)\times \{\kappa \}\right).$$
Also we denote $`X_0=D^\tau \times \kappa `$ and $`X_1=Z\times \{\kappa \}=\{(z_\alpha ,\kappa ):\alpha <\tau \}`$. Then $`X=X_0X_1`$.
It is clear that $`X_1`$ is closed in $`X`$ and discrete, so $`e(X)\tau `$.
It remains to prove that $`we(X)=\omega `$. First, note that $`X_0`$ is countably compact, hence star-Lindelöf. So it remains to prove that $`X_1`$ is relatively star-Lindelöf in $`X`$, i.e. for every open cover $`𝒰`$ of $`X`$ there is a countable $`AX`$ such that $`St(A,𝒰\}X_1`$.
Let $`𝒰`$ be an open cover of $`X`$. For every $`\alpha <\tau `$ choose an $`U_\alpha 𝒰`$ so that $`(z_\alpha ,\kappa )U_\alpha `$. Further, for every $`\alpha <\tau `$ choose $`\xi _\alpha <\kappa `$ and $`B_\alpha `$, an element of the standard base of $`D^\tau `$, so that $`(z_\alpha ,\kappa )(B_\alpha \times (\xi _\alpha ,\kappa ])XU_\alpha `$. It remains to check that
($`+`$) there is a countable $`CD^\tau `$ such that $`B_\alpha C\mathrm{}`$ for every $`\alpha <\tau `$.
Indeed, since $`\mathrm{cf}(\kappa )>\tau `$, there is a $`\gamma <\kappa `$ such that $`\gamma >\xi _\alpha `$ for all $`\alpha <\tau `$. Put $`A=C\times \{\gamma \}`$. Then $`U_\alpha A\mathrm{}`$ for all $`\alpha <\tau `$, so $`X_1St(A,𝒰)`$.
Now we check ($`+`$). The set $`B_\alpha `$ has the form
$$B_\alpha =\{xD^\tau :x(\alpha )=1\text{ and }x(\alpha ^{})=0\text{ for all }\alpha ^{}A_\alpha \}$$
where $`A_\alpha `$ is some finite subset of $`\tau \{\alpha \}`$. Consider the set mapping $`f`$ that assigns $`A_\alpha `$ to $`\alpha `$ for each $`\alpha <\tau `$. By Fodor’s theorem, there is a countable, $`f`$-free family $`=\{H_n:n\omega \}`$ of subsets of $`\tau `$ such that $`=\tau `$. For each $`n\omega `$, denote by $`c_n`$ the indicator function of $`H_n`$, i.e. $`c_n(\alpha )=1`$ iff $`\alpha H_n`$. Since $`H_n`$ is $`f`$-free, $`B_\alpha c_n`$ for all $`\alpha H_n`$. Put $`C=\{c_n:n\omega \}`$. Then $`B_\alpha C\mathrm{}`$ for every $`\alpha <\tau `$, i.e. ($`+`$) holds.
Pseudocompactness of $`X`$ follows from the fact that $`X`$ contains a dense countably compact subspace $`X_0`$. $`\mathrm{}`$
Now we are going to show that in the normal case the extent of a space of countable weak extent is not greater than $`𝐜`$. In fact, we will prove a slightly more general statement. Recall that a family of sets is linked if every two elements have nonempty intersection. The linked-Lindelöf number of $`X`$ is the cardinal $`ll(X)=\mathrm{min}\{\tau :`$ every open cover of $`X`$ has a subcover representable as the union of at most $`\tau `$ many linked subfamilies$`\}`$ . A space $`X`$ with $`ll(X)=\omega `$ is called linked-Lindelöf . It is easy to see that $`ll(X)we(X)`$ for every $`X`$.
###### Theorem 3
For every normal space $`X`$, $`e(X)2^{ll(X)}`$.
Proof: Let $`\tau `$ be an infinite cardinal, $`K`$ a closed discrete subspace of a normal space $`X`$ and $`|K|=k>2^\tau `$. We have to show that $`ll(X)>\tau `$. It is easy to construct a family $`𝒜`$ of subsets of $`K`$ such that $`|𝒜|=k`$ and for every nonempty finite subfamily of $`𝒜`$, say $`A_1,\mathrm{},A_n,A_{n+1},\mathrm{},A_{n+m}`$,
$$()|A_1\mathrm{}A_n(KA_{n+1})\mathrm{}(KA_{n+m})|=k.$$
For every $`A𝒜`$ pick a continuous function $`f_A:XI`$ such that $`f(A)=\{1\}`$ and $`f(KA)=\{0\}`$. Denote $`=\{f_A:A𝒜\}`$ and $`F=\mathrm{\Delta }:XI^{}`$. Then $`||=k`$. Note that $`F(K)D^{}`$. It follows from $`()`$ that $`F(K)`$ is dense in $`D^{}`$, moreover, every open set in $`D^{}`$ containes $`k`$ elements of $`F(K)`$. There is therefore a bijection $`\phi :K`$, where $``$ is the standard base of $`D^{}`$, such that $`\phi (z)F(z)`$ for every $`zK`$. Every element $`B`$ has the form
$$B=B_{f_1\mathrm{}f_n}^{i_1\mathrm{}i_n}=\{xD^{}:x(f_1)=i_1,\mathrm{},x(f_n)=i_n\}$$
where $`n`$, $`f_1,\mathrm{}f_n`$ and $`i_1,\mathrm{},i_nD`$. Denote
$$U(B)=\{xI^{}:j\{1,\mathrm{},n\}\left(\begin{array}{c}x(f_j)>\frac{1}{2}\text{ if }i_j=1\\ x(f_j)<\frac{1}{2}\text{ if }i_j=0\end{array}\right)\}.$$
Further, for every $`zK`$ put $`\stackrel{~}{\phi }(z)=U(\phi ^1(z))`$. Then $`\stackrel{~}{\phi }(z)`$ is a neighbourhood of $`F(z)`$ in $`I^{}`$. Note that
$$()\stackrel{~}{\phi }(z)\stackrel{~}{\phi }(z^{})\mathrm{}\text{ iff }\phi (z)\phi (z^{})\mathrm{}$$
Let $`𝒢`$ denote the family of all continuous functions form $`X`$ to $`I`$, $`G=\mathrm{\Delta }𝒢:XI^𝒢`$, $`\pi :I^𝒢I^{}`$ is the natural projection. For each $`zK`$ denote $`\stackrel{~}{\stackrel{~}{\phi }}(z)=\pi ^1(\phi (z))`$. Then $`\stackrel{~}{\stackrel{~}{\phi }}(z)`$ is a neighbourhood of $`G(z)`$ in $`I^𝒢`$ and
$$()\stackrel{~}{\stackrel{~}{\phi }}(z)\stackrel{~}{\stackrel{~}{\phi }}(z)\mathrm{}\text{ iff }\stackrel{~}{\phi }(z)\stackrel{~}{\phi }(z^{})\mathrm{}.$$
Last, for every $`zK`$ put $`\stackrel{~}{\stackrel{~}{\stackrel{~}{\phi }}}(z)=(\stackrel{~}{\stackrel{~}{\phi }}(z)G(K\{z\}))G(X)`$. Then $`\stackrel{~}{\stackrel{~}{\stackrel{~}{\phi }}}(z)`$ is a neighbourhood of $`G(z)`$ in $`G(X)`$ and
$$(\mathrm{v})\stackrel{~}{\stackrel{~}{\stackrel{~}{\phi }}}(z)\stackrel{~}{\stackrel{~}{\stackrel{~}{\phi }}}(z^{})\mathrm{}\text{ iff }\stackrel{~}{\stackrel{~}{\phi }}(z)\stackrel{~}{\stackrel{~}{\phi }}(z)\mathrm{}.$$
Put $`𝒰_0=\{\stackrel{~}{\stackrel{~}{\phi }}(z):zK\}`$. Since $`G`$ is a homeomorphic embedding, $`G(K)`$ is closed in $`G(X)`$, so $`O=G(X)G(K)`$ is open and hence $`𝒰=𝒰_0\{O\}`$ is an open cover of $`G(X)`$.
Since $`w(D^{})>2^\tau `$, $``$, a base of $`D^{}`$, is not representable as the union of at most $`\tau `$ many linked subfamilies (see e.g. ). By $`()`$, $`()`$ and $`(\mathrm{v})`$ the same can be said about the family $`𝒰_0`$. Note that for every $`zK`$, $`\stackrel{~}{\stackrel{~}{\stackrel{~}{\phi }}}(z)`$ is the only element of $`𝒰`$ that containes $`z`$. So $`𝒰`$ does not have a subcover representable as the union of at most $`\tau `$ many linked subfamilies and thus $`ll(X)=ll(G(X))>\tau `$.$`\mathrm{}`$
It is not clear whether the inequality in the previous theorem can be made strict, even with star-Lindelöf number instead of linked-Lindelöf.
Asknowlegement. The author expresses his gratitude to Angelo Bella and to Marion Scheepers for usefull discussions, in particular, Angelo Bella has drawn authors attention to the paper and Marion Scheepers has drawn author’s attention to the book .
The paper was written while the author was visiting the University of California, Davis. The author expresses his gratitude to colleagues from UC Davis for their kind hospitality.
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# Spin-measurement retrodiction revisited
## 1 Introduction
In contrast to classical mechanics, quantum mechanics does not allow all observables of a physical system to be measured simultaneously. Only for a set of mutually commuting observables definite values can be attributed to a physical system. Therefore it is one of the predictions of quantum mechanics that two spin components of a spin-$`\frac{1}{2}`$ particle cannot be measured simultaneously. That is, a given spin state cannot be an eigenstate of two or more spin operators $`\widehat{𝝈}𝒏_l`$, with the usual Pauli matrices $`\widehat{𝝈}=(\widehat{\sigma }_x,\widehat{\sigma }_y,\widehat{\sigma }_z)`$ and unit vectors $`𝒏_l`$. The result of a future measurement on a state can only be known in advance if the state is an eigenstate of the measurement operator. Now, if the spin is to be measured along a set of axes $`\{𝒏_l\}`$, then at most the result of one of the measurements can be predicted, namely when the spin state is an eigenstate of $`\widehat{𝝈}𝒏_i`$ for a specific $`i`$.
However, we can ask the question, whether it is possible to retrodict the results of a spin measurement along more than one possible axis. To be more specific we have the following problem: Alice prepares an initial spin-$`\frac{1}{2}`$ particle and gives it to Bob. Bob now performs one spin measurement along one of a set of axes the directions of which are known to Alice. However, Alice does not know which axis out of the set Bob chooses. After his measurement Bob gives the spin-$`\frac{1}{2}`$ particle back to Alice. Now Alice can perform another measurement on the particle, which should enable her to retrodict the value Bob obtained in his measurement. The fact that Alice does not yet know the axis along which Bob actually measures is most important. It means that she must be able to infer all possible results as a function of the axis that was chosen by Bob. After her measurement Bob tells her along which axis he actually measured and she has to tell him the result he got with certainty. This means that Bob can even cheat and tell Alice that he measured along axis 1 and Alice tells him which result (she believes!) he got. But then Bob confesses that he had lied and in fact measured along axis 2 and again Alice should be able to tell Bob the result he obtained if indeed he had measured along axis 2. And now Bob might have lied again and measured along another axis, and so on. This means that the measurement Alice makes has to extract information about all possible measurement axes out of the system, even though only one measurement was performed. The question now is: how should Alice prepare the initial spin-$`\frac{1}{2}`$ particle and what measurement does she have to perform on the particle Bob gives back to her?
To state the problem more clearly we put it into mathematical language. Alice prepares a spin-$`\frac{1}{2}`$ particle in a (possibly entangled) state $`|\psi _{AB}_{AB}:=_A_B`$. $`_B`$ is the two-dimensional Hilbert space of the spin-$`\frac{1}{2}`$ particle on which Bob can perform a spin measurement and $`_A`$ is an arbitrary Hilbert space. The state $`|\psi _{AB}`$ is measured by Bob along an axis $`𝒏_l\{𝒏_l:1lm\}`$, i.e. Bob applies the measurement operator $`\text{1}\text{1}(\widehat{𝝈}𝒏_l)`$. This measurement projects the original state onto an eigenstate
$$|\varphi _{\eta _l}(𝒏_l)=\frac{1}{2}(\text{1}\text{1}\text{1}\text{1}+\eta _l\text{1}\text{1}(\widehat{𝝈}𝒏_l))|\psi _{AB}$$
(1)
of the spin operator with eigenvalue $`\eta _l=\pm 1`$. Note that these states are not normalized if we assume normalisation for the initial state $`|\psi _{AB}`$. Now Alice can perform a measurement on $`|\varphi _{\eta _l}(𝒏_l)`$ and afterwards has to know the result $`\eta _l`$ Bob obtained, along whatever axis he measured. We denote the operator Alice applies to project onto a basis by $`\widehat{}=_j\lambda _j|\varphi _j\varphi _j|`$ with $`\{|\varphi _j\}`$ a basis of $`_{AB}`$. The task is to find $`|\psi _{AB}`$ and $`\widehat{}`$.
The whole process can be visualized like this<sup>1</sup><sup>1</sup>1The fact that Alice ‘sends’ her state to Bob should not be taken literally. She can as well prepare a state, leave and then Bob comes and looks at the state. Then spin states do not have to be parallely transported and thus complications can be avoided.:
Let us illustrate the problem by looking at the simple examples in which $`m`$=1 and $`m`$=2, where $`|\psi _{AB}`$ is unentangled and $`_A`$ is trivial. Therefore $`\text{1}\text{1}(\widehat{𝝈}𝒏_l)`$ reduces to $`\widehat{𝝈}𝒏_l`$. The $`m`$=1 case is trivially solved by either Alice preparing $`|\psi _{AB}`$ as an eigenstate of $`\widehat{𝝈}𝒏_1`$ or by Alice measuring along $`𝒏_1`$ afterwards. In the first case the result is not even retrodicted but can be predicted.
The $`m`$=2 can be solved easily as well. If Alice prepares an eigenstate of $`\widehat{𝝈}𝒏_1`$, she knows in advance which result Bob gets if he applies $`\widehat{𝝈}𝒏_1`$. To infer the value Bob obtains if he measures along $`𝒏_2`$, Alice applies the measurement operator $`\widehat{𝝈}𝒏_2`$ on the state Bob gave back to her. This procedure enables her to tell the value Bob obtained along whatever axis he measured. In fact, here we have a mixture of a prediction and a retrodiction, as the result of one possible measurement is known from the beginning.
At first glance the knowledge of spin components along different axes seems to contradict quantum mechanics. This led Vaidman, Aharanov and Albert (VAA) who first stated and solved the problem for three orthogonal axes to use the provocative title “How to Ascertain the Values of $`\widehat{\sigma }_x,\widehat{\sigma }_y`$ and $`\widehat{\sigma }_z`$ of a Spin-$`\frac{1}{2}`$ Particle” . But a closer examination of the problem reveals that the kind of information Alice has about different spin components is not the one forbidden by quantum mechanics. Alice, of course, does not really know two or more eigenvalues of non-commuting operators but only extracts conditional information. The values Alice gets are a function of Bob’s measurement axis and not until she knows along which axis Bob actually measured do these values have any physical meaning. Only when the additional information about the choice of axis is revealed, one value is singled out. This one now does have physical significance, namely it is the one eigenvalue Bob actually obtained. This also points out that if Bob cheats and tells Alice that he had measured along an axis different from the one he actually used, the result Alice tells him is meaningless.
Obviously, the Hilbert space $`_A`$ was not used in the $`m`$=1,2 examples given above and it was sufficient to prepare an unentangled state in the two-dimensional space $`_B`$. A state in a two-dimensional Hilbert space is usually called a qubit and is written as $`a|+b|`$, where $`|`$ and $`|`$ denote an orthonormal basis. (For spin-$`\frac{1}{2}`$ particles these states are often eigenstates of spin measurements along a specific axis. E.g. for the z-axis they are denoted by $`|_z`$ and $`|_z`$.)
The problem no longer has a simple solution if Bob can measure along three or more axes, which may or may not be orthogonal. Let us look at the case where Bob can measure along the three orthogonal axes $`x`$, $`y`$ and $`z`$, i.e. one of the measurement operators $`\text{1}\text{1}\widehat{\sigma }_x,\text{1}\text{1}\widehat{\sigma }_y,\text{1}\text{1}\widehat{\sigma }_z`$ is applied. We can convince ourselves that it is no longer sufficient to prepare an unentangled qubit if $`m`$ 3. Hence, if the problem can be solved at all, $`_A`$ can no longer be trivial.
The problem as stated above was first presented by VAA in , who solved it for the $`m`$=3 orthogonal case. Alice’s initial state was prepared as
$$|\psi _{AB}=\frac{1}{\sqrt{2}}\left(|_z|_z+|_z|_z\right)$$
(2)
and Alice’s measurement operator $`\widehat{}=_j\lambda _j|\varphi _j\varphi _j|`$ with $`|\varphi _j`$ as given in the appendix. In the next section a similar solution for the case that $`|\psi _{AB}`$ is the singlet state will be derived.
In Ben-Menahem used a more algebraic way to generalise the problem and showed that there are solutions for the cases $`m`$=3 and $`m`$=4, but not for $`m5`$. However for $`m`$=4, the four axes Bob may measure along are no longer independent but have to satisfy the condition $`_{l=1}^4𝒏_l=0`$ for there to be a solution. In this case, the dimension of the Hilbertspace $`_{AB}`$ is shown to be six. This space, however, is not sufficient to solve the $`m`$=3 non-orthogonal case, where $`_{AB}`$ has to be eight-dimensional.
Very recently, this method of spin-measurement retrodiction was used to set up a secure key distribution in quantum cryptography .
## 2 Solution for three orthogonal axes
After having stated the problem and given some simple examples with $`m`$=1 and $`m`$=2, we can now try to construct a solution for higher dimensional cases. For the moment we restrict ourselves to the case in which Bob can measure along three orthogonal axes, that is either of the measurement operators $`\text{1}\text{1}\widehat{\sigma }_x,\text{1}\text{1}\widehat{\sigma }_y,\text{1}\text{1}\widehat{\sigma }_z`$ can be applied to the original state $`|\psi _{AB}`$. Each of the measurements has two possible outcomes $`\eta _l=\pm 1`$, so there are six possible results altogether, which are denoted by $`|\varphi _{\eta _l}(𝒏_l)\{|,x`$, $`|,x`$, $`|,y`$, $`|,y`$, $`|,z`$, $`|,z\}_{AB}`$. Note that these are different from the basis states of the two-dimensional Hilbert space $`_B`$ which are denoted by $`\{|_x`$, $`|_x\}`$, $`\{|_y`$, $`|_y\}`$ or $`\{|_z`$,$`|_z\}`$. They will only coincide if the original state $`|\psi _{AB}`$ is an unentangled qubit.<sup>2</sup><sup>2</sup>2Even then they will only be isomorphic, as $`|_x_B`$ and $`|,x_{AB}`$ with $`_A`$ trivial.
Now Alice measures one of these states $`|\varphi _{\eta _l}(𝒏_l)`$, i.e. she applies the operator $`\widehat{}`$ to it and therefore projects the state onto a basis $`\{|\varphi _j\}`$ of $`_{AB}`$ with distinct eigenvalues $`\lambda _j`$ respectively. To ensure that her measurement gives the result she wants, Alice has to construct a look-up table. This table must tell her that if she measures e.g. $`\lambda _2`$ Bob got ‘up’ if he measured along x, ‘up’ if he measured along y and ‘down’ if he measured along z. That is, Alice needs a table like this:
| | x | y | z |
| --- | --- | --- | --- |
| $`\lambda _1`$ | $``$ | $``$ | $``$ |
| $`\lambda _2`$ | $``$ | $``$ | $``$ |
| $`\lambda _3`$ | $``$ | $``$ | $``$ |
| $`\lambda _4`$ | $``$ | $``$ | $``$ |
So if she measures $`\lambda _2`$, for example, she immediately knows that Bob sent her one of $`\{|,x,|,y,|,z\}`$. If Bob now tells her the axis along which he measured, Alice is able to tell which result he got. This leads to a geometric way of solving the problem by visualizing the $`|\varphi _{\eta _l}(𝒏_l)`$ and $`|\varphi _j`$ states as vectors in $`_{AB}`$. Alice projects the vector she gets from Bob onto her basis $`\{|\varphi _j\}`$ and gets a result $`\lambda _i`$. But then she has only enough information if each of her basis vectors is orthogonal to three of the six possible vectors Bob can send. Only then can she exclude these three and only the three other vectors can lead to the specific $`\lambda _i`$ she measured. If she chooses the states orthogonal to each of her basis vectors to be one of each of the three pairs $`\{\{|,x,|,x\},\{|,y,|,y\},\{|,z,|,z\}\}`$ the problem is now solved. If Bob tells her the axis along which he measured, Alice knows the state he obtained.
This geometric point of view yields the defining equations for the basis $`\{|\varphi _j\}`$ for a given $`|\psi _{AB}`$. As mentioned above, Alice cannot send an unentangled qubit to solve the problem. This suggests preparing $`|\psi _{AB}`$ as an entangled state of two qubits yielding a four-dimensional Hilbert space $`_{AB}`$. The entanglement of the states is necessary as Bob will measure only in a two-dimensional subspace of $`_{AB}`$. This means that, if an unentangled state was prepared, the Hilbertspace $`_A`$ would not be used at all and could as well be omitted. In fact the four-dimensional space of two qubits turns out to be necessary and sufficient to solve this case of the problem.
If Alice makes a guess and prepares $`|\psi _{AB}`$ to be the singlet state
$$|\psi _{AB}=\frac{1}{\sqrt{2}}(|_z|_z)$$
(3)
which - up to a global phase factor - looks the same in all bases. Here, the short hand notation $`|=||`$ is used. Then $`|\varphi _{\eta _l}(𝒏_l)`$ can be specified to be proportional to $`\{|_x,|_x,|_y,|_y,|_z,|_z\}`$. Now Alice wants to have the following look-up table:
$`x`$ $`y`$ $`z`$ $`\lambda _1`$ $``$ $``$ $``$ $`\lambda _2`$ $``$ $``$ $``$ $`\lambda _3`$ $``$ $``$ $``$ $`\lambda _4`$ $``$ $``$ $``$
That is, she wants
$$|\varphi _1\{|,x,|,y,|,z\},$$
$$|\varphi _2\{|,x,|,y,|,z\},$$
$$|\varphi _3\{|,x,|,y,|,z\},$$
$$|\varphi _4\{|,x,|,y,|,z\}.$$
(4)
and this choice of orthogonality conditions (4) leads to the four orthonormal vectors
$$|\varphi _1=\frac{1}{\sqrt{2}}|+\frac{1}{2}\left[\right|e^{\frac{i\pi }{4}}|e^{\frac{i\pi }{4}}]$$
$$|\varphi _2=\frac{1}{\sqrt{2}}|\frac{1}{2}\left[\right|e^{\frac{i\pi }{4}}|e^{\frac{i\pi }{4}}]$$
$$|\varphi _3=\frac{1}{\sqrt{2}}|+\frac{1}{2}\left[\right|e^{\frac{i\pi }{4}}|e^{\frac{i\pi }{4}}]$$
$$|\varphi _4=\frac{1}{\sqrt{2}}|\frac{1}{2}\left[\right|e^{\frac{i\pi }{4}}|e^{\frac{i\pi }{4}}]$$
(5)
all expressed in the z-basis. (5) together with the state (3) and the above look-up table solves the problem. (5) is very similar to the result that is given in and which is included in the appendix. In fact all possible bases can be obtained from one by a unitary transformation, as we shall see.
Now one could ask whether Alice really needs to have four distinct eigenvalues $`\lambda _j`$ and therefore a big look-up table, or whether the number of eigenvalues can be reduced. However, the reduction of distinct eigenvalues turns out to be impossible, whatever Hilbertspace $`_A`$ is used. To prove this, we assume $`_A`$ to be $`k`$-dimensional with basis $`\{|i:1ik\}`$. Then the state Alice prepares can be written as
$`|\psi _{AB}`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{k}{}}}{\displaystyle \underset{\alpha \{,\}}{}}a_{j,\alpha }|j|\alpha _z`$ (6)
$`=`$ $`{\displaystyle \underset{j}{}}a_{j,}|j|_z+{\displaystyle \underset{j}{}}a_{j,}|j|_z`$
$`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left({\displaystyle \underset{j}{}}a_{j,}|j\left(|_x+|_x\right)+{\displaystyle \underset{j}{}}a_{j,}|j\left(|_x|_x\right)\right)`$
$`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left({\displaystyle \underset{j}{}}a_{j,}|j\left(|_y+|_y\right)+{\displaystyle \underset{j}{}}a_{j,}|j(i)(|_y|_y)\right).`$
Bob’s measurement leads to either of
$`|\varphi _{+1}(𝒏_z)`$ $`=`$ $`{\displaystyle \underset{i}{}}a_{i,}|i|_z`$ (7)
$`|\varphi _1(𝒏_z)`$ $`=`$ $`{\displaystyle \underset{i}{}}a_{i,}|i|_z`$ (8)
$`|\varphi _{+1}(𝒏_x)`$ $`=`$ $`\left({\displaystyle \underset{i}{}}a_{i,}|i+{\displaystyle \underset{i}{}}a_{i,}|i\right){\displaystyle \frac{1}{\sqrt{2}}}\left(|_z+|_z\right)`$ (9)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left({\displaystyle \underset{i}{}}a_{i,}|i,_z+{\displaystyle \underset{i}{}}a_{i,}|i,_z+{\displaystyle \underset{i}{}}a_{i,}|i,_z+{\displaystyle \underset{i}{}}a_{i,}|i,_z\right)`$
$`|\varphi _1(𝒏_x)`$ $`=`$ $`\left({\displaystyle \underset{i}{}}a_{i,}|i{\displaystyle \underset{i}{}}a_{i,}|i\right){\displaystyle \frac{1}{\sqrt{2}}}(|_z|_z)`$ (10)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left({\displaystyle \underset{i}{}}a_{i,}|i,_z{\displaystyle \underset{i}{}}a_{i,}|i,_z{\displaystyle \underset{i}{}}a_{i,}|i,_z+{\displaystyle \underset{i}{}}a_{i,}|i,_z\right)`$
$`|\varphi _{+1}(𝒏_y)`$ $`=`$ $`\left({\displaystyle \underset{i}{}}a_{i,}|ii{\displaystyle \underset{i}{}}a_{i,}|i\right){\displaystyle \frac{1}{\sqrt{2}}}\left(|_z+i|_z\right)`$ (11)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left({\displaystyle \underset{i}{}}a_{i,}|i,_z+i{\displaystyle \underset{i}{}}a_{i,}|i,_zi{\displaystyle \underset{i}{}}a_{i,}|i,_z+{\displaystyle \underset{i}{}}a_{i,}|i,_z\right)`$
$`|\varphi _1(𝒏_y)`$ $`=`$ $`\left({\displaystyle \underset{i}{}}a_{i,}|i+i{\displaystyle \underset{i}{}}a_{i,}|i\right){\displaystyle \frac{1}{\sqrt{2}}}(|_zi|_z)`$ (12)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left({\displaystyle \underset{i}{}}a_{i,}|i,_zi{\displaystyle \underset{i}{}}a_{i,}|i,_z+i{\displaystyle \underset{i}{}}a_{i,}|i,_z+{\displaystyle \underset{i}{}}a_{i,}|i,_z\right).`$
We see that (12)=$`\sqrt{2}`$(7)+$`\sqrt{2}`$(8)-(11) and (10)=$`\sqrt{2}`$(7)+$`\sqrt{2}`$(8)-(9). So we showed that the six results that can be obtained by Bob’s measurement are not linearly independent but lie in a four-dimensional subspace of $`_{AB}`$. This however tells us that the number of distinct $`\lambda _j`$ has to be four. This can be easily understood in the geometric picture adopted above. Assume there are only three distinct $`\lambda _j`$. However big $`_{AB}`$ is, the subspace that contains the states (7)-(12) is four-dimensional and thus $`_{AB}`$ can be assumed to be four-dimensional as well. To ensure that Alice gets the look-up table she wants, two of the four basis vectors $`|\varphi _j`$ have to be orthogonal to three states $`\{|\varphi _{\eta _x}(𝒏_x),|\varphi _{\eta _y}(𝒏_y),|\varphi _{\eta _z}(𝒏_z)\}`$ with $`\eta _l=1`$ or $`1`$. The latter however are linearly independent and hence span a three-dimensional space. Therefore Alice cannot find a two-dimensional subspace orthogonal to all three. This proves that Alice really needs four disinct values $`\lambda _j`$ and the size of the look-up table cannot be reduced.
We also note that as soon as $`m`$3 we no longer have a mixture of prediction and retrodiction. Alice can only retrodict the results Bob obtained after having projected onto her final basis. If she prepared $`|\psi _{AB}`$ to be an eigenstate of one of the $`\text{1}\text{1}\widehat{𝝈}𝒏_l`$ she would have to find out the result Bob obtained by measuring along one of the two remaining axes by a single measurement, which is impossible. On the other hand Alice cannot send an eigenstate of two distinct spin operators even if entangled states are used, as we have $`[\text{1}\text{1}\widehat{\sigma }_i,\text{1}\text{1}\widehat{\sigma }_j]=2i\epsilon _{ijk}\text{1}\text{1}\widehat{\sigma }_k`$, and thus the measurement operators still do not commute.
## 3 Non-orthogonal axes
After having looked at the orthogonal case from as geometric point of view, we will now tackle the case in which Bob can measure along non-orthogonal axes and see that the number of axes can even be increased. For that purpose we shall take a more algebraic route, following .
In the previous case in section 2, Alice’s basis states $`|\varphi _j`$, and therefore the operator $`\widehat{}`$, could be derived from the condition that each of Alice’s basis states $`|\varphi _i`$ is orthogonal to three of Bob’s possible results $`|\varphi _{\eta _l}(𝒏_l)`$. Alice just guessed that $`|\psi _{AB}`$ might either be (3) or (2), hence calculated $`|\varphi _{\eta _l}(𝒏_l)`$ and via (4) found a basis for $`\widehat{}`$. However, to solve the problem in general, both $`|\psi _{AB}`$ and $`\widehat{}`$ should be calculated, i.e. $`|\psi _{AB}`$ should not be found by trial and error. To attack this problem, we first note that Alice’s initial state can be written as
$$|\psi _{AB}=\underset{j=1}{\overset{2k}{}}b_j|\varphi _j,$$
(13)
where $`\{|\varphi _j\}`$ is a basis of $`_{AB}`$. Without loss of generality the $`b_j`$ can be assumed to be real. Then the problem reduces to finding the basis states $`|\varphi _j`$ together with the values $`b_j`$, as these give both $`\widehat{}`$ and $`|\psi _{AB}`$. To find these we look at the conditions they have to satisfy, to ensure Alice can retrodict Bob’s results. First of all $`|\psi _{AB}`$ should be normalised, which gives:
$$\underset{j}{}b_j^2=1.$$
(14)
A second condition is that, in order to enable Alice to retrodict Bob’s results $`\eta _l`$ deterministically, the states $`|\varphi _{+1}(𝒏_l)`$ and $`|\varphi _1(𝒏_l)`$ must lie in the span of disjoint subsets of $`\left\{|\varphi _j\right\}`$. Otherwise the result $`\lambda _i`$, which Alice gets if she measures using $``$, would not tell her $`\eta _l`$. As this must be true for all axes $`l`$ along which Bob can measure, there are $`m`$ pairs of disjoint subsets. Therefore, let us introduce $`2m`$ sets of indices $`S_{\eta _l}(𝒏_l)`$ that indicate which basis vectors span the vectors
$$|\varphi _{\eta _l}(𝒏_l)=\underset{jS_{\eta _l}(𝒏_l)}{}b_j|\varphi _j.$$
(15)
Here the $`b_j`$ are the same as in (13), which can be seen from
$$|\psi _{AB}=|\varphi _+(𝒏_l)+|\varphi _{}(𝒏_l).$$
(16)
This equation only states that Bob always gets either ‘up’ or ‘down’. So the $`S_{\eta _l}(𝒏_l)`$ provide a look-up table for Alice. In the case she measures $`\lambda _i`$ corresponding to $`|\varphi _i`$, Alice checks whether $`iS_{+1}(𝒏_l)`$ or $`iS_1(𝒏_l)`$ which tells her which result Bob got if he measured along $`𝒏_l`$.
For further convenience it is helpful to introduce a sign function for the partitions $`S_{\eta _l}(𝒏_l)`$ in the following way:
$$\epsilon _j^{(l)}=\{\begin{array}{c}+1\text{for}jS_{+1}(𝒏_l)\\ 1\text{for}jS_1(𝒏_l)\end{array}$$
(17)
For example, if $`|\varphi _{+1}(𝒏_l)=b_1|\varphi _1+b_3|\varphi _3+b_6|\varphi _6`$, $`\epsilon _1^{(l)}`$, $`\epsilon _3^{(l)}`$ and $`\epsilon _6^{(l)}`$ give +1, the others give -1.
Now some algebraic manipulations have to be done to obtain sufficient information for Alice to construct her states. The basic task is to find some constraints for the look-up table that Alice chooses. Once a “good” table is chosen, the basis states can be constructed. From (16) we get (recalling (1)):
$$\text{1}\text{1}(\widehat{\sigma }𝒏_l)|\psi _{AB}=\underset{j}{}ϵ_j^{(l)}b_j|\varphi _j\text{for}\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}lm.$$
(18)
Furthermore the Pauli-matrices satisfy<sup>3</sup><sup>3</sup>3Note that $`(\widehat{\text{1}\text{1}𝝈})𝒏_a=\text{1}\text{1}(\widehat{𝝈}𝒏_a)`$.:
$$((\widehat{\text{1}\text{1}𝝈})𝒏_a)((\widehat{\text{1}\text{1}𝝈})𝒏_b)=𝒏_a𝒏_b+i(𝒏_a\times 𝒏_b)\widehat{\text{1}\text{1}𝝈}.$$
(19)
The expectation value of (19) yields:
$${}_{AB}{}^{}\psi |𝒏_a𝒏_b|\psi _{AB}^{}+i_{AB}\psi |(𝒏_a\times 𝒏_b)\widehat{\text{1}\text{1}𝝈}|\psi _{AB}=\underset{j}{}\epsilon _j^{(a)}\epsilon _j^{(b)}b_j^2$$
(20)
and therefore the two equations,
$$0=_{AB}\psi |(𝒏_l(\widehat{\text{1}\text{1}𝝈}))|\psi _{AB}=\underset{j}{}\epsilon _j^{(l)}b_j$$
(21)
and
$$𝒏_l𝒏_j=\underset{s}{}\epsilon _s^{(l)}\epsilon _s^{(j)}b_s^2,$$
(22)
for the real and imaginary components. From (21) we obtain:
$$\underset{jS_\eta }{}b_j^2=\frac{1}{2},$$
(23)
which tells us that the square of the coefficients of the basis states $`|\varphi _j`$ have to add to $`\frac{1}{2}`$ for each disjoint subset. Finally we note that in the case in which $`m>3`$, of course three axes are sufficient to span all the others:
$$𝒏_{k+3}=\underset{l=1}{\overset{3}{}}c_l^{(k)}𝒏_l\text{for}\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}km3,$$
(24)
and thus from (18):
$$\epsilon _s^{(k+3)}=\underset{l=1}{\overset{3}{}}c_l^{(k)}\epsilon _s^{(l)}\text{for}\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}s2k.$$
(25)
It turns out that these relations enable Alice to calculate both the $`b_j`$ and the $`|\varphi _j`$. It was shown in that as a necessary and sufficient condition for the construction of a basis we need (i) $`\epsilon _l^{(l)}`$ and $`b_j`$ that solve (14), (21) and (22); and (ii) if $`m`$3, there exist numbers $`c_l^{(k)}`$ $`ll3`$ and $`lkm3`$ such that (25) holds. Furthermore it was proved that no solutions exist for $`m>4`$ and that if $`m=4`$ the directions have to satisfy $`_{l=1}^4𝒏_l=0`$. But all this only means that the look-up table has to satisfy certain conditions. This, however, was anticipated already as not any arbitrary look-up table should lead to a solution.
However, once Alice chooses a suitable look-up table, the basis states can be derived. The orientation of Bob’s axes (22), together with Alice condition (23), lead to a unique set of components $`b_j`$. The crucial equation (18) now tells us the action of Bob‘s measurement on $`|\psi _{AB}`$. These can be inverted to obtain $`\{|\varphi _j\}`$. The details of this process can best be understood by looking at an example.
In the following we will explicitely construct the basis for the case m=4, using this method. The non-orthogonal $`m=3`$ case can be solved equivalently and a possible basis is given in the appendix.
The problem is solved using a six-dimensional Hilbert space $`_{AB}`$. A basis of this space is denoted by $`\{|\mathrm{\hspace{0.17em}2},1,|\mathrm{\hspace{0.17em}2},1,\mathrm{},|\mathrm{\hspace{0.17em}0},1`$ where $`|\rho ,\zeta =|\rho |\zeta `$ with $`\rho =1,2,3`$ and $`\zeta =\pm 1`$ is the tensor product of a qutrit and a qubit.
We use the partition provided in , namely
$$S_+(𝒏_1)=\{1,2,3\},$$
$$S_+(𝒏_2)=\{1,5,6\},$$
$$S_+(𝒏_3)=\{3,4,6\}.$$
(26)
which satisfy the necessary conditions. From (25) we get $`\epsilon _1^{(4)}=\epsilon _3^{(4)}=\epsilon _6^{(4)}=1;\epsilon _2^{(4)}=\epsilon _4^{(4)}=\epsilon _5^{(4)}=1`$. This is nothing but the following look-up table for Alice:
$`𝒏_1`$ $`𝒏_2`$ $`𝒏_3`$ $`𝒏_4`$ $`\lambda _1`$ $``$ $``$ $``$ $``$ $`\lambda _2`$ $``$ $``$ $``$ $``$ $`\lambda _3`$ $``$ $``$ $``$ $``$ $`\lambda _4`$ $``$ $``$ $``$ $``$ $`\lambda _5`$ $``$ $``$ $``$ $``$ $`\lambda _6`$ $``$ $``$ $``$ $``$
From (23) this table gives immediatly
$$b_1^2+b_2^2+b_3^2=b_4^2+b_5^2+b_6^2=b_1^2+b_5^2+b_6^2=b_3^2+b_4^2+b_6^2=\frac{1}{2}.$$
(27)
These equations can be solved in terms of the two free parameters $`b_5`$ and $`b_6`$:
$$b_1^2=b_4^2=\frac{1}{2}b_5^2b_6^2;b_2^2=b_6^2;b_3^2=b_5^2.$$
(28)
(22) yields:
$$𝒏_1𝒏_2=14b_5^24b_6^2;𝒏_2𝒏_3=4b_6^21;𝒏_3𝒏_1=4b_5^21.$$
(29)
Therefore the two parameters $`b_5,b_6`$ specify the mutual orientation of the axes. These are all the preliminaries we need to construct our basis. To specify the axes we make the symmetric choice:
$$b_j=\frac{1}{\sqrt{6}}j.$$
The crucial step is that (18) tells us the action of $`\widehat{\sigma }_x,\widehat{\sigma }_y,\widehat{\sigma }_z`$ on $`|\psi _{AB}`$, giving four orthonormal states:
$`|\psi _{AB}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{6}}}(|\varphi _1+|\varphi _2+|\varphi _3+|\varphi _4+|\varphi _5+|\varphi _6)`$ (30)
$`(\text{1}\text{1}\widehat{\sigma }_x)|\psi _{AB}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{6}}}(|\varphi _1+|\varphi _2+|\varphi _3|\varphi _4|\varphi _5|\varphi _6)`$ (31)
$`(\text{1}\text{1}\widehat{\sigma }_y)|\psi _{AB}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{12}}}(2|\varphi _1|\varphi _2|\varphi _32|\varphi _4+|\varphi _5+|\varphi _6)`$ (32)
$`(\text{1}\text{1}\widehat{\sigma }_z)|\psi _{AB}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(|\varphi _2+|\varphi _3|\varphi _5+|\varphi _6)`$ (33)
These are orthonormal and can be extended to a basis by
$$|\chi _1=\frac{1}{2}(|\varphi _1+|\varphi _2|\varphi _4+|\varphi _6)$$
(34)
and
$$|\chi _2=\frac{1}{12}(|\varphi _1|\varphi _2+2|\varphi _3|\varphi _4+2|\varphi _5|\varphi _6).$$
(35)
These equations have to be inverted to get explicit expressions for the basis states $`|\varphi _j`$. In order to express them in terms of the basis $`\{|\mathrm{\hspace{0.17em}2},1,|\mathrm{\hspace{0.17em}2},1,\mathrm{},|\mathrm{\hspace{0.17em}0},1`$ of the six-dimensional vector space $`_{AB}`$, we note that $`\frac{1}{\sqrt{2}}((\text{1}\text{1}\widehat{\sigma }_x)+i\zeta (\text{1}\text{1}\widehat{\sigma }_y))|\psi _{AB}`$ is an eigenstate of $`(\text{1}\text{1}\widehat{\sigma }_z)`$ with eigenvalue $`\zeta =\pm 1`$, as well as $`\frac{1}{\sqrt{2}}((\text{1}\text{1}\text{1}\text{1}+\zeta (\text{1}\text{1}\widehat{\sigma }_y))|\psi _{AB}`$ and ($`|\chi _1\delta _{1\zeta }+|\chi _2\delta _{1\zeta }`$), with the usual Kronecker symbol $`\delta _{ij}`$. Hence these states form a basis of the three-dimensional eigenspace of the operator $`\text{1}\text{1}\widehat{\sigma }_z`$ with eigenvalue $`\zeta `$. However $`|2,\zeta ,|1,\zeta ,|0,\zeta ,`$ is also a basis of this subspace, therefore there has to be a unitary transformation, such that
$$\left(\begin{array}{c}\frac{1}{\sqrt{2}}((\text{1}\text{1}\widehat{\sigma }_x)+i\zeta (\text{1}\text{1}\widehat{\sigma }_y))|\psi _{AB}\\ \frac{1}{\sqrt{2}}(\text{1}\text{1}\text{1}\text{1}+\zeta (\text{1}\text{1}\widehat{\sigma }_y))|\psi _{AB}\\ |\chi _1\delta _{1,\zeta }+|\chi _2\delta _{1,\zeta }\end{array}\right)=e^{i\widehat{\theta }_\zeta \widehat{\tau }+i\mathrm{\Lambda }_\zeta }\left(\begin{array}{c}|2,\zeta \\ |1,\zeta \\ |0,\zeta \end{array}\right).$$
(36)
Here $`\widehat{\theta }_\pm `$ are arbitrary eight-vectors, and $`\widehat{\tau }`$ is a vector consisting of the eight generators of SU(3); $`\mathrm{\Lambda }_\pm `$ are arbitrary numbers. Thus, the most general solution is characterised by a set of 24 numbers $`(𝜽_\pm ,\mathrm{\Lambda }_\pm ,\{\lambda _j\}).`$ Equation (36) allows us to express the basis states $`|\varphi _j`$ in terms of the computational basis of $`_{AB}`$. Without loss of generality, we choose $`𝜽_\zeta `$=0 and $`\mathrm{\Lambda }_\zeta `$=0 for both $`\zeta `$ and finally get:
$`|\varphi _1`$ $`=`$ $`\left({\displaystyle \frac{1}{\sqrt{1}2}}{\displaystyle \frac{i}{\sqrt{6}}}\right)|2,+{\displaystyle \frac{1}{\sqrt{1}2}}|1,{\displaystyle \frac{1}{2}}|0,+\left({\displaystyle \frac{1}{\sqrt{1}2}}+{\displaystyle \frac{i}{\sqrt{6}}}\right)|2,+{\displaystyle \frac{1}{\sqrt{1}2}}|1,{\displaystyle \frac{1}{\sqrt{1}2}}|0,`$
$`|\varphi _2`$ $`=`$ $`\left({\displaystyle \frac{1}{\sqrt{1}2}}+{\displaystyle \frac{i}{\sqrt{6}}}\right)|2,+{\displaystyle \frac{1}{\sqrt{1}2}}|1,{\displaystyle \frac{1}{2}}|0,+\left({\displaystyle \frac{1}{\sqrt{1}2}}{\displaystyle \frac{i}{\sqrt{6}}}\right)|2,+{\displaystyle \frac{1}{\sqrt{1}2}}|1,{\displaystyle \frac{1}{\sqrt{1}2}}|0,`$
$`|\varphi _3`$ $`=`$ $`\left({\displaystyle \frac{1}{\sqrt{1}2}}+{\displaystyle \frac{i}{\sqrt{2}4}}\right)|2,+\left({\displaystyle \frac{1}{\sqrt{1}2}}+{\displaystyle \frac{1}{\sqrt{8}}}\right)|1,+\left({\displaystyle \frac{1}{\sqrt{1}2}}{\displaystyle \frac{i}{\sqrt{2}4}}\right)|2,+\left({\displaystyle \frac{1}{\sqrt{1}2}}{\displaystyle \frac{1}{\sqrt{8}}}\right)|1,+{\displaystyle \frac{1}{\sqrt{3}}}|0,`$
$`|\varphi _4`$ $`=`$ $`\left({\displaystyle \frac{1}{\sqrt{1}2}}{\displaystyle \frac{i}{\sqrt{2}4}}\right)|2,+\left({\displaystyle \frac{1}{\sqrt{1}2}}{\displaystyle \frac{1}{\sqrt{8}}}\right)|1,+\left({\displaystyle \frac{1}{\sqrt{1}2}}+{\displaystyle \frac{i}{\sqrt{2}4}}\right)|2,+\left({\displaystyle \frac{1}{\sqrt{1}2}}+{\displaystyle \frac{1}{\sqrt{8}}}\right)|1,+{\displaystyle \frac{1}{\sqrt{3}}}|0,`$
$`|\varphi _5`$ $`=`$ $`\left({\displaystyle \frac{1}{\sqrt{1}2}}{\displaystyle \frac{i}{\sqrt{2}4}}\right)|2,+\left({\displaystyle \frac{1}{\sqrt{1}2}}+{\displaystyle \frac{1}{\sqrt{8}}}\right)|1,+{\displaystyle \frac{1}{2}}|0,+\left({\displaystyle \frac{1}{\sqrt{1}2}}+{\displaystyle \frac{i}{\sqrt{2}4}}\right)|2,+\left({\displaystyle \frac{1}{\sqrt{1}2}}{\displaystyle \frac{1}{\sqrt{8}}}\right)|1,{\displaystyle \frac{1}{\sqrt{1}2}}|0,`$
$`|\varphi _6`$ $`=`$ $`\left({\displaystyle \frac{1}{\sqrt{1}2}}+{\displaystyle \frac{i}{\sqrt{2}4}}\right)|2,+\left({\displaystyle \frac{1}{\sqrt{1}2}}+{\displaystyle \frac{1}{\sqrt{8}}}\right)|1,+{\displaystyle \frac{1}{2}}|0,+\left({\displaystyle \frac{1}{\sqrt{1}2}}{\displaystyle \frac{i}{\sqrt{2}4}}\right)|2,+\left({\displaystyle \frac{1}{\sqrt{8}}}+{\displaystyle \frac{1}{\sqrt{1}2}}\right)|1,{\displaystyle \frac{1}{\sqrt{1}2}}|0,,`$
all expressed in the z-basis.
So we have explicitly constructed a basis for a special symmetric orientation of our axes. This technique can be applied to all possible orientations that satisfy the condition $`_{l=1}^4𝒏_l=0`$.
## 4 The Quantum Network
After having explicitly constructed some measurement bases, we now want to know how Alice’s measurements can actually be performed. Thus, we want to construct a device that prepares Alice’s initial state, and the operator $`\widehat{}`$ that projects onto one of the bases that were just determined.
In principle, this can be achieved by constructing a quantum network which is a kind of building plan for Alice’s devices. For that purpose we regard the Hilbert space $`_{AB}`$ as a subspace of the $`2^N`$-dimensional space $`=_B_B\mathrm{}_B`$. The basis of this space is taken to be the tensor product of the qubit basis of $`_B`$, i.e. the states $`|\mathrm{},|\mathrm{}`$, … <sup>4</sup><sup>4</sup>4It is convenient to use $`|\mathrm{\hspace{0.17em}0}:=|`$ and $`|\mathrm{\hspace{0.17em}1}:=|`$ in this section.. In this space of qubits we can now perform certain operations and in that way try to construct the whole experimental device out of a few basic transformations which we call quantum gates. As we demand reversibility, quantum gates have to be unitary. (We could also say, that as we have quantum gates which obey the rules of quantum mechanics, we have unitary transformations and therefore reversibility). It was shown in that there is a set of universal quantum gates from which every evolution of a quantum system can be approximated with arbitrary precision. A set of universal gates can be taken to be the Phase Shift $`P`$($`\varphi `$), defined by
$`|0`$ $``$ $`|0`$
$`|1`$ $``$ $`e^{i\varphi }|1,`$ (37)
the Controlled Not ($`CNOT`$),
$`|00`$ $``$ $`|00`$
$`|01`$ $``$ $`|01`$
$`|10`$ $``$ $`|11`$
$`|11`$ $``$ $`|10,`$ (38)
and the Hadamard gate ($`H`$),
$`|0`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}(|0+|1)`$
$`|1`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}(|0|1).`$ (39)
If we use the representation,
$$|\mathrm{\hspace{0.17em}0}=\left(\begin{array}{cc}1& \\ 0& \end{array}\right),|\mathrm{\hspace{0.17em}1}=\left(\begin{array}{cc}0& \\ 1& \end{array}\right),|\mathrm{\hspace{0.17em}00}=|\mathrm{\hspace{0.17em}0}|\mathrm{\hspace{0.17em}0}=\left(\begin{array}{c}1\\ 0\\ 0\\ 0\end{array}\right),\mathrm{},$$
(40)
the basic transformations can be written as:
$$P(\varphi )=\left(\begin{array}{cc}1& 0\\ 0& e^{i\varphi }\end{array}\right),H=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right),CNOT=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 1& 0\end{array}\right).$$
(41)
Note that these gates only act on one or two qubits. The fact that they are nevertheless universal is very important as these gates can, in principle, be built and therefore even complicated calculations can be done without using gates that act on many qubits at once.
These gates can also be visualized in an elegant way. We denote the Phase Shift by
$$\text{ }\text{ }\text{}\text{ },$$
(42)
the Hadamard gate by
(43)
and the Controlled Not by
(44)
where on the left hand side of the diagram the input state is given, the output on the right, and $`x,y\{0,1\}`$. Finally we note that the Controlled Phase Shift CP($`\varphi `$) can be built out of the above states. It is denoted by
$`CP(\varphi )=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& e^{i\varphi }\end{array}\right)\text{ }\text{ }\text{}\text{ }e^{ixy\varphi }|x|y.`$ (49)
After having set up the basic framework we will present networks that prepare Alice’s initial state and project onto the basis $`|\varphi _j`$. Let us first look at the case in which the symmetric state (2) is prepared, which is projected onto (58). In this case we have the following network:
(50)
We use the convention that the original state is described by $`|\mathrm{𝟘}=|0|0\mathrm{}|0(_B)^N=`$. Then the first part of the network gives:
$$|\mathrm{𝟘}=|0|0\stackrel{H_1}{}\frac{1}{\sqrt{2}}(|00+|10)\stackrel{CN_2}{}\frac{1}{\sqrt{2}}(|00+|11).$$
(51)
Here the indices denote the qubit on which the transformation is performed. To prove that the second part projects onto the basis (58), we have to show that the set of basis vectors in (58) is mapped bijectively onto $`\{|00,|01,|10,|11\}`$. For that purpose we introduce yet another gate, the Controlled Unitary gate ($`CU`$):
$$\text{ }\text{ }\text{}CU=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& \frac{e^{i\frac{3\pi }{4}}}{\sqrt{2}}& i\frac{e^{i\frac{3\pi }{4}}}{\sqrt{2}}\\ 0& 0& i\frac{e^{i\frac{3\pi }{4}}}{\sqrt{2}}& \frac{e^{i\frac{3\pi }{4}}}{\sqrt{2}}\end{array}\right).$$
(52)
The Contolled Hadamard ($`CH`$) can be built directly.
Now we can easily check:
$`|\varphi _{1,2}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|00\pm {\displaystyle \frac{1}{2}}\left[|01e^{\frac{i\pi }{4}}+|10e^{\frac{i\pi }{4}}\right]`$
$`\stackrel{P(\pi )_1}{}`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|00\pm {\displaystyle \frac{1}{2}}\left[|01e^{\frac{i\pi }{4}}|10e^{\frac{i\pi }{4}}\right]`$
$`\stackrel{CN_1}{}`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|00\pm {\displaystyle \frac{1}{2}}\left[|11e^{\frac{i\pi }{4}}|10e^{\frac{i\pi }{4}}\right]`$
$`\stackrel{CU_2}{}`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|00\pm {\displaystyle \frac{1}{2}}[{\displaystyle \frac{1}{\sqrt{2}}}(|11+i|10)e^{\frac{i\pi }{2}}{\displaystyle \frac{1}{\sqrt{2}}}(|10+i|11)e^{i\pi }]={\displaystyle \frac{1}{\sqrt{2}}}|00\pm |10)`$
$`\stackrel{H_1}{}`$ $`\{\begin{array}{c}|00\text{for}|\varphi _1\\ |10\text{for}|\varphi _2\end{array}`$
$`|\varphi _{3,4}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|\mathrm{\hspace{0.17em}11}\pm {\displaystyle \frac{1}{2}}\left[|\mathrm{\hspace{0.17em}01}e^{\frac{i\pi }{4}}+|\mathrm{\hspace{0.17em}10}e^{\frac{i\pi }{4}}\right]`$
$`\stackrel{P(\pi )}{}`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|11\pm {\displaystyle \frac{1}{2}}\left[|01e^{\frac{i\pi }{4}}|10e^{\frac{i\pi }{4}}\right]`$
$`\stackrel{CN_1}{}`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|01\pm {\displaystyle \frac{1}{2}}\left[|11e^{\frac{i\pi }{4}}|10e^{\frac{i\pi }{4}}\right]`$
$`\stackrel{CU_2}{}`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|01\pm {\displaystyle \frac{1}{2}}[{\displaystyle \frac{1}{\sqrt{2}}}(|11+i|10)e^{i\frac{\pi }{2}}{\displaystyle \frac{1}{\sqrt{(}2)}}(|10+i|11)e^{i\pi }]={\displaystyle \frac{1}{\sqrt{2}}}|00\pm |10)`$
$`\stackrel{H_1}{}`$ $`\{\begin{array}{c}|01\text{for}|\varphi _3\\ |11\text{for}|\varphi _4\end{array}`$
If Alice prepares the singlet state (3), the gate looks slightly more complicated:
(53)
Here the Not gate denotes<sup>5</sup><sup>5</sup>5In fact a Not gate could also be built directly.:
(54)
Again it can be proved that this gate prepares the correct state $`|\psi _{AB}`$ and projects onto $`|\varphi _j`$ as given in (5). The first part gives:
$`|\mathrm{𝟘}=|0|0\stackrel{Not_2}{}|0|1\stackrel{CN_1}{}|1|1\stackrel{H_2}{}{\displaystyle \frac{1}{\sqrt{2}}}(|10|11)\stackrel{CN_1}{}{\displaystyle \frac{1}{\sqrt{2}}}(|10|01).`$ (55)
To prove that the second part projects onto (5), we first look at the first three steps in the diagram:
$`|\varphi _{1,2}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|01\pm {\displaystyle \frac{1}{2}}\left[|11e^{\frac{i\pi }{4}}|00e^{\frac{i\pi }{4}}\right]`$
$`\stackrel{CP(\pi )}{}`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|01\pm {\displaystyle \frac{1}{2}}\left[|11e^{\frac{i\pi }{4}}|00e^{\frac{i\pi }{4}}\right]`$
$`\stackrel{Not_2}{}`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|00\pm {\displaystyle \frac{1}{2}}\left[|10e^{\frac{i\pi }{4}}|01e^{\frac{i\pi }{4}}\right]`$
$`\stackrel{P(\frac{\pi }{2})_1,P(\frac{\pi }{2})_2}{}`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|00\pm {\displaystyle \frac{1}{2}}\left[|01e^{\frac{i\pi }{4}}+|10e^{\frac{i\pi }{4}}\right]`$
$`|\varphi _{3,4}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|10\pm {\displaystyle \frac{1}{2}}\left[|11e^{\frac{i\pi }{4}}|00e^{\frac{i\pi }{4}}\right]`$ (56)
$`\stackrel{CP(\pi )}{}`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|10\pm {\displaystyle \frac{1}{2}}\left[|11e^{\frac{i\pi }{4}}|00e^{\frac{i\pi }{4}}\right]`$
$`\stackrel{Not_2}{}`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|11\pm {\displaystyle \frac{1}{2}}\left[|10e^{\frac{i\pi }{4}}|01e^{\frac{i\pi }{4}}\right]`$
$`\stackrel{P(\frac{\pi }{2})_1,P(\frac{\pi }{2})_2}{}`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|11\pm {\displaystyle \frac{1}{2}}\left[|01e^{\frac{i\pi }{4}}+|10e^{\frac{i\pi }{4}}\right].`$
But this is the basis (58). Therefore, these gates transform one basis into the other. The rest of the gate is the same as in (50) and hence the gate projects onto the correct basis.
## 5 Appendix
In the appendix we give some basic formulae as well as a different basis for the $`m`$=3 orthogonal case and the basis for the $`m`$=3 non-orthogonal case.
(i), Basic Transformations:
$`|_x`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(|_z+|_z)`$
$`|_x`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(|_z|_z)`$
$`|_y`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(|_z+i|_z)`$
$`|_y`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(|_zi|_z).`$ (57)
(ii), In the state (2) was used and the following basis was obtained:
$`|\varphi _1`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|_z+{\displaystyle \frac{1}{2}}\left[\right|_ze^{\frac{i\pi }{4}}+|_ze^{\frac{i\pi }{4}}]`$
$`|\varphi _2`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|_z{\displaystyle \frac{1}{2}}\left[\right|_ze^{\frac{i\pi }{4}}+|_ze^{\frac{i\pi }{4}}]`$
$`|\varphi _3`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|_z+{\displaystyle \frac{1}{2}}\left[\right|_ze^{\frac{i\pi }{4}}+|_ze^{\frac{i\pi }{4}}]`$
$`|\varphi _4`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|_z{\displaystyle \frac{1}{2}}\left[\right|_ze^{\frac{i\pi }{4}}+|_ze^{\frac{i\pi }{4}}].`$ (58)
This is equivalent to (5) and can be obtained by a unitary transformation.
(iii), The $`m=3`$ non-orthogonal case:
In this case Alice can choose a look-up table like:
$`𝒏_1`$ $`𝒏_2`$ $`𝒏_3`$ $`\lambda _1`$ $``$ $``$ $``$ $`\lambda _2`$ $``$ $``$ $``$ $`\lambda _3`$ $``$ $``$ $``$ $`\lambda _4`$ $``$ $``$ $``$ $`\lambda _5`$ $``$ $``$ $``$ $`\lambda _6`$ $``$ $``$ $``$ $`\lambda _7`$ $``$ $``$ $``$ $`\lambda _8`$ $``$ $``$ $``$
A special choice of the orientation of axes with $`b_1^2=b_4^2=b_5^2=b_6^2=b_7^2=b_8^2=\frac{1}{8}`$ and $`b_2=0,b_3^2=\frac{1}{8}`$ leads to:
$`|\varphi _1`$ $`=`$ $`|`$
$`|\varphi _2`$ $`=`$ $`{\displaystyle \frac{1}{4}}(1+i\sqrt{3})|{\displaystyle \frac{1}{4}}(1i\sqrt{3})|+{\displaystyle \frac{1}{4}}\left(\right|+|)+{\displaystyle \frac{\sqrt{35}}{10}}|{\displaystyle \frac{\sqrt{10}}{20}}|`$
$`|\varphi _3`$ $`=`$ $`+{\displaystyle \frac{1}{4}}(1+i\sqrt{3})|+{\displaystyle \frac{1}{4}}(1i\sqrt{3})|+{\displaystyle \frac{1}{4}}\left(\right|+|)+{\displaystyle \frac{\sqrt{35}}{10}}|{\displaystyle \frac{\sqrt{10}}{20}}|`$
$`|\varphi _4`$ $`=`$ $`({\displaystyle \frac{\sqrt{2}}{4}}+i{\displaystyle \frac{\sqrt{6}}{12}})|+({\displaystyle \frac{\sqrt{2}}{4}}i{\displaystyle \frac{\sqrt{(}6)}{12}})|+({\displaystyle \frac{\sqrt{2}}{4}}{\displaystyle \frac{\sqrt{3}}{6}})|+({\displaystyle \frac{\sqrt{2}}{4}}+{\displaystyle \frac{\sqrt{3}}{6}})|+`$
$`+{\displaystyle \frac{\sqrt{7}}{7}}|{\displaystyle \frac{1}{35}}|+{\displaystyle \frac{\sqrt{5}}{10}}|`$
$`|\varphi _5`$ $`=`$ $`({\displaystyle \frac{1}{4}}i{\displaystyle \frac{\sqrt{3}}{12}})|+({\displaystyle \frac{1}{4}}+i{\displaystyle \frac{\sqrt{3}}{12}})|+({\displaystyle \frac{1}{4}}+{\displaystyle \frac{\sqrt{6}}{12}})|+({\displaystyle \frac{1}{4}}{\displaystyle \frac{\sqrt{6}}{12}})|+{\displaystyle \frac{\sqrt{10}}{4}}|`$
$`|\varphi _6`$ $`=`$ $`({\displaystyle \frac{1}{4}}i{\displaystyle \frac{\sqrt{3}}{12}})|+({\displaystyle \frac{1}{4}}+i{\displaystyle \frac{\sqrt{3}}{12}})|+({\displaystyle \frac{1}{4}}+{\displaystyle \frac{\sqrt{6}}{12}})|+({\displaystyle \frac{1}{4}}{\displaystyle \frac{\sqrt{6}}{12}})|{\displaystyle \frac{\sqrt{14}}{7}}|`$
$`{\displaystyle \frac{2}{35}}\sqrt{35}|{\displaystyle \frac{3}{20}}\sqrt{10}|`$
$`|\varphi _7`$ $`=`$ $`({\displaystyle \frac{1}{4}}i{\displaystyle \frac{\sqrt{3}}{12}})|+({\displaystyle \frac{1}{4}}+i{\displaystyle \frac{\sqrt{3}}{12}})|+({\displaystyle \frac{1}{4}}+{\displaystyle \frac{\sqrt{6}}{6}})|+({\displaystyle \frac{1}{4}}{\displaystyle \frac{\sqrt{6}}{6}})|{\displaystyle \frac{\sqrt{14}}{7}}|`$
$`{\displaystyle \frac{3}{70}}\sqrt{35}|{\displaystyle \frac{\sqrt{10}}{20}}|`$
$`|\varphi _8`$ $`=`$ $`({\displaystyle \frac{1}{4}}+i{\displaystyle \frac{\sqrt{3}}{12}})|+({\displaystyle \frac{1}{4}}i{\displaystyle \frac{\sqrt{3}}{12}})|+({\displaystyle \frac{1}{4}}+{\displaystyle \frac{\sqrt{6}}{6}})|+({\displaystyle \frac{1}{4}}{\displaystyle \frac{\sqrt{6}}{6}})|{\displaystyle \frac{\sqrt{14}}{7}}|`$
$`{\displaystyle \frac{3}{70}}\sqrt{35}|{\displaystyle \frac{\sqrt{10}}{20}}|,`$
all expressed in the z-basis. As in (36), all parameters of the unitary transformation were chosen to be zero.
## 6 Acknowledgements
I would like to thank Prof. Artur Ekert and Daniel Oi for their assistance and support as well as for very many interesting discussions. Furthermore, I would like to thank the authors of for permission to use the diagrams in (42), (43), (44) and (49).
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# Approximate expressions for polar gap electric field of pulsars
## 1 Introduction
Accurate models for $`E_{}`$ — the electric field component parallel to the local magnetic field $`B`$ above polar caps of rotating magnetized neutron stars — are of great theoretical interest in the context of magnetospheric activity of pulsars. In the framework of space charge limited flow model (originally introduced by Arons & Scharlemann (as (1979))), Harding & Muslimov (hm (1998)) (hereafter HM98) derived $`E_{}`$ including the general relativistic frame dragging effect, worked out by Muslimov & Tsygan (mt (1992)) (hereafter MT92). HM98 considered the case with $`E_{}`$ screened at both a lower and an upper boundary of acceleration region i.e. at the star surface and at a pair formation front. Since the full expression for $`E_{}`$ is too cumbersome for practical applications, HM98 offered simple analytic expressions valid in various limiting cases. In this paper we revise and extend their results. In particular, we show that in the most frequently considered case, i.e. when the gap accelerator length is comparable to the size of the polar cap radius, the approximation derived by HM98 overestimates $`E_{}`$ by a ratio of the neutron star radius to the polar cap radius.
## 2 Formulae for $`E_{}`$
In the following, $`h`$ denotes the altitude above the neutron star surface and $`h_c`$ is the gap height above which a pair formation front screens $`E_{}`$ (i.e. $`E_{}=0`$ for $`hh_c`$). Altitudes $`zh/R_{\mathrm{ns}}`$, $`z_ch_c/R_{\mathrm{ns}}`$ and radial distances $`\eta 1+z`$, $`\eta _c1+z_c`$ scaled with the star radius $`R_{\mathrm{ns}}`$ will also be used. The magnetic colatitude $`\xi \theta /\theta (\eta )`$ is scaled with the half-opening angle of the polar magnetic flux tube $`\theta (\eta )`$.
A solution of Poisson’s equation for the polar gap with $`h_cR_{\mathrm{ns}}`$ as derived by HM98 reads
$`E_{}^{(1)}E_0\theta _0^3(1ϵ)^{1/2}\{{\displaystyle \frac{3}{2}}\kappa S_1\mathrm{cos}\chi +`$ (1)
$`\text{ }+{\displaystyle \frac{3}{8}}\theta _0H(1)\delta (1)S_2\mathrm{sin}\chi \mathrm{cos}\varphi \}`$
where,
$$\begin{array}{c}S_1=\underset{i=1}{\overset{\mathrm{}}{}}\frac{8J_0(k_i\xi )}{k_i^4J_1(k_i)}(z,z_c,\gamma _i),\hfill \\ \\ S_2=\underset{i=1}{\overset{\mathrm{}}{}}\frac{16J_1(\stackrel{~}{k}_i\xi )}{\stackrel{~}{k}_i^4J_2(\stackrel{~}{k}_i)}(z,z_c,\stackrel{~}{\gamma }_i),\hfill \end{array}$$
(2)
$`(z,z_c,\gamma )=[a_1(\gamma \eta 1)e^{\gamma z}+a_2(\gamma \eta +1)e^{\gamma z}+`$ (3)
$`\text{ }+a_1(1\gamma )a_2(1+\gamma )]/(a_1+a_2),`$
$$a_1=(\gamma \eta _c+1)e^{\gamma z_c}\gamma 1,a_2=\gamma 1(\gamma \eta _c1)e^{\gamma z_c}$$
(4)
and
$$\gamma _i\frac{k_i}{\theta _0(1ϵ)^{1/2}},\mathrm{and}\stackrel{~}{\gamma }_i\frac{\stackrel{~}{k}_i}{\theta _0(1ϵ)^{1/2}},$$
(5)
where $`k_i`$ and $`\stackrel{~}{k}_i`$ are the positive roots of the Bessel functions $`J_0`$ and $`J_1`$, respectively, $`H(1)\delta (1)1`$, and $`\chi `$ is a tilt angle between the rotation and the magnetic dipole axes. Values of $`E_0`$, $`\theta _0`$, $`ϵ`$, and $`\kappa `$ are given by:
$$\begin{array}{c}E_0B_{\mathrm{pc}}\frac{\mathrm{\Omega }R_{\mathrm{ns}}}{c},\text{ }\theta _0=\left(\frac{\mathrm{\Omega }R_{\mathrm{ns}}}{cf(1)}\right)^{1/2},\hfill \\ \\ ϵ\frac{2GM}{R_{\mathrm{ns}}c^2},\text{ }\kappa \frac{ϵI}{MR_{\mathrm{ns}}^2},\hfill \end{array}$$
(6)
where $`B_{\mathrm{pc}}`$ is the magnetic field strength at the magnetic pole, $`\mathrm{\Omega }`$ is the pulsar rotation frequency, and $`I`$, $`M`$ are the moment of inertia, and mass of the neutron star, respectively. The function $`f(\eta )`$ reads
$$f(\eta )=3x^3\left[\mathrm{ln}(1x^1)+(1+(2x)^1)/x\right],$$
(7)
with $`x=\eta /ϵ`$.
The electric field $`E_{}^{(1)}`$ as a function of height $`h`$ is drawn in Fig. 1 for several values of gap height $`h_c`$. It can be seen that $`E_{}`$ saturates in a twofold way: First, for a fixed $`h_c`$ it assumes a constant value at $`h`$ exceeding the polar cap radius $`r_{\mathrm{pc}}`$ (after a linear increase with $`h`$). Second, for a fixed $`hh_c`$ it initially increases linearly with $`h_c`$ to become constant for $`h_c>r_{\mathrm{pc}}`$. As noted by Harding & Muslimov this behaviour can be reproduced with simple approximations of Eq. (1) for different regimes of validity.
Since $`k_i2.4`$ it follows form Eq. (5) that $`\gamma _iz`$ becomes smaller than 1 for $`hr_{\mathrm{pc}}/3`$. Thus, in the limit $`hr_{\mathrm{pc}}/3`$ and $`h_cr_{\mathrm{pc}}/3`$ (or $`\gamma _iz1`$ and $`\gamma _iz_c1`$) the function $`(z,z_c,\gamma _i)`$ may be approximated with
$$(z,z_c,\gamma _i)\frac{1}{2}\gamma _i^3\left(1\frac{z}{z_c}\right)zz_c$$
(8)
which leads to
$`E_{}^{(2)}3{\displaystyle \frac{\mathrm{\Omega }R_{\mathrm{ns}}}{c}}{\displaystyle \frac{B_{\mathrm{pc}}}{1ϵ}}(1{\displaystyle \frac{z}{z_c}})zz_c[\kappa \mathrm{cos}\chi +`$ (9)
$`\text{ }+{\displaystyle \frac{1}{2}}\theta _0\xi H(1)\delta (1)\mathrm{sin}\chi \mathrm{cos}\varphi ],`$
(HM98), where the relations
$$\underset{i=1}{\overset{\mathrm{}}{}}\frac{2J_0(k_i\xi )}{k_iJ_1(k_i)}=1,\mathrm{and}\underset{i=1}{\overset{\mathrm{}}{}}\frac{2J_1(\stackrel{~}{k}_i\xi )}{\stackrel{~}{k}_iJ_2(\stackrel{~}{k}_i)}=\xi $$
(10)
have been used (see eg. MT92). As can be seen in Fig 1, in practice, Eq. (9) (open circles) reproduces Eq. (1) (thick solid line) for $`h_cr_{\mathrm{pc}}/3`$ (and $`hh_c`$).
Approximate behaviour of $`E_{}`$ for $`h_cr_{\mathrm{pc}}`$ can be determined by considering the opposite limiting case for the accelerator height: $`h_cr_{\mathrm{pc}}/3`$. Assuming $`\gamma _iz_c1`$, $`h<(h_cr_{\mathrm{pc}})`$, (and $`h_cR_{\mathrm{ns}}`$) in Eq. (1) one obtains
$$(z,z_c,\gamma _i)\gamma _i\left(1e^{\gamma _iz}\right),$$
(11)
which results in the same formula for $`E_{}`$ as the one derived by MT92 for the case with no upper gap boundary (with the condition $`E_{}`$ = 0 fulfilled only at the star surface):
$`E_{}^{(3)}3E_0\theta _0^2\{\kappa \mathrm{cos}\chi {\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{4J_0(k_i\xi )}{k_i^3J_1(k_i)}}[1e^{\gamma _iz}]+`$ (12)
$`+\theta _0H(1)\delta (1)\mathrm{sin}\chi \mathrm{cos}\varphi {\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{2J_1(\stackrel{~}{k}_i\xi )}{\stackrel{~}{k}_i^3J_2(\stackrel{~}{k}_i)}}[1e^{\stackrel{~}{\gamma }_iz}]\}`$
Nevertheless, for $`r_{\mathrm{pc}}h_cR_{\mathrm{ns}}`$, Eq. (12) reproduces $`E_{}^{(1)}`$ (Eq. 1) almost over the entire acceleration length (see the thin solid line in Fig. 1) except its very final part, where $`h>(h_cr_{\mathrm{pc}})`$.
Taylor expansion of (11) reveals the linear increase of $`E_{}`$ with $`h`$ for $`h0.1r_{\mathrm{pc}}`$ and the saturation above $`hr_{\mathrm{pc}}`$:
$`E_{}^{(4)}3{\displaystyle \frac{\mathrm{\Omega }R_{\mathrm{ns}}}{c}}{\displaystyle \frac{B_{\mathrm{pc}}}{(1ϵ)^{1/2}}}\theta _0z[\kappa f_1(\xi )\mathrm{cos}\chi +`$ (13)
$`\text{ }+{\displaystyle \frac{1}{4}}\theta _0f_2(\xi )H(1)\delta (1)\mathrm{sin}\chi \mathrm{cos}\varphi ],`$
for $`h0.1r_{\mathrm{pc}}`$ and $`0.5r_{\mathrm{pc}}h_cR_{\mathrm{ns}}`$, and
$`E_{}^{(5)}{\displaystyle \frac{3}{2}}{\displaystyle \frac{\mathrm{\Omega }R_{\mathrm{ns}}}{c}}B_{\mathrm{pc}}\theta _0^2(1\xi ^2)[\kappa \mathrm{cos}\chi +`$ (14)
$`\text{ }+{\displaystyle \frac{1}{4}}\theta _0\xi H(1)\delta (1)\mathrm{sin}\chi \mathrm{cos}\varphi ],`$
for $`r_{\mathrm{pc}}h(h_cr_{\mathrm{pc}})`$ and $`2r_{\mathrm{pc}}h_cR_{\mathrm{ns}}`$. The magnetic colatitude profiles
$$f_1(\xi )=\underset{i=1}{\overset{\mathrm{}}{}}\frac{4J_0(k_i\xi )}{k_i^2J_1(k_i)},\mathrm{and}f_2(\xi )=\underset{i=1}{\overset{\mathrm{}}{}}\frac{8J_1(\stackrel{~}{k}_i\xi )}{\stackrel{~}{k}_i^2J_2(\stackrel{~}{k}_i)}$$
(15)
are shown in Fig. 2 and the relations
$$\underset{i=1}{\overset{\mathrm{}}{}}\frac{8J_0(k_i\xi )}{k_i^3J_1(k_i)}=1\xi ^2,\underset{i=1}{\overset{\mathrm{}}{}}\frac{16J_1(\stackrel{~}{k}_i\xi )}{\stackrel{~}{k}_i^3J_2(\stackrel{~}{k}_i)}=\xi (1\xi ^2)$$
(16)
have been used in derivation of (14) (eg. MT92). Approximation (13) is shown in Fig. 1 as a dot-dashed line.
For $`h_cR_{\mathrm{ns}}`$ the symmetry
$$E_{}^{(1)}(h)E_{}^{(1)}(h_ch)$$
(17)
may easily be proven to hold. Therefore, to extend the validity of Eq. (13) to $`hh_c`$ it is useful to construct the approximation
$$E_{}^{(6)}=E_{}^{(4)}\left(1\frac{z}{z_c}\right)$$
(18)
which works reasonably well for $`0.5r_{\mathrm{pc}}h_cr_{\mathrm{pc}}`$ (dashed lines in Fig. 1).
The expression derived by HM98 for the case $`h_cr_{\mathrm{pc}}`$ clearly overestimates $`E_{}`$ by a factor $`\theta _0^1(r_{\mathrm{pc}}/R_{\mathrm{ns}})^1`$ (cf. eq. (A3) in HM98, also eq. (1) in Harding & Muslimov (1998a ), or a formula used by Zhang, Harding & Muslimov (zhm (2000))). Moreover, unlike in eq. (A3) of HM98, $`E_{}^{(4)}`$ depends on the magnetic colatitude $`\xi `$ via $`f_1(\xi )`$. Our approximations are correct because they reproduce the “exact” formula for $`E_{}^{(1)}`$ (Eq. 1). Moreover, $`E_{}^{(4)}`$ of Eq. (13) is identical to the formula of MT92 for the case with no screening at $`h_c`$ (taken in the limit $`\gamma _iz1`$). This identity should hold since for $`r_{\mathrm{pc}}<h_cR_{\mathrm{ns}}`$ both gap boundaries influence $`E_{}`$ only within corresponding adjacent regions of height $`r_{\mathrm{pc}}`$.
Since Eq. (17) holds for $`h_cR_{\mathrm{ns}}`$, $`E_{}`$ in the upper part of accelerator ($`r_{\mathrm{pc}}<hh_c`$) may be approximated with
$$E_{}^{(7)}(h)E_{}^{(3)}(h_ch)\{\begin{array}{c}E_{}^{(4)}(h_ch),\hfill \\ \mathrm{for}(h_cr_{\mathrm{pc}})hh_c\hfill \\ \\ E_{}^{(5)},\hfill \\ \mathrm{for}r_{\mathrm{pc}}<h(h_cr_{\mathrm{pc}}).\hfill \end{array}$$
(19)
In the limit where $`h_cr_{\mathrm{pc}}`$ and $`hr_{\mathrm{pc}}`$, a solution of Poisson’s equation for the elongated polar gap reads
$`E_{}^{(8)}E_0\theta _0^2\{{\displaystyle \frac{3}{2}}{\displaystyle \frac{\kappa }{\eta ^4}}\mathrm{cos}\chi [(1\xi ^2)+`$ (20)
$`\left({\displaystyle \frac{\eta _c}{\eta }}\right)^4{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{8J_0(k_i\xi )}{k_i^3J_1(k_i)}}e^{\gamma _i(\eta _c)(\eta _c\eta )}]+`$
$`+{\displaystyle \frac{3}{8}}g(\eta )\mathrm{sin}\chi \mathrm{cos}\varphi [\xi (1\xi ^2)+`$
$`{\displaystyle \frac{\eta _c}{\eta }}{\displaystyle \frac{g(\eta _c)}{g(\eta )}}{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{16J_1(\stackrel{~}{k}_i\xi )}{\stackrel{~}{k}_i^3J_2(\stackrel{~}{k}_i)}}e^{\stackrel{~}{\gamma }_i(\eta _c)(\eta _c\eta )}]\}`$
(MT92), where now
$$\begin{array}{c}\gamma _i\frac{k_i}{\theta (\eta _c)\eta _c(1ϵ/\eta _c)^{1/2}},\hfill \\ \\ \stackrel{~}{\gamma }_i\frac{\stackrel{~}{k}_i}{\theta (\eta _c)\eta _c(1ϵ/\eta _c)^{1/2}},\hfill \end{array}$$
(21)
$`g(\eta )=\theta (\eta )\delta (\eta )H(\eta )`$, $`\theta (\eta )=\theta _0(\eta f(1)/f(\eta ))^{1/2}`$, and the functions $`\delta `$, and $`H`$ are defined in MT92 or HM98. Since $`\gamma _i1`$, the two summation terms which reflect the screening effect of the pair formation front at $`h_c`$ contribute to $`E_{}`$ only at $`h`$ very close to $`h_c`$ (see Fig. 3) so that $`\gamma _i(\eta _c)(\eta _c\eta )1`$. Thus, for $`h_cr_{\mathrm{pc}}`$ and $`r_{\mathrm{pc}}hh_c`$ one may simply use a formula derived in MT92 for the case with no screening at $`h_c`$:
$`E_{}^{(9)}E_0\theta _0^2\{{\displaystyle \frac{3}{2}}{\displaystyle \frac{\kappa }{\eta ^4}}(1\xi ^2)\mathrm{cos}\chi +`$ (22)
$`\text{ }+{\displaystyle \frac{3}{8}}g(\eta )\xi (1\xi ^2)\mathrm{sin}\chi \mathrm{cos}\varphi \}.`$
For $`P=0.1`$ s the approximation $`E_{}^{(9)}`$ matches the approximation $`E_{}^{(3)}`$ (which is valid for low altitudes (Eq. 12)) at $`h0.6r_{\mathrm{pc}}`$. This altitude is much lower than estimated by Muslimov & Harding (mh (1997)). Consequently, below $`0.6r_{\mathrm{pc}}`$ either $`E_{}^{(3)}`$ or $`E_{}^{(4)}`$ should be used.
## 3 Conclusions
We find that for accelerator height $`h_c`$ approaching the polar cap radius $`r_{\mathrm{pc}}`$ the accelerating electric field (for nonorthogonal rotators) may be approximated according to Eq. (18):
$$E_{}=1.46\frac{B_{12}}{P^{3/2}}h\left(1\frac{h}{h_c}\right)f_1(\xi )\mathrm{cos}\chi \mathrm{Gauss},$$
(23)
where $`B_{12}=B_{\mathrm{pc}}/(10^{12}G)`$, $`P`$ is the rotation period in seconds, $`h`$ is in cm, for a star with $`M=1.4M_{}`$, $`R_{\mathrm{ns}}=10^6`$ cm. This result differs from the approximation derived in HM98 by a considerable period-dependent factor $`r_{\mathrm{pc}}/R_{\mathrm{ns}}`$. Moreover, unlike the HM98’ formula it depends on the magnetic colatitude $`\xi `$ through the factor $`f_1(\xi )`$. It may be of particular importance for a possibility of limiting acceleration by the resonant inverse compton scattering (see Dyks & Rudak dr (2000)).
We emphasize that the corrections presented in this research note do not affect the results of the numerical calculations presented by Harding & Muslimov (hm (1998)) because their numerical procedures include the exact expressions given by Eqs. (1) and (20). However, their analytic estimates of the height and width of the acceleration zone and of the maximum particle energy need to be revised.
Since for an elongated polar gap ($`h_cr_{\mathrm{pc}}`$) the upper gap boundary (pair formation front) influences the electric field only within the negligible ($`h_c`$) upper part of accelerator, the formulae derived by Muslimov and Tsygan (mt (1992)) for no upper boundary at $`h_c`$ may be used within the entire gap with no significant overestimates of electron energies.
###### Acknowledgements.
This work was supported by KBN grant 2P03D 02117. We appreciate informative contacts with Alice Harding and Bing Zhang.
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# STATISTICS OF SMOOTHED COSMIC FIELDS IN PERTURBATION THEORY I: Formulation and useful formulas in second-order perturbation theory
## 1 INTRODUCTION
Any conceivable theory of structure formation in the universe predicts statistical properties of observable quantities. Therefore, the analysis of the present inhomogeneity of the universe is inevitably statistical. However, what kind of statistics is useful is not obvious, since we can invent infinite number of statistics to be analyzed. Whether we can adopt better statistical descriptions is of great importance in this sense. Each statistic has both advantages and disadvantages. The power spectrum, for example, can fully characterize the random Gaussian fields, while it does not contain any information on the non-Gaussianity, which contains significant information in the gravitational instability theory. Among various statistical quantities, there is a promising class of statistics which utilizes smoothed cosmic fields. The smoothed field has less noisy property than the actual (unsmoothed, or raw) cosmic fields, such as galaxy distributions, temperature fluctuations of cosmic microwave background (CMB), shear fields of the gravitational lensing, and so forth. Perhaps the simplest example of such statistics is the variance $`\delta _{R}^{}{}_{}{}^{2}`$ of smoothed density contrast $`\delta _R`$, which is a function of smoothing length $`R`$. Similarly, higher-order cumulants $`\delta _{R}^{}{}_{}{}^{N}_\mathrm{c}`$ ($`N=3,4,\mathrm{}`$) are also simple statistics. The density probability distribution function (PDF) $`P(\delta _R)`$, which in principle can be constructed from the hierarchy of cumulants (e.g., Balian & Schaeffer, 1989), is an another example of popular statistics of smoothed cosmic fields.
Rather recently, more complex statistics of smoothed cosmic fields have become popular in cosmology, such as the genus statistic (Gott, Melott & Dickinson, 1986), density peak statistics (Bardeen et al., 1986), area, length, level-crossing statistics (Ryden, 1988a), Minkowski functionals (Schmalzing & Buchert, 1997), etc. These statistics provide assuring characterizations of the clustering pattern that can not be perceived only by the hierarchy of cumulants or by the PDF. The genus statistic is a powerful measure of the morphology in the 3-dimensional (Gott et al., 1989; Moore et al., 1992; Park, Gott & da Costa, 1992; Beaky, Scherrer & Villumsen, 1992; Rhoads, Gott & Postman, 1994; Vogeley et al., 1994; Protogeros & Weinberg, 1997; Sahni, Sathyaprakash & Shandarin, 1997; Canavezes et al., 1998; Springel et al., 1998; Colley et al., 2000), and 2-dimensional (Melott et al., 1989; Gott, Mao, Park & Lahav, 1992; Coles & Plionis, 1991; Plionis, Valdarnini & Coles, 1992; Davies & Coles, 1993; Coles et al., 1993; Colley, 1997) clustering of galaxies and clusters, and also of the pattern of temperature fluctuations of CMB radiation (Bond & Efstathiou, 1987; Torres, 1994; Smoot et al., 1994; Torres, Cayon, Martinez-Gonzalez & Sanz, 1995; Park et al., 1998). The peak statistics of the 3-dimensional density field are frequently used in connection with the statistics of the collapsing object (Kaiser, 1984; Mann, Heavens & Peacock, 1993; Croft & Efstathiou, 1994; Watanabe, Matsubara & Suto, 1994; Cen, 1998; Gabrielli, Labini & Durrer, 2000), while the peak statistics of the CMB fluctuations (Sazhin, 1985; Bond & Efstathiou, 1987; Coles & Barrow, 1987; Vittorio & Juszkiewicz, 1987; Kogut et al., 1995; Cayon & Smoot, 1995; Fabbri & Torres, 1996; Heavens & Sheth, 1999) and of the weak lensing fields (van Waerbeke, 2000; Jain & Van Waerbeke, 2000) are suitable for constraining cosmological models. The area, length, and level-crossing statistics directly quantify the amount of contour surfaces (Ryden, 1988b; Ryden et al., 1989; Torres, 1994). The Minkowski functionals (Mecke, Buchert & Wagner, 1994; Sahni, Sathyaprakash, & Shandarin, 1998; Sathyaprakash, Sahni, & Shandarin, 1998; Kerscher, Schmalzing, Buchert & Wagner, 1998; Bharadwaj et al., 2000), which are closely related to the above statistics, are also applied to smoothed cosmic fields (Winitzki & Kosowsky, 1997; Naselsky & Novikov, 1998; Schmalzing & Gorski, 1998; Schmalzing et al., 1999; Schmalzing & Diaferio, 2000). These statistics of smoothed density fields are considered as powerful descriptors of the statistical information of the universe. It is therefore essential to establish theoretical predictions of the behavior of such statistics so that we can effectively and ideally analyze the data of our universe.
Recent developments of the perturbation theory (for review, see Bernardeau, Colombi, Gaztanaga, & Scoccimarro, 2002) in calculating the variance, the cumulants, and the PDF are remarkable. The perturbation theory becomes more and more useful because the recent developments of observations enable us to have widely covered sample volume of the universe, which can minimize undesirable strongly nonlinear effects which we do not understand well. The direct comparison of the theoretical predictions of the perturbation theory with the actual data is a promising field of research in this sense. In the case of top-hat smoothing function, Juszkiewicz, Bouchet & Colombi (1993) and Bernardeau (1994a) used the tree-level (i.e., lowest-order) perturbation theory to obtain the third- and forth-order cumulants, i.e., the skewness and the kurtosis. The same quantities with Gaussian smoothing are calculated by Goroff et al. (1986), Matsubara (1994) and Łokas et al. (1995). Quite cleverly, Bernardeau (1994b) took advantage of special properties of the top-hat smoothing function, and succeeded to obtain full hierarchy of higher order moments in the tree-level perturbation theory. He also obtained the PDF from this hierarchy of cumulants, which remarkably describes the nonlinear behavior of the gravitational instability in numerical simulations. Beyond the tree-level calculation, the perturbation theory with loop-corrections has been also developed (Juszkiewicz, 1981; Vishniac, 1983; Juszkiewicz, Sonoda & Barrow, 1984; Coles, 1990; Suto & Sasaki, 1991; Makino, Sasaki & Suto, 1992; Jain & Bertschinger, 1994; Baugh & Efstathiou, 1994; Scoccimarro & Frieman, 1996a, b). Because of the simplicity of the statistics, calculating the variance, the cumulants and the PDF were primarily the playground of perturbation-theorists. The evaluation of other statistics by the perturbation theory is less trivial. A quite useful technique is the Edgeworth expansion, which was first applied to cosmology in calculating the PDF from seeds perturbations by Scherrer & Bertschinger (1991), and from non-linear perturbation theory by Juszkiewicz et al. (1995), and Bernardeau & Kofman (1995). The analytic expression of the genus statistics in the perturbation theory is derived by Matsubara (1994), whose technique corresponds to multivariate version of the Edgeworth expansion. In some literatures the Edgeworth expansion is used in connecting the statistics and dynamics of the universe (Chodorowski & Lokas, 1997; Łokas, 1998; Taruya & Soda, 1999). The purpose of this paper is to give a comprehensive description of the formalism by which the perturbative evaluation is possible for wide range of non-trivial statistics of smoothed cosmic fields in general.
Since the number of spatial dimensions of our universe is three, the statistics of large-scale cosmic fields are defined in either one-, two-, or three-dimensional space. For example, the space of the density field in a redshift map of galaxies or quasars is three-dimensional. A projected galaxy map on the sky, a shear field of gravitational lensing, and temperature fluctuations of CMB on the sky are fields in two-dimensional space. The absorption lines of quasar spectra and pencil-beam surveys of galaxies define fields in one-dimensional space. Thus it is useful to develop our statistical method in $`d`$-dimensions for generality. As illustrative examples of applications of our method, we calculate level-crossing statistic (or equivalently length and area statistics), genus statistics, density-extrema statistics and Minkowski functionals. As cosmic fields, we consider density and velocity fields in three dimensions, the projected density field of galaxies in two dimensions. The basic formalism and results of the second order theory are presented in this paper. We will give results of the third order theory in a future paper, which are technically more involved.
The primary purpose of the present paper is to fully describe the basic formalism for the statistics of smoothed cosmic fields in perturbation theory. Applications of the second order theory to popular statistics and cosmic fields are systematically presented. Thus, this paper is in some way a mixture of the new results and comprehensive review of the old results. The major new results in this paper are:
* Explicit formulas of the lowest non-Gaussian correction in various smoothed cosmic fields in $`d`$-dimensions, including all the Minkowski functionals. They are presented as functions of both density threshold and rescaled threshold by volume-fraction.
* Relations among several statistics of the smoothed field are clarified.
* Derivatives of skewness parameter of the velocity field and projected 2D field are derived in second order perturbation theory of gravitational instability theory. Especially, skewness parameters with Gaussian smoothing are detailed. These quantities are particularly important to perturbatively evaluate the statistics of smoothed field.
* Genus curve against the scaled threshold in perturbation theory is discussed and a theory-motivated new asymmetry parameter in genus curve is introduced.
This paper is structured in the following way. §2 and §3 are carrying out a mathematical exercise of expressing high-order terms in terms of skewness parameters; the physics of gravity only enters once one actually calculates the skewness parameters in §4. Thus §2 and §3 are a derivation of an extension of the Edgeworth expansion, and the basic results are equations (2.22) and (2.85). In §3, popular statistics of smoothed fields are examined. The second-order expressions in terms of the skewness parameters are given for PDF, level-crossing statistics, 2-dimensional and 3-dimensional genus statistics, 2-dimensional weighted extrema, and Minkowski functionals in §3.1–§3.6. These quantities are re-expressed as functions of the volume-fraction threshold in §3.7. The skewness parameters for several cosmic fields are calculated in §4, applying the second order perturbation theory. Detailed calculations of the simple hierarchical model, the 3-dimensional density field, the velocity field, the 2-dimensional projected density field, and the weak lensing field are given in §4.1–§4.5. The effect of biasing on the skewness parameters are discussed in §4.6. We discuss implications of our second-order results in §5. We introduce a theory-motivated asymmetry factor to characterize the shift of the genus curve in this section. The conclusions are given in §6. Useful Gaussian integrals including Hermite polynomials are given in Appendix A. The bispectrum in a two-dimensional projected field is reviewed in Appendix B. The symbols in this paper are summarized in Appendix C.
## 2 PERTURBATIVE EXPANSION OF STATISTICS
### 2.1 Smoothed fields
We consider a cosmic random field $`f(𝒙)`$ which represents any field constructed from observable quantities of the universe, such as three-dimensional density field, velocity field, or two-dimensional projected density field, shear or convergence field of weak lensing, temperature fluctuations of CMB, etc. The coordinates $`𝒙`$ can be either three-, two-, or one-dimensions.
We assume the field $`f`$ is already smoothed by a smoothing function $`W_R`$ with smoothing length $`R`$ which cuts the high frequency fluctuations which suffer strongly nonlinear effects:
$`f(𝒙)={\displaystyle d^dx^{}W_R(|𝒙𝒙^{}|)f_{\mathrm{raw}}(𝒙^{})},`$ (2.1)
where $`d=1,2,3`$ is the dimension of the space $`𝒙`$, and $`f_{\mathrm{raw}}`$ is a raw, unsmoothed field. Two of the most popular 3D smoothing functions are the tophat smoothing function
$`W_R(x)={\displaystyle \frac{3}{4\pi R^3}}\mathrm{\Theta }(Rx),`$ (2.2)
and the Gaussian smoothing function
$`W_R(x)={\displaystyle \frac{1}{(2\pi )^{3/2}R^3}}\mathrm{exp}\left({\displaystyle \frac{x^2}{2R^2}}\right).`$ (2.3)
We also assume that the mean value of the field $`f`$ is zero, and that the existence the variance $`\sigma _{0}^{}{}_{}{}^{2}`$:
$`f=0,f^2=\sigma _{0}^{}{}_{}{}^{2}.`$ (2.4)
It is convenient to introduce a normalized field $`\alpha `$ which has a unit variance as follows:
$`\alpha {\displaystyle \frac{f}{\sigma _0}},\alpha ^2=1.`$ (2.5)
### 2.2 Expressing the Non-Gaussian Statistics by Gaussian Integration
The statistics of a smoothed cosmic field we are interested in are the functions of the field $`\alpha `$ and its spatial derivatives as we will see in the following sections. We denote the series of spatial derivatives by a set of variables $`(A_\mu )`$ which is defined by, for example, in 3D case,
$`(A_\mu )=(\alpha ,_1\alpha ,_2\alpha ,_3\alpha ,_{1}^{}{}_{}{}^{2}\alpha ,_{2}^{}{}_{}{}^{2}\alpha ,_{3}^{}{}_{}{}^{2}\alpha ,_1_2\alpha ,_1_3\alpha ,_2_3\alpha ,\mathrm{}),`$ (2.6)
where $`_i=/x_i`$. For convenience, the index is denoted as $`\mu =0`$, 1, 2, 3, (11), (22), (33), (12), (13), (23), $`\mathrm{}`$ in such cases. In this example of equation (2.6), only the value of the field and of the derivatives on a single point is considered, but in general, more than two points can be considered. The set $`A_\mu `$ forms multivariate random fields, which is denoted as an $`N`$-dimensional vector $`𝑨`$ in the following. The dimension $`N`$ is the total number of derivatives which appear in the definition of the statistics we are interested in. The statistical information is described by the multivariate PDF, $`P(𝑨)`$. The Fourier transform of the PDF is the partition function:
$`Z(𝑱)={\displaystyle _{\mathrm{}}^{\mathrm{}}}d^NAP(𝑨)\mathrm{exp}(i𝑱𝑨).`$ (2.7)
At this point, the cumulant expansion theorem (e.g., Ma, 1985) is very useful. This theorem states that $`\mathrm{ln}Z`$ is the generating function of the cumulants, $`M_{\mu _1\mathrm{}\mu _n}^{(n)}A_{\mu _1}\mathrm{}A_{\mu _n}_\mathrm{c}`$:
$`\mathrm{ln}Z(𝑱)={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{i^n}{n!}}{\displaystyle \underset{\mu _1=1}{\overset{N}{}}}\mathrm{}{\displaystyle \underset{\mu _n=1}{\overset{N}{}}}M_{\mu _1\mathrm{}\mu _n}^{(n)}J_{\mu _1}\mathrm{}J_{\mu _n}.`$ (2.8)
It follows from $`f=0`$ that $`A_\mu =0`$, and first several cumulants are given by
$`M_\mu ^{(1)}=0,`$ (2.9)
$`M_{\mu _1\mu _2}^{(2)}=A_{\mu _1}A_{\mu _2},`$ (2.10)
$`M_{\mu _1\mu _2\mu _3}^{(3)}=A_{\mu _1}A_{\mu _2}A_{\mu _3},`$ (2.11)
$`M_{\mu _1\mu _2\mu _3\mu _4}^{(4)}=A_{\mu _1}A_{\mu _2}A_{\mu _3}A_{\mu _4}A_{\mu _1}A_{\mu _2}A_{\mu _3}A_{\mu _4}A_{\mu _1}A_{\mu _3}A_{\mu _2}A_{\mu _4}A_{\mu _1}A_{\mu _4}A_{\mu _2}A_{\mu _3},`$ (2.12)
and so forth. From equations (2.8) and (2.9), the partition function is given by
$`Z(𝑱)=\mathrm{exp}\left({\displaystyle \frac{1}{2}}𝑱^\mathrm{T}𝑴𝑱\right)\mathrm{exp}\left({\displaystyle \underset{n=3}{\overset{\mathrm{}}{}}}{\displaystyle \frac{i^n}{n!}}{\displaystyle \underset{\mu _1,\mathrm{}\mu _n}{}}M_{\mu _1\mathrm{}\mu _n}^{(n)}J_{\mu _1}\mathrm{}J_{\mu _n}\right),`$ (2.13)
where $`𝑴`$ is an $`N\times N`$ matrix whose components are given by $`M_{\mu \nu }^{(2)}`$. On the other hand, the equation (2.7) is inverted as
$`P(𝑨)={\displaystyle \frac{1}{(2\pi )^N}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}d^NJZ(𝑱)\mathrm{exp}\left(i𝑱𝑨\right).`$ (2.14)
Substituting $`J_\mu i/A_\mu `$ in the last term of equation (2.13), the distribution function of equation (2.14) can be transformed in a form
$`P(𝑨)=\mathrm{exp}\left({\displaystyle \underset{n=3}{\overset{\mathrm{}}{}}}{\displaystyle \frac{()^n}{n!}}{\displaystyle \underset{\mu _1,\mathrm{}\mu _n}{}}M_{\mu _1\mathrm{}\mu _n}^{(n)}{\displaystyle \frac{^n}{A_{\mu _1}\mathrm{}A_{\mu _n}}}\right)P_\mathrm{G}(𝑨),`$ (2.15)
where
$`P_\mathrm{G}(𝑨)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^N}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}d^NJ\mathrm{exp}\left(i𝑱𝑨{\displaystyle \frac{1}{2}}𝑱^\mathrm{T}𝑴𝑱\right)`$ (2.16)
$`=`$ $`{\displaystyle \frac{1}{(2\pi )^{N/2}\sqrt{det𝑴}}}\mathrm{exp}\left({\displaystyle \frac{1}{2}}𝑨^\mathrm{T}𝑴^1𝑨\right),`$ (2.17)
is the multivariate Gaussian distribution function characterized by the correlation matrix $`𝑴`$.
Any statistical quantity of a smoothed cosmic field is expressed by an average $`F`$ of a certain function $`F(𝑨)`$ as we will see in the following sections. Thus, from equation (2.15), we obtain
$`F`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}d^NAP(𝑨)F(𝑨)`$ (2.18)
$`=`$ $`\mathrm{exp}\left({\displaystyle \underset{n=3}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}{\displaystyle \underset{\mu _1,\mathrm{}\mu _n}{}}M_{\mu _1\mathrm{}\mu _n}^{(n)}{\displaystyle \frac{^n}{A_{\mu _1}\mathrm{}A_{\mu _n}}}\right)F(𝑨)_\mathrm{G},`$ (2.19)
where
$`\mathrm{}_\mathrm{G}{\displaystyle _{\mathrm{}}^{\mathrm{}}}d^NAP_\mathrm{G}(𝑨)F(𝑨)`$ (2.20)
denotes the averaging by the Gaussian distribution function of equation (2.17).
This form, equation (2.19), is useful when the deviation from the Gaussian distribution is not large. In principle, this equation reduces the general statistical averaging procedure to Gaussian integrations. However, it contains the infinite series, thus we have to truncate this expression by some criteria. In most cases of interest, the weakly nonlinear evolution of cosmic fields satisfy $`M^{(n)}𝒪(\sigma _{0}^{}{}_{}{}^{n2})`$, as we will see in § 4. When this relation holds, we can expand the distribution function to arbitrary order in $`\sigma _0`$. In the following, we assume this relation, and introduce the normalized cumulants:
$`\widehat{M}_{\mu _1\mathrm{}\mu _n}^{(n)}={\displaystyle \frac{M_{\mu _1\mathrm{}\mu _n}^{(n)}}{\sigma _{0}^{}{}_{}{}^{n2}}},`$ (2.21)
which are assumed to be of order one in terms of $`\sigma _0`$. In this case, the equation (2.19) is expanded as, up to $`𝒪(\sigma _{0}^{}{}_{}{}^{2})`$,
$`F=F_\mathrm{G}+{\displaystyle \frac{1}{3!}}{\displaystyle \widehat{M}_{\mu _1\mu _2\mu _3}^{(3)}F_{,\mu _1\mu _2\mu _3}_\mathrm{G}\sigma _0}`$
$`+\left[{\displaystyle \frac{1}{4!}}{\displaystyle \widehat{M}_{\mu _1\mu _2\mu _3\mu _4}^{(4)}F_{,\mu _1\mu _2\mu _3\mu _4}_\mathrm{G}}+{\displaystyle \frac{1}{2(3!)^2}}{\displaystyle \widehat{M}_{\mu _1\mu _2\mu _3}^{(3)}\widehat{M}_{\nu _1\nu _2\nu _3}^{(3)}F_{,\mu _1\mu _2\mu _3\nu _1\nu _2\nu _3}_\mathrm{G}}\right]\sigma _{0}^{}{}_{}{}^{2}+𝒪(\sigma _{0}^{}{}_{}{}^{3}),`$ (2.22)
where we introduce the notation, $`F_{,\mu _1\mu _2\mu _3}^3F/A_{\mu _1}A_{\mu _2}A_{\mu _3}`$ etc. The calculation of the factors $`F_{,\mu \mathrm{}}_\mathrm{G}`$ is performed for individual statistic which gives the explicit form of the function $`F`$.
### 2.3 Two-point Correlations
In the expansion of equation (2.22), we need to calculate the Gaussian average of derivatives of the function $`F`$:
$`F_{,\mu _1\mu _2\mathrm{}}_\mathrm{G}={\displaystyle d^NAP_\mathrm{G}(𝑨)F_{,\mu _1\mu _2\mathrm{}}(𝑨)}={\displaystyle \frac{1}{(2\pi )^{N/2}\sqrt{det𝑴}}}{\displaystyle d^NA\mathrm{exp}\left(\frac{1}{2}𝑨^\mathrm{T}𝑴^1𝑨\right)F_{,\mu _1\mu _2\mathrm{}}(𝑨)}.`$ (2.23)
In most cases, the evaluation is analytically feasible because only Gaussian integration is needed. For the evaluation of such integration, we need the correlation matrix $`𝑴`$. Throughout this paper, we consider statistically homogeneous, and isotropic fields, in which case, the correlation matrix $`𝑴`$ is simplified because of the symmetry. In fact, the correlation matrix takes the following forms:
$`\alpha \alpha =1,`$ (2.24)
$`\alpha \alpha _{,i}=0,`$ (2.25)
$`\alpha \alpha _{,ij}={\displaystyle \frac{1}{d}}{\displaystyle \frac{\sigma _{1}^{}{}_{}{}^{2}}{\sigma _{0}^{}{}_{}{}^{2}}}\delta _{ij},`$ (2.26)
$`\alpha _{,i}\alpha _{,j}={\displaystyle \frac{1}{d}}{\displaystyle \frac{\sigma _{1}^{}{}_{}{}^{2}}{\sigma _{0}^{}{}_{}{}^{2}}}\delta _{ij},`$ (2.27)
$`\alpha _{,i}\alpha _{,jk}=0,`$ (2.28)
$`\alpha _{,ij}\alpha _{,kl}={\displaystyle \frac{1}{d(d+2)}}{\displaystyle \frac{\sigma _{2}^{}{}_{}{}^{2}}{\sigma _{0}^{}{}_{}{}^{2}}}\left(\delta _{ij}\delta _{kl}+\delta _{ik}\delta _{jl}+\delta _{il}\delta _{jk}\right),`$ (2.29)
where $`\alpha _{,ij}=\alpha /x_ix_j`$ etc. and we define the following quantities:
$`\sigma _{1}^{}{}_{}{}^{2}=f^2f=\alpha ^2\alpha \sigma _{0}^{}{}_{}{}^{2}.`$ (2.30)
$`\sigma _{2}^{}{}_{}{}^{2}=^2f^2f=^2\alpha ^2\alpha \sigma _{0}^{}{}_{}{}^{2}.`$ (2.31)
In most cases of interest, the derivatives of order higher than two do not appear in the definition of the statistics, so we do not give explicit form of those correlations.
In the following, we use the notation $`\eta _i=\alpha _{,i}`$ and $`\zeta _{ij}=\alpha _{,ij}`$ following Bardeen et al. (1986). As is often the case, when the second order derivatives $`\zeta _{ij}`$ appears in $`F`$ as simple polynomials, the following transform is particularly useful:
$`\stackrel{~}{\zeta }_{ij}=\zeta _{ij}+{\displaystyle \frac{1}{d}}{\displaystyle \frac{\sigma _{1}^{}{}_{}{}^{2}}{\sigma _{0}^{}{}_{}{}^{2}}}\delta _{ij}\alpha .`$ (2.32)
This transform erase the correlation between $`\alpha `$ and $`\stackrel{~}{\zeta }_{ij}`$, and the non-zero correlations are only
$`\alpha ^2=1,`$ (2.33)
$`\eta _{1}^{}{}_{}{}^{2}=\eta _{2}^{}{}_{}{}^{2}=\eta _{3}^{}{}_{}{}^{2}={\displaystyle \frac{1}{d}}{\displaystyle \frac{\sigma _{1}^{}{}_{}{}^{2}}{\sigma _{0}^{}{}_{}{}^{2}}},`$ (2.34)
$`\stackrel{~}{\zeta }_{11}^{}{}_{}{}^{2}=\stackrel{~}{\zeta }_{22}^{}{}_{}{}^{2}=\stackrel{~}{\zeta }_{33}^{}{}_{}{}^{2}={\displaystyle \frac{3}{d(d+2)}}{\displaystyle \frac{\sigma _{2}^{}{}_{}{}^{2}}{\sigma _{0}^{}{}_{}{}^{2}}}\left(1{\displaystyle \frac{d+2}{3d}}\gamma ^2\right),`$ (2.35)
$`\stackrel{~}{\zeta }_{11}\stackrel{~}{\zeta }_{22}=\stackrel{~}{\zeta }_{11}\stackrel{~}{\zeta }_{33}=\stackrel{~}{\zeta }_{22}\stackrel{~}{\zeta }_{33}={\displaystyle \frac{1}{d(d+2)}}{\displaystyle \frac{\sigma _{2}^{}{}_{}{}^{2}}{\sigma _{0}^{}{}_{}{}^{2}}}\left(1{\displaystyle \frac{d+2}{d}}\gamma ^2\right),`$ (2.36)
$`\stackrel{~}{\zeta }_{12}^{}{}_{}{}^{2}=\stackrel{~}{\zeta }_{13}^{}{}_{}{}^{2}=\stackrel{~}{\zeta }_{23}^{}{}_{}{}^{2}={\displaystyle \frac{1}{d(d+2)}}{\displaystyle \frac{\sigma _{2}^{}{}_{}{}^{2}}{\sigma _{0}^{}{}_{}{}^{2}}},`$ (2.37)
where
$`\gamma ={\displaystyle \frac{\sigma _{1}^{}{}_{}{}^{2}}{\sigma _0\sigma _2}}.`$ (2.38)
The Gaussian integration is straightforward if the function $`F`$ is expressed by polynomials of $`\stackrel{~}{\zeta }_{ij}`$.
If the function $`F`$ is more complicated in which $`\zeta _{ij}`$ are not simply given by polynomials, it is useful to completely diagonalize the correlation matrix $`𝑴`$ of equations (2.24)–(2.24). We introduce the following transform, which is quite similar one in Bardeen et al. (1986):
$`x={\displaystyle \frac{\sigma _0}{\sigma _2}}\left({\displaystyle \underset{i}{}}\zeta _{ii}+{\displaystyle \frac{\sigma _{1}^{}{}_{}{}^{2}}{\sigma _{0}^{}{}_{}{}^{2}}}\alpha \right),`$ (2.39)
$`y={\displaystyle \frac{\sigma _0}{\sigma _2}}{\displaystyle \frac{\zeta _{11}\zeta _{22},}{2}}`$ (2.40)
$`z={\displaystyle \frac{\sigma _0}{\sigma _2}}{\displaystyle \frac{\zeta _{11}+\zeta _{22}2\zeta _{33}}{2}}.`$ (2.41)
If $`d=1`$, we ignore the variables $`y,z`$, and similarly, if $`d=2`$, we ignore the variable $`z`$ in the above equations and in the following. The above transform is similar to Bardeen et al. (1986) but notice it is not identical. This transform completely diagonalize the correlation matrix:
$`\alpha ^2=1,`$ (2.42)
$`\eta _{1}^{}{}_{}{}^{2}=\eta _{2}^{}{}_{}{}^{2}=\eta _{3}^{}{}_{}{}^{2}={\displaystyle \frac{1}{d}}{\displaystyle \frac{\sigma _{1}^{}{}_{}{}^{2}}{\sigma _{0}^{}{}_{}{}^{2}}},`$ (2.43)
$`x^2=1\gamma ^2,`$ (2.44)
$`y^2={\displaystyle \frac{1}{d(d+2)}},`$ (2.45)
$`z^2={\displaystyle \frac{3}{d(d+2)}},`$ (2.46)
$`\zeta _{12}^{}{}_{}{}^{2}=\zeta _{13}^{}{}_{}{}^{2}=\zeta _{23}^{}{}_{}{}^{2}={\displaystyle \frac{1}{d(d+2)}}{\displaystyle \frac{\sigma _{2}^{}{}_{}{}^{2}}{\sigma _{0}^{}{}_{}{}^{2}}},`$ (2.47)
and all non-diagonal correlations are zero. For later convenience, we write the inverse relation of the transform of $`x`$, $`y`$, and $`z`$:
$`\zeta _{11}={\displaystyle \frac{\sigma _2}{4\sigma _0}}(x+4y+2z+\gamma \alpha ),`$ (2.48)
$`\zeta _{22}={\displaystyle \frac{\sigma _2}{4\sigma _0}}(x4y+2z+\gamma \alpha ),`$ (2.49)
$`\zeta _{33}={\displaystyle \frac{\sigma _2}{2\sigma _0}}(x2z+\gamma \alpha ),`$ (2.50)
If $`\zeta _{33}`$ does not appear in the function $`F`$, the variable $`z`$ should be omitted in the above equations. If both $`\zeta _{22}`$ and $`\zeta _{33}`$ do not appear in the function $`F`$, the variables $`y`$ and $`z`$ should be omitted.
After expressing the function $`F`$ in terms of the diagonalized variables $`\alpha `$, $`\eta _i`$, $`x`$, $`y`$, $`z`$ and $`\zeta _{ij}(i<j)`$, the calculation of the Gaussian integration of equation (2.23) is performed.
### 2.4 Three-point Correlations
In this paper, only first two terms in equation (2.22) are considered. Thus we evaluate $`M^{(3)}`$ here. When the spatial symmetry is taken into account, the quantity $`M^{(3)}`$ reduces to the following expressions:
$`\alpha ^2\alpha _{,i}=0,`$ (2.51)
$`\alpha ^2\alpha _{,ij}={\displaystyle \frac{\alpha ^2^2\alpha }{d}}\delta _{ij},`$ (2.52)
$`\alpha \alpha _{,i}\alpha _{,j}={\displaystyle \frac{\alpha ^2^2\alpha }{2d}}\delta _{ij},`$ (2.53)
$`\alpha \alpha _{,i}\alpha _{,jk}=0,`$ (2.54)
$`\alpha \alpha _{,ij}\alpha _{,kl}={\displaystyle \frac{\alpha ^2\alpha ^2\alpha }{d(d+2)}}\left(\delta _{ij}\delta _{kl}+\delta _{ik}\delta _{jl}+\delta _{il}\delta _{jk}\right){\displaystyle \frac{3(\alpha \alpha )^2\alpha }{d(d1)(d+2)}}\left[\delta _{ij}\delta _{kl}{\displaystyle \frac{d}{2}}\left(\delta _{ik}\delta _{jl}+\delta _{il}\delta _{jk}\right)\right],`$ (2.55)
$`\alpha _{,i}\alpha _{,j}\alpha _{,k}=0,`$ (2.56)
$`\alpha _{,i}\alpha _{,j}\alpha _{,kl}={\displaystyle \frac{(\alpha \alpha )^2\alpha }{d(d1)}}\left[\delta _{ij}\delta _{kl}{\displaystyle \frac{1}{2}}\left(\delta _{ik}\delta _{jl}+\delta _{il}\delta _{jk}\right)\right],`$ (2.57)
$`\alpha _{,i}\alpha _{,jk}\alpha _{,lm}=0,`$ (2.58)
etc. To prove the above equations, the following identities for isotropic fields are useful:
$`\alpha _{,i}\alpha _{,j}\alpha _{,ij}={\displaystyle \frac{1}{2}}\alpha \alpha )^2\alpha ,`$ (2.59)
$`\alpha \alpha _{,ij}\alpha _{,ij}={\displaystyle \frac{3}{2}}\alpha \alpha )^2\alpha +\alpha (^2\alpha )(^2\alpha ).`$ (2.60)
Although the above equations are valid for $`d0,1`$, the case $`d=1`$ is obtained by just ignoring the terms with $`(d1)^1`$. Generally, more complicated quantities can be appeared for $`M^{(3)}`$, but the above relations are sufficient for our applications in this paper.
According to the spatial symmetry of the above equations, the third-order correlations $`\widehat{M}_{\mu \nu \lambda }^{(3)}`$ are explicitly given. At this point, it is useful to define the following quantities:
$`S^{(0)}={\displaystyle \frac{f^3}{\sigma _{0}^{}{}_{}{}^{4}}}={\displaystyle \frac{\alpha ^3}{\sigma _0}},`$ (2.61)
$`S^{(1)}={\displaystyle \frac{3}{4}}{\displaystyle \frac{f^2(^2f)}{\sigma _{0}^{}{}_{}{}^{2}\sigma _{1}^{}{}_{}{}^{2}}}={\displaystyle \frac{3}{4}}{\displaystyle \frac{\alpha ^2(^2\alpha )\sigma _0}{\sigma _{1}^{}{}_{}{}^{2}}},`$ (2.62)
$`S^{(2)}={\displaystyle \frac{3d}{2(d1)}}{\displaystyle \frac{(ff)(^2f)}{\sigma _{1}^{}{}_{}{}^{4}}}={\displaystyle \frac{3d}{2(d1)}}{\displaystyle \frac{(\alpha \alpha )(^2\alpha )\sigma _{0}^{}{}_{}{}^{3}}{\sigma _{1}^{}{}_{}{}^{4}}},`$ (2.63)
$`S_2^{(2)}={\displaystyle \frac{f(^2f)(^2f)}{\sigma _{1}^{}{}_{}{}^{4}}}={\displaystyle \frac{\alpha (^2\alpha )(^2\alpha )\sigma _{0}^{}{}_{}{}^{3}}{\sigma _{1}^{}{}_{}{}^{4}}}.`$ (2.64)
We call quantities $`S^{(a)}`$ skewness parameters. The first one $`S^{(0)}`$ is usually called as skewness, and the others are its derivatives. They are to be calculated from usual perturbation theories in § 4. Using these quantities, the third-order correlations are given by
$`\widehat{M}_{000}^{(3)}=S^{(0)},`$ (2.65)
$`\widehat{M}_{00i}^{(3)}=0,`$ (2.66)
$`\widehat{M}_{00(ii)}^{(3)}={\displaystyle \frac{4}{3d}}{\displaystyle \frac{\sigma _{1}^{}{}_{}{}^{2}}{\sigma _{0}^{}{}_{}{}^{2}}}S^{(1)},`$ (2.67)
$`\widehat{M}_{00(ij)}^{(3)}=0,(i<j),`$ (2.68)
$`\widehat{M}_{0ii}^{(3)}={\displaystyle \frac{2}{3d}}{\displaystyle \frac{\sigma _{1}^{}{}_{}{}^{2}}{\sigma _{0}^{}{}_{}{}^{2}}}S^{(1)},`$ (2.69)
$`\widehat{M}_{0ij}^{(3)}=0,(ij),`$ (2.70)
$`\widehat{M}_{0i(jk)}^{(3)}=0,`$ (2.71)
$`\widehat{M}_{0(ii)(ii)}^{(3)}={\displaystyle \frac{1}{d^2(d+2)}}{\displaystyle \frac{\sigma _{1}^{}{}_{}{}^{4}}{\sigma _{0}^{}{}_{}{}^{4}}}[2(d1)S^{(2)}3dS_2^{(2)}),`$ (2.72)
$`\widehat{M}_{0(ii)(jj)}^{(3)}={\displaystyle \frac{1}{d^2(d+2)}}{\displaystyle \frac{\sigma _{1}^{}{}_{}{}^{4}}{\sigma _{0}^{}{}_{}{}^{4}}}\left[2S^{(2)}+dS_2^{(2)}\right],(ij),`$ (2.73)
$`\widehat{M}_{0(ii)(jk)}^{(3)}=0,(j<k),`$ (2.74)
$`\widehat{M}_{0(ij)(ij)}^{(3)}={\displaystyle \frac{1}{d(d+2)}}{\displaystyle \frac{\sigma _{1}^{}{}_{}{}^{4}}{\sigma _{0}^{}{}_{}{}^{4}}}\left[S^{(2)}S_2^{(2)}\right](i<j),`$ (2.75)
$`\widehat{M}_{0(ij)(kl)}^{(3)}=0,(i<j,k<l,ik),`$ (2.76)
$`\widehat{M}_{ijk}^{(3)}=0,`$ (2.77)
$`\widehat{M}_{ii(ii)}^{(3)}=0,`$ (2.78)
$`\widehat{M}_{ii(jj)}^{(3)}={\displaystyle \frac{2}{3d^2}}{\displaystyle \frac{\sigma _{1}^{}{}_{}{}^{4}}{\sigma _{0}^{}{}_{}{}^{4}}}S^{(2)},(ij),`$ (2.79)
$`\widehat{M}_{ii(jk)}^{(3)}=0,(j<k),`$ (2.80)
$`\widehat{M}_{ij(kk)}^{(3)}=0,(ij),`$ (2.81)
$`\widehat{M}_{ij(ij)}^{(3)}={\displaystyle \frac{1}{3d^2}}{\displaystyle \frac{\sigma _{1}^{}{}_{}{}^{4}}{\sigma _{0}^{}{}_{}{}^{4}}}S^{(2)},(i<j),`$ (2.82)
$`\widehat{M}_{ij(kl)}^{(3)}=0,(i<j,k<l,ik),`$ (2.83)
$`\widehat{M}_{i(jk)(lm)}^{(3)}=0,`$ (2.84)
and so forth, where repeated indices are not summed over in the above equations, and $`\widehat{M}^{(3)}`$ is symmetric under permutation of its indices. Thus, denoting $`F_{\mu \nu \lambda }F_{,\mu \nu \lambda }_\mathrm{G}`$ for simplicity,
$`{\displaystyle \underset{\mu ,\nu ,\lambda }{}}\widehat{M}_{\mu \nu \lambda }^{(3)}F_{\mu \nu \lambda }=\widehat{M}_{000}^{(3)}F_{000}+3{\displaystyle \underset{i}{}}\widehat{M}_{00(ii)}^{(3)}F_{00(ii)}+3{\displaystyle \underset{i}{}}\widehat{M}_{0ii}^{(3)}F_{0ii}+3{\displaystyle \underset{i}{}}\widehat{M}_{0(ii)(ii)}^{(3)}F_{0(ii)(ii)}`$
$`+\mathrm{\hspace{0.17em}6}{\displaystyle \underset{i<j}{}}\widehat{M}_{0(ii)(jj)}^{(3)}F_{0(ii)(jj)}+3{\displaystyle \underset{i<j}{}}\widehat{M}_{0(ij)(ij)}^{(3)}F_{0(ij)(ij)}+3{\displaystyle \underset{ij}{}}\widehat{M}_{ii(jj)}^{(3)}F_{ii(jj)}+6{\displaystyle \underset{i<j}{}}\widehat{M}_{ij(ij)}^{(3)}F_{ij(ij)}+\mathrm{}`$
$`=S^{(0)}F_{000}{\displaystyle \frac{2S^{(1)}}{d}}{\displaystyle \frac{\sigma _{1}^{}{}_{}{}^{2}}{\sigma _{0}^{}{}_{}{}^{2}}}\left(2{\displaystyle \underset{i}{}}F_{00(ii)}{\displaystyle \underset{i}{}}F_{0ii}\right)`$
$`{\displaystyle \frac{2S^{(2)}}{d^2(d+2)}}{\displaystyle \frac{\sigma _{1}^{}{}_{}{}^{4}}{\sigma _{0}^{}{}_{}{}^{4}}}\left[3(d1){\displaystyle \underset{i}{}}F_{0(ii)(ii)}6{\displaystyle \underset{i<j}{}}F_{0(ii)(jj)}+{\displaystyle \frac{3d}{2}}{\displaystyle \underset{i<j}{}}F_{0(ij)(ij)}+(d+2){\displaystyle \underset{ij}{}}F_{ii(jj)}(d+2){\displaystyle \underset{i<j}{}}F_{ij(ij)}\right]`$
$`+{\displaystyle \frac{3S_2^{(2)}}{d(d+2)}}{\displaystyle \frac{\sigma _{1}^{}{}_{}{}^{4}}{\sigma _{0}^{}{}_{}{}^{4}}}\left(3{\displaystyle \underset{i}{}}F_{0(ii)(ii)}+2{\displaystyle \underset{i<j}{}}F_{0(ii)(jj)}+{\displaystyle \underset{i<j}{}}F_{0(ij)(ij)}\right)+\mathrm{}.`$ (2.85)
This equation gives the second-order correction term of the statistical quantity $`F`$ through equation (2.22). Once the field $`f`$ is specified, the skewness parameters $`S^{(a)}`$ are calculated by dynamical perturbation theory of the field $`f`$. The remaining factors in the above equation are the Gaussian integrations of the derivatives of the function $`F`$, i.e., $`F_{,\mu \nu \lambda }_\mathrm{G}`$. These factors can be calculated once the function $`F`$ is given. In the next section, we calculate the latter factors for individual statistics.
## 3 STATISTICS OF SMOOTHED COSMIC FIELD
In this section we calculate the factor $`F_{,\mu _1\mu _2\mathrm{}}_\mathrm{G}`$ for each statistic. Some of the results in this section have been previously presented. The Edgeworth expansion of the PDF in §3.1 is a familiar result. The result of 3D genus statistic in §3.4 was already given by Matsubara (1994). We include these old results for completeness. Other subsections present new results.
### 3.1 Probability Distribution Function
Perhaps the simplest yet non-trivial statistic is the PDF, $`P(f)`$. The perturbative expansion of the PDF is known as the Edgeworth expansion (Scherrer & Bertschinger, 1991; Juszkiewicz et al., 1995; Bernardeau & Kofman, 1995). As the simplest example, we re-derive the known Edgeworth expansion from the point of view of our general formalism above (see also Matsubara, 1995a).
Since the PDF is simply given by $`P(f)=\delta (f^{}f)_f^{}`$, where $`\delta `$ is the Dirac’s delta function, the function $`F`$ in the previous section for PDF $`P(f)`$ is given by
$`F={\displaystyle \frac{1}{\sigma _0}}\delta (\alpha f/\sigma _0)`$ (3.1)
Since this form of $`F`$ does not depend on derivatives of $`\alpha `$, only $`F_{000}`$ survives in the equation (2.85). From equation (A8) with $`n=0`$ and $`k=3`$, $`F_{000}=(2\pi )^{1/2}e^{\nu ^2/2}H_3(\nu )`$, where $`\nu =f/\sigma _0`$. Thus, the PDF is derived from the equation (2.22):
$`P(f)={\displaystyle \frac{e^{\nu ^2/2}}{\sqrt{2\pi }\sigma _0}}\left[1+\sigma _0{\displaystyle \frac{S^{(0)}}{6}}H_3(\nu )+𝒪(\sigma _{0}^{}{}_{}{}^{2})\right],`$ (3.2)
which reproduces the well-known result.
There is less advantage of applying our formalism to this simple statistic which can be treated by standard methods. Our formalism has advantages when more non-trivial statistics are considered as shown below.
### 3.2 Level Crossing, Length and Area Statistics
Next three statistics we consider here are the level-crossing statistic $`N_1`$, the length statistic $`N_2`$ and the area statistic $`N_3`$. The level-crossing statistic is defined by the mean number of intersection of a straight line and threshold contours of the field. The length statistic is defined by the mean length of intersection of a 2D surface and the threshold contours of the field. The area statistic is defined by the mean area of the contour surface in a 3D space (Ryden, 1988a; Ryden et al., 1989; Matsubara, 1996). The level-crossing statistic is defined for 1D, 2D, and 3D cosmic fields, the length statistic is defined for 2D and 3D cosmic fields, while the area statistic is defined only for 3D cosmic fields. For statistically isotropic fields, those three statistics are proportional to each other.
In general, a statistic of the smoothed field $`f`$ is a function of the threshold $`f_\mathrm{t}`$, or of the normalized threshold $`\nu =f_\mathrm{t}/\sigma _0`$. The explicit expressions of statistics $`N_1`$, $`N_2`$, and $`N_3`$ are given by (Ryden, 1988a)
$`N_1(\nu )=\delta (\alpha \nu )\left|\eta _1\right|,`$ (3.3)
$`N_2(\nu )=\delta (\alpha \nu )\left[(\eta _1)^2+(\eta _2)^2\right]^{1/2},`$ (3.4)
$`N_3(\nu )=\delta (\alpha \nu )\left[(\eta _1)^2+(\eta _2)^2+(\eta _3)^2\right]^{1/2}.`$ (3.5)
For isotropic fields, these statistics are actually equivalent (Ryden, 1988a). In fact, the distribution function of $`\eta _i\alpha _{,i}`$ for fixed $`\alpha =\nu `$ is the function of only the magnitude $`|𝜼|`$. Thus, using spherical coordinates for $`𝜼`$, one can see
$`N_1(\nu )=N_2(\nu ){\displaystyle \frac{d\varphi }{2\pi }|\mathrm{cos}\varphi |}=N_3(\nu ){\displaystyle \frac{d\mathrm{\Omega }}{4\pi }|\mathrm{sin}\theta \mathrm{cos}\varphi |},`$ (3.6)
i.e.,
$`N_1(\nu )={\displaystyle \frac{2}{\pi }}N_2(\nu )={\displaystyle \frac{1}{2}}N_3(\nu ).`$ (3.7)
Thus we only need to consider $`N_1`$ which has the simplest expression and the rests of the statistics are automatically given by equation (3.7). However, if the field is anisotropic, such as the density field in redshift space (Matsubara, 1996), the equation (3.7) no longer holds, and equations (3.3)–(3.5) should be used for each statistic. The equation (2.85) only holds for statistically isotropic fields, but the equation (2.22) is applicable even for anisotropic fields.
Now we calculate the factor $`F_{,\mu _1\mu _2\mathrm{}}_\mathrm{G}`$ for the particular statistic $`N_1`$. The indices $`\mu _1,\mu _2,\mathrm{}`$ only take $`0`$ and $`1`$ for $`N_1`$ statistic. Let the number of $`0`$ be $`k`$ and the number of $`1`$ be $`l`$. Then the factor is given by
$`F_{,\mu _1\mu _2\mathrm{}}_\mathrm{G}=R(k,l)\left({\displaystyle \frac{}{\alpha }}\right)^k\left({\displaystyle \frac{}{\eta _1}}\right)^l\delta (\alpha \nu )|\eta _1|_\mathrm{G}.`$ (3.8)
Since the variables $`\alpha `$ and $`\eta _i`$ are uncorrelated in the Gaussian averaging, we can use equations (A8) and (A11) in appendix A. Thus the above Gaussian integration results in
$`R(k,l)={\displaystyle \frac{h_{l2}}{\pi }}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^{1l}e^{\nu ^2/2}H_k(\nu ),`$ (3.9)
where $`h_l`$ is given by equation (A7). Now, calculation of equation (2.85) is straightforward:
$`{\displaystyle \underset{\mu ,\nu ,\lambda }{}}\widehat{M}_{\mu \nu \lambda }^{(3)}F_{\mu \nu \lambda }=S^{(0)}R(3,0)+2S^{(1)}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^2R(1,2)={\displaystyle \frac{1}{\pi }}{\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}e^{\nu ^2/2}\left[S^{(0)}H_3(\nu )+2S^{(1)}H_1(\nu )\right].`$ (3.10)
On the other hand, the Gaussian contribution is simply given by
$`F_\mathrm{G}=R(0,0)={\displaystyle \frac{1}{\pi }}{\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}e^{\nu ^2/2}.`$ (3.11)
Thus, the perturbative expansion of equation (2.22) up to second order is finally given by
$`N_1(\nu )={\displaystyle \frac{1}{\pi }}{\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}e^{\nu ^2/2}\left\{1+\left[{\displaystyle \frac{S^{(0)}}{6}}H_3(\nu )+{\displaystyle \frac{S^{(1)}}{3}}H_1(\nu )\right]\sigma _0+𝒪(\sigma _{0}^{}{}_{}{}^{2})\right\}.`$ (3.12)
The second order formulas for area and length statistics are given by equation (3.7) with the above equation. To evaluate the above formula, the factor $`S^{(0)}`$ and $`S^{(1)}`$ should be known. Those factors are evaluated by usual perturbation theory in the following section.
### 3.3 2D Genus Statistic
The next statistics we consider are the genus statistics. The genus statistics have been attractive because it has a geometrical meaning of the clustering as well as the cosmological significance. The 2D genus statistic $`G_2`$ is defined in a two-dimensional plane $`S`$ in a $`d`$-dimensional space, so that $`d2`$ is required. In this plane $`S`$, there are contours corresponding to each threshold $`\nu `$. The 2D genus statistic is defined by the number of contours surrounding regions higher than the threshold value minus the number of contours surrounding regions lower than the threshold value (Adler, 1981; Coles, 1988; Melott et al., 1989; Gott et al., 1990). This definition is intuitive, but an alternative, equivalent definition is more useful: first we set an arbitrary, fixed direction on the plane $`S`$. Then there are maxima and minima of contours according to that direction. These points are classified into upcrossing points and downcrossing points with respect to that chosen direction. The number of those points are used to define the 2D genus statistic as the following way:
$`G_2={\displaystyle \frac{1}{2}}[(\text{\# of upcrossing minima})(\text{\# of upcrossing maxima})`$
$`(\text{\# of downcrossing minima})+(\text{\# of downcrossing maxima})],`$ (3.13)
per unit area of the surface. According to this definition, the explicit expression of the 2D genus statistic is given by
$`G_2(\nu )={\displaystyle \frac{1}{2}}\delta (\alpha \nu )\delta (\eta _1)|\eta _2|\zeta _{11}.`$ (3.14)
For this statistic, the indices $`\mu _1,\mu _2,\mathrm{}`$ only take $`0`$, $`1`$, $`2`$ and $`(11)`$. Let the number of $`0`$ be $`k`$, $`1`$ be $`l_1`$, $`2`$ be $`l_2`$, and $`(11)`$ be $`m`$. Then we need to calculate the following quantity:
$`F_{,\mu _1\mu _2\mathrm{}}_\mathrm{G}=R(k;l_1,l_2;m){\displaystyle \frac{1}{2}}\left({\displaystyle \frac{}{\alpha }}\right)^k\left({\displaystyle \frac{}{\eta _1}}\right)^{l_1}\left({\displaystyle \frac{}{\eta _2}}\right)^{l_2}\left({\displaystyle \frac{}{\zeta _{11}}}\right)^m\delta (\alpha \nu )\delta (\eta _1)|\eta _2|\zeta _{11}_\mathrm{G}.`$ (3.15)
Since the second derivative $`\zeta _{11}`$ appears as a polynomial, we just use the transform of equation (2.32). Then, from equations (A8), (A11) and (A12), the above equation reduces to
$`R(k;l_1,l_2;m)={\displaystyle \frac{h_{l_1}h_{l_22}}{(2\pi )^{3/2}}}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^{2l_1l_22m}e^{\nu ^2/2}\left[H_{k+1}(\nu )\delta _{m0}H_k(\nu )\delta _{m1}\right].`$ (3.16)
Thus, from equation (2.85),
$`{\displaystyle \underset{\mu ,\nu ,\lambda }{}}\widehat{M}_{\mu \nu \lambda }^{(3)}F_{\mu \nu \lambda }`$ $`=`$ $`2S^{(1)}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^2\left[2R(2;0,0;1)R(1;2,0;0)R(1;0,2;0)\right]2S^{(2)}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^4R(0;0;2,1)`$ (3.17)
$`=`$ $`{\displaystyle \frac{1}{(2\pi )^{3/2}}}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^2e^{\nu ^2/2}\left[S^{(0)}H_4(\nu )+4S^{(1)}H_2(\nu )+2S^{(2)}\right].`$
On the other hand, the Gaussian contribution is simply given by
$`F_\mathrm{G}=R(0;0,0;0)={\displaystyle \frac{1}{(2\pi )^{3/2}}}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^2e^{\nu ^2/2}H_1(\nu ).`$ (3.18)
Thus, the perturbative expansion of equation (2.22) up to second order is finally given by
$`G_2(\nu )={\displaystyle \frac{1}{(2\pi )^{3/2}}}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^2e^{\nu ^2/2}\left\{H_1(\nu )+\left[{\displaystyle \frac{S^{(0)}}{6}}H_4(\nu )+{\displaystyle \frac{2S^{(1)}}{3}}H_2(\nu )+{\displaystyle \frac{S^{(2)}}{3}}\right]\sigma _0+𝒪(\sigma _{0}^{}{}_{}{}^{2})\right\}.`$ (3.19)
### 3.4 3D Genus Statistic
The second order formula for 3D genus statistic was already derived by Matsubara (1994), but the detailed derivation was omitted. For completeness, we revisit the same quantity from our general point of view here. While the 2D genus statistic is defined by the number of contour lines in 2D surface, the 3D genus statistic (Gott, Melott & Dickinson, 1986) is defined by the number of contour surfaces and the number of handles in 3D space. Thus the 3D genus is defined only for cosmic fields of $`d=3`$. The 3D genus statistic $`G_3`$ is defined by
$`G_3=\left[(\text{\# of handles of contours})(\text{\# of isolated contours})\right],`$ (3.20)
per unit volume of the 3D space. This quantity is mathematically equivalent to $`1/2`$ times Euler characteristic of the contour surfaces, and thus is proportional to the total surface integral of local curvature of contours from the Gauss-Bonnet theorem. Although those definition is intuitive, an alternative, equivalent definition is more useful as in the 2D genus case. We set an arbitrary direction in the 3D space. Then there are maxima, minima and saddle points according to that direction. From the number of these points, the 3D genus is defined by
$`G_3={\displaystyle \frac{1}{2}}\left[(\text{\# of maxima})+(\text{\# of minima})(\text{\# of saddle points})\right],`$ (3.21)
per unit volume. According to this definition, the explicit expression of the 3D genus statistic is given by (Doroshkevich, 1970; Adler, 1981; Bardeen et al., 1986)
$`G_3(\nu )={\displaystyle \frac{1}{2}}\delta (\alpha \nu )\delta (\eta _1)\delta (\eta _2)|\eta _3|\left(\zeta _{11}\zeta _{22}\zeta _{12}^{}{}_{}{}^{2}\right).`$ (3.22)
We need to calculate the following quantity:
$`F_{,\mu _1\mu _2\mathrm{}}_\mathrm{G}`$ $`=`$ $`R(k;l_1,l_2,l_3;m_{11},m_{22},m_{12})`$ (3.23)
$``$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{}{\alpha }}\right)^k{\displaystyle \underset{i=1}{\overset{3}{}}}\left({\displaystyle \frac{}{\eta _i}}\right)^{l_i}{\displaystyle \underset{ij}{\overset{2}{}}}\left({\displaystyle \frac{}{\zeta _{ij}}}\right)^{m_{ij}}\delta (\alpha \nu )\delta (\eta _1)\delta (\eta _2)|\eta _3|\left(\zeta _{11}\zeta _{22}\zeta _{12}^{}{}_{}{}^{2}\right)_\mathrm{G}.`$
After applying the transform of equation (2.32), the Gaussian integration of $`\stackrel{~}{\zeta }_{ij}`$ fixing $`\alpha `$ is given by
$`{\displaystyle \underset{ij}{\overset{2}{}}}\left({\displaystyle \frac{}{\zeta _{ij}}}\right)^{m_{ij}}\left(\zeta _{11}\zeta _{22}\zeta _{12}^{}{}_{}{}^{2}\right)|\alpha _\mathrm{G}`$
$`=\left({\displaystyle \frac{\sigma _{1}^{}{}_{}{}^{2}}{d\sigma _{0}^{}{}_{}{}^{2}}}\right)^{2_{ij}m_{ij}}\left[H_2(\alpha )J_0^{(2)}(\{m_{ij}\})H_1(\alpha )J_1^{(2)}(\{m_{ij}\})+J_2^{(2)}(\{m_{ij}\})\right],`$ (3.24)
where $`J_m^{(2)}`$ is defined in Table 1.
Then, from equations (A8), (A11) and (A12), the equation (3.23) reduces to
$`R(k;l_1,l_2,l_3;m_{11},m_{22},m_{12})`$
$`={\displaystyle \frac{1}{(2\pi )^2}}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^{3_{i=1}^3l_i2_{ij}^2m_{ij}}h_{l_1}h_{l_2}h_{l_32}e^{\nu ^2/2}`$
$`\times \left[J_0^{(2)}(\{m_{ij}\})H_{k+2}(\nu )J_1^{(2)}(\{m_{ij}\})H_{k+1}(\nu )+J_2^{(2)}(\{m_{ij}\})H_k(\nu )\right].`$ (3.25)
Since the 3D genus is only defined for $`d3`$ and our universe has the spatial dimension 3, only $`d=3`$ is meaningful for actual cosmic fields. Nevertheless, we preserve the general dimension $`d`$ for some flavor of generality. Thus, from equation (2.85),
$`{\displaystyle \underset{\mu ,\nu ,\lambda }{}}\widehat{M}_{\mu \nu \lambda }^{(3)}F_{\mu \nu \lambda }={\displaystyle \frac{1}{(2\pi )^2}}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^3e^{\nu ^2/2}\left[S^{(0)}H_5(\nu )+6S^{(1)}H_3(\nu )+6S^{(2)}H_1(\nu )\right].`$ (3.26)
The Gaussian contribution is given by (Doroshkevich, 1970; Hamilton, Gott & Weinberg, 1986)
$`F_\mathrm{G}=R(0;0,0,0;0,0,0)={\displaystyle \frac{1}{(2\pi )^2}}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^3e^{\nu ^2/2}H_2(\nu ).`$ (3.27)
Thus, the perturbative expansion of equation (2.22) up to second order is finally given by
$`G_3(\nu )={\displaystyle \frac{1}{(2\pi )^2}}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^3e^{\nu ^2/2}\left\{H_2(\nu )+\left[{\displaystyle \frac{S^{(0)}}{6}}H_5(\nu )+S^{(1)}H_3(\nu )+S^{(2)}H_1(\nu )\right]\sigma _0+𝒪(\sigma _{0}^{}{}_{}{}^{2})\right\}.`$ (3.28)
The above equation with $`d=3`$ agrees with Matsubara (1994)<sup>1</sup><sup>1</sup>1The notations $`S`$, $`T`$, and $`U`$ in Matsubara (1994) are related to the notations here by $`S=S^{(0)}`$, $`T=2S^{(1)}/3`$, $`U=S^{(2)}/3`$.
### 3.5 2D Weighted Extrema Density
Next, we consider the weighted extrema density above a threshold $`\nu `$. The field extrema is defined to be points where all the first-order spatial derivatives of the field $`f`$ vanish: $`f/x_i=0`$. The weight $`\pm 1`$ is associated to each extrema according to the number of negative eigenvalues of spatial second-order derivatives of the field. When the threshold is high enough, the weighted extrema are approximately identified with the field peaks.
Mathematically, this statistic is equivalent to the 2D genus statistics, according to the Morse’s theorem (Morse, & Cairns, 1969; Adler, 1981). Therefore, we do not need to calculate separately to obtain the result on this statistic. However, we present the derivation of the extrema density as an alternative calculation to the 2D genus.
The weighted extrema in 2D field is given by
$`\rho _{\mathrm{e2}}(\nu )=\theta (\alpha \nu )\delta (\eta _1)\delta (\eta _2)\left(\zeta _{11}\zeta _{22}\zeta _{12}^{}{}_{}{}^{2}\right).`$ (3.29)
Following the similar calculation of previous examples, and using equations in Appendix A, we obtain
$`F_{,\mu _1\mu _2\mathrm{}}_\mathrm{G}`$ $`=`$ $`R(k;l_1,l_2;m_{11},m_{22},m_{12})`$ (3.30)
$``$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{}{\alpha }}\right)^k{\displaystyle \underset{i=1}{\overset{2}{}}}\left({\displaystyle \frac{}{\eta _i}}\right)^{l_i}{\displaystyle \underset{ij}{\overset{2}{}}}\left({\displaystyle \frac{}{\zeta _{ij}}}\right)^{m_{ij}}\delta (\alpha \nu )\delta (\eta _1)\delta (\eta _2)\left(\zeta _{11}\zeta _{22}\zeta _{12}^{}{}_{}{}^{2}\right)_\mathrm{G}`$
$`=`$ $`{\displaystyle \frac{1}{(2\pi )^{3/2}}}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^{2_{i=1}^2l_i2_{ij}^2m_{ij}}h_{l_1}h_{l_2}e^{\nu ^2/2}`$
$`\times \left[J_0^{(2)}(\{m_{ij}\})H_{k+1}(\nu )J_1^{(2)}(\{m_{ij}\})H_k(\nu )+J_2^{(2)}(\{m_{ij}\})H_{k1}(\nu )\right],`$
and the perturbative expansion of equation (2.22) up to second order is finally given by
$`\rho _{\mathrm{e2}}(\nu )={\displaystyle \frac{1}{(2\pi )^{3/2}}}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^2e^{\nu ^2/2}\left\{H_1(\nu )+\left[{\displaystyle \frac{S^{(0)}}{6}}H_4(\nu )+{\displaystyle \frac{2S^{(1)}}{3}}H_2(\nu )+{\displaystyle \frac{S^{(2)}}{3}}\right]\sigma _0+𝒪(\sigma _{0}^{}{}_{}{}^{2})\right\}.`$ (3.31)
This result is in fact equivalent to the 2D genus statistics given by equation (3.19).
### 3.6 Minkowski Functionals
Most of the Minkowski functionals are closely related to statistics discussed above. In this subsection, we comprehensively describe the exact relation between the Minkowski functionals $`V_k^{(d)}`$ of a smoothed field and the statistical quantities considered above.
The Minkowski functional of $`k=0`$ is simply the volume fraction of the excursion set $`K`$ which is defined by high-density regions above a given threshold $`\nu `$:
$`V_0^{(d)}(\nu )={\displaystyle \frac{1}{V}}{\displaystyle _K}𝑑V.`$ (3.32)
The other functionals with $`k=1,2,\mathrm{},d`$ are defined by the integral of the curvatures on isodensity surfaces of the threshold $`\nu `$ (Schmalzing & Buchert, 1997; Schmalzing & Gorski, 1998). In 3-dimensions, $`d=3`$, they are evaluated by a surface integration averaged over whole system of volume $`V`$ (Schmalzing & Buchert, 1997), i.e.,
$`V_k^{(3)}(\nu )={\displaystyle \frac{1}{V}}{\displaystyle _K}d^2A(𝒙)v_k^{(3)}(\nu ,𝒙),`$ (3.33)
where the local Minkowski functionals,
$`v_1^{(3)}(\nu ,𝒙)={\displaystyle \frac{1}{6}},`$ (3.34)
$`v_2^{(3)}(\nu ,𝒙)={\displaystyle \frac{1}{6\pi }}\left({\displaystyle \frac{1}{R_1}}+{\displaystyle \frac{1}{R_2}}\right),`$ (3.35)
$`v_3^{(3)}(\nu ,𝒙)={\displaystyle \frac{1}{4\pi }}{\displaystyle \frac{1}{R_1R_2}},`$ (3.36)
are defined by the principal curvatures $`1/R_1`$ and $`1/R_2`$ of the surface oriented toward lower density values. In 2-dimensions, $`d=2`$, the Minkowski functionals of $`k=1,2`$ are evaluated by a line integration averaged over whole system of 2D volume (surface) $`V`$ (Schmalzing & Gorski, 1998), i.e.,
$`V_k^{(2)}(\nu )={\displaystyle \frac{1}{V}}{\displaystyle _K}𝑑L(𝒙)v_k^{(2)}(\nu ,𝒙),`$ (3.37)
where the local Minkowski functionals,
$`v_1^{(2)}(\nu ,𝒙)={\displaystyle \frac{1}{4}},`$ (3.38)
$`v_2^{(2)}(\nu ,𝒙)={\displaystyle \frac{1}{2\pi }}{\displaystyle \frac{1}{R_1}},`$ (3.39)
are defined by the principal curvature $`1/R_1`$ of the line oriented toward lower density values.
All the Minkowski functionals for a Gaussian random field are analytically derived by Tomita (1986):
$`V_k^{(d)}(\nu )={\displaystyle \frac{1}{(2\pi )^{(k+1)/2}}}{\displaystyle \frac{\omega _d}{\omega _{dk}\omega _k}}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^ke^{\nu ^2/2}H_{k1}(\nu ),`$ (3.40)
where the factor $`\omega _k=\pi ^{k/2}/\mathrm{\Gamma }(k/2+1)`$ is the volume of the unit ball in $`k`$ dimensions, so that $`\omega _0=1`$, $`\omega _1=2`$, $`\omega _2=\pi `$, $`\omega _3=4\pi /3`$ (Schmalzing & Buchert, 1997).
It turns out that the Minkowski functionals in 2- and 3-dimensions are identical to the statistics $`N_1`$, $`G_2`$ and $`G_3`$ in each dimensions, except for normalization factors. In fact, from the Crofton’s formula (Crofton, 1868), the $`k`$-th Minkowski functional is given by
$`V_k^{(d)}={\displaystyle \frac{\omega _d}{\omega _{dk}\omega _k}}{\displaystyle _{_k^{(d)}}}𝑑\mu _k(E)\chi ^{(k)}(KE).`$ (3.41)
In this formula for body $`K`$ in $`d`$ dimensions, we consider an arbitrary $`k`$-dimensional hypersurface $`E`$ and calculate the Euler characteristic $`\chi ^{(k)}`$ of the intersection $`KE`$ in $`k`$ dimensions. This quantity is integrated over the space $`_k^{(d)}`$ of all conceivable hypersurfaces. The integration measure $`d\mu _k(E)`$ is normalized to give $`_{_k^{(d)}}𝑑\mu _k(E)=1`$. From this formula, we can see the statistics $`G_3`$, $`G_2`$, $`N_1`$ are proportional to the Minkowski functionals of $`V_3^{(d)}`$, $`V_2^{(d)}`$, $`V_1^{(d)}`$. In fact, $`\chi ^{(3)}`$ is given by $`1`$ times the 3D genus (or, equivalently, $`1/2`$ times the Euler number of boundaries, $`\chi ((KE))`$), $`\chi ^{(2)}`$ is identical to the 2D genus, and $`\chi ^{(1)}`$ is just $`1/2`$ times the number of level-crossing points (Adler, 1981). Thus, Minkowski functionals are given by
$`V_3^{(d)}(\nu )={\displaystyle \frac{\omega _d}{\omega _{d3}\omega _3}}G_3(\nu ),`$ (3.42)
$`V_2^{(d)}(\nu )={\displaystyle \frac{\omega _d}{\omega _{d2}\omega _2}}G_2(\nu ),`$ (3.43)
$`V_1^{(d)}(\nu )={\displaystyle \frac{\omega _d}{2\omega _{d1}\omega _1}}N_1(\nu ),`$ (3.44)
where the boundary of the body $`K`$ is identified with the isodensity contours of threshold $`\nu `$ and the statistics on right hand sides are defined in $`d`$-dimensions. Therefore, we have already obtained the weakly non-Gaussian expressions for Minkowski functionals, i.e., the Minkowski functionals of $`k=1,2,3`$ are given by the above equations and equations (3.12), (3.19), and (3.28). The Gaussian parts of the above equations exactly reproduce the Tomita’s formula (3.40).
The remaining Minkowski functional is the volume functional
$`V_0^{(d)}(\nu )=\theta (\nu \alpha ).`$ (3.45)
In this case, from equation (A8) in Appendix A,
$`R(k)\left({\displaystyle \frac{}{\alpha }}\right)^k\theta (\alpha \nu )_\mathrm{G}={\displaystyle \frac{e^{\nu ^2/2}}{\sqrt{2\pi }}}H_{k1}(\nu ).`$ (3.46)
so that equation (2.85) reduces to
$`V_0^{(d)}(\nu )={\displaystyle \frac{1}{2}}\mathrm{erfc}\left({\displaystyle \frac{\nu }{\sqrt{2}}}\right)+{\displaystyle \frac{e^{\nu ^2/2}}{\sqrt{2\pi }}}{\displaystyle \frac{S^{(0)}}{6}}H_2(\nu )\sigma _0+𝒪(\sigma _{0}^{}{}_{}{}^{2}).`$ (3.47)
The equivalent form can be obtained by integrating the Edgeworth expansion of PDF, equation (3.2).
All the formulas of Minkowski functionals derived above for $`0k3`$ fit into a single expression:
$`V_k^{(d)}(\nu )`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^{(k+1)/2}}}{\displaystyle \frac{\omega _d}{\omega _{dk}\omega _k}}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^ke^{\nu ^2/2}`$ (3.48)
$`\times \left\{H_{k1}(\nu )+\left[{\displaystyle \frac{1}{6}}S^{(0)}H_{k+2}(\nu )+{\displaystyle \frac{k}{3}}S^{(1)}H_k(\nu )+{\displaystyle \frac{k(k1)}{6}}S^{(2)}H_{k2}(\nu )\right]\sigma _0+𝒪(\sigma _{0}^{}{}_{}{}^{2})\right\}.`$
### 3.7 Rescaling the Threshold Density by Volume Fractions
The density threshold $`\nu `$ so far is simply defined so that isodensity surface is identified by $`f=\nu \sigma _0`$. However, the horizontal shift of the nonlinear genus curve etc. is considerably attributed to the nonlinear shift of probability distribution of the density field (e.g., Gott, Weinberg & Melott, 1987; Matsubara & Yokoyama, 1996). In order to cancel the latter shift, the threshold $`\stackrel{~}{\nu }`$ is defined so that the volume fraction $`f_V`$ on the high-density side of the isodensity surface equals to
$`f_V={\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle _{\stackrel{~}{\nu }}^{\mathrm{}}}𝑑te^{t^2/2}.`$ (3.49)
In fact, most of the work on genus analysis uses the genus curve plotted against the volume-fraction threshold $`\stackrel{~}{\nu }`$. Recently, Seto (2000) re-expressed the weakly non-Gaussian formula of genus curve (Matsubara, 1994) in terms of $`\stackrel{~}{\nu }`$, using perturbative expansion of the probability distribution function of the density field (Juszkiewicz et al., 1995). We follow this method to re-express the weakly non-Gaussian formulas of various statistical quantities derived above in terms of the volume-fraction threshold $`\stackrel{~}{\nu }`$.
The relation of $`\nu `$ and $`\stackrel{~}{\nu }`$ of weakly non-Gaussian field is simply given by equating the two equations (3.47) and (3.49). Up to first order in $`\sigma _0`$, the relation reduces to
$`\nu =\stackrel{~}{\nu }+{\displaystyle \frac{S^{(0)}}{6}}H_2(\stackrel{~}{\nu })\sigma _0+𝒪(\sigma _{0}^{}{}_{}{}^{2}).`$ (3.50)
It is straightforward to rewrite the various analytical formulas we derived above. The results for level-crossing statistic, 2D and 3D genus are respectively given by
$`N_1(\stackrel{~}{\nu })`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}e^{\stackrel{~}{\nu }^2/2}\left\{1+\left[{\displaystyle \frac{1}{3}}\left(S^{(1)}S^{(0)}\right)H_1(\stackrel{~}{\nu })\right]\sigma _0+𝒪(\sigma _{0}^{}{}_{}{}^{2})\right\},`$ (3.51)
$`G_2(\stackrel{~}{\nu })`$ $`=`$ $`\rho _{\mathrm{e2}}(\stackrel{~}{\nu })={\displaystyle \frac{1}{(2\pi )^{3/2}}}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^2e^{\stackrel{~}{\nu }^2/2}`$ (3.52)
$`\times \left\{H_1(\stackrel{~}{\nu })+\left[{\displaystyle \frac{2}{3}}\left(S^{(1)}S^{(0)}\right)H_2(\stackrel{~}{\nu })+{\displaystyle \frac{1}{3}}\left(S^{(2)}S^{(0)}\right)\right]\sigma _0+𝒪(\sigma _{0}^{}{}_{}{}^{2})\right\},`$
$`G_3(\stackrel{~}{\nu })`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^2}}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^3e^{\stackrel{~}{\nu }^2/2}`$ (3.53)
$`\times \left\{H_2(\stackrel{~}{\nu })+\left[\left(S^{(1)}S^{(0)}\right)H_3(\stackrel{~}{\nu })+\left(S^{(2)}S^{(0)}\right)H_1(\stackrel{~}{\nu })\right]\sigma _0+𝒪(\sigma _{0}^{}{}_{}{}^{2})\right\},`$
The results for Minkowski functionals of $`k=1,2,3`$ are again given by these equations and equations (3.42)–(3.44). All the Minkowski functionals for $`0k3`$ fit into a single expression:
$`V_k^{(d)}(\stackrel{~}{\nu })`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^{(k+1)/2}}}{\displaystyle \frac{\omega _d}{\omega _{dk}\omega _k}}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^ke^{\stackrel{~}{\nu }^2/2}`$ (3.54)
$`\times \left\{H_{k1}(\stackrel{~}{\nu })+\left[{\displaystyle \frac{k}{3}}\left(S^{(1)}S^{(0)}\right)H_k(\stackrel{~}{\nu })+{\displaystyle \frac{k(k1)}{6}}\left(S^{(2)}S^{(0)}\right)H_{k2}(\stackrel{~}{\nu })\right]\sigma _0+𝒪(\sigma _{0}^{}{}_{}{}^{2})\right\}.`$
Remarkably, the highest-order Hermite polynomial in the non-Gaussian correction terms vanishes in each case. In addition, the skewness parameters only appear as combinations of the form, $`S^{(a)}S^{(0)}`$, which makes the result simpler compared with the original form with the direct threshold $`\nu `$. As we will see in the following sections, the numerical values of $`S^{(a)},(a=0,1,2)`$ are quite close, or even identical in some special models. This means that the non-Gaussian corrections of the above statistics are smaller with the rescaled threshold $`\stackrel{~}{\nu }`$ than with the original threshold $`\nu `$. This is one of the central results in this paper. This tendency is in agreement with the analyses of numerical simulations.
## 4 SKEWNESS PARAMETERS FOR SMOOTHED FIELDS
We need to know the skewness parameters, $`S^{(a)}`$ for the evaluation of the second-order perturbative terms of equation (2.85). These quantities can be calculated by usual perturbation theory, once we specify the cosmic field, $`f`$. In each example of the previous section, the quantity $`S_2^{(2)}`$ does not appear so that we evaluate $`S^{(0)}`$, $`S^{(1)}`$, and $`S^{(2)}`$ in this section. The other kinds of skewness parameters like $`S_2^{(2)}`$ can be similarly evaluated without difficulty.
### 4.1 Hierarchical Model
Before we explore actual cosmic fields, we consider a simple, phenomenological statistical model, i.e., the hierarchical model of higher order correlation functions (e.g., Peebles, 1980). In this model, $`N`$-point correlation function is a sum of $`N1`$ products of the two-point correlation function. Specifically, the three-point correlation function is given by
$`f(𝒙_1)f(𝒙_2)f(𝒙_3)=Q[f(𝒙_1)f(𝒙_2)f(𝒙_2)f(𝒙_3)+f(𝒙_2)f(𝒙_3)f(𝒙_3)f(𝒙_1)`$
$`+f(𝒙_3)f(𝒙_1)f(𝒙_1)f(𝒙_2)].`$ (4.1)
We assume the field $`f`$ in the above equation is already smoothed. In this case, skewness parameters, equations (2.61)–(2.64) are given by straightforward calculations (Matsubara, 1994):
$`S^{(0)}=S^{(1)}=S^{(2)}=3Q.`$ (4.2)
These values depend on a hierarchical amplitude, $`Q`$, which is a free parameter of this model. The relative amplitudes among $`S^{(a)}`$ are not freely adjusted in the above equation. In the case of volume-fraction threshold, the first non-Gaussian correction of the various statistics considered in the previous section is absent, since they depend only on $`S^{(a)}S^{(0)}`$.
### 4.2 3D Density Field
Next, we consider the three-dimensional density field. The skewness parameters of this field in perturbation theory are already calculated by Matsubara (1994) and Matsubara & Suto (1996), using Fourier transforms of the field. We comprehensively review this calculation here for completeness. There is an alternative way to calculate the skewness not depending on Fourier transforms (Buchalter & Kamionkowski, 1999).
The cosmic field $`f`$ is identified with the 3D density contrast, $`\rho /\overline{\rho }1`$, where $`\rho `$ is the density field. The dimension this field is defined in is three, $`d=3`$. The Fourier transform of the field is useful in the following:
$`\stackrel{~}{f}(𝒌)={\displaystyle d^3xe^{i𝒌𝒙}f(𝒙)}.`$ (4.3)
In this notation, two- and three-point correlations in Fourier space have the forms,
$`\stackrel{~}{f}(𝒌_1)\stackrel{~}{f}(𝒌_2)=(2\pi )^2\delta ^3(𝒌_1+𝒌_2)P(k_1),`$ (4.4)
$`\stackrel{~}{f}(𝒌_1)\stackrel{~}{f}(𝒌_2)\stackrel{~}{f}(𝒌_3)=(2\pi )^3\delta ^3(𝒌_1+𝒌_2+𝒌_3)B(k_1,k_2,k_3).`$ (4.5)
The above forms are the consequence of the statistical homogeneity of the space, where the functions $`P`$ and $`B`$ are the power spectrum and the bispectrum, respectively. Thus, the variance and its variants \[Eqs. (2.4), (2.30), and (2.31)\], in their Fourier representation, are given by
$`\sigma _{j}^{}{}_{}{}^{2}={\displaystyle \frac{k^2dk}{2\pi ^2}k^{2j}P(k)},`$ (4.6)
and the skewness parameters of equations (2.61)–(2.64) are given by
$`S^{(0)}={\displaystyle \frac{1}{\sigma _{0}^{}{}_{}{}^{4}}}{\displaystyle \frac{d^3k_1}{(2\pi )^3}\frac{d^3k_2}{(2\pi )^3}B(k_1,k_2,|𝒌_1+𝒌_2|)},`$ (4.7)
$`S^{(1)}={\displaystyle \frac{3}{4\sigma _{0}^{}{}_{}{}^{2}\sigma _{1}^{}{}_{}{}^{2}}}{\displaystyle \frac{d^3k_1}{(2\pi )^3}\frac{d^3k_2}{(2\pi )^3}|𝒌_1+𝒌_2|^2B(k_1,k_2,|𝒌_1+𝒌_2|)},`$ (4.8)
$`S^{(2)}={\displaystyle \frac{9}{4\sigma _{1}^{}{}_{}{}^{4}}}{\displaystyle \frac{d^3k_1}{(2\pi )^3}\frac{d^3k_2}{(2\pi )^3}(𝒌_1𝒌_2)|𝒌_1+𝒌_2|^2B(k_1,k_2,|𝒌_1+𝒌_2|)}.`$ (4.9)
For initial random Gaussian density field, the second order perturbation theory predicts the power spectrum and the bispectrum as follows (e.g., Peebles, 1980; Fry, 1984; Bouchet et al., 1992; Bernardeau, 1994a):
$`P(k)=P_{\mathrm{LIN}}(k)W^2(kR)+𝒪(\sigma _{0}^{}{}_{}{}^{4}),`$ (4.10)
$`B(k_1,k_2,k_3)=\left[1+E+\left({\displaystyle \frac{k_2}{k_1}}+{\displaystyle \frac{k_1}{k_2}}\right){\displaystyle \frac{𝒌_1𝒌_2}{k_1k_2}}+(1E)\left({\displaystyle \frac{𝒌_1𝒌_2}{k_1k_2}}\right)^2\right]P_{\mathrm{LIN}}(k_1)P_{\mathrm{LIN}}(k_2)W(k_1R)W(k_2R)W(k_3R)`$
$`+\mathrm{cyc}.(1,2,3)+𝒪(\sigma _{0}^{}{}_{}{}^{6}),`$ (4.11)
where $`P_{\mathrm{LIN}}(k)`$ is the linear power spectrum, and $`E`$ is a weak function of cosmology (Bernardeau, 1994a; Bernardeau et al., 1995). The field smoothing corresponds to the multiplication of the window function $`W(kR)`$ which is the three-dimensional Fourier transform of the smoothing function $`W_R`$. It is a good approximation to use the value of $`E`$ for an Einstein-de Sitter universe, $`E=3/7`$, in most cases. The explicit form of the function $`E`$ in terms of cosmological parameters $`\mathrm{\Omega }_0`$ and $`\lambda _0`$ is given by Matsubara (1995b), which is accurately fitted by
$`E{\displaystyle \frac{3}{7}}\mathrm{\Omega }_{0}^{}{}_{}{}^{1/30}{\displaystyle \frac{\lambda _0}{80}}\left(1{\displaystyle \frac{3}{2}}\lambda _0\mathrm{log}_{10}\mathrm{\Omega }_0\right).`$ (4.12)
The perturbation theory is considered to be an expansion by a parameter $`\sigma _0`$. In this respect, the power spectrum $`P`$ and the bispectrum $`B`$ is of order $`\sigma _{0}^{}{}_{}{}^{2}`$ and $`\sigma _{0}^{}{}_{}{}^{4}`$, respectively.
Substituting equation (4.10) into equation (4.6), we obtain
$`\sigma _{j}^{}{}_{}{}^{2}(R)={\displaystyle \frac{k^2dk}{2\pi ^2}k^{2j}P_{\mathrm{LIN}}(k)W^2(kR)},`$ (4.13)
up to the lowest order. Similarly, substituting equation (4.11) into equations (4.7)–(4.9), and introducing new integration variables, $`l_1|𝒌_1|R`$, $`l_2|𝒌_2|R`$, and $`\mu =𝒌_1𝒌_2/(k_1k_2)`$, we obtain
$`S^{(a)}(R)={\displaystyle \frac{1}{\sigma _{0}^{}{}_{}{}^{4}}}\left({\displaystyle \frac{\sigma _0}{\sigma _1R}}\right)^{2a}{\displaystyle \frac{l_{1}^{}{}_{}{}^{2}dl_1}{2\pi ^2R^3}\frac{l_{2}^{}{}_{}{}^{2}dl_2}{2\pi ^2R^3}P_{\mathrm{LIN}}\left(\frac{l_1}{R}\right)P_{\mathrm{LIN}}\left(\frac{l_2}{R}\right)W^2(l_1)W^2(l_2)\stackrel{~}{S}^{(a)}(l_1,l_2)},`$ (4.14)
where $`a=0,1,2`$ and
$`\stackrel{~}{S}^{(a)}(l_1,l_2)={\displaystyle \frac{3}{4}}{\displaystyle _1^1}𝑑\mu \left[1+E+\left({\displaystyle \frac{l_2}{l_1}}+{\displaystyle \frac{l_1}{l_2}}\right)\mu +(1E)\mu ^2\right]{\displaystyle \frac{W\left(\sqrt{l_{1}^{}{}_{}{}^{2}+l_{2}^{}{}_{}{}^{2}+2l_1l_2\mu }\right)}{W(l_1)W(l_2)}}`$
$`\times \{\begin{array}{cc}2,\hfill & a=0,\hfill \\ l_{1}^{}{}_{}{}^{2}+l_{2}^{}{}_{}{}^{2}+l_1l_2\mu ,\hfill & a=1,\hfill \\ 3l_{1}^{}{}_{}{}^{2}l_{2}^{}{}_{}{}^{2}\left(1\mu ^2\right),\hfill & a=2,\hfill \end{array}`$ (4.18)
So far the smoothing function is arbitrary. For a general smoothing function, the above equations can be numerically integrated to obtain the skewness parameters for each model of the power spectrum. For some smoothing functions, further analytical reductions of the above equations are possible. As a popular example, we consider the Gaussian smoothing, $`W(l)=\mathrm{exp}(l^2/2)`$, which is frequently adopted for practical purposes. We follow the similar technique of Łokas et al. (1995), in which they derived the skewness of the density field with Gaussian smoothing.
For the Gaussian smoothing,
$`{\displaystyle \frac{W\left(\sqrt{l_{1}^{}{}_{}{}^{2}+l_{2}^{}{}_{}{}^{2}+2l_1l_2\mu }\right)}{W(l_1)W(l_2)}}=e^{l_1l_2\mu }.`$ (4.19)
In this case, the following formula of the modified Bessel function $`I_\nu (z)`$ is useful:
$`{\displaystyle _1^1}𝑑\mu P_l(\mu )e^{\mu z}=(1)^l\sqrt{{\displaystyle \frac{2\pi }{z}}}I_{l+1/2}(z),`$ (4.20)
where $`P_l`$ is the $`l`$-th Legendre polynomial. From this formula, the angular integration of $`\mu `$ in equation (4.18) can be analytically performed and the result is
$`\stackrel{~}{S}^{(0)}=\sqrt{{\displaystyle \frac{2\pi }{l_1l_2}}}\left[(2+E)I_{1/2}(l_1l_2){\displaystyle \frac{3}{2}}\left({\displaystyle \frac{l_2}{l_1}}+{\displaystyle \frac{l_1}{l_2}}\right)I_{3/2}(l_1l_2)+(1E)I_{5/2}(l_1l_2)\right],`$ (4.21)
$`\stackrel{~}{S}^{(1)}=\sqrt{2\pi l_1l_2}\{{\displaystyle \frac{5+2E}{4}}({\displaystyle \frac{l_2}{l_1}}+{\displaystyle \frac{l_1}{l_2}})I_{1/2}(l_1l_2)[{\displaystyle \frac{3(9+E)}{10}}+{\displaystyle \frac{l_{2}^{}{}_{}{}^{2}}{l_{1}^{}{}_{}{}^{2}}}+{\displaystyle \frac{l_{1}^{}{}_{}{}^{2}}{l_{2}^{}{}_{}{}^{2}}}]I_{3/2}(l_1l_2)`$
$`+{\displaystyle \frac{2E}{2}}({\displaystyle \frac{l_2}{l_1}}+{\displaystyle \frac{l_1}{l_2}})I_{5/2}(l_1l_2){\displaystyle \frac{3(1E)}{10}}I_{7/2}(l_1l_2)\},`$ (4.22)
$`\stackrel{~}{S}^{(2)}=\sqrt{2\pi }(l_1l_2)^{3/2}[{\displaystyle \frac{3(3+2E)}{5}}I_{1/2}(l_1l_2){\displaystyle \frac{9}{10}}({\displaystyle \frac{l_2}{l_1}}+{\displaystyle \frac{l_1}{l_2}})I_{3/2}(l_1l_2)`$
$`{\displaystyle \frac{3(3+4E)}{7}}I_{5/2}(l_1l_2)+{\displaystyle \frac{9}{10}}({\displaystyle \frac{l_2}{l_1}}+{\displaystyle \frac{l_1}{l_2}})I_{7/2}(l_1l_2){\displaystyle \frac{18(1E)}{35}}I_{9/2}(l_1l_2)].`$ (4.23)
At this point, it is useful to define the following quantity:
$`S_m^{\alpha \beta }(R){\displaystyle \frac{\sqrt{2\pi }}{\sigma _{0}^{}{}_{}{}^{4}}}\left({\displaystyle \frac{\sigma _0}{\sigma _1R}}\right)^{\alpha +\beta 2}{\displaystyle \frac{l_{1}^{}{}_{}{}^{2}dl_1}{2\pi ^2R^3}\frac{l_{2}^{}{}_{}{}^{2}dl_2}{2\pi ^2R^3}P_{\mathrm{LIN}}\left(\frac{l_1}{R}\right)P_{\mathrm{LIN}}\left(\frac{l_2}{R}\right)e^{l_{1}^{}{}_{}{}^{2}l_{2}^{}{}_{}{}^{2}}l_{1}^{}{}_{}{}^{\alpha 3/2}l_{2}^{}{}_{}{}^{\beta 3/2}I_{m+1/2}(l_1l_2)}.`$ (4.24)
In the above equation, variance parameters $`\sigma _0`$, $`\sigma _1`$ are given by equation (4.13). The nonlinear correction for $`\sigma _j`$ is not needed because our estimate of $`S^{(a)}`$ is only lowest order in $`\sigma _0`$. When the higher order corrections, e.g., third-order perturbation corrections, are estimated, one should be sure that all the necessary nonlinear corrections are properly taken into account. With this quantity, the skewness parameters are given by
$`S^{(0)}(R)=(2+E)S_0^{11}3S_1^{02}+(1E)S_2^{11},`$ (4.25)
$`S^{(1)}(R)={\displaystyle \frac{3}{2}}\left[{\displaystyle \frac{5+2E}{3}}S_0^{13}{\displaystyle \frac{9+E}{5}}S_1^{22}S_1^{04}+{\displaystyle \frac{2(2E)}{3}}S_2^{13}{\displaystyle \frac{1E}{5}}S_3^{22}\right],`$ (4.26)
$`S^{(2)}(R)=9\left[{\displaystyle \frac{3+2E}{15}}S_0^{33}{\displaystyle \frac{1}{5}}S_1^{24}{\displaystyle \frac{3+4E}{21}}S_2^{33}+{\displaystyle \frac{1}{5}}S_3^{24}{\displaystyle \frac{2(1E)}{35}}S_4^{33}\right].`$ (4.27)
Thus the lowest order estimates of skewness parameters are given by the above equations (4.24)–(4.27). For each given power spectrum, the integration of equation (4.24) is straightforward.
The resulting skewness parameters are independent on the amplitude of the power spectrum. When the power spectrum is given by a CDM-like model, $`P_{\mathrm{LIN}}(k)kT_{\mathrm{CDM}}^{}{}_{}{}^{2}(k/\mathrm{\Gamma })`$, where
$`T_{\mathrm{CDM}}(p)={\displaystyle \frac{\mathrm{ln}(1+2.34p)}{2.34p}}\left[1+3.89p+(16.1p)^2+(5.46p)^3+(6.71p)^4\right]^{1/4},`$ (4.28)
is the CDM-like transfer function fitted by Bardeen et al. (1986), and $`\mathrm{\Gamma }`$ is the shape parameter of this model, then the skewness parameters are functions of $`\mathrm{\Gamma }R`$. In Table 2, we give the values of skewness parameters $`S^{(a)}`$ for CDM-like models for several values of $`\mathrm{\Gamma }R`$.
In this table, the value of $`E`$ is approximated by $`3/7`$. Since the skewness parameters are weak functions of $`\mathrm{\Gamma }R`$ as seen from the table, one can interpolate the values in this table to obtain the values of arbitrary scales for practical purposes.
When the power spectrum is given by a power-law form,
$`P_{\mathrm{LIN}}(k)=Ak^{n_\mathrm{s}},`$ (4.29)
the integration of equation (4.24) can be analytically performed. First, the simple Gaussian integration gives
$`\sigma _{j}^{}{}_{}{}^{2}={\displaystyle \frac{A}{4\pi ^2R^{n_\mathrm{s}+2j+3}}}\mathrm{\Gamma }\left({\displaystyle \frac{n_\mathrm{s}+2j+3}{2}}\right).`$ (4.30)
Second, we expand the modified Bessel function as
$`I_\nu (z)=\left({\displaystyle \frac{z}{2}}\right)^\nu {\displaystyle \underset{r=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(z/2)^{2r}}{r!\mathrm{\Gamma }(\nu +r+1)}}.`$ (4.31)
Then the equation (4.24) reduces to
$`S_m^{\alpha \beta }={\displaystyle \frac{2}{(2m+1)!!}}\left({\displaystyle \frac{n_\mathrm{s}+3}{2}}\right)^{1(\alpha +\beta )/2}\left({\displaystyle \frac{n_\mathrm{s}+3}{2}}\right)_{(\alpha +m1)/2}\left({\displaystyle \frac{n_\mathrm{s}+3}{2}}\right)_{(\beta +m1)/2}`$
$`\times F({\displaystyle \frac{n_\mathrm{s}+\alpha +m+2}{2}},{\displaystyle \frac{n_\mathrm{s}+\beta +m+2}{2}},m+{\displaystyle \frac{3}{2}};{\displaystyle \frac{1}{4}}),`$ (4.32)
where $`(\alpha )_n=\mathrm{\Gamma }(\alpha +n)/\mathrm{\Gamma }(\alpha )=\alpha (\alpha +1)\mathrm{}(\alpha +n1)`$, and $`F`$ is the Gauss hypergeometric function:
$`F(\alpha ,\beta ,\gamma ;z)={\displaystyle \underset{r=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\alpha )_r(\beta )_r}{(\gamma )_r}}{\displaystyle \frac{z^r}{r!}}.`$ (4.33)
The equations (4.25)–(4.27) and equation (4.32) give the skewness parameters $`S^{(a)}`$. The following recursion relations for hypergeometric function
$`\alpha F(\alpha +1,\beta ,\gamma +1;z)=\gamma F(\alpha ,\beta ,\gamma ;z)+(\alpha \gamma )F(\alpha ,\beta ,\gamma +1;z),`$ (4.34)
$`\alpha \beta F(\alpha +1,\beta +1,\gamma +1;z)={\displaystyle \frac{\gamma (\gamma 1)}{z}}\left[F(\alpha ,\beta ,\gamma 1;z)F(\alpha ,\beta ,\gamma ;z)\right],`$ (4.35)
$`(\gamma \alpha )(\gamma \beta )zF(\alpha ,\beta ,\gamma +1;z)+\gamma [(\alpha +\beta 2\gamma +1)z+\gamma 1]F(\alpha ,\beta ,\gamma ;z)`$
$`+\gamma (\gamma 1)(z1)F(\alpha ,\beta ,\gamma 1;z)=0.`$ (4.36)
simplify the result:
$`S^{(0)}(n_\mathrm{s})=S^{(1)}(n_\mathrm{s})=3F({\displaystyle \frac{n_\mathrm{s}+3}{2}},{\displaystyle \frac{n_\mathrm{s}+3}{2}},{\displaystyle \frac{3}{2}};{\displaystyle \frac{1}{4}})(n_\mathrm{s}+22E)F({\displaystyle \frac{n_\mathrm{s}+3}{2}},{\displaystyle \frac{n_\mathrm{s}+3}{2}},{\displaystyle \frac{5}{2}};{\displaystyle \frac{1}{4}}),`$ (4.37)
$`S^{(2)}(n_\mathrm{s})=3F({\displaystyle \frac{n_\mathrm{s}+5}{2}},{\displaystyle \frac{n_\mathrm{s}+5}{2}},{\displaystyle \frac{5}{2}};{\displaystyle \frac{1}{4}}){\displaystyle \frac{3}{5}}(n_\mathrm{s}+44E)F({\displaystyle \frac{n_\mathrm{s}+5}{2}},{\displaystyle \frac{n_\mathrm{s}+5}{2}},{\displaystyle \frac{7}{2}};{\displaystyle \frac{1}{4}}).`$ (4.38)
In this power-law case, skewness parameters do not depend on scales $`R`$ but only on power-law index, $`n_\mathrm{s}`$. Incidentally, $`S^{(0)}(n_\mathrm{s})`$ and $`S^{(1)}(n_\mathrm{s})`$ are identical. This is just the coincidence and is not generally the case when the power spectrum is not given by power-law. Several numerical values are shown in Table 3, where $`E=3/7`$ is assumed.
### 4.3 3D Velocity Field
Next, we consider the 3D velocity field as the cosmic field $`f`$. Since the rotational components of the velocity field are decaying modes of gravitational evolution in perturbation theory, we only consider the rotation-free component. The cosmic field $`f`$ is identified with the dimensionless scalar field,
$`f(𝒙)={\displaystyle \frac{1}{H}}𝒗(𝒙),`$ (4.39)
where $`H`$ is the Hubble parameter. One can also consider other quantities like radial component of the velocity field, $`V=𝒏𝒗`$, where $`𝒏`$ is the line-of-sight normal vector. Those quantities are more complicated than a simple divergence. We illustrate only the simplest case in this paper. The second order perturbation theory predicts power spectrum and bispectrum of the velocity field as follows (e.g., Bernardeau, 1994a):
$`P(k)=g_{f}^{}{}_{}{}^{2}P_{\mathrm{LIN}}(k)W^2(kR)+𝒪(\sigma _{0}^{}{}_{}{}^{4}),`$ (4.40)
$`B(k_1,k_2,k_3)=g_{f}^{}{}_{}{}^{3}\left[1+E_v+\left({\displaystyle \frac{k_2}{k_1}}+{\displaystyle \frac{k_1}{k_2}}\right){\displaystyle \frac{𝒌_1𝒌_2}{k_1k_2}}+(1E_v)\left({\displaystyle \frac{𝒌_1𝒌_2}{k_1k_2}}\right)^2\right]P_{\mathrm{LIN}}(k_1)P_{\mathrm{LIN}}(k_2)W(k_1R)W(k_2R)W(k_3R)`$
$`+\mathrm{cyc}.(1,2,3)+𝒪(\sigma _{0}^{}{}_{}{}^{6}).`$ (4.41)
In the above equation, the factor $`g_f`$ is a logarithmic derivative of the growth factor, given by
$`g_f(\mathrm{\Omega }_0,\lambda _0)={\displaystyle \frac{d\mathrm{ln}D}{d\mathrm{ln}a}}|_0\mathrm{\Omega }_{0}^{}{}_{}{}^{4/7}+{\displaystyle \frac{\lambda _0}{70}}\left(1+{\displaystyle \frac{\mathrm{\Omega }_0}{2}}\right),`$ (4.42)
(Lightman & Schechter, 1990; Lahav et al., 1991) and $`D`$ is the linear growth factor and $`a`$ is the expansion factor. The logarithmic derivative is evaluated at present. The factor $`E_v`$ is a weak function of cosmology. It is a good approximation to use the value for an Einstein-de Sitter universe, $`E_v=1/7`$. The explicit form of the function $`E_v`$ in terms of $`\mathrm{\Omega }`$ and $`\lambda `$ is given by Matsubara (1995b), which is accurately fitted by
$`{\displaystyle \frac{E_v+1}{2}}{\displaystyle \frac{3}{7}}\mathrm{\Omega }_{0}^{}{}_{}{}^{11/200}{\displaystyle \frac{\lambda _0}{70}}\left(1{\displaystyle \frac{7}{3}}\lambda _0\mathrm{log}_{10}\mathrm{\Omega }_0\right).`$ (4.43)
The similarity of the equations (4.40) and (4.41) for velocity field to the equations (4.10), (4.11) for density field is obvious. We can easily see the skewness parameters for the velocity field are obtained by similar equations as (4.25)–(4.27):
$`S^{(0)}(R)={\displaystyle \frac{1}{g_f}}\left[(2+E_v)S_0^{11}3S_1^{02}+(1E_v)S_2^{11}\right],`$ (4.44)
$`S^{(1)}(R)={\displaystyle \frac{3}{2g_f}}\left[{\displaystyle \frac{5+2E_v}{3}}S_0^{13}{\displaystyle \frac{9+E_v}{5}}S_1^{22}S_1^{04}+{\displaystyle \frac{2(2E_v)}{3}}S_2^{13}{\displaystyle \frac{1E_v}{5}}S_3^{22}\right],`$ (4.45)
$`S^{(2)}(R)={\displaystyle \frac{9}{g_f}}\left[{\displaystyle \frac{3+2E_v}{15}}S_0^{33}{\displaystyle \frac{1}{5}}S_1^{24}{\displaystyle \frac{3+4E_v}{21}}S_2^{33}+{\displaystyle \frac{1}{5}}S_3^{24}{\displaystyle \frac{2(1E_v)}{35}}S_4^{33}\right],`$ (4.46)
where $`S_m^{\alpha \beta }`$ is given by equation (4.24) without any modification. In Table 4, we show the values of skewness parameters of the velocity field for the CDM-like models for several values of $`\mathrm{\Gamma }R`$.
In this table, the value of $`E_v`$ is set as $`1/7`$.
For the power spectrum of the power-law form,
$`S^{(0)}(n_\mathrm{s})=S^{(1)}(n_\mathrm{s})={\displaystyle \frac{1}{g_f}}\left[3F({\displaystyle \frac{n_\mathrm{s}+3}{2}},{\displaystyle \frac{n_\mathrm{s}+3}{2}},{\displaystyle \frac{3}{2}};{\displaystyle \frac{1}{4}})(n_\mathrm{s}+22E_v)F({\displaystyle \frac{n_\mathrm{s}+3}{2}},{\displaystyle \frac{n_\mathrm{s}+3}{2}},{\displaystyle \frac{5}{2}};{\displaystyle \frac{1}{4}})\right],`$ (4.47)
$`S^{(2)}(n_\mathrm{s})={\displaystyle \frac{3}{g_f}}\left[F({\displaystyle \frac{n_\mathrm{s}+5}{2}},{\displaystyle \frac{n_\mathrm{s}+5}{2}},{\displaystyle \frac{5}{2}};{\displaystyle \frac{1}{4}}){\displaystyle \frac{1}{5}}(n_\mathrm{s}+44E_v)F({\displaystyle \frac{n_\mathrm{s}+5}{2}},{\displaystyle \frac{n_\mathrm{s}+5}{2}},{\displaystyle \frac{7}{2}};{\displaystyle \frac{1}{4}})\right].`$ (4.48)
Several numerical values are shown in Table 5, where $`E_v=1/7`$ is assumed.
### 4.4 2D Projected Density Field
The projection of the density field $`\rho _\mathrm{p}`$ defines the 2D cosmic fields on the sky ($`d=2`$). Here, we derive the skewness parameters for this field. The 2D projected density field with top-hat kernel is investigated by Bernardeau (1995). Here, we are interested in Gaussian kernel for our purpose.
In a Friedman-Lemaître universe, the comoving angular diameter distance at a comoving distance $`\chi `$ is given by
$`S_K(\chi )=\{\begin{array}{cc}{\displaystyle \frac{\mathrm{sinh}\left(\chi \sqrt{K}\right)}{\sqrt{K}}},\hfill & K<0,\hfill \\ \chi ,\hfill & K=0,\hfill \\ {\displaystyle \frac{\mathrm{sin}\left(\chi \sqrt{K}\right)}{\sqrt{K}}},\hfill & K>0,\hfill \end{array}`$ (4.52)
depending on the sign of the spatial curvature $`K=H_0^{\mathrm{\hspace{0.17em}2}}(\mathrm{\Omega }_0+\lambda _01)`$. Thus, in spherical coordinates, projected density field in 2D is given by
$`\rho _\mathrm{p}(\theta ,\varphi )={\displaystyle 𝑑\chi S_{K}^{}{}_{}{}^{2}(\chi )n(\chi )\rho (\chi ,\theta ,\varphi ;\tau _0\chi )},`$ (4.53)
where $`n(\chi )`$ is the selection function without volume factor, normalized as $`𝑑\chi S_K^{\mathrm{\hspace{0.17em}2}}(\chi )n(\chi )=1`$, and $`\rho (\chi ,\theta ,\varphi ;\tau )`$ is the 3D comoving density field<sup>2</sup><sup>2</sup>2The comoving density field is defined by the density per unit comoving volume and thus satisfy $`\overline{\rho }=`$ constant.. The present value of the conformal time $`\tau =𝑑t/a`$ is $`\tau _0`$.
The projected density contrast is defined by $`\rho _\mathrm{p}/\overline{\rho }_\mathrm{p}1`$, where one can see $`\overline{\rho }_\mathrm{p}=\overline{\rho }`$. We identify the projected density contrast with the 2D ($`d=2`$) field $`f`$. Since the smoothing angle $`\theta _\mathrm{f}`$ is much smaller than $`\pi `$ in most of the interested cases, we consider the small patch of the sky of the vicinity of the polar axis, $`\theta 1`$. With this approximation, we introduce the variables $`\theta _1=\theta \mathrm{cos}\varphi `$, and $`\theta _2=\theta \mathrm{sin}\varphi `$, which are considered as 2D Euclidean coordinates, $`𝜽`$. Therefore, the projection equation is given by
$`f(𝜽)={\displaystyle 𝑑\chi S_K^{\mathrm{\hspace{0.17em}2}}(\chi )n(\chi )\delta _{3\mathrm{D}}(\chi ,S_K(\chi )𝜽;\tau _0\chi )},`$ (4.54)
where $`\delta _{3\mathrm{D}}(𝒙,\tau )=\rho /\overline{\rho }1`$ is the density contrast at comoving coordinates $`𝒙`$ and conformal time $`\tau `$.
The power spectrum and the bispectrum for the above projected field are given by the Limber’s equations (B2) and (B10) in Appendix B, with $`q(\chi )=n(\chi )`$:
$`P_{2\mathrm{D}}(\omega )={\displaystyle 𝑑\chi S_K^{\mathrm{\hspace{0.17em}2}}(\chi )n^2(\chi )P_{3\mathrm{D}}(\frac{\omega }{S_K(\chi )};\tau _0\chi )},`$ (4.55)
$`B_{2\mathrm{D}}(\omega _1,\omega _2,\omega _3)={\displaystyle 𝑑\chi S_K^{\mathrm{\hspace{0.17em}2}}(\chi )n^3(\chi )B_{3\mathrm{D}}(\frac{\omega _1}{S_K(\chi )},\frac{\omega _2}{S_K(\chi )},\frac{\omega _2}{S_K(\chi )};\tau _0\chi )},`$ (4.56)
where $`P`$, $`B`$ are 2D projected power spectrum and bispectrum, respectively, of the field $`f`$, and $`P_{3\mathrm{D}}`$, $`B_{3\mathrm{D}}`$ are 3D power spectrum and bispectrum, respectively. The 3D power spectrum and the 3D bispectrum are evaluated by the second order perturbation theory. They are similar to equations (4.10) and (4.11), but we have to take into account the time-dependence here. They are given by
$`P_{3\mathrm{D}}(k;\tau _0\chi )=D^2(\chi )P_{\mathrm{LIN}}(k)`$ (4.57)
$`B_{3\mathrm{D}}(k_1,k_2,k_3;\tau _0\chi )=D^4(\chi )\left[1+E(\chi )+\left({\displaystyle \frac{k_2}{k_1}}+{\displaystyle \frac{k_1}{k_2}}\right){\displaystyle \frac{𝒌_1𝒌_2}{k_1k_2}}+(1E(\chi ))\left({\displaystyle \frac{𝒌_1𝒌_2}{k_1k_2}}\right)^2\right]P_{\mathrm{LIN}}(k_1)P_{\mathrm{LIN}}(k_2)`$
$`+\mathrm{cyc}.(1,2,3),`$ (4.58)
where $`D(\chi )`$ is the linear growth factor at conformal lookback time $`\chi `$, (i.e., at conformal time $`\tau =\tau _0\chi `$), which is normalized as $`D(0)=1`$. The following fitting formula (Carroll, Press & Turner, 1992) is useful:
$`D{\displaystyle \frac{a\mathrm{\Omega }}{\mathrm{\Omega }_0}}{\displaystyle \frac{\mathrm{\Omega }_{0}^{}{}_{}{}^{4/7}\lambda _0+(1+\mathrm{\Omega }_0/2)(1+\lambda _0/70)}{\mathrm{\Omega }^{4/7}\lambda +(1+\mathrm{\Omega }/2)(1+\lambda /70)}},`$ (4.59)
where $`\mathrm{\Omega }`$ and $`\lambda `$ are time-dependent cosmological parameters at conformal lookback time $`\chi `$. The variable $`E(\chi )`$ is a weak function of time and cosmology, and for Einstein-de Sitter universe, $`E=3/7`$. This quantity $`E`$ is the same we used in 3D density field, but here we also take into account the time-dependence. It is accurately approximated by
$`E{\displaystyle \frac{3}{7}}\mathrm{\Omega }^{1/30}{\displaystyle \frac{\lambda }{80}}\left(1{\displaystyle \frac{3}{2}}\lambda \mathrm{log}_{10}\mathrm{\Omega }\right).`$ (4.60)
The variance parameters of the smoohted projected field are given by
$`\sigma _{j}^{}{}_{}{}^{2}(\theta _\mathrm{f})={\displaystyle \frac{\omega d\omega }{2\pi }\omega ^{2j}P(\omega )W^2(\omega \theta _\mathrm{f})}={\displaystyle \frac{1}{\theta _{\mathrm{f}}^{}{}_{}{}^{2j+2}}}{\displaystyle 𝑑\chi S_K^{\mathrm{\hspace{0.17em}2}}(\chi )n^2(\chi )D^2(\chi )\mathrm{\Sigma }_{j}^{}{}_{}{}^{2}\left[S_K(\chi )\theta _\mathrm{f}\right]},`$ (4.61)
where
$`\mathrm{\Sigma }_{j}^{}{}_{}{}^{2}(R)=R^{2j+2}{\displaystyle \frac{kdk}{2\pi }k^{2j}P_{\mathrm{LIN}}(k)W^2(kR)}={\displaystyle \frac{ldl}{2\pi }l^{2j}P_{\mathrm{LIN}}\left(\frac{l}{R}\right)W^2(l)}.`$ (4.62)
The skewness parameters of the smoothed projected field are given by
$`S^{(a)}(\theta _\mathrm{f})={\displaystyle \frac{1}{\sigma _{0}^{}{}_{}{}^{4}\theta _{\mathrm{f}}^{}{}_{}{}^{4}}}\left({\displaystyle \frac{\sigma _0}{\sigma _1\theta _\mathrm{f}}}\right)^{2a}{\displaystyle 𝑑\chi S_{K}^{}{}_{}{}^{2}(\chi )n^3(\chi )D^4(\chi )\mathrm{\Sigma }_{0}^{}{}_{}{}^{42a}\left[S_K(\chi )\theta _\mathrm{f}\right]\mathrm{\Sigma }_{1}^{}{}_{}{}^{2a}\left[S_K(\chi )\theta _\mathrm{f}\right]C^{(a)}\left[S_K(\chi )\theta _\mathrm{f}\right]},`$ (4.63)
where
$`C^{(a)}(R)={\displaystyle \frac{1}{\mathrm{\Sigma }_{0}^{}{}_{}{}^{4}}}\left({\displaystyle \frac{\mathrm{\Sigma }_0}{\mathrm{\Sigma }_1}}\right)^{2a}{\displaystyle \frac{l_1dl_1}{2\pi }\frac{l_2dl_2}{2\pi }P_{\mathrm{LIN}}\left(\frac{l_1}{R}\right)P_{\mathrm{LIN}}\left(\frac{l_2}{R}\right)W^2(l_1)W^2(l_2)\stackrel{~}{C}^{(a)}(l_1,l_2)},`$ (4.64)
and
$`\stackrel{~}{C}^{(a)}(l_1,l_2)={\displaystyle \frac{3}{2\pi }}{\displaystyle _1^1}{\displaystyle \frac{d\mu }{\sqrt{1\mu ^2}}}\left[1+E+\left({\displaystyle \frac{l_2}{l_1}}+{\displaystyle \frac{l_1}{l_2}}\right)\mu +(1E)\mu ^2\right]{\displaystyle \frac{W\left(\sqrt{l_1^{\mathrm{\hspace{0.17em}2}}+l_2^{\mathrm{\hspace{0.17em}2}}+2l_1l_2\mu }\right)}{W(l_1)W(l_2)}}`$
$`\times \{\begin{array}{cc}2,\hfill & a=0,\hfill \\ l_{1}^{}{}_{}{}^{2}+l_{2}^{}{}_{}{}^{2}+l_1l_2\mu ,\hfill & a=1,\hfill \\ 4l_{1}^{}{}_{}{}^{2}l_{2}^{}{}_{}{}^{2}\left(1\mu ^2\right),\hfill & a=2,\hfill \end{array}`$ (4.68)
So far the smoothing function is arbitrary. For a general smoothing function, the above equations can be numerically integrated to obtain the skewness parameters for each model of the power spectrum. In the following, we adopt the Gaussian smoothing function, $`W(l)=\mathrm{exp}(l^2/2)`$. For this smoothing function, the equation (4.19) holds even for this 2D case. In this case, the following integral representation of the modified Bessel function $`I_\nu (z)`$ for $`\nu =0`$ is useful:
$`I_0(z)={\displaystyle \frac{1}{\pi }}{\displaystyle _1^1}{\displaystyle \frac{d\mu }{\sqrt{1\mu ^2}}}e^{z\mu }.`$ (4.69)
Actually, the derivatives of the above equation,
$`{\displaystyle \frac{1}{\pi }}{\displaystyle _1^1}{\displaystyle \frac{d\mu }{\sqrt{1\mu ^2}}}\mu ^me^{z\mu }=\left({\displaystyle \frac{d}{dz}}\right)^mI_0(z),`$ (4.70)
are sufficient to perform the angular integration in equation (4.68). Moreover, one can use the property of the Bessel function, $`I_0^{}=I_1`$, and $`I_m^{}=(I_{m1}+I_{m+1})/2`$ to obtain formulas,
$`{\displaystyle \frac{1}{\pi }}{\displaystyle _1^1}{\displaystyle \frac{d\mu }{\sqrt{1\mu ^2}}}\mu e^{z\mu }=I_1(z),`$ (4.71)
$`{\displaystyle \frac{1}{\pi }}{\displaystyle _1^1}{\displaystyle \frac{d\mu }{\sqrt{1\mu ^2}}}\mu ^2e^{z\mu }={\displaystyle \frac{1}{2}}I_0(z)+{\displaystyle \frac{1}{2}}I_2(z),`$ (4.72)
$`{\displaystyle \frac{1}{\pi }}{\displaystyle _1^1}{\displaystyle \frac{d\mu }{\sqrt{1\mu ^2}}}\mu ^3e^{z\mu }={\displaystyle \frac{3}{4}}I_1(z){\displaystyle \frac{1}{4}}I_3(z),`$ (4.73)
$`{\displaystyle \frac{1}{\pi }}{\displaystyle _1^1}{\displaystyle \frac{d\mu }{\sqrt{1\mu ^2}}}\mu ^4e^{z\mu }={\displaystyle \frac{3}{8}}I_0(z)+{\displaystyle \frac{1}{2}}I_2(z)+{\displaystyle \frac{1}{8}}I_4(z).`$ (4.74)
From these formulas, equation (4.68) reduce to
$`\stackrel{~}{C}^{(0)}={\displaystyle \frac{3}{2}}\left[(3+E)I_0(l_1l_2)2\left({\displaystyle \frac{l_2}{l_1}}+{\displaystyle \frac{l_1}{l_2}}\right)I_1(l_1l_2)+(1E)I_2(l_1l_2)\right],`$ (4.75)
$`\stackrel{~}{C}^{(1)}={\displaystyle \frac{3}{4}}\{(4+E)(l_1^{\mathrm{\hspace{0.17em}2}}+l_2^{\mathrm{\hspace{0.17em}2}})I_0(l_1l_2)[{\displaystyle \frac{15+E}{2}}l_1l_2+2({\displaystyle \frac{l_1^{\mathrm{\hspace{0.17em}3}}}{l_2}}+{\displaystyle \frac{l_2^{\mathrm{\hspace{0.17em}3}}}{l_1}})]I_1(l_1l_2)`$
$`+(2E)(l_1^{\mathrm{\hspace{0.17em}2}}+l_2^{\mathrm{\hspace{0.17em}2}})I_2(l_1l_2){\displaystyle \frac{1E}{2}}l_1l_2I_3(l_1l_2)\},`$ (4.76)
$`\stackrel{~}{C}^{(2)}={\displaystyle \frac{3}{4}}[(5+3E)l_1^{\mathrm{\hspace{0.17em}2}}l_2^{\mathrm{\hspace{0.17em}2}}I_0(l_1l_2)2l_1l_2(l_1^{\mathrm{\hspace{0.17em}2}}+l_2^{\mathrm{\hspace{0.17em}2}})I_1(l_1l_2)4(1+E)l_1^{\mathrm{\hspace{0.17em}2}}l_2^{\mathrm{\hspace{0.17em}2}}I_2(l_1l_2)`$
$`+\mathrm{\hspace{0.33em}2}l_1l_2(l_1^{\mathrm{\hspace{0.17em}2}}+l_2^{\mathrm{\hspace{0.17em}2}})I_3(l_1l_2)(1E)l_1^{\mathrm{\hspace{0.17em}2}}l_2^{\mathrm{\hspace{0.17em}2}}I_4(l_1l_2)].`$ (4.77)
At this point, defining
$`C_m^{\alpha \beta }(R)={\displaystyle \frac{1}{\mathrm{\Sigma }_{0}^{}{}_{}{}^{4}}}\left({\displaystyle \frac{\mathrm{\Sigma }_0}{\mathrm{\Sigma }_1}}\right)^{\alpha +\beta 2}{\displaystyle \frac{l_1dl_1}{2\pi }\frac{l_2dl_2}{2\pi }P_{\mathrm{LIN}}\left(\frac{l_1}{R}\right)P_{\mathrm{LIN}}\left(\frac{l_2}{R}\right)e^{l_1^{\mathrm{\hspace{0.17em}2}}l_2^{\mathrm{\hspace{0.17em}2}}}l_{1}^{}{}_{}{}^{\alpha 1}l_{2}^{}{}_{}{}^{\beta 1}I_m(l_1l_2)},`$ (4.78)
the equation (4.64) reduces to
$`C^{(0)}={\displaystyle \frac{3}{2}}\left[(3+E)C_0^{11}4C_1^{02}+(1E)C_2^{11}\right],`$ (4.79)
$`C^{(1)}={\displaystyle \frac{3}{4}}\left[2(4+E)C_0^{13}{\displaystyle \frac{15+E}{2}}C_1^{22}4C_1^{04}+2(2E)C_2^{13}{\displaystyle \frac{1E}{2}}C_3^{22}\right],`$ (4.80)
$`C^{(2)}={\displaystyle \frac{3}{4}}\left[(5+3E)C_0^{33}4C_1^{24}4(1+E)C_2^{33}+4C_3^{24}(1E)C_4^{33}\right].`$ (4.81)
The two-dimensional integration of equation (4.78) is performed only once as a function of $`R`$. The result is stored as a table, and is used in one-dimensional numerical integration of equation (4.63) for finally obtaining the skewness parameters in 2D projected density fields. Functions $`C^{(a)}`$ for CDM-like models are given in Table 6 which are functions of $`\mathrm{\Gamma }R`$.
When the 3D power spectrum is given by a power-law form of equation (4.29), the integration by wave length $`l_1`$, $`l_2`$ can be analytically performed as in the 3D case. In fact, the parameter $`\mathrm{\Sigma }_j`$ of equation (4.62) with Gaussian smoothing $`W(l)=e^{l^2/2}`$ is given by
$`\mathrm{\Sigma }_{j}^{}{}_{}{}^{2}(R)={\displaystyle \frac{A}{4\pi R^{n_\mathrm{s}}}}\mathrm{\Gamma }\left({\displaystyle \frac{n_\mathrm{s}+2j+2}{2}}\right),`$ (4.82)
and the variance parameters are given by
$`\sigma _{j}^{}{}_{}{}^{2}(\theta _\mathrm{f})={\displaystyle \frac{A}{4\pi \theta _{\mathrm{f}}^{}{}_{}{}^{n_\mathrm{s}+2j+2}}}\mathrm{\Gamma }\left({\displaystyle \frac{n_\mathrm{s}+2j+2}{2}}\right){\displaystyle 𝑑\chi S_{K}^{}{}_{}{}^{2n_\mathrm{s}}(\chi )n^2(\chi )D^2(\chi )}.`$ (4.83)
With similar technique used in 3D density field,
$`C_m^{\alpha \beta }={\displaystyle \frac{1}{2^mm!}}\left({\displaystyle \frac{n+2}{2}}\right)^{1(\alpha +\beta )/2}\left({\displaystyle \frac{n+2}{2}}\right)_{(\alpha +m1)/2}\left({\displaystyle \frac{n+2}{2}}\right)_{(\beta +m1)/2}`$
$`\times F({\displaystyle \frac{n+\alpha +m+1}{2}},{\displaystyle \frac{n+\beta +m+1}{2}},m+1;{\displaystyle \frac{1}{4}}).`$ (4.84)
The equations (4.79)–(4.81) and the above equation finally give the values of $`C^{(a)}`$. The recursion relations of equations (4.34) and (4.35) simplify the result:
$`C^{(0)}(n_\mathrm{s})=C^{(1)}(n_\mathrm{s})=3F({\displaystyle \frac{n_\mathrm{s}+2}{2}},{\displaystyle \frac{n_\mathrm{s}+2}{2}},1;{\displaystyle \frac{1}{4}}){\displaystyle \frac{3}{2}}(n_\mathrm{s}+1E)F({\displaystyle \frac{n_\mathrm{s}+2}{2}},{\displaystyle \frac{n_\mathrm{s}+2}{2}},2;{\displaystyle \frac{1}{4}}),`$ (4.85)
$`C^{(2)}(n_\mathrm{s})=3F({\displaystyle \frac{n_\mathrm{s}+4}{2}},{\displaystyle \frac{n_\mathrm{s}+4}{2}},2;{\displaystyle \frac{1}{4}}){\displaystyle \frac{3(n_\mathrm{s}+33E)}{4}}F({\displaystyle \frac{n_\mathrm{s}+4}{2}},{\displaystyle \frac{n_\mathrm{s}+4}{2}},3;{\displaystyle \frac{1}{4}}).`$ (4.86)
These functions are independent on scale $`R`$ in the power-law case, but dependent on power-law index, $`n_\mathrm{s}`$. Again, $`C^{(0)}(n_\mathrm{s})`$ and $`C^{(1)}(n_\mathrm{s})`$ are identical only in the power-law case. Numerical values are given in Table 7, where $`E=3/7`$ is assumed.
For the power-law case, the skewness parameters of equation (4.63) reduces to the following simple form:
$`S^{(a)}(n_\mathrm{s})={\displaystyle \frac{{\displaystyle 𝑑\chi \left[S_K(\chi )\right]^{22n_\mathrm{s}}n^3(\chi )D^4(\chi )}}{\left\{{\displaystyle 𝑑\chi \left[S_K(\chi )\right]^{2n_\mathrm{s}}n^2(\chi )D^2(\chi )}\right\}^2}}C^{(a)}(n_\mathrm{s}).`$ (4.87)
It is interesting to compare the results with those of top-hat smoothing. According to Bernardeau (1995), the top-hat smoothing gives $`C^{(0)}(n_\mathrm{s})=36/7+3(n_\mathrm{s}+2)/2`$ for power-law case. Comparing this expression with our Table 7, they roughly agree with each other, and have similar behavior with spectral index. In detail, the Gaussian smoothing gives more or less larger values. For example, top-hat smoothing gives $`C^{(0)}=5.14,3.64,2.14`$ for $`n_\mathrm{s}=2,1,0`$, respectively, while Gaussian smoothing gives $`C^{(0)}=5.14,3.89,3.01`$, respectively. It is reasonable that the $`C^{(0)}`$ of top-hat smoothing differs from that of Gaussian smoothing for large spectral index, because the top-hat smoothing gathers more power on small scales than Gaussian smoothing.
### 4.5 Weak Lensing Field
The local convergence field of the weak lensing is commonly used for studying the large-scale structure of the universe (e.g., Kaiser, 1998; Bartelmann & Schneider, 1999). Assuming the situation where Limber’s equation (see Appendix B) and also the Born approximation (Kaiser, 1998) hold, the following correspondence between the convergence field $`\kappa `$ and the 3D density contrast $`\delta _{3\mathrm{D}}`$ is useful (e.g., Mellier, 1999):
$`\kappa (𝜽)={\displaystyle \frac{3}{2}}H_{0}^{}{}_{}{}^{2}\mathrm{\Omega }_0{\displaystyle _0^{\mathrm{}}}𝑑\chi ^{}S_{K}^{}{}_{}{}^{2}(\chi ^{})n(\chi ^{}){\displaystyle _0^\chi ^{}}𝑑\chi {\displaystyle \frac{S_K(\chi )S_K(\chi ^{}\chi )}{a(\chi )S_K(\chi ^{})}}\delta _{3\mathrm{D}}[\chi ^{},𝜽S_K(\chi ^{});\tau _0\chi ^{}],`$ (4.88)
where $`a(\chi )`$ is the scale factor at conformal lookback time $`\chi =\tau _0\tau `$. The above equation is reduced to exactly the same form as the projected field of equation (4.54), but with the substitution,
$`n(\chi )n_{\mathrm{wl}}(\chi )={\displaystyle \frac{3}{2}}{\displaystyle \frac{H_{0}^{}{}_{}{}^{2}\mathrm{\Omega }_0}{a(\chi )}}{\displaystyle _\chi ^{\mathrm{}}}𝑑\chi ^{}S_{K}^{}{}_{}{}^{2}(\chi ^{}){\displaystyle \frac{S_K(\chi ^{}\chi )}{S_K(\chi ^{})S_K(\chi )}}n(\chi ^{}).`$ (4.89)
However, one should note that this substitution is valid only under the assumption of the Born approximation. Although the effect of the Born approximation on the skewness is known to be weak (Bernardeau et al., 1997), the validity of the Born approximation in general situation has not been tested in detail (Mellier, 1999). There is some subtlety which could arise when various combinations of skewness are considered.
### 4.6 The Biases
The above expressions of the skewness parameters are for unbiased fields. The skewness parameters for biased fields are non-trivial. They depend on the details of the biasing scheme in the real universe which is poorly known so far. However, the skewness parameters are well-defined quantities, so that they are calculated from the first principle once the biasing scheme is given.
Perhaps, one of the simplest, yet non-trivial case is the local, deterministic biasing. In this case, we can follow Fry & Gaztanaga (1993) to obtain the perturbative expansion of the biasing:
$`\delta _\mathrm{g}=b\delta +{\displaystyle \frac{b_2}{2}}\left(\delta ^2\delta ^2\right)+\mathrm{},`$ (4.90)
where $`\delta _\mathrm{g}`$ and $`\delta `$ are the galaxy and mass density contrast, respectively. The spatial dimension is arbitrary, so that $`\delta _\mathrm{g}`$, $`\delta `$ can be functions in either 1D, 2D, or 3D space. The biased variance in lowest order is given by
$`\sigma _{0,\mathrm{g}}=b\sigma _0.`$ (4.91)
It is also straightforward to calculate skewness parameters. After some algebra, all skewness parameters are shown to transform in a same way:
$`S_\mathrm{g}^{(a)}={\displaystyle \frac{S^{(a)}}{b}}+{\displaystyle \frac{3b_2}{b^2}},`$ (4.92)
irrespective of spatial dimensions, $`d`$, and of kinds of skewness parameters, $`a=0,1,2`$. In this framework, the parameters $`b`$ and $`b_2`$ are needed. A possible way to determine these parameters is to measure the variance $`\sigma _{\mathrm{g}}^{}{}_{}{}^{2}`$ and the skewness $`S_\mathrm{g}`$ from the observation. Theoretical models predict the variance $`\sigma _{\mathrm{m}}^{}{}_{}{}^{2}`$ and the skewness $`S_\mathrm{m}`$ of the mass distribution. We obtain biasing parameters by $`b=\sigma _\mathrm{g}/\sigma _\mathrm{m}`$ and $`b_2=b(bS_\mathrm{g}S_\mathrm{m})/3`$. The derivatives of the skewness, $`S_\mathrm{g}^{(1)}`$, $`S_\mathrm{g}^{(2)}`$, are then obtained from equation (4.92).
In case of level-crossing, and genus statistics as functions of scaled threshold $`\stackrel{~}{\nu }`$, the second-order corrections in equation (3.51)–(3.53) only depends on the difference of the skewness parameters times the variace. Quite remarkably, this type of second-order correction term does not depend on bias parameter at all:
$`\left[S_\mathrm{g}^{(a)}S_\mathrm{g}^{(0)}\right]\sigma _{0,\mathrm{g}}=\left[S^{(a)}S^{(0)}\right]\sigma _0`$ (4.93)
Thus, the second-order nonlinear corrections for level-crossing and genus statistics of locally biased field are exactly the same as that of unbiased mass density field. This can be considered as an advantage of the scaled threshold $`\stackrel{~}{\nu }`$ in these statistics, because the biasing is one of the most embarrasing uncertainty in the analysis of galaxy distribution. However, one should note this result is derived under the local biasing scheme. So far the locality of the bias is not guaranteed in general.
Stochastic argument (Dekel & Lahav, 1999) of the biasing is more complicated for derivative skewness $`S^{(1)}`$ and $`S^{(2)}`$ than for usual skewness $`S^{(0)}`$, because they involve the correlation between field derivatives. The phenomenological nature of the stochastic biasing requires many parameters which is not calculable from first principles. Therefore, the stochastic biasing sheme does not effectively work for our problem. More physical treatment of the biasing schemes of galaxy formation is needed, in which case the nonlocality of the bias could also be important (Matsubara, 1999).
## 5 IMPLICATIONS OF SECOND-ORDER RESULTS
The perturbative calculations offer valuable aspects of weakly non-linear evolution of the various statistics without laborious parameter survey by numerical simulations. In this paper, we obtain the lowest nonlinear corrections to relatively popular statistics of smoothed cosmic fields. Using these results, it is interesting to see how the weakly nonlinear effect tends to distort the Gaussian prediction for those exemplified statistics.
In Figure 1, various statistics for 3D density field is shown.
The amplitude of each statistic is appropriately normalized as we are interested in the deviation from the Gaussian prediction. If we neglect the normalization, Minkowski functionals of $`k=1,2,3`$ are equivalent to the statistics $`N_1`$, $`G_2`$, and $`G_3`$, respectively, so that they degenerate in this figure (we should note the sign of $`V_3^{(3)}`$ is inverted).
The rms $`\sigma _0`$, which is considered as a weakly nonlinear parameter, is set $`\sigma _0=0.3`$. A limit $`\sigma _00`$ corresponds to the prediction of the linear theory, which is given by thin solid lines in the figure. This linear prediction is equivalent to the Gaussian fluctuations, because we assume the initial density field is random Gaussian.
In general, the curves of statistics plotted against the direct density threshold, $`\nu `$ (dotted lines), exhibit considerable deviations from Gaussian predictions. The overall tendency does not depend much on the shape of the spectrum we consider here, i.e., power-law spectrum with index $`2`$, $`1`$, $`0`$, and CDM model with smoothing length $`R=4/\mathrm{\Gamma }`$, where $`\mathrm{\Gamma }`$ is the shape parameter of the CDM spectrum. They are consistent with the so-called meat-ball shift, which means that there are more isolated regions in a nonlinear field than in a Gaussian field for a fixed threshold. In fact, $`N_1`$, $`G_2`$, $`G_3`$ of high value of threshold, e.g., $`\nu 2`$, virtually correspond to the number of isolated regions, and each figure shows that the number is indeed increased by weakly nonlinear evolution.
The weakly nonlinear formula for the genus curve against the density threshold, $`G_3(\nu )`$, which was first derived by Matsubara (1994) has been compared with numerical simulations in literatures. Matsubara & Suto (1996) shows the good agreement of the analytic prediction with the simulations for various spectra. Colley et al. (2000) compared the prediction of the genus curve against direct $`\nu `$ with the simulated SDSS data. Unfortunately, in their published paper, they have transcribing errors, which incorrectly made the perturbation theory considerably disagree with their data. There are also ambiguity in their comparison on biasing which alter the values of skewness parameters. One can guess the biasing effect on skewness parameters by equation (4.90). They chose the peak particles as galaxies, and the linear biasing parameter is inferred as $`b1.3`$, but the nonlinear parameter $`b_2`$ is not obvious in their work. Some literatures indicate the skewness $`S^{(0)}`$ of peaks are roughly given by 1–2 (Watanabe, Matsubara & Suto, 1994; Plionis & Valdarnini, 1995) but for highly biased peaks $`b2`$. If we adopt $`b_2=0.5`$, the skewness is given by $`S^{(0)}=1.8`$ which is not unreasonable. If it is the case for their simulation, the perturbative prediction and their data completely agree with each other. The $`\chi ^2`$-value per degrees of freedom reduces to only 1.03 (private communication with W. N. Colley & D. H. Weinberg). Obviously, we have to further investigate the biasing effects in numerical simulations to obtain a conclusive result.
Most of the topological analyses of the previous work use the scaled threshold $`\stackrel{~}{\nu }`$. The dashed lines in Fig. 1 shows the corresponding curves. The deviations from the linear theory is dramatically reduced. This fact is empirically known by the analyses of numerical simulations (Gott, Melott & Dickinson, 1986; Gott, Weinberg & Melott, 1987). The reason for this reduction is mathematically due to the closeness of the values of skewness parameters $`S^{(a)}`$, ($`a=0,1,2`$), since all the terms of the nonlinear corrections in equations (3.51)–(3.53) depend only on $`S^{(a)}S^{(0)}`$ ($`a=1,2`$). For the hierarchical model of equation (4.1), they are exactly zero, that means there is not any (second order) nonlinear correction for the hierarchical model. Since the hierarchical model is known to roughly approximate the nonlinear evolution, it is not surprising that more realistic fields have only small corrections of nonlinearity if they are plotted against the volume-fraction threshold, $`\stackrel{~}{\nu }`$.
For power-law models, $`S^{(0)}`$ and $`S^{(1)}`$ are exactly the same. That makes the nonlinear correction for $`N_1`$ or $`V_1^{(3)}`$ exactly vanishes. Thus, the nonlinear corrections for other statistics are arisen by the difference between $`S^{(0)}`$ and $`S^{(2)}`$. For the CDM model, there still is a difference between $`S^{(0)}`$ and $`S^{(1)}`$, but it is relatively small as seen from the Table 2.
As for the topological statistics, $`G_2`$, $`G_3`$, $`V_2^{(3)}`$ and $`V_3^{(3)}`$, deviations of curves against $`\stackrel{~}{\nu }`$ from the linear theory prediction depend on the underlying spectrum through differences of skewness parameters. The weakly nonlinear effect for redder spectrum of $`n_\mathrm{s}=2`$ induces a sponge-like shift, that means the number of holes in isolated regions increases. On the other hand, the bluer spectrum of $`n_\mathrm{s}=0`$ indicates a meat-ball shift. These tendencies are qualitatively in agreement with the numerical results (e.g. Ryden et al., 1989; Melott et al., 1989; Park & Gott, 1991).
It is not trivial to estimate errors expected in observationally estimating statistics of smoothed cosmic fields. The observational errors mainly consists of the shot noise and the cosmic variace. The systematic comparison with simulations should be used for the error estimates. Unfortunately, previous earlier simulations does not have enough resolution to be quantitatively compared in weakly nonlinear regime. Canavezes et al. (1998) use much larger simulations than those earlier work, and gives the genus curve in weakly nonlinear regime in their Fig. 9. Relative depths of left and right troughs in the genus curve show slightly meat-ball shift, which is consistent with our prediction. Although their variation of smoothing length is rather limited, their plots ensures the slight meat-ball shift predicted by the perturbation theory can definitely be observed. Colley et al. (2000) plot the 2D genus curve, but they use the thin slice which suffers strongly nonlinear effects so that the quantitative comparison is not appropriate with weakly nonlinear results. They report the meat-ball shift in 2D genus, which is the same direction of the weakly nonlinear correction besides amplitude. We find there is stll neccesity of proper, systematic comparison between the perturbation theory and numerical simulations in weakly nonlinear regime.
The decisive factor between sponge-like shift and meat-ball shift is the sign and amplitude of $`S^{(2)}S^{(0)}`$, since $`S^{(1)}S^{(0)}`$ is exactly or approximately zero for wide range of models like power-law, or CDM-like models. If the factor $`S^{(2)}S^{(0)}`$ is positive, the meat-ball shift takes place. If this factor is negative, the sponge-like shift occurs. For power-law case, that factor is positive for $`n_\mathrm{s}>1.4`$, and is negative for $`n_\mathrm{s}<1.4`$. The amplitude of this factor times the amplitude of the nonlinearity parameter, $`(S^{(2)}S^{(0)})\sigma _0`$ determines the amplitude of the meat-ball shift in the genus curve.
Prominent shifts are seen around the two troughs of the genus curve at $`\stackrel{~}{\nu }=\pm \sqrt{3}`$. At these troughs, the factor $`S^{(1)}S^{(0)}`$ does not contribute to the genus curve since the accompanying factor $`H_3(\stackrel{~}{\nu })=\stackrel{~}{\nu }^3\stackrel{~}{\nu }`$ vanishes. Therefore, one can almost completely characterize the shifts by a factor $`(S^{(2)}S^{(0)})\sigma _0`$, which we call the ’genus asymmetry parameter’, around the troughs. We propose this genus asymmetry parameter as a theory-motivated parameter of the genus asymmetry. This factor is observationally determined by fitting the shape of the genus curve by a form,
$`G(\stackrel{~}{\nu })H_2(\stackrel{~}{\nu })+AH_3(\stackrel{~}{\nu })+BH_1(\stackrel{~}{\nu }),`$ (5.1)
where $`B`$ is the genus asymmetry parameter. In this fitting, $`A`$ is much smaller if the non-Gaussian features are from purely gravitational evolution of the initial Gaussian field, and if the power spectrum is smooth enough as CDM models. The genus asymmetry parameter for CDM models with shape parameter $`\mathrm{\Gamma }`$, normalized by $`\sigma _8`$ is listed for various smoothing length in Table 8.
Since the factor $`S^{(2)}S^{(0)}`$ increase with the spectral index, as seen from Table 3, this factor also increase with the smoothing length for CDM model on scales of interest. On the other hand, the factor $`\sigma _0`$ is a decreasing function of the smoothing length. As a result, the genus asymmetry parameter is not a monotonic function of the smoothing length. For example, the genus asymmetry parameter in the $`\mathrm{\Gamma }=0.2`$, $`\sigma _8=1`$ CDM model is approximately $`0.08`$ on wide range of smoothing length of $`10h^1\mathrm{Mpc}\stackrel{<}{}R\stackrel{<}{}20h^1\mathrm{Mpc}`$.
From the second order formula for the genus of equation (3.53), the genus at the troughs is proportional to $`2\pm \sqrt{3}(S^{(2)}S^{(0)})\sigma _0`$. Therefore, the fraction of the deviation from the Gaussian prediction at $`\stackrel{~}{\nu }=\pm \sqrt{3}`$ is $`\pm \sqrt{3}(S^{(2)}S^{(0)})\sigma _0/2`$. When the genus asymmetry parameter is $`0.08`$, this fraction is about $`\pm 7\%`$. Thus, the perturbation theory predicts that the values of genus at positive and negative troughs differ $`14\%`$ for the $`\mathrm{\Gamma }=0.3`$, $`\sigma _8=1`$ CDM model, and so on. This difference increases with the shape parameter $`\mathrm{\Gamma }`$. The degree of deviation is qualitatively consistent with the numerical result in Fig. 9 of Canavezes et al. (1998), although quantitative comparison is still difficult because of the noise in the simulation.
The statistics for the velocity field is plotted in Figure 2.
Since the the sign of the skewness parameters is negative in this case, the weakly nonlinear evolution of statistics against the density threshold $`\nu `$ indicates the sponge-like shift. In terms of the volume-fraction threshold $`\stackrel{~}{\nu }`$, on the other hand, meat-ball shifts are observed even for relatively bluer spectrum with $`n_\mathrm{s}0`$.
The 2D projected galaxy statistics are dependent on the selection function of galaxies and cosmological models. As an example, we assume the APM luminosity function (Laveday et al., 1992) for galaxies with B band magnitude limit $`m_{\mathrm{lim}}=19`$. The differential number count $`dN/dz(z)`$ for this sample is plotted in Figure 3.
The resulting mean redshift is $`z=0.12`$. The relation between the differential number count and the normalized mean number density $`n(\chi )`$ in comoving coordinates are given by
$`n\left[\chi (z)\right]={\displaystyle \frac{H(z){\displaystyle \frac{dN}{dz}}(z)}{\left\{S_K\left[\chi (z)\right]\right\}^2{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dN}{dz}}𝑑z}},`$ (5.2)
where
$`H(z)=H_0\sqrt{(1+z)^3\mathrm{\Omega }_0+(1+z)^2(1\mathrm{\Omega }_0\lambda _0)+\lambda _0},`$ (5.3)
$`\chi (z)={\displaystyle _0^z}{\displaystyle \frac{dz^{}}{H(z^{})}},`$ (5.4)
are the Hubble parameter and the comoving distance at redshift $`z`$, respectively. Once the selection function $`n(\chi )`$ is given by equation (5.2), the integration of equations (4.61)–(4.63), and the interpolation of tabulated values of $`C^{(a)}`$ in Table 6 give the skewness parameters. For power-law power spectra, the skewness parameters are given by the simpler integration of equation (4.87) and the values of $`C^{(a)}`$ in Table 7.
In Figure 4, 2-dimensional statistics, $`N_1`$, $`G_2`$, $`V_1^{(2)}`$ and $`V_2^{(2)}`$ are plotted with assumed cosmological parameters, $`\mathrm{\Omega }_0=0.3`$, $`\lambda _0=0.7`$, and APM luminosity function with limiting magnitude 19.
The nonlinear parameter is assumed as $`\sigma _0=0.2`$. For the CDM model, the shape parameter is assumed as $`\mathrm{\Gamma }=0.25`$, and we take the smoothing angle as $`\theta _\mathrm{f}=1^{}`$. With this smoothing angle the value $`\sigma _0=0.2`$ corresponds to the normalization $`\sigma _8=1.22`$. Basic features for these 2-dimensional statistics are the same as for 3-dimensional density field, except that the smoothing scale in CDM model we adopt corresponds to smaller scale than in the example of 3-dimensional density field of the Figure 1.
## 6 CONCLUSIONS
In this work, we comprehensively presented a basic formalism to treat the statistics of smoothed cosmic fields in perturbation theory. This formalism provides a methodology on evaluating how various statistics deviate from the prediction of simple random Gaussian fields. As long as the non-Gaussianity is weak, the behavior of the statistics caused by non-Gaussianity is predicted by our formalism, which enable us to quantitatively compare the statistical quantities and the source of the non-Gaussianity. This method is considered as an extension of the Edgeworth expansion, which has been proven to be useful in various fields of research, as long as the non-Gaussianity is weak. In this paper, we derive useful formulas and relations focusing on application of the second-order perturbation theory to various cosmic fields.
Several examples of statistics of cosmic fields in second-order perturbation theory are investigated in datail, including level-crossing statistics, 2D and 3D genus statistics, 2D extrema statistics, and the Minkowski functionals, which are extensively used in cosmology. More complicated statistics, such as 2D and 3D density peaks, can also be calculated, although they are more tedius.
A particular interest in cosmology of our method is in the application to the cosmic fields. Even if the cosmic field was random Gaussian at the initial stage, the gravitational evolution induces the non-Gaussianity. The gravitational instability is a well-defined process, so that we can evaluate the non-Gaussianity without any ambiguity when the evolution remains in the quasi-linear regime provided that the biasing from the non-gravitational process is simple enough on large scales. Therefore, we performed the perturbative analysis to obtain the necessary skewness parameters based on the gravitational instability theory. We considered the 3D density field, the 3D velocity field, the 2D projected density field. In the application of second order theory to various statistics of smoothed cosmic fields, three types of skewness parameters are commonly useful, i.e., $`S^{(0)}`$, $`S^{(1)}`$, an $`S^{(2)}`$. Extensive calculations of these parameters for various cosmic fields are one of the new results of this paper. It would be true that other skewness parameters are needed when other complex statistics are considered. Such other parameters, if needed, are similarly calculated by the method we outlined in this paper.
We find the lowest order deviations from the Gaussian predictions of various statistics of smoothed cosmic fields depend only on the differences of the skewness parameters when we use a threshold $`\stackrel{~}{\nu }`$, which is rescaled by a volume-fraction of the smoothed field. This rescaling makes the lowest deviations much smaller than in the case of direct threshold $`\nu `$. This is because three skewness parameters $`S^{(a)}`$ take similar values if it is arisen from the gravitational evolution. For the phenomenological hierarchical model, these parameters are identical. In this case, the weakly nonlinear correction of the statistical quantities in terms of the volume-fraction threshold vanishes. When evaluated by the second-order perturbation theory of density fluctuations, those three types of skewness parameter are still close to each other. This fact explains the smallness of the deviations from Gaussian predictions of statistical quantities like genus, level-crossing, or Minkowski functionals when the rescaled threshold by volume-fraction is used.
We discussed small, but detectable deviations from Gaussian prediction of the 3D genus curve against rescaled threshold in detail. In the framework of the second-order perturbation theory, a prominent deviation of the genus curve occurs at the two troughs of the curve. Relative depths of left and right troughs in the genus curve show slightly meat-ball shift for CDM-like models. We found the degree of this asymmetry is proportional to the combination $`(S^{(2)}S^{(0)})\sigma _0`$. We call this factor as a genus asymmetry parameter, which we propose as theory-motivated parameter that characterize the asymmetry of the genus curve. The genus asymmetry parameter can observationally be obtained by fitting the genus curve by Hermite polynomials as equation (5.1). Qualitative comparison with the numerical simulations in literatures suggests that such asymmetry can actually be observable. We have not estimated over what dynamic range in nonlinearity parameter like $`\sigma _0`$, $`\nu \sigma _0`$, or $`\stackrel{~}{\nu }\sigma _0`$ are each perturbative expressions valid. This estimation requires a systematic comparison with large N-body simulations, which is beyond the scope of this paper, and will be given in a subsequent paper of the series.
In principle, any order in the perturbation theory can be calculated as further as one would like. Although the computation of the higher-order theory becomes more and more tedious, the necessity of the comparison with large-scale cosmological observations is a good reason to perform such computation as further as we can. One of the spectacular example of the detailed comparison between perturbation theory and observations is the fine structure constant in quantum electrodynamics (e.g., Kinoshita, 1996). Our analysis in this paper will be extended to the third-order perturbation theory in a subsequent paper of the series. The present time is in an unique decade when the observations of cosmic fields are in unforeseen progress, like large-scale redshift surveys, detailed mapping of CMB fluctuations, gravitational lensing surveys, and so forth. Statistics of smoothed density field with higher-order perturbation theory will provide an unique method to analyze those high-precision data. The precision cosmology is undoubtedly providing clues to unlock the door to the origin of the universe.
I would like to thank M. Kerscher and B. Jain for discussions. I wish to acknowledge support from JSPS Postdoctoral Fellowships for Research Abroad, and from the Ministry of Education, Culture, Sports, Science, and Technology, Grant-in-Aid for Encouragement of Young Scientists, 13740150, 2001.
## Appendix A USEFUL GAUSSIAN INTEGRALS
In this appendix, we give Gaussian integrals which are useful in this paper. In the following, $`H_n`$ is the Hermite polynomials,
$`H_n(\nu )=e^{\nu ^2/2}\left({\displaystyle \frac{}{\nu }}\right)^ne^{\nu ^2/2},`$ (A1)
and we further employ the notation
$`H_1(\nu )e^{\nu ^2/2}{\displaystyle _\nu ^{\mathrm{}}}𝑑\nu e^{\nu ^2/2}=\sqrt{{\displaystyle \frac{\pi }{2}}}e^{\nu ^2/2}\mathrm{erfc}\left({\displaystyle \frac{\nu }{\sqrt{2}}}\right).`$ (A2)
Several Hermite polynomials are
$`H_0(\nu )=1,H_1(\nu )=\nu ,H_2(\nu )=\nu ^21,H_3(\nu )=\nu ^33\nu ,`$
$`H_4(\nu )=\nu ^46\nu ^2+3,H_5(\nu )=\nu ^510\nu ^3+15\nu .`$ (A3)
The Hermite polynomial at zero is given by
$`H_n(0)`$ $`=`$ $`\{\begin{array}{cc}0,\hfill & (n:\mathrm{odd})\hfill \\ (1)^{n/2}(n1)!!\hfill & (l:\mathrm{even})\hfill \end{array}`$ (A6)
$``$ $`h_n.`$ (A7)
We generalize the above definition of $`h_n`$ to the case $`n<0`$ by interpreting $`(n1)!!`$ as the appropriate gamma function so that $`(1)!!=1`$, $`(3)!!=1`$, etc. For example, $`h_2=1`$, $`h_0=1`$, $`h_2=1`$, $`h_4=3`$, and so forth. In this appendix, we give useful Gaussian averages $`\mathrm{}_\mathrm{G}`$ of equation (2.20) for normalized cosmic field $`\alpha `$ defined by equation (2.5) and its spatial derivatives $`\eta _i=\alpha _{,i}`$. In the following, $`\eta `$ represents any one of the components $`\eta _i`$. First, concerning $`\alpha `$,
$`{\displaystyle \frac{^k\delta (\alpha \nu )}{\alpha ^k}}H_n(\alpha )_\mathrm{G}={\displaystyle \frac{e^{\nu ^2/2}}{\sqrt{2\pi }}}H_{k+n}(\nu ),`$ (A8)
$`{\displaystyle \frac{^k\theta (\alpha \nu )}{\alpha ^k}}H_n(\alpha )_\mathrm{G}={\displaystyle \frac{e^{\nu ^2/2}}{\sqrt{2\pi }}}H_{k+n1}(\nu ).`$ (A9)
Second, concerning $`\eta `$, in the notation of equation (A7),
$`{\displaystyle \frac{^l|\eta |}{\eta ^l}}_\mathrm{G}=\sqrt{{\displaystyle \frac{2}{\pi }}}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^{1l}h_{l2},`$ (A11)
$`{\displaystyle \frac{^l\delta (\eta )}{\eta ^l}}_\mathrm{G}=\sqrt{{\displaystyle \frac{1}{2\pi }}}\left({\displaystyle \frac{\sigma _1}{\sqrt{d}\sigma _0}}\right)^{l1}h_l.`$ (A12)
## Appendix B LIMBER’S EQUATION FOR BISPECTRUM
The correlation functions on projected sky is expressible by the 3-dimensional correlation functions. The explicit relation for the two-point correlation function is given by Limber’s equation (Limber, 1954). The Fourier-space version of Limber’s equation is given by Kaiser (Kaiser, 1998) and is somewhat simpler. His argument was generalized to higher-order correlation functions and their Fourier transforms (Scoccimarro, Zaldarriaga, & Hui, 1999; Buchalter, Kamionkowski, & Jaffe, 2000). Since the higher-order Limber’s equation in Fourier space was discussed in the context of weak lensing field in literatures, here we review the derivation in a way more useful in this paper. Let 2D projected field $`f`$ be the projection of a time-dependent 3D field $`F(𝒙;\tau )`$ along a light-cone:
$`f(𝜽)={\displaystyle 𝑑\chi S_K^{\mathrm{\hspace{0.17em}2}}(\chi )q(\chi )F(\chi ,𝜽S_K(\chi );\tau _0\chi )},`$ (B1)
where $`q(\chi )`$ is some radial weighting function and $`\chi `$ is the radial comoving distance, and $`\tau _0`$ is the conformal time at the observer. The past light-cone of the observer is specified by the equation $`\chi =\tau _0\tau `$. The comoving angular distance $`S_K(\chi )`$ is defined by equation (4.52).
In the following, we explicitly derive the relation for 3-point correlation function, and bispectrum. From the Limber’s equation, power spectrum is already derived by Kaiser (1998):
$`P_f(\omega )={\displaystyle 𝑑\chi S_K^{\mathrm{\hspace{0.17em}2}}(\chi )q^2(\chi )P_F(\frac{\omega }{S_K(\chi )};\tau _0\chi )},`$ (B2)
where $`P_f`$ and $`P_F`$ is the power spectrum of fields $`f`$ and $`F`$, respectively. We generalize this equation to the one for 3-point statistics below. The generalization of the following derivation to higher-order statistics is straightforward. The angular 3-point correlation function $`w_f^{(3)}`$ of $`f`$ is
$`w_f^{(3)}(𝜽_1,𝜽_2,𝜽_3)`$ $`=`$ $`{\displaystyle 𝑑\chi _1S_K^{\mathrm{\hspace{0.17em}2}}(\chi _1)q(\chi _1)𝑑\chi _2S_K^{\mathrm{\hspace{0.17em}2}}(\chi _2)q(\chi _2)𝑑\chi _3S_K^{\mathrm{\hspace{0.17em}2}}(\chi _3)q(\chi _3)}`$ (B3)
$`\times F(\chi _1,𝜽_1S_K\left(\chi _1\right);\tau _0\chi )F(\chi _2,𝜽_2S_K\left(\chi _2\right);\tau _0\chi )F(\chi _3,𝜽_3S_K\left(\chi _3\right);\tau _0\chi )`$
$``$ $`{\displaystyle 𝑑\chi S_{K}^{}{}_{}{}^{6}(\chi )q^3(\chi )𝑑\chi _1𝑑\chi _2\zeta _F(\chi _1,𝜽_1S_K(\chi _1);\chi _2,𝜽_2S_K(\chi _2);\chi ,𝜽_3S_K(\chi );\tau _0\chi )},`$
where $`\zeta _F(𝒙_1;\mathrm{};𝒙_3;\tau )`$ is the spatial 3-point correlation function of the field $`F`$ with the 3-point configuration $`(𝒙_1;𝒙_2;𝒙_3)`$ at conformal time $`\tau `$. According to the spirit of the Limber’s equation, we assume that $`S_K^{\mathrm{\hspace{0.17em}2}}(\chi )q(\chi )`$ is slowly varying compared to the scale of the fluctuations of interest and also that these fluctuations occur on a scale much smaller than the curvature scale. The equation (B3) is the generalization of the Limber’s equation to higher-order correlation functions.
Now we transform equation (3.18) to obtain the 2D bispectrum. We use the following convention of the Fourier transforms
$`\stackrel{~}{f}(𝝎)={\displaystyle d^2\theta f(𝜽)e^{i𝝎𝜽}},`$ (B4)
$`\stackrel{~}{F}(𝒌;\tau )={\displaystyle d^3xF(𝒙;\tau )e^{i𝒌𝒙}},`$ (B5)
the bispectrum $`B_f`$ of the 2D field $`f`$, and $`B_F`$ of the 3D field $`F`$ are defined by
$`\stackrel{~}{f}(𝝎_1)\stackrel{~}{f}(𝝎_2)\stackrel{~}{f}(𝝎_3)=(2\pi )^2\delta ^2(𝝎_1+𝝎_2+𝝎_3)B_f(\omega _1,\omega _2,\omega _3),`$ (B6)
$`\stackrel{~}{F}(𝒌_1;\tau )\stackrel{~}{F}(𝒌_2;\tau )\stackrel{~}{F}(𝒌_3;\tau )=(2\pi )^3\delta ^3(𝒌_1+𝒌_2+𝒌_3)B_F(k_1,k_2,k_3;\tau ),`$ (B7)
where $`\omega _i=|𝝎_i|`$ and $`k_i=|𝒌_i|`$. The Dirac’s delta function comes from the translational invariance of statistics. From the relations,
$`B_f(\omega _1,\omega _2,\omega _3)={\displaystyle d^2\theta _1d^2\theta _2v_f(𝜽_1,𝜽_2,\mathrm{𝟎})e^{i𝝎_1𝜽_1i𝝎_2𝜽_2}},(\omega _3=|𝝎_1+𝝎_2|),`$ (B8)
$`\zeta _F(𝒙_1,𝒙_2,𝒙_3;\tau )={\displaystyle \frac{d^3k_1}{(2\pi )^3}\frac{d^3k_2}{(2\pi )^3}B_F(k_1,k_2,k_3;\tau )e^{i𝒌_1(𝒙_1𝒙_3)+i𝒌_2(𝒙_2𝒙_3)}},`$ (B9)
the Fourier transform of the equation (B3) reduces to a simple equation,
$`B_f(\omega _1,\omega _2,\omega _3)={\displaystyle 𝑑\chi S_K^{\mathrm{\hspace{0.17em}2}}(\chi )q^3(\chi )B_F(\frac{\omega _1}{S_K(\chi )},\frac{\omega _2}{S_K(\chi )},\frac{\omega _2}{S_K(\chi )};\tau _0\chi )}.`$ (B10)
This is a bispectrum version of the Limber’s equation and the generalization of the Kaiser’s equation for power spectrum (B2).
## Appendix C Symbol Index
In Table 9, the quantites used in this paper are listed.
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# 1 Introduction
## 1 Introduction
The study of magnetic monopoles started with the Dirac’s paper in 1948 where for the first time the possibility of existence of this kind of particles in nature was analysed.
Some years later t’Hooft and Polyakov again found the magnetic monopoles but in a different context, the so called non-abelian monopoles had a scalar field as a source.
More recently Wali has studied the Bogomol’nyi equations and he again obtained the magnetic monopole type of solution. The minimization of the functional of energy gives us the conditions for analysis of Bogomol’nyi equations.
Finally Christiansen at all have studied the Bogomol’nyi equations in Gauge theory and they verify conditions on functional of energy which is suitable for obtaining finite energy for a static configuration of field (gauge and scalar field) in $`D=2+1`$ dimensions.
In our case the objective is to study the Einstein-gauge-field-Higgs-fermions system in $`D=3+1`$ dimensions and, following of Christiansen and Wali, we wish to know the contribution for magnetic field and the potential that minimizes the energy functional when we include fermions in the model. Second we verify the possibility of the existence of non abelian magnetic monopoles in $`D=3+1`$ dimension with fermion contributions.
The paper is written with the following outline:
1. First we consider the complete system Higgs-gauge-fermions and coupling between gravitation and scalar field. We have established equations of motion for systems with and without fermions.
2. We review some results of Wali at al for Bogomol’nyi equations and for energy functional when fermions and gravitation are not present.
3. Next we consider the contribution of fermions alone but without gravitation. There is no coupling between the scalar field and gravitation. We obtain here the magnetic field with fermionic contributions and the potential function that minizes the energy functional.
4. We verify if there is a structure for the magentic monopole when we consider the covariant divergence associated with the gauge field.
5. Finally we return to the original problem and consider the lagrangean with all fields and coupling between gravitation and scalar field. All the equations are solved for this case, but there remain two problems that have no solutions.
The first problem is about the new magnetic field and the potential function when we account for the fermion contributions and the background of gravitation.
The questions is: what is the new magnetic field and potential function that minimizes the functional of energy if we put fermions and gravitation?
Second, what is the meaning of the presence of spinors in the initial lagrangean? What is the physical interpretation for spinors here?
In reality we have been speaking about “fermions”. However in a classical problem, the correct description is “spinors”. The question is what is the meaning of the spinors in our problem?
Anyway we have a partial solution for this problem and it will appear elesewhere.
Let us start with a non Abelian Higgs-gauge-“fermions” and coupling between scalar field and gravitation in $`3+1`$ dimensions. We would like to verify the behaviour of magnetic monopole of t’Hooft-Polyakov type in a background of gravitation with presence of fermions.
We take the action given by Wali for a system of Einstein-Yang-Mills-Higgs in the form:
$$S=d^4x\sqrt{g}\left[\frac{1}{16\pi Gv^2}R\varphi ^2\frac{1}{4}\left(F_{\mu \nu }^a\right)^2+\frac{1}{2}\left(D_\mu \varphi ^a\right)^2\frac{\lambda }{4}\left(\varphi ^2v^2\right)^2\right]$$
(1.1)
and now we will consider the Dirac Lagrangean
$$=\overline{\psi }_\alpha ^i\left(i\gamma _{\alpha \beta }^\mu 𝒟_\mu m\delta _{\alpha \beta }\right)\psi _\beta ^i+\chi \overline{\psi }_\alpha ^i\varphi ^a\left(T^a\right)_{ij}\psi _\alpha ^j$$
(1.2)
for fermion contribution.
The matrix $`T^a`$ are hermitian with null trace. It describes the three generators of $`SU(2)`$ group, here $`a=1,2,3`$ in the $`N`$-dimensional representation. The indices $`i,j=1,\mathrm{}N`$, where $`N`$ represent the dimension of given irreducible representation of $`SU(2)`$. The $`\chi `$ is only a constant.
The signature of metric is given by $`(+)`$. The indices $`\mu ,\nu `$ take values from $`\underset{¯}{0}`$ to $`\underset{¯}{3}`$ and the indices $`i,j`$ assume values between $`1`$ and $`3`$. The indices $`a,b`$ vary with the representation of a given gauge group.
We define $`R`$ as the scalar curvature, $`F_{\mu \nu }^a`$ represents the stress field associated with the gauge field $`A_\mu ^a`$ given as
$$F_{\mu \nu }^a=_\mu A_\nu ^a_\nu A_\mu ^a+gf^{abc}A_\mu ^bA_\nu ^c$$
(1.3)
where $`f^{abc}`$ represent the structure constants of the group. The gauge covariant derivative associated with Higgs field $`\varphi ^a`$ is given by
$$D_\mu \varphi ^a=_\mu \varphi ^a+gf^{abc}A_\mu ^b\varphi ^c.$$
(1.4)
A new covariant derivative associated with the spinor field is necessary here. The fermion field derivative $`\psi _\beta ^i`$ is given by
$$𝒟_\mu \psi _\beta ^i=_\mu \psi _\beta ^iigA_\mu ^a\left(T^a\right)_{ij}\psi _\beta ^j$$
(1.5)
where $`g`$ is the coupling constant.
The term $`{\displaystyle \frac{\lambda }{4}}\left(\varphi ^2v^2\right)^2`$ represents the possibility of gauge symmetry breaking. We choose the unit system such that $`4\pi Gv^2=1`$. Thus the action is written as
$`S`$ $`=`$ $`{\displaystyle }d^4x\sqrt{g}[{\displaystyle \frac{1}{4}}R\varphi ^2{\displaystyle \frac{1}{4}}\left(F_{\mu \nu }^a\right)+{\displaystyle \frac{1}{2}}\left(D_\mu \varphi ^a\right)^2+U\left(\varphi _i\right)+`$
$`\overline{\psi }_\alpha ^i(i\gamma _{\alpha \beta }^\mu 𝒟_\mu m\delta _{\alpha \beta })\psi _\beta ^i+\chi \overline{\psi }_\alpha ^i\varphi ^a\left(T^a\right)_{ij}\psi _\alpha ^j].`$
Here we use $`U\left(\varphi _i\right)`$; $`i=1,2,3`$ in the form given by (1.1).
The equations of motions for Yang-Mills and Higgs fields when fermion are not considered are given respectively by
$`{\displaystyle \frac{1}{\sqrt{g}}}D_\mu \left(\sqrt{g}F^{\mu \nu a}\right)=gf^{abc}\varphi ^bD^\nu \varphi ^c,`$ (1.7)
$`{\displaystyle \frac{1}{\sqrt{g}}}D_\mu \left(\sqrt{g}D^\mu \varphi ^a\right)=\left[{\displaystyle \frac{R}{2}}+\lambda \left(\varphi ^2v^2\right)\right]\varphi ^a`$ (1.8)
and the Einstein equations are given by
$$G_{\mu \nu }=R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R=\frac{2}{\varphi ^2}T_{\mu \nu }$$
(1.9)
where the stress energy-momentum tensor for Einstein Yang-Mills-Higgs is written as
$$T_{\mu \nu }=g_{\mu \nu }TF_{\mu \rho }^aF_\nu ^{\rho a}+D_\mu \varphi ^aD_\nu \varphi ^a\frac{1}{2}\varphi ^2\left(R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R\right)$$
(1.10)
and
$$T=\frac{1}{4}F_{\mu \nu }^aF^{\mu \nu a}=\frac{1}{2}D_\mu \varphi ^aD_\nu \varphi ^a+\frac{\lambda }{4}\left(\varphi ^2v^2\right)^2$$
(1.11)
The contribution to the energy-momentum due the fermions is given by
$`T_{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{1}{2}}i(\overline{\psi }_\alpha ^i\gamma _{\nu \alpha \beta }𝒟_\mu \psi _\beta ^i+\overline{\psi }_\alpha ^i\gamma _{\mu \alpha \beta }𝒟_\nu \psi _\beta ^i)+g_{\mu \nu }(m\overline{\psi }_\alpha ^i\psi _\beta ^i\delta ^{\alpha \beta }+`$ (1.12)
$``$ $`ie_a^\chi \psi _\alpha ^i\gamma _{\alpha \beta }^a𝒟_\chi \psi _\beta ^i\chi \overline{\psi }_\alpha ^i\varphi ^a\left(T^a\right)_{ij}\psi _\alpha ^i)`$
Now with fermions, the equations of motion for fields $`A_\mu ^a`$ $`\varphi ^a`$ and $`\psi _\beta ^i`$ are given respectively by
$`_\beta F_{\alpha \beta }^d=gf^{acd}\left[F_{\alpha \nu }^aA_\nu ^c\left(D_\alpha \varphi ^a\right)\varphi ^c\right]+g\psi _\epsilon ^i\gamma _{\epsilon \gamma }^\mu \delta _{\mu \alpha }\left(T^d\right)_{ij}\psi _\gamma ^j,`$
$`\left({\displaystyle \frac{1}{\sqrt{g}}}_\mu \sqrt{g}\right)D_\mu \varphi ^a=D_\mu \left(D_\mu \varphi ^a\right)+\left[{\displaystyle \frac{R}{2}}+\lambda \left(\varphi ^2v^2\right)\right]\varphi ^a\chi \overline{\psi }_\alpha ^i\left(T^a\right)_{ij}\psi _\alpha ^j,`$
(1.13)
$`{\displaystyle \frac{1}{\sqrt{g}}}\left(_\mu \sqrt{g}\right)\overline{\psi }_\alpha ^ki\gamma _{\alpha \gamma }^\mu =\gamma _{\alpha \gamma }^\mu 𝒟_\mu \overline{\psi }_\alpha ^k+m\overline{\psi }_\alpha ^i\delta _\alpha ^\gamma \delta _k^i\chi \overline{\psi }_\gamma ^i\varphi ^a\left(T^a\right)_{ik},`$ (1.14)
## 2 Static Equations and Bogomol’nyi Conditions
####
We wish to get only static solutions (time independent) for the system described by (1.6). Thus, using the technique of Bogomol’nyi, with appropriate boundery conditions ,
$$D\varphi =0,\varphi ^2=v^2.$$
(2.1)
The gauge field without fermions is given by
$$A_i^a=\frac{1}{gv^2}\epsilon ^{abc}\varphi ^b_i\varphi ^c+\frac{1}{v}\varphi ^aA_i$$
(2.2)
where $`A_i`$ is arbitrary and $`F_{ij}^a`$ satisfy
$$\varphi ^aF_{ij}^a=\frac{\varphi ^2}{v}_{ij}$$
(2.3)
and the field $`_{ij}`$ is given by
$$_{ij}=\frac{1}{gv\varphi ^2}\epsilon ^{abc}\varphi ^a_i\varphi ^b_j\varphi ^c+_iA_j_jA_i.$$
(2.4)
Only the static abelian gauge field will survive for long distances. Then, we define the “magnetic field” $`_i`$ associated with the monopole for long distance as
$$^i=\frac{1}{2}\epsilon ^{abc}_{jk}=\frac{1}{2gv^3}\epsilon ^{ijk}\epsilon ^{abc}\varphi ^a_j\varphi ^b_k\varphi ^c+\epsilon ^{ijk}_jA_k.$$
(2.5)
The magnetic charge of the configuration is given by
$$g=\frac{1}{4\pi }d^3x_i^i=\frac{1}{8\pi gv^3}_{s_{\mathrm{}}^2}𝑑\sigma _i\epsilon ^{ijk}\epsilon ^{abc}\varphi ^a_j\varphi ^b_k\varphi ^c=\frac{n}{g}$$
(2.6)
where
$$n=\frac{1}{8\pi v^3}_{s_{\mathrm{}}^2}𝑑\sigma _i\epsilon ^{ijk}\epsilon ^{abc}\varphi ^a_j\varphi ^b_k\varphi ^c.$$
(2.7)
which is a topological number.
Taking the limit $`{\displaystyle \frac{\lambda }{g}}0`$, the functional of energy obtained from (1.6), without fermions for the flat spacetime is given by
$$\epsilon =d^3x\left(\frac{1}{4}F_{ij}^aF^{ija}\frac{1}{2}D_i\varphi ^aD^i\varphi ^a\right)$$
(2.8)
that is an energy of Higgs-Yang-Mills field.
For the case with the fermions, we have the energy functional written as,
$`\epsilon `$ $`=`$ $`{\displaystyle }d^3x[{\displaystyle \frac{1}{4}}F_{ij}^aF^{ija}{\displaystyle \frac{1}{2}}D_i\varphi ^aD^i\varphi ^a+\overline{\psi }_\alpha ^i(m\delta _{\alpha \beta }i\gamma _{\alpha \beta }^k𝒟_k)\psi _\beta ^i+`$ (2.9)
$``$ $`\chi \overline{\psi }_\alpha ^i\varphi ^a\left(T^a\right)_{ij}\psi _\alpha ^j+U\left(\varphi _i\right)].`$
We are not considering the coupling between the scalar and gravitation here.
¿From eq. (2.8) we can define the electric and magnetic fields as:
$`E_i^a=D_i\varphi ^a,`$ (2.10)
$`B_i^a={\displaystyle \frac{1}{2}}\epsilon _{ijk}F^{jka}`$ (2.11)
for the case when the fermions are not present and when the flat spacetime is considered $`(R=0)`$; in other words, the eq. (2.10) and (2.11) are good definitions for the case when we don’t consider the gravitation background. If we use the radiation gauge, $`A_0^a=0`$, and for the static solution $`D_0\varphi ^a=0`$, from eq. (1.3) and (1.4) it follows that
$`F_{oi}^a`$ $`=`$ $`_oA_i^a_iA_o^a+gf^{abc}A_o^bA_i^c,`$ (2.12)
$`F_{oi}^a=0=D_i\varphi ^a.`$
Then we have obtained a solution (null electricaly), because the electric field is zero in the whole space.
We have still the following inequality.
$`\epsilon `$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^3x\left[\left(E_i^a\right)^2+\left(B_i^a\right)^2\right]^2}`$ (2.13)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^3x\left(E_i^aB_i^a\right)\left(E_i^aB_i^a\right)}\pm {\displaystyle d^3xE_i^aB_i^a}\pm {\displaystyle d^3xE_i^aB_i^a}.`$
Using now the Bianchi indetity for $`F_{ij}^a`$, we can write
$$\pm d^3xE_i^aB_i^a=\pm d^3x_i\left(\frac{1}{2}\epsilon ^{ijk}F_{jk}^a\varphi ^a\right)$$
(2.14)
Comparing eq. (2.14) with eq. (2.3) the surface integral in eq. (2.14) asymptotically takes the value
$$\pm vd^3x_iB^i=\frac{4\pi nv}{g}.$$
(2.15)
On the other hand, if the Bogomol’nyi equations
$$E_i^a=B_i^a$$
(2.16)
are satisfied, then the functional of energy is definitely minimized. It was shown in , that this follows naturely from the relations (2.13) – (2.16).
Now, consider the case described by (2.9) where the fermions are present but without the background of gravitation.
We can still write suitably the energy functional (2.9) in terms of fields $`E_i^a`$ and $`B_i^a`$ with the same arguments. However, here it will not have the electric component due to the presence of the fermions. The magnetic field will be different since it shall have the contribution of fermions.
The scalars fields will be treated as a condensate of fermions. We shall define the following quantities:
$`\eta `$ $`=`$ $`m\overline{\psi }_\alpha ^i\psi _\beta ^i\delta _{\alpha \beta },`$
$`\xi `$ $`=`$ $`i\overline{\psi }_\alpha ^i\gamma _{\alpha \beta }^k𝒟_k\psi _\beta ^i,`$ (2.17)
$`\text{and}\mathrm{\Delta }`$ $`=`$ $`\chi \overline{\psi }_\alpha ^i\varphi ^a\left(T^a\right)_{ij}\psi _\alpha ^j.`$
We have created three new scalar fields $`(\eta ,\xi ,\mathrm{\Delta })`$. So, our Lagrangean (1.6) with the scalar curvature $`R=0`$ or its form eq. (2.8) with the fermions present, it is not possible anymore to obtain an energy functional that is satured by $`\lambda \varphi ^4`$ potential as in .
Using now the same prescription as in it is conjectured that it is possible to verify that eq. (2.9) in $`D=3+1`$ dimensions may be reduced to the following form:
$`\epsilon ={\displaystyle \frac{ev^2}{2}}\mathrm{\Phi }_{_{B_k^a}}\pm {\displaystyle \frac{1}{2e}}{\displaystyle }d\sigma _iJ_j^a\epsilon ^{ijk}\epsilon ^{abc}\varphi ^b_k\varphi ^c{\displaystyle }d^3x\{{\displaystyle \frac{1}{2}}(B_k^a_k\varphi ^a\sqrt{2U})^2`$
$`{\displaystyle }[{\displaystyle \frac{e}{2}}(v^2|\varphi |^2)(\eta +\xi +\mathrm{\Delta })\left(_k\varphi ^a\right)\sqrt{2U}]B_k^a_k\varphi ^a+{\displaystyle \frac{1}{2}}\left|(D_1\pm iD_2)\varphi ^a|^2\right\}.`$ (2.18)
Here $`\mathrm{\Phi }_{_{B_k^a}}`$ means the non-abelian magnetic flux. The second term is a surface term and goes to zero at infinity since all fields go to zero at infinity.
We need to discover what magnetic field $`B_k^a`$ and potential $`U(\varphi _i)`$ will minimize the energy functional or in other words, what are $`B_k^a`$ and $`U(\varphi _i)`$ which will saturate the functional, $`\epsilon `$?
If we use the duality condition as in the form
$$D_1\varphi ^a=iD_2\varphi ^a$$
(2.19)
and noting that we are considering only configurations with finite energy the surface integral of current vector is null.
In the Bogolomol’nyi limit given by eq. (2.18) the minimum of energy is obtained exactly if
$$\epsilon =\frac{ev^2}{2}\mathrm{\Phi }_{_{B_k^a}}.$$
(2.20)
Since we wish to saturate the functional $`\epsilon `$ the magnetic field now carrying the fermion’s contribution will be given by
$$_k^a=\frac{e}{2}\left(v^2|\varphi |^2\right)\left(\eta +\xi +\delta \right)\left(_k\varphi ^a\right).$$
(2.21)
and the new potential will be written as
$$U(\varphi ^a,\eta ,\xi ,\mathrm{\Delta })=\frac{e}{8}\left(|\varphi |^2v^2\right)^2\left(\eta +\xi +\mathrm{\Delta }\right)^2\left(_k\varphi ^a\right)^4$$
(2.22)
Clearly the potential is now of type $`\lambda \varphi ^8`$ because of fermion contribution.
With the definition of the covariant derivative given by
$$D_\mu =_\mu +i\stackrel{~}{g}A_\mu $$
we get
$$D_\mu F^{\mu \nu a}=_\mu F^{\mu \nu a}+i\stackrel{~}{g}[A_\mu ,F^{\mu \nu }]^a.$$
Thus, the covariant divergence is given as
$$D_\mu F^{\mu \nu a}=\stackrel{}{}B^a+\stackrel{~}{g}f^{abc}\stackrel{}{A}_b\stackrel{}{B}_c=0$$
(2.23)
or still in a compact form,
$$D_\mu F^{\mu \nu a}=\stackrel{}{D}\stackrel{}{B}=0$$
(2.24)
where
$$D=\stackrel{}{}+\stackrel{~}{g}f^{abc}\stackrel{}{A}.$$
(2.25)
Since $`B_k^a`$ is given by eq. (2.21) and $`|\varphi |^2`$ is written as
$$|\varphi |^2=\varphi ^a\varphi _a$$
(2.26)
a suitably ansatz for scalar field may be chosen such as
$$\varphi ^a=F\frac{r^a}{r}(r\mathrm{}).$$
(2.27)
Using eq. (2.27) and eq. (2.26) in equation eq. (2.23) or eq. (2.24) above we get
$$\stackrel{}{}\stackrel{}{B}_k^a\frac{g}{r^2}.$$
(2.28)
This gives us a structure of the magnetic monopole where $`g`$ represents the source of the magnetic field (non-Abelian magnetic field) whose source has the origin in the scalar field, gauge field and condensate of fermions $`(\varphi ^a,A_\mu ^a,\eta ,\xi ,\mathrm{\Delta })`$. It’s sufficient to choose or to fix the “$`F`$” function such as that form to obtain the Gauss law from eq. (2.28).
## The energy functional in a curved spacetime
####
Now the same problem is proposed in a curved spacetime.
The energy functional (static functional) that is obtained from eq. (1.6) in a curved spacetime is described by.
$`\epsilon `$ $`=`$ $`{\displaystyle }d^3x\sqrt{g}[{\displaystyle \frac{1}{4}}R\varphi ^2+{\displaystyle \frac{1}{4}}F_{ij}^aF^{ija}{\displaystyle \frac{1}{2}}D_i\varphi ^aD^i\varphi ^a+{\displaystyle \frac{\lambda }{4}}(\varphi ^2v^2)^2+`$ (2.29)
$`\overline{\psi }_\alpha ^i(m\delta _{\alpha \beta }i\gamma _{\alpha \beta }^k𝒟_k)\psi _\beta ^i\chi \overline{\psi }_\alpha ^i\varphi ^a\left(T^a\right)_{ij}\psi _\alpha ^j].`$
The energy functional can, in principle, be minimized with the gravitation as a background field only if we introduce a third covariante derivative associated with gravitation, but we need also to transfer the dynamics from the metric to vierbein and to write an appropriate covariante derivative carrying the spin connection such as:
$$\stackrel{~}{D}_\mu =_\mu +\frac{1}{8}i\left[\gamma _a\gamma _b\right]B_\mu ^{ab}$$
where
$$\gamma ^\mu =e_a^\mu (x)\gamma ^a$$
and
$$g_{\mu \nu }=e_\mu ^ae_\mu ^b\eta _{ab}$$
are Dirac’s matrices and metric respectively written in a local Lorentz coordinate system, $`e_\mu ^a`$ are the vierbeins and $`B_\mu ^{ab}(x)`$ are spin connections, completely determined by vierbeins. Now the gravitation is considered under flat space-time given by $`\eta _{ab}`$ in Minkowski space.
## 3 The basic Equations
####
We choose the spherically symmetric static metric as
$$ds^2=A^2(r)dt^2B^2(r)dr^2C^2(r)r^2d\mathrm{\Omega }^2.$$
(3.1)
The Einstein tensor and the scalar curvature are easily obtained,
$`G_{00}`$ $`=`$ $`{\displaystyle \frac{A^2}{B^2}}\left[{\displaystyle \frac{1}{r^2}}\left({\displaystyle \frac{B^2}{C^2}}1\right)+{\displaystyle \frac{2}{r}}{\displaystyle \frac{B^{}}{B}}+2{\displaystyle \frac{B^{}}{B}}{\displaystyle \frac{C^{}}{C}}{\displaystyle \frac{2C^{\prime \prime }}{C}}{\displaystyle \frac{6C^{}}{rC}}\left({\displaystyle \frac{C^{}}{C}}\right)\right],`$
$`G_{rr}`$ $`=`$ $`\left({\displaystyle \frac{C^{}}{C}}\right)^2+{\displaystyle \frac{2}{r}}{\displaystyle \frac{C^{}}{C}}+{\displaystyle \frac{1}{r^2}}\left(1{\displaystyle \frac{B^2}{C^2}}\right)+{\displaystyle \frac{2C^{}}{C}}{\displaystyle \frac{A^{}}{A}}+{\displaystyle \frac{2}{r}}{\displaystyle \frac{A^{}}{A}},`$ (3.2)
$`G_{\theta \theta }`$ $`=`$ $`{\displaystyle \frac{C^2r^2}{B^2}}\left({\displaystyle \frac{A^{\prime \prime }}{A}}+{\displaystyle \frac{1}{r}}{\displaystyle \frac{A^{}}{A}}{\displaystyle \frac{A^{}}{A}}{\displaystyle \frac{B^{}}{B}}+{\displaystyle \frac{C^{}}{C}}{\displaystyle \frac{A^{}}{A}}{\displaystyle \frac{1}{r}}{\displaystyle \frac{B^{}}{B}}{\displaystyle \frac{B^{}}{B}}{\displaystyle \frac{C^{}}{C}}+{\displaystyle \frac{C^{\prime \prime }}{C}}+{\displaystyle \frac{2}{r}}{\displaystyle \frac{C^{}}{C}}\right),`$
$`G_{\varphi \varphi }`$ $`=`$ $`sen^2\theta G_{\theta \theta },`$
$`R`$ $`=`$ $`{\displaystyle \frac{2}{r^2}}({\displaystyle \frac{1}{B^2}}{\displaystyle \frac{1}{C^2}})+{\displaystyle \frac{2}{B^2}}[{\displaystyle \frac{A^{\prime \prime }}{A}}+{\displaystyle \frac{2}{r}}{\displaystyle \frac{A^{}}{A}}{\displaystyle \frac{B^{}}{B}}{\displaystyle \frac{A^{}}{A}}{\displaystyle \frac{2}{r}}{\displaystyle \frac{B^{}}{B}}{\displaystyle \frac{2B^{}}{B}}{\displaystyle \frac{C^{}}{C}}+`$ (3.3)
$`{\displaystyle \frac{2C^{}}{C}}{\displaystyle \frac{A^{}}{A}}+{\displaystyle \frac{2C^{\prime \prime }}{C}}+{\displaystyle \frac{6C^{}}{rC}}+{\displaystyle \frac{C^{}}{C}}]`$
Then if we consider the eq. (1.10) and (1.11) together with the eq. (3.2) we take components of the energy momentum tensor for the global system of Einstein-gauge-Higgs-fermions.
The components are given by:
$`\stackrel{~}{T}_{00}`$ $`=`$ $`{\displaystyle \frac{\varphi ^2A^2}{B^2}}[{\displaystyle \frac{2C^{}}{C}}{\displaystyle \frac{A^{}}{A}}+{\displaystyle \frac{2}{r}}{\displaystyle \frac{A^{}}{A}}{\displaystyle \frac{A^{}}{A}}{\displaystyle \frac{B^{}}{B}}+{\displaystyle \frac{A^{\prime \prime }}{A}}{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{A^{}}{A}}\right)^2+{\displaystyle \frac{\lambda v^2}{4}}{\displaystyle \frac{\left(h^21\right)^2}{h^2}}+`$
$`{\displaystyle \frac{2h^{\prime \prime }}{h}}+\left({\displaystyle \frac{2h^{}}{h}}\right)+{\displaystyle \frac{2h^{}}{h}}({\displaystyle \frac{2C^{}}{C}}+{\displaystyle \frac{2}{r}}{\displaystyle \frac{B^{}}{B}})+{\displaystyle \frac{3A^{}}{A}}{\displaystyle \frac{h^{}}{h}}]+\tau _{00},`$
$`\stackrel{~}{T}_{rr}`$ $`=`$ $`\varphi ^2\left[{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{A^{}}{A}}\right)^2{\displaystyle \frac{A^{}}{A}}{\displaystyle \frac{h^{}}{h}}+{\displaystyle \frac{2B^{}}{B}}{\displaystyle \frac{h^{}}{h}}{\displaystyle \frac{\lambda v^2}{4}}B^2{\displaystyle \frac{\left(h^21\right)^2}{h^2}}\right]+\tau _{11},`$
$`\stackrel{~}{T}_{\theta \theta }`$ $`=`$ $`{\displaystyle \frac{C^2r^2}{B^2}}\varphi ^2[{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{A^{}}{A}}\right)^2{\displaystyle \frac{h^{\prime \prime }}{h}}\left({\displaystyle \frac{h^{}}{h}}\right)^2{\displaystyle \frac{h^{}}{h}}({\displaystyle \frac{2C^{}}{C}}+{\displaystyle \frac{2}{r}}{\displaystyle \frac{B^{}}{B}})+`$ (3.4)
$`{\displaystyle \frac{\lambda v^2}{4}}{\displaystyle \frac{B^2\left(h^21\right)^2}{h^2}}]+\tau _{22},`$
$`\stackrel{~}{T}_{\varphi \varphi }`$ $`=`$ $`sen^2\theta \stackrel{~}{T}_{\theta \theta }+\tau _{33}`$
where $`h=h(r)={\displaystyle \frac{\varphi ^a}{vr^a}}`$ and $`\tau _{00},\tau _{11},\tau _{22},\tau _{33}`$ are the components of the energy-momentum tensor for fermions given by
$`\tau _{00}`$ $`=`$ $`i\overline{\psi }_\alpha ^i\gamma _{0\alpha \beta }𝒟_0\psi _\beta ^i+g_{00}\left(m\overline{\psi }_\alpha ^i\psi _\beta ^i\delta _\beta ^\alpha ie_a^\chi \psi _\alpha ^i\gamma _{\alpha \beta }^a𝒟_\chi \psi _\beta ^i\chi \psi _\alpha ^i\varphi ^a\left(T^a\right)_{ij}\psi _\alpha ^j\right),`$
$`\tau _{11}`$ $`=`$ $`i\overline{\psi }_\alpha ^i\gamma _{1\alpha \beta }𝒟_1\psi _\beta ^i+g_{11}\left(m\overline{\psi }_\alpha ^i\psi _\beta ^i\delta _\beta ^\alpha ie_a^\chi \psi _\alpha ^i\gamma _{\alpha \beta }^a𝒟_\chi \psi _\beta ^i\chi \overline{\psi }_\alpha ^i\varphi ^a\left(T^a\right)_{ij}\psi _\alpha ^j\right),`$
$`\tau _{22}`$ $`=`$ $`i\overline{\psi }_\alpha ^i\gamma _{2\alpha \beta }𝒟_2\psi _\beta ^i+g_{22}\left(m\overline{\psi }_\alpha ^i\psi _\beta ^i\delta _\beta ^\alpha ie_a^\chi \psi _\alpha ^i\gamma _{\alpha \beta }^a𝒟_\chi \psi _\beta ^i\chi \overline{\psi }_\alpha ^i\varphi ^a\left(T^a\right)_{ij}\psi _\alpha ^j\right),`$
$`\tau _{33}`$ $`=`$ $`i\overline{\psi }_\alpha ^i\gamma _{3\alpha \beta }𝒟_3\psi _\beta ^i+g_{33}\left(m\overline{\psi }_\alpha ^i\psi _\beta ^i\delta _\beta ^\alpha ie_a^\chi \psi _\alpha ^i\gamma _{\alpha \beta }^a𝒟_\chi \psi _\beta ^i\chi \overline{\psi }_\alpha ^i\varphi ^a\left(T^a\right)_{ij}\psi _\alpha ^j\right)`$ (3.5)
where $`e_a^\chi `$ imply the tetrads or vierbeins.
This set of equations together with the Bogomol’nyi equations are the basic equations for our system, Einstein-gauge-Higgs-fermions.
Now the next step is to find solutions for the Einstein equations with the objective to find the monopoles appearing in this gravitational background when the fermions are present. But on the other hand when we have interaction between spinors and gravitation field the unique way to consider that coupling in our case is only if we introduce a local Lorentz coordinate system. Thus, we do not have curved space-time anymore.
## Conclusions
####
We have analysed two types of systems. One of them with Yang-Mills and scalar fields in flat spacetime and other consisting of the Yang-Mills-scalar field and “fermions”.
We treated “the fermions” as a condensate of scalars fields. In this case the magnetic field that saturated the energy functional has contributions from fermions. The potential which minimized the same functional is of the kind $`\lambda \varphi ^8`$ and thus different from whose potential is of type $`\lambda \varphi ^4`$ for $`2+1`$ dimension case.
For the case of Yang-Mills-scalar field the Bogomol’nyi equations have a simple solution in flat space time. In the present case our conjecture in eq. (2.18) gives us the non abelian magnetic field much more complicated and potential with fermion’s contribution. A structure of magnetic monopoles can be found with a field given by eq. (2.28).
Finally we have assumed the system Einstein-Higgs-scalar field-gauge and we have obtained all the equations for this system but solutions to the Einstein equations are lacking. Magnetic monopoles appear in a new context and it’s still necessary to find in a new magnetic field $`B_k^a`$ and potential $`U(\varphi _i)`$ that will be suitable for minimising the functional of energy when the background gravitational field is considered. At this point we have a conflict between a curve space-time and a local lorentz manifold. On the first way we have the the complete set of equations plus Bogolmo’ni equations being the basic equations for Einstein-gauge-Higgs-fermions system- and, in principle, the energy functional could be minimized with the graviation field as a background, but the other hand because the spinors we need to introduce a local coordinates system and to proceed the minimization of the functional of energy via vierbeins and spin connection. This will be a part of a separate analysis.
### Acknowledgements:
####
I would like to thank the Department of Physics, University of Alberta for their hospitality. This work was supported by CNPq (Governamental Brazilian Agencie for Research.
I would like to thank also Dr. Don N. Page for his kindness and attention with me at Univertsity of Alberta.
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# The HST Survey of BL Lacertae Objects. III. Morphological Properties of Low-Redshift Host Galaxies
## 1 Introduction
The Hubble Space Telescope (HST) snapshot imaging survey of BL Lac objects<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555. has produced a large, homogeneous data set of over one hundred high-resolution images for the study of the environments of BL Lac nuclei. The targets were drawn from various samples and include sources in a wide range of redshift ($`z`$ from 0.03 to $``$ 1.4) as well as objects with unknown $`z`$. Results from the survey are given in Scarpa et al. (2000; hereafter, Paper I), and Urry et al. (2000; hereafter, Paper II). Previous HST imaging of BL Lacs had been done only for a small number of sources (see Falomo et al. 1997; Jannuzi et al. 1997; Urry et al. 1999) but demonstrated the capabilities of HST for studying the environs of BL Lacs.
The HST snapshot images have two main advantages with respect to ground-based images (e.g., Falomo 1996; Wurtz, Stocke, & Yee 1996; Falomo & Kotilainen 1999): an order of magnitude improved spatial resolution and very good homogeneity of data. On the other hand the snapshot images were obtained with relatively short exposure times (typically minutes) and therefore are not very deep. This sets limits to the detection of very faint extended features, a well-known problem for HST images (see Falomo & Kotilainen 1999 for a direct comparison of HST and ground-based images of BL Lacs).
Our first papers on the overall properties of the BL Lac objects derived from HST data focused on the determination of the absolute magnitude, scale-length, and type of the host galaxies, as well as the luminosity of the nuclear source (see Paper II). This was based on accurate decomposition of the azimuthally averaged luminosity radial profiles which allowed the study of the sources down to surface brightness $`\mu _R`$25 mag arcsec<sup>-2</sup>. Additional morphological properties of the surrounding nebulosity, such as ellipticity, isophotal twisting, isophote shape (disky versus boxy), presence of dust, etc., can only be derived from a full two-dimensional approach, which is possible for a subset of bright nearby host galaxies.
For elliptical galaxies it has been shown that isophote shapes are related to internal kinematics and to the radio and X-ray emission properties (e.g. Bender, Döbereiner, & Möllenhoff 1987). Moreover significant distortions of the isophotes and/or twisting may be produced by strong galaxy interactions and/or merger events. It is therefore of interest to investigate these observable features in the galaxies hosting BL Lac nuclei in order to understand the (possible) role of galaxy structure on the nuclear activity. This kind of analysis has been performed so far only for a small number of nearby BL Lac objects (e.g., Falomo 1991,1996; Heidt et al. 1999) because the presence of the bright nucleus often hinders a fully detailed morphological study.
In this paper we present the results of a two-dimensional analysis for 30 nearby ($`z0.2`$) snapshot BL Lacs, i.e., those with the highest signal-to-noise ratios and the largest extent (see Table 1). For 25 of the 110 objects observed in the HST snapshot survey (Paper I) the redshift is unknown. Of these 17 are spatially unresolved and are therefore most probably at $`z>0.2`$. For the remaining 8 resolved objects with unknown redshift, based on the apparent magnitudes of the nebulosities (see Paper I) and assuming the luminosity of the host is consistent with that of other snap BL Lacs (see Fig. 1), one expects that at most a few can be at $`z<0.2`$. Therefore the objects discussed here represent practically the totality of $`z<0.2`$ BL Lacs from our original sample of 110 objects.
For these low-$`z`$ objects it is possible with HST to investigate features and structures in the host galaxy that are undetectable with ground-based observations. The host galaxy can be investigated down to less than 1 kpc from the nucleus (corresponding to 0$`^{\prime \prime }`$.3 at $`z=0.2`$). In particular one can search for slight off-centering of the nucleus with respect to the galaxy, determine the presence of sub-components in the host galaxy, investigate the detailed shape of the isophotes, and search for dusty features and close companions.
Because the two broad components in BL Lac spectral energy distributions (SEDs) have peak power outputs ($`\nu L_\nu `$) at wavelengths that increase systematically with luminosity (e.g., Ulrich, Maraschi & Urry 1997), BL Lacs found in radio and X-ray surveys have generally different SEDs. Radio-selected objects tend to have SEDs peaking at infrared-optical wavelengths and in the MeV-GeV gamma-ray band (low-frequency-peaked BL Lacs, or LBL) and exhibit luminosities approaching those of quasars. In contrast, X-ray-selected BL Lacs have SEDs peaking at UV-X-ray wavelengths and again at TeV energies (high-frequency-peaked BL Lacs, or HBL) and are generally less luminous. Of the 30 sources in the present study, 6 are of LBL type , the rest being HBL (Padovani & Giommi 1995). In the following analysis we have not addressed whether morphological properties are different for the two types of BL Lacs, having found from our earlier analysis that the overall host galaxy characteristics do not depend on type (Paper II).
In § 2 we briefly describe the observations and our data analysis. § 3 gives the results on the host galaxies properties and shows the comparison with non-active ellipticals and radio galaxies. Finally in § 4 we summarize the results and draw the main conclusions from this study. Throughout the paper we used H<sub>0</sub> = 50 km s<sup>-1</sup> kpc<sup>-1</sup> and q$`{}_{0}{}^{}=0`$.
## 2 Observations and Data Analysis
The snapshot observations of BL Lacs have been described in detail in Paper I, therefore only the main steps are reviewed here. All objects considered here were observed with the WFPC2 camera of the HST in the F702W filter. Each final image was obtained from the combination of typically three exposures for a total integration time of about 5-20 minutes (see Table 1 in Paper I for details). The Point Spread Function (PSF) used is a combination of a model derived from Tiny Tim package (Krist 1995) plus a diffuse exponential tail that accounts for the first order scattered light (see Urry et al. 1999 and Paper I for details). Photometric calibration was taken from Holtzman et al. (1995).
We performed a two-dimensional surface photometry analysis using an interactive numerical mapping package (AIAP; Fasano 1994), which produces an isophotal map of the image. This map was masked in order to avoid regions contaminated by companion objects, diffraction spikes, and any other extraneous features visible in the images. The isophotes were then fitted with ellipses down to $`\mu _R`$ = 22–23 mag arcsec<sup>-2</sup> (see Fig. 2 for an example). There are five geometric free parameters per isophote: semi-major axis ($`a`$), center position, ellipticity ($`ϵ`$), and position angle (PA), which allow us to characterize the morphological and the photometric properties of the galaxy.
From this analysis we then derived the average position angles and ellipticity as a function of the generalized radius, $`r=a\times (1ϵ)^{1/2}`$, and investigated their dependence on $`r`$ (isophote twisting and ellipticity profiles). In addition we examined the possible displacement of the centers of isophotes and the centering of the nucleus with respect to the galaxy.
## 3 Results
### 3.1 Luminosity Profiles and Two Dimensional Modeling
For each object we derived a luminosity profile from the two-dimensional isophotal analysis and compared with that derived from the azimuthal average reported in Paper I. This was possible down to $`\mu _R`$ 22–23 mag arcsec<sup>-2</sup> since the fainter outer regions are too noisy to be investigated with the isophotal map. For all sources we find that the luminosity profile derived from the fit with ellipses of the isophotes is always in excellent agreement with that extracted from the azimuthal average (see example in Fig. 3).
We also fitted the isophotes with a two-dimensional galaxy model (either a de Vaucouleurs law or an exponential disk) convolved with the PSF plus a point source modeled by a scaled PSF. In all cases the de Vaucouleurs law gives a better fit to the data than the exponential disk model, in agreement with Papers I and II. This result is supported by the lack of spiral structures seen in the WFPC2 images and by the diskiness of the isophotes (see discussion in the following paragraphs). Absolute magnitudes and scale lengths of the host galaxies derived from the two-dimensional fitting are also in excellent agreement with the results of one-dimensional profile fitting (Paper I), and are consistent with the results from our deep HST images for a small number of BL Lac objects (Falomo et al. 1997).
### 3.2 Ellipticity
Since the ellipticity $`ϵ`$ is in general dependent on the distance from the galaxy center, we report its value at the effective radius. However, because of the relatively high surface brightness limit, in some objects the ellipticity at the effective radius is not measurable and in these cases we took the maximum ellipticity.
The values of $`ϵ`$ are reported in Table 1. The error on observed $`ϵ`$ for isophotes far from the nuclear regions is typically 0.02 (see Fasano & Bonoli 1990 for discussion). We show in Figure 4 the distribution of ellipticity for our 30 low-redshift BL Lacs compared with that of 200 nearby normal (radio-quiet) ellipticals (Fasano & Vio 1991) and of a sample of 79 low-redshift radio galaxies (Govoni et al. 2000). Both normal ellipticals and radio galaxies exhibit the same ellipticity distribution: $`<ϵ>_{Ell}`$= 0.22 $`\pm `$ 0.13 and $`<ϵ>_{RG}`$= 0.21 $`\pm `$ 0.12. The average ellipticity for BL Lac hosts, $`<ϵ>_{BLL}`$ = 0.16 $`\pm `$ 0.09, is slightly smaller than that of non-radio ellipticals and radio galaxies but is still consistent with being drawn from the same population ($`P_{KS}0.1`$).
### 3.3 Isophotal Displacement and Nuclear De-Centering
Displacement of isophotes with respect to the center may represent either a global distortion of the host galaxy, due to recent tidal interactions with close companions (e.g., Aguilar & White 1986), or a consequence of gravitational microlensing (Ostriker and Vietri 1985). In the first case one would expect non-concentric isophotes to exhibit an asymmetry of distribution. To quantify this asymmetry we have computed the dimensionless parameter $`\delta \sqrt{(X_cX_o)^2+(Y_cY_o)^2}/r`$, where $`X_c,Y_c`$ are the centers of isophotal ellipses, $`X_o,Y_o`$ is the location of the nucleus, and $`r`$ is the radial distance to the particular elliptical isophote (defined earlier). Since $`\delta `$ may slightly change with $`r`$ we took the value at the effective radius R<sub>e</sub> of the galaxy or, if R<sub>e</sub> is beyond our measurements, we took the average value excluding the values at $`r<1`$ arcsec.
In Figure 5 we show the distribution of $`\delta `$ derived from our sample of low-redshift BL Lacs compared with that of radio galaxies from Govoni et al. (2000). The average value for our sample is: $`<\delta >_{BLL}=0.027\pm 0.01`$ (median $`\delta `$ is 0.030), marginally consistent with zero but suggestive of small distortions in the galaxy isophotes. This is similar to the value reported for low $`z`$ radio galaxies, $`<\delta >_{RG}=0.03\pm 0.04`$ (Govoni et al. 2000), using a homogeneous procedure, while for 40 normal ellipticals observed by Sparks et al. (1991) from the ground, Colina & de Juan (1995) give $`<\delta >_{Ell}`$ = 0.02 (see in Govoni et al. discussion and comparison with other samples of radio galaxies). This is consistent with the idea that galaxies hosting BL Lac nuclei and/or radio sources have indistinguishable structures and have undergone a similar degree of interaction.
In the microlensing scenario (Ostriker & Vietri 1985, 1990), what appears to be the host galaxy is actually an intervening galaxy whose stars are microlensing a background quasar. Since the alignment will not be perfect the galaxy will generally be off-center with respect to the point source, typically by $`0.5`$ arcsec (Merrifield 1992). To investigate this possibility we computed for each object the limit of off-centering, $`\mathrm{\Delta }=|(x,y)_{PSF}(x,y)_{gal}|`$, comparing the PSF location with that of the galaxy (derived by averaging of isophotal centers after subtraction of a scaled nuclear point source and exclusion of the central circle of radius 0.5 arcsec). With HST images the off-centering can be computed with unprecedented accuracy given the significantly better spatial resolution. The centers of nucleus and the host galaxy have a typical uncertainty of 0.2 and 0.4 pixels, respectively, so the net uncertainty is typically $`0.05`$ arcsec.
None of the objects in our (low-redshift) sample show significant ($`>2\sigma `$) displacement of the point source with respect to the galaxy (values of $`\mathrm{\Delta }`$ are given in Table 1); the observed displacements are always $`<0.1`$ arcsec and therefore consistent with zero off-centering. Since the low-redshift BL Lacs were the prime candidates for microlensing according to Ostriker & Vietri (1985), our result strongly rules out the microlensing hypothesis for BL Lac objects in general.
### 3.4 Isophotal Twisting
For each host galaxy having ellipticity larger than 0.15 we have estimated the isophote twisting, $`\mathrm{\Delta }\theta `$, computing the difference of position angles over the whole range of observed surface brightness (see Table 1). For the others the small ellipticity introduces large errors on the position angle so that it is not possible to get a reliable measurement of the isophotal twist.
The largest observed twist is $`30^{}`$ in 1853+671; this value is however uncertain because of the relatively small ellipticity and short exposure time. We found five other objects (corresponding to $``$20% of the sample) with isophote twists larger than 15. This is consistent (see Fig. 6) with what was found for a sample of 43 isolated (normal) ellipticals (Fasano & Bonoli 1989) and for 79 low-redshift radio galaxies (Govoni et al. 2000) using the same method to derive $`\mathrm{\Delta }`$PA.
A small amount of isophote twisting may be explained simply by the triaxiality of the galaxies, while larger twists may be due to the effects of tidal interaction with companion galaxies (Kormendy 1982). The lack of large twists in our sample is consistent with the overall smoothness (unperturbed shape) of the host galaxies.
### 3.5 Isophotal Shapes: Disky and Boxy
To check for small deviations from purely elliptical isophotes we analyzed the amplitude of the fourth cosine component, $`C_4`$, of the Fourier fit to the isophotes (e.g., Bender & Saglia 1998). A significant positive value of $`C_4`$ corresponds to disky-shaped isophotes while a negative value indicates a boxy-shaped structure.
We found that 80% of the sources show a $`C_4`$ amplitude (at the effective radius) smaller than 1% while no object exhibits $`|C_4|`$ larger than 3%. The distribution of $`C_4`$, shown in Figure 7, is symmetric around zero (mean value $`<C_4>_{BLL}=0.03`$%, rms dispersion 0.7%) and since the size of the error is around 0.5% we conclude that in these galaxies there is no systematic deviation of isophotal shape from a simple ellipse. Similar distributions are observed in a sample of low-redshift radio galaxies using the same method ($`<C_4>_{RG}=0.04`$% $`\pm 1.2`$%; Govoni et al. 2000) and also for radio galaxies in rich clusters (Ledlow and Owen 1995), as well as for normal ellipticals (Jorgensen et al. 1995).
### 3.6 Structures in the Host Galaxies
A first look at the raw images of BL Lacs as well as the surface photometry analysis above indicates that galaxies hosting BL Lacs are rather smooth and unperturbed (see images in Paper I). However, faint sub-structures could be hidden by the smooth contribution of the stellar component, so we subtracted a two-dimensional model of the galaxy using the best-fit parameters from the surface photometry analysis. This procedure enhances faint structures such as companions, jet-like features, or any other high-contrast feature superposed onto the image of the galaxy (e.g., Faundez-Abans & Oliveira-Abans 1998, and references therein). Clearly any region masked out during isophotal analysis (including the diffraction pattern of PSF spikes) will come out from this procedure. Moreover due to the sharpness of the nucleus and the insufficient sampling of the PC it was not possible to model the PSF adequately with this technique and some residuals (typically diffraction spikes) remain close to the nucleus. This has very little impact on our conclusions since structures very close to the nucleus are not considered.
In Figure 8 we show three examples of the images of the BL Lacs before (left) and after (right) subtraction of the galaxy model. Panel (a) shows the optical jet of PKS 0521-36 after subtraction of the galaxy-plus-point-source model. Panel (b) shows the image of BL Lac itself (2200+420) which, after subtraction of the model, does not exhibit significant features. Finally panel (c) shows an example (H 2356-309) of a companion at $`1`$ arcsec south-west of the nucleus. In all three subtracted images the emission from the diffraction spikes (not modeled) is clearly visible.
The BL Lac objects in the present sample span the redshift range $`0.03z0.2`$. The usable field of view of the PC camera corresponds thus to projected linear size of $`25`$ to $`135`$ kpc. At HST resolution therefore we are practically able to explore structures in the host galaxies as close as 0.2 – 0.9 kpc (depending on $`z`$) to the nucleus (we assume that within 5 pixels from the center the signal is strongly dominated by the nuclear point source).
In three cases (0521-36, 3C 371 and 2201+044) an optical jet is clearly visible (see Scarpa et al. 1999b for a detailed study of these objects). One object (0806+524) has arc-like emission at $`2`$ arcsec from the nucleus, which could be a shell around the galaxy, a feature not uncommon among ellipticals (see Scarpa et al. 1999a for details). No other faint structures (shells, bars or X features) are detected in the observed sources.
HST observations of nearby radio galaxies have shown that dusty features are present in the central regions of most of the observed objects (Kleijn et al. 2000). These features often take the form of disks or lanes but sometimes the distribution is irregular. The observed dust pattern is generally confined within a region typically less than 1.5 kpc (in most cases the physical sizes of disks and lanes range from 0.3 to 1.5 kpc). Since BL Lacs are believed to be radio galaxies with the jet pointing toward us, we should observe similar dusty features also in our objects. However, two difficulties make this detection less likely in our sample objects. First, the bright nucleus tends to out-shine the central regions of the host galaxy. Second, the average distance of our objects is larger than that of comparison samples of radio galaxies observed with HST, making the angular size of dusty features smaller. At $`z=0.05`$, 1 kpc corresponds to $`0.75`$ arcsec; at this distance from the center the light from the nucleus is normally fainter than that from the host galaxy (see Fig. 1 in Paper I) so that these features could in principle be observable at the very lowest redshift but become progressively more difficult to detect at higher $`z`$.
In order to explore this possibility we have therefore considered only the 9 objects at $`z<0.1`$ and searched for dusty features in the central regions after subtraction of the bright point source from the original image. In only one object (1959+650), which is at redshift $`z=0.048`$, is a significant dust lane observed at $`0.8`$ arcsec (corresponding to $`1`$ kpc) north of the nucleus and roughly aligned with the major axis of the galaxy. This feature was also noted by Heidt et al. (1999) on ground-based images and is discussed in our previous work (Scarpa et al. 1999a) on peculiar objects observed in the snapshot survey. For the rest of the eight $`z<0.1`$ sources no clear signature of dust is found. Note that four of these objects are at a redshift similar to or smaller than that of 1959+650. This lack of detection suggests that the relevance of dusty features may be different in BL Lacs and radio galaxies, with the caveat of the small statistics considered. Moreover, note that looking at features close to the nucleus requires not only high spatial resolution but also a PSF with sufficient sampling and faint wings. HST with the new Advanced Camera for Surveys High Resolution Camera offers two-times better sampling, and has a coronograph that depresses the PSF wings by factors of a few (and obscures the central 0.9 arcsec), so it could contribute to elucidating this point.
## 4 Close Companions
In a number of objects we noted the presence of faint companions around the target, a feature that seems also common to the HST images of quasars. We have systematically explored the near environments of the 30 low-redshift BL Lac sources with the aim of identifying potential companion objects. For three cases (0706+592, 1440+122 and 2356-309) a very close ($`\mathrm{\Delta }<`$ 1$`^{\prime \prime }`$.2) compact companion is detected. These companions are located at projected distances of the order of 1–5 kpc from the nucleus. For five other BL Lac objects (see Table 2) diffuse or compact companions are found within a projected distance of $`20`$ kpc (if they are at the same redshift as the BL Lac).
In order to evaluate the statistical significance of the occurrence of these close companions we derived the expected number of objects within a given radius from the BL Lac based on the number counts observed in our PC frames well away from the BL Lac. The observed average density of objects with $`19.5<m_R<23.5`$ mag is $`9\pm 5`$ arcmin<sup>-2</sup>, in good agreement with the observed galaxy density in several high galactic latitude fields (Metcalfe et al. 1991).
At these faint magnitudes the counts at high galactic latitudes are dominated by galaxies. The average galactic latitude of our fields is $`|b|=45^{\mathrm{deg}}`$ but two fields (2200+420, 2344+514) are at low galactic latitude ($`|b|10^{\mathrm{deg}}`$) and one (1514–241) is at $`|b|27^{\mathrm{deg}}`$ but close to the direction toward the galactic center. The observed average density of stellar (unresolved) objects including all PC frames is $`3.2\pm 4.2`$ arcmin<sup>-2</sup>, while excluding the three fields mentioned above yields a star density of $`2.1\pm 2.6`$ arcmin<sup>-2</sup>. This is consistent with other estimates of the stellar density at high galactic latitude (Hintzen, Romanishin & Valdes 1991; Metcalfe et al. 1991).
Assuming the observed average density of faint ($`19.5<m_R<23.5`$ mag) objects, 9 arcmin<sup>-2</sup>, around 30 targets we would expect to find 0.06 companions at $`r<`$0$`^{\prime \prime }`$.5, while one (a compact companion at 0.3 arcsec from 1440+122) is observed (Poisson probability P=0.06). Within 1$`^{\prime \prime }`$.5 we would expect to observe 0.6 objects while 3 are detected (P=0.02). At larger radii ($`r<`$5<sup>′′</sup>) we observe 11 companions when 6 are expected (P=0.02). To test the reality of this apparent excess, we counted the number of detected objects within these same radii for random positions in each PC frame. In those cases we found 0, 1, and 7 companions, respectively, in agreement with the expected numbers of faint field objects.
Although the numbers are small and consequently the statistics are not very good, we conclude from our analysis that there is some indication that an association between BL Lacs and companions is not dominated by chance alignments.
The identity or role of these companions is not yet clear, however. In other BL Lac fields, optical spectroscopy of the companion has sometimes proved the physical association with the active galaxy (Falomo 1996); in other cases, they have proved to be stars or foreground or background galaxies, fortuitously aligned with the BL Lac (e.g., Stickel et al. 1991). For most cases here, spectra of the companion objects are not available and only statistical arguments suggest some are likely associated with the BL Lac object.
Some of the close companions are resolved faint galaxies with absolute magnitudes (assuming they are at the distance of the BL Lac object) in the range $`21<m_R<19`$ mag (Table 2). This is consistent with previous evidence that BL Lac objects inhabit environments with higher-than-average galaxy densities, as in groups or poor clusters (Pesce, Falomo & Treves 1995; Wurtz et al. 1996). In the two cases with more than one companion, 1229+645 and 1440+122, the two companion objects are clearly different from one another, one being resolved and the other unresolved,
We note that a significant number of the observed companions are compact and appear point-like even at HST resolution. Based on the counts of Galactic stars, we would expect only 1 or 2 chance alignments with stars while 6 compact companions are observed (Table 2). The excess of compact companions is in fact more significant than the excess of resolved companions (galaxies), especially since none of the compact sources belong to the low galactic latitude fields where the star density is higher than the average value.
The nature of these unresolved companions is unclear. If they were at the redshift of the BL Lac object, their magnitudes would be in the range $`18<m_R<15`$ mag, too luminous for globular clusters. These magnitudes are consistent with dwarf ellipticals but the observed sources are probably too compact, unless they contain weak active nuclei (but they would have to be very weak indeed). One can ask whether some of these compact sources could be supernovae of type Ia. The SN Ia rate in early type galaxies is estimated to be about 0.1–0.2 SN/100yr/10<sup>10</sup> L$`{}_{}{}^{B}{}_{}{}^{}`$ (Cappellaro et al. 1999). This means that in a massive galaxy of $`M_R`$ -23.7 mag we expect to have about 1-2 SN/100yr. Because the SN Ia at maximum reach $`M_R`$ -19.4 mag and our limiting magnitude for stellar objects is $`m_R24.5`$ mag (it is $`m_R24`$ mag for extended objects), at $`z0.1`$ we should be able to detect such SN for about 100 days. Therefore we should find 0.2 SN in our 30 galaxies observed. Unless the SN rate in these active galaxies is significantly higher than the average for normal ellipticals, it is highly unlikely that the compact companions are SN.
We believe the nature of all these companions, resolved or unresolved, will remain uncertain until spectroscopy is performed. Whatever their identity, there is little evidence of tidal interaction in the BL Lac host galaxy (appreciable isophote twisting, non-concentric isophotes, etc.), which suggests the companions are not strongly perturbing it. Gravitational interactions might have had a greater role in the past in the formation and fueling of the active nucleus, and what we see now could be a leftover of an earlier (close) interaction. One interesting possibility suggested by numerical simulations (Bekki 1999) is that the close companions are the product of a past major merger between gas-rich galaxies and also provide fuel for the nuclear activity but are not necessarily linked to the formation of a massive black hole in the nucleus, or to the sustenance of an active relativistic jet.
## 5 Summary and Conclusions
We have presented detailed morphological analysis of HST images for 30 galaxies hosting BL Lac sources at redshift $`z0.2`$. The galaxies were investigated with an order of magnitude better spatial resolution than it is possible from the ground.
Our analysis indicates that in spite of the presence of the active nucleus the host galaxy appears in most cases to be a completely “normal” elliptical. Both the ellipticity and the isophotal shape distributions are similar to those for radio galaxies and radio-quiet ellipticals. This suggests that tidal interactions are very infrequent or are short-lived with respect to the nuclear activity time scale.
We find no indication of displacement and/or off-centering of the galaxy isophotes with respect to the nucleus, meaning the unresolved nuclear source truly sits in the center of the galaxy. This rules out the microlensing hypothesis for BL Lacs, which predicts frequent off-centering of the nucleus. This does not exclude the possibility that some objects could be lensed (see Scarpa et al. 1999a), but it cannot be a widespread explanation of the BL Lac phenomenon.
An interesting result is the finding of excess close companions in several cases. A high incidence of companions of BL Lacs was first suggested by (Falomo, Melnick & Tanzi 1991) from sub-arcsec resolution ground-based imaging, and was later confirmed by Pesce, Falomo and Treves (1995), Falomo (1996) and Heidt et al. (1999). Close companions seem also to be common around quasars imaged by HST (Bahcall et al. 1997; Disney et al. 1995).
We thank G. Fasano for useful suggestions and support performing isophotal (AIAP) analysis and E. Cappellaro and M. Turatto for discussion about the SN rate. This work was partly supported by the Italian Ministry for University and Research (MURST) under grant Cofin98-02-32 and Cofin 98-02-15. Support for CMU and RS was provided by NASA through grant number GO-06363.01-95A from the Space Telescope Science Institute, which is operated by AURA, Inc., under NASA contract NAS 5-26555.
| Table 1 | | | | | | |
| --- | --- | --- | --- | --- | --- | --- |
| Structural Properties of Low-Redshift Host Galaxies | | | | | | |
| Name | z | $`ϵ^{(a)}`$ | $`C_4^{(b)}`$ | $`\delta ^{(c)}`$ | $`\mathrm{\Delta }^{(d)}`$ | $`\mathrm{\Delta }\theta ^{(e)}`$ |
| 0145+138 | 0.124 | 0.09 | -0.40 | 0.04 | 0.05 | 22.0 |
| 0229+200 | 0.139 | 0.16 | -0.50 | 0.02 | 0.03 | 10.0 |
| 0347-121 | 0.188 | 0.03 | 0.30 | 0.03 | 0.06 | $`\mathrm{}`$ |
| 0350-371 | 0.165 | 0.25 | -1.20 | 0.02 | 0.02 | 5.0 |
| 0521-365 | 0.055 | 0.30 | 0.00 | 0.02 | 0.04 | 4.0 |
| 0548-322 | 0.069 | 0.20 | -0.01 | 0.03 | 0.04 | 12.0 |
| 0706+591 | 0.125 | 0.15 | -0.80 | 0.02 | 0.04 | 16.0 |
| 0806+524 | 0.137 | 0.08 | 1.20 | 0.03 | 0.05 | $`\mathrm{}`$ |
| 0829+046 | 0.180 | 0.10 | 0.20 | 0.04 | 0.04 | $`\mathrm{}`$ |
| 0927+500 | 0.188 | 0.24 | 0.60 | 0.03 | 0.04 | 7.0 |
| 1104+384 | 0.031 | 0.19 | 0.00 | 0.03 | 0.04 | 15.0 |
| 1136+704 | 0.045 | 0.04 | 0.00 | 0.01 | 0.01 | $`\mathrm{}`$ |
| 1212+078 | 0.136 | 0.04 | -0.10 | 0.01 | 0.01 | $`\mathrm{}`$ |
| 1218+304 | 0.182 | 0.11 | -1.50 | 0.04 | 0.05 | 23.0 |
| 1229+643 | 0.164 | 0.16 | -0.30 | 0.02 | 0.03 | 12.0 |
| 1255+244 | 0.141 | 0.09 | 0.60 | 0.02 | 0.03 | $`\mathrm{}`$ |
| 1418+546 | 0.152 | 0.30 | -0.30 | 0.02 | 0.02 | 15.0 |
| 1426+428 | 0.129 | 0.33 | -0.30 | 0.02 | 0.03 | 3.0 |
| 1440+122 | 0.162 | 0.20 | 1.40 | 0.03 | 0.01 | 10.0 |
| 1514-241 | 0.049 | 0.02 | 0.20 | 0.02 | 0.03 | $`\mathrm{}`$ |
| 1728+502 | 0.055 | 0.06 | 0.40 | 0.01 | 0.01 | $`\mathrm{}`$ |
| 1807+698 | 0.051 | 0.05 | -0.80 | 0.02 | 0.04 | $`\mathrm{}`$ |
| 1853+671 | 0.212 | 0.15 | 0.40 | 0.05 | 0.04 | 30.0 |
| 1959+650 | 0.048 | 0.20 | 1.50 | 0.03 | 0.07 | 17.0 |
| 2005-489 | 0.071 | 0.26 | 0.80 | 0.03 | 0.05 | 12.0 |
| 2200+420 | 0.069 | 0.25 | -0.90 | 0.02 | 0.05 | 4.0 |
| 2201+044 | 0.027 | 0.06 | -0.60 | 0.05 | 0.00 | $`\mathrm{}`$ |
| 2326+174 | 0.213 | 0.12 | -0.90 | 0.04 | 0.03 | 23.0 |
| 2344+514 | 0.044 | 0.24 | -0.60 | 0.02 | 0.02 | 5.0 |
| 2356-309 | 0.165 | 0.25 | 0.70 | 0.04 | 0.07 | 6.0 |
| <sup>(a)</sup> Ellipticity of the host galaxy at the effective radius. | | | | | | |
| Typical error $`\pm `$ 0.02 far from the nucleus. | | | | | | |
| <sup>(b)</sup> Fourier coefficient 100\*$`C_4`$ describing the isophote | | | | | | |
| shape (see text); typical uncertainty is 0.5. | | | | | | |
| <sup>(c)</sup> Relative displacement of isophotes (dimensionless, | | | | | | |
| see definition in § 3.3); typical uncertainty is $`\pm `$ 0.01. | | | | | | |
| <sup>(d)</sup> Off-centering (arcsec) of the nucleus with respect to the host | | | | | | |
| galaxy (see definition in § 3.3); typical uncertainty is 0.05 arcsec. | | | | | | |
| <sup>(e)</sup> Maximum observed twisting of isophotes (deg); | | | | | | |
| measured only for host galaxies with $`ϵ>`$ 0.1. | | | | | | |
Table 2
Properties of Close Environments of BL Lac Objects Name Feature $`m_R`$ $`M_R^{(a)}`$ $`\mathrm{\Delta }^{(b)}`$ PA<sup>(c)</sup> (mag) (mag) (<sup>′′</sup>/kpc) (deg) 0521–365 Optical jet 19.9 $`\mathrm{}`$ 1.8 / 2.7 $``$ 305 0706+592 Compact companion 24.8 -15 1.14 / 3.5 170 0829+046 Companion galaxy 19.9 -20.5 4.9 / 20 145 1229+645 Companion galaxy 19.2 -20.9 3.4 / 13 210 Compact companion 21.3 -18.8 4.5 / 17 3 1426+428 Compact companion 21.7 -17.9 3.9 / 12 15 1440+122 Companion galaxy 16.7 -23.4 2.5 / 9.4 260 Compact companion 19.5 -20.6 0.3 / 1.1 70 1959+650 Dust lane $`\mathrm{}`$ $`\mathrm{}`$ 1.2 / 1.6 $``$ 20 1807+698 Optical jet 21.7 $`\mathrm{}`$ 3.1 / 4.2 $``$ 210 1853+671 Companion galaxy 21.8 -18.9 2.1 / 9.8 323 2005–489 Diffuse companion 22.7 -15.5 8.5 / 15.6 15 2201+044 Optical jet 24.2 $`\mathrm{}`$ 2.1 / 1.6 $``$ 315 2326+174 Compact companion 23.2 -17.5 3.2 / 15 150 2356–309 Compact companion 22.5 -17.3 1.2 / 4.6 113 (a) Absolute magnitude assuming the same redshift as the BL Lac object. (b) Projected distance from the nucleus (in arcsec and kpc) assuming assuming the same redshift as the BL Lac object. (c) Position angle of the feature.
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# The hydration number of Na+ in liquid water
## I Introduction
Solvation of simple ions in aqueous solution is not yet fully understood despite its fundamental importance to chemical and biological processes. For example, disagreement persists regarding the hydration number of the Na<sup>+</sup> ion in liquid water. A pertinent problem of current interest centers on the selectivity of biological ion channels; it seems clear that the selective transport of K<sup>+</sup> relative to Na<sup>+</sup> ions in potassium channelskchannel:98 ; guidoni:99 ; laio:99 depends on details of the ion hydration that might differ for K<sup>+</sup> relative to Na<sup>+</sup>.
Experimental efforts to define the hydration structure of Na<sup>+</sup>(aq) using diffractioncaminiti:80 ; skipper:89 and spectroscopicMichaellian:78 methods produce a hydration number ranging between four and sixOhtaki:93 . Simulations have obtained a range of values, but most predict six water molecules in the inner hydration sphere of the Na<sup>+</sup> ionHeinzinger:79 ; Mezei:81 ; Impey:83 ; Chandrasekhar:84 ; bounds:85 ; wilson:85 ; Zhu:91 ; Heinzinger:93 ; lee:94 ; toth:96 ; obst:96 ; Koneshan:98b . An ‘ab initio’ molecular dynamics simulation produced five inner shell water molecules neighboring Na<sup>+</sup>(aq)schwegler:00 .
An important limitation of theoretical studies of ion hydration concerns the sufficiency of model force fields used in classical statistical mechanical calculations. In the most customary approaches, interatomic force fields used in theories or simulations derive from empirical fits of a parameterized model to a variety of experimental data. ‘Ab initio’ molecular dynamics avoids this intermediate modeling step by approximate solution of the electronic Schroedinger equation for each configuration of the nucleimarx:99 ; alfe:2000 . This technique thus goes significantly beyond conventional simulations regarding the accuracy of the force fields. It also augments theories built more directly on electronic structure studies of ion-water complexes by adopting approximate descriptions of the solution environment of those complexesrempe:99 ; pratt:98 ; feature ; martin:97 ; pratt:99 ; G98 .
Relative to conventional simulations, ‘ab initio’ molecular dynamics simulations also have some important limitations due to the high computational demand. Applications of the method have been restricted to small systems simulated for short times. For example, an ‘ab initio’ molecular dynamics studyschwegler:00 of the Na<sup>+</sup>(aq) ion comparable to the present work obtained a thermal trajectory lasting 3 ps after minimal thermal aging. The present work, though still limited to relatively small systems, pushes such calculations to longer times that might permit more precise determination for Na<sup>+</sup>(aq) of primitive hydration properties. The analysis here utilizes the last half of a 12 ps thermal trajectory. The quasi-chemical theoryrempe:99 ; pratt:98 ; feature ; martin:97 ; pratt:99 and separate electronic structure calculations on Na(H<sub>2</sub>O)$`{}_{n}{}^{}^+`$ complexes assist in this analysis.
## II Methods
The system consisted of one Na<sup>+</sup> ion and 32 water molecules in a cubic box with edge 9.86518 Å and periodic boundary conditions. The dimensions of the box correspond to a water density of 1 g/cm<sup>3</sup> and zero partial molar volume for the solute. Initial conditions were obtained as in an earlier ‘ab initio’ molecular dynamics simulation on Li<sup>+</sup>(aq)rempe:99 . In that earlier work, an optimized structure for the inner sphere Li(H<sub>2</sub>O)$`{}_{6}{}^{}^+`$ complex was equilibrated with 26 water molecules under conventional simulation conditions for liquid water, utilizing a current model force field and assuming a partial molar volume of zero. In the present calculation, the same pre-equilibrated system was used as an initial configuration for the ‘ab initio’ molecular dynamics except that an optimized structure for the inner sphere Na(H<sub>2</sub>O)$`{}_{6}{}^{}^+`$ complex replaced the hexa-coordinated Li<sup>+</sup> structure. Constant pressure or constant water activity simulations, defined by intensive rather than extensive variables, probably would produce a more useful characterization of the solvent thermodynamic state for these small systems, but those alternatives are currently impractical.
Molecular dynamics calculations based upon a gradient-corrected electron density functional description of the electronic structure and interatomic forces were carried out on this Na<sup>+</sup>(aq) system utilizing the VASP programvasp1 . The ions were represented by ultrasoft pseudopotentialsvasp2 and a kinetic energy cutoff of 31.5 Rydberg limited the plane wave basis expansions of the valence electronic wave functions. The equations of motion were integrated in time steps of 1 fs, which is small enough to sample the lowest vibrational frequency of water. A thermostat constrained the system temperature to 300 K during the first 4.3 ps of simulation time. After removing the thermostat, the temperature rose slightly and then leveled off by 6 ps to an average of 344 $`\pm `$ 24 K. During the simulation, the initial $`n`$=6 hydration structure relaxed into $`n`$=4 and $`n`$=5 alternatives, such as those shown in Fig. 1. All analyses reported here rely on the trajectory generated subsequent to the 6 ps of aging with the system at a temperature elevated from room temperature.
## III Results
The ion-oxygen radial distribution function is shown in Fig. 2. The first maximum occurs at a radius of 2.35 Å from the Na<sup>+</sup> ion and the minimum at radius 3.12 Å demarcates the boundary of the first and innermost hydration shell. An average of $`n`$=4.6 water molecules occupy the inner hydration shell. Fig. 3 tracks the instantaneous number of water oxygen atoms found within the first hydration shell of the Na<sup>+</sup>, defined by radius r$``$3.12 Å for the upper panel. The fractions $`x_4`$ and $`x_5`$ of four-coordinate and five-coordinate hydration structures, respectively, constitute $`x_4`$=40% and $`x_5`$=56% of the last 6 ps of the simulation. Structures in which the Na<sup>+</sup> ion acquires six innershell water molecules occur with a 4% frequency, while structures with three and seven innershell water molecules occur less than 1% of the time. Analysis of the mean-square displacement of the Na<sup>+</sup> ion (Fig. 4) produces a self-diffusion constant of 1.0$`\times `$10<sup>-5</sup> cm<sup>2</sup>/s, which agrees reasonably well with an experimental result of 1.33x10<sup>-5</sup>cm<sup>2</sup>/shertz:73 .
These results correspond coarsely with an ‘ab initio’ molecular dynamics calculation on this system carried-out independentlyschwegler:00 . The most probable inner shell occupancy found there was also five, but the probabilities of n=4 and n=6 were reversed from what we find here. This difference may be associated with the lower temperature used in Ref schwegler:00 .
One motivation for this study arises from the quasi-chemical theory of solutions. According to this formulation, $`x_0`$ contributes a ‘chemical’ contribution to $`\mu _{Na^+}^{ex}`$, the excess chemical potential or absolute hydration free energy of the ion in liquid waterpratt:99 ,
$`\beta \mu _{\mathrm{Na}^+}{}_{}{}^{ex}=\mathrm{ln}x_0\mathrm{ln}\left[e^{\beta \mathrm{\Delta }U}{\displaystyle \underset{j}{}}(1b_{\mathrm{Na}^+j})_0\right].`$ (1)
Here the inner shell is defined by specifying a function $`b_{\mathrm{Na}^+j}`$ that is equal to one (1) when solvent molecule j is inside the defined inner shell and zero (0) otherwise; $`\mathrm{\Delta }U`$ is the interaction energy of the solvent with the solute Na<sup>+</sup> that is treated as a test particle, $`\beta ^1`$=k<sub>B</sub>T, and the subscript zero associated with $`\mathrm{}_0`$ indicates a test particle average pratt:99 . The second term on the right-hand side of Eq. (1) is the excess chemical potential of the solute lacking inner shell solvent molecules whereas the first term is the free energy of allowing solvent molecules to occupy the inner shell. The validity of Eq. (1) has been established elsewherepratt:99 . The second term on the right of Eq. 1 is the outer sphere contribution to the excess chemical potential in contrast to the first or chemical term.
The utility of this quasi-chemical formulation is the suggestionhummer:cp:2000 of more detailed study of the $`x_n`$, the fractions of n-coordinate hydration structures found in solution, on the basis of the equilibria forming inner shell complexes of different aggregation number:
$`\mathrm{Na}(\mathrm{H}_2\mathrm{O})_{\mathrm{m}=0}{}_{}{}^{+}+\mathrm{nH}_2\mathrm{O}\mathrm{Na}(\mathrm{H}_2\mathrm{O})_\mathrm{n}{}_{}{}^{+}.`$ (2)
Utilizing the chemical equilibrium ratios
$`K_n={\displaystyle \frac{\rho _{\mathrm{Na}(\mathrm{H}_2\mathrm{O})_\mathrm{n}^+}}{\rho _{\mathrm{H}_2\mathrm{O}}{}_{}{}^{n}\rho _{\mathrm{Na}(\mathrm{H}_2\mathrm{O})_{\mathrm{m}=0}^+}^{}}},`$ (3)
the normalized $`x_n`$ can be expressed as
$`x_n={\displaystyle \frac{K_n\rho _{\mathrm{H}_2\mathrm{O}}^n}{\underset{m0}{}K_m\rho _{\mathrm{H}_2\mathrm{O}}^m}}.`$ (4)
The $`\rho _\sigma `$ are the number densities and, in particular, $`\rho _{\mathrm{H}_2\mathrm{O}}`$ is the molecule number density of liquid water. If the medium external to the clusters is neglected, the equilibrium ratios, denoted as $`K_n^{(0)}`$, can be obtained from electronic structure calculations on the complexes, assuming for the thermal motion of the atoms the harmonic approximation evaluated at the calculated minimum energy configuration. Finally utilization of a dielectric continuum approximation for the outer sphere contributions to the chemical potential gives a natural, though approximate, quasi-chemical modelrempe:99 ; pratt:98 ; feature ; martin:97 ; pratt:99 ; G98 .
For the present problem, the quasi-chemical approximation was implemented following precisely the procedures of the earlier study of Li<sup>+</sup>(aq)rempe:99 , except that the sodium ion cavity radius for the dielectric model calculation was assigned as R$`_{Na^+}`$=3.1 Å, the distance of the first minimum of the radial distribution function of Fig. 2. The temperature and density used were 344 K and 1.0 g/cm<sup>3</sup> and the value of the bulk dielectric constant was 65.3uematsu:80 .
Results of the calculations are summarized in Fig. 5. The electronic structure results are consonant with those found previously for the Li<sup>+</sup> ion. The n=4 inner sphere gas-phase complex has the lowest free energy. Although outer sphere placements are obtained for additional water molecules in the minimum energy structures of larger clusters, attention is, nevertheless, here restricted to inner sphere structures. The mean occupation number predicted by this quasi-chemical model is $`n`$ = 4.0; the computed absolute hydration free energy of the Na<sup>+</sup> ion under these conditions is -103 kcal/mol, not including any repulsive force (packing) contributions. An experimental value for Na<sup>+</sup> ion in liquid water at room temperature is -87 kcal/molmarcus:94 .
Because of the significance of $`x_0`$ \[Eq. 1\], we fitted several model distributions {$`x_n`$} based on different ideas to the ‘ab initio’ molecular dynamics results. The varying success of those models in inferring $`x_0`$ was enlightening. An instructive selection of those models is shown in Fig. 6 and we describe those results here.
First, we note that though the preceeding quasi-chemical approximation does not agree closely with the ’ab initio’ molecular dynamics simulation, the populations obtained from the quasi-chemical approximation, $`\widehat{x}_n`$, can serve as a default model for a maximum entropy inference of $`x_n`$hummer:cp:2000 . In this approach we model
$$\mathrm{ln}x_j=\mathrm{ln}\widehat{x}_j\lambda _0j\lambda _1j(j1)\lambda _2/2\mathrm{},$$
(5)
with Lagrange multipliers $`\lambda _k`$ adjusted to conform to available moment information
$$\left(\genfrac{}{}{0pt}{}{n}{j}\right)=\underset{k}{}x_k\left(\genfrac{}{}{0pt}{}{k}{j}\right)$$
(6)
for j = 0, 1, 2, …. In view of the limited data available, use of more than two moments produced operationally ill-posed fitting problems.
One difficulty with this specific approach is that the ‘ab initio’ molecular dynamics produced $`x_7>`$0 in contrast to the electronic structure methods that found no minimum energy hepta-coordinated inner-sphere clusters. Since the observed $`x_7`$ is likely to be relatively less accurate and is furthest away from the desired n=0 element, we excluded n=7 configurations of the ‘ab initio’ molecular dynamics, renormalized the probabilities $`x_n`$, and recalculated the moments. As the upper panel in Fig. 6 shows, this simple maximum entropy model is qualitatively satisfactory although not quantitatively convincing. The fitted model significantly disagrees with the observed $`x_3`$. The chemical contribution suggested by Fig. 6 is approximately -70 kcal/mol. Using the Born formula, $`q^2(11/ϵ)/2R`$ with R=3.12 Å, to estimate the outer sphere contributions represented by the last term in Eq. 1, then the net absolute hydration free energy falls in the neighborhood of -115 kcal/mol. Since experimental values for the absolute hydration free energy at room temperature center around -90 kcal/mol, this comparison shows that the present free energy results are not to be interpreted quantitatively, but rather as indicative of the present state of the theory.
A second approach focused on testing a default model that supplies a nonzero $`\widehat{x}_7`$; we used the Gibbs default model $`\widehat{x}_n1/n!`$ that would give the correct answer for an ideal gas ‘solvent.’ This model has the additional and heuristic advantage of being significantly broader. Our experience has been that these maximum entropy fitting procedures work better when the default model is broader than the distribution sought. The results, illustrated in the middle panel of Fig. 6, show an improved fit. Here the chemical contribution to the free energy is -23 kcal/mol, yielding a net absolute hydration free energy of -68 kcal/mol when the same Born formula is used to estimate the outer sphere contributions.
A third fitting possibility was based on a suggestion from a previous ‘ab initio’ molecular dynamics calculation on K<sup>+</sup>(aq): that the innermost four water molecules have a special statusramaniah:98 . In fact, the quasi-chemical approximation above and the fitting of the upper panel of Fig. 6 suggests also that the $`x_n`$ results for n$``$4 and for n$``$5 display different behaviors. The radial distribution function of Fig. 2, somewhat better resolved than heretofore, is relevant to this issue and, in contrast, doesn’t directly support a hypothesis of two populations of water molecules in the inner shell. Nevertheless, that g(r) does not rule out the possibility that the structures might become more flexible as the inner shell nears maximum capacity with lower incremental binding energies.
To clarify these possibilities, we reduced the radius defining the inner sphere to R=2.68 Å, for which $`<`$n$`>`$ is close to 4 (see bottom panel of Fig. 2) and reanalysed the ‘ab initio’ molecular dynamics trajectory to extract the appropriate alternative moment information. Again using the Gibbs default model, we obtained the results shown in the lowest panel of Fig. 6. The inferred chemical contribution is $`RT\mathrm{ln}x_0`$ -13 kcal/mol. Using again the Born approximation for the outer sphere contribution, this time with R=2.68 Å, we obtain an absolute hydration free energy estimate of -65 kcal/mol.
The insensitivity of these latter results to choice of inner sphere radius deserves emphasis and further discussion. From a formal point of view, the inner sphere radius R serves only a bookkeeping role; the left side of Eq. 1 should be strictly unaffected by changes in R. Nevertheless, the terms on the right side of that equation are individually affected by changes in R. Thus, we might take the insensitivity of the sum of those individual terms as an indication that the inevitable approximations are reasonably balanced. Values of R for which the sum Eq. 1 is insensitive are pragmatic values given the approximations made. Fig. 7 illustrates these points and establishes the pragmatic value R$``$3.06Å for the current application. The similarity of this value with the radius of the inner shell suggested by Fig. 2 (3.12Å) is encouraging. The value -68 kcal/mol is then suggested for the hydration free energy in the absence of any account of packing or van der Waals interactions.
## IV Conclusions
The ‘ab initio’ molecular dynamics simulation predicts the most probable occupancy of the inner shell of Na<sup>+</sup>(aq) to be 5 and the mean occupancy to be 4.6 water molecules at infinite dilution, T=344 K, and a nominal water density of 1 g/cm<sup>3</sup>. The simulation produces both a satisfactory Na-O radial distribution function and self-diffusion coefficient for Na<sup>+</sup>, but these satisfactory results required more care with thermalization and averaging time than is most common with these demanding calculations. Recently, this point has been separately emphasized in the context of ‘ab initio’ simulation of watersorenson:00 .
The complementary calculation framed in terms of quasi-chemical theory based on electronic structure results for ion-water clusters, the harmonic approximation for cluster motion, and a dielectric continuum model for outer sphere contributions underestimates the inner shell water molecule occupancies for Na<sup>+</sup> in liquid water. Maximum entropy fitting of the inner shell occupancy distribution shows that the ion-water cluster results yield a distribution significantly narrower than that obtained from the simulations. For this reason, naive inference of the absolute hydration free energy of Na<sup>+</sup>(aq) based on the cluster electronic structure results and utilizing information gleaned from the ‘ab initio’ molecular dynamics was unsuccessful. The electronic structure calculations found minimum energy Na(H<sub>2</sub>O)$`{}_{7}{}^{}^+`$ clusters only with obvious outer sphere placements of some water molecules though hepta-coordinate inner sphere clusters were observed in the ‘ab initio’ molecular dynamics with the most natural cluster definition. These results suggest that the anharmonicities and large amplitude motion are serious concerns, particularly for the larger clusters, and that the approximate theory utilized for outer sphere contributions must treat cluster conformations differently from minimum energy structures of the isolated clusters.
Abandonment of the cluster electronic structure results in favor of a broader default model improved the modeling of the $`x_n`$ distribution on the basis of the information extracted from the simulation. A sequence of more aggressive fits eventually suggested the value -68 kcal/mol for the hydration free energy at this somewhat elevated temperature on the basis of the quasi-chemical perspective of inner sphere occupancies but in the absence of any account of packing or van der Waals interactions.
We acknowledge helpful discussions of many related issues with Gerhard Hummer and Joel Kress. This work was supported by the US Department of Energy under contract W-7405-ENG-36 and the LDRD program at Los Alamos.
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# Wigner-Dyson Statistics from the Keldysh 𝜎-Model
In the past two decades, the field theory of disordered electronic systems, a matrix version of the nonlinear $`\sigma `$–model, has attracted a great deal of theoretical attention . The $`\sigma `$-model provides the perhaps most rigorous and comprehensive theoretical framework by which interacting and non–interacting disordered electronic systems can be explored. Three different, albeit closely related versions of this model exist: The supersymmetric (SUSY) , replica and dynamic or Keldysh $`\sigma `$-model. Due to different microscopic starting points these three formulations cover partly complementary application areas: whereas SUSY is tailor made to the analysis of non-perturbative phenonmena, the replica, and in particular the Keldysh formalism are strong at problems involving interactions – about which SUSY has nothing to say. This constellation amounts to some theoretical vacuum because sooner or later non-perturbative interaction related problems will come into focus.
An important step towards improving the situation was made recently (first within the framework of SUSY) when it became clear that the large energy asymptotic of Wigner-Dyson (WD) level statistics – the non-perturbative prototype problem – may be obtained from a careful analysis of the saddle point structure of the $`\sigma `$–model (in contrast to a full integration over the field manifold which is out of question for the non-SUSY representants). Subsequently these ideas were adopted to the replica $`\sigma `$–model , which proved to be useful for the solution of simple interacting problems . Yet, the application spectrum of the replica theory, excluding non-equilibrium phenomena, is significantly narrower than that of its Keldysh counterpart. Furthermore, replica analyses categorically rely on a formidable analytic continuation procedure, which tends to obscure the underlying phyiscs and essentially complicates practical applications. It is therefore desirable to formulate a prescription whereby non-perturbative information can be obtained from the conceptually more transparent Keldysh formalism. This is the subject of the present Letter.
All the theories mentioned above share the normalization condition $`Z=1`$, where $`Z`$ is the functional partition function. It is this universal normalzation property which makes the models useful in the analysis of disorder properties. The normalization in turn is due to a global internal symmetry underlying the model: supersymmetry, replica permutation symmetry or, within the Keldysh formulation, a symmetry principle related to the causality of the theory. Building on these structures, the saddle point analysis of any $`\sigma `$–model must consist of three parts: (i) identification of all saddle points and their action, (ii) proof that $`Z=1`$ holds within the stationary phase scheme (iii) saddle point evaluation of any specific correlation function. Using WD spectral statistics as a example we are going to apply this program to the Keldysh model of unitary (broken time reversal invariance) symmetry. The generalization to different symmetries and observables will be discussed elsewhere .
We wish to compute the two-point correlation function $`R(ϵ_1,ϵ_2)=\mathrm{\Delta }^2\rho (ϵ_1)\rho (ϵ_2)=\mathrm{\Delta }^2/(2\pi ^2)\mathrm{}\mathrm{tr}\{G(ϵ_1^+)\}\mathrm{tr}\{G(ϵ_2^{})\}`$ of the density of states $`\rho (ϵ)`$ for a model of free electrons with broken time reversal invariance and subject to a random potential. Here $`\mathrm{\Delta }`$ is the mean level spacing and $`\mathrm{}`$ stands for a (cumulative) disorder average. Our starting point is the low energy effective partition function $`Z=𝒟Qe^{iS[Q]}`$, where
$$iS[Q]=\frac{\pi }{4\mathrm{\Delta }}\frac{d^d𝐫}{L^d}\mathrm{tr}\left\{DQQ+4i\widehat{ϵ}Q\right\},$$
(1)
$`D`$ is the diffusion constant and $`L`$ the system size. The matrix–valued field $`Q=\{Q_{ϵ,ϵ^{}}^{l,l^{}}\}`$ acts in a product space defined through an index $`l=1,2`$ labelling the forward and backward Keldysh time contour, $`𝒞^l=\{t(\mathrm{},\pm \mathrm{})\}`$, and the energy indices, $`ϵ`$, Fourier conjugate to the time variables $`t`$ on $`𝒞^l`$. For all what follows it will be convenient to discretize the (a priori continuous) energy variables in units of some small spacing $`\delta _ϵ`$. Since the application range of the effective action is limited to energies $`|ϵ|<\tau ^1`$ , this manipulation implies that the $`Q`$’s become matrices of finite dimension $`2K`$, where $`K=2/(\delta _ϵ\tau )`$. The internal structure of these matrices is defined through (a) the constraint $`Q^2(𝐫)=𝟙`$ and (b) Hermiticity, $`Q^{}=Q`$. Finally, $`\widehat{ϵ}=\mathrm{diag}(\widehat{ϵ}^1,\widehat{ϵ}^2)`$ is a $`2K`$-dimensional diagonal matrix where $`\widehat{ϵ}^l=\mathrm{diag}(ϵ_1,\mathrm{}ϵ_n,\mathrm{}ϵ_K)i(1)^l0`$.
The definition of the low energy effective action (1) implies one more important structural element which is not explicit in the notation: by definition of the trace ’tr’, the energy integration/summation over any continuous and analytically benign function of the energy indices $`ϵ`$ produces zero (e.g. $`\mathrm{tr}\{\widehat{ϵ}\}=0`$). For lack of better terminology we will refer to this feature as ’causality’. That the causality criterion is not explicit in the definition of the trace is not just a matter of notational convenience. The crux is that to rigorously retrieve the full causality properties of the microscopic Keldysh partition function, all energies, $`ϵ`$, including energies in excess of the width of the spectrum of the microscopic Hamiltonian, have to be taken into account. The philosophy behind declaring the causality principle to an intrinsic feature of the effective action (1) is that $`S`$ merely represents the low energy sector of some larger parent theory, $`S+S_{\mathrm{high}}`$. After the inclusion of the high energy sector $`S_{\mathrm{high}}`$, all energy summations could, in principle, be extended to infinity and the correct spectral structures of the Keldysh partition function would be retrieved.
This anticipation has been at the root of previous effective action formulations of the Keldysh approach, and, needless to say, is of relevance for all matters related to spectral statistics. We will therefore subdivide our analysis into two parts. Taking a pragmatic point of view we will first show how the spectral correlation function can be obtained from the low energy model, Eq. (1), once the causality assumption has been made. In a second part we will then show how causality can be made manifest, on the expense of including large energies. We emphasize that part II of the analysis has been included for reasons of conceptual completeness; in practical applications, the large energy sector will normally not play a role.
Part I: The simplest way of computing $`R(ϵ_1,ϵ_2)`$ from the functional integral over the effective action $`iS[Q]`$, is to couple the energy vector $`\widehat{ϵ}`$ to sources. This can be done by generalizing $`\widehat{ϵ}^l\widehat{ϵ}^{l,\kappa }\widehat{ϵ}^l+\widehat{\kappa }^l`$, where the energy diagonal matrices $`\widehat{\kappa }^l=\{\delta _{nn^{}}\delta _{nn_l}\kappa _l\}`$ and $`ϵ_{n_l}`$ are the discrete energies closest to the arguments $`ϵ_l`$, $`l=1,2`$. It is then a straightforward matter to verify that the definition of the Keldysh partition function implies $`R(ϵ_1,ϵ_2)=\mathrm{\Delta }^2/(2\pi ^2)\mathrm{Re}_{\kappa _1\kappa _2}^2|_{\kappa =0}Z[\widehat{\kappa }]\mathrm{dis}`$, where ’dis’ stands for the disconnected part of the functional average. To keep the presentation simple we will focus in the quantitative analysis of this expression on the contribution of the spatial zero mode $`Q(𝐫)Q=\mathrm{const}.`$ (which ought to reproduce WD statistics). The inclusion of the spatially fluctuating modes, which is straightforward and does not introduce conceptually new elements, will be briefly discussed in the end.
The key to understanding the structure of the zero mode integration lies in the observation that the (source-free) effective action $`iS[Q]=i\pi \mathrm{\Delta }^1\mathrm{tr}\{\widehat{ϵ}Q\}`$ possesses a multitude of $`2^{2K}`$ isolated saddle points. Indeed, any of the (energy diagonal) configurations $`\mathrm{\Lambda }=\mathrm{diag}(\pm 1,\pm 1,\mathrm{},\pm 1)`$ solves the stationary phase equation of the model, $`\delta _QS[Q]=\pi \mathrm{\Delta }^1[\widehat{ϵ},Q]=0`$. In what follows we are going to show that a Gaussian integration around these saddle points produces WD statistics.
To prepare the integration, consider the contribution of any saddle point $`\mathrm{\Lambda }`$ with $`K+p`$ entries $`+1`$ and $`Kp`$ entries $`1`$. We first re-order these elements (through some global unitary transformation) such that $`\mathrm{\Lambda }`$ assumes the form $`\mathrm{\Lambda }=\mathrm{diag}(1,\mathrm{},1,1,\mathrm{}1)`$. Next, a set of field configurations weakly fluctuating around $`\mathrm{\Lambda }`$ is introduced through $`Q=T\mathrm{\Lambda }T^1`$, where the unitary rotation matrices $`T=\mathrm{exp}\left(\genfrac{}{}{0pt}{}{0^{}B}{B^{}\mathrm{\hspace{0.17em}0}}\right)`$, $`B`$ is a $`(K+p)\times (Kp)`$-dimensional complex generator matrix, and the block decomposition corresponds to the signature of $`\mathrm{\Lambda }`$.
In principle we should now proceed by expanding the action to second order in $`B`$ and integrate. Fortunately, however, there is no need to carry out this program for all $`2^{2K}`$ saddle points explicitly. The reason is that among the entity of saddle points $`\mathrm{\Lambda }`$, there is one element $`\mathrm{\Lambda }_0^{l,l^{}}(1)^{l+1}\delta ^{l,l^{}}`$ that plays a distinguished role. What makes $`\mathrm{\Lambda }_0`$ special is that, unlike the other $`\mathrm{\Lambda }`$’s, its structure is compatible with the signature of the imaginary increments of $`\widehat{ϵ}`$. Building on this feature, previous analyses of the Keldysh $`\sigma `$-model indeed focused on a perturbative expansion around $`\mathrm{\Lambda }_0`$ and did not take the other saddle points into account.
As a warm-up, let us outline how an integration around the standard saddle point produces the unit–normalization of the source–free partition function $`Z[\widehat{\kappa }=0]=1`$. Substituting $`Q=T\mathrm{\Lambda }_0T^1`$ into the zero mode action and expanding to second order in $`B`$ we obtain the quadratic action
$$iS_{\mathrm{\Lambda }_0}^{(2)}[B]=\frac{i\pi }{\mathrm{\Delta }}\left[\mathrm{tr}\{\widehat{ϵ}\mathrm{\Lambda }_0\}\underset{n,n^{}}{}(ϵ_n^+ϵ_n^{}^{})|B_{nn^{}}|^2\right].$$
(2)
Integration over $`B`$ then leads to
$$Z_0=\mathrm{const}.\times e^{i\pi \mathrm{\Delta }^1\mathrm{tr}\{\widehat{ϵ}\mathrm{\Lambda }_0\}}F_0;F_0=\underset{n,n^{}}{}(ϵ_n^+ϵ_n^{}^{})^1$$
as the contribution of $`\mathrm{\Lambda }_0`$ to $`Z[0]`$. Due to the causality property, $`\mathrm{tr}\{\widehat{ϵ}\mathrm{\Lambda }_0\}=0`$. Similarly, $`F_0=\mathrm{exp}\{_{nn^{}}\mathrm{ln}(ϵ_n^+ϵ_n^{}^{})\}=\mathrm{exp}\{0\}=1`$, where the presence of the imaginary increments guarantees that the branch cut singularity of the logarithm is not touched. (We re-emphasize that, at this stage, the causality rule has the mere status of a working assumption. In part II we will make up for its precise formulation, and show that the proper ultraviolet extension of the fluctuation determinant fixes the unspecified ’const.’ to unity.) Combining factors, we find $`Z_0=1`$. As a corollary we remark that $`Z_0=1`$ implies vanishing of the total contribution of all non–standard saddle points to $`Z[0]`$. Before turning to these other saddle points, let us discuss the contribution $`R_0`$ of the standard saddle point to the spectral correlation function. A straightforward expansion of $`Z[\widehat{\kappa }]`$ to first order in $`\kappa _1`$ and $`\kappa _2`$ yields $`R_0=\frac{1}{2}_{nn^{}}B_{n_1n}B_{nn_1}^{}B_{n_2n^{}}^{}B_{n^{}n_2}_B`$, where $`\mathrm{}_B`$ stands for the Gaussian average over the action (2). Integration over $`B`$ then leads to $`R_0=1/(2s^2)`$, where $`s=\pi \omega ^+/\mathrm{\Delta }`$ and $`\omega ^+ϵ_{n_1}^+ϵ_{n_2}^{}`$.
We next turn to the discussion of the other saddle points. In fact we will focus on just a single non-standard saddle point $`\stackrel{~}{\mathrm{\Lambda }}`$, namely the one where the signs of the two entries corresponding to our reference energies $`ϵ_{n_{1,2}}`$, are flipped: $`\stackrel{~}{\mathrm{\Lambda }}\mathrm{\Lambda }_02\delta ^{l1}\delta _{nn_1}+2\delta ^{l2}\delta _{nn_2}`$. From the analysis of $`\stackrel{~}{\mathrm{\Lambda }}`$, the role of all the other saddle points will become clear. The re-ordering needed to bring $`\stackrel{~}{\mathrm{\Lambda }}`$ into the canonical form implies that the action $`iS_{\stackrel{~}{\mathrm{\Lambda }}}^{(2)}`$ differs from $`iS_{\mathrm{\Lambda }_0}^{(2)}`$ in two respects: first, the contribution of the saddle point itself $`iS[\stackrel{~}{\mathrm{\Lambda }}]=i\pi \mathrm{\Delta }^1\mathrm{tr}\{\widehat{ϵ}\stackrel{~}{\mathrm{\Lambda }}\}=i\pi \mathrm{\Delta }^12\omega ^+`$ no longer vanishes. Second, in the fluctuation contribution, the two energy arguments $`ϵ_{n_2}`$ and $`ϵ_{n_1}`$ are exchanged. Given these structural changes we find it more convenient to reverse the order of the evaluation of the functional integral: first integrate out fluctuations, then expand in the sources. The integration over $`B_{nn^{}}`$ leads to a fluctuation factor $`\stackrel{~}{F}`$ similar to $`F_0`$ above, only that $`ϵ_{n_1}`$ and $`ϵ_{n_2}`$ are exchanged and coupled to the respective sources $`\kappa _{1,2}`$:
$$\stackrel{~}{F}=F_0\frac{ϵ_{n_1}ϵ_{n_2}}{ϵ_{n_2}^\kappa ϵ_{n_1}^\kappa }\underset{nn_1}{}\frac{ϵ_nϵ_{n_2}}{ϵ_nϵ_{n_1}^\kappa }\underset{n^{}n_2}{}\frac{ϵ_{n_1}ϵ_n^{}}{ϵ_{n_2}^\kappa ϵ_n^{}}.$$
(3)
We next add to the products the “missing” factors $`n=n_{1,2}`$ and use that the now unconstrained products over $`n,n^{}`$, as well as the factor $`F_0`$, give unity (causality). This leads to the result
$$\stackrel{~}{F}=\frac{ϵ_{n_1}ϵ_{n_2}}{ϵ_{n_2}^\kappa ϵ_{n_1}^\kappa }\frac{ϵ_{n_1}ϵ_{n_1}^\kappa }{ϵ_{n_1}ϵ_{n_2}}\frac{ϵ_{n_2}^\kappa ϵ_{n_2}}{ϵ_{n_1}ϵ_{n_2}}=\frac{\kappa _1\kappa _2}{\omega ^{+2}},$$
(4)
where the last equality is valid to leading non-vanishing order in the sources, $`\kappa _1,\kappa _2`$. Notice the proportionality of $`\stackrel{~}{F}`$ to $`\kappa _1\kappa _2`$. This implies (a) that the sources in the action $`S[\stackrel{~}{\mathrm{\Lambda }}]`$ may be set to zero (we are differentiating the functional at $`\kappa _l=0`$) and (b) that only the non–standard saddle point $`\stackrel{~}{\mathrm{\Lambda }}`$ contributes to the correlation function. Indeed, for any other non-standard saddle point one or several signatures corresponding to energy arguments $`ϵ_nϵ_{n_1,n_2}`$ are changed. These energies are not coupled to sources. Repeating the steps outlined above one finds that the $`B`$-integral around these saddle points gives zero (i.e. the $`\kappa 0`$ limit of the factor $`\stackrel{~}{F}`$ above). The same argument also shows that the non-standard saddle points do not contribute to $`Z[0]`$. Differentiating Eq. (4) w.r.t. $`\kappa _{1,2}`$ and adding the contribution of the standard saddle point, we obtain the well known result
$$R(\omega )=\frac{1}{2}\mathrm{Re}\frac{1\mathrm{exp}(2is)}{s^2}=\left(\frac{\mathrm{sin}s}{s}\right)^2$$
(5)
for the two point correlation function of the zero mode theory. We finally mention that the inclusion of spatially fluctuating modes into the formalism (a) does not change the saddle point structure and (b), after Gaussian integration, leads to a renormalization $`\stackrel{~}{F}𝒟(\omega )\stackrel{~}{F}`$, where
$$𝒟(\omega )\underset{𝐪0}{}\frac{(D𝐪^2)^2}{(D𝐪^2)^2+\omega ^2}$$
(6)
and $`𝐪`$ are the quantized momenta associated to fluctuations in a finite size system. Combining this with the contribution of $`\mathrm{\Lambda }_0`$, we reproduce the familiar result for the level statistics of the unitary disordered electron gas. Since this result was obtained in a saddle point approximation, its validity is restricted to energies $`\omega \mathrm{\Delta }`$. Indeed, in the opposite limit the fluctuation modes $`B`$ become too light to be treated in the Gaussian approximation. The fact that the zero mode result Eq. (5) is actually exact for any $`\omega `$ is a “coincidence” (see however Ref. ), specific to the unitary ensemble.
Part II: The analysis above crucially relied on the causality postulate: summations over continuous functions of $`ϵ`$ produce zero. As mentioned above, this feature could not be made explicit because the validity of the action (1) is limited to low energies. In this second part of the Letter we are going to construct an energetically enlarged formulation whereby large and small $`ϵ`$ are treated on the same footing. This will make the causality property manifest. At the same time we will see why the pragmatic scheme employed above is sufficient for the calculation of low energy observables.
In order to not unnecessarily complicate the discussion we will formulate this part of the analysis for an $`N`$-dimensional random matrix theory (RMT) Hamiltonian $`H=\{H_{\mu \nu }\}`$ defined through the Gaussian correlation law $`H_{\mu \nu }=0`$ and $`H_{\mu \nu }H_{\nu ^{}\mu ^{}}=N^1\delta _{\mu \mu ^{}}\delta _{\nu \nu ^{}}`$. The advantage gained is that $`H`$ has neatly defined universal large energy properties, i.e. that we will not need to consider the non-universal UV asymptotics of the free–electron Hamiltonian underlying Eq. (1). Later on we will argue that, as far as the present discussion is concerned, the specific modeling of the Hamiltonian is of no relevance. Further, to deal with a manifestly $`UV`$-regularized model, we compactify our energy variables. This can be done by discretizing the temporal Keldysh contours $`𝒞^l`$ to a lattice of small spacing $`\delta _t`$. As a result, the energy variables $`ϵ_n\{\delta _ϵ,2\delta _ϵ,\mathrm{},K\delta _ϵ\}`$, where $`K2\pi /(\delta _ϵ\delta _t)1`$.
The effective Keldysh action for the RMT model can be obtained by a straightforward adaptation of previous derivations of the $`\sigma `$–model for RMT Hamiltonians to the specifics of the Keldysh $`\sigma `$-model. As an intermediate result one obtains the partition function $`Z[\widehat{\kappa }]=𝒟Q\mathrm{exp}\{iS[Q,\widehat{\kappa }]\}`$ where the action $`iS[Q,\widehat{\kappa }]=N\left(\frac{1}{2}\mathrm{tr}(Q^2)\mathrm{tr}\mathrm{ln}\left[\widehat{z}+Q\right]\right)`$ with $`\widehat{z}=\sigma _3\delta _t^1(e^{i\sigma _3\delta _t\widehat{ϵ}}1)`$ ($`\sigma _3`$ acts in the Keldysh $`2\times 2`$ space). At this stage $`Q=\{Q_{ϵ,ϵ^{}}^{l,l^{}}\}`$ is a $`2K`$-dimensional matrix that has been introduced to decouple the $`H`$-averaged action (no constraint $`Q^2=𝟙`$ as yet). The unusual phase–type appearance of the energy argument is due to the time discretization. However, in the limit $`ϵ\delta _t1`$, $`\widehat{z}i\widehat{ϵ}`$ and we retrieve the standard form of the RMT $`\sigma `$-model action . We next subject the action to a saddle point analysis (stabilized by the parameter $`N1`$). Variation of the action w.r.t. $`Q`$ yields the quadratic equation $`Q=(\widehat{z}+Q)^1`$. The $`2^{2K}`$–fold degenerate energy diagonal set of solutions, $`\mathrm{\Lambda }^{(l)}(ϵ_n)\mathrm{\Lambda }_{ϵ_n,ϵ_n}^{l,l}`$, is given by
$$\mathrm{\Lambda }^{(l)}(ϵ_n)=\frac{1}{2}\left[z_l(ϵ_n)\pm i(z_l^2(ϵ_n)4)^{1/2}\right].$$
(7)
This expression determines the entire spectral structure of the model. First, notice that for energies $`ϵ_n1\delta _t^1`$, $`\mathrm{\Lambda }^{(l)}(ϵ_n)=\pm 1+𝒪(ϵ_n)`$. This means that the solutions $`\mathrm{\Lambda }^{(l)}(ϵ_n)`$ represent an UV extension of the saddle points $`\mathrm{\Lambda }`$ discussed in part I. We next ask whether the two sign alternatives in Eq. (7) are equivalent or whether the model has a preferred choice. Indeed, the latter is the case: for energies $`ϵ_n1`$ greatly in excess of the width of the spectrum, $`\mathrm{\Lambda }^{(l)}(ϵ_n)`$ must approach zero – the free Gaussian saddle point of the non–disordered model. This condition determines $`\mathrm{\Lambda }_0(ϵ_n)=\frac{1}{2}[\widehat{z}(ϵ_n)+\sigma _3(\widehat{z}^2(ϵ_n)4)^{1/2}]`$ as the canonical solution. For low energies $`\mathrm{\Lambda }_0(ϵ_n)`$ reduces to the saddle point $`\mathrm{\Lambda }_0`$ discussed above. This saddle point has the important property, $`_n[\mathrm{\Lambda }_0(ϵ_n)]^k=0`$, $`k`$ a positive integer. The outline of the proof is as follows: (due to the presence of a finite imaginary increment) the summation over $`ϵ_n`$ is equivalent to an integration of the variable $`w=\mathrm{exp}\{i\delta _tϵ_n\}`$ over the complex unit circle. It is straightforward to verify that for $`|w|1`$, $`\mathrm{\Lambda }_0(w)`$ is analytic (the branch cut singularity of the square root lies inside the unit circle). Further, for $`|w|1`$, the integrand decays as $`w^{(k+1)}`$. From Cauchy’s theorem we conclude that the summation gives zero $`\mathrm{}`$. Summarizing, we have found a complex saddle point structure which extends the low energy saddle points $`\mathrm{\Lambda }=\pm 1`$ discussed in part I into the UV regime. The complex structure of the theory entails the existence of a ’natural’ saddle point $`\mathrm{\Lambda }_0`$. Existence and behavior of $`\mathrm{\Lambda }_0`$ vitally depend on the large energy, $`ϵ1`$, asymptotics of the theory.
To more explicitly establish contact with the low energy regime discussed in part I, we next introduce fluctuations around the saddle points (7). Defining $`Q=\mathrm{\Lambda }+P`$, where $`P`$ is some Hermitian matrix, and expanding the action $`S[\mathrm{\Lambda }+P]`$ to second order in $`P`$ we obtain $`iS_\mathrm{\Lambda }^{(2)}[P]=\frac{N}{2}\mathrm{tr}(P^2+P\mathrm{\Lambda }P\mathrm{\Lambda })`$. This is the UV-extension of the low energy action discussed in part I. Indeed, substituting the small $`ϵ`$ asymptotics of $`\mathrm{\Lambda }`$ and using that for the RMT model, $`\mathrm{\Delta }=\pi /N`$ , we find that for energies $`ϵ1`$, $`S_\mathrm{\Lambda }^{(2)}`$ reduces to the action, Eq. (2).
One can now step by step repeat the analysis that led to the correlation function of part I. The only difference is that instead of energy denominators $`(ϵ_n^+ϵ_n^{}^{})`$ constructions like $`1+\mathrm{\Lambda }_0^{(1)}(ϵ_n)\mathrm{\Lambda }_0^{(2)}(ϵ_n^{})`$ appear. Due to the compact phase-type appearance of the energy arguments in $`\mathrm{\Lambda }_0(ϵ_n)`$ all energy summations converge. Further, the properties of $`\mathrm{\Lambda }_0`$ discussed above imply the vanishing of expressions like $`_nf(\mathrm{\Lambda }_0(ϵ_n))`$ where $`f`$ may be any analytic function. This implements the causality principle. In parentheses we note that the detailed execution of this program yields a unit normalization of the partition function, without any undetermined prefactors.
The discussion above provokes the obvious question whether the conclusions drawn from the high energy asymptotics of the action are specific to the random matrix model? We believe that the answer is negative. Recapitulating the sequence of arguments, it is evident that everything hinges on the absence of singularities outside the complex unit circle defined through the compactified energy argument. This in turn is a guaranteed feature as long as the single particle retarded SCBA Green function of the model system has a well defined pole structure below the real axis. In practice, for condensed matter systems with non-universal high energy behavior, it may be difficult to find closed solutions of the mean field equations that manifestly display this feature. We believe, however, that this is a practical rather than a principle difficulty. Summarizing, the main goal of part II was to show that the causality feature underlying this and previous analyses of the effective Keldysh action can be made an explicit ingredient of the model, on the expense of including large energy asympotics. In practical applications of the formalism, one will normally use causality feature pragmatically, as exemplified in part I. Yet it cannot be excluded that situations arise, where large scale spectral structures become essential.
To conclude, on the example of level statistics, we have demonstrated how non–perturbative quantum effects may be incorporated into the framework of the dynamic Keldysh $`\sigma `$–model. In many respects, the Keldish scheme appears to be simpler and more transparent than its relatives, replica and SUSY. We expect our methods to be useful in the analysis of interaction phenomena in disordered electronic systems.
We benefited a lot from numerous discussions with A. Andreev and M. Janssen. A.K. was partially supported by the BSF-9800338 grant.
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# Tritium Beta Decay, Neutrino Mass Matrices and Interactions Beyond the Standard Model
## 1 Introduction
One of the outstanding problems of subatomic physics is that of neutrino mass. While this problem is being attacked on many fronts, the most direct approach would seem to be the use of Tritium beta decay to make a direct measurement of the mass of the electron anti-neutrino . However, as was known long before the advent of the Standard Model (SM), the exact form of the electron spectrum depends on the Lorentz structure of the weak current involved in beta decay. This observation was recently revisited , wherein it was shown that possible interference between the usual SM $`SU(2)_W`$ current and a (weaker) non-SM current having a different Lorentz structure<sup>1</sup><sup>1</sup>1Note that, here and throughout this paper, we are looking at extensions of the SM which are distinguishable at low energies. This does not include additional Left-chiral interactions . can be exacerbated by the presence of a CKM-like mixing in the lepton sector. As we shall show in this paper, such effects could produce a signal that might be interpreted as a negative value of the square of the neutrino mass in an analysis which did not take interference into account. Whatever the eventual resolution of the current experimental situation, our discussion makes clear the fact that experimental analyses should not prejudice the result by assuming solely a SM structure.
The concept of using a detailed measurement of the high energy portion of the electron spectrum from nuclear beta decay to determine the mass of the neutrino was introduced in Fermi’s original paper. In that work, he assumed that the interaction was due to a vector current. This gives rise (as we discuss below) to a different dependence on neutrino mass than the $`VA`$ current of the SM, as evidenced qualitatively by the curves shown in his paper. As the possibility of other Lorentz current structures was considered, Koefed-Hansen pointed out the need to know the exact nature of the currents to properly extract a neutrino mass. Following the establishment of the dominance of $`VA`$, Enz recast the problem in terms of possible interference between that current and other (weaker) Lorentz currents and Jackson, Treiman and Wyld extended the discussion to include possible time reversal invariance violation.
These earlier discussions were based on only one neutrino. From LEP data we now know that there is a net of three light neutrinos with full SM coupling to the $`Z^0`$. As we show below, both because the possible interferences depend on the neutrino mass and because Tritium beta decay experiments do not measure the neutrinos, it is appropriate to frame the discussion in terms of mass eigenstates. As this is more easily visualized with Dirac neutrinos, we begin our discussion with them, expanding to the more general and more widely accepted case of Majorana neutrinos in a following section. We then examine the effect that any such interference will have on the extraction of the mass parameters.
In a subsequent section, we present the effects of a possible background scalar field, interacting only with neutrinos , on the Tritium spectrum. We also examine the possibility that the interference effects could manifest themselves at the upper end of the neutrino spectrum from other nuclear beta decays . Following this, we examine the implications of the new interaction currents for neutrino neutral current scattering, which devolve from the necessity of “sterile” neutrino interactions. In the final section we reiterate our conclusions and that experimentalists should use our functional form for data analysis. We emphasize that this form must be used to avoid introducing unwarranted theoretical prejudice into the interpretation of the experimental results.
## 2 Dirac neutrinos
To present the general concepts that underlie our calculation in the most accessible form, we first do the analysis for Dirac neutrinos. In later sections we shall deal with the more general case of Majorana neutrinos. By dealing with Dirac neutrinos, all of the manipulations to calculate total rates from Feynman diagrams are standard and can be found in any text on relativistic quantum mechanics.
### 2.1 Mass eigenstates
For the purpose of the discussion in this section, we assume the existence of three massive Dirac neutrinos, $`\nu _i^D`$, with Dirac masses $`m_i^D`$. All three masses are less (really much less) than half the mass of the $`Z^0`$. We allow for the possibility that, as in the quark sector, current eigenstates and mass eigenstates are not necessarily identical. In particular, we shall refer to the $`\nu _f`$, $`f=(e,\mu ,\tau )`$ as those linear combinations of mass eigenstates which are coupled to the corresponding charged leptons by the $`SU(2)_L`$ currents of the SM. These current eigenstates are expressed in terms of the mass eigenstates as
$$\nu _f=\underset{i}{}\mathrm{cos}\theta _f^i\nu _i$$
where the $`\mathrm{cos}\theta _f^i`$ are the direction cosines in the coordinate system spanned by the mass eigenstates. (If two or more masses are degenerate, we choose an orthogonal set.) Under these assumptions, the current eigenstates also form an orthogonal set. This is depicted in Fig. 1.
If there are additional objects, besides the known massive weak vector bosons, that produce charge changing currents which couple to both the quarks and the leptons, we may define neutral leptons which are current eigenstates of this new interaction in exactly the same way, denoted by $`\widehat{\nu }_f`$, $`f=(e,\mu ,\tau )`$, where now
$$\widehat{\nu }_f=\underset{i}{}\mathrm{cos}\widehat{\theta }_f^i\nu _i$$
and one would naturally expect that, in general,
$$\widehat{\theta }_f^i\theta _f^i.$$
### 2.2 Interaction Hamiltonian
Interactions beyond the SM must be weaker, at low energies, than the usual Left-chiral SU(2) to avoid serious conflict with existing data. Presumably this is due to the boson mediating the interaction being much heavier than the known W’s and Z, and/or the coupling constant being smaller than that for the SM. Absent a particular Grand Unified Theory in which one wishes to embed the SM and the new interactions, that is all one can say.
For low energy physics, like nuclear $`\beta `$-decay in general and Tritium $`\beta `$-decay in particular, such new interactions can only appear as effective currents in the four fermion formulation of the theory with the usual space time structures of $`S,P,T,V`$ or $`A`$. Given the dominance of the SM, it is reasonable to recast this as $`S_L,S_R,T,R`$ or $`L`$, where $`R(V+A)`$ and $`L(VA)`$, with a similar construction for $`S_L`$ and $`S_R`$ from $`S`$ and $`P`$. The effect on the spectrum displayed below only occurs if there are additional currents in the lepton sector that are different from $`L`$. Corresponding changes in the hadron currents affect only the scale of each new contribution. To emphasize this fact, we employ an unconventional notation; more conventional descriptions are given, for example, by Enz .
The effective low energy Hamiltonian for semi-leptonic decays is
$$H_I=\underset{\alpha ,\beta =S_L,S_R,R,L,T}{}G^{\alpha \beta }\underset{f}{}(J_{h\alpha }^{}J_{f\beta }+h.c.)$$
(1)
where, for example,
$$J_{f\alpha }=\overline{\psi _f}\mathrm{\Gamma }_\alpha \psi _{\nu _f}$$
(2)
with $`\psi _f`$ representing a charged lepton of a given flavor, $`\psi _{\nu _f}`$ a neutral lepton associated with that charged lepton through the particular interaction (see the discussion below for more detail). A similar construction can be made on the hadron side. Note that, in this convention, the first Greek index in Eq.(1) refers to the hadron current. Explicitly,
$`\mathrm{\Gamma }_{S_L}`$ $`=`$ $`(1\gamma ^5)`$
$`\mathrm{\Gamma }_{S_R}`$ $`=`$ $`(1+\gamma ^5)`$
$`\mathrm{\Gamma }_R`$ $`=`$ $`\gamma ^\mu (1+\gamma ^5)`$
$`\mathrm{\Gamma }_L`$ $`=`$ $`\gamma ^\mu (1\gamma ^5)`$
$`\mathrm{\Gamma }_T`$ $`=`$ $`[\gamma ^\mu ,\gamma ^\nu ]/2.`$ (3)
Most off diagonal combinations ($`\alpha \beta `$) in the sum vanish, the exceptions being for $`(S_R,S_L)`$ and $`(R,L)`$. For nuclear beta decay in the SM, only $`\beta =L`$ and $`\alpha =L,R`$ survive. The relation to the basic parameters of the SM is given by
$$G^{RL}+G^{LL}=V_{ud}\frac{\pi \alpha _W}{\sqrt{2}M_W^2},$$
(4)
where $`V_{ud}`$ is the appropriate element of the hadronic CKM matrix , $`\alpha _W`$ is the fine structure constant for the $`SU(2)_W`$ of the SM, $`\alpha _W=\frac{g_W^2}{4\pi }`$, where $`g_W`$ is the coupling constant, and $`M_W`$ is the mass of the usual $`W^\pm `$. To complete the connection with conventional notation, note that
$$(G^{RL}G^{LL})=\frac{G_A}{G_V}(G^{RL}+G^{LL}).$$
(5)
The reason that the SM $`SU(2)_L`$ induces an hadronic Right-chiral coupling is that the quarks are confined and their wave functions cause $`G_A`$ to be renormalized with respect to $`G_V`$. If there is an $`SU(2)_R`$, mediated by a heavier vector boson or having weaker coupling constant, or both, one still expects the hadronic current to be modified. However, in parallel to the case of the SM interaction, the leptonic currents would be expected to be pure $`V+A`$. In fact, one also expects hadronic renormalization effects for $`S`$, $`P`$ and $`T`$. In addition, there will be, in general, a separate “CKM” matrix in the quark sector for each interaction, mirroring the discussion of mass and current eigenstates for neutrinos given above.
The various currents in the effective interaction Hamiltonian may be generated by several different assumptions about physics beyond the SM. The most prominent are the exchange of a charge-changing scalar, arising, for example, in supersymmetric models , or the existence of a vector boson coupled to Right-chiral fermions which may or may not mix with the W bosons of the SM. These are shown in Fig. 2, where we have included those cases which may impact Tritium beta decay. Although antisymmetric tensor interactions have been proposed in some (mostly experimental) contexts , they would normally be expected to arise in higher-loop order or by Fierz transformation from scalar leptoquark interactions and so be exceptionally weak. In particular, Voloshin has obtained a very stringent bound on the strength of such a tensor interaction ($`10^4G_F`$), which applies unless a precise orthogonality holds for electron neutrino states defined by the different interactions.
As we will show below, it will be difficult to distinguish amongst the effective currents. Given an effect in Tritium from some such current, it will be even harder to trace that effect to a particular diagram in Fig. 2. However, this does mean that it makes sense to pursue such effects in Tritium beta decay whatever one’s prejudice about a particular source of “new physics”.
In terms of the diagrams of Fig. 2, we may define a new fine structure constant, $`\widehat{\alpha }`$ (to go with b) or the first diagram of c)) as $`\widehat{\alpha }=\frac{\widehat{g}^2}{4\pi }`$. Then the ratio of the effective strength of this interaction at low energy to that of the SM $`SU(2)_L`$ is given by $`\widehat{\rho }_X=\frac{\widehat{g}^2}{g^2}\frac{M_W^2}{M_X^2}`$, where $`M_X`$ is the mass of the non-SM boson being exchanged.
### 2.3 Interference terms
If there are additional bosons that couple to the Left-chiral Vector fermion current , they can only renormalize $`L`$ and cannot lead to any new structure near the end point. For this study of Tritium beta decay all hadronic currents are evaluated in the approximation of no nuclear recoil, which eliminates a pure Pseudo-scalar current.
Nontheless, we allow for an arbitrary mixture of Scalar and Pseudo-scalar, expressed as Left- and Right-chiral Scalar currents. Both couple hadronically through S and couple independently to the lepton current. Consequently, we need to evaluate the effect on the Tritium beta spectrum of possible interferences between the SM Left-chiral Vector current and a Right-chiral Vector current, a Right-chiral Scalar current and a Left-chiral Scalar current. In this work, a Left-chiral Scalar current is defined as that Scalar current for which the Left-chiral projector acts on the neutrino field for the current defined as in Eq.(2), and a corresponding definition for the Right-chiral case.
Since the hadronic currents must be renormalized and the fitting procedure for any Tritium beta decay experiment fits the overall rate, we only quote here the relative factors from the lepton traces. In doing so, we make use of the fact that, when contracted with the hadron traces and integrated over the outgoing neutrino directions, the only terms that survive are those proportional to factors of the hadron mass or energy, which differ negligibly due to the small size of the momentum transferred to the final state hadrons (“no recoil” approximation):
$$\begin{array}{cc}LL(SM)\hfill & \hfill E_\nu E_\beta \\ RR\hfill & \hfill E_\nu E_\beta \\ LR+RL\hfill & \hfill 2m_\nu m_e\\ S_LS_L\hfill & \hfill E_\nu E_\beta \\ S_RS_R\hfill & \hfill E_\nu E_\beta \\ S_LS_R+S_RS_L\hfill & \hfill 2m_\nu m_e\\ TT\hfill & \hfill [E_\nu E_\beta m_\nu m_e]\\ LS_R+S_RL\hfill & \hfill 2m_\nu E_\beta \\ LS_L+S_LL\hfill & \hfill 2E_\nu m_e\\ LT+TL\hfill & \hfill 2[E_\nu m_em_\nu E_\beta ]\end{array}$$
The negative sign of the terms with a factor of $`m_\nu `$ arises from the fact that the neutral lepton in Tritium beta decay is an anti-neutrino.
For the interactions involving L (and by symmetry, R) currents, the Lorentz inner product of the hadron trace with the lepton trace produces an overall factor of $`(G_V^2+3G_A^2)`$, and this same factor occurs for the interference terms between the $`W_L`$ and $`W_R`$ exchanges.
The case of of $`W_LW_R`$ mixing is more complicated. It must include a factor like this for the interference term with the SM in which the $`W_R`$ couples directly to both leptons and hadrons as above, but is more complicated and also involves other combinations, such as $`(G_V^23G_A^2)`$ for the term involving the interference between the SM and the amplitude in which the mixed propagator couples via $`W_L`$ to the hadrons and $`W_R`$ to the leptons.
The corresponding result for the interference terms of a scalar current with the SM currents is proportional to $`G_V^2`$ alone. This means that the scalar interactions are somewhat less efficient at producing new effects for the same relative strength as Right-chiral interactions.
Note that all interference effects of tensor interactions with the SM (which effects are proportional to factors of $`3G_A^2`$) are encompassed by the scalar terms and so we will not explicitly discuss tensor contributions further in this paper.
For our main purpose here, we need only consider the possibilities of interference between the SM current and either a Right-chiral Vector current (for Dirac neutrinos) or a Right-chiral Scalar current for either Dirac or Majorana neutrinos. (The effect of interference with the Left-chiral Scalar current appears in Sec.6, and an additional tensor interaction is just a linear combination of the Left-chiral Scalar and Right-chiral Scalar terms.)
### 2.4 Pre-existing limits
The best limits on the relative strengths of additional currents often make use of interference effects , and are deduced assuming that the same neutrino is produced with the electron in all cases. As we discussed above, this is not a necessarily valid assumption. In fact, the interference effects need to be examined for each mass eigenstate and, depending upon the details of the interaction, could vanish over most of the spectrum even if the effective coupling does not.
As an example, consider the possible effect on the spectrum due to the interference between a Left-chiral Scalar and the SM Left-chiral interaction. As shown in the previous section, for each mass eigenstate $`\nu _i`$, this interference is proportional to $`\widehat{\rho }_{S_L}E_\nu m_e\mathrm{cos}\theta _e^i\mathrm{cos}\widehat{\theta _e^i}`$. For most applications, the neutrino masses are negligible compared to $`E_\nu `$, so $`q_\nu E_\nu `$ and thresholds may be ignored. This means that, when one sums over the mass eigenstates, the result is proportional to $`\widehat{\rho }_{S_L}\times \mathrm{cos}(x)`$ where $`x`$ is the angle between $`\widehat{\nu _e}`$ and $`\nu _e`$, not simply to $`\widehat{\rho }_{S_L}`$.
Recent limits are given as
$`\widehat{\rho }_R`$ $``$ $`.07`$
$`\widehat{\rho }_{S_R}`$ $``$ $`.1`$
$`\widehat{\rho }_{S_L}`$ $``$ $`.01`$
where the analysis has implicitly assumed unit value for $`\mathrm{cos}(x)`$.
## 3 Majorana neutrinos
In this section we discuss the more general possibility, arising from the lack of any conserved charge for the neutrinos (at least at the low energy scales below one-half the mass of the $`Z_0`$) that the mass eigenstates correspond to Majorana neutrinos.
### 3.1 Massive Weyl neutrinos
We follow the development presented by Ramond . The basic object is a two complex component Weyl spinor which can be represented under the Lorentz group in either of two inequivalent irreducible representations, labelled conventionally as $`(1/2,0)`$ and $`(0,1/2)`$. These two representations of the same Weyl field are isomorphic; the isomorphism connecting them is the operation of Charge conjugation. The usual Dirac spinor in the Weyl (or chiral) representation is constructed from one Weyl spinor in the $`(1/2,0)`$ representation and an independent Weyl spinor in the $`(0,1/2)`$ representation having the correct charges under appropriate interactions. A (Dirac) mass term coupling such representations allows for a common global or local phase and, in particular, allows for a Noether current that guarantees the conservation of particle number.
In this same representation, the four component representation of a Majorana spinor is constructed by taking, for the Weyl spinor in the $`(0,1/2)`$ representation, plus or minus times the Charge conjugate of the Weyl spinor in the $`(1/2,0)`$ representation. This produces even or odd Charge conjugation eigenstates, each of which contains exactly the information contained in the original Weyl spinor. The same procedure as used for a Dirac mass term now produces a Majorana mass term, only now, by construction, the two pieces of the spinor have conjugate behavior under a global or a local phase and one cannot construct a Noether current related to fermion number conservation. Also, the interaction charges are opposite in the two parts of the spinor, so this construction is not appropriate for any fermion with a conserved charge. Since both Left chiral $`(1/2,0)`$ and Right chiral $`(0,1/2)`$ representations are present, this Majorana spinor is adequate to describe all of the SM weak interactions and, if the mass is non-zero, can mediate neutrinoless double beta decay.
### 3.2 Pure Weyl or “see-saw” neutrinos
Appealing only to the SM without extensions, one need only consider three Weyl neutrino fields associated with three (Majorana) mass eigenstates, possibly rotated from the current eigenstates. By convention, these masses are labelled $`m_L`$. With respect to Tritium beta decay, there can be no interference with Right-chiral Vector currents or with Right-chiral Scalar currents. On the other hand, interference with the Left-chiral Scalar current is possible, but, since that current is the most constrained by other data, the possibility of describing the data is less viable.
There is another nagging problem with pure, massive Weyl neutrinos. This has to do with the possibility, inherent in the SM, that the $`SU(2)_L`$ symmetry could be restored at high temperature. In this scenario, the SM Higgs loses its vacuum expectation value ($`vev`$), all four weak bosons become massless and the weak charge is conserved. The last statement is contradicted by the existence of the Majorana mass, which is certainly not connected to the standard Higgs.
To achieve such a mass, one generally introduces a new scalar field, the Majoron , the sole purpose of which is generate a Majorana neutrino mass via the $`vev`$ of the Majoron field. As this $`vev`$ must be small to avoid distortion of the ratio of $`W`$ and $`Z`$ masses beyond experimental constraints, it may also be expected to evaporate as well when the $`SU(2)_L`$ symmetry is restored, although this is not guaranteed.
At some level, this is addressed by the usual see-saw . In that case one builds two distinct Majorana neutrinos, one from the Weyl field described above and one from a Weyl field that could be used with the first to build Dirac neutrinos. One then assumes that the usual Higgs produces some Dirac masses ($`m_D`$) (presumably of the order of charged lepton masses, although that is not critical) and that the new Weyl field, which is sterile to all known interactions, develops a Majorana mass through some means that does not affect the SM. That mass is conventionally termed $`M_R`$, and is actually a $`3\times 3`$ (or larger, depending on the full model) matrix in flavor space.
This has the well known pleasant consequence that the masses of the light (active) neutrinos are of the order of $`\frac{m_D^2}{M_R}`$, leading to very small neutrino masses. However, it should clearly be noted that there is no principle leading to this construction. There is also no guarantee that the rank of $`M_R`$ is three, so there could be vastly different patterns in different generations. In this description (assuming rank three), the three light neutrinos are essentially indistinguishable from the pure Weyl case insofar as Tritium beta decay is concerned. However, since $`m_D`$ is assumed to be proportional to the Higgs vacuum expectation value, the problem with symmetry restoration is avoided. We do note, however, that $`LR`$ symmetric models with a (relatively) low lying Right-chiral scale face the same problem of symmetry restoration.
### 3.3 Pseudo-Dirac neutrinos
For many years there has been discussion of the possibilities associated with a very small Majorana mass, either compared with Dirac masses or masses only off-diagonal in flavor space, the so-called “Pseudo-Dirac” case . In this case, which has recently been revisited in the literature , for Tritium beta decay there is really no distinction from the pure Dirac case. Any time evolution will be so slow that the implications for the spectra considered in this paper are negligible.
## 4 Implications for Tritium beta decay
We now examine the effect of the possibilities we have discussed on the analysis of the beta spectra in Tritium beta decay. We recognize that no attempt should be made to extract any parameters from the published literature, as proper analysis requires inclusion of specific experimental details. What we shall show is that, for parameter ranges which are not in conflict with other experimental results, the combination of non-vanishing neutrino masses and an interference at very low neutrino energy can cause an analysis, which assumes no interference, to produce negative values of $`m_\nu ^2`$. Furthermore, there is a dependence of the extracted value on the range of $`\beta `$ energies used for the fit (for differential spectra) or the $`\beta `$ bias energy (for integral spectra) which dependence matches that reported by some experiments .
### 4.1 Differential spectra
Our discussion in this paper is directed at those experiments using molecular Tritium as a source. Other sources require well known modifications to the discussion which are no different for our case than for the SM.
For the SM only, the differential spectrum, $`\frac{dN}{dE_\beta }`$, is already a sum of several individual spectra for each possible end point corresponding to a particular final energy of the $`{}_{}{}^{3}HeT`$ molecular ionic system. Strictly speaking, the sum should include an integral over the continuum of ionic breakup states, although the analysis may find that it is adequate to represent that continuum with a finite sum. This may be expressed as
$$\left(\frac{dN}{dE_\beta }\right)_{SM}=\underset{i}{}P_i\left(\frac{dN}{dE_\beta }(_0^i)\right)_{SM}$$
(6)
where $`P_i`$ is the probability of leaving the ionic system in the $`i^{th}`$ final state and $`_0^i`$ is the maximum energy available to the $`\beta `$ for that final state, assuming that $`m_\nu =0`$. At present, the $`P_i`$ are calculated and current experiments estimate that the theoretical uncertainty which this introduces constitutes a small part of the error budget. The end point of the spectrum (for $`m_\nu =0`$) is then $`_0^0`$ where $`i=0`$ denotes the ground state of the molecular ion. Of course, the expression in Eq.(6) is theoretical only. In fitting to a particular experimental spectrum, any energy dependent experimental corrections must either be included explicitly or allowed to vary, within reason, in the fitting procedure.
In this paper we shall illustrate the effects of possible interferences on the theoretical spectrum appropriate to $`_0^0`$. The extension to the full spectrum follows Eq.(6). For zero mass neutrinos, where the SM weak current eigenstate may be taken as the only neutrino produced, we have
$`\left({\displaystyle \frac{dN}{dE_\beta }}(_0^0)\right)_{SM}`$ $`=`$ $`KF(E_\beta )q_\beta q_\nu E_\beta E_\nu `$ (7)
$`=`$ $`KF(E_\beta )q_\beta E_\beta (_0^0E_\beta )^2,`$
where $`K`$ includes all those quantities, such as the nuclear matrix element, the source strength and the overall experimental efficiency, subsumed by the measured normalization, $`F(E_\beta )`$ is the Fermi function, $`q_\beta `$ and $`E_\beta `$ are the momentum and relativistic energy of the $`\beta `$ and $`q_\nu `$ and $`E_\nu `$ are the momentum and relativistic energy of the antineutrino. The second line reflects the assumptions that $`E_\nu =_0^0E_\beta `$ and that $`q_\nu =E_\nu `$ for a massless neutrino.
The simplest extension to the SM consists of assuming that there are no additional interactions, that there is still only one relevant neutrino (strictly, $`\overline{\nu }_e`$), but that this neutrino may have a mass. In this case we make the replacement $`q_\nu =\sqrt{E_\nu ^2m_\nu ^2}`$ and fit
$$\left(\frac{dN}{dE_\beta }\right)_1=\underset{i}{}P_i\left(\frac{dN}{dE_\beta }(_0^i)\right)_1$$
(8)
where, for example,
$`\left({\displaystyle \frac{dN}{dE_\beta }}(_0^0)\right)_1`$ $`=`$ $`KF(E_\beta )q_\beta E_\beta E_\nu \sqrt{E_\nu ^2m_\nu ^2}\mathrm{\Theta }(_0^0E_\beta m_\nu )`$ (9)
$`=`$ $`KF(E_\beta )q_\beta E_\beta (_0^0E_\beta )^2\sqrt{1{\displaystyle \frac{m_\nu ^2}{(_0^0E_\beta )^2}}}`$
$`\times \mathrm{\Theta }(_0^0E_\beta m_\nu ),`$
treating $`m_\nu ^2`$ as an additional parameter. $`\mathrm{\Theta }(z)`$ is the Heavyside function. As is well known, the reported values of $`m_\nu ^2<0`$ are an indication of an excess of counts near the endpoint over the expectation for zero mass, with all other parameters determined from the robust spectrum far below the end point.
The first obvious extension is to include the possibility of different mass eigenstates, with the SM current eigenstates being mixtures of those mass eigenstates denoted by $`\nu _k`$,
$$\overline{\nu }_e=\underset{k}{}\mathrm{cos}\theta _e^k\overline{\nu }_k$$
(10)
giving
$`\left({\displaystyle \frac{dN}{dE_\beta }}(_0^0)\right)_2`$ $`=`$ $`KF(E_\beta )q_\beta E_\beta (_0^0E_\beta )^2{\displaystyle \underset{k}{}}\mathrm{cos}^2\theta _e^k\sqrt{1{\displaystyle \frac{m_k^2}{(_0^0E_\beta )^2}}}`$ (11)
$`\times \mathrm{\Theta }(_0^0E_\beta m_k),`$
which does not improve the fit near the endpoint .
The next extension, which is the point of this paper, is to include the possibility of interference with additional interactions which appear, at the low energies associated with nuclear beta decay, as currents with a different character under Lorentz transformations. As stated in Sec.2.3, we confine ourselves to Scalar currents, which can play a role for any neutrinos, and Right-chiral Vector currents which affect only Dirac and pseudo-Dirac neutrinos.
Such currents will produce direct effects on the rate, proportional to $`\widehat{\rho }_X^2`$ as well as interference terms proportional to $`\widehat{\rho }_X`$. The former will be hard to observe and, for Tritium beta decay, will be absorbed into normalization factors, but are included here for completeness.
In subsections 2.1 and 2.2 and in Eq.(10) above, we used the notation $`\mathrm{cos}\widehat{\theta }_f^i`$ to refer to the direction cosines for new current eigenstates coupled to the charged lepton $`f`$ relative to the mass eigenstate labelled by $`i`$. In this case, the only flavor eigenstate considered is $`e`$ and we want to consider different possible currents, so we change the notation to read $`\mathrm{cos}\theta _{iX}`$ where $`i`$ refers to the mass eigenstate and $`X`$ refers to the current $`R`$, $`S_R`$ or $`S_L`$ (and is omitted for the SM current). We further define $`\rho _X=\widehat{\rho }_X(ME)_X`$ where $`(ME)_X`$ accounts for the ratio of the hadronic matrix element of the designated current relative to that of the SM in the particular nuclear system, here $`A=3`$. This includes the (quark) CKM matrix elements appropriate to that non-SM interaction.
At the risk of being overly pedantic, we present the discussion for both of the cases under consideration first separately and then combine them. For a Right-chiral Vector current (without L-R mixing), we make the substitution
$$\mathrm{cos}^2\theta _iE_\beta E_\nu \mathrm{cos}^2\theta _iE_\beta E_\nu +\mathrm{cos}^2\theta _{iR}\rho _R^2E_\beta E_\nu 2\mathrm{cos}\theta _i\mathrm{cos}\theta _{iR}\rho _Rm_em_i.$$
Defining
$$ϵ_{iR}=\rho _R\frac{\mathrm{cos}\theta _{iR}}{\mathrm{cos}\theta _i}$$
we may recombine this as
$$\mathrm{cos}^2\theta _iE_\beta E_\nu (12ϵ_{iR}\frac{m_e}{E_\beta }\frac{m_i}{E_\nu }+ϵ_{iR}^2).$$
We remind the reader that the sign of the middle term arises from the fact that the neutral lepton in Tritium beta decay, in terms of Dirac fields, is an anti-neutrino. Since $`ϵ_{iR}`$ contains the ratio of direction cosines, it also may have either sign. This allows for positive interference near the end point. Similar remarks apply throughout this section.
Turning now to the other case of interest to us here, of a Right-chiral Scalar interaction (as defined in Secs.2.2,2.3), we make the substitution
$$\mathrm{cos}\theta _i^2E_\beta E_\nu \mathrm{cos}\theta _i^2E_\beta E_\nu +\left(\frac{G_V^2}{G_V^2+3G_A^2}\right)$$
$$\times \left[\mathrm{cos}^2\theta _{iS_R}\rho _{S_R}^2E_\beta E_\nu 2\mathrm{cos}\theta _i\mathrm{cos}\theta _{iS_R}\rho _{S_R}E_\beta m_i\right]$$
And defining
$$ϵ_{iS_R}=\rho _{S_R}\frac{\mathrm{cos}\theta _{iS_R}}{\mathrm{cos}\theta _i}$$
we obtain the expression
$$\mathrm{cos}^2\theta _iE_\beta E_\nu \left[1+\left(\frac{G_V^2}{G_V^2+3G_A^2}\right)\left(2ϵ_{iS_R}\frac{m_i}{E_\nu }+ϵ_{iS_R}^2\right)\right].$$
Finally, combining both of these possibilities (and keeping in mind that any tensor interaction contributions are encompassed by scalar terms), we obtain
$$\left(\frac{dN}{dE_\beta }\right)=\underset{i}{}P_i\left(\frac{dN}{dE_\beta }(_0^i)\right)$$
(12)
with
$`\left({\displaystyle \frac{dN}{dE_\beta }}(_0^i)\right)`$ $`=`$ $`KF(E_\beta )q_\beta (_0^iE_\beta ){\displaystyle \underset{k}{}}\mathrm{\Theta }(_0^iE_\beta m_k)`$ (13)
$`\times E_\beta (_0^iE_\beta )\sqrt{1{\displaystyle \frac{m_k^2}{(_0^iE_\beta )^2}}}`$
$`\times ([\mathrm{cos}^2\theta _k+\mathrm{cos}^2\theta _{kR}\rho _R^2+\left({\displaystyle \frac{G_V^2}{G_V^2+3G_A^2}}\right)\mathrm{cos}^2\theta _{kS_R}\rho _{S_R}^2]`$
$`2m_em_k[\mathrm{cos}\theta _k\mathrm{cos}\theta _{kR}\rho _R]`$
$`2m_kE_\beta \left({\displaystyle \frac{G_V^2}{G_V^2+3G_A^2}}\right)[\mathrm{cos}\theta _k\mathrm{cos}\theta _{kS_R}\rho _{S_R}]).`$
While this expression has everything in it, it is not particularly useful for fitting experimental data. We can cast it into a more useful form by noticing that, near the end point, the dependence on $`E_\beta `$ is very gentle and one is really interested in the dependence on $`E_\nu =(_0^iE_\beta )`$. Furthermore, the product $`F(E_\beta )q_\beta `$ is nearly constant for $`\beta `$ energies near the end point. Note also that the ratio $`\frac{m_e}{E_\beta }`$ varies by less than $`2`$ parts in $`10^3`$ as the kinetic energy of the $`\beta `$ varies from $`17.5keV`$ to $`18.5keV`$, a much wider range than is used in modern fits . These observations suggest the following approximations.
First, absorb $`F(E_\beta )q_\beta E_\beta `$ into a new “constant”
$$K^{}=KF(E_\beta )q_\beta E_\beta .$$
Second, for some average value of $`E_\beta `$, take $`\frac{m_e}{E_\beta }=\frac{m_e}{<E_\beta >}`$ as constant<sup>2</sup><sup>2</sup>2Since it is buried in a parameter to be fit, the actual value is not important. and define
$$\epsilon _k=ϵ_{kR}^2+ϵ_{kS_R}^2\left(\frac{G_V^2}{G_V^2+3G_A^2}\right)$$
and
$$\varphi _k=2\frac{m_k}{(1+\epsilon _k)}\left[\frac{m_e}{<E_\beta >}ϵ_{kR}+ϵ_{kS_R}\left(\frac{G_V^2}{G_V^2+3G_A^2}\right)\right].$$
Then
$`{\displaystyle \frac{dN}{dE_\beta }}(_0^i)`$ $``$ $`K^{}{\displaystyle \underset{k}{}}\mathrm{cos}^2\theta _k(_0^iE_\beta )^2(1+\epsilon _k)\left[1+{\displaystyle \frac{\varphi _k}{(_0^iE_\beta )}}\right]`$ (14)
$`\times \sqrt{1{\displaystyle \frac{m_k^2}{(_0^iE_\beta )^2}}}\mathrm{\Theta }(_0^iE_\beta m_k)`$
### 4.2 Integral spectra
While many of the experiments performed over the last few decades are differential measurements , the two ongoing experiments are integral measurements which accept all $`\beta `$s with energy above some cutoff energy . To obtain the theoretical form for the expected count rate for a given set of parameters, we should integrate Eq.(13) over $`E_\beta `$ from $`E_\beta ^C`$ to $`\mathrm{}`$. This daunting prospect requires numerical treatments. However, the form given in Eq.(14) is both a very good approximation and amenable to analytic integration. The integral needs to be evaluated for each endpoint energy $`_0^i`$ and for each mass eigenvalue $`m_k`$. Let us explicate the procedure for $`_0^0`$ and one mass, $`m_k`$.
We need to evaluate
$$K^{}\mathrm{cos}^2\theta _k_{E^C}^{\mathrm{}}𝑑E_\beta (_0^0E_\beta )^2\sqrt{1\frac{m_k^2}{(_0^0E_\beta )^2}}\mathrm{\Theta }(_0^0E_\beta m_k)$$
$$\times (1+\epsilon _k)\left[1+\frac{\varphi _k}{(_0^0E_\beta )}\right].$$
Changing variables to $`E_\nu =(_0^0E_\beta )`$ and defining $`E_\nu ^C=(_0^0E_\beta ^C)`$, we get
$`N(E_\nu ^C)`$ $`=`$ $`K^{}\mathrm{cos}^2\theta _k{\displaystyle _{m_k}^{E_\nu ^C}}𝑑E_\nu \sqrt{1{\displaystyle \frac{m_k^2}{E_\nu ^2}}}E_\nu ^2(1+\epsilon _k)\left(1+{\displaystyle \frac{\varphi _k}{E_\nu }}\right)`$ (15)
$`=`$ $`K^{}\mathrm{cos}^2\theta _k(1+\epsilon _k)[{\displaystyle \frac{1}{3}}((E_\nu ^C)^2m_k^2)^{3/2}`$
$`+{\displaystyle \frac{\varphi _k}{2}}(E_\nu ^C\sqrt{(E_\nu ^C)^2m_k^2}m_k^2\mathrm{ln}({\displaystyle \frac{E_\nu ^C+\sqrt{(E_\nu ^C)^2m_k^2}}{m_k}}))].`$
### 4.3 Fitting fitted differential spectra
We emphasize, in this subsection and the next, that we attempt neither to fit data nor to extract reliable values of parameters, as that must be done with full knowledge of all experimental details. What we shall do is treat the published values of $`<m^2>_{fit}`$ as a representation of the data. We then ask what parameters will produce similar values of $`<m^2>_{fit}`$ if a fit were to be made, using only the SM expression, to our theoretical spectrum which includes interference effects. The purpose of doing that, in this paper, is to demonstrate that the interference terms apparently do allow a representation of the experimental data.
Furthermore, to simplify both the discussion and our task, we assume that the various experimental groups have correctly performed the weighted sums over the final states of the molecular ionic system. Therefore, we need only study the effects of interference compared with the SM analysis on one branch, which we shall take to be the ground state branch. In fact, the entire analysis of these next two subsections goes through unchanged for other branches, but the equations become needlessly cumbersome for simply the demonstration of our point.
Let us now further assume for simplicity that only one mass eigenstate is important for the fit near the end point. Let that mass eigenvalue be denoted as $`m_1`$.
Since, for Tritium beta decay, the range of $`E_\beta `$ is
$$m_eE_\beta 1.035m_e$$
one will not be able to distinguish between interference terms involving $`E_\beta `$ and $`m_e`$. On the other hand, since we are interested in effects where $`E_\nu `$ varies significantly compared to $`m_\nu `$, those terms will have very different effects on the spectrum.
Now scale $`E_\nu `$ and the $`\varphi _i`$ in units of $`m_1`$,
$`x`$ $`=`$ $`E_\nu /m_1`$
$`f_i`$ $`=`$ $`\varphi _i/m_1`$ (16)
In this case, defining $`Y_1=\mathrm{cos}^2\theta _1(1+\epsilon _1)`$, the differential spectrum becomes
$$\frac{dN}{dE_\beta }=K^{}Y_1x^2(1+\frac{f_1}{x})\sqrt{1\frac{1}{x^2}}\mathrm{\Theta }(x1).$$
(17)
If we were to fit this with a formula for the spectrum which is derived under the assumption that there is only one neutrino and that it has a mass extracted from the spectrum as $`<m^2>_{fit}`$, we would fit Eq.(17) with the function
$$\frac{dN}{dE_\beta }=K^{}Y_Tx^2\sqrt{1\frac{<r^2>_{fit}}{x^2}}$$
(18)
where $`<m^2>_{fit}=m_1^2<r^2>_{fit}`$ is the extracted (apparent) value of the neutrino mass-squared, and $`Y_T=_k\mathrm{cos}^2\theta _k(1+\epsilon _k)`$. The reason that $`Y_T`$ appears is that we assume that the experimental normalization of the data is taken from a region of lower electron energies where neutrino mass effects are presumed to be negligible and all mass eigenstates (and currents) contribute.
Setting Eqs.(17) and (18) equal, at some particular value of $`x`$, gives an equation for $`<r^2>_{fit}`$,
$`<r^2>_{fit}`$ $`=`$ $`x^2[1\zeta ^4]2xf_1\zeta ^4`$ (19)
$`+[1f_1^2]\zeta ^4+2x^1f_1\zeta ^4+x^2f_1^2\zeta ^4`$
where we have defined $`\zeta ^2=Y_1/Y_T`$ for convenience. At $`x=1`$, precisely at the end point of the physical spectrum, Eq.(19) gives the result $`<r^2>_{fit}=1`$.
It is instructive to consider the case in which $`\zeta ^2=1`$ and the first term in Eq.(19) vanishes. (This can occur, for example, if $`\overline{\nu _e}`$ is a mass eigenstate so that $`\mathrm{cos}\theta _1=1`$ and $`\mathrm{cos}\theta _2=\mathrm{cos}\theta _3=0`$.) In this case, no matter how small the interference is, the fit value of $`m^2`$ will eventually become negative far enough away from the end point, assuming all other quantities were perfectly assigned. In fact, this could mean that a value of the mass much less than an $`eV`$ could affect the fit to the spectrum several $`eV`$ below the end point. As we remarked above when discussing pseudo-Dirac neutrinos, a very small splitting, as suggested by the solar neutrino problem, would not appear in this discussion. Other possible flavor mixings would destroy this special condition. Existing reports from accelerator measurements and reactor experiments limit the size of such mixing, but allow it to be non-zero.
If $`\overline{\nu }_e`$ is nearly a mass eigenstate, then the second term in Eq.(19) can still dominate, leading to $`<r^2>_{fit}<0`$ for some values of $`x`$. The value of $`x`$ at which $`<r^2>_{fit}`$ again becomes positive is a sensitive function of $`f_1`$ and $`\zeta ^4`$. For example, taking $`\zeta =0.9996`$, which corresponds to a minimum value of $`\mathrm{sin}^2(2\theta _1)0.003`$ (obtained by setting $`\zeta =\mathrm{cos}\theta _1`$, i.e., ignoring possible corrections due to the small, but generally non-negligible values of the $`\epsilon _k`$) in a two flavor mixing scenario, and $`f_1=0.075`$, which is within the limit on $`X_{S_R}`$ , $`<r^2>_{fit}`$ reaches $`2.5`$ at $`x=45`$ and turns positive a bit beyond $`x=85`$. In such a case, the observed structure at a given point in the spectrum reflects a mass on a much smaller scale.
On the other hand, if the quantity $`\zeta ^4`$ is small compared to $`1`$, $`<r^2>_{fit}`$ will grow as $`x`$ increases unless $`f_1`$ is so large that $`(2f_1+f_1^2)\zeta ^4>(1\zeta ^4)`$, so that the derivative with respect to $`x`$ is negative at $`x=1`$. (In fact, for $`f_1=1`$, analysis shows that $`\zeta ^4`$ must be greater than $`(16/27)`$ to get a negative value of $`<r^2>_{fit}`$ for any value of $`x`$.) For this case to be interesting, a second mass eigenstate must be involved with the interference term being destructive, so that, far from the end point, the interference cancels. Note that, for the cases at hand, the interference is proportional to the mass so that the more massive eigenstate can have a smaller admixture.
Since the modern integral experiments provide the best data available, we defer numerical examples to the next subsection.
### 4.4 Fitting fitted integral spectra
To obtain sufficient statistics, experiments that measure differential spectra make a global fit to data over some range from the endpoint up to some value of $`E_\nu `$, which translates, in practice, to fitting over a range of $`E_\beta `$ down to some cut-off value. This is done automatically in those experiments that measure an integral spectrum.
For fitting to sums of differential spectra, the fitting procedure, viewed in terms of theoretical constructs only and ignoring essential experimental details, corresponds to finding the value of $`<r^2>_{fit}`$ that minimizes the integral
$$I=_1^{x_c}𝑑xx^4\left[Y_1(1+\frac{f_1}{x})\sqrt{1\frac{1}{x^2}}Y_T\sqrt{1\frac{<r^2>_{fit}}{x^2}}\right]^2$$
(20)
where, as in the previous subsection, we have assumed, for simplicity of presentation, that one mass eigenstate with mass $`m_1`$ is important, $`x`$,$`<r^2>_{fit}`$, and $`f_1`$ are as defined previously and $`x_c=E_\nu ^C/m_1`$.
This was the form presented in Ref.() and was appropriate for the earlier differential experiments . However, for the analysis of the integral experiments, one wants to use Eq.(15). The appropriately scaled version is
$$N(x_c)=K^{}Y_1\left(\frac{1}{3}(x_{c}^{}{}_{}{}^{2}1)^{3/2}+\frac{f_1}{2}\left[x_c\sqrt{x_{c}^{}{}_{}{}^{2}1}ln(x_c+\sqrt{x_{c}^{}{}_{}{}^{2}1})\right]\right).$$
(21)
Fitting this with the usual SM formula, where $`K^{}Y_T`$ is taken from the fit far from the end point (experimental normalization) and only one mass eigenstate is included, corresponds to equating
$`N(x_c)`$ $`=`$ $`K^{}Y_T{\displaystyle _1^{x_c}}𝑑xx^2\sqrt{1{\displaystyle \frac{<r^2>_{fit}}{x^2}}}`$ (22)
$`=`$ $`K^{}Y_T{\displaystyle \frac{1}{3}}[((x_{c}^{}{}_{}{}^{2}<r^2>_{fit})^{3/2}(1<r^2>_{fit})^{3/2}]`$
to the same expression in Eq.(21). This procedure gives an implicit equation for $`<r^2>_{fit}`$. In fact the lower limit in Eq.(22) makes a very small contribution for most values of $`x_c`$ of interest, so that, to a good approximation, one may drop that term and obtain an explicit equation for $`<r^2>_{fit}`$.
We used that equation to search through parameter space to find sets that gave a reasonable representation of the latest Mainz data .<sup>3</sup><sup>3</sup>3The experimentalists report substantial consistency of their results ; we have simply found it more convenient to make use of the Mainz presentation of their data. We emphasize again that we have not constructed a fit, even to the representation in terms of $`<m^2>_{fit}`$, but rather have selected a set of parameters that give a reasonable representation of the published results. We have done this to demonstrate that the inclusion of possible interference effects of non-SM currents shows no obvious evidence of being at odds with the data.
To that end we offer a comparison in Figs. 3-6. The data points were read from a reproduction of Fig. 3 of a preprint of Ref.() and should not be taken as a precisely accurate representation. In our Fig. 3, open circles correspond to the Q2 data set of their Fig. 3; in our Fig. 4, open squares correspond to their Q3 data set; in our Fig. 5, open diamonds correspond to their Q4 data set and in our Fig. 6, asterisks correspond to their Q5 data set. No attempt was made to generate fits; rather the parameter sets were chosen to give a maximum negative $`<m^2>_{fit}`$ about $`30eV`$ below the endpoint, $`_0^0`$, and to turn $`<m^2>_{fit}`$ positive at $`E_\beta 18350eV`$. In the figures, the various calculated curves are labelled by the value of $`f_1`$. In Table I, we list the full parameter set for each curve, where the parameters are $`f_1`$, $`m_1`$ and $`\zeta ^2`$. The final row, labelled $`\mathrm{sin}^22\theta `$, is the minimum value one would deduce for that quantity in a two component oscillation formula, obtained by taking the value of $`\zeta ^2`$ to equal $`\mathrm{cos}^2\theta _1`$, i.e., treating the $`\epsilon _k`$ as negligibly small.
Table I. Interference parameters for theoretical curves plotted in Figs. 3-6.
| $`f_1`$ | .05 | .07 | .09 | .11 | .13 |
| --- | --- | --- | --- | --- | --- |
| $`m_1(eV)`$ | 1.603 | 2.228 | 2.888 | 3.431 | 4.334 |
| $`\zeta ^2`$ | .99954 | .99910 | .99850 | .99783 | .99675 |
| $`\mathrm{sin}^22\theta `$ | .00185 | .00360 | .00600 | .00870 | .01300 |
Since the authors of chose not to combine these data sets, for reasons they describe, we also do not attempt to combine them. In fact, as we discuss in the next section, there are theoretical reasons to avoid combining runs taken at different times. In spite of that restriction, we believe that Figs. 3-6 make the point that, for values of the parameters which are within reason, the possibility of interference between the SM current and a weaker, non-SM current can give the observed negative values of $`<m^2>_{fit}`$ when the data is analyzed under the assumption that only the SM current is present.
### 4.5 Discussion
As we showed in Sec.(4.4), it is possible to generate values of $`<m^2>_{fit}`$ that resemble the published data for values of $`f_1`$ between $`.05`$ and $`.13`$. The lower end of this range is easily accommodated by existing limits, even without invoking the possibility, discussed in Sec.(2.4) that these limits are only appropriate for a sum over the mass eigenstates. The upper end of the range may be accommodated by invoking this last point or by noting that more than one non-SM current may be affecting the Tritium data in combination.
While we have not carried out an exhaustive analysis of the allowed parameter space, it is worth noting that, by studying Eq.(19), it is possible to get some sense for the allowed parameter ranges. For example, requiring $`<m^2>_{fit}=10eV^2`$ at $`E_\nu =100eV`$ relates $`m_1`$ to $`f_1`$. In particular, if $`f_1<0.11`$, then $`1eV<m_1<10eV`$. The lower limit comes from the obvious effect that the strength is $`m_1f_1`$ and the upper limit from the fact that the kinematic effect of $`q_\nu `$ begins to close the phase space. We do not mean to emphasize $`f_1=0.11`$ particularly; our point is that the present results strongly suggest that fitting to data using our functional form will give $`m_1`$ near this approximate range.
The required mass eigenvalues, on the other hand, would appear to be in conflict with recent limits from double beta decay . However, those limits are on the expectation of the Majorana mass and, as Wolfenstein pointed out , for a pseudo-Dirac neutrino it is only the difference between the mass eigenstates that is limited by bounds from neutrinoless double beta decay experiments. Taken together with the tiny $`\mathrm{\Delta }m^2`$ inferred from solar neutrino studies , the values of $`m_1`$ (in the few $`eV`$ range) that appear in this analysis strongly imply that the electron neutrino is either a Dirac particle or a pseudo-Dirac particle. The fact that $`\theta _1`$ is very small (inferred from the small value of $`1\zeta ^2`$) in all of the preferred parameter sets is consistent with a small amount of flavor mixing, but the fact that it cannot be zero also indicates that some flavor mixing is required.
We note here that, if this scenario obtains, then there do exist additional neutrino degrees of freedom beyond those active in SM interactions. Such degrees of freedom are usually termed “sterile”; however, our analysis of Tritium beta decay clearly implies that these “sterile” neutrinos necessarily participate in new interactions beyond the SM. In Sec (7) we discuss other implications of this observation.
## 5 Environmental Effects
We have previously studied the consequences that may arise if there is a very light scalar particle coupled only to neutrinos.<sup>4</sup><sup>4</sup>4X.-G. He et. al. have shown that such a scenario is possible in extended gauge theories. The primary effect is that, wherever in space the classical scalar field arising from such interactions is non-zero, the effective mass of the neutrino is altered from its vacuum mass (the mass a neutrino would have in a region of space devoid of other neutrinos, but including all field theoretic contributions). It is this effective mass at the site of a Tritium beta decay event which will govern the interference effects discussed above.
A second result reported in Ref.() is that, for a wide range of parameters which are not in conflict with known data, neutrinos will form clouds during an early stage in the development of the Universe. The extent and density of such clouds will depend on the details of the neutrino vacuum masses, the mass of the scalar and the size of the coupling between the scalar and neutrinos. However, there is no known impediment to considering a length scale between a solar radius and the size of a solar system, with the range of the scalar field being somewhat smaller. If the scale of the particular cloud in which the Solar system developed is on the order of the Earth’s orbit, it would be possible that the strength of the scalar field sampled at the Earth would vary with the Earth’s position along its orbit, leading to a time dependence of the effective mass.
Alternatively, the cloud itself could be undergoing various forms of collective motion (rotations of an ellipsoidal shape or vibrations in various modes), which will be reflected in variations of the strength of the scalar field.
The exact form of the variation of the effective mass as observed on Earth would further depend on unknown details. For example, if the electron anti-neutrino is mostly aligned with the heaviest vacuum mass eigenstate, then the stronger the scalar field, the smaller the effective mass. If, however, as we showed in Ref.(), the electron antineutrino is mostly aligned with one of the lighter vacuum mass eigenstates, as the scalar field gets larger, the magnitude of the effective mass increases. What is uniformly true is that, as the scalar field varies, the effective mass varies. Further, if the effective mass is a small fraction of the vacuum mass, small percentage variations of the scalar field can lead to larger percentage variations of the effective neutrino mass.
Another consequence that follows from this conjecture, as discussed in detail in , is that neutrinos could not constitute a hot dark matter component of the Universe during the formation of large scale structures, since they would be “locked up” within massive clouds with non-relativistic net kinetic energies. Thus a number of the reported limits on neutrino mass from cosmology would need to be revisited.
While there is, of course, no proof that such a background scalar field exists in the region of the Earth (or at all), the fact that it cannot be ruled out suggests that there is a value to independently analyzing experiments done at different times. If there is a time dependence in the effective mass matrix, this may imply a time dependence to the resulting direction cosines. In the formulation presented here, that would affect both $`\mathrm{cos}\theta _1`$ and the value of $`f_1`$ through this and its dependence on $`\mathrm{cos}\widehat{\theta }_1`$. Hence all the parameters of the fit may demonstrate a time dependence, and the correlations are very hard to predict a priori. As bizarre as this possibility may seem, it is clearly prudent to await further experimental information before discarding it.
## 6 Neutrino endpoint effects
The interference terms discussed above display a general symmetry between the $`\beta `$ and the neutrino. M. Goldhaber has raised the question of possible observable effects at the opposite end of the spectrum, where the neutrino carries away all of the available energy. In fact, the only terms that would be observable are complimentary to those discussed above, since any terms proportional to the neutrino mass eigenvalues will be completely dominated by $`E_\nu `$. This leaves the term proportional to $`\rho _{S_R}\frac{m_e}{E_\beta }`$. Since the scale of variation with energy is set by $`m_e`$, we would expect the variation to occur over several $`MeV`$. Consequently, Tritium is not the place to look for such effects.
As the Fermi function varies rapidly at small $`E_\beta `$, the analysis must make use of the full form given above. Furthermore, given the enhancement of the very low energy electron spectrum, compared with the suppression for a positron spectrum, electron emitters may be preferred. This particular current is the most severely constrained, and this part of the spectrum corresponds to large $`E_\nu `$ so that threshold effects are not relevant. Nonetheless, the fact that the scale of the variation is set by $`m_e`$ may make it possible to see such an effect in neutrino spectra of future experiments.
Although the observation of interference effects in such a transition would not directly impact the interpretation of Tritium beta decay, it would demonstrate the existence of non-SM currents and would be a very interesting result in its own right.
## 7 Implications for neutral currents
As we discussed in section (4.5), the combination of our analysis of Tritium beta decay and other experimental results strongly suggests that the electron neutrino is a pseudo-Dirac object with new interaction(s) involving the components which are sterile in the SM. If these new interactions devolve from Higgs mediated scalar currents, as in supersymmetric extensions of the SM , or from Right-chiral Vector currents, as in Left-Right symmetric models , there will be associated new neutral current interactions for the “sterile” components of the neutrino fields. In fact, such neutral currents occur in most proposed extensions of the SM.
These new interactions could affect the interpretation of any experiment sensitive to neutral currents. For example, in the case of the Sudbury Neutrino Observatory experiment (SNO) , oscillation of solar neutrinos into a “sterile” component could produce a neutral current signal intermediate in strength between that expected for oscillations among active neutrino components only and the reduction of signal observed for the charged current interactions. Absent any special quantum coherence effects, the bounds on the strength of any new charged current interactions suggest that the additional neutral current signal provided by this mechanism is not likely to be sufficiently large to confuse an oscillation into “sterile” components with those among active neutrinos of different flavors. Nonetheless, the SNO experimental group must specify the bounds on “sterile” neutral current strengths assumed in the analysis of their data.
Similarly, if, as one expects, these considerations apply to other flavors as well, experimentalists such as those involved in the SuperKamiokande experiment need to consider the possible strength of neutral current interactions of “sterile” neutrinos. That particular case is exacerbated by the fact that the effect of interest is proportional to forward scattering, hence depends on amplitudes rather than rates.
## 8 Conclusions
In this paper we have examined the possibility that interference between the SM Left-chiral current and a weaker, non-SM current with a different Lorentz character may be the origin of the “anomaly” in the Tritium beta decay spectrum near the end point. On general theoretical grounds, it is expected that such currents must exist, the only uncertainty being their strength. To avoid the unwarranted prejudice that only the SM interaction contributes to the decay, experimentalists should include these interference terms when fitting data.
Given the small variation in the electron energy, $`E_\beta `$, we have presented formulas for differential spectra, Eq.(14), and for integral spectra, Eq.(15), which are good approximations appropriate for fitting to data. At a minimum, these should be employed to obtain reasonable parameter values from which to initiate searches with more complete expressions.
Explicitly, we recommend the use of Eq.(15), with $`k=1`$ and $`\varphi _1=f_1m_1`$, for the analysis of ongoing integral experiments , with all of sums over molecular end points implied by Eqs.(12 and 13) and experimental details included. This analysis should be carried out independently of other experimentally derived constraints as those generally reflect differing parameter combinations. In particular, this experiment is uniquely affected by neutrino mass thresholds.
Using a characterization of the data in terms of $`<m^2>_{fit}`$, we have shown that our interpretation of the negative values reported to date is possible for parameter values which are not, in fact, in conflict with other experimental constraints. For consistency with experimental results on neutrino oscillations and neutrinoless double beta decay, the electron neutrino must be either a Dirac or a pseudo-Dirac object with a small CKM-like mixture to other flavor eigenstates.
Finally, we examined three other related considerations: One is the possibility that neutrino endpoint effects might appear to be time dependent due to environmental factors. The second is the complementary effect of new interactions, as discussed here, on the high neutrino energy (low electron energy) end of beta spectra . Lastly, and perhaps most dramatically, we noted that our analysis, when combined with other results, strongly implies the existence of new interactions involving “sterile” neutrinos. This last, in turn, has important implications for the interpretation of experiments studying neutral current interactions.
## 9 Acknowledgments
We are happy to acknowledge valuable conversations on these topics with Maurice Goldhaber, Bill Louis and Peter Herczeg. We thank Christian Weinheimer for careful reading of and comments on drafts of this paper. This research is partially supported by the Department of Energy under contract W-7405-ENG-36, by the National Science Foundation and by the Australian Research Council. One of us (GJS) acknowledges the hospitality of the Institute for Nuclear Theory at the University of Washington on several occasions, where portions of this work were carried out.
## 10 Figure captions
Figure 1. Current eigenstate neutrinos displayed in the space spanned by the three mass eigenstates. The vectors denoted by $`\nu _f`$ are current eigenstates for the SM $`SU(2)_L`$ current, those denoted by $`\widehat{\nu }_f`$ are current eigenstates for whatever other charged current one is considering.
Figure 2. Diagrams contributing to nuclear beta decay under various scenarios. a) SM interaction. Hadronic renormalization produces both L and R hadronic currents in the effective Hamiltonian. b) Scalar exchange. Recoil effects suppress the hadronic Pseudo-scalar coupling. In principle the two leptonic couplings, $`S_L`$ and $`S_R`$ can be different. c) Direct Right-chiral Vector couplings. d) Possible mixing between the $`X_R^{}`$ and $`W^{}`$. The analogous diagram in which the $`W^{}`$ couples to the lepton current gives no observable change in the Tritium spectrum. Again, hadronic renormalization will lead to both L and R effective hadronic currents in both diagrams c) and d).
Figure 3. Neutrino mass-squared extracted from Mainz data set Q2 vs. integral cutoff on electron energy and corresponding results from the SM model analysis of a spectrum including interference effects, for various parameter values. (Details in text.)
Figure 4. The same as Fig.3 for Mainz data set Q3.
Figure 5. The same as Fig.3 for Mainz data set Q4.
Figure 6. The same as Fig.3 for Mainz data set Q5.
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# Transverse momentum dependence of directed particle flow at 160A GeV
## I Introduction
The collective flow of hadrons in ultrarelativistic heavy-ion collisions is a very useful signal to probe the evolution of hot and dense nuclear matter from the onset of its formation . Since the development of flow is closely related to the equation of state (EOS) of nuclear matter, the investigation of the flow can shed light on the transition to a new phase of matter, the so-called quark-gluon plasma (QGP), and its subsequent hadronization . If the transition from the QGP to hadronic phase is of first order, the vanishing of the pressure gradients in the mixed phase leads to the so-called softening of the EOS . The latter should be distinctly seen in the behaviour of the excitation function of the collective flow.
At present, the Fourier expansion technique is widely employed to study collective flow phenomena . Namely, the invariant distribution $`E{\displaystyle \frac{d^3N}{d^3p}}`$ is presented as
$$E\frac{d^3N}{d^3p}=\frac{1}{\pi }\frac{d^2N}{dp_t^2dy}\left[1+2\underset{n=1}{\overset{\mathrm{}}{}}v_n\mathrm{cos}(n\varphi )\right],$$
(1)
where $`p_t`$ and $`y`$ are the transverse momentum and rapidity, and $`\varphi `$ is the azimuthal angle between the momentum of the particle and the reaction plane. The first two Fourier coefficients in Eq.(1), $`v_1`$ and $`v_2`$, are dubbed directed and elliptic flow, respectively. Since both types of anisotropic flow depend on rapidity and transverse momentum, one is able to study double differential distributions
$$v_n(p_t,\mathrm{\Delta }y)=\underset{y_1}{\overset{y_2}{}}\mathrm{cos}(n\varphi )\frac{d^2N}{dp_t^2dy}𝑑y/\underset{y_1}{\overset{y_2}{}}\frac{d^2N}{dp_t^2dy}𝑑y$$
(2)
and
$$v_n(y,\mathrm{\Delta }p_t)=\underset{p_t^{(1)}}{\overset{p_t^{(2)}}{}}\mathrm{cos}(n\varphi )\frac{d^2N}{dp_t^2dy}𝑑p_t^2/\underset{p_t^{(1)}}{\overset{p_t^{(2)}}{}}\frac{d^2N}{dp_t^2dy}𝑑p_t^2.$$
(3)
Model calculations show that elliptic flow is built up at the early phase of nuclear collisions , whereas directed flow develops until the late stage of the reaction . But it is well known that the particles with high transverse momentum are emitted at the onset of the collective expansion, i.e., their directed flow can carry information about the EOS of the dense nuclear phase. Our first goal is to check the emission times of high-$`p_t`$ hadrons and to study the transverse momentum dependence of directed flow in heavy-ion collisions at SPS (160A GeV) energies.
Apparently, one would expect that the directed flow drops to zero when the transverse momentum decreases. For protons such a behaviour has been observed at SIS energies ($`1`$AGeV) . Experimental results on pion and proton directed flow at both AGS and SPS energies, show a qualitatively different picture: $`v_1(p_t)`$ is positive at high $`p_t`$ and slightly but negative at low transverse momenta, i.e., it approaches zero from the negative side. One of the possible explanations of such a peculiar behaviour has been proposed in within the framework of a thermal model. Here the interplay of the radial expansion of a thermalized source and the directed flow has been discussed. It was shown that the $`v_1(p_t)`$ of protons became negative at small values of the transverse momentum provided the transverse expansion velocity of a thermalized source was $`\beta 0.55c`$ or higher. However, the significance of, e.g., an initial anisotropy of the geometrical configuration in non-central collision for the development of anisotropic flow is not completely understood at energies above 10A GeV. The aim of the present paper is also to elaborate the role of non-dynamic, i.e. geometrical effects, which can cause the preferential emission of particles in the direction opposite to that of the normal flow (the so-called antiflow) at small $`p_t`$. For this purpose the quark-gluon string model (QGSM) is chosen.
The paper is organised as follows. A brief description of the microscopic model is given in Sect. II. Rapidity and transverse momentum dependences of frozen nucleons and pions, calculated in Pb+Pb central collisions at 160A GeV, are also discussed. Section III presents a systematic study of the directed pion and nucleon flow at SPS energies as a function of $`p_t`$ and $`y`$. A comparison with experimental data is performed as far as those are available. The directed flow of nucleons, which is developed alongside of the normal flow, grows with rising transverse momentum, while the directed flow of pions changes its orientation from antiflow at low $`p_t`$’s to normal flow at high transverse momenta. In order to compare to the predictions of the thermal approach as well, the calculations for the proton flow are fitted by an expanding thermalized source. Finally, conclusions are drawn in Sect. IV.
## II Features of particle production and freeze-out in the model
The QGSM, which treats the elementary hadronic interactions on the basis of Gribov-Regge theory, is based on the $`1/N_c`$ (where $`N_c`$ is the number of quark colours or flavours) topological expansion of the amplitude for processes in quantum chromodynamics and string phenomenology of particle production in inelastic binary collisions of hadrons. The model incorporates the production of particles via string excitation and subsequent fragmentation, as well as the formation of resonances and hadron rescattering. As independent degrees of freedom the QGSM includes octet and nonet vector and pseudoscalar mesons, and octet and decuplet baryons, and their antiparticles. The model simplifies the in-medium effects and focuses mainly on the development of an intranuclear cascade. Further details on the QGSM can be found elsewhere .
For the simulations at SPS energies, $`E_{lab}=160`$A GeV, light <sup>32</sup>S+<sup>32</sup>S and heavy <sup>208</sup>Pb+<sup>208</sup>Pb symmetric systems have been chosen. According to the QGSM predictions the mean number of interactions per hadron, $`N_{int}^h`$, equals 2 even for central sulphur-sulphur collisions and 9 for central lead-lead collisions at 160A GeV. Due to a significant increase of $`N_{int}^h`$ with rising mass number, one can study the role of the intranuclear cascade on the formation of transverse collective flow. Although the light S+S system, or rather a part of it, cannot be treated as a thermalized source, the formation of thermally equilibrated matter in Pb+Pb collisions is not ruled out. The system expands until all interactions and decays in the reaction have ceased. This stage corresponds to the conditions of the thermal freeze-out. Note also that the system of final particles in the course of model simulations may be well approximated by a core and a halo structure. The halo contains frozen particles already decoupled from the main system, and the core consists of hadrons intensively interacting, both elastically and inelastically, with each other. In other words, there is no sharp freeze-out picture in the microscopic model like the QGSM or the relativistic quantum molecular dynamics (RQMD) model , in contrast to macroscopic hydrodynamic models (see, e.g., and references therein).
The evolution of the number of frozen particles with time $`t`$ is shown in Fig. 1(a) for pions and nucleons in Pb+Pb central collisions. The contour plots (dashed areas) correspond to the $`dN/dt`$ distribution in different rapidity intervals. One can see that the pionic distributions peak at $`t8`$ fm/$`c`$, and the nucleon distributions reach their maxima later, at $`t15`$ fm/$`c`$. It is interesting that the positions of the maxima on the time scale are not shifted when the rapidity range is enlarged. Many hadrons with high rapidity, especially pions, are emitted from the very beginning of the nuclear collision. These particles usually have rather high transverse momentum as well. To illustrate this idea, the time evolution of the transverse mass distributions of nucleons and pions at the freeze-out is shown in Fig. 1(b). We see that nucleons with maximal transverse momenta in lead-lead collisions are coming either from the very beginning of the reaction or from intermediate times with a maximum at $`t13`$ fm/$`c`$. In contrast to nucleons pions with highest $`p_t`$ are produced in inelastic primary NN collisions in heavy-ion reactions, while soft particles are emitted during the whole evolution time. This is a general trend in the production of soft and hard particles in relativistic heavy-ion collisions. Since the excitation function of the transverse particle flow is very sensitive to pressure gradients in the system, the directed flow of both pions and nucleons might change its behaviour with increasing $`p_t`$. Obviously, pions coming from primary NN collisions cannot carry information about properties of hot and dense nuclear matter, nor about the relaxation process. These particles are just produced in the surface regions of the touching nuclei. The admixture of such pions will severely distort the spectrum of pions stemming from the (nearly) thermalized source. In contrast, the fraction of high-$`p_t`$ nucleons emitted promptly after the primary collisions is relatively small compared to total number of high-$`p_t`$ nucleons. Therefore, the spectrum of nucleons with large transverse momentum in heavy-ion collisions at SPS energies might be even more useful to study the early stage of the fireball evolution than that of pions.
## III Directed flow of hadrons
### A Comparison with experimental data
First, directed flow and elliptic flow of pions and protons, calculated for Pb+Pb minimum bias events with the maximum impact parameter $`b_{max}=11`$ fm at SPS energies, are compared in Fig. 2 with the experimental data of the NA49 Collaboration . Both for pions and for protons the agreement between the microscopic calculations and the experimental data, as well as with the relativistic quantum molecular dynamics (RQMD) model results (see Fig. 1 of ), is quite reasonable. The directed flow of protons has a positive slope in the midrapidity range, which becomes steeper as the rapidity window is shifted towards projectile/target rapidity. Also, the pionic directed flow, which has a characteristic negative slope of $`v_1(y)`$ in the range $`1y5`$, drops to zero and even becomes positive at $`yy_{max}`$. This behaviour can be understood, provided the fast particles are formed on the leading quarks (mesons) and diquarks (baryons) at the early times of the collision. The number of secondary interactions per particle with high rapidity is small, thus, the flow in the fragmentation regions is basically determined by the initial geometry of the system . The centrality dependence of anisotropic flow in the QGSM has been studied in . However, as was discussed in Sect. II, the directed flow of soft particles should differ from that of hard particles. Therefore, the $`v_1(y,\mathrm{\Delta }p_t)`$ distribution given by Eq. (3) is used to study the transverse momentum dependence of directed flow.
### B $`v_1(y)`$ in $`p_t`$ intervals
Figure 3 depicts the directed flow of nucleons and pions in two $`p_t`$ intervals, $`0.3<p_t<0.6`$ GeV/$`c`$ and $`0.6<p_t<0.9`$ GeV/$`c`$, for Pb+Pb collisions with different centrality. The maximum impact parameter for a symmetric system is $`b_{max}=2R_A`$. The value of the reduced impact parameter $`\stackrel{~}{b}=b/b_{max}`$ in the simulations varies from 0.15 (central collisions) up to 0.9 (most peripheral collisions). At $`p_t<0.6`$ GeV/$`c`$ the pionic flow exhibits the typical antiflow in both, semicentral and peripheral, collisions. In the same $`p_t`$ interval the nucleon flow increases as the reaction becomes more peripheral. But at $`b8`$ fm the flow becomes softer in the midrapidity range. In very peripheral collisions the directed flow of nucleons shows an antiflow behaviour which is similar to that of the pionic directed flow. As was shown in , see also , such a transformation of the nucleon flow is explained merely by shadowing: It is well known that the presence of even a small amount of quark-gluon plasma leads to a softening of the equation of state, which results in a significant reduction of the directed flow. However, since the QGP is expected to be produced primarily in central heavy-ion collisions, the effect should be most pronounced in central collisions. In contrast, shadowing causes the disappearance of nucleon directed flow and the development of antiflow in the midrapidity region especially in semiperipheral and peripheral collisions, as well as in light systems.
The behaviour of directed flow changes drastically in the transverse momentum range $`0.6<p_t<0.9`$ GeV/$`c`$, presented in Fig. 3(b). Although the nucleon directed flow decreases in the midrapidity range at $`b10`$ fm, its normal component still dominates over the antiflow counterpart. Moreover, even high-$`p_t`$ pions prefer the direction of normal flow, distinctly seen in semiperipheral events with $`4b6`$ fm. This reflects again that hadrons with high transverse momenta ($`p_t0.6`$ GeV/$`c`$) are produced mainly at the early stage of nuclear collisions, as was discussed in Sect. II, see Fig. 1(b).
The transition of the directed flow of pions from antiflow to normal flow with rising transverse momentum has recently been observed in Au+Au collisions at 1 AGeV . The effect is stronger in peripheral collisions and at target rapidities. Likely, this can be explained by the effective shadowing caused by the spectator matter, in accord with quantum molecular dynamics (QMD) transport calculations. A similar feature is seen at SPS at projectile rapidities (see below). However, there is a difference in correlations between early freeze-out times and high transverse momentum of hadrons in heavy-ion collisions at 1A GeV and 160A GeV: At 1A GeV high-$`p_t`$ pions are emitted within the first 15-20 fm/$`c`$ of the reaction and can be used as a “time clock” for the reaction which probes the high density phase . However, at SIS energies high-$`p_t`$ pions have experienced in average more than two reactions by the formation and subsequent decay of $`\mathrm{\Delta }`$-resonances , which is enough for at least partial thermalization of their spectrum. In contrast to this, high-$`p_t`$ pions as a probe of the hot and dense phase in Pb+Pb collisions at 160A GeV should be handled with care because of the lack of rescattering for extremely energetic pions in the latter case.
The role of rescattering in the formation of directed flow is illustrated in Fig. 4. Here the directed flow as a function of rapidity in several $`p_t`$ intervals is compared in minimum bias S+S and Pb+Pb events at 160A GeV. Lacking a sufficiently large amount of secondary interactions per hadron, the directed flow of both nucleons and pions in the light S+S system varies very weakly with rising transverse momentum from $`p_t0.3`$ GeV/$`c`$ to $`0.6<p_t<0.9`$ GeV/$`c`$. Nevertheless, the change of the slope of pionic directed flow from antiflow to normal flow with rising $`p_t`$ is distinctly seen for both reactions. In the heavy ion system the directed flow of hadrons, especially nucleons, strongly depends on the $`p_t`$ range. Note that, since the bulk amount of particles is produced with transverse momenta less than 300 MeV/$`c`$, the distribution, given by Eq. (3) integrated over the whole $`p_t`$ interval, is very close to that for $`0<p_t<0.3`$ GeV/$`c`$.
### C $`v_1(p_t)`$ in rapidity intervals
Figure 5 depicts the directed flow of pions and nucleons as a function of their transverse momentum in different rapidity windows. Both, in S+S and Pb+Pb collisions the directed flow of nucleons weakly depends on $`p_t`$ in the midrapidity range, $`3<y<4`$. It starts rising with increasing $`p_t`$ near the projectile/target rapidity. The directed flow of pions seems to increase, especially in lead-lead collisions, with rising transverse momentum. At $`p_t0.3`$ GeV/$`c`$ the flow of pions in both reactions is small but negative in all three rapidity intervals. A similar picture is observed for low-$`p_t`$ nucleons with the rapidity $`y5`$ in S+S collisions. As already mentioned in the introduction, one possible explanation for the occurrence of antiflow at low $`p_t`$ is the collective motion of a group of particles, which are in (local) thermal equilibrium. Using the expression for the directed flow of non-relativistic particles emitted isotropically by a transversely expanding thermal source :
$`v_1(p_t)`$ $`=`$ $`{\displaystyle \frac{p_t\beta _a}{2T}}\left[1{\displaystyle \frac{m\beta _0}{p_t}}{\displaystyle \frac{I_1(\xi )}{I_0(\xi )}}\right],`$ (4)
$`\xi `$ $`=`$ $`{\displaystyle \frac{\beta _0p_t}{T}},`$ (5)
where $`T`$ is the temperature, $`\beta _a`$ is the collective velocity along the directed flow axis, $`\beta _0`$ is the transverse expansion velocity of the source, and $`I`$ is the modified Bessel function, one can also reproduce the negative values of $`v_1(p_t)`$ in the low $`p_t`$ region. Note, that the validity of Eq. (4) for protons is restricted to the range of $`p_t0.5`$ GeV/$`c`$. A relativistic generalisation of Eq. (4) to an expansion in three dimensions leads to an increase of about 15% in the expansion velocity needed to describe the low $`p_t`$ dip in the $`v_1(p_t)`$ distribution. Results for a fit of Eq. (4) to the low-$`p_t`$ part ($`p_t0.4`$ GeV/$`c`$) of the $`v_1^N(p_t)`$ distribution in the rapidity interval $`4<y<5`$ are plotted in Fig. 5 also. The fitting parameters are $`\beta _a=0.1c`$, $`T=140`$ MeV, and $`\beta _0=0.60(0.40)c`$ for S+S (Pb+Pb) collisions. However, the collective flow in sulphur-sulphur system is weak , and the whole system is far from being in thermal equilibrium. Therefore, the plausible explanation of negative values of $`v_1(p_t)`$ is shadowing.
It is worth noting that the flow of $`\mathrm{\Delta }`$ resonances follows the nucleon flow (see, e.g., ). Pions coming from the decays of $`\mathrm{\Delta }`$’s behave similar to pions emitted from a moving thermal source , i.e., the pionic flow which originates from $`\mathrm{\Delta }`$ resonance decays should be positive at high $`p_T`$ and negative at low transverse momenta. Also, in lead-lead collisions local thermal equilibrium can be reached at least in the central zone of semicentral collisions with impact parameter $`b4`$ fm . In such events directed flow can be affected by a radial isotropic expansion, the formation of long-lived $`\mathrm{\Delta }`$ resonance matter, and nuclear shadowing. In peripheral collisions the possible formation of a thermalized source becomes less important. Here the development of nucleon antiflow in the midrapidity region, as well as in the low $`p_t`$ interval, is completely determined by shadowing.
To compare results of the microscopic calculations with the experimental data, the directed flow of protons and pions in the rapidity range $`4<y<5`$ is shown separately in Fig. 6(a). One sees that the model describes the flow of pions reasonably well, but predicts much stronger signal for protons at $`p_t0.25`$ GeV/$`c`$. In this context it should be noticed that the $`p_t`$ dependence measured by the NA49 Collaboration substantially deviates from the systematics seen at both SIS and AGS energies, where almost a linear increase of the proton directed flow with $`p_t`$ has been observed in Au+Au collisions. In the latter case this behaviour is well reproduced by the RQMD simulations (see Fig. 20 of ). Thus, at SPS energies the QGSM predictions are in line with the SIS/AGS flow systematics, whereas the NA49 Collaboration observes almost a zero proton flow signal as a function of transverse momentum in the interval $`0.1p_t0.9`$ GeV/$`c`$. The present minimum bias calculations are directly compared with the data in the corresponding rapidity window. If this comparison is not biased by (unknown to us) efficiency cuts (in no special cuts were mentioned), this fact, in principle, can be taken as indication of a new dynamical feature not included into the current version of the QGSM. It is interesting that in the target fragmentation region, where, for instance, the formation of the QGP is quite unlikely, the model predictions for the directed flow $`v_1^p(p_t)`$ are close to the experimental data of the WA98 Collaboration , as shown in Fig. 6(b). But the production of even a small amount of quark-gluon plasma should reduce the strength of the low-$`p_t`$ pion flow too. A few explanations which can lead to revision of the experimental data have been proposed recently . The main argument is as follows: Whereas the orientation of the reaction plane is well defined in the microscopic calculations, its orientation in the experimental event has to be somehow reconstructed. Two-particle correlations in the azimuthal plane are usually applied for this purpose. However, as was discussed in detail in , there are several other sources of azimuthal particle correlations not related to the flow. Experimental data corrected for transverse momentum and Hanbury-Brown-Twiss (HBT) correlations for pions, and for $`p_t`$-correlations and correlations from $`\mathrm{\Delta }`$ decays for protons , are plotted onto the results of microscopic calculations in Fig. 6(a) also. We see that the directed flow of protons is almost not changed due to substantial mutual cancellations of $`p_t`$ correlations and $`\mathrm{\Delta }`$-decay correlations, working in the opposite direction. The agreement between the model calculations and corrected data for the directed flow of low-$`p_t`$ pions becomes even worse. On the other hand, comparison between the VENUS raw data sample and the data filtered through the GEANT model of the NA49 detector seems to reveal that the sources of non-flow correlations are insignificant . This important problem, definitely, needs further investigations. It would be very interesting also to compare model predictions with the coming data on sulphur-sulphur collisions .
## IV Conclusions
In summary, we have for the first time performed a systematic study of the directed flow of nucleons and pions as a function of rapidity in different $`p_t`$ intervals at SPS energies. It is shown that the slope of the directed flow of nucleons with $`p_t0.6`$ GeV/$`c`$ is positive (normal flow) in semicentral and semiperipheral collisions, and negative (antiflow) in very peripheral ones, where $`b/b_{max}0.7`$. At higher transverse momenta the slopes of both, pion and nucleon directed flow become positive. Our findings agree with recent experimental results which report a change of sign for the pion directed flow with increasing $`p_t`$ in Au+Au collisions at SIS energies . There the effect was attributed solely to shadowing by the spectator matter. It was found also that high-$`p_t`$ pions ($`p_t0.4`$ GeV/$`c`$) in heavy-ion reactions at 1A GeV are emitted within the first 13 fm/$`c`$ or even earlier. The microscopic analysis of the particle freeze-out conditions in heavy-ion collisions at SPS energies shows that particles with maximal $`p_t`$ are produced essentially either in primary nucleon-nucleon collisions (mesons and nucleons) or in the early phase of the reaction within the first 17 fm/$`c`$ (nucleons). In the yield of high-$`p_t`$ nucleons the admixture of those emitted within the first two fm/$`c`$’s is small. This means that high-$`p_t`$ nucleons might be very useful to study the nuclear equation of state at high temperatures and densities whereas a large amount of high-$`p_t`$ pions is produced even too early to probe this phase.
The pion and nucleon directed flow in S+S collisions, as well as pion directed flow in Pb+Pb collisions, indicates negative values as particle transverse momenta approach zero. Such behaviour can be explained by the interference between a radially expanded thermalized system and anisotropic flow . However, since the formation of a rapidly expanding thermal source in peripheral heavy-ion collisions and in light S+S collisions is rather unlikely, the effect must be caused by nuclear shadowing. Hadrons emitted at small rapidity in the antiflow direction can propagate freely, while hadrons emitted in the normal flow direction still remain within the expanding subsystem (or core) of interacting particles. To study the interplay between the isotropic radial flow and anisotropic directed flow one has to subtract shadowing from the analysis of experimental data.
The microscopic model calculations are in reasonable agreement with most experimental data on anisotropic flow, both directed and elliptic, in minimum bias Pb+Pb events. The model is able to reproduce quantitatively the strong negative flow of pions, $`v_1^\pi (p_t,4<y<5)`$, at low transverse momenta, $`p_t0.25`$ GeV/$`c`$ which is seen by the NA49 Collaboration. Concerning the $`p_t`$ dependence of the proton flow the model calculations show the same systematics as observed at SIS and AGS energies, namely, a linear rise of the $`v_1`$ with $`p_t`$, whereas the very weak signal at $`p_t0.9`$ GeV/$`c`$ is observed in this rapidity window. If the comparison with data is not biased by unknown efficiency cuts, there is left space for the speculation about dynamical features not incorporated into the microscopic model. To clarify this point and to make more definite conclusions it would be interesting also to compare the model calculations with forthcoming S+S data at the same energy.
###### Acknowledgements.
We are thankful to L. Csernai, P. Danielewicz, D. Röhrich, D. Strottman, and H. Wolter for the interesting discussions and fruitful comments. This work was supported in part by the Bundesministerium für Bildung und Forschung (BMBF) under contract 06TÜ986.
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# Brane-World Cosmology in Higher Derivative Gravity or Warped Compactification in the Next-to-leading Order of AdS/CFT Correspondence
## 1 Introduction
In brane-world scenarios one lives on the boundary (observable Universe) where gravity is trapped . Such brane is embedded in higher dimensional bulk space. The investigation of cosmological aspects of brane-worlds (see also references therein) shows that at some circumstances the inflationary Universe could be realized on the brane. Even the experimental tests to search for higher dimensional deviations of our world due to bulk/boundary structure may be proposed . However, the study of such brane cosmology has been done so far almost exclusively for Einstein or dilatonic gravities. On the same time, higher derivative (HD) gravity represents very natural generalization of general relativity. It enjoys various nice features like renormalizability and asymptotic freedom in four dimensions (see for introduction and review), possibility of self-consistent compactification on quantum and classical levels (see examples in ), sufficiently small corrections to Newton potential at reasonable range of parameters, etc. Moreover, HD terms are typical for string effective action in the derivatives expansion . Hence, HD gravity in higher dimensions represents the interesting model where brane-world cosmology should be investigated and the trapping of gravity should be discussed. Note that propagator in such theory is qualitatively different from the one in general relativity that is why some new phenomena may be expected.
The important remark is in order. In the widely accepted versions of brane-world scenario one studies $`d+1`$-dimensional gravity coupled to a brane in the formalism where normally two free parameters (bulk cosmological constant and brane tension) present. Adding of $`d`$-dimensional Gibbons-Hawking boundary action and brane cosmological constant term to action one can get de Sitter brane in AdS bulk even in Einstein gravity. The position of such brane is fixed in terms of brane tension.
Our ideology is somehow different, in the spirit of refs.. Namely, one considers the addition of surface counterterms which make variational procedure to be well-defined and eliminate the leading divergences of action. Brane cosmological constant (or brane tension) is not considered as free parameter (as it was in original brane-world scenario) but it is fixed by condition of finiteness of spacetime when brane goes to infinity. In such approach, the possibility of cosmological de Sitter brane-world in Einstein theory is eliminated. However, as we explain below in theories different from Einstein gravity this possibility is not ruled out thanks to other free parameters of theory.
Our purpose in the present work will be to study the brane-world cosmology in d5 HD gravity without (or with) quantum corrections. We consider general model with arbitrary coefficients (next section) and derive bulk and boundary equations of motion. The explicit structure of surface counterterms is very important in such derivation.
The solution of bulk equation of motion gives d5 AdS. Then, the brane equation of motion is discussed. It gives the restrictions to HD gravity parameters from the condition of realization of spherical, hyperbolic or flat branes on the boundary. The corresponding radius is derived when it exists. The specific versions of HD gravity like Weyl or Gauss-Bonnet theories appear in this formalizm as particular examples.
It is very important that there may be twofold point of view to HD gravity. From one side, this is just alternative to general relativity. From another side, within AdS/CFT correspondence some versions of HD gravity represent SG duals where Einstein and cosmological terms are of leading order and HD terms are of next-to-leading order in large $`N`$ approximation. The explicit example of that sort is presented in next section (SG dual to $`𝒩=2`$ SCFT). Brane-world cosmology for such theory naturally appears as warped compactification in the next-to-leading order of AdS/CFT correspondence. It is shown that in such situation the creation of spherical brane is impossible in leading as well as in next-to-leading order of large $`N`$ expansion. On the same time, the account of next-to-leading order terms makes possible the creation of hyperbolic brane living in d5 AdS bulk (this effect was prohibited in the leading order of AdS/CFT correspondence).
In section three we investigate the modification of above scenario when quantum CFT is living on the brane. This is done via the anomaly induced effective action. (The corresponding study for bulk Einstein gravity has been done in refs..<sup>2</sup><sup>2</sup>2In usual 4d world the anomaly driven inflation has been proposed in refs.. In fact, this puts world-brane scenario in form of warped compactification in AdS/CFT set-up as the corresponding RG flow.) We show that for range of HD gravity parameters where classical consideration does not give inflationary or hyperbolic branes, the quantum brane matter effects improve the situation: purely quantum creation of inflationary or hyperbolic branes in d5 AdS occurs (like in Einstein gravity). On the same time if the phenomenon already existed on classical level then quantum corrections do not destroy it. The bulk clearly is not modified.
For SG dual of $`𝒩=2`$ SCFT theory such picture may be naturally understood in terms of holographic RG description in AdS/CFT set-up. That is two dual descriptions (SG dual and QFT dual) are matched together in some energy region into the global representation of some RG flow. It is interesting that both sides are given in next-to-leading order of large $`N`$ expansion as conformal anomaly coefficients for SCFT under discussion include not only quadratic but also linear terms on $`N`$. Then, it is shown that in such theory the spherical brane in d5 AdS is created (the role of conformal anomaly coefficients is dominant). Without quantum CFT living on the brane it was impossible. For hyperbolic brane creation the qualitative results are not changed if compare with previous section. Finally, in the last section we present brief summary of results and mention some possible developments along the direction under consideration.
## 2 Brane-World Cosmology in HD gravity
We consider 5d spacetime whose boundary is 4 dimensional sphere S<sub>4</sub>, which can be identified with a D3-brane, four-dimensional hyperboloid H<sub>4</sub>, or four dimensional flat space R<sub>4</sub>. The bulk part is given by 5 dimensional Euclidean Anti-de Sitter space $`\mathrm{AdS}_5`$
$$ds_{\mathrm{AdS}_5}^2=dy^2+\mathrm{sinh}^2\frac{y}{l}d\mathrm{\Omega }_4^2.$$
(1)
Here $`d\mathrm{\Omega }_4^2`$ is given by the metric of S<sub>4</sub>, H<sub>4</sub> or R<sub>4</sub> with unit radius. One also assumes the boundary (brane) lies at $`y=y_0`$ and the bulk space is given by gluing two regions given by $`0y<y_0`$.
One starts with the following action:
$$S=d^5x\sqrt{\widehat{G}}\left\{a\widehat{R}^2+b\widehat{R}_{\mu \nu }\widehat{R}^{\mu \nu }+c\widehat{R}_{\mu \nu \xi \sigma }\widehat{R}^{\mu \nu \xi \sigma }+\frac{1}{\kappa ^2}\widehat{R}\mathrm{\Lambda }\right\}.$$
(2)
Here the conventions of curvatures are given by
$`R`$ $`=`$ $`g^{\mu \nu }R_{\mu \nu }`$
$`R_{\mu \nu }`$ $`=`$ $`\mathrm{\Gamma }_{\mu \lambda ,\kappa }^\lambda +\mathrm{\Gamma }_{\mu \kappa ,\lambda }^\lambda \mathrm{\Gamma }_{\mu \lambda }^\eta \mathrm{\Gamma }_{\kappa \eta }^\lambda +\mathrm{\Gamma }_{\mu \kappa }^\eta \mathrm{\Gamma }_{\lambda \eta }^\lambda `$
$`\mathrm{\Gamma }_{\mu \lambda }^\eta `$ $`=`$ $`{\displaystyle \frac{1}{2}}g^{\eta \nu }\left(g_{\mu \nu ,\lambda }+g_{\lambda \nu ,\mu }g_{\mu \lambda ,\nu }\right).`$ (3)
When $`a=b=c=0`$, the action (2) becomes that of the Einstein gravity:
$$S=_{M_{d+1}}d^{d+1}x\sqrt{\widehat{G}}\left\{\frac{1}{\kappa ^2}\widehat{R}\mathrm{\Lambda }\right\}.$$
(4)
If we choose
$$a=\frac{1}{6}\widehat{c},b=\frac{4}{3}\widehat{c},c=\widehat{c}$$
(5)
the HD part of action is given by the square of the Weyl tensor $`C_{\mu \nu \rho \sigma }`$ :
$$S=d^5x\sqrt{\widehat{G}}\left\{\widehat{c}\widehat{C}_{\mu \nu \xi \sigma }\widehat{C}^{\mu \nu \xi \sigma }+\frac{1}{\kappa ^2}\widehat{R}\mathrm{\Lambda }\right\}.$$
(6)
It is interesting that the string theory dual to $`𝒩=2`$ superconformal field theory is presumably IIB string on $`\mathrm{AdS}_5\times X_5`$ where $`X_5=S^5/Z_2`$. (The $`𝒩=2`$ $`Sp(N)`$ theory arises as the low-energy theory on the world volume on $`N`$ D3-branes sitting inside 8 D7-branes at an O7-brane). Then in the absence of Weyl term, $`\frac{1}{\kappa ^2}`$ and $`\mathrm{\Lambda }`$ are given by
$$\frac{1}{\kappa ^2}=\frac{N^2}{4\pi ^2},\mathrm{\Lambda }=\frac{12N^2}{4\pi ^2}.$$
(7)
This defines the bulk gravitational theory dual to super YM theory with two supersymmetries. The Riemann curvature squared term in the above bulk action may be deduced from heterotic string via heterotic-type I duality , which gives $`𝒪(N)`$ correction:
$$a=b=0,c=\frac{6N}{2416\pi ^2}.$$
(8)
Hence, HD gravity with above coefficients defines SG dual of super Yang-Mills theory (with two supersymmetries) in next-to-leading order of AdS/CFT correspondence .
Using field redefinition ambiguity one can suppose that there exists the scheme where $`R_{\mu \nu \alpha \beta }^2`$ may be modified to $`C_{\mu \nu \alpha \beta }^2`$ in the same way as in ref.. Then, the action (4) is presumably the bulk action (in another scheme) dual to $`𝒩=2`$ SCFT.
Let us start from the bulk equations of motion. First we investigate if the equations of motion for the general action (2) have a solution which describes anti de Sitter space, whose metric is given by
$$ds^2=\widehat{G}_{\mu \nu }^{(0)}dx^\mu dx^\nu =\frac{l^2}{4}\rho ^2d\rho d\rho +\underset{i=1}{\overset{4}{}}\rho ^1\eta _{ij}dx^idx^j.$$
(9)
When we assume the metric in the form (9), the scalar, Ricci and Riemann curvatures are given by
$$\widehat{R}^{(0)}=\frac{20}{l^2},\widehat{R}_{\mu \nu }^{(0)}=\frac{4}{l^2}G_{\mu \nu }^{(0)},\widehat{R}_{\mu \nu \rho \sigma }^{(0)}=\frac{1}{l^2}\left(G_{\mu \rho }^{(0)}G_{\nu \sigma }^{(0)}G_{\mu \sigma }^{(0)}G_{\nu \rho }^{(0)}\right),$$
(10)
which tell that these curvatures are covariantly constant. Then in the equations of motion from the action (2), the terms containing the covariant derivatives of the curvatures vanish and the equations have the following form:
$`0`$ $`=`$ $`{\displaystyle \frac{1}{2}}G_{\zeta \xi }^{(0)}\left\{a\widehat{R}^{(0)2}+b\widehat{R}_{\mu \nu }^{(0)}\widehat{R}^{(0)\mu \nu }+c\widehat{R}_{\mu \nu \rho \sigma }^{(0)}\widehat{R}^{(0)\mu \nu \rho \sigma }+{\displaystyle \frac{1}{\kappa ^2}}\widehat{R}^{(0)}\mathrm{\Lambda }\right\}`$ (11)
$`+2aR^{(0)}R_{\zeta \xi }^{(0)}+2b\widehat{R}_{\mu \zeta }^{(0)}\widehat{R}_{}^{(0)\mu }{}_{\xi }{}^{}+2c\widehat{R}_{\zeta \mu \nu \rho }^{(0)}\widehat{R}_\xi ^{(0)\mu \nu \rho }+{\displaystyle \frac{1}{\kappa ^2}}\widehat{R}_{\zeta \xi }^{(0)}.`$
Then substituting Eqs.(10) into (11), one gets
$`0`$ $`=`$ $`{\displaystyle \frac{80a}{l^4}}+{\displaystyle \frac{16b}{l^4}}`$ (12)
$`+{\displaystyle \frac{8c}{l^4}}{\displaystyle \frac{12}{\kappa ^2l^2}}\mathrm{\Lambda }.`$
The equation (12) can be solved with respect to $`l^2`$ if
$$\frac{144}{\kappa ^4}16\left\{20a+4b+2c\right\}\mathrm{\Lambda }0$$
(13)
which can been found from the determinant in (12). Then we obtain
$$l^2=\frac{\frac{12}{\kappa ^2}\pm \sqrt{\frac{144}{\kappa ^4}16\left\{20a+4b+2c\right\}\mathrm{\Lambda }}}{2\mathrm{\Lambda }}.$$
(14)
The sign in front of the root in the above equation may be chosen to be positive which corresponds to the Einstein gravity ($`a=b=c=0`$). For SG dual of $`𝒩=2`$ $`Sp(N)`$ theory, we find from (8)
$$\frac{1}{l^2}=1+\frac{1}{24N}+𝒪\left(N^2\right).$$
(15)
Now, let us discuss the surface terms in HD gravity on the chosen background
$$ds^2\widehat{G}_{\mu \nu }dx^\mu dx^\nu =\frac{l^2}{4}\rho ^2d\rho d\rho +\underset{i=1}{\overset{d}{}}\widehat{g}_{ij}dx^idx^j,\widehat{g}_{ij}=\rho ^1g_{ij}.$$
(16)
If the boundary of AdS<sub>5</sub> lies at $`\rho =\rho _0`$, the variation $`\delta S`$ contains the derivative of $`\delta \widehat{g}^{ij}`$ with respect to $`\rho `$, which makes the variational principle ill-defined. In order that the variational principle is well-defined on the boundary, the variation of the action should be written in the form of
$$\delta S=d^5x\sqrt{\widehat{G}}\delta \widehat{g}^{ij}\times \left(\text{eq. of motion}\right)+_{\rho =\rho _0}d^4x\sqrt{\widehat{g}}\delta \widehat{g}^{ij}\left\{\mathrm{}\right\}$$
(17)
after using the partial integration. If we put $`\left\{\mathrm{}\right\}=0`$ for $`\left\{\mathrm{}\right\}`$ in (17), we could obtain the boundary condition. If the variation of the action on the boundary contains $`(\delta \widehat{g}^{ij})^{}`$, however, we cannot partially integrate it with respect to $`\rho `$ on the boundary to rewrite the variation in the form of (17) since $`\rho `$ is the coordinate expressing the direction perpendicular to the boundary. Therefore the “minimum” of the action is ambiguous. Such a problem was well studied by Gibbons and Hawking in for the Einstein gravity ($`a=b=c=0`$). The boundary term was added to the action, which cancels the variation :
$$S_b^{\mathrm{GH}}=\frac{2}{\stackrel{~}{\kappa }^2}_{\rho =\rho _0}d^4x\sqrt{\widehat{g}}D_\mu n^\mu .$$
(18)
Here $`n^\mu `$ is the unit vector normal to the boundary. In the coordinate choice (16), the action (18) has the form
$$S_b^{\mathrm{GH}}=\frac{2}{\stackrel{~}{\kappa }^2}_{\rho =\rho _0}d^4x\sqrt{\widehat{g}}\frac{\rho }{l}\widehat{g}_{ij}\left(\widehat{g}_{ij}\right)^{}.$$
(19)
Then the variation over the metric $`\widehat{g}_{ij}`$ gives
$$\delta S_b^{\mathrm{GH}}=\frac{2}{\stackrel{~}{\kappa }}_{\rho =\rho _0}d^4x\sqrt{\widehat{g}}\frac{\rho }{l}\left[\delta \widehat{g}^{ij}\left\{\widehat{g}_{ik}\widehat{g}_{il}\left(\widehat{g}_{kl}\right)^{}\frac{1}{2}\widehat{g}_{ij}\widehat{g}_{kl}\left(\widehat{g}_{kl}\right)^{}\right\}+\widehat{g}_{ij}\left(\delta \widehat{g}_{ij}\right)^{}\right].$$
(20)
From the other side, the surface terms in the variation of the bulk Einstein action ($`a=b=c=0`$ in (4)) have the form
$`\delta S^{\mathrm{Einstein}}`$ $`=`$ $`{\displaystyle d^5x\sqrt{\widehat{G}}\delta \widehat{g}^{ij}\times \left(\text{Einstein equation}\right)}`$ (21)
$`+{\displaystyle \frac{1}{\kappa ^2}}{\displaystyle _{\rho =\rho _0}}d^4x\sqrt{\widehat{g}}{\displaystyle \frac{2\rho }{l}}\left[\widehat{g}_{ij}^{}\delta \widehat{g}^{ij}+\widehat{g}_{ij}\left(\delta \widehat{g}^{ij}\right)^{}\right].`$
Then we find the terms containing $`\left(\delta \widehat{g}^{ij}\right)^{}`$ in (20) and (21) are cancelled with each other.
We also need the counterterms, besides Gibbons-Hawking term (18), to cancell the divergence coming from the infinite volume of AdS. Such a kind of counterterms can be given by the local quantities on the 4 dimensional boundary. In , the surface counterterms are discussed for higher derivative gravities in all detail. Note that they are relevant also for quantum cosmology.
We also add the surface terms $`S_b^{(1)}`$ corresponding to Gibbons-Hawking term (18) and $`S_b^{(2)}`$ which is the leading counterterm corresponding to the vacuum energy on the brane:
$`S_b`$ $`=`$ $`S_b^{(1)}+S_b^{(2)}`$
$`S_b^{(1)}`$ $`=`$ $`{\displaystyle }d^4x\sqrt{\widehat{g}}[4\stackrel{~}{a}\widehat{R}D_\mu n^\mu +2\stackrel{~}{b}(n_\mu n_\nu \widehat{R}^{\mu \nu }D_\sigma n^\sigma +\widehat{R}_{\mu \nu }D^\mu n^\nu )`$
$`+8\stackrel{~}{c}n_\mu n_\nu \widehat{R}^{\mu \tau \nu \sigma }D_\tau n_\sigma {\displaystyle \frac{2}{\stackrel{~}{\kappa }^2}}D_\mu n^\mu ]`$
$`S_b^{(2)}`$ $`=`$ $`\eta {\displaystyle d^4x\sqrt{\widehat{g}}}.`$ (22)
In , in order to cancell the leading order divergence, which appears when the brane goes to infinity, we got
$$\eta =\frac{32T}{l^3}+\frac{8}{l\kappa ^2}+\frac{4\stackrel{~}{T}}{l^3}\frac{2}{l\stackrel{~}{\kappa }^2}.$$
(23)
Here
$$T=10a+2b+c,\stackrel{~}{T}=10\stackrel{~}{a}+2\stackrel{~}{b}+\stackrel{~}{c}.$$
(24)
Note that unlike to standard brane-world scenarios $`\eta `$ is not free parameter.
The metric of $`\mathrm{S}_4`$ with the unit radius is given by
$$d\mathrm{\Omega }_4^2=d\chi ^2+\mathrm{sin}^2\chi d\mathrm{\Omega }_3^2.$$
(25)
Here $`d\mathrm{\Omega }_3^2`$ is described by the metric of 3 dimensional unit sphere. If we change the coordinate $`\chi `$ to $`\sigma `$ by
$$\mathrm{sin}\chi =\pm \frac{1}{\mathrm{cosh}\sigma },$$
(26)
one obtains
$$d\mathrm{\Omega }_4^2=\frac{1}{\mathrm{cosh}^2\sigma }\left(d\sigma ^2+d\mathrm{\Omega }_3^2\right).$$
(27)
On the other hand, the metric of the 4 dimensional flat Euclidean space is given by
$$ds_{4\mathrm{E}}^2=d\zeta ^2+\zeta ^2d\mathrm{\Omega }_3^2.$$
(28)
Then by changing the coordinate as
$$\zeta =\mathrm{e}^\sigma ,$$
(29)
one gets
$$ds_{4\mathrm{E}}^2=\mathrm{e}^{2\sigma }\left(d\sigma ^2+d\mathrm{\Omega }_3^2\right).$$
(30)
For the 4 dimensional hyperboloid with the unit radius, the metric is given by
$$ds_{\mathrm{H4}}^2=d\chi ^2+\mathrm{sinh}^2\chi d\mathrm{\Omega }_3^2.$$
(31)
Changing the coordinate $`\chi `$ to $`\sigma `$
$$\mathrm{sinh}\chi =\frac{1}{\mathrm{sinh}\sigma },$$
(32)
one finds
$$ds_{\mathrm{H4}}^2=\frac{1}{\mathrm{sinh}^2\sigma }\left(d\sigma ^2+d\mathrm{\Omega }_3^2\right).$$
(33)
Motivated by (27), (30) and (33), one takes the metric of 5 dimensional space time as follows:
$$ds^2=dz^2+\mathrm{e}^{2A(z,\sigma )}\underset{i,j=1}{\overset{4}{}}\stackrel{~}{g}_{ij}dx^idx^j,\stackrel{~}{g}_{\mu \nu }dx^\mu dx^\nu l^2\left(d\sigma ^2+d\mathrm{\Omega }_3^2\right).$$
(34)
Here the coordinate $`z`$ is related the coordinate $`\rho `$ in (9) by
$$\rho =\mathrm{e}^{\frac{2z}{l}}.$$
(35)
Under the choice of metric in (34), the curvatures have the following forms:
$`R_{zizj}`$ $`=`$ $`\mathrm{e}^{2A}\left(A_{,zz}\left(A_{,z}\right)^2\right)\stackrel{~}{g}_{ij}`$
$`R_{zAz\sigma }`$ $`=`$ $`l^2\mathrm{e}^{2A}A_{,z\sigma }g_{AB}^s`$
$`R_{\sigma A\sigma B}`$ $`=`$ $`\left(l^2\mathrm{e}^{2A}A_{,\sigma \sigma }l^2\mathrm{e}^{4A}\left(A_{,z}\right)^2\right)g_{AB}^s`$
$`R_{ABCD}`$ $`=`$ $`\left(l^2\mathrm{e}^{2A}l^2\mathrm{e}^{2A}\left(A_{,\sigma }\right)^2l^4\mathrm{e}^{4A}\left(A_{,z}\right)^2\right)\left(g_{AC}^sg_{BD}^sg_{AD}^sg_{BC}^s\right)`$
$`R_{zz}`$ $`=`$ $`4\left(A_{,zz}\left(A_{,z}\right)^2\right)`$
$`R_{z\sigma }`$ $`=`$ $`3A_{,z\sigma }`$
$`R_{\sigma \sigma }`$ $`=`$ $`l^2l^2\mathrm{e}^{2A}\left(A_{,zz}4\left(A_{,z}\right)^2\right)3A_{,\sigma \sigma }`$
$`R_{AB}`$ $`=`$ $`\left(l^2\mathrm{e}^{2A}\left(A_{,zz}4\left(A_{,z}\right)^2\right)A_{,\sigma \sigma }2\left(A_{,\sigma }\right)^2+2\right)g_{AB}^s`$
$`R`$ $`=`$ $`8A_{,zz}20\left(A_{,z}\right)^2+l^2\mathrm{e}^{2A}\left(6A_{,\sigma \sigma }6\left(A_{,sigma}\right)^2+6\right).`$ (36)
Here $`_{,\mu \nu \mathrm{}}\frac{}{x^\mu }\frac{}{x^\nu }\mathrm{}()`$. Other curvatures except those obtained by permutating the indeces of the above curvatures vanish. We also write the metric of S<sub>3</sub> in the following form:
$$d\mathrm{\Omega }_3^2=\underset{A,B=1}{\overset{3}{}}g_{AB}^sdx^Adx^B.$$
(37)
One gets that $`n^\mu `$ and the covariant derivative of $`n^\mu `$ are
$$n^\mu =\delta _\rho ^\mu ,D_in^j=\delta _i^jA_{,z}(\text{others}=0).$$
(38)
Then the actions $`S`$ in (2) and $`S_b`$ in (2) have the following forms:
$`S`$ $`=`$ $`l^4{\displaystyle }d^5x\mathrm{e}^{4A}\sqrt{g^s}[a\{64\left(A_{,zz}\right)^2+320A_{,zz}\left(A_{,z}\right)^2+400\left(A_{,z}\right)^4`$ (39)
$`+36l^4\mathrm{e}^{4A}\left(A_{,\sigma \sigma }\right)^2+72l^4\mathrm{e}^{4A}A_{,\sigma \sigma }\left(A_{,\sigma }\right)^2+36l^4\mathrm{e}^{4A}\left(A_{,\sigma }\right)^4`$
$`+l^2\mathrm{e}^{2A}\left(96A_{,zz}+240\left(A_{,z}\right)^2\right)\left(A_{,\sigma \sigma }+\left(A_{,\sigma }\right)^21\right)`$
$`+l^4\mathrm{e}^{4A}(72A_{,\sigma \sigma }72\left(A_{,\sigma }\right)^2+36)\}`$
$`+b\{20\left(A_{,zz}\right)^2+64A_{,zz}\left(A_{,z}\right)^2+80\left(A_{,z}\right)^4+18l^2\mathrm{e}^{2A}\left(A_{,z\sigma }\right)^2`$
$`+12l^4\mathrm{e}^{4A}\left(A_{,\sigma \sigma }\right)^2+12l^4\mathrm{e}^{4A}A_{,\sigma \sigma }\left(A_{,\sigma }\right)^2+12l^4\mathrm{e}^{4A}\left(A_{,\sigma }\right)^4`$
$`+l^2\mathrm{e}^{2A}\left(12A_{,zz}+48\left(A_{,z}\right)^2\right)\left(A_{,\sigma \sigma }+\left(A_{,\sigma }\right)^21\right)`$
$`+l^4\mathrm{e}^{4A}(12A_{,\sigma \sigma }24\left(A_{,\sigma }\right)^2+12)\}`$
$`+c\{16\left(A_{,zz}\right)^2+32A_{,zz}\left(A_{,z}\right)^2+40\left(A_{,z}\right)^4+24l^2\mathrm{e}^{2A}\left(A_{,z\sigma }\right)^2`$
$`+12l^4\mathrm{e}^{4A}\left(A_{,\sigma \sigma }\right)^2+12l^4\mathrm{e}^{4A}\left(A_{,\sigma }\right)^4`$
$`+24l^2\mathrm{e}^{2A}\left(A_{,z}\right)^2(A_{,\sigma \sigma }+\left(A_{,\sigma }\right)^21)+12l^4\mathrm{e}^{4A}((2A_{,\sigma })^2+1)\}`$
$`+{\displaystyle \frac{1}{\kappa ^2}}\{(8A_{,zz}20\left(A_{,z}\right)^2)`$
$`+(6A_{,\sigma \sigma }6\left(A_{,\sigma }\right)^2+6)\mathrm{e}^{2A}\}+\mathrm{\Lambda }]`$
$`S_b`$ $`=`$ $`l^4{\displaystyle }d^4x\mathrm{e}^{4A}\sqrt{g^s}[16\stackrel{~}{a}\{(8A_{,zz}20\left(A_{,z}\right)^2)`$ (40)
$`+(6A_{,\sigma \sigma }6\left(A_{,\sigma }\right)^2+6)\mathrm{e}^{2A}\}A_{,z}`$
$`+2\stackrel{~}{b}\left\{\left(20A_{,zz}32\left(A_{,z}\right)^2\right)+\left(6A_{,\sigma \sigma }6\left(A_{,\sigma }\right)^2+6\right)\mathrm{e}^{2A}\right\}A_{,z}`$
$`+32\stackrel{~}{c}(A_{,zz}\left(A_{,z}\right)^2)+\eta ]`$
¿From the variation over $`A`$, one obtains the following equation on the brane, which lies at $`z=z_0`$:
$`\delta \left(S+2S_b\right)`$ (41)
$`=`$ $`2V_3l^4{\displaystyle }d\sigma \mathrm{e}^{4A}[(32\stackrel{~}{T}+24\stackrel{~}{U})A_{,z}\delta A_{,zz}`$
$`+\{(32T24U32\stackrel{~}{T}+24\stackrel{~}{U})A_{,zz}+\{(32T24U96\stackrel{~}{T}+72\stackrel{~}{U})\left(A_{,z}\right)^2`$
$`+12(U\stackrel{~}{U})l^2\mathrm{e}^{2A}(A_{,\sigma \sigma }+\left(A_{,\sigma }\right)^21){\displaystyle \frac{8}{\kappa ^2}}+{\displaystyle \frac{8}{\stackrel{~}{\kappa }^2}}\}\delta A_{,z}`$
$`+\{(64T128\stackrel{~}{T}+96\stackrel{~}{U})A_{,zz}A_{,z}`$
$`+\left(160T128\stackrel{~}{T}\right)\left(A_{,z}\right)^3`$
$`\left(48T24\stackrel{~}{U}\right)l^2\mathrm{e}^{2A}\left(A_{,\sigma \sigma }+\left(A_{,\sigma }\right)^21\right)A_{,z}`$
$`+\left(36b96c12\stackrel{~}{U}\right)l^2\mathrm{e}^{2A}A_{,z\sigma \sigma }+\left(72b192c\right)l^2\mathrm{e}^{2A}A_{,\sigma }A_{,z\sigma }`$
$`+({\displaystyle \frac{8}{\kappa ^2}}+{\displaystyle \frac{32}{\stackrel{~}{\kappa }^2}})A_{,z}4\eta \}\delta A].`$
The factors $`2`$ in front of $`S_b`$ and $`V_3`$ come from the fact that we are considering two bulk space (corresponding to $`B_5^{(1,2)}`$ in ) which have one common boundary (S<sub>4</sub> in ) or brane.
Here $`T`$ and $`\stackrel{~}{T}`$ are defined in (24) and
$$U=8a+b,\stackrel{~}{U}=8\stackrel{~}{a}+\stackrel{~}{b}$$
(42)
and $`V_3`$ is the volume of the unit 3 sphere:
$$V_3=d^3x_A\sqrt{g^s}=2\pi ^2.$$
(43)
In order that the variational principle is well-defined, the coefficients of $`\delta A_{,zz}`$ and $`\delta A_{,z}`$ should vanish. For general $`A`$, only one solution is given by the Weyl gravity in (5)
$`a={\displaystyle \frac{1}{6}}\widehat{c},b={\displaystyle \frac{4}{3}}\widehat{c},c=\widehat{c}`$
$`\stackrel{~}{a}={\displaystyle \frac{1}{6}}\stackrel{~}{\widehat{c}},\stackrel{~}{b}={\displaystyle \frac{4}{3}}\stackrel{~}{\widehat{c}},\stackrel{~}{c}=\stackrel{~}{\widehat{c}}`$
$`\stackrel{~}{\kappa }^2=\kappa ^2.`$ (44)
This remarkable property of Weyl gravity indicates to some natural connection between such version of HD gravity and brane physics.
When $`A`$ is given by the AdS<sub>5</sub> and the brane is S<sub>4</sub>,
$$A=\mathrm{ln}\mathrm{sinh}\frac{z}{l}\mathrm{ln}\mathrm{cosh}\sigma ,$$
(45)
Eq.(41) has the following form:
$`\delta \left(S+2S_b\right)`$ (46)
$`=`$ $`2V_3l^4{\displaystyle }d\sigma {\displaystyle \frac{\mathrm{sinh}^4\frac{z_0}{l}}{\mathrm{cosh}^4\sigma }}[(32\stackrel{~}{T}+24\stackrel{~}{U}){\displaystyle \frac{\mathrm{coth}\frac{z_0}{l}}{l}}\delta A_{,zz}`$
$`+\left\{{\displaystyle \frac{64\stackrel{~}{T}}{l^2}}\mathrm{coth}^2{\displaystyle \frac{z_0}{l}}+{\displaystyle \frac{32(T\stackrel{~}{T})}{l^2}}+{\displaystyle \frac{8}{\stackrel{~}{\kappa }^2}}{\displaystyle \frac{8}{\kappa ^2}}\right\}\delta A_{,z}`$
$`+\{{\displaystyle \frac{48\stackrel{~}{U}}{l^2\mathrm{sinh}^2\frac{z_0}{l}}}({\displaystyle \frac{8}{\kappa ^2}}+{\displaystyle \frac{32}{\stackrel{~}{\kappa }^2}}){\displaystyle \frac{\mathrm{coth}\frac{z_0}{l}}{l}}4\eta \}\delta A].`$
Then in order that the variational principle is well-defined, we obtain
$`0`$ $`=`$ $`32\stackrel{~}{T}+24\stackrel{~}{U}`$
$`0`$ $`=`$ $`\stackrel{~}{T}`$
$`0`$ $`=`$ $`{\displaystyle \frac{32(T\stackrel{~}{T})}{l^2}}+{\displaystyle \frac{8}{\stackrel{~}{\kappa }^2}}{\displaystyle \frac{8}{\kappa ^2}}`$ (47)
or
$$0=\stackrel{~}{T}=\stackrel{~}{U},\frac{1}{\stackrel{~}{\kappa }^2}=\frac{1}{\kappa ^2}\frac{4T}{l^2}.$$
(48)
The above results are consistent with those in . Then since $`\eta `$ in (23) is given by
$$\eta =\frac{6}{l\kappa ^2}\frac{24T}{l^3}$$
(49)
the equation of motion (in terms of the coefficients) has the following form:
$$0=\left(\frac{24}{\kappa ^2}+\frac{32T}{l^2}\right)\frac{\mathrm{coth}\frac{z_0}{l}}{l}\frac{24}{l\kappa ^2}+\frac{96T}{l^3}.$$
(50)
For the pure Einstein case ($`a=b=c=0`$ or $`T=0`$), the equation (50) reproduces the previous equation in by putting $`\kappa ^2=16\pi G`$. In the pure Einstein case, there is no solution of Eq.(50). Then for the case, we need to add the quantum correction coming from the trace anomaly of the matter fields on the brane in order that the equation corresponding to (50) has a non-trivial solution.
In case of the higher derivative gravity in (50), there can be a solution in general. The r.h.s. in Eq.(50) goes to positive infinity when $`z_0+0`$ if $`\frac{24}{\kappa ^2}+\frac{32T}{l^2}>0`$. On the other hand, the r.h.s. becomes $`\frac{128T}{l^2}`$ when $`z_0`$ goes to positive infinity. Then if $`T<0`$, there can be a solution in (50) without the quantum correction on the brane. As the r.h.s. is the monotonically increasing function of $`z_0`$, there is only one solution if $`T<0`$. We should also note that there does not appear corrections from $`R^2`$ gravity terms for the Weyl gravity (5), where $`T=0`$. For SG dual of $`𝒩=2`$ $`Sp(N)`$ theory, from (8), we find
$$T=\frac{6N}{2416\pi ^2}.$$
(51)
As $`T>0`$, there is no classical solution for spherical brane.
If we rewrite (50) as
$$\frac{\mathrm{coth}\frac{z_0}{l}}{l}=\frac{\frac{24}{l\kappa ^2}\frac{96T}{l^3}}{\frac{24}{\kappa ^2}+\frac{32T}{l^2}},$$
(52)
the r.h.s. is the monotonically increasing function of the absolute value $`|T|`$ of $`T`$ if $`T<0`$. Since $`\mathrm{coth}\frac{z_0}{l}`$ is the monotonically decreasing function of $`z_0`$, the radius $``$ of S<sub>4</sub>, which is given by
$$=l\mathrm{e}^{\stackrel{~}{A}(y_0)}=l\mathrm{sinh}\frac{z_0}{l},$$
(53)
decreases if $`|T|`$ increases when $`T<0`$ and $`l`$ is fixed. We should note that $`l`$ can be a function of $`T`$ since $`l`$ is given by (14), which is given in terms of $`T`$ as follows:
$$l^2=\frac{\frac{12}{\kappa ^2}\pm \sqrt{\frac{144}{\kappa ^4}32T\mathrm{\Lambda }}}{2\mathrm{\Lambda }}.$$
(54)
If we fix $`\mathrm{\Lambda }`$ instead of $`l`$, the situation becomes very complicated.
Using $``$ in (53), we can rewrite Eq.(52) in the following form:
$$0=\left(\frac{24}{l\kappa ^2}+\frac{32T}{l^3}\right)\sqrt{1+\frac{l^2}{^2}}\frac{24}{l\kappa ^2}+\frac{96T}{l^3}.$$
(55)
For SG dual of $`𝒩=2`$ $`Sp(N)`$ theory, using (15) and (51), one gets
$$0=\sqrt{1+\frac{1}{^2}}1+\frac{1}{48N}\sqrt{1+\frac{1}{^2}}\left(5\frac{1}{^2+1}\right)+\frac{22}{48N}+𝒪\left(N^2\right).$$
(56)
In this case, there is no any solution for $``$. It is remarkable that warped compactification to spherical brane is not realistic in leading (Einstein theory) as well as in next-to-leading order of AdS/CFT correspondence.
Instead of the brane of S<sub>4</sub> in (45), we can consider the brane of H<sub>4</sub>, where $`A`$ is given by
$$A=\mathrm{ln}\mathrm{cosh}\frac{z}{l}\mathrm{ln}\mathrm{sinh}\sigma .$$
(57)
By the similar calculation as for S<sub>4</sub>, we again obtain the Eqs.(48) and (49). The equation corresponding to (50) has the following form:
$$0=\left(\frac{24}{\kappa ^2}+\frac{32T}{l^2}\right)\frac{\mathrm{tanh}\frac{z_0}{l}}{l}\frac{24}{l\kappa ^2}+\frac{96T}{l^3}.$$
(58)
In case of pure Einstein gravity, there is no solution. When $`z_0=0`$, the r.h.s. in Eq.(58) becomes $`\frac{24}{l\kappa ^2}+\frac{96T}{l^3}`$, which can be regarded as negative. On the other hand, when $`z_0`$ goes to positive infinity, the r.h.s. becomes $`\frac{128T}{l^2}`$. Then if $`T>0`$, which is different from the case of the S<sub>4</sub> brane, there can be a solution in (58) without the quantum correction on the brane. Rewriting Eq.(58) in the form
$$\frac{\mathrm{tanh}\frac{z_0}{l}}{l}=\frac{\frac{24}{l\kappa ^2}\frac{96T}{l^3}}{\frac{24}{\kappa ^2}+\frac{32T}{l^2}},$$
(59)
we find the radius $`_\mathrm{H}`$ of H<sub>4</sub>, which is defined by
$$_\mathrm{H}=l\mathrm{e}^{\stackrel{~}{A}(y_0)}=l\mathrm{cosh}\frac{z_0}{l},$$
(60)
The radius $`_\mathrm{H}`$ is monotonically decreasing function of $`|T|`$ again if $`T>0`$ and $`l`$ is fixed since the l.h.s. in (59) is the monotonically increasing function of $`z_0`$ and the r.h.s. is the monotonically decreasing function of $`|T|`$ if $`T>0`$.
Using $`_\mathrm{H}`$ in (60), one can present Eq.(59) in the following form:
$$0=\left(\frac{24}{l\kappa ^2}+\frac{32T}{l^3}\right)\sqrt{1\frac{l^2}{_\mathrm{H}^2}}\frac{24}{l\kappa ^2}+\frac{96T}{l^3}.$$
(61)
For SG dual of $`𝒩=2`$ $`Sp(N)`$ theory, using (15) and (51), we have
$$0=\sqrt{1\frac{1}{_{\mathrm{H}}^{}{}_{}{}^{2}}}1+\frac{1}{48N}\sqrt{1\frac{1}{^2}}\left(5\frac{1}{_{\mathrm{H}}^{}{}_{}{}^{2}1}\right)+\frac{22}{48N}+𝒪\left(N^2\right).$$
(62)
For large $`_\mathrm{H}`$, Eq.(56) has the following form:
$$0=\frac{1}{2_{\mathrm{H}}^{}{}_{}{}^{2}}+\frac{2}{3N}+𝒪\left(_{\mathrm{H}}^{}{}_{}{}^{4}\right)+𝒪\left(N^2\right)+𝒪\left(N^1_{\mathrm{H}}^{}{}_{}{}^{2}\right)$$
(63)
or
$$\frac{1}{_\mathrm{H}^2}=\frac{4}{3N}+𝒪\left(N^2\right).$$
(64)
Thus, we demonstrated that next-to-leading order of AdS/CFT correspondence may qualitatively change the results on brane-world cosmology in the leading order. Indeed, in the leading order (Einstein theory) the warped compactification as 5d AdS with hyperbolic brane was impossible. On the same time, account of next-to-leading terms (on the example of particular SCFT dual) improves the situation: creation of hyperbolic brane in 5d AdS space becomes possible.
Let us discuss the situation where the higher derivative gravity in five dimensions corresponds to the Gauss-Bonnet combination which is topological invariant in four deimensions. Then $`a`$, $`b`$ and $`c`$ are given by
$$a=c=\widehat{a},b=4\widehat{a}.$$
(65)
One gets
$$T_{\mathrm{GB}}=3\widehat{a}.$$
(66)
Then all the discussion given above can be used by replacing $`T`$ by $`3\widehat{a}`$ (compare with independent calculation in ref.).
The situation is changed for the case that the brane is R<sub>4</sub>, where $`A`$ is given by
$$A=\frac{z}{l}+\sigma .$$
(67)
Since $`A_{,zz}=A_{,\sigma \sigma }=0`$, the coefficient of $`\delta A_{,z}`$ in (41) is given by
$$0=\frac{32T24U96\stackrel{~}{T}+72\stackrel{~}{U}}{l^2}\frac{8}{\kappa ^2}+\frac{8}{\stackrel{~}{\kappa }^2}$$
(68)
Then we obtain equations weaker than (48):
$$\stackrel{~}{U}=\frac{4}{3}\stackrel{~}{T},\frac{1}{\stackrel{~}{\kappa }^2}=\frac{1}{\kappa ^2}\frac{4T3U}{l^2},$$
(69)
and $`\eta `$ in (23) is given by
$$\eta =\frac{6}{l\kappa ^2}\frac{24T+6U4\stackrel{~}{T}}{l^3}.$$
(70)
The brane equation, which is the coefficient of $`\delta A`$ in (41) has the following form:
$`0`$ $`=`$ $`\left(160T128\stackrel{~}{T}\right){\displaystyle \frac{1}{l^3}}+\left({\displaystyle \frac{8}{\kappa ^2}}+{\displaystyle \frac{32}{\stackrel{~}{\kappa }^2}}\right){\displaystyle \frac{1}{l}}4\eta `$ (71)
$`=`$ $`{\displaystyle \frac{128T+120U144\stackrel{~}{T}}{l^3}}.`$
Then we have
$$\stackrel{~}{T}=\frac{8}{9}T+\frac{5}{6}U.$$
(72)
As one sees it admits the number of solutions for very large range of HD terms coefficients. Actually, chosing the suitable surface term the flat brane solution always exists. Thus, we explicitly showed that brane-world cosmology with spherical or hyperbolic or flat brane is possible for big class of HD gravities. The corresponding restrictions to HD terms coefficients are explicitly obtained. The version of HD gravity corresponding to next-to-leading order of AdS/CFT correspondence for specific SCFT is naturally included as sub-class of such theory.
## 3 Brane-Worlds with Account of Brane Quantum Matter
In the present section we will discuss the modification of the above scenario in the situation when quantum matter lives on the brane. Of course, bulk dynamics is not touched by brane quantum effects. It is interesting to remark also that in case of AdS/CFT correspondence the explanation of presence of such quantum brane matter effective action naturally appears via holographic renormalization group . In other words, the two dual descritions (SG dual and QFT one) could be patched together into the unique global description of some RG flow . Of course, we will consider general situation when for general HD gravity with arbitrary coefficients some quantum CFT lives on the brane.
The quantum correction induced by the trace anomaly of the free conformally invariant matter fields on the brane can be realized by adding the following effective action $`W`$ to $`S+S_b`$:
$`W`$ $`=`$ $`\widehat{b}{\displaystyle d^4x\sqrt{\stackrel{~}{g}}\stackrel{~}{F}A}`$ (73)
$`+b^{}{\displaystyle }d^4x\sqrt{\stackrel{~}{g}}\{A[2\stackrel{~}{\mathrm{\Delta }}^2+\stackrel{~}{R}_{\mu \nu }\stackrel{~}{}_\mu \stackrel{~}{}_\nu {\displaystyle \frac{4}{3}}\stackrel{~}{R}\stackrel{~}{\mathrm{\Delta }}^2+{\displaystyle \frac{2}{3}}(\stackrel{~}{}^\mu \stackrel{~}{R})\stackrel{~}{}_\mu ]A`$
$`+(\stackrel{~}{G}{\displaystyle \frac{2}{3}}\stackrel{~}{\mathrm{\Delta }}\stackrel{~}{R})A\}`$
$`{\displaystyle \frac{1}{12}}\left\{b^{\prime \prime }+{\displaystyle \frac{2}{3}}(b+b^{})\right\}{\displaystyle d^4x\sqrt{\stackrel{~}{g}}\left[\stackrel{~}{R}6\stackrel{~}{\mathrm{\Delta }}A6(\stackrel{~}{}_\mu A)(\stackrel{~}{}^\mu A)\right]^2}.`$
In (73), one chooses the 4 dimensional boundary metric as
$$g_{(4)}^{}{}_{\mu \nu }{}^{}=\mathrm{e}^{2A}\stackrel{~}{g}_{\mu \nu }$$
(74)
and we specify the quantities with $`\stackrel{~}{g}_{\mu \nu }`$ by using $`\stackrel{~}{}`$. $`G`$ ($`\stackrel{~}{G}`$) and $`F`$ ($`\stackrel{~}{F}`$) are the Gauss-Bonnet invariant and the square of the Weyl tensor:
$`G`$ $`=`$ $`R^24R_{ij}R^{ij}+R_{ijkl}R^{ijkl}`$
$`F`$ $`=`$ $`{\displaystyle \frac{1}{3}}R^22R_{ij}R^{ij}+R_{ijkl}R^{ijkl}.`$ (75)
In the effective action (73), with $`N`$ scalar, $`N_{1/2}`$ spinor, $`N_1`$ vector fields, $`N_2`$ ($`=0`$ or $`1`$) gravitons and $`N_{\mathrm{HD}}`$ higher derivative conformal scalars, $`\widehat{b}`$, $`b^{}`$ and $`b^{\prime \prime }`$ are
$`\widehat{b}`$ $`=`$ $`{\displaystyle \frac{N+6N_{1/2}+12N_1+611N_28N_{\mathrm{HD}}}{120(4\pi )^2}}`$
$`b^{}`$ $`=`$ $`{\displaystyle \frac{N+11N_{1/2}+62N_1+1411N_228N_{\mathrm{HD}}}{360(4\pi )^2}},`$
$`b^{\prime \prime }`$ $`=`$ $`0.`$ (76)
As usually, $`b^{\prime \prime }`$ may be changed by the finite renormalization of local counterterm in gravitational effective action. As we shall see later, the term proportional to $`\left\{b^{\prime \prime }+\frac{2}{3}(\widehat{b}+b^{})\right\}`$ in (73), and therefore $`b^{\prime \prime }`$, does not contribute to the equations of motion. For $`𝒩=4`$ $`SU(N)`$ SYM theory
$$\widehat{b}=b^{}=\frac{N^21}{4(4\pi )^2},$$
(77)
and for $`𝒩=2`$ $`Sp(N)`$ theory
$$\widehat{b}=\frac{12N^2+18N2}{24(4\pi )^2},b^{}=\frac{12N^2+12N1}{24(4\pi )^2}.$$
(78)
Notice that due to the structure of conformal anomaly the next-to-leading term for $`𝒩=4`$ super Yang-Mills theory dual is zero. Non-trivial term proportional to third power on curvatures appears in gravitational action as next-to-next-to-leading term. We should also note that $`W`$ in (73) is defined up to conformally invariant functional, which cannot be determined from only the conformal anomaly. The conformally flat space is an example where anomaly induced effective action is defined uniquely. However, one can argue that such conformally invariant functional is irrrelevant for us because it does not contribute to brane dynamics (does not depend on $`A`$).
In the choice of the metric (34), we find $`\stackrel{~}{F}=\stackrel{~}{G}=0`$, $`\stackrel{~}{R}=\frac{6}{l^2}`$ etc. and (73) looks
$`W`$ $`=`$ $`V_3{\displaystyle }d\sigma [b^{}A(2A_{,\sigma \sigma \sigma \sigma }8A_{,\sigma \sigma })`$ (79)
$`2(b+b^{})(1A_{,\sigma \sigma }(A_{,\sigma })^2)^2].`$
Under the variation over $`A`$, the change of $`W`$ is given by
$`\delta W`$ $`=`$ $`V_3l^4{\displaystyle }d\sigma \{4b^{}(A_{,\sigma \sigma \sigma \sigma }4A_{,\sigma \sigma })4(\widehat{b}+b^{})(A_{,\sigma \sigma \sigma \sigma }+2A_{,\sigma \sigma }`$ (80)
$`6(A_{,\sigma })^2A_{,\sigma \sigma })\}\delta A.`$
Then by substituting the solution (45), we find Eq.(50) is changed as
$$0=2\left\{\left(\frac{24}{\kappa ^2}+\frac{32T}{l^2}\right)\frac{\mathrm{coth}\frac{z_0}{l}}{l}\frac{24}{l\kappa ^2}+\frac{96T}{l^3}\right\}\mathrm{sinh}^4\frac{z_0}{l}+24b^{}.$$
(81)
Using the radius $``$ of S<sub>4</sub>, which is given in (53), Eq.(81) is rewritten
$$0=2\left\{\left(\frac{24}{l\kappa ^2}+\frac{32T}{l^3}\right)\sqrt{1+\frac{l^2}{^2}}\frac{24}{l\kappa ^2}+\frac{96T}{l^3}\right\}^4+24b^{}.$$
(82)
For H<sub>4</sub> brane, using the radius $`_\mathrm{H}`$ in (60), one gets
$$0=2\left\{\left(\frac{24}{l\kappa ^2}+\frac{32T}{l^3}\right)\sqrt{1\frac{l^2}{_\mathrm{H}^2}}\frac{24}{l\kappa ^2}+\frac{96T}{l^3}\right\}_\mathrm{H}^4+2b^{}.$$
(83)
For R<sub>4</sub> brane, the equation (71) or (72) is not changed.
For the case of S<sub>4</sub>, the l.h.s. of (82) goes to $`24b^{}`$ when $`0`$ and behaves as $`\frac{256T}{l^3}^4`$ for large $``$ if $`T0`$. Furthermore if $`T>0`$, the l.h.s. of (82) is the monotonically increasing function of $``$. Then if $`T>0`$, $`b^{}<0`$, there is a unique solution. Since when $`T>0`$ there is no solution for the classical case in (50), the solution for $`T>0`$ is generated by the quantum brane matter effects. We should note that even if $`b^{}>0`$ there is always a solution again if $`T<0`$. Here, HD gravity plays the essential role. On the other hand even if $`b^{}<0`$ and $`T<0`$, the l.h.s. of (82) has a unique maximum as a function of $``$. Then if the value of the maximum is positive, there are two solutions for $``$, which satisfies (82). Since the l.h.s. of (82) is the monotonically increasing function of $``$ when $`T>0`$, there is no any solution if $`T>0`$ and $`b^{}>0`$.
For the case of H<sub>4</sub>, the situation does not change if compare with S<sub>4</sub> if $`T0`$. This is because the behavior of the r.h.s. in (83) is again governed by the sign of $`T`$ when $`_\mathrm{H}`$ is large. Then if $`b^{}<0`$ and $`T>0`$ or if $`b^{}>0`$ and $`T<0`$, there is always a solution. The solution for $`T<0`$ is generated by the brane matter quantum effects. If $`T<0`$ and $`b^{}<0`$, there can be two (quantum) solutions. If $`b^{}<0`$ and $`T<0`$, there is no any solution.
The interesting example is provided by $`𝒩=2`$ $`Sp(N)`$ theory in the situation when SG dual and QFT descriptions are matched together via holographic RG, in both cases in next-to-leading order of AdS/CFT correspondence. In other words, using (15), (51) (SG dual up to next-to-leading order) and (78) (conformal anomaly for SCFT), in Eqs. (82) and (83) leads to
$`0`$ $`=`$ $`\left\{\sqrt{1+{\displaystyle \frac{1}{^2}}}1+{\displaystyle \frac{1}{48N}}\sqrt{1+{\displaystyle \frac{1}{^2}}}\left(5{\displaystyle \frac{1}{^2+1}}\right)+{\displaystyle \frac{22}{48N}}\right\}^4`$ (84)
$`{\displaystyle \frac{1}{384}}{\displaystyle \frac{1}{384N}}+𝒪\left(N^2\right),`$
$`0`$ $`=`$ $`\left\{\sqrt{1{\displaystyle \frac{1}{_{\mathrm{H}}^{}{}_{}{}^{2}}}}1+{\displaystyle \frac{1}{48N}}\sqrt{1{\displaystyle \frac{1}{^2}}}\left(5{\displaystyle \frac{1}{_{\mathrm{H}}^{}{}_{}{}^{2}1}}\right)+{\displaystyle \frac{22}{48N}}\right\}_{\mathrm{H}}^{}{}_{}{}^{4}`$ (85)
$`{\displaystyle \frac{1}{384}}{\displaystyle \frac{1}{384N}}+𝒪\left(N^2\right).`$
Since $`T>0`$ and $`b^{}<0`$, there are always solutions in (84) and (85) for finite $``$ and $`_\mathrm{H}`$. When the brane is S<sub>4</sub> in (84), $``$ becomes $`𝒪(1)`$ and the higher derivative terms give a correction of $`𝒪\left(N^1\right)`$. When $`N\mathrm{}`$, the solution of (84) is numerically given by $`^2=0.020833\mathrm{}`$. Substituting this value into the r.h.s. of (84), we find that it takes a negative value of $`\frac{0.00215076}{N}`$. Since $`T>0`$, the r.h.s. is monotonically increasing function of $``$ and goes negative when $`0`$. Then the above result tells that the correction of $`𝒪\left(N^1\right)`$ makes $``$ large. Since $`\frac{1}{}`$ corresponds to the rate of inflation when we Wick-rotate S<sub>4</sub> to de Sitter space, the correction makes the rate small. This is some indication that realistic inflationary cosmology may not be comfortable with warped compactification in AdS/CFT correspondence.
For the brane of H<sub>4</sub>, the quantum correction of $`𝒪(N^2)`$ to $`\frac{1}{_\mathrm{H}^2}`$ of the classical solution in (64) exists but since further higher derivative gravity terms like $`R^4`$ also give the contribution of $`𝒪(N^2)`$, the correction is beyond the control.
Hence, we demonstrated that role of quantum brane matter may be in the significant change of bulk/boundary structure. As we saw there exists the range of HD terms coefficients for which the creation of inflationary or hyperbolic Universe living in d5 AdS is caused exclusively by brane quantum effects. It could be also relevant in frames of AdS/CFT set-up where correct holographic RG description shows the necessity of anomaly induced effective action of brane CFT. In its own turn, the corresponding quantum effects change the brane structure and indicate (despite the negative results of leading order analysis) to the possibility of quantum creation of inflationary brane in d5 AdS space in the next-to-leading order of AdS/CFT correspondence.
## 4 Discussion
In summary, we investigated brane-world Universe solutions (of special form) for five-dimensional higher derivative gravity. It is shown that such Universe occurs for range of theory parameters. As brane part may be given by de Sitter space which after analytical continuation to Lorentzian signature represents ever expanding inflationary Universe then such configuration could be relevant to observable world. The particular examples of Weyl, Gauss-Bonnet or SG dual to some SCFT are also examinated. The role of brane quantum CFT is investigated in the quantum creation of spherical or hyperbolic brane Universes.
There are few interesting topics which may be left for future studies. First of all, one has to investigate the structure of HD propagator near brane in order to understand in detail how HD gravity is trapped. In other words, graviton profile and corrections to Newton potential should be estimated. Second, the dilaton may be included into the analysis of this paper. However, the number of HD terms in dilatonic gravity grows significally. As a result, the analysis is getting too complicated technically. Nevertheless, it could be done at least for some truncated versions of HD dilatonic gravity (say, conformally invariant theory or dilatonic Gauss-Bonnet). Note also that some versions of such theory may be considered as SG duals for non-commutative (super) Yang-Mills theory (presumbly in next-to-leading order). Third, other cosmologies may be considered in the same fashion where bulk and (or) boundary is modified. In particular, the situation where bulk is AdS black hole and boundary is some FRW Universe (or vice-versa) deserves careful study. Fourth, it would be interesting to discuss the cosmological perturbations around our background and the details of late-time inflation. For example, in Einstein gravity the domain wall CFT significally suppresses the metric perturbations . What will be the role of HD gravitational terms in such phenomenon?
## Acknoweledgements
We thank O. Obregon and V. Tkach for helpful discussions. The work by SDO has been supported in part by CONACyT (CP, ref.990356 and grant 28454E) and in part by RFBR.
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# Introduction
## Introduction
Unimodular integer matrices induce interesting dynamical systems on the torus, such as Arnold’s famous cat map \[4, Ch. 1, Ex. 1.16\]. This is an example of a hyperbolic dynamical system that is ergodic and mixing , and also a topological Anosov system. Therefore, with a suitable metric, it makes the 2-torus into a Smale space, see \[28, Thm. 1.2.9\] for details. Although induced from a linear system of ambient space, the dynamics on the torus is rather complicated, and these systems serve as model systems in symbolic dynamics and in many applications. Very recently, also cat maps on the 4-torus (and their quantizations) have begun to be studied .
Hyperbolic toral automorphisms also play a prominent role in the theory of quasicrystals through their appearance as inflation symmetries, see and references therein. In particular, in the one-dimensional case, these symmetries give rise to interesting non-linear dynamical systems called trace maps that have extensively been used to study the physical properties of one-dimensional quasicrystals.
It is always helpful to know the symmetries and reversing symmetries of a dynamical system , and it was perhaps a little surprising that in the case of $`\mathrm{GL}(2,)`$ and $`\mathrm{PGL}(2,)`$ a rather complete classification could be given, see for a detailed account or for a summary. The main difficulty when the matrix entries are restricted to integers is that one can no longer refer to the usual normal forms of matrices over $``$ or $``$, but has to use discrete methods instead. Fortunately, there is a strong connection with algebraic number theory, see for an introduction, and this connection is certainly not restricted to the 2D situation.
It is thus the aim of this article to extend the results of our earlier article to the setting of matrices in $`\mathrm{GL}(n,)`$. The answers will be less complete and also less explicit, but the connection to unit groups in orders of algebraic number fields is still strong enough to give quite a number of useful and general results, both on symmetries and reversing symmetries. From a purely algebraic point of view, the results derived below are actually rather straight-forward. However, these results, and the methods used to derive them, are not at all common in the dynamical systems community. Therefore, this article is also intended to introduce some of these techniques, and we try to spell out the details or give rather precise references at least. Furthermore, as with our article , the results of this paper have relevance to both the dynamics community (e.g. the dynamics of hyperbolic toral automorphisms generated by (symplectic) $`\mathrm{SL}(4,)`$ matrices ) and to the quasicrystal community (e.g. inflation symmetries of planar point sets projected from 4D lattices, where the symmetries are generated by $`\mathrm{GL}(4,)`$ matrices ).
The article is organized as follows. We start with a section on the background material, including the group theoretic setup we use and a recollection of those results from algebraic number theory that we will need later on. Section 2 is the main part of this article. Here, we derive the structure of the symmetry group of toral automorphisms with simple spectrum and discuss reversibility. Section 3 extends the set of possible symmetries to affine transformations and summarizes the analogous problem for projective matrices. It also discusses the extension of (reversing) symmetries to matrices that are no longer unimodular.
## 1 Setting the scene
In this Section, we explain in more detail what we mean by symmetries and reversing symmetries, and we also recall some results from algebra and algebraic number theory that we will need.
### 1.1 Symmetries and reversing symmetries
For a general setting, consider some (topological) space $`\mathrm{\Omega }`$, and let $`\mathrm{Aut}(\mathrm{\Omega })`$ be its group of homeomorphisms or, more generally, a subgroup of homeomorphisms of $`\mathrm{\Omega }`$ which preserve some additional structure of $`\mathrm{\Omega }`$. Consider now an element $`F\mathrm{Aut}(\mathrm{\Omega })`$ which, by definition, is invertible. Then, the group
$$𝒮(F):=\{G\mathrm{Aut}(\mathrm{\Omega })|GF=FG\}$$
(1)
is called the symmetry group of $`F`$ in $`\mathrm{Aut}(\mathrm{\Omega })`$. In group theory, it is called the centralizer of $`F`$ in $`\mathrm{Aut}(\mathrm{\Omega })`$, denoted by $`\mathrm{cent}_{\mathrm{Aut}(\mathrm{\Omega })}(F)`$. This group certainly contains all powers of $`F`$, but often more.
Quite frequently, one is also interested in mappings $`R\mathrm{Aut}(\mathrm{\Omega })`$ that conjugate $`F`$ into its inverse,
$$RFR^1=F^1.$$
(2)
Such $`R`$ is called a reversing symmetry of $`F`$, and when such an $`R`$ exists, we call $`F`$ reversible. We will, in general, not use different symbols for symmetries and reversing symmetries from now on, because together they form a group,
$$(F):=\{G\mathrm{Aut}(\mathrm{\Omega })|GFG^1=F^{\pm 1}\},$$
(3)
the so-called reversing symmetry group of $`F`$, see for details. If $`F`$ denotes the group generated by $`F`$, $`(F)`$ is a subgroup of the normalizer of $`F`$ in $`\mathrm{Aut}(\mathrm{\Omega })`$.
There are two possibilities: either $`(F)=𝒮(F)`$ (if $`F`$ is an involution or if it has no reversing symmetry) or $`(F)`$ is a $`C_2`$-extension (the cyclic group of order 2) of $`𝒮(F)`$ which means that $`𝒮(F)`$ is a normal subgroup of $`(F)`$ and the factor group is
$$(F)/𝒮(F)C_2.$$
(4)
The underlying algebraic structure has fairly strong consequences. One is that reversing symmetries cannot be of odd order , another one is the following product structure \[10, Lemma 2\].
###### Fact 1
If $`F`$ $`(`$with $`F^2Id)`$ has an involutory reversing symmetry $`R`$, the reversing symmetry group of $`F`$ is given by
$$(F)=𝒮(F)\times _sC_2,$$
(5)
i.e. it is a semi-direct product.<sup>1</sup><sup>1</sup>1We use $`N\times _sH`$ for the semi-direct product of two groups $`N`$ and $`H`$, with $`N`$ being the normal subgroup. $`\mathrm{}`$
We can say more about the structure of $`(F)`$ if we restrict the possibilities for $`𝒮(F)`$, e.g. if we assume that $`𝒮(F)C_{\mathrm{}}`$ or $`𝒮(F)C_{\mathrm{}}\times C_2`$ with the $`C_2`$ being a subgroup of the centre of $`\mathrm{Aut}(\mathrm{\Omega })`$, compare also for some group theoretic discussion. This situation will appear frequently below.
### 1.2 (Reversing) symmetries of powers of a mapping
In what follows, we summarize some of the concepts and results of Ref. and, in particular, Ref. . It may happen that some power of $`F`$ has more symmetries than $`F`$ itself (we shall see examples later on), i.e. $`𝒮(F^k)`$ (for some $`k>1`$) is larger than $`𝒮(F)`$ which is contained as a subgroup. The analogous possibility exists for $`(F^k)`$ versus $`(F)`$. If such a situation occurs, we say that $`F`$ possesses additional (reversing) $`k`$-symmetries. Let us make this a little more precise.
It is trivial that mappings $`F`$ of finite order (with $`F^k=Id`$, say) possess the entire group $`\mathrm{Aut}(\mathrm{\Omega })`$ as $`k`$-symmetry group. Let us thus concentrate on mappings $`F\mathrm{Aut}(\mathrm{\Omega })`$ of infinite order. We denote by $`𝒮_{\mathrm{}}(F)`$ the set of automorphisms that commute with some positive power of $`F`$. This set can be seen as the inductive limit of $`𝒮(F^k)`$ as $`k\mathrm{}`$, with divisibility as partial order on $``$, and $`𝒮_{\mathrm{}}(F)`$ is thus a subgroup of $`\mathrm{Aut}(\mathrm{\Omega })`$. Let $`\mathrm{\#}_F(G)`$ denote the minimal $`k`$ such that $`GF^k=F^kG`$. Then
$$𝒮_{\mathrm{}}(F)=\{G\mathrm{Aut}(\mathrm{\Omega })|\mathrm{\#}_F(G)<\mathrm{}\}.$$
(6)
Of course it may happen that $`\mathrm{\#}_F(G)1`$ on $`𝒮_{\mathrm{}}(F)`$ which means that no power of $`F`$ has additional symmetries. On the other hand, $`\mathrm{\#}_F(G)`$ might be larger than one in which case we call $`G`$ a genuine or true $`k`$-symmetry. $`G`$ is a true<sup>2</sup><sup>2</sup>2Although the distinction between true and other $`k`$-symmetries is necessary in general, we shall usually drop the attribute “true” whenever misunderstandings are unlikely. $`k`$-symmetry of $`F`$ if and only if the mapping $`GFGF^1`$ generates a proper $`k`$-cycle. We shall meet this phenomenon later on.
Quite similarly, one defines reversing $`k`$-symmetries and their orbit structure , but we will not expand on that here.
### 1.3 Some recollections from algebraic number theory
Much of what we state and prove below can be seen as an application of several well-known results from algebraic number theory. The starting point is the connection between algebraic number theory and integral matrices, see for an introduction.
To fix notation, let $`\mathrm{Mat}(n,)`$ denote the ring of integer<sup>3</sup><sup>3</sup>3Here, and in what follows, integer means rational integer, i.e. an integer in $``$. Other kinds of integers, such as algebraic, will be specified explicitly. $`n\times n`$-matrices. An element $`M`$ of it is called unimodular if $`det(M)=\pm 1`$, and the subset of all unimodular matrices forms the group $`\mathrm{GL}(n,)`$. For the characteristic polynomial of a matrix $`M`$, we will use the convention
$$P(x):=det(x\mathrm{𝟏}M)=\underset{i=1}{\overset{n}{}}(x\lambda _i)$$
(7)
where $`\lambda _1,\mathrm{},\lambda _n`$ denote the eigenvalues of $`M`$. With this convention, $`P(x)`$ is monic, i.e. its leading coefficient is $`1`$. If $`M`$ is an integer matrix, $`P(x)`$ has integer coefficients only, so all eigenvalues of $`M`$ are algebraic integers. Conversely, the set of algebraic integers, which we denote by $`𝒜`$, consists of all numbers that appear as roots of monic integer polynomials.
To show the intimate relation more clearly, let us recall the following property (see item (b) on p. 306 of ):
###### Fact 2
Let $`P(x)=x^n+a_{n1}x^{n1}+\mathrm{}+a_1x+a_0`$ be a monic polynomial with integer coefficients $`a_i`$ that is irreducible over $``$. Let $`\alpha `$ be any root of it, and $`A`$ an integer matrix that has $`P(x)`$ as its characteristic polynomial. Then, the rings $`[\alpha ]`$ and $`[A]`$ are isomorphic. $`\mathrm{}`$
We write $`[x]`$ for the ring of polynomials in $`x`$ with coefficients in $``$, see \[23, p. 90\] for details. Clearly, Fact 2 also extends to the isomorphism of the rings $`[\alpha ]`$ and $`[A]`$. For background material on polynomial rings, we refer to \[23, Ch. IV\]. Let us only add that a polynomial in $`[x]`$ is irreducible over $``$ if and only if it is irreducible over $``$, see \[23, Thm. IV.2.3\].
Let $`P(x)`$ be any (i.e. not necessarily irreducible) monic integer polynomial of degree $`n`$. In general, there are many different matrices $`A`$ which have $`P(x)`$ as their characteristic polynomial, and even different matrix classes (we say that $`A,B\mathrm{Mat}(n,)`$ belong to the same matrix class if they are conjugate by a $`\mathrm{GL}(n,)`$ matrix $`C`$, i.e. $`A=CBC^1`$). Let us recall the following helpful result on the number of matrix classes, see \[37, Thm. 5\] for details:
###### Fact 3
Let $`P(x)`$ be a monic polynomial of order $`n`$ with integer coefficients that is irreducible over $``$, and let $`\alpha `$ be any of its roots. Then the number of matrix classes generated by matrices $`A\mathrm{Mat}(n,)`$ with $`P(A)=0`$ equals the number of ideal classes $`(`$or class number, for short$`)`$ of the order $`[\alpha ]`$. In particular, this class number is finite, and it is larger than or equal to the class number of the maximal order $`𝒪_{\mathrm{max}}`$ of $`(\alpha )`$. $`\mathrm{}`$
Let us explain some of the terms used here. If $`\alpha `$ is an algebraic number, $`(\alpha )`$ denotes<sup>4</sup><sup>4</sup>4Note the difference between the meaning of $`[\alpha ]`$ and $`(\alpha )`$. the smallest field extension of the rationals that contains $`\alpha `$. Its degree, $`n`$, is the degree of the irreducible monic integer polynomial that has $`\alpha `$ as its root. The set $`𝒪_{\mathrm{max}}:=(\alpha )𝒜`$ is the ring of (algebraic) integers in $`(\alpha )`$, and is called its maximal order. More generally, a subring $`𝒪`$ of $`𝒪_{\mathrm{max}}`$ is called an order, if it contains 1 and if its rational span, $``$$`𝒪`$, is all of $`(\alpha )`$. A subset of $`𝒪`$ is called an ideal if it is both a $``$-module (i.e. closed under addition and subtraction) and closed under multiplication by arbitrary numbers from $`𝒪`$. The ideals come in classes that are naturally connected to the matrix classes introduced above, see for further details.
An element $`\epsilon (\alpha )`$ is called a unit (or, more precisely<sup>5</sup><sup>5</sup>5This distinction is useful if orders other than $`𝒪_{\mathrm{max}}`$ appear., a unit in $`𝒪_{\mathrm{max}}`$) if both $`\epsilon `$ and its inverse, $`\epsilon ^1`$, are algebraic integers and hence are in $`𝒪_{\mathrm{max}}`$. This happens if and only if the corresponding matrix is unimodular, i.e. if any monic integer polynomial $`P(x)`$ that has $`\epsilon `$ as a root has coefficient $`a_0=\pm 1`$. So, matrices in $`\mathrm{GL}(n,)`$ and units in algebraic number fields are two facets of the same coin. The units of $`𝒪_{\mathrm{max}}`$ form a group under multiplication, denoted by $`𝒪_{\mathrm{max}}^\times `$ in the sequel. Similarly, if $`𝒪𝒪_{\mathrm{max}}`$ is any order of $`(\alpha )`$, we write $`𝒪^\times `$ for its group of units, which is then a subgroup of $`𝒪_{\mathrm{max}}^\times `$.
Independently of whether the monic polynomial $`P(x)`$ is irreducible or not, there is always at least one matrix which has characteristic polynomial $`P(x)`$, namely the so-called (left) companion matrix:
$$A^{(\mathrm{})}=\left(\begin{array}{ccccc}0& 1& & & \mathrm{𝟎}\\ 0& 0& 1& & \\ \mathrm{}& & \mathrm{}& \mathrm{}& \\ 0& & & 0& 1\\ a_0& & & a_{n2}& a_{n1}\end{array}\right).$$
(8)
Consequently, the class number is $`1`$ (if $`P(x)`$ is irreducible, the companion matrix actually corresponds to the principal ideal class, see \[37, Thm. 9\]). Another obvious choice, always belonging to the same matrix class, is the (right) companion matrix, $`A^{(r)}`$, obtained from $`A^{(\mathrm{})}`$ by reflection in both diagonal and anti-diagonal, i.e.
$$A^{(r)}=RA^{(\mathrm{})}R$$
(9)
with the involution
$$R=\left(\begin{array}{ccc}\mathrm{𝟎}& & 1\\ & \text{ }\text{ }& \\ 1& & \mathrm{𝟎}\end{array}\right).$$
(10)
This matrix will reappear several times in what follows.
## 2 $`\mathrm{GL}(n,)`$ matrices and toral automorphisms
Let us generally assume that $`n2`$. The toral automorphisms of the $`n`$-torus $`𝕋^n:=^n/^n`$ can be represented by the unimodular $`n\times n`$-matrices with integer coefficients which form the group $`\mathrm{GL}(n,)`$. It now plays the role of $`\mathrm{Aut}(\mathrm{\Omega })`$ from Section 1.1. Note that the elements of $`\mathrm{GL}(n,)`$ preserve the linear structure of the torus.
### 2.1 Symmetries
The first thing we will look at, given a toral automorphism $`MGL(n,)`$, is its symmetry group within the class of toral automorphisms. So, we want to determine the centralizer of $`M`$ in $`\mathrm{GL}(n,)`$,
$$𝒮(M)=\text{cent}_{\mathrm{GL}(n,)}(M)=\{G\mathrm{GL}(n,)|MG=GM\}.$$
(11)
To be more precise, we are mainly interested in the structure of the symmetry group rather than in explicit sets of generators and relations. This is invariant under conjugation, i.e. if we know it for an element $`M`$, we also know it for any other element of the form $`BMB^1`$ because
$$𝒮(BMB^1)=B𝒮(M)B^1.$$
(12)
A given integer matrix $`M\mathrm{GL}(n,)`$ determines its characteristic polynomial (7) which is monic and has integer coefficients, so its roots are algebraic integers. Now, two principal situations can occur for the characteristic polynomial: it is either reducible over $``$ (which happens if and only if at least one eigenvalue of $`M`$ is an algebraic integer of degree less than $`n`$) or it is irreducible. In the latter case, since we are working over the field $``$, we know that the roots must be pairwise distinct. So we have
###### Fact 4
Let $`M`$ be an integer matrix with irreducible characteristic polynomial. Then $`M`$ is simple and hence diagonalizable over $``$. $`\mathrm{}`$
Here, $`M`$ is called simple if it has no repeated eigenvalues (which is also called separable elsewhere). This case will be dealt with completely.
If the characteristic polynomial is reducible, the matrix can still be simple, and we will see the general answer for this case, too. Beyond that, $`M`$ can either be semi-simple (i.e. diagonalizable over $``$) or not, and we will not be able to say much about this case. This is really not surprising, as this situation is closely related to the rather difficult classification problem of crystallographic point groups, see for answers in dimensions $`4`$ and for a recent survey.
Let us now state one further prerequisite for tackling the symmetry question. In view of later extensions, we do this in slightly more generality than needed in the present Section. Recall that an $`n\times n`$-matrix $`M`$, acting on a vector space $`V`$, is called cyclic, if a vector $`vV`$ exists such that $`\{v,Mv,M^2v,\mathrm{},M^{n1}v\}`$ is a basis of $`V`$. Also, the monic polynomial $`Q`$ of minimal degree that annihilates $`M`$, i.e. $`Q(M)=0`$, is called the minimal polynomial of $`M`$. By the Cayley-Hamilton Theorem, it always is a factor of the characteristic polynomial of $`M`$.
###### Fact 5
Let $`M\mathrm{Mat}(n,)`$ be a rational matrix, with characteristic polynomial $`P(x)`$ and minimal polynomial $`Q(x)`$. Then the following assertions are equivalent.
* The matrix $`M`$ is cyclic.
* The degree of $`Q(x)`$ is $`n`$.
* $`P(x)=Q(x)`$.
* $`G\mathrm{Mat}(n,)`$ commutes with $`M`$ $``$ $`G[M]`$.
Proof: A convenient source is \[18, Ch. III\]. The equivalence of statements ($`a`$) – ($`c`$) is a consequence of Thm. III.2. The equivalence of ($`a`$) with ($`d`$) follows from Thm. III.17, and the Corollary following it, together with the Corollary of Ch. III.17. Alternatively, see \[3, Cor. 5.5.16\]. $`\mathrm{}`$
###### Lemma 1
Let $`M\mathrm{GL}(n,)`$ have a characteristic polynomial $`P(x)`$ that is irreducible over $``$, and let $`\lambda `$ be a root of $`P(x)`$. Then the centralizer of $`M`$ in $`\mathrm{GL}(n,)`$ is isomorphic to a subgroup of finite index of the unit group in the ring of integers $`𝒪_{\mathrm{max}}`$ of the algebraic number field $`(\lambda )`$.
Proof: By assumption, $`P(x)`$ is also the minimal polynomial of $`M`$, so any $`\mathrm{GL}(n,)`$-matrix which commutes with $`M`$ is, by Fact 5, a polynomial in $`M`$ with rational coefficients. Consequently, $`𝒮(M)`$ is isomorphic to a subset of $`[M]`$ that forms a group under (matrix) multiplication. So, we have to analyze $`[M]`$ to find out what this group is.
Let $`(P(x))`$ denote the ideal in $`[x]`$ generated by our polynomial $`P(x)`$. Then $`[x]/(P(x))[\lambda ]`$, see \[23, p. 224\], and $`[\lambda ][M]`$, by Fact 2 resp. the remark following it. Since $`\lambda `$ is algebraic over $``$ and $`P(x)`$ is irreducible, we know by \[23, Prop. V.1.4\] that $`[\lambda ]=(\lambda )`$ is an algebraic number field, of degree $`n`$ over $``$. Under the isomorphism $`[M](\lambda )`$, $`\mathrm{GL}(n,)`$-matrices correspond to units in $`(\lambda )`$, hence $`𝒮(M)`$ must be isomorphic to a subgroup of $`𝒪_{\mathrm{max}}^\times `$, the unit group of the maximal order of $`(\lambda )`$.
Observe that every matrix in $`[M]`$ commutes with $`M`$, in particular those of $`[M]\mathrm{GL}(n,)`$, which form a subgroup of $`𝒮(M)`$. But $`[M][\lambda ]`$ means that this subgroup is isomorphic to the unit group $`[\lambda ]^\times `$. So, if we identify $`𝒮(M)`$ with its image in $`(\lambda )`$ under the isomorphism, it is sandwiched between $`[\lambda ]^\times `$ and $`𝒪_{\mathrm{max}}^\times `$. Note that $`[\lambda ]𝒪_{\mathrm{max}}`$ is an order.
Finally, recall that the unit group of an order $`𝒪`$ is a finitely generated Abelian group, and that it is always of maximal rank, see \[25, Thm. I.12.12\] or \[11, Sec. 4, Thm. 5\], i.e. its rank equals that of the unit group of the maximal order $`𝒪_{\mathrm{max}}`$. In particular, the group-subgroup index $`[𝒪_{\mathrm{max}}^\times :[\lambda ]^\times ]`$ is finite, and $`𝒮(M)`$ must then also be of finite index in $`𝒪_{\mathrm{max}}^\times `$. $`\mathrm{}`$
Let us comment on this result. First of all, it does not matter which root $`\lambda _i`$ of $`P(x)`$ we choose, as all the $`n`$ (possibly different) realizations $`(\lambda _i)`$ are mutually isomorphic, and so are their unit groups. Explicit isomorphisms are given by the elements of the Galois group of the splitting field $`𝒦=(\lambda _1,\mathrm{},\lambda _n)`$ of $`P(x)`$, see \[23, Ch. VI.2\] for details. Note also that, in general, $`(\lambda )`$ will be a true subfield of the splitting field $`𝒦`$ – so, it is really the unit group of $`(\lambda )`$ that matters, and not the unit group of $`𝒦`$.
Another way to view the result, in a more matrix oriented way (and similar to our approach in ), is to look at the diagonalization of $`M`$,
$$UMU^1=\mathrm{diag}(\lambda _1,\mathrm{},\lambda _n).$$
(13)
Here, $`U^1`$ can be arranged to have its $`j`$-th column in the field $`(\lambda _j)`$ because one can solve the corresponding eigenvector equation in the smallest field extension of $``$ that contains $`\lambda _j`$. In fact, we only have to do this for the first column — the others are then obtained by applying appropriate Galois automorphisms to the first one.
Any other matrix $`G\mathrm{GL}(n,)`$ with $`[G,M]=GMMG=0`$ must now also fulfil
$$[UGU^1,UMU^1]=U[G,M]U^1=\mathrm{\hspace{0.33em}0}.$$
But only diagonal matrices can commute with $`\mathrm{diag}(\lambda _1,\mathrm{},\lambda _n)=UMU^1`$ because the eigenvalues are pairwise distinct. So, we must have
$$UGU^1=\mathrm{diag}(\mu _1,\mathrm{},\mu _n),$$
with all $`\mu _i(\lambda _i)`$ units. They are, however, not independent but obtained from one another by the same set of Galois automorphisms that were used to link the columns of the matrix $`U^1`$, which is why we get the result.
Lemma 1 raises the question: What is the unit group of the maximal order in $`K=(\lambda )`$? The answer is given by Dirichlet’s unit theorem, see \[13, Sec. 11.C\] or \[30, p. 334\]. Group the roots of the irreducible polynomial $`P(x)`$ into $`n_1`$ real roots and $`n_2`$ pairs of complex conjugate roots, so that $`n=n_1+2n_2`$. (In other words: we have $`n_1`$ real and $`n_2`$ pairs of complex conjugate realizations of the abstract number field $`(\lambda )`$).
###### Fact 6
Let $`\lambda `$ be an algebraic number of degree $`n=n_1+2n_2`$, with $`n_1`$ and $`n_2`$ as described above. Then, the units in the maximal order $`𝒪_{\mathrm{max}}^\times `$ of the algebraic number field $`K=(\lambda )`$ form the group
$$E(K)=𝒪_{\mathrm{max}}^\times T\times ^{n_1+n_21}$$
(14)
where $`T=𝒪_{\mathrm{max}}\{\text{roots of unity}\}`$ is a finite Abelian group and cyclic. $`\mathrm{}`$
In particular, this means that $`T`$, which is also called the torsion subgroup of $`E(K)`$, is generated by one element. In many cases below, we will simply find $`TC_2`$. Combining now Lemma 1 with Fact 6, we immediately obtain
###### Proposition 1
Under the assumptions of Lemma 1, the symmetry group $`𝒮(M)\mathrm{GL}(n,)`$ is a subgroup of $`E(K)`$ of $`(\text{14})`$ of maximal rank, i.e. we have
$$𝒮(M)T^{}\times ^{n_1+n_21}$$
where $`T^{}`$ is a subgroup of the torsion group $`T`$ as it appears in $`(\text{14})`$. $`\mathrm{}`$
Note that Proposition 1 does not imply that the torsion-free parts of $`𝒮(M)`$ and $`E(K)`$ are the same, only that they are isomorphic. In fact, a typical situation will be that they are different in the sense that $`E(K)`$ is generated by the fundamental units, but $`𝒮(M)`$ only by suitable powers thereof.
What, in turn, can we say about the torsion group $`T^{}`$ in Proposition 1? Whenever the characteristic polynomial $`P(x)`$ of $`M\mathrm{GL}(n,)`$ is irreducible and has at least one real root (e.g. if $`n`$ is odd), $`\alpha `$ say, then $`K=(\alpha )`$ is real, and $`(\alpha )S^1=\{\pm 1\}`$, where $`S^1`$ is the unit circle. Consequently, the torsion subgroup of $`E(K)`$ in this case is $`T=\{\pm 1\}C_2`$. Since a toral automorphism always commutes with $`\pm \mathrm{𝟏}`$, we obtain
###### Corollary 1
If, under the assumptions of Lemma 1, one root of the irreducible polynomial $`P(x)`$ is real, the torsion group in Proposition 1 is $`T^{}C_2`$. In particular, this is the case whenever the degree of $`P(x)`$ is odd. $`\mathrm{}`$
Let us look at two examples in $`\mathrm{GL}(3,)`$, namely
$$M_1=\left(\begin{array}{ccc}0& 0& 1\\ 1& 0& 0\\ 0& 1& 1\end{array}\right)\text{ and }M_2=\left(\begin{array}{ccc}1& 1& 0\\ 1& 0& 1\\ 1& 1& 1\end{array}\right),$$
(15)
which are taken from \[24, Eqs. (5.21) and (5.3)\]. They have been studied thoroughly in the context of inflation generated one-dimensional quasicrystals with a cubic irrationality as inflation factor. We have $`det(M_1)=1`$, $`det(M_2)=1`$, and the characteristic polynomials are $`P_1(x)=x^3x^21`$ and $`P_2(x)=x^32x^2x+1`$, both irreducible over $``$. Both matrices are hyperbolic, and the largest eigenvalue in each case is a Pisot-Vijayaraghavan number, i.e. an algebraic integer $`>1`$ all algebraic conjugates of which lie inside the unit circle, see for details.
Now, $`M_1`$ has one real and a pair of complex conjugate roots, so our above results lead to $`𝒮(M_1)C_2\times `$, where the infinite cyclic group is actually generated by $`M_1`$ itself because its real root is a fundamental unit. $`M_2`$, in turn, has three real roots, and we thus get $`𝒮(M_2)C_2\times ^2`$. As generators of $`^2`$, one may choose $`M_2`$ and $`M_2^{}`$ where
$$M_2^{}=\left(\begin{array}{ccc}0& 1& 0\\ 1& \text{-}1& 1\\ 1& 1& 0\end{array}\right)$$
which can be checked explicitly. This example is of relevance in connection with planar quasicrystals with sevenfold symmetry, see \[24, Sec. 5.2\] for details, where the three-dimensional toral automorphism $`M_2`$ shows up in the cut and project description of special directions in the quasicrystal. Other examples related to planar quasicrystals with 8-, 10- and 12-fold symmetry, in which the torsion subgroup $`T^{}`$ of Proposition 1 is different from $`C_2`$, will be given in Section 2.2.
Let us now return to the general discussion and extend the previous results to the case that $`M`$ is simple, and hence diagonalizable (over $``$) with pairwise different eigenvalues. Since diagonal matrices with pairwise different entries only commute with diagonal matrices, we have:
###### Corollary 2
If $`M\mathrm{GL}(n,)`$ is simple, $`𝒮(M)\mathrm{GL}(n,)`$ is Abelian. $`\mathrm{}`$
Note that the converse is not true: even if $`M`$ is only semi-simple, or not even that (i.e. not diagonalizable), $`𝒮(M)`$ can still be Abelian, e.g. if $`M`$ observes the conditions of Fact 5. As far as we are aware, not even the Abelian subgroups of $`\mathrm{GL}(n,)`$ are fully classified, see and references given there for background material.
Let $`P(x)`$ be the characteristic polynomial of a simple matrix $`M`$. If it is reducible over $``$, it factorizes as $`P(x)=_{i=1}^{\mathrm{}}P_i(x)`$ into irreducible monic polynomials $`P_i(x)`$.
###### Theorem 1
Let $`M\mathrm{GL}(n,)`$ be simple and let its characteristic polynomial be $`P(x)=_{i=1}^{\mathrm{}}P_i(x)`$, with $`P_i(x)`$ irreducible over $``$. Then, the symmetry group of $`M`$, $`𝒮(M)\mathrm{GL}(n,)`$, is a finitely generated Abelian group of the form
$$𝒮(M)=T\times ^r$$
(16)
where $`T`$ is a finite Abelian group of even order, with at most $`\mathrm{}`$ generators. Furthermore, if the irreducible component $`P_i(x)`$ has $`n_1^{(i)}`$ real roots and $`n_2^{(i)}`$ pairs of complex conjugate roots, the rank $`r`$ of the free Abelian group in $`(\text{16})`$ is given by
$$r=\underset{i=1}{\overset{\mathrm{}}{}}(n_1^{(i)}+n_2^{(i)}1).$$
(17)
Proof: Since $`M`$ is simple, the degree of its minimal polynomial is $`n`$ and Fact 5 tells us that $`\mathrm{cent}_{\mathrm{Mat}(n,)}=[M]`$. As $`P(x)`$ has no repeated factors (so that $`[M]`$ contains no radicals), we get, by \[18, Thm. III.4\],
$$[M][\alpha _1]\mathrm{}[\alpha _{\mathrm{}}]$$
(18)
where $`\alpha _i`$ is any root<sup>6</sup><sup>6</sup>6Note that the $`\alpha _i`$ have pairwise different minimal polynomials by assumption, but that $`[\alpha _i][\alpha _j]`$ for $`ij`$ is still possible. of $`P_i(x)`$, for $`1i\mathrm{}`$.
Under the assumptions made, each $`[\alpha _i]=(\alpha _i)`$ is a field. Since a $`\mathrm{GL}(n,)`$-matrix in $`[M]`$ will correspond to a unit in each of the $`(\alpha _i)`$, we can now apply Lemma 1 to each component, giving (16) as the direct sum of $`\mathrm{}`$ unit groups. The rank in (17) follows now from Proposition 1.
The torsion part $`T`$ is a finitely generated Abelian group, with (at most) one generator per irreducible component of $`P(x)`$, of which there are $`\mathrm{}`$. Clearly, $`𝒮(M)`$ always contains the elements $`\pm \mathrm{𝟏}`$, so $`\{\pm 1\}C_2`$ is a subgroup of $`T`$. The order of $`T`$ is then divisible by $`2`$, hence even. $`\mathrm{}`$
Although Theorem 1 does not give the general answer to the question for the symmetry group $`𝒮(M)`$, it certainly gives the generic answer, because the property of $`M`$ having simple spectrum is generic. But what about the remaining cases? Without further elaborating on this, let us summarize a few aspects and otherwise refer to the literature for a summary of methods to actually determine the precise centralizer.
If $`M\mathrm{GL}(n,)`$ is semi-simple, but not simple, its characteristic polynomial contains a square, and whether or not $`𝒮(M)`$ is still Abelian (and then of the above form) depends on whether or not the minimal polynomial of $`M`$ has degree $`n`$, see Fact 5. Note, in particular, that the following situation can emerge. If $`P(x)`$ has a repeated factor, but the corresponding matrix $`M`$ is a block matrix, then the two blocks giving the same factor of $`P(x)`$ can still be inequivalent, if the corresponding class number is larger than one which equals the number of different matrix classes, see Fact 3.
If $`M`$ is not even semi-simple, things get even more involved. We can still have Abelian symmetry groups, e.g. if $`M`$ is a Jordan block such as $`\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)`$, compare the results of \[10, Sec. 2.1.2\] on parabolic automorphisms of $`𝕋^2`$. Clearly, this also follows from Fact 5: if $`M\mathrm{GL}(n,)`$ is conjugate to a single Jordan block, its minimal polynomial has degree $`n`$ and all $`\mathrm{GL}(n,)`$-matrices which commute with $`M`$ are in $`[M]`$. This remains true if such a block occurs in a matrix that otherwise has simple spectrum disjoint from $`1`$.
###### Corollary 3
Let $`M\mathrm{GL}(n,)`$. If the minimal polynomial of $`M`$ has degree $`n`$, then $`𝒮(M)\mathrm{GL}(n,)`$ is Abelian. $`\mathrm{}`$
In a wider setting for symmetries, a stronger statement can be formulated, see Proposition 5 below and the comments following it.
The general classification, however, and the non-Abelian cases in particular, gets increasingly difficult with growing $`n`$ and has been completed only for small $`n`$, see and references given there. Nevertheless, for any given $`M`$, the centralizer can be determined explicitly by means of various algorithmic program packages.
Let us, at the end of this part and before we illustrate some of the above results by further examples, give a particular case of one matrix written as a polynomial of another.
###### Fact 7
Let $`K`$ be a field and $`M\mathrm{GL}(n,K)`$ be an invertible matrix with characteristic polynomial $`P(x)=_{\mathrm{}=0}^na_{\mathrm{}}x^{\mathrm{}}`$, where $`a_n=1`$ and $`a_00`$. Then, the inverse matrix is given by
$$M^1=\frac{1}{a_0}\underset{\mathrm{}=0}{\overset{n1}{}}a_{\mathrm{}+1}M^{\mathrm{}}.$$
Proof: Observe that $`P(M)=0`$ from the Cayley-Hamilton Theorem. The verification of $`M^1M=\mathrm{𝟏}`$ is then a straight-forward calculation. $`\mathrm{}`$
Before we discuss the connection of our approach to quasicrystallography in a separate Section, let us illustrate Theorem 1 with a recent example of a 4D cat map taken from \[31, Eq. 3.21\], namely
$$M=\left(\begin{array}{cccc}0& 0& 1& 0\\ 0& 0& 0& 1\\ 1& 0& 2& 1\\ 0& 1& 1& 2\end{array}\right).$$
(19)
The characteristic polynomial is $`P(x)=x^44x^3+5x^24x+1`$ which is reducible over $``$ and splits as
$$P(x)=P_1(x)P_2(x)=(x^23x+1)(x^2x+1)$$
into $``$-irreducible polynomials. Since $`P_1`$ has two real and $`P_2`$ one pair of complex conjugate roots, Theorem 1 gives $`𝒮(M)T^{}\times `$, with $`T^{}`$ a subgroup of $`T=C_2\times C_6`$. Also, since no root of $`P_1`$ is a fundamental unit of the corresponding maximal order (which is $`[\tau ]`$ with $`\tau =(1+\sqrt{5})/2`$), the generator of the infinite cyclic group in $`𝒮(M)`$ could still differ from $`M`$.
To determine the details, one easily checks that the most general matrix to commute with $`M`$ is
$$G=\left(\begin{array}{cccc}a& b& c& d\\ b& a& d& c\\ c& d& a+2c+d& a+c+2d\\ d& c& a+c+2d& a+2c+d\end{array}\right).$$
A necessary condition for $`G`$ to be in $`\mathrm{GL}(4,)`$ is then $`a,b,c,d`$. This allows to exclude the existence of a root of $`M`$ in $`𝒮(M)`$, and also no element of third order is possible. So we obtain
$$𝒮(M)=C_2\times C_2\times M.$$
We will revisit this example below in the context of reversibility.
### 2.2 Three examples from planar quasicrystallography
Planar tilings with 8-, 10- and 12-fold symmetry play an important role in the description of so-called quasicrystalline T-phases, see for background material. They are of interest also in the present context because hyperbolic toral automorphisms show up through their inflation symmetry.
For the 8-fold case, consider the polynomial
$$P(x)=x^4+1$$
(20)
which has $`\xi `$, $`\xi ^3`$, $`\xi ^5`$, and $`\xi ^7`$ as roots, $`\xi =e^{2\pi i/8}`$, which are primitive. So, $`P(x)`$ is irreducible over $``$, and Lemma 1 and Proposition 1 apply. In fact, $`(\xi )`$ here is a cyclotomic field with class number one, maximal order $`[\xi ]`$ and unit group $`[\xi ]^\times C_8\times `$.
If we denote the actual matrices that represent the generators for the groups $`C_8`$ and $``$ by $`M`$ and $`G`$, respectively, it is natural to take the companion matrix of $`P(x)`$ for $`M`$ and to choose $`G`$ accordingly, resulting in
$$M=\left(\begin{array}{cccc}0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\\ 1& 0& 0& 0\end{array}\right),G=\left(\begin{array}{cccc}1& 1& 0& 1\\ 1& 1& 1& 0\\ 0& 1& 1& 1\\ 1& 0& 1& 1\end{array}\right).$$
(21)
By construction, $`M`$ is a matrix of order 8. So, $`𝒮(M)=M,GC_8\times `$. What is more, anticipating the next Section, $`M`$ turns out to be reversible with the matrix $`R`$ of (10) as reversing symmetry. Using Fact 1 and observing $`[G,R]=0`$, this means that $`(M)=M,G,RD_8\times `$. This, together with two similar examples, is summarized in Table 1.
Note that the case of 12-fold symmetry is more complicated because the fundamental unit in $`[\xi ]`$, for $`\xi =e^{2\pi i/12}`$, is the square root of $`(2+\sqrt{3})\xi `$, and hence not a simple homothety. This means that the representing matrix does not commute with $`R`$. The reversing symmetry group of this case, $`(C_{12}\times )\times _sC_2`$, does contain a subgroup of the form $`D_{12}\times `$ though – it is generated by $`M`$, $`G^{}=M^1G^2`$ and $`R`$, where $`G^{}`$ corresponds to the non-fundamental unit $`2+\sqrt{3}`$.
### 2.3 Reversibility
The examples of Section 2.2 were reversible, i.e. they fulfilled $`GMG^1=M^1`$ for some $`G\mathrm{GL}(4,)`$, in particular for the involution $`R\mathrm{GL}(4,)`$ of (10). It is easy to check that this is also true for $`M`$ of (19). However, as we will see below, neither of the examples of (15) are reversible in $`\mathrm{GL}(3,)`$.
In this Section, we are concerned with determining when reversibility can occur in $`\mathrm{GL}(n,)`$, and what we can say about the nature of the reversing symmetry $`G\mathrm{GL}(n,)`$, e.g. whether it can be taken to be an involution so that, by Fact 1, the reversing symmetry group $`(M)\mathrm{GL}(n,)`$ is a semi-direct product. Note that if $`M\mathrm{GL}(n,)`$ is reversible, it has been shown that there always exists an involutory reversing symmetry \[36, Thm. 2.1\]. Already for $`\mathrm{GL}(2,)`$, this is no longer true : the matrix $`M=\left(\begin{array}{cc}5& 7\\ 7& 10\end{array}\right)`$ is reversible in $`\mathrm{GL}(2,)`$ with the reversing symmetry $`G=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ of order 4, but with no involutory reversing symmetry in $`\mathrm{GL}(2,)`$.
While reversibility was still a frequent phenomenon in $`\mathrm{GL}(2,)`$, see , it becomes increasingly restrictive with growing $`n`$. To see this, let us consider necessary conditions for reversibility. If $`M\mathrm{GL}(n,)`$ is reversible, then $`M^1=GMG^1`$ for some matrix $`G`$, and $`M`$ and $`M^1`$ must have the same characteristic polynomial, $`P(x)`$. On the other hand, the spectrum of $`M`$ must then be self-reciprocal, i.e. with $`\lambda `$ also $`1/\lambda `$ must be an eigenvalue, with matching multiplicities. Recall that $`P(x)=_{i=1}^n(x\lambda _i)`$ and observe that
$$\underset{i=1}{\overset{n}{}}\left(x\frac{1}{\lambda _i}\right)=\frac{(1)^nx^n}{det(M)}\underset{i=1}{\overset{n}{}}\left(\frac{1}{x}\lambda _i\right).$$
(22)
But by assumption, $`_{i=1}^n(x\lambda _i)=_{i=1}^n(x\frac{1}{\lambda _i})`$, so we arrive at
###### Proposition 2
A necessary condition for the matrix $`M\mathrm{GL}(n,)`$ to be reversible is $`\mathrm{spec}(M)=\mathrm{spec}(M^1)`$ and thus the equation
$$P(x)=\frac{(1)^nx^n}{det(M)}P(1/x)$$
(23)
which we call the self-reciprocity of $`P(x)`$. $`\mathrm{}`$
One immediate consequence is the following.
###### Corollary 4
If $`M\mathrm{GL}(n,)`$ is reversible, $`det(M)=\pm 1`$. If, in addition, the multiplicity of the eigenvalue $`\lambda =1`$ is even $`(`$allowing for multiplicity $`0`$ if $`1`$ is not an eigenvalue of $`M`$$`)`$, then $`det(M)=1`$.
Proof: Observe that the reversibility of $`M`$ implies $`G=MGM`$. Taking determinants gives the first assertion because neither $`M`$ nor $`G`$ is singular.
Next, note that $`\lambda =\pm 1`$ are the only complex numbers with $`\lambda =1/\lambda `$. All other eigenvalues come in reciprocal pairs. Since the determinant is the product over all eigenvalues, the second statement follows. $`\mathrm{}`$
It might be instructive to reformulate Proposition 2 and Corollary 4 in terms of elementary symmetric polynomials. Let $`S_k`$, $`k=0,1,\mathrm{},n`$, denote the $`k`$-th elementary symmetric polynomial in $`n`$ indeterminates. The $`S_k`$ are given by $`S_01`$ and
$$S_k(x_1,\mathrm{},x_n)=\underset{i_1<\mathrm{}<i_k}{}x_{i_1}\mathrm{}x_{i_k}.$$
(24)
They are algebraically independent over $``$ and have the generating function
$$\underset{k=0}{\overset{n}{}}S_k(x_1,\mathrm{},x_n)t^k=\underset{i=1}{\overset{n}{}}(1+x_it)$$
(25)
where $`t`$ is another indeterminate.
Now observe that, for a characteristic polynomial $`P(x)`$, we have
$$P(x)=x^n+\underset{k=1}{\overset{n}{}}(1)^kS_k(\lambda _1,\mathrm{},\lambda _n)x^{nk}.$$
Consequently,
$$x^nP(1/x)=\mathrm{\hspace{0.33em}1}+\underset{k=1}{\overset{n}{}}(1)^kS_k(\lambda _1,\mathrm{},\lambda _n)x^k$$
and a comparison with Proposition 2 reveals that
$$S_k(\lambda _1,\mathrm{},\lambda _n)=det(M)S_{nk}(\lambda _1,\mathrm{},\lambda _n)$$
(26)
for all $`0kn`$. For $`k=0`$, this is just the statement that $`det(M)=\pm 1`$. Note that the elementary symmetric polynomials, when evaluated at the roots of $`P(x)`$, reproduce (up to a sign) the entries of the last row of the left companion matrix (8).
Turning now to the reversibility of matrices in $`\mathrm{GL}(n,)`$, we first observe that, generically, the reversible cases can only occur when $`n`$ is even and $`det(M)=+1`$:
###### Proposition 3
Consider $`M\mathrm{GL}(n,)`$ and let $`P(x)`$ be the characteristic polynomial of $`M`$. If $`n>1`$ is odd or $`det(M)=1`$ we have:
* if $`M`$ is reversible in $`\mathrm{GL}(n,)`$, $`P(x)`$ is reducible over $``$, and the spectrum of $`M`$ contains $`1`$ or $`1`$;
* if $`P(x)`$ is irreducible over $``$, $`M`$ cannot be reversible in $`\mathrm{GL}(n,)`$.
Proof: If $`M`$ is reversible, $`\lambda \mathrm{spec}(M)`$ implies $`1/\lambda \mathrm{spec}(M)`$, so the eigenvalues are either $`\pm 1`$ or have to come in pairs, $`\lambda 1/\lambda `$. If $`n`$ is odd, we must have at least one eigenvalue that is $`\pm 1`$, and that gives a factor $`(x1)`$ in $`P(x)`$. On the other hand, if $`det(M)=1`$, we must have at least one eigenvalue that is $`1`$ which gives a factor $`(x+1)`$ in $`P(x)`$. In both cases, one notes that whenever $`\pm 1`$ is a zero of a polynomial over $``$, factoring out $`(x1)`$ can be done over $``$.
Conversely, if $`P(x)`$ is irreducible over $``$, $`\mathrm{spec}(M)`$ cannot contain an eigenvalue of the form $`\pm 1`$, and $`n`$ odd or $`det(M)=1`$ is then incompatible with $`M`$ being reversible. $`\mathrm{}`$
It is clear from this that reversibility is rather restrictive. If $`P(x)`$ splits into irreducible components $`P_i(x)`$, then each is subject to the constraints described above, or has to be matched with its reciprocal partner polynomial – if that would be an integer polynomial at all. In particular, if $`P(x)`$ is reducible but contains an isolated irreducible factor of odd order $`3`$, or of even order with constant term $`1`$, reversibility of $`M`$ is ruled out. For example, this confirms that $`M_1`$ and $`M_2`$ of (15) are not reversible (in fact, they are not even reversible in $`\mathrm{GL}(3,)`$).
The key problem in deciding upon similarity of $`M`$ and $`M^1`$ in $`\mathrm{GL}(n,)`$ is that $``$ is not a field. But it is clear that the corresponding similarity within $`\mathrm{GL}(n,)`$ (the matrix entries now belonging to the field of rationals) is both a necessary condition and a much easier problem. It would not help to further extend $``$ to $``$ due to the following result, see \[23, Cor. XIV.2.3\].
###### Fact 8
A matrix $`M\mathrm{GL}(n,)`$ is similar to $`M^1`$ within the group $`\mathrm{GL}(n,)`$ if and only if this is already the case in $`\mathrm{GL}(n,)`$. $`\mathrm{}`$
In the light of this, let us first recall some facts about normal forms over $``$, where similarity is (in theory) a decideable problem. The normal form of a matrix $`M`$ is based on its polynomial invariants, or invariants for short, see \[23, Sec. XIV.2\]. They are often also called the invariant factors of $`M`$ (or, more explicitly, of $`(x\mathrm{𝟏}M)`$), compare \[3, Def. 4.4.6\], meaning certain polynomials that derive from the matrix $`(x\mathrm{𝟏}M)`$, see below. The following result is a direct consequence of \[23, Thm. XIV.2.6\] or \[3, Thm. 5.3.3\].
###### Fact 9
Two matrices in $`\mathrm{Mat}(n,)`$ are similar in $`\mathrm{GL}(n,)`$ if and only if they have the same polynomial invariants. In particular, this applies to $`M`$ and $`M^1`$ for any $`M\mathrm{GL}(n,)`$. $`\mathrm{}`$
Let us briefly recall how the polynomial invariants $`q_1,\mathrm{},q_r`$ of a matrix $`M\mathrm{Mat}(n,)`$ can be found, where $`rn`$ is a uniquely determined integer that depends on $`M`$. We formulate this for integer matrices, but it applies, with little change, also to rational ones. Set $`p_0=1`$ and let $`p_k`$ (for $`1kn`$) be the greatest common divisor of all minors of $`(x\mathrm{𝟏}M)`$ of order $`k`$, so that $`p_k`$ clearly divides $`p_{k+1}`$, and $`p_n=P(x)=det(x\mathrm{𝟏}M)`$. Let $`\mathrm{}`$ denote the largest integer $`k`$ for which $`p_k=1`$ and define $`q_i=p_{\mathrm{}+i}/p_{\mathrm{}+i1}`$, where $`1ir=n\mathrm{}`$. These polynomials over $``$ are the polynomial invariants of $`M`$ and satisfy the following divisibility property:
$$q_i|q_{i+1}.$$
(27)
The prime factors of $`q_i`$ over $``$, taken with their multiplicity, are called its elementary divisors. The product of the invariant factors of $`M`$ (equivalently, the product of all their elementary divisors) gives the characteristic polynomial $`P(x)`$ of $`M`$. Furthermore, the minimum polynomial $`Q(x)`$ of $`M`$ is given by $`q_r`$, or, equivalently, by the characteristic polynomial $`P(x)`$ divided by $`p_{n1}`$.
Note that a systematic way to find the invariant factors of $`M`$ is to bring $`(x\mathrm{𝟏}M)`$, seen as a matrix over the principal ideal domain $`[x]`$, into its so-called Smith normal form, see \[3, Ch. 5.3\] for details. This is a diagonal matrix in of the form $`\mathrm{diag}(1,\mathrm{},1,q_1(x),\mathrm{},q_r(x))`$. For large $`n`$, calculating this form can be a computationally difficult exercise; for small $`n`$, the Smith normal form can be found from algebraic program packages. Nevertheless, significantly, the invariant factors completely determine the Frobenius normal form of the matrix $`M`$:
###### Fact 10
Let $`M\mathrm{Mat}(n,)`$ have polynomial invariants $`q_1,\mathrm{},q_r`$ of degrees $`n_1,\mathrm{},n_r`$, with $`n_1+\mathrm{}+n_r=n`$. Then $`M`$ is similar, in $`\mathrm{GL}(n,)`$, to a block diagonal matrix $`[B_1,\mathrm{},B_r]`$ where $`B_i`$ is the $`n_i\times n_i`$ left companion matrix of the polynomial $`q_i`$. $`\mathrm{}`$
The existence of a block diagonal matrix similar to $`M`$ is equivalent to the statement that $`M`$ leaves invariant a set of (cyclic) subspaces of $`^n`$ with respective dimensions $`n_1,\mathrm{},n_r`$, see \[23, Thm. XIV.2.1\] for details. One can actually give more refined normal forms by using the elementary divisors of each invariant to replace the diagonal blocks $`B_i`$ with subblock decompositions based upon the elementary divisors and their multiplicities. Combining Fact 9 and Fact 10, matrices with the same polynomial invariants can both be brought to the same normal form and thus are similar.
The normal form of Fact 10 highlights the left companion matrices $`B_i`$. For what follows, we are interested in the reversibility of such matrices. In this respect, let $`M^{(\mathrm{})}`$ and $`M^{(r)}`$ be the left and right companion matrices corresponding to a polynomial $`P(x)`$. Suppose $`P(x)`$ conforms to the reciprocity condition (23) of Proposition 2. Then one can check that $`M^{(r)}`$ is the inverse of $`M^{(\mathrm{})}`$. But we already know from (9) that $`M^{(r)}=RM^{(\mathrm{})}R^1`$ where $`R=R^1`$ is the involution from (10). Combining this with the normal form above, we obtain:
###### Theorem 2
Let $`M\mathrm{GL}(n,)`$. Then, $`M`$ is reversible in $`\mathrm{GL}(n,)`$ if and only if each of the polynomial invariants of $`M`$ satisfies the reciprocity condition $`(\text{23})`$ separately. In this situation, the reversing symmetry can be chosen to be an involution.
Proof: By Fact 10, $`M=SDS^1`$ where $`S\mathrm{GL}(n,)`$ and $`D\mathrm{GL}(n,)`$ is a block diagonal matrix of the form $`D=[B_1,\mathrm{},B_r]`$, where $`r1`$ and $`B_i\mathrm{GL}(n_i,)`$ is the left companion matrix corresponding to the invariant $`q_i`$ of degree $`n_i`$. It follows that $`M^1=SD^1S^1`$ where $`D^1=[B_1^1,\mathrm{},B_r^1]`$. Consequently, $`M`$ and $`M^1`$ are similar if and only if $`D`$ and $`D^1`$ are similar.
Suppose that each of the polynomial invariants satisfies the condition (23). Then, from the remark before Theorem 2, $`B_i^1`$ is the right companion matrix corresponding to $`q_i`$ and is similar to $`B_i`$ via the involution $`R_i\mathrm{GL}(n_i,)`$ which consists of $`1`$’s on its anti-diagonal, as in (10). It follows that $`D`$ is similar to $`D^1`$ via the block diagonal involution $`R:=[R_1,\mathrm{},R_r]`$, and so $`M`$ and $`M^1`$ are similar by the involution $`SRS^1`$.
On the other hand, suppose that $`M`$ is similar to $`M^1`$ in $`\mathrm{GL}(n,)`$. Hence, the corresponding block diagonal matrices $`D`$ and $`D^1`$ are also similar in $`\mathrm{GL}(n,)`$, via some element $`G`$. Now, to each block $`B_i`$ of $`D`$ corresponds an invariant vector subspace $`V_i`$ of $`^n`$ of dimension $`n_i`$. A subspace $`V_i`$ is thus either mapped by $`G`$ to itself (it is a symmetric subspace) or to another subspace $`V_j`$ of the same dimension. In the first case, $`B_i`$ must be conjugate to its inverse via the restriction of $`G`$ to $`V_i`$. This means that the characteristic polynomial of $`B_i`$, which is $`q_i`$, must satisfy the condition (23). On the other hand, if $`V_i`$ is mapped to $`V_j`$ by $`G`$ with $`n_i=n_j`$, it follows that the invariants $`q_i`$ and $`q_j`$ differ by at most a sign. They thus share the same eigenvalues and must each satisfy condition (23) on their own. $`\mathrm{}`$
If $`M`$ has only one non-trivial invariant, $`q_1(x)`$, it follows from the above discussion that its characteristic polynomial $`P(x)`$ coincides with its minimal polynomial $`Q(x)`$ and both equal $`q_1(x)`$ (so $`M`$ is cyclic from Fact 5). Conversely, $`M`$ cyclic means it has only one invariant. The previous Theorem now gives:
###### Corollary 5
If $`M\mathrm{GL}(n,)`$ has only one polynomial invariant, in particular if the characteristic polynomial $`P(x)`$ is irreducible over $``$, then $`M`$ is reversible in $`\mathrm{GL}(n,)`$ if and only if its characteristic polynomial $`P(x)`$ satisfies the reciprocity condition $`(\text{23})`$. $`\mathrm{}`$
Note that Theorem 2 and Corollary 5 do not extend to requiring, for reversible $`M`$, that the elementary divisors within an invariant polynomial satisfy (23). For example, any matrix in $`\mathrm{GL}(4,)`$ with one invariant polynomial $`q_1(x)=P(x)=(x^2x1)(x^2+x1)`$ is reversible in $`\mathrm{GL}(4,)`$ although the elementary divisors separately violate (23).
Let us give some illustrations of the use of Theorem 2 and Corollary 5 for small values of $`n`$. These results show that all $`M\mathrm{SL}(2,)`$ are reversible in $`\mathrm{GL}(2,)`$ because their invariant factors fall into one of the following cases:
1. $`q_1(x)=q_2(x)=(x\pm 1)`$, so $`r=2`$ and $`M=\mathrm{𝟏}`$;
2. $`q_1(x)=P(x)=x^2\mathrm{tr}(M)x+1`$, so $`r=1`$ and $`P(x)`$ is self-reciprocal.
Yet, we know from that $`\mathrm{SL}(2,)`$ matrices exist that are not reversible in $`\mathrm{GL}(2,)`$ (see also Section 3.3 below for further discussion). Furthermore, if $`M\mathrm{GL}(2,)`$ with $`detM=1`$, then it can only have one polynomial invariant, $`q_1(x)=P(x)=x^2\mathrm{tr}(M)x1`$. By Proposition 3 or Corollary 5, $`M`$ is reversible in $`\mathrm{GL}(n,)`$ if and only if $`M`$ has eigenvalues $`\lambda =\pm 1`$ and $`q_1(x)`$ factors into $`(x1)(x+1)`$, in which case $`M\mathrm{GL}(2,)`$ is an involution and reversible, with reversing symmetry as itself. This approach gives another way of retrieving some of the results of on the reversibility in $`\mathrm{GL}(2,)`$.
Turning to $`\mathrm{GL}(3,)`$, Proposition 3 implies that if $`M`$ is reversible then $`P(x)`$ must have a factor $`(x\pm 1)`$. In other words, $`(x\pm 1)^i`$, for some $`1i3`$, is an elementary divisor of $`P(x)`$. If $`M`$ has more than one polynomial invariant, then the divisibility property (27) implies that $`P(x)`$ completely decomposes into a product of 3 factors of the form $`(x\pm 1)`$. Generically, however, this will not happen and instead $`M`$ is reversible in $`\mathrm{GL}(3,)`$ if and only if it has only one invariant of the form $`q_1(x)=P(x)=(x\pm 1)(x^2(\mathrm{tr}(M)\pm 1)x+1)`$.
For $`\mathrm{GL}(4,)`$, elements with three or four polynomial invariants are diagonal matrices with $`+1`$’s or (an even number of) $`1`$’s on the diagonal. They are all reversible. Reversible elements with two invariants must have $`q_1(x)=(x\pm 1)`$ and $`q_2(x)=(x\pm 1)(x^2(\mathrm{tr}(M)\pm 2)x+1)`$, or $`q_1(x)=q_2(x)`$, a monic quadratic with constant term $`+1`$.
As $`n`$ increases, Theorem 2 can exclude many matrices from being reversible in $`\mathrm{GL}(n,)`$ because they are not reversible in $`\mathrm{GL}(n,)`$. In particular, we can ask for an example $`M`$ with more than one polynomial invariant where the characteristic polynomial $`P(x)`$ satisfies (23), yet $`M`$ is irreversible in $`\mathrm{GL}(n,)`$ because it violates Theorem 2. If we take $`n2`$ and even, the first possibility appears in $`\mathrm{GL}(8,)`$. For example, we can take a matrix with invariants $`q_1(x)=x^2x1`$ and $`q_2(x)=(x^2x1)(x^2+x1)^2`$. Both polynomials violate the reciprocity condition (23) but they compensate each other so that their product, the characteristic polynomial, does satisfy the condition. The Frobenius normal form with these invariants is the block diagonal matrix $`M=[B_1,B_2]`$,
$$M=\left(\begin{array}{cccccccc}0& 1& 0& 0& 0& 0& 0& 0\\ 1& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& \text{-}0& \text{-}1& \text{-}1& \text{-}4& \text{-}3& \text{-}4& \text{-}1\end{array}\right).$$
(28)
This matrix is not reversible in $`\mathrm{GL}(8,)`$ (although its square is, see Section 3.2 below).
The remaining problem is now to find possible reversibility in $`\mathrm{GL}(n,)`$ of unimodular matrices that are already reversible in $`\mathrm{GL}(n,)`$. For $`n=2`$, we were able to solve this problem using a special algebraic structure (the amalgamated free product) of $`\mathrm{PGL}(2,)`$. This structure is not available for $`n3`$.
Important examples of unimodular integer matrices which are reversible in $`\mathrm{GL}(n,)`$ are the symplectic matrices $`\mathrm{Sp}(2n,)\mathrm{SL}(2n,)`$. Recall that a symplectic matrix $`M\mathrm{Sp}(2n,)`$ satisfies $`M^tJM=J`$ where $`M^t`$ denotes the transpose of $`M`$ and $`J`$ is the $`2n\times 2n`$ integer block matrix
$$J=\left(\begin{array}{cc}\mathrm{𝟎}& \mathrm{𝟏}\\ \text{-}\mathrm{𝟏}& \mathrm{𝟎}\end{array}\right)$$
of order $`4`$. Since in general
$$M^tJM=JM^t=JM^1J^1,$$
(29)
it follows that $`M\mathrm{Sp}(2n,)`$ is reversible if and only if $`M`$ is similar to $`M^t`$ in $`\mathrm{GL}(2n,)`$. But any invertible square matrix with entries in a field $`F`$ is similar to its transpose in $`\mathrm{GL}(m,F)`$, see \[3, Prop. 5.3.7\] (but this need not be true e.g. in $`\mathrm{GL}(m,)`$). In particular, $`M\mathrm{Sp}(2n,)`$ is reversible in $`\mathrm{GL}(2n,)`$ and its invariant factors will all satisfy (23). Also, it is clear from (29) that if $`M`$ is symplectic and symmetric, then $`M^t=M=JM^1J^1`$ so that $`M`$ is actually reversible in $`\mathrm{GL}(2n,)`$.
Orthogonal integer matrices $`U\mathrm{GL}(n,)`$ which satisfy $`UU^t=\mathrm{𝟏}`$ are other examples of unimodular integer matrices which are always reversible in $`\mathrm{GL}(n,)`$ (since $`U^1=U^t`$ and, as before, $`U^t`$ and $`U`$ are similar in $`\mathrm{GL}(n,)`$).
We make some further remarks on the problem of deciding reversibility within $`\mathrm{GL}(n,)`$. Firstly, note that in the proof of Theorem 2 above, when we used the reversibility of the left companion matrices with characteristic polynomials satisfying (23), this reversibility was in $`\mathrm{GL}(n,)`$ itself. Furthermore, the reversing symmetry $`R`$ was an involution, so by Fact 1 we have the following result.
###### Theorem 3
For each integer polynomial $`P(x)`$ of degree $`n`$ that satisfies the necessary self-reciprocity condition $`(\text{23})`$ for reversibility, there is at least one reversible matrix class in $`\mathrm{GL}(n,)`$, represented by the left companion matrix $`M^{(\mathrm{})}`$ of $`P(x)`$, and we have $`(M^{(\mathrm{})})=𝒮(M^{(\mathrm{})})\times _sC_2`$. $`\mathrm{}`$
Secondly, by Fact 3, the number of representing matrix classes of an irreducible characteristic polynomial $`P(x)`$ equals the class number of the order $`[\alpha ]`$, with $`\alpha `$ any of the roots of $`P(x)`$. If this class number is one, there is only the class represented by the companion matrix. If the class number is two, one is the companion matrix class which we know to be reversible if the spectrum is self-reciprocal. But then, the other class must also be reversible because there is no further partner left.
Class numbers are widely studied in algebraic number theory, and one can find both extensive tables in books (e.g. see ) and also various program packages to calculate them, e.g. the program package KANT<sup>7</sup><sup>7</sup>7See http://www.math.TU-Berlin.de/$``$kant/..
Let us add another example, of rather different flavour, and look at an interesting class of algebraic integers, the so-called Salem numbers. They are the algebraic integers $`\alpha >1`$ with all conjugates $`\alpha ^{}`$ having modulus $`|\alpha ^{}|1`$ and with at least one conjugate on the unit circle, see \[35, Ch. III.3\] for details. Salem’s Theorem then says that their degree is always even, that $`\alpha `$ and $`1/\alpha `$ are the only real conjugates, and that all other conjugates are on the unit circle. In particular, a Salem number is a unit. Putting our above results to work, we get
###### Corollary 6
Each Salem number occurs as the eigenvalue of a reversible toral automorphism. If $`\alpha `$ is a Salem number of degree $`n=2m`$, and $`M`$ a corresponding $`\mathrm{GL}(n,)`$ matrix, then $`𝒮(M)C_2\times ^m`$. $`\mathrm{}`$
The polynomial $`P(x)=x^42x^32x^22x+1`$ provides one of the simplest examples. Its roots are $`\tau \pm \sqrt{\tau }`$ (both real) and $`(\tau 1)\pm i\sqrt{\tau 1}`$ (both on the unit circle), where $`\tau =\left(1+\sqrt{5}\right)/2`$ is the golden number. The corresponding companion matrix $`M`$ is not of finite order, and the reversing symmetry group is thus $`(M)(C_2\times ^2)\times _sC_2`$.
Let us come back to the general discussion and ask for the properties of reversing symmetries. Let $`G`$ be a reversing symmetry of $`M`$, so that $`GM=M^1G`$ and hence also $`GM^n=M^nG`$ for all $`n`$. If $`p(x)`$ is any polynomial, we then also get $`Gp(M)=p(M^1)G`$. Since $`G`$ is a reversing symmetry, $`G^2`$ is a symmetry. If we now assume that $`M\mathrm{GL}(n,)`$ has minimal polynomial of degree $`n`$, we get $`G^2[M]`$ from Fact 5, i.e. $`G^2=q(M)`$ for some $`q`$ with coefficients in $``$. Consequently, $`Gq(M)G^1=q(M^1)`$, but also $`Gq(M)G^1=G^2=q(M)`$, so that $`q(M)=q(M^1)`$.
If $`q(M)`$ is a monomial, i.e. $`q(M)=M^{\mathrm{}}`$ for some $`\mathrm{}`$, then $`M^2\mathrm{}=1`$ and hence $`G^4=1`$. This case is also discussed in \[16, Prop. 2(i)\]. It clearly extends to the situation that $`q(M^1)=\left(q(M)\right)^1`$, which is more general. Apart from this, we recall the following result from \[16, p. 21\] (its proof, which was only contained in the preprint version of , is a coset counting argument).
###### Fact 11
Let $`M`$ be of infinite order. If the factor group $`𝒮(M)/M`$ is finite, then any reversing symmetry $`G`$ of $`M`$ must be of finite order, and $`G^{2k}=\mathrm{𝟏}`$ for some integer $`k`$ that divides the order of the factor group. $`\mathrm{}`$
Let us only add that, in line with our above argument, one first obtains $`G^{2k}M`$ and hence $`G^{4k}=\mathrm{𝟏}`$. But $`M`$ by assumption, so it cannot have a subgroup of order 2, and thus already $`G^{2k}=\mathrm{𝟏}`$.
Fact 11 certainly applies to our scenario whenever $`M`$ is not of finite order, but $`𝒮(M)`$ Abelian and or rank 1. This type of result is helpful because it restricts the search for reversing symmetries to one among elements of finite order. It is certainly possible to extend the result to other cases, but in general one has to expect reversing symmetries of infinite order, in particular if the rank of $`𝒮(M)`$ is $`2`$. Even then some results are possible because it would be sufficient to know whether reversibility implied the existence of some reversing symmetries of finite order. However, this question is more involved and thus postponed.
## 3 Extensions and further directions
In this Section, we will summarize some additional aspects of our analysis, namely the extension of symmetries to affine mappings, the modifications needed to treat the related situation of the projective matrix group $`\mathrm{PGL}(n,)`$, and the extension of (reversing) symmetries from a group setting to that of (matrix) rings or semi-groups.
### 3.1 Extension to affine transformations
So far, we have mainly discussed linear transformations (w.r.t. the torus), but it is an interesting question what happens if one extends the search for (reversing) symmetries to the group of affine transformations. Since both arguments and results are the exact analogues of those for the case $`n=2`$ as derived in , we will be very brief here.
In Euclidean $`n`$-space, the group of affine transformations is the semi-direct product $`𝒢_a=^n\times _s\mathrm{GL}(n,)`$, with $`^n`$ being the normal subgroup. Elements are written as $`(t,M)`$ with $`t^n`$ and $`M\mathrm{GL}(n,)`$, and the product of two transformations is $`(t,M)(t^{},M^{})=(t+Mt^{},MM^{})`$. The neutral element is $`(0,\mathrm{𝟏})`$, and we have $`(t,M)^1=(M^1t,M^1)`$.
If we now observe that $`𝕋^n=^n/^n`$, it is immediately clear that the affine transformations of $`𝕋^n`$ form the group
$$𝒢_a^{𝕋^n}=𝕋^n\times _s\mathrm{GL}(n,)$$
(30)
which is still a semi-direct product. Here, $`𝕋^n`$ can be written as $`[0,1)^n`$ with addition mod. 1, and the product of transformations is modified accordingly.
If we now ask for an affine (reversing) symmetry of a matrix $`M`$ (now being identified with the element $`(0,M)𝒢_a^{𝕋^n}`$) we find
###### Proposition 4
The affine transformation $`(t,G)`$ is a $`(`$reversing$`)`$ symmetry of the toral automorphism $`(0,M)`$ if and only if
* $`G`$ is a $`(`$reversing$`)`$ symmetry of $`M`$ in $`\mathrm{GL}(n,)`$ and
* $`Mt=t`$ $`(\text{mod }1)`$.
Proof: We have $`(t,G)(0,M)=(t,GM)`$ and also $`(0,M^{\pm 1})(t,G)=(M^{\pm 1}t,M^{\pm 1}G)`$. But then, the statement follows from the uniqueness of factorization in semi-direct products. $`\mathrm{}`$
From the condition $`Mt=t`$ (mod 1) it is clear that we need not consider all translations in $`𝕋^n`$ but only those with rational components, which we denote as $`\mathrm{\Lambda }_{\mathrm{}}`$. For many concrete problems, it would actually be even more appropriate to restrict to discrete sublattices, e.g. to the so-called $`q`$-division points $`\mathrm{\Lambda }_q(C_q)^n`$ which consists of all rational points with denominator $`q`$. We will not follow this idea here, however.
From the above result, it is clear that we can get (reversing) $`k`$-symmetries (recall the definitions from Section 1.2). In fact, the equation $`M^kt=t`$ on the torus has $`a_k=|det(M^k\mathrm{𝟏})|`$ different solutions provided no eigenvalue of $`M^k`$ is 1. Clearly,
$$a_k=\underset{\mathrm{}|k}{}\mathrm{}c_{\mathrm{}}$$
(31)
where $`c_{\mathrm{}}`$ counts the true orbits of length $`\mathrm{}`$, and the Möbius inversion formula gives
$$c_k=\frac{1}{k}\underset{\mathrm{}|k}{}\mu (\frac{k}{\mathrm{}})a_{\mathrm{}}$$
(32)
with the Möbius function $`\mu (m)`$ \[13, p. 29\]. If $`c_k`$ is positive for some $`k`$, we get a $`k`$-symmetry (and, hence, eventually a reversing $`k`$-symmetry) of $`M`$. These numbers can easily be calculated explicitly, where a very natural tool is provided by the so-called dynamical or Artin-Mazur $`\zeta `$-functions . Here, the $`a_k`$’s can be extracted from the series expansion of the logarithm of the $`\zeta `$-function, while the $`c_k`$’s appear as exponents of the factors of the Euler product expansion of the $`\zeta `$-function itself.
### 3.2 The case of $`\mathrm{PGL}(n,)`$
Let us start by the observation that $`\mathrm{PGL}(n,)`$ can be described via quotienting w.r.t. $`\{\pm \mathbf{\hspace{0.17em}1}\}`$, i.e.
$$\mathrm{PGL}(n,)\mathrm{GL}(n,)/\{\pm \mathbf{\hspace{0.17em}1}\}.$$
In other words, rather than consider single matrices $`M`$, one has to consider pairs, $`[M]:=\{\pm M\}`$. Let us write $`𝒮[M]`$ for the new $`\mathrm{PGL}`$ case and keep the old notation for the $`\mathrm{GL}`$ situation treated above.
The modification needed for the symmetry analysis given above is then actually fairly trivial, as we always had $`\pm \mathrm{𝟏}`$ among them, and we can simply factor that out. So, we get
$$𝒮[M]𝒮(M)/\{\pm \mathrm{𝟏}\}.$$
(33)
The case of reversing symmetries, however, requires some care. Since we now calculate mod $`\pm \mathrm{𝟏}`$, a projective matrix $`[M]`$ can also be reversible through $`GMG^1=M^1`$. But if this happens, the square, $`M^2`$, is again reversible in the old sense. This mechanism can (and will) give rise to reversing 2-symmetries (recall Section 1.2 for the definition) in $`\mathrm{GL}(n,)`$. Whereas gave examples in $`\mathrm{GL}(2,)`$, Eq. (29) shows that skew-symmetric symplectic matrices satisfying $`M^t=M`$ are not reversible in $`\mathrm{GL}(2n,)`$ whereas their squares are reversible. Also, the example (28) is reversible in $`\mathrm{PGL}(8,)`$ since a $`G\mathrm{GL}(8,)`$ can be found, by direct calculation, that satisfies the relation $`GMG^1=M^1`$.
Let us finally check what happens in the extension to affine transformations. In complete analogy to the case $`n=2`$, see , one can show that the corresponding affine group is the semidirect product $`\mathrm{\Lambda }_2\times _s\mathrm{PGL}(n,)`$ with $`\mathrm{\Lambda }_2`$ the 2-division points. This really is the consequence of identifying $`x`$ with $`x`$ on $`𝕋^n`$, and $`\mathrm{\Lambda }_2`$ is the set of $`2^n`$ translations that satisfy the condition $`t=t\text{(mod 1)}`$.
Now, the above Proposition 4 applies to the case of $`\mathrm{PGL}`$-matrices in very much the same way, just the possible translations $`t`$ are restricted to the 2-division points.
### 3.3 Symmetries among general integer matrices
For most of this article, we have focused on matrices in $`\mathrm{GL}(n,)`$ and their symmetries within the same group. However, none of the proofs given above depends on that restriction, and one can indeed also treat the case that both $`M`$ and its symmetries are allowed to live in the larger set $`\mathrm{Mat}(n,)`$ which is no longer a group w.r.t. multiplication, but a ring. Nevertheless, we will continue to use the symbol $`𝒮(M)`$, now meaning
$$𝒮(M):=\{G\mathrm{Mat}(n,)[M,G]=0\}.$$
The most obvious extended symmetries which one gets in $`\mathrm{Mat}(n,)`$ are the integer multiples of the identity, but there really is a hierarchy of objects to look at, and it is most transparent if one phrases the situation for a matrix in $`\mathrm{Mat}(n,)`$ first:
###### Fact 12
Let $`M\mathrm{Mat}(n,)`$ and let its characteristic polynomial $`P(x)`$ be irreducible over $``$. Let $`\lambda `$ be any of its roots and $`K=(\lambda )`$ the corresponding algebraic number field. Then the following statements hold.
* $`\mathrm{cent}_{\mathrm{Mat}(n,)}(M)K`$.
* $`\mathrm{cent}_{\mathrm{GL}(n,)}(M)K^{}`$, where $`K^{}=K\{0\}`$.
* $`\mathrm{cent}_{\mathrm{Mat}(n,)}(M)𝒪`$, where $`𝒪`$ is an order in $`K`$.
* $`\mathrm{cent}_{\mathrm{GL}(n,)}(M)𝒪^\times `$. $`\mathrm{}`$
The proof is a slight variation of what we did for Fact 5 and Lemma 1, and need not be spelled out again. It is important to note that the order $`𝒪`$ appearing here, as mentioned before, in general is not the maximal order of $`K`$, though it contains $`[\lambda ]`$. The following is now an immediate consequence.
###### Proposition 5
Let $`M`$ be an integer matrix with irreducible characteristic polynomial $`P(x)`$. Let $`\lambda `$ be a root of $`P(x)`$, and let $`𝒪_{\mathrm{max}}`$ be the maximal order in $`(\lambda )`$. Then, $`𝒮(M)`$ is both a $``$-module and a ring, and isomorphic to an order $`𝒪`$ that satisfies $`[\lambda ]𝒪𝒪_{\mathrm{max}}`$. $`\mathrm{}`$
Note that there is now also a natural extension to the case of simple matrices $`M`$, compare Theorem 1 and Eq. (18), but we omit further details here. Also, from Fact 5 it is clear that $`𝒮(M)`$ is Abelian if and only if the minimal polynomial of $`M`$ has degree $`n`$.
As to reversibility, this new point of view requires some thought. By Corollary 4, $`M`$ reversible implies $`det(M)=\pm 1`$, so reversible integer matrices are restricted to $`\mathrm{GL}(n,)`$. It would then not be unnatural to also insist on the existence of at least one unimodular matrix $`G`$ with $`M^1=GMG^1`$, and reversibility is basically as above, except that, if we enlarge the symmetries of $`M`$ from subgroups of $`\mathrm{GL}(n,)`$ to subrings of $`\mathrm{Mat}(n,)`$, reversing symmetries get enlarged accordingly. Note, however, that the ring structure is lost: the sum of a symmetry and a reversing symmetry is not a meaningful operation in this context. Together, they only form a monoid, i.e. a semi-group with unit element.
To go one step further, one could then also rewrite the reversibility condition as
$$G=MGM$$
and only demand that $`G`$ is non-singular, to avoid pathologies with projections to subspaces and to keep the statement of Corollary 4. Note that this is a slightly weaker form of reversibility, as one does not assume that $`G^1`$ is a meaningful mapping in this context. In particular, it is clear that there is then no need any more to restrict $`G`$ to unimodular integer matrices, so that now $`G\mathrm{GL}(n,)`$. This will, in general, lead to new cases of (weak) reversibility, as is apparent from the explicit constructions in .
To give a concrete example, consider the matrices
$$M=\left(\begin{array}{cc}4& 9\\ 7& 16\end{array}\right)\text{and}G=\left(\begin{array}{cc}3& \mathrm{\hspace{0.17em}0}\\ 4& 3\end{array}\right).$$
(34)
Then $`M`$ is the matrix from \[10, Ex. 2\] that was shown to be irreversible in $`\mathrm{GL}(2,)`$. In fact, the automorphism of $`𝕋^2`$ induced by $`M`$ is irreversible even in the larger group of homeomorphisms of the 2-torus, see and \[1, p. 9\] for details on the connection between general and linear homeomorphisms. Nevertheless, one can check that $`G=MGM`$, where $`det(G)=9`$. In other words, $`M`$ is reversible in $`\mathrm{GL}(2,)`$, as are all elements of $`\mathrm{SL}(2,)`$ by our previous discussion following Corollary 5.
Note that $`G`$ does not induce a homeomorphism of the 2-torus because $`G^1`$ is not an integer matrix. However, $`G`$ does induce an automorphism on any lattice of the torus of the form
$$\mathrm{\Lambda }_q:=\{(\frac{m}{q},\frac{n}{q})^t0m,n<q\}$$
for which $`det(G)0`$ (mod $`q`$). This is relevant as a recent study shows: the quantum map which corresponds to $`M`$ (which, in turn, corresponds to the action of $`M`$ on a (Wigner) lattice of the torus) showed an eigenvalue statistics according to the circular orthogonal ensemble (COE) rather than the unitary one (CUE). So, even though $`M`$ does not have a reversing symmetry in the sense of Section 2.3, the presence of “pseudo-symmetries” such as $`G`$ still leave their mark!
### Acknowledgements
M. B. would like to thank Peter A. B. Pleasants and Alfred Weiss for several clarifying discussions, Gabriele Nebe for helpful advice on the literature, and Wilhelm Plesken for suggesting a number of improvements. This work was supported by the German and Australian Research Councils (DFG and ARC).
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# Measurement of the Relative Branching Fraction of Υ(4𝑆) to Charged and Neutral 𝐵-Meson Pairs
## Abstract
We analyze $`9.7\times 10^6B\overline{B}`$ pairs recorded with the CLEO detector to determine the production ratio of charged to neutral $`B`$-meson pairs produced at the $`\mathrm{{\rm Y}}`$(4S) resonance. We measure the rates for $`B^0J/\psi K^{()0}`$ and $`B^+J/\psi K^{()+}`$ decays and use the world-average $`B`$-meson lifetime ratio to extract the relative widths $`\frac{f_+}{f_{00}}=\frac{\mathrm{\Gamma }(\mathrm{{\rm Y}}(4\mathrm{S})\mathrm{B}^+\mathrm{B}^{})}{\mathrm{\Gamma }(\mathrm{{\rm Y}}(4\mathrm{S})\mathrm{B}^0\overline{\mathrm{B}}^0)}=1.04\pm 0.07(stat)\pm 0.04(syst)`$. With the assumption that $`f_++f_{00}=1`$, we obtain $`f_{00}=0.49\pm 0.02(stat)\pm 0.01(syst)`$ and $`f_+=0.51\pm 0.02(stat)\pm 0.01(syst)`$. This production ratio and its uncertainty apply to all exclusive $`B`$-meson branching fractions measured at the $`\mathrm{{\rm Y}}`$(4S) resonance.
preprint: CLNS 00/1670 CLEO 00-7
J. P. Alexander,<sup>1</sup> R. Baker,<sup>1</sup> C. Bebek,<sup>1</sup> B. E. Berger,<sup>1</sup> K. Berkelman,<sup>1</sup> F. Blanc,<sup>1</sup> V. Boisvert,<sup>1</sup> D. G. Cassel,<sup>1</sup> M. Dickson,<sup>1</sup> P. S. Drell,<sup>1</sup> K. M. Ecklund,<sup>1</sup> R. Ehrlich,<sup>1</sup> A. D. Foland,<sup>1</sup> P. Gaidarev,<sup>1</sup> L. Gibbons,<sup>1</sup> B. Gittelman,<sup>1</sup> S. W. Gray,<sup>1</sup> D. L. Hartill,<sup>1</sup> B. K. Heltsley,<sup>1</sup> P. I. Hopman,<sup>1</sup> C. D. Jones,<sup>1</sup> D. L. Kreinick,<sup>1</sup> M. Lohner,<sup>1</sup> A. Magerkurth,<sup>1</sup> T. O. Meyer,<sup>1</sup> N. B. Mistry,<sup>1</sup> E. Nordberg,<sup>1</sup> J. R. Patterson,<sup>1</sup> D. Peterson,<sup>1</sup> D. Riley,<sup>1</sup> J. G. Thayer,<sup>1</sup> P. G. Thies,<sup>1</sup> B. Valant-Spaight,<sup>1</sup> A. Warburton,<sup>1</sup> P. Avery,<sup>2</sup> C. Prescott,<sup>2</sup> A. I. Rubiera,<sup>2</sup> J. Yelton,<sup>2</sup> J. Zheng,<sup>2</sup> G. Brandenburg,<sup>3</sup> A. Ershov,<sup>3</sup> Y. S. Gao,<sup>3</sup> D. Y.-J. Kim,<sup>3</sup> R. Wilson,<sup>3</sup> T. E. Browder,<sup>4</sup> Y. Li,<sup>4</sup> J. L. Rodriguez,<sup>4</sup> H. Yamamoto,<sup>4</sup> T. Bergfeld,<sup>5</sup> B. I. Eisenstein,<sup>5</sup> J. Ernst,<sup>5</sup> G. E. Gladding,<sup>5</sup> G. D. Gollin,<sup>5</sup> R. M. Hans,<sup>5</sup> E. Johnson,<sup>5</sup> I. Karliner,<sup>5</sup> M. A. Marsh,<sup>5</sup> M. Palmer,<sup>5</sup> C. Plager,<sup>5</sup> C. Sedlack,<sup>5</sup> M. Selen,<sup>5</sup> J. J. Thaler,<sup>5</sup> J. Williams,<sup>5</sup> K. W. Edwards,<sup>6</sup> R. Janicek,<sup>7</sup> P. M. Patel,<sup>7</sup> A. J. Sadoff,<sup>8</sup> R. Ammar,<sup>9</sup> A. Bean,<sup>9</sup> D. Besson,<sup>9</sup> R. Davis,<sup>9</sup> N. Kwak,<sup>9</sup> X. Zhao,<sup>9</sup> S. Anderson,<sup>10</sup> V. V. Frolov,<sup>10</sup> Y. Kubota,<sup>10</sup> S. J. Lee,<sup>10</sup> R. Mahapatra,<sup>10</sup> J. J. O’Neill,<sup>10</sup> R. Poling,<sup>10</sup> T. Riehle,<sup>10</sup> A. Smith,<sup>10</sup> J. Urheim,<sup>10</sup> S. Ahmed,<sup>11</sup> M. S. Alam,<sup>11</sup> S. B. Athar,<sup>11</sup> L. Jian,<sup>11</sup> L. Ling,<sup>11</sup> A. H. Mahmood,<sup>11,</sup><sup>*</sup><sup>*</sup>*Permanent address: University of Texas - Pan American, Edinburg, TX 78539. M. Saleem,<sup>11</sup> S. Timm,<sup>11</sup> F. Wappler,<sup>11</sup> A. Anastassov,<sup>12</sup> J. E. Duboscq,<sup>12</sup> E. Eckhart C. Gwon,<sup>12</sup> T. Hart,<sup>12</sup> K. Honscheid,<sup>12</sup> D. Hufnagel,<sup>12</sup> H. Kagan,<sup>12</sup> R. Kass,<sup>12</sup> T. K. Pedlar,<sup>12</sup> H. Schwarthoff,<sup>12</sup> J. B. Thayer,<sup>12</sup> E. von Toerne,<sup>12</sup> M. M. Zoeller,<sup>12</sup> S. J. Richichi,<sup>13</sup> H. Severini,<sup>13</sup> P. Skubic,<sup>13</sup> A. Undrus,<sup>13</sup> S. Chen,<sup>14</sup> J. Fast,<sup>14</sup> J. W. Hinson,<sup>14</sup> J. Lee,<sup>14</sup> N. Menon,<sup>14</sup> D. H. Miller,<sup>14</sup> E. I. Shibata,<sup>14</sup> I. P. J. Shipsey,<sup>14</sup> V. Pavlunin,<sup>14</sup> D. Cronin-Hennessy,<sup>15</sup> Y. Kwon,<sup>15,</sup>Permanent address: Yonsei University, Seoul 120-749, Korea. A.L. Lyon,<sup>15</sup> E. H. Thorndike,<sup>15</sup> C. P. Jessop,<sup>16</sup> H. Marsiske,<sup>16</sup> M. L. Perl,<sup>16</sup> V. Savinov,<sup>16</sup> D. Ugolini,<sup>16</sup> X. Zhou,<sup>16</sup> T. E. Coan,<sup>17</sup> V. Fadeyev,<sup>17</sup> Y. Maravin,<sup>17</sup> I. Narsky,<sup>17</sup> R. Stroynowski,<sup>17</sup> J. Ye,<sup>17</sup> T. Wlodek,<sup>17</sup> M. Artuso,<sup>18</sup> R. Ayad,<sup>18</sup> C. Boulahouache,<sup>18</sup> K. Bukin,<sup>18</sup> E. Dambasuren,<sup>18</sup> S. Karamov,<sup>18</sup> G. Majumder,<sup>18</sup> G. C. Moneti,<sup>18</sup> R. Mountain,<sup>18</sup> S. Schuh,<sup>18</sup> T. Skwarnicki,<sup>18</sup> S. Stone,<sup>18</sup> G. Viehhauser,<sup>18</sup> J.C. Wang,<sup>18</sup> A. Wolf,<sup>18</sup> J. Wu,<sup>18</sup> S. Kopp,<sup>19</sup> S. E. Csorna,<sup>20</sup> I. Danko,<sup>20</sup> K. W. McLean,<sup>20</sup> Sz. Márka,<sup>20</sup> Z. Xu,<sup>20</sup> R. Godang,<sup>21</sup> K. Kinoshita,<sup>21,</sup>Permanent address: University of Cincinnati, Cincinnati, OH 45221 I. C. Lai,<sup>21</sup> S. Schrenk,<sup>21</sup> G. Bonvicini,<sup>22</sup> D. Cinabro,<sup>22</sup> S. McGee,<sup>22</sup> L. P. Perera,<sup>22</sup> G. J. Zhou,<sup>22</sup> E. Lipeles,<sup>23</sup> M. Schmidtler,<sup>23</sup> A. Shapiro,<sup>23</sup> W. M. Sun,<sup>23</sup> A. J. Weinstein,<sup>23</sup> F. Würthwein,<sup>23,</sup><sup>§</sup><sup>§</sup>§Permanent address: Massachusetts Institute of Technology, Cambridge, MA 02139. D. E. Jaffe,<sup>24</sup> G. Masek,<sup>24</sup> H. P. Paar,<sup>24</sup> E. M. Potter,<sup>24</sup> S. Prell,<sup>24</sup> V. Sharma,<sup>24</sup> D. M. Asner,<sup>25</sup> A. Eppich,<sup>25</sup> T. S. Hill,<sup>25</sup> R. J. Morrison,<sup>25</sup> R. A. Briere,<sup>26</sup> B. H. Behrens,<sup>27</sup> W. T. Ford,<sup>27</sup> A. Gritsan,<sup>27</sup> J. Roy,<sup>27</sup> and J. G. Smith<sup>27</sup>
<sup>1</sup>Cornell University, Ithaca, New York 14853
<sup>2</sup>University of Florida, Gainesville, Florida 32611
<sup>3</sup>Harvard University, Cambridge, Massachusetts 02138
<sup>4</sup>University of Hawaii at Manoa, Honolulu, Hawaii 96822
<sup>5</sup>University of Illinois, Urbana-Champaign, Illinois 61801
<sup>6</sup>Carleton University, Ottawa, Ontario, Canada K1S 5B6
and the Institute of Particle Physics, Canada
<sup>7</sup>McGill University, Montréal, Québec, Canada H3A 2T8
and the Institute of Particle Physics, Canada
<sup>8</sup>Ithaca College, Ithaca, New York 14850
<sup>9</sup>University of Kansas, Lawrence, Kansas 66045
<sup>10</sup>University of Minnesota, Minneapolis, Minnesota 55455
<sup>11</sup>State University of New York at Albany, Albany, New York 12222
<sup>12</sup>Ohio State University, Columbus, Ohio 43210
<sup>13</sup>University of Oklahoma, Norman, Oklahoma 73019
<sup>14</sup>Purdue University, West Lafayette, Indiana 47907
<sup>15</sup>University of Rochester, Rochester, New York 14627
<sup>16</sup>Stanford Linear Accelerator Center, Stanford University, Stanford, California 94309
<sup>17</sup>Southern Methodist University, Dallas, Texas 75275
<sup>18</sup>Syracuse University, Syracuse, New York 13244
<sup>19</sup>University of Texas, Austin, TX 78712
<sup>20</sup>Vanderbilt University, Nashville, Tennessee 37235
<sup>21</sup>Virginia Polytechnic Institute and State University, Blacksburg, Virginia 24061
<sup>22</sup>Wayne State University, Detroit, Michigan 48202
<sup>23</sup>California Institute of Technology, Pasadena, California 91125
<sup>24</sup>University of California, San Diego, La Jolla, California 92093
<sup>25</sup>University of California, Santa Barbara, California 93106
<sup>26</sup>Carnegie Mellon University, Pittsburgh, Pennsylvania 15213
<sup>27</sup>University of Colorado, Boulder, Colorado 80309-0390
Measurements of exclusive $`B`$-decay branching fractions from $`e^+e^{}`$ collider operation at the $`\mathrm{{\rm Y}}`$(4S) resonance assume equal production rates of charged and neutral $`B`$-meson pairs . In the literature, the uncertainty in a specific branching fraction due to a lack of knowledge of the production ratio is often ignored.
Any physics based upon comparisons of absolute decay rates of charged and neutral $`B`$ mesons will profit from a more precise knowledge of the $`B`$-production ratio, $`f_+/f_{00}\mathrm{\Gamma }(\mathrm{{\rm Y}}(4S)B^+B^{})/\mathrm{\Gamma }(\mathrm{{\rm Y}}(4S)B^0\overline{B}^0)`$. For example, a comparison of the branching fractions of two-body hadronic decays can be used to obtain information on the relative contributions from external and internal spectator decays . For all exclusive decay modes studied, the $`B^+`$ branching fraction was found to be larger than the corresponding $`B^0`$ branching fraction, indicating constructive interference between the external and internal spectator amplitudes. This is in contrast to the destructive interference observed in hadronic charm decay. The magnitude of the constructively interfering fraction depends on the value of $`f_+/f_{00}`$. Another application of the $`f_+/f_{00}`$ ratio arises in the use of ratios of charmless hadronic $`B`$-decay rates to set bounds on the angle $`\gamma `$, the phase of the CKM matrix element $`V_{ub}`$ . The uncertainty on $`f_+/f_{00}`$ contributes to the systematic uncertainty of the $`\gamma `$ bound.
A better measurement of $`f_+/f_{00}`$ would also allow a more meaningful comparison with theoretical predictions of the relative $`B^+B^{}`$ and $`B^0\overline{B}^0`$ production rates at the $`\mathrm{{\rm Y}}`$(4S) resonance. If there are no other important differences between the two $`\mathrm{{\rm Y}}`$(4S) decays, such as $`B^+B^0`$ mass splitting or isospin-violating form factors in the decay amplitude, Coulomb corrections to $`B^+B^{}`$ production near threshold are not negligible, giving rise to $`\frac{\mathrm{\Gamma }(\mathrm{{\rm Y}}(4S)B^+B^{})}{\mathrm{\Gamma }(\mathrm{{\rm Y}}(4S)B^0\overline{B}^0)}1.18`$ . Other authors argue that the $`B`$-meson substructure cannot be ignored and strongly reduces the Coulomb effect in the $`B`$-production ratio to 1.05$``$1.07, depending on the $`B`$ masses and momenta.
Existing measurements of the admixture ratio of charged to neutral $`B`$ mesons produced at the $`\mathrm{{\rm Y}}(4S)`$ resonance have an uncertainty of $``$15%. One measurement used the branching-fraction ratio of $`(B^+J/\psi K^{()+})`$ to $`(B^0J/\psi K^{()0})`$ to yield $`\frac{f_+}{f_{00}}\times \frac{\tau _{B^+}}{\tau _{B^0}}=1.15\pm 0.17\pm 0.06`$, where the first uncertainty is statistical, the second is systematic, and $`\tau _B`$ denotes the $`B`$ lifetime. Another measurement used a ratio of $`BD^{}l\nu `$ decays to extract $`\frac{f_+}{f_{00}}\times \frac{\tau _{B^+}}{\tau _{B^0}}=1.14\pm 0.14\pm 0.13`$.
In the present analysis, we study the decays $`BJ/\psi K^{()}`$, which are isospin conserving transitions, since the $`J/\psi `$ daughter is an iso-singlet and the $`B`$ and $`K^{()}`$ mesons are both iso-doublets. The decays $`B^+J/\psi K^{()+}`$ and $`B^0J/\psi K^{()0}`$ must therefore have equal partial widths and we can extract $`R\frac{f_+}{f_{00}}\times \frac{\tau _{B^+}}{\tau _{B^0}}=\frac{𝒩(B^+J/\psi K^{()+})}{𝒩(B^0J/\psi K^{()0})}`$, where $`𝒩`$ is the efficiency-corrected signal yield. Using the ratio of two similar decay rates to extract $`R`$, we exploit the cancellation of common experimental uncertainties. Throughout this Letter, reference to charge conjugate states is implicit.
The data analyzed in this study were recorded at the Cornell Electron Storage Ring (CESR) with two configurations of the CLEO detector, CLEO II and CLEO II.V. The data consist of an integrated luminosity of $`9.2`$fb<sup>-1</sup> of $`e^+e^{}`$ annihilations recorded at the $`\mathrm{{\rm Y}}(4S)`$ resonance and of $`4.6`$fb<sup>-1</sup> taken in the continuum, 60 MeV below the $`\mathrm{{\rm Y}}(4S)`$ energy. The results in this Letter are based upon $`9.7\times 10^6`$ $`B\overline{B}`$ candidates and supersede those of Ref. .
The components of the CLEO detector most relevant to this analysis are the charged-particle tracking system, the 7800-crystal CsI electromagnetic calorimeter, and the muon chambers. The first third of the data were collected with the CLEO II detector , which measured the momenta of charged particles in a tracking system consisting of an inner 6-layer straw-tube chamber, a 10-layer precision drift chamber, and a 51-layer main drift chamber, all operating inside a 1.5 T solenoidal magnet. The main drift chamber also provided a measurement of the specific ionization loss ($`dE/dx`$) used in particle identification. Two thirds of the data were taken with the CLEO II.V configuration, for which the innermost straw-tube chamber was replaced with a 3-layer silicon vertex detector , and the argon-ethane gas of the main drift chamber was replaced with a helium-propane mixture. The muon identification system in both the CLEO II and CLEO II.V configurations consisted of proportional counters placed at various depths in the return yoke of the magnet.
Since the backgrounds for $`BJ/\psi K^{()}`$ decays are very low, track and photon quality requirements have been designed to maximize signal yield. We reconstruct $`BJ/\psi K^{()}`$ candidates in the data samples taken at the $`\mathrm{{\rm Y}}(4S)`$ energy.
Candidate $`J/\psi `$ mesons are reconstructed in their leptonic decay modes, requiring $`J/\psi `$ lepton daughter tracks to have momenta larger than 800 MeV/$`c`$.
For $`J/\psi `$ reconstruction in the muon channel, one of the muon candidates was required to penetrate the steel absorber to a depth greater than three nuclear interaction lengths. For the opposite sign daughter candidate, no muon detection requirement was imposed.
Electron candidates were identified based on the ratio of the track momentum to the associated shower energy in the CsI calorimeter and specific ionization loss in the drift chamber. Bremsstrahlung produces a radiative tail in the $`e^+e^{}`$ invariant mass distribution below the $`J/\psi `$ pole. We recovered some of the resultant efficiency loss by detecting the radiated photon. We selected photon candidates ($`E_\gamma >`$ 10 MeV) with the smallest angle to the $`e^\pm `$ track, provided this angle did not exceed 5. The $`J/\psi e^+e^{}`$ efficiency was increased by $``$20%, without adding background.
We reconstructed $`\mathrm{15\hspace{0.17em}900}\pm 700`$ inclusive $`J/\psi l^+l^{}`$ candidates (Fig. 1), about equally shared in the two dilepton reconstruction modes. The resolution in the $`J/\psi `$ invariant mass was $``$13 MeV.
We required the dimuon invariant mass to be within 50 MeV of the world-average $`J/\psi `$ mass , corresponding to a $``$3.5 standard deviation ($`\sigma `$) selection. For the dielectron invariant mass we required $`150`$ MeV $`<(m_{ee}m_{J/\psi })<`$ 50 MeV to allow for the radiative tail. The $`J/\psi `$ energy resolution was improved by a factor $``$4 after performing a kinematic fit of the dilepton invariant mass to the $`J/\psi `$ mass. We required $`J/\psi `$ candidates to have momenta below 2 GeV/$`c`$, which is near the kinematic limit for $`J/\psi `$ mesons originating from a $`B`$ meson nearly at rest.
The $`K_S^0\pi ^+\pi ^{}`$ candidates were selected from pairs of tracks forming well-measured displaced vertices. The resolution in $`\pi ^+\pi ^{}`$ invariant mass is approximately 2.5 MeV. Due to very low background in $`BJ/\psi K_S^0`$ candidates, we only require that neutral kaon candidates have a normalized mass within $`10\sigma `$ (because the $`K_S^0`$ mass distribution has non-negligible non-Gaussian tails) and a normalized flight distance greater than zero.
Charged kaon and pion candidates are required to have a measured $`dE/dx`$ within $`3\sigma `$ of the energy loss expected for the given particle type. Neutral pions are reconstructed from photon pairs detected within the barrel region of the CsI calorimeter, $`\mathrm{cos}\theta _\gamma <0.71`$, where $`\theta _\gamma `$ is the polar angle of the candidate photon with respect to the $`e^+e^{}`$ beam axis. The photons must have a minimum energy of 30 MeV and their normalized invariant mass is required to be within $`2.5\sigma `$ of the $`\pi ^0`$ mass. This diphoton invariant mass is then kinematically constrained to the $`\pi ^0`$ mass. Charged and neutral pions and kaons are used to reconstruct the four $`K^{}`$ decay modes. Candidate $`K^{}`$ mesons are required to have a $`K\pi `$ invariant mass within 75 MeV of the world-average $`K^{}`$ mass .
We fully reconstruct $`B`$-meson candidates by employing the kinematics of a $`B\overline{B}`$ pair produced almost at rest. We use the energy difference $`\mathrm{\Delta }EE(J/\psi )+E(K^{()})E_{\mathrm{beam}}`$ as well as the beam-constrained mass $`M(B)\sqrt{E_{\mathrm{beam}}^2p^2(B)}`$ as selection observables. The resolution in $`\mathrm{\Delta }E`$ is 15 MeV for $`J/\psi K^{}`$ with a $`\pi ^0`$ candidate in the final state and 9$``$11 MeV for the other modes. We find the resolution in $`M(B)`$ to be $``$2.5 MeV, which is dominated by the beam energy spread. We select signal candidates by requiring $`5.2`$ GeV$`<M(B)<`$ 5.3 GeV and $`\mathrm{\Delta }E<3\sigma _{\mathrm{\Delta }E}`$. The beam-constrained mass distributions for events within the $`\mathrm{\Delta }E`$ signal region are shown in Fig. 2.
We extract the signal yield in each mode by performing a binned maximum-likelihood fit to the $`M(B)`$ projection, where the signal is given by a single Gaussian distribution with fixed mean of 5.28 GeV and fixed width of 2.5 MeV. The background is fit to a first-order polynomial joined with an elliptic function to fit the threshold nature of the beam-constrained mass distribution. The $`M(B)`$ distributions in the $`\mathrm{\Delta }E`$ sideband regions exhibit a slope consistent with zero. These sideband regions are at least $`4\sigma _{\mathrm{\Delta }E}`$ and less than one pion mass away from the $`\mathrm{\Delta }E`$ signal region. We fix the slope of the background shape to zero and allow the level of the combinatoric background to be determined from the fit to the $`M(B)`$ projection of the $`\mathrm{\Delta }E`$ signal region.
We must account for the individual final states being reconstructed in a different channel (cross-feed), since for such candidates both the total energy and the beam-constrained mass lie near the signal region. We evaluate the reconstruction efficiency, as well as the amount of cross-feed from a given channel i to another channel j, using a sample of simulated $`BJ/\psi K_i^{()}`$ events to generate a $`6\times 6`$ efficiency matrix for the $`J/\psi e^+e^{}`$ and $`J/\psi \mu ^+\mu ^{}`$ cases, as well as for CLEO II and CLEO II.V, separately. The CLEO detector simulation is based upon GEANT . Simulated events are processed in a manner similar to that for the data. There is negligible cross-feed between the $`J/\psi K`$ and the $`J/\psi K^{}`$ modes. The cross-feed into $`J/\psi K^{}`$ modes with a charged-pion $`K^{}`$ daughter is near 5%, whereas cross-feed into $`J/\psi K^{}`$ modes with a neutral-pion $`K^{}`$ daughter ranges between $`830`$% of the raw yield. Efficiencies and cross-feed-corrected yields are listed in Table I. As a cross check, we also quote the branching fractions computed for the analyzed $`BJ/\psi K^{()}`$ decays.
We extract our result from the cross-feed and reconstruction-efficiency corrected yields using the world-average values for the respective daughter decay branching fractions . We obtain four independent measurements of $`R`$ listed in Table II.
We evaluate the uncertainties in the reconstruction efficiency due to track finding, track fitting, charged hadron identification, $`K_S^0`$ finding, and $`\pi ^0`$ finding. Since we use the ratio of two decay rates that each involve $`J/\psi l^+l^{}`$ candidates, uncertainties in lepton identification are negligible. We estimate the full systematic bias due to daughter reconstruction efficiency uncertainties by taking into account correlations between the different final states in the numerator and denominator, resulting in some cancellation. Propagating these uncertainties through the weighted average of the results in Table II, we arrive at a systematic uncertainty on $`R`$ due to the understanding of reconstruction efficiencies of $`{}_{1.5\%}{}^{}{}_{}{}^{+1.0\%}`$. The polarization of the decay $`BJ/\psi K^{}`$ is modeled in our simulation with a longitudinal polarization fraction of $`\mathrm{\Gamma }_L/\mathrm{\Gamma }`$ = 0.52 in accordance with Ref. . We estimate the impact of the value used for $`\mathrm{\Gamma }_L/\mathrm{\Gamma }`$ on the $`BJ/\psi K^{}`$ efficiencies by generating signal events with the nominal polarization varied by $`\pm 1\sigma =\pm 0.08`$. The central value for $`R`$ changes by less than 0.8% due to this variation. We vary the $`B`$ candidate signal width by $`\pm `$0.2 MeV and estimate the systematic uncertainty from this source to be less than 0.5%. We extract the signal yield using a background function that allows for a slope in the non-signal region of the beam-constrained mass distribution and assign a systematic bias of 3% on the central value for $`R`$ from our assumption of flat background. We attribute a 1.1% uncertainty to limited statistics of the simulated event samples used to extract the efficiency matrices. Adding all contributions in quadrature we obtain a total systematic uncertainty of $`{}_{3.7\%}{}^{}{}_{}{}^{+3.5\%}`$ on $`R`$.
We weight the results of Table II with their statistical uncertainty and, combined with the estimated systematic uncertainty, we extract
$`R\frac{f_+}{f_{00}}\times \frac{\tau _{B^+}}{\tau _{B^0}}=1.11\pm 0.07\pm 0.04,`$
where the first uncertainty is statistical and the second is systematic.
Using the world-average lifetime ratio of charged and neutral $`B`$ mesons, $`1.066\pm 0.024`$ , we obtain a measurement of the production ratio
$`\frac{f_+}{f_{00}}=\frac{\mathrm{\Gamma }(\mathrm{{\rm Y}}(4\mathrm{S})\mathrm{B}^+\mathrm{B}^{})}{\mathrm{\Gamma }(\mathrm{{\rm Y}}(4\mathrm{S})\mathrm{B}^0\overline{\mathrm{B}}^0)}=1.04\pm 0.07\pm 0.04,`$
and, assuming $`f_++f_{00}=1`$, we also extract $`f_{00}=0.49\pm 0.02\pm 0.01`$ and $`f_+=0.51\pm 0.02\pm 0.01`$.
We have measured the ratio of charged to neutral production of $`B`$ mesons at the $`\mathrm{{\rm Y}}(4S)`$ resonance to be $`1.04\pm 0.07\pm 0.04`$, which is consistent with unity within an error of 8%. This is the most precise measurement of $`f_+/f_{00}`$. Our result is also consistent with theoretical predictions of greater charged than neutral $`B`$-meson production in $`\mathrm{{\rm Y}}(4S)`$ decays near threshold. We emphasize that the ratio $`f_+/f_{00}`$ and its uncertainty must be taken into account when performing measurements that compare charged and neutral $`B`$ decays at the $`\mathrm{{\rm Y}}(4S)`$ resonance.
We gratefully acknowledge the effort of the CESR staff in providing us with excellent luminosity and running conditions. This work was supported by the National Science Foundation, the U.S. Department of Energy, the Research Corporation, the Natural Sciences and Engineering Research Council of Canada, the A.P. Sloan Foundation, the Swiss National Science Foundation, and the Alexander von Humboldt Stiftung.
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# Local models in the ramified case I. The EL-case
## 1 Introduction
In the arithmetic theory of Shimura varieties it is of interest to have a model over $`\mathrm{Spec}𝒪_E`$, where $`E`$ is the completion of the reflex field at some finite prime of residue characteristic $`p`$. If the Shimura variety is the moduli space of abelian varieties of PEL type and the level structure at $`p`$ is of parahoric type, such a model was defined in \[RZ\] by posing the moduli problem over $`\mathrm{Spec}𝒪_E`$. In \[RZ\] it was conjectured that this model is flat over $`\mathrm{Spec}𝒪_E`$, but in \[P\] it was shown that this is not always true. It still seems reasonable to expect this flatness property when the group $`G`$ defining the Shimura variety splits over an unramified extension of $`𝐐_p`$, and this is supported by the theorem of Görtz \[G\] in the EL-case. On the other hand, in \[P\] it is shown that this conjecture fails in general if the group $`G`$ defining the Shimura variety has localization over $`𝐐_p`$ isomorphic to the group of unitary similitudes corresponding to a ramified quadratic extension of $`𝐐_p`$. Furthermore, a modified moduli problem was proposed in loc. cit. which defines a closed subscheme of the original model and which stands a better chance of being flat over $`\mathrm{Spec}𝒪_E`$.
Our main purpose of the present paper is to come to grips with the phenomenon of non-flatness by investigating the simplest case in which it can occur. As usual, the problem can be reduced to the consideration of the associated local model which locally for the étale topology around each point of the special fiber coincides with the model of the Shimura variety. In the rest of the paper we only consider the local models which can be defined in terms of linear algebra as schemes over the spectrum of a complete discrete valuation ring with perfect residue field. However, in view of the fact that the models proposed in \[RZ\] are not “the right ones” in general we shall term them naive local models and reserve the name of local models for certain closed subschemes of the naive local models defined in the body of the paper. Both of them have as generic fibers a closed subscheme of a Grassmannian, and as special fiber a closed subvariety of the affine partial flag variety over the residue field corresponding to the fixed parahoric.
As a by-product of our investigations, as they concern the special fibers, we also obtain several results on the structure of Schubert varieties in affine Grassmannians and their relation to nilpotent orbit closures.
The simplest case of a naive local model occurs for the standard models. Let us define them. Let $`F_0`$ be a complete discretely valued field with ring of integers $`𝒪_{F_0}`$ and perfect residue field. Let $`F/F_0`$ be a totally ramified extension of degree $`e`$, with ring of integers $`𝒪_F`$. Let $`V`$ be a $`F`$-vector space of dimension $`d`$, and let $`\mathrm{\Lambda }`$ be a $`𝒪_F`$-lattice in $`V`$. Choose for each embedding $`\phi `$ of $`F`$ into a separable closure $`F_0^{\mathrm{sep}}`$ of $`F_0`$ an integer $`r_\phi `$ with
(1.1)
$$0r_\phi d,\phi .$$
Put $`r=_\phi r_\phi `$. Then the standard model for $`GL_d`$ and $`𝐫=(r_\phi )`$, denoted $`M(\mathrm{\Lambda },𝐫)`$, parametrizes the points in the Grassmannian of subspaces $``$ of rank $`r`$ of $`\mathrm{\Lambda }`$ which are $`𝒪_F`$-stable and on which the representation of $`𝒪_F`$ is prescribed in terms of $`𝐫`$ (comp. (2)). It is defined over $`\mathrm{Spec}𝒪_E`$ where $`E=E(V,𝐫)`$ is the reflex field, the field of definition of the prescribed representation of $`𝒪_F`$. Over $`F_0^{\mathrm{sep}}`$ the scheme $`M=M(\mathrm{\Lambda },𝐫)`$ is isomorphic to the product over all $`\phi `$ of Grassmannians of subspaces of dimension $`r_\phi `$ in a vector space of dimension $`d`$. Over the residue field $`k`$ of $`𝒪_E`$ the scheme $`\overline{M}=M_{𝒪_E}k`$ is a closed subvariety of the affine Grassmannian of $`GL_d`$ over $`k`$, and is in fact a union of Schubert strata which are enumerated by the following dominant coweights of $`GL_d`$
(1.2)
$$𝒮^0(r,e,d)=\{𝐬=s_1\mathrm{}s_d;es_1,s_d0,\underset{i}{}s_i=r\}.$$
An easy dimension count shows that the dimension of the generic fiber and the special fiber are unequal, unless all integers $`r_\phi `$ differ by at most 1. In this last case we conjecture, and often can prove that $`M`$ is flat over $`\mathrm{Spec}𝒪_E`$. In particular, in the Hilbert-Blumenthal case ($`d=2`$, and $`r_\phi =1`$, $`\phi `$) the standard model coincides with the flat model constructed in this case in \[DP\]. In all cases when two $`r_\phi `$ differ from each other by more than one, $`M`$ is not flat over $`\mathrm{Spec}𝒪_E`$.
To analyze the standard model we construct a diagram
(1.3)
$$M\stackrel{\pi }{}\stackrel{~}{M}\stackrel{\varphi }{}N.$$
Here $`\stackrel{~}{M}`$ is the $`GL_r`$-torsor over $`M`$ which fixes a basis of the variable subspace $``$. The morphism $`\varphi `$ is defined as follows. Let $`\pi `$ be a uniformizer of $`𝒪_F`$ which satisfies an Eisenstein polynomial $`Q(\pi )=0`$. Let
(1.4)
$$N=\{A\mathrm{Mat}_{r\times r};\mathrm{det}(TIA)\underset{\phi }{}(T\phi (\pi ))^{r_\phi },Q(A)=0\}.$$
By considering the action of $`\pi `$ on the variable subspace $``$ and expressing it as a matrix in terms of the fixed basis of $``$, we obtain the morphism $`\varphi `$. We show the following result (Theorem 4.1).
Theorem A: The morphism $`\varphi `$ is smooth of relative dimension $`rd`$.
Using Theorem A many structure problems on $`M`$ can be reduced to corresponding questions on $`N`$. After a finite extension $`𝒪_E𝒪_K`$, the variety $`N`$ can be seen as a rank variety in the sense of Eisenbud and Saltman (\[ES\]), whereas the special fiber $`\overline{N}=N_{𝒪_E}k`$ is a subscheme of the nilpotent variety,
(1.5)
$$\overline{N}=\{A\mathrm{Mat}_{r\times r};\mathrm{det}(TIA)T^r,A^e=0\}.$$
Using now results of Mehta and van der Kallen (\[M-vdK\]) on the structure of the closures of nilpotent conjugacy classes and basing ourselves on the methods of Eisenbud and Saltman, we use this reduction procedure to prove the following result (Theorem 5.4).
Theorem B: Let $`M^{\mathrm{loc}}`$ be the scheme theoretic closure of $`M_{𝒪_E}E`$ in $`M`$. Then
1. $`M^{\mathrm{loc}}`$ is normal and Cohen-Macaulay.
2. The special fiber $`\overline{M}^{\mathrm{loc}}`$ is reduced, normal with rational singularities and is the union of all those strata of $`\overline{M}`$ which correspond to $`𝐬`$ in (1.2) with $`𝐬𝐫^{}`$ (dual partition to $`𝐫`$).
3. If the scheme $`\overline{N}`$ is reduced, then $`M^{\mathrm{loc}}=M`$ provided that all $`r_\phi `$ differ by at most 1.
We conjecture that the hypothesis made in (iii) is automatically satisfied, i.e. that $`\overline{N}`$ is always reduced. This is true by a classical result of Kostant if $`re`$ (in which case the second condition in (1.5) is redundant). For $`e<r`$ it seems a difficult problem. In a companion paper to ours, J. Weyman proves our conjecture for $`e=2`$, and for arbitrary $`e`$ when char $`k=0`$.
Recall that Lusztig \[L\] has interpreted certain Schubert varieties in the affine Grassmannian of $`GL_r`$ as a compactification of the nilpotent variety of $`GL_r`$ (namely the Schubert variety corresponding to the coweight $`(r,0,\mathrm{},0)`$), compatible with the orbit stratifications of both varieties. In particular, as used by Lusztig in his paper, all singularities of nilpotent orbit closures occur in certain Schubert varieties in the affine Grassmannians. However, the main thrust of Theorem A (when restricted to the special fibers) goes in the other direction. Namely, we prove:
Theorem C: Any Schubert variety in the affine Grassmannian of $`GL_d`$ is smoothly equivalent to a nilpotent orbit closure for $`GL_r`$, for suitable $`r`$. In particular, it is normal with rational singularities.
We point out the recent preprint by Faltings \[F\] in which he proves basic results (like normality) on Schubert varieties in affine flag varieties for arbitrary reductive groups.
The disadvantage of $`M^{\mathrm{loc}}`$ compared to $`M`$ is that $`M^{\mathrm{loc}}`$ is not defined by a moduli problem in general. However, assume that $`e=2`$ and order the embeddings so that $`r_{\phi _1}r_{\phi _2}`$. If $`r_{\phi _1}=r_{\phi _2}`$, then $`E=F_0`$ and it follows from (iii) above and Weyman’s result that $`M^{\mathrm{loc}}=M`$. If $`r_{\phi _1}>r_{\phi _2}`$, we may use $`\phi _1`$ to identify $`E`$ with $`F`$. Then using work of Strickland (see Cor. 5.10) we can see that $`M^{\mathrm{loc}}`$ is defined inside $`M`$ by the following condition on $``$,
(1.6)
$$^{r_{\phi _2}+1}(\pi \phi _1(\pi )\mathrm{Id}|)=0.$$
In general, there is a connection to a conjecture of De Concini and Procesi describing the ideal of the closure of a nilpotent conjugacy class. Assuming this conjecture to be true and under a technical hypothesis we can write down a number of conditions which would define $`M^{\mathrm{loc}}`$ inside $`M`$. Since these conditions are highly redundant, the interest of such a description may be somewhat limited, however.
Let $`N^{\mathrm{loc}}`$ be the scheme theoretic closure of $`N_{𝒪_E}E`$ in $`N`$. Then the special fiber $`\overline{N}^{\mathrm{loc}}`$ can be identified with the closure in $`\overline{N}`$ of the nilpotent conjugacy class corresponding to $`𝐫^{}`$. Let $`K`$ be the Galois closure of $`F/F_0`$ with ring of integers $`𝒪_K`$ and residue field $`k^{}`$. Then $`N_{𝒪_E}𝒪_K`$ has a canonical resolution of singularities,
(1.7)
$$\mu :𝒩N_{𝒪_E}𝒪_K,$$
i.e. a morphism with source a smooth $`𝒪_K`$-scheme $`𝒩`$ which is an isomorphism on the generic fiber and whose special fiber can be identified with the Springer-Spaltenstein resolution of $`\overline{N}^{\mathrm{loc}}_kk^{}`$. Using the fact that $`\mu _{𝒪_K}k^{}`$ is semi-small, one can calculate the direct image of the constant perverse sheaf on $`\overline{𝒩}`$ (Borho, MacPherson \[BM\], Braverman, Gaitsgory \[BG\]). Transporting back the result to $`M`$ via Theorem A, we obtain the following result (comp. Theorem 7.1).
Theorem D: Let $`R\psi ^{}`$ denote the complex of nearby cycles of the $`𝒪_K`$-scheme $`M_{𝒪_E}𝒪_K`$. Then there is the following identity of perverse sheaves pure of weight 0 on $`\overline{M}_kk^{}`$,
$$R\psi ^{}[\mathrm{dim}\overline{M}](\frac{1}{2}\mathrm{dim}\overline{M})=\underset{𝐬𝐫^{}}{}K_{𝐫^{},𝐬}IC_{M_𝐬_kk^{}}.$$
Here $`K_{𝐫^{},𝐬}`$ is a Kostka number.
We refer to section 7 for an explanation of the notation and for the question of descending this result from $`𝒪_K`$ to $`𝒪_E`$.
Via Theorem A, the resolution of singularities (1.7) is closely related to an analogous resolution of singularities of $`M_{𝒪_E}𝒪_K`$ whose special fiber can be identified with the Demazure resolution of the affine Schubert variety corresponding to the coweight $`𝐫^{}`$ of $`GL_d`$. As pointed out to us by Ngô, from this identification one obtains another formula for the complex of nearby cycles $`R\psi ^{}`$.
This concludes the discussion of our results on the standard models for $`GL_d`$ and $`𝐫`$. The general naive local models of $`EL`$-type are projective schemes $`M^{\mathrm{naive}}`$ over $`\mathrm{Spec}𝒪_E`$ which at least over $`\mathrm{Spec}𝒪_{\stackrel{˘}{E}}`$ (where $`\stackrel{˘}{E}`$ is the completion of the maximal unramified extension of $`E`$) are closed subschemes of products of standard models. We define closed subschemes $`M^{\mathrm{loc}}`$ of $`M^{\mathrm{naive}}`$ by demanding that the projection into any standard model lies in the flat closure considered above. The modification of the flatness conjecture of \[RZ\] in the present case is that $`M^{\mathrm{loc}}`$ is flat over $`\mathrm{Spec}𝒪_E`$. We have nothing to say about this conjecture, except to point out that the description of $`M^{\mathrm{loc}}(k)`$ implicit in \[KR\], in terms of $`\mu `$-permissible elements in the Iwahori-Weyl group is correct, as follows from Theorem B, (ii).
In the original version of this paper Theorem A was proved by exhibiting affine charts around the worst singularities of $`M`$ over which the morphisms $`\pi `$ and $`\varphi `$ can be made completely explicit. The present simple proof of Theorem A is based on ideas from \[FGKV\]. We thank D. Gaitsgory for pointing it out to us. We would also like to thank G. Laumon, T. Haines and B.C. Ngô for helpful discussions on the material of sections 6 and 7. We are grateful to J. Weyman for his interest in our conjecture and for making his results available in a companion paper to ours. We thank G. Pfister for computer calculations using the symbolic algebra package REDUCE, which gave us the courage to elevate an initially naive question to the rank of a conjecture. We also thank the Max-Planck-Institut Bonn for its hospitality and support. The first named author was also partially supported by NSF grant DMS99-70378 and by a Sloan Research Fellowship.
General notational convention: If $`X`$ is a scheme over Spec $`R`$ and $`R^{}`$ is an $`R`$-algebra, we often write $`X_RR^{}`$ or $`X_R^{}`$ for $`X\times _{\mathrm{Spec}R}\mathrm{Spec}R^{}`$.
## 2 Standard models for $`GL_d`$
In this section and the sections 37 we will use the following notation. Let $`F_0`$ be a complete discretely valued field with ring of integers $`𝒪_{F_0}`$ and uniformizer $`\pi _0`$, and perfect residue field. Let $`F`$ be a totally ramified separable extension of degree $`e`$ of $`F_0`$, with ring of integers $`𝒪_F`$. Let $`\pi `$ be a uniformizer of $`𝒪_F`$ which is a root of the Eisenstein polynomial
(2.1)
$$Q(T)=T^e+\underset{k=0}{\overset{e1}{}}b_kT^k,b_0\pi _0𝒪_{F_0}^\times ,b_k(\pi _0).$$
Let $`V`$ be an $`F`$-vector space of dimension $`d`$ and $`\mathrm{\Lambda }`$ an $`𝒪_F`$-lattice in $`V`$. We fix a separable closure $`F_0^{\mathrm{sep}}`$ of $`F_0`$. Finally, we choose for each embedding $`\phi :FF_0^{\mathrm{sep}}`$ an integer $`r_\phi `$ with
(2.2)
$$0r_\phi d.$$
Associated to these data we have the reflex field $`E`$, a finite extension of $`F_0`$ contained in $`F_0^{\mathrm{sep}}`$ with
(2.3)
$$\mathrm{Gal}(F_0^{\mathrm{sep}}/E)=\{\sigma \mathrm{Gal}(F_0^{\mathrm{sep}}/F_0);r_{\sigma \phi }=r_\phi ,\phi \}.$$
Let $`𝒪_E`$ be the ring of integers in $`E`$. We now formulate a moduli problem on $`(\mathrm{Sch}/\mathrm{Spec}𝒪_E)`$:
$`M(S)`$ $`=`$ $`\{\mathrm{\Lambda }_{𝒪_{F_0}}𝒪_S;\text{ a }𝒪_F_{𝒪_{F_0}}𝒪_S\text{-submodule,}`$
which is locally on $`S`$ a direct summand as
$`𝒪_S`$-module, with $`\mathrm{det}(a|)=\underset{\phi }{}\phi (a)^{r_\phi }\}.`$
The last identity is meant as an identity of polynomial functions on $`𝒪_F`$ (comp. \[K\], \[RZ\]).
It is obvious that this functor is representable by a projective scheme $`M=M(\mathrm{\Lambda },𝐫)`$ over $`\mathrm{Spec}𝒪_E`$. This scheme is called the standard model for $`GL_d`$ corresponding to $`𝐫=(r_\phi )_\phi `$ (and to $`F_0,F`$ and $`\pi `$).
Let us analyze the geometric general fiber and the special fiber of $`M`$. We have a decomposition
(2.5)
$$F_{F_0}F_0^{\mathrm{sep}}=\underset{\phi :FF_0^{\mathrm{sep}}}{}F_0^{\mathrm{sep}}.$$
Correspondingly we get a decomposition of $`V_{F_0}F_0^{\mathrm{sep}}`$ into $`F_0^{\mathrm{sep}}`$-vector spaces
(2.6)
$$V_{F_0}F_0^{\mathrm{sep}}=\underset{\phi }{}V_\phi .$$
Each summand is of dimension $`d`$. The determinant condition in (2) can now be interpreted as saying that $`M_{𝒪_E}F_0^{\mathrm{sep}}`$ parametrizes subspaces $`_\phi `$ of $`V_\phi `$, one for each $`\phi `$, of dimension $`r_\phi `$. In other words,
(2.7)
$$M_{𝒪_E}F_0^{\mathrm{sep}}=\underset{\phi }{}\mathrm{Grass}_{r_\phi }(V_\phi ).$$
In particular,
(2.8)
$$dimM_{𝒪_E}E=\underset{\phi }{}r_\phi (dr_\phi ).$$
Denote by $`k`$ the residue field of $`𝒪_E`$. Let us consider $`\overline{M}=M_{𝒪_E}k`$. Put
(2.9)
$$W=\mathrm{\Lambda }_{𝒪_{F_0}}k,\mathrm{\Pi }=\pi \mathrm{id}_k.$$
Then $`W`$ is a $`k`$-vector space of dimension $`de`$ and $`\mathrm{\Pi }`$ is a nilpotent endomorphism with $`\mathrm{\Pi }^e=0`$. The conditions that the subspace $`W`$ gives a point of $`\overline{M}`$ translate into the following: $``$ is $`\mathrm{\Pi }`$-stable, $`dim_k=r:=_\phi r_\phi `$ and $`det(T\mathrm{\Pi }|)T^r`$. In other words $`M_{𝒪_E}k`$ is the closed subscheme of $`\mathrm{\Pi }`$-stable subspaces $``$ in $`\mathrm{Grass}_r(W)`$ which satisfy the above condition on the characteristic polynomial. We point out that the $`k`$-scheme $`\overline{M}=M_{𝒪_E}k`$ only depends on $`r`$, not on the partition $`(r_\phi )`$ of $`r`$.
## 3 Relation to the affine Grassmannian
We denote by $`\stackrel{~}{\mathrm{Grass}}_k`$ the affine Grassmannian over $`k`$ associated to $`GL_d`$. Recall (\[BL\]) that this is the Ind-scheme over $`\mathrm{Spec}k`$ whose $`k`$-rational points parametrize the $`k[[\mathrm{\Pi }]]`$-lattices in $`k((\mathrm{\Pi }))^d`$. Here $`k[[\mathrm{\Pi }]]`$ denotes the power series ring in the indeterminate $`\mathrm{\Pi }`$ over $`k`$. On $`\stackrel{~}{\mathrm{Grass}}_k`$ we have an action of the group scheme $`\stackrel{~}{𝒢}`$ over $`k`$ with $`k`$-rational points equal to $`GL_d(k[[\mathrm{\Pi }]])`$. The orbits of this action are finite-dimensional irreducible locally closed subvarieties (with the reduced scheme structure) which are parametrized by the dominant coweights of $`GL_d`$, i.e. by $`d`$-tuples of integers $`𝐬=(s_1,\mathrm{},s_d)`$ with
(3.1)
$$s_1\mathrm{}s_d.$$
Furthermore (\[BL\]), if $`𝒪_𝐬`$ denotes the orbit corresponding to $`𝐬`$, we have
(3.2)
$$\mathrm{dim}𝒪_𝐬=𝐬,2\varrho \text{and}𝒪_𝐬^{}\mathrm{closure}(𝒪_𝐬)𝐬^{}𝐬.$$
Here $`2\varrho =(d1,d3,\mathrm{},1d)`$ and $`,`$ denotes the standard scalar product on $`𝐑^d`$. Furthermore, $`𝐬^{}𝐬`$ denotes the usual partial order on dominant coweights, i.e.
$$𝐬^{}𝐬s_1^{}s_1,s_1^{}+s_2^{}s_1+s_2,\mathrm{},s_1^{}+\mathrm{}+s_d^{}=s_1+\mathrm{}+s_d.$$
Let us fix an isomorphism of $`k[[\mathrm{\Pi }]]`$-modules
(3.3)
$$\mathrm{\Lambda }_{𝒪_{F_0}}k(k[[\mathrm{\Pi }]]/\mathrm{\Pi }^e)^d.$$
Then we obtain a closed embedding (of Ind-schemes),
(3.4)
$$\iota :M_{𝒪_E}k\stackrel{~}{\mathrm{Grass}}_k.$$
On $`k`$-rational points, $`\iota `$ sends a point of $`M(k)`$, corresponding to a $`k[\mathrm{\Pi }]`$-submodule $``$ of $`\mathrm{\Lambda }_{𝒪_{F_0}}k=(k[[\mathrm{\Pi }]]/\mathrm{\Pi }^e)^d`$, to its inverse image $`\stackrel{~}{}`$ in $`k[[\mathrm{\Pi }]]^d`$, which is a $`k[[\mathrm{\Pi }]]`$-lattice contained in $`k[[\mathrm{\Pi }]]^d`$,
(3.5)
$$\begin{array}{ccc}\stackrel{~}{}& & k[[\mathrm{\Pi }]]^d\\ & & \\ & & (k[[\mathrm{\Pi }]]/\mathrm{\Pi }^e)^d& .\end{array}$$
The embedding $`\iota `$ is equivariant for the action of $`\stackrel{~}{𝒢}`$ in the following sense. Consider the smooth group scheme $`𝒢`$ over $`\mathrm{Spec}𝒪_{F_0}`$,
(3.6)
$$𝒢=\underset{¯}{\mathrm{Aut}}_{𝒪_F}(\mathrm{\Lambda }).$$
In fact, we will only need the base change of $`𝒢`$ to $`\mathrm{Spec}𝒪_E`$ which we denote by the same symbol. The group scheme $`𝒢`$ acts on $`M`$ by
(3.7)
$$(g,)g().$$
Let
$`\overline{𝒢}=𝒢_{𝒪_E}k`$ $`=`$ $`\underset{¯}{\mathrm{Aut}}_{k[\mathrm{\Pi }]/\mathrm{\Pi }^e}(\mathrm{\Lambda }_{𝒪_{F_0}}k)`$
$``$ $`GL_d(k[[\mathrm{\Pi }]]/\mathrm{\Pi }^e).`$
In this way $`\overline{𝒢}`$ becomes a factor group of $`\stackrel{~}{𝒢}`$ and the equivariance of $`\iota `$ means that the action of $`\stackrel{~}{𝒢}`$ stabilizes the image of $`\iota `$, that the action on this image factors through $`\overline{𝒢}`$ and that $`\iota `$ is $`\overline{𝒢}`$-equivariant.
A point of $`M`$ with values in a field extension $`k^{}`$ of $`k`$, corresponding to a $`\mathrm{\Pi }`$-stable subspace $``$ of $`\mathrm{\Lambda }_{𝒪_{F_0}}k^{}`$, has image in $`𝒪_𝐬_{}`$, where $`𝐬_{}`$ is the Jordan type of the nilpotent endomorphism $`\mathrm{\Pi }`$. It follows that the orbit decomposition of $`\overline{M}=M_{𝒪_E}k`$ under the action of $`\overline{𝒢}`$ has the form
(3.9)
$$\overline{M}=\underset{𝐬}{}M_𝐬.$$
Here $`M_𝐬𝒪_𝐬`$ via $`\iota `$ and $`𝐬`$ ranges over the subset of (3.1) given by
(3.10)
$$𝒮^0(r,e,d)=\{𝐬=s_1s_2\mathrm{}s_d;es_1,s_d0,\underset{i}{}s_i=r\},$$
i.e. the partitions of $`r`$ into at most $`d`$ parts bounded by $`e`$. In all of the above we have ignored nilpotent elements.
###### Proposition 3.1
The special fibre $`\overline{M}=M_{𝒪_E}k`$ is irreducible of dimension $`drec^2(2c+1)f`$.
Here we have written $`r=ce+f`$ with $`0f<e`$.
Proof. Among the coweights in $`𝒮^0(r,e,d)`$ there is a unique maximal one,
(3.11)
$$𝐬_{\mathrm{max}}=𝐬_{\mathrm{max}}(r,e)=(e,\mathrm{},e,f,0\mathrm{}0)=(e^c,f).$$
Hence $`M_{𝐬_{\mathrm{max}}}`$ is open and dense in $`\overline{M}`$. Its dimension is equal to $`𝐬_{\mathrm{max}},2\varrho `$, which gives the result.
Sometimes for convenience we shall number the embeddings $`\phi `$ in such a way that the $`r_i=r_{\phi _i}`$ form a decreasing sequence $`𝐫=(r_1r_2\mathrm{}r_e)`$. Then $`𝐫`$ is a partition of $`r`$ into at most $`e`$ parts bounded by $`d`$.
Let $`𝐫_{\mathrm{min}}=𝐫_{\mathrm{min}}(r,e)=𝐬_{\mathrm{max}}^{}`$ be the dual partition to $`𝐬_{\mathrm{max}}`$, i.e.
(3.12)
$$𝐫_{\mathrm{min}}=(c+1,\mathrm{},c+1,c,\mathrm{},c)=((c+1)^f,c^{ef}).$$
###### Proposition 3.2
We have
$$\mathrm{dim}M(\mathrm{\Lambda },𝐫)_{𝒪_E}E\mathrm{dim}M(\mathrm{\Lambda },𝐫)_{𝒪_E}k,$$
with equality if and only if $`𝐫=𝐫_{\mathrm{min}}(r,e)`$ (after renumbering $`𝐫`$), i.e. iff all $`r_\phi `$ differ by at most one.
Proof. If $`𝐫=𝐫_{\mathrm{min}}(r,e)`$, then
$`\mathrm{dim}M(\mathrm{\Lambda },𝐫)_{𝒪_E}E`$ $`=`$ $`dr{\displaystyle \underset{\phi }{}}r_\phi ^2=drf(c+1)^2(ef)c^2`$
$`=`$ $`drf((c+1)^2c^2)ec^2`$
$`=`$ $`drec^2(2c+1)f`$
$`=`$ $`\mathrm{dim}M(\mathrm{\Lambda },𝐫)_{𝒪_E}k.`$
Here we used (2.8) in the first line and the previous proposition in the last line. Now let $`𝐫`$ be arbitrary and let $`𝐭=𝐫^{}`$ be the dual partition to $`𝐫`$, i.e.
(3.13)
$$t_1=\mathrm{\#}\{\phi ;r_\phi 1\},t_2=\mathrm{\#}\{\phi ;r_\phi 2\},\text{etc.}$$
By (2.2) the partition $`𝐭`$ lies in $`𝒮^0(r,e,d)`$. Hence $`𝐭𝐬_{\mathrm{max}}(r,e)`$ with equality only if $`𝐫=𝐫_{\mathrm{min}}(r,e)`$. Hence, if $`𝐫𝐫_{\mathrm{min}}(r,e)`$, we have
(3.14)
$$\mathrm{dim}M_𝐭<\mathrm{dim}M_{𝐬_{\mathrm{max}}}=\mathrm{dim}\overline{M}.$$
Now
$`\mathrm{dim}M_𝐭`$ $`=`$ $`𝐭,2\varrho =t_1(d1)+t_2(d3)+\mathrm{}+t_d(1d)`$
$`=`$ $`d{\displaystyle \underset{i}{\overset{d}{}}}t_i{\displaystyle \underset{i=1}{\overset{d}{}}}(2i1)t_i.`$
But $`_it_i=_\phi r_\phi =r`$. The second sum on the right hand side can be written as a sum of contributions of each $`\phi `$. Each fixed $`\phi `$ contributes $`_{j=1}^{r_\phi }(2j1)=r_\phi ^2`$. Hence
(3.15)
$$\mathrm{dim}M_𝐭=dr\underset{\phi }{}r_\phi ^2=\mathrm{dim}M(\mathrm{\Lambda },𝐫)_{𝒪_E}E.$$
Taking into account (3.14), the result follows.
###### Corollary 3.3
If $`𝐫𝐫_{\mathrm{min}}(r,e)`$, the corresponding standard model $`M`$ is not flat over $`\mathrm{Spec}𝒪_E`$.
There remains the question whether if $`𝐫=𝐫_{\mathrm{min}}(r,e)`$, the corresponding standard model is flat over $`𝒪_E`$; we will return to it in section 5.
We end this section with the following remark. Among all strata of $`\overline{M}`$ enumerated by $`𝒮^0(r,e,d)`$ there is a unique minimal one,
(3.16)
$$𝐬_{\mathrm{min}}=𝐬_{\mathrm{min}}(r,d)=((u+1)^j,u^{dj}).$$
Here we have written $`r=ud+j,0j<d`$. The corresponding stratum $`M_{𝐬_{\mathrm{min}}}`$ is closed and lies in the closure of any other stratum.
## 4 Relation to the nilpotent variety
In this section we describe the connection of the standard models for $`GL_d`$ with the nilpotent variety. Consider the $`GL_r`$-torsor $`\pi :\stackrel{~}{M}M`$ where
(4.1)
$$\stackrel{~}{M}(S)=\{(,\psi );M(S),\psi :\stackrel{}{}𝒪_S^r\}$$
and $`GL_r`$ acts via $`(,\psi )(,\gamma \psi )`$. Then $`𝒢`$ acts on $`\stackrel{~}{M}`$ via
$$(g,(,\psi ))(g(),\psi g^1),$$
and the morphism $`\pi `$ is equivariant.
Note that $`_\phi (T\phi (\pi ))^{r_\phi }`$ has coefficients in $`𝒪_E`$; let us define a scheme $`N=N(𝐫)`$ over $`\mathrm{Spec}𝒪_E`$ via the functor which to the $`𝒪_E`$-algebra $`R`$ associates the set
(4.2)
$$\{AM_{r\times r}(R);\mathrm{det}(TIA)\mathrm{\Pi }_\phi (T\phi (\pi ))^{r_\phi },Q(A)=0\}.$$
The group scheme $`GL_r`$ acts on $`N`$ via conjugation $`A\gamma A\gamma ^1`$. There is a morphism
(4.3)
$$\varphi :\stackrel{~}{M}N,\varphi ((,\psi ))=\psi (\pi |)\psi ^1,$$
which is $`GL_r`$-equivariant. We have for $`g𝒢`$,
$$\varphi (g(,\psi ))=\psi g^1(\pi |g())(\psi g^1)^1=\psi (\pi |)\psi ^1=\varphi ((,\psi ))$$
since $`g`$ commutes with $`\pi `$. Therefore, the morphism $`\varphi `$ is $`𝒢`$-equivariant with trivial $`𝒢`$-action on the target $`N`$. We therefore obtain a diagram of morphisms
(4.4)
$$M\stackrel{\pi }{}\stackrel{~}{M}\stackrel{\varphi }{}N,$$
which is equivariant for the action of $`𝒢\times GL_r`$, where the first factor acts trivially on the right hand target and the second factor acts trivially on the left hand target.
Since $`Q(T)=_\phi (T\phi (\pi ))`$ with all elements $`\phi (\pi )`$ pairwise distinct, we see that the generic fiber of $`N`$ consists of all semisimple matrices with eigenvalues $`\phi (\pi )`$ with multiplicity $`r_\phi `$. It follows that $`GL_r`$ acts transitively on $`N_{𝒪_E}E`$, which is smooth of dimension $`r^2\mathrm{\Sigma }_\phi r_\phi ^2`$ over $`\mathrm{Spec}E`$. The special fiber $`\overline{N}`$ of $`N`$ is the subscheme of $`\mathrm{Mat}_{r\times r}`$ defined by the equations
(4.5)
$$\overline{N}=\{A\mathrm{Mat}_{r\times r};A^e=0,\mathrm{det}(TIA)T^r\}.$$
This is a closed subscheme of the variety of nilpotent $`r\times r`$ matrices (which is defined by $`\mathrm{det}(TIA)T^r`$).
The following result exhibits a close connection between the standard model and the nilpotent variety.
###### Theorem 4.1
The morphism $`\varphi :\stackrel{~}{M}N`$ is smooth of relative dimension rd.
Proof. Let $`𝔐𝔬𝔡=𝔐𝔬𝔡(𝒪_F,𝐫)`$ be the algebraic stack over $`\mathrm{Spec}𝒪_E`$ given by the fibered category of $`𝒪_F_{𝒪_{F_0}}𝒪_S`$-modules $``$ which are locally free $`𝒪_S`$-modules of rank $`r`$ (comp. \[LMoB\] 3.4.4, 4.6.2.1) and for which
$$\mathrm{det}(TI\pi |)\underset{\varphi }{}(T\varphi (\pi ))^{r_\varphi }$$
as polynomials. There is an isomorphism
$$𝔐𝔬𝔡[N/GL_r]$$
(where the quotient stack is for the conjugation action) given by $`𝒢`$ the conjugation $`GL_r`$-torsor of matrices $`A`$ giving the action of $`\pi `$ on $`𝒢`$. By the definition of the quotient stack the diagram (4.4) corresponds to a morphism
(4.6)
$$\overline{\varphi }:M[N/GL_r].$$
The morphism $`\varphi `$ is smooth if and only if $`\overline{\varphi }`$ is a smooth morphism of algebraic stacks. Under the identification above the morphism $`\overline{\varphi }`$ becomes
$$\overline{\varphi }:M𝔐𝔬𝔡;(\mathrm{\Lambda }_{𝒪_{F_0}}𝒪_S).$$
Now consider the morphism $`\varphi ^{}:\stackrel{~}{M}N`$ obtained by composing $`\varphi `$ with the automorphism of $`N`$ given by $`A{}_{}{}^{t}A`$; $`\varphi `$ is smooth if and only if $`\varphi ^{}`$ is smooth. Set $`^{}:=Hom_{𝒪_S}(,𝒪_S)`$ which is also naturally an $`𝒪_F_{𝒪_{F_0}}𝒪_S`$-module. By regarding $`\stackrel{~}{M}`$ as the $`GL_r`$-torsor over $`M`$ giving the $`𝒪_S^r`$-trivializations of the dual $`^{}`$ we see that $`\varphi ^{}`$ descends to
$$\overline{\varphi }^{}:M𝔐𝔬𝔡;(\mathrm{\Lambda }_{𝒪_{F_0}}𝒪_S)^{}.$$
As before, it is enough to show that $`\overline{\varphi }^{}`$ is smooth.
The stack $`𝔐𝔬𝔡`$ supports the universal $`𝒪_F_{𝒪_{F_0}}𝒪_{𝔐𝔬𝔡}`$-module $``$. Let us consider the $`𝒪_{𝔐𝔬𝔡}`$-module
$$𝒯=Hom_{𝒪_F_{𝒪_{F_0}}𝒪_{𝔐𝔬𝔡}}(\mathrm{\Lambda }_{𝒪_{F_0}}𝒪_{𝔐𝔬𝔡},)^d.$$
This defines a vector bundle $`𝔙(𝒯^{})`$ over $`𝔐𝔬𝔡`$ (see \[LMoB\] 14.2.6). The structure morphism $`𝔙(𝒯^{})𝔐𝔬𝔡`$ is representable and smooth. Its effect on objects is given by
$$(,\mathrm{\Lambda }_{𝒪_{F_0}}𝒪_S\stackrel{f}{}).$$
Consider now the perfect $`𝒪_{F_0}`$-bilinear pairing
$$(,):𝒪_F\times 𝒪_F𝒪_{F_0};(x,y)=\mathrm{Tr}_{F/F_0}(\delta ^1xy)$$
where $`\delta `$ is an $`𝒪_F`$-generator of the different $`𝒟_{F/F_0}`$. This pairing gives an $`𝒪_F`$-module isomorphism $`𝒪_F^{}=Hom_{𝒪_{F_0}}(𝒪_F,𝒪_{F_0})𝒪_F`$. Now choose an $`𝒪_F`$-module isomorphism $`\mathrm{\Lambda }𝒪_F^d`$. We obtain functorial $`𝒪_F_{𝒪_{F_0}}𝒪_S`$-module isomorphisms $`(\mathrm{\Lambda }_{𝒪_{F_0}}𝒪_S)^{}\mathrm{\Lambda }_{𝒪_{F_0}}𝒪_S`$ and a morphism $`i:M𝔙(𝒯^{})`$ given by
$$(\mathrm{\Lambda }_{𝒪_{F_0}}𝒪_S)(^{},\mathrm{\Lambda }_{𝒪_{F_0}}𝒪_S(\mathrm{\Lambda }_{𝒪_{F_0}}𝒪_S)^{}^{}).$$
We therefore see that $`i`$ is representable and an open immersion (it gives an isomorphism between $`M`$ and the open substack of $`𝔙(𝒯^{})`$ given by the full subcategory of objects for which the morphism $`f`$ is surjective). The morphism $`\overline{\varphi }^{}`$ is the composition
$$M𝔙(𝒯^{})𝔐𝔬𝔡$$
and therefore is smooth. The relative dimension of $`\varphi `$ is equal to the relative dimension of $`\overline{\varphi }`$; this, in turn, is equal to the relative dimension of the composition above. However, this is the same as the relative dimension of the vector bundle $`𝔙(𝒯^{})𝔐𝔬𝔡`$ which is equal to $`rd`$.
###### Remarks 4.2
(i) The proof of Theorem 4.1 given above follows ideas which appear in the paper \[FGKV\] (see especially loc.cit. §4.2) and were brought to the attention of the authors by D. Gaitsgory. A previous version of the paper contained a more complicated proof which used explicit matrix calculations to describe affine charts for the scheme $`M`$.
(ii) Let us consider the conjugation action of $`GL_r`$ on the special fibre $`\overline{N}`$. The orbits of this action are parametrized by
(4.7)
$$𝒮(r,e)=\{𝐬=(s_1\mathrm{}s_r);es_1,s_r0,\mathrm{\Sigma }_is_i=r\}.$$
We denote the corresponding orbit by $`N_𝐬`$. Again $`N_𝐬^{}`$ lies in the closure of $`N_𝐬`$ if and only if $`𝐬^{}𝐬`$. Obviously $`𝒮^0(r,e,d)𝒮(r,e)`$ and the $`𝒢\times GL_r`$-equivariance of the diagram (4.4) shows that for $`𝐬𝒮^0(r,e,d)`$,
(4.8)
$$\varphi (\pi ^1(M_𝐬))=N_𝐬.$$
In particular, the image of $`\stackrel{~}{M}_{𝒪_E}k`$ under $`\varphi `$ is the union of orbits corresponding to $`𝐬𝒮^0(r,e,d)`$ and the diagram (4.4) induces an injection of the set of $`𝒢`$-orbits in $`\overline{M}`$ into the set of $`GL_r`$-orbits in $`\overline{N}`$. It is easy to see that the complement of $`𝒮^0(r,e,d)`$ in $`𝒮(r,e)`$ is closed under the partial order $``$ on $`𝒮(r,e)`$, i.e. corresponds to a closed subset of $`\overline{N}`$.
(iii) The dimension of $`N_𝐬`$ is given by the formula
(4.9)
$$\mathrm{dim}N_𝐬=r^2\underset{i=1}{\overset{e}{}}r_i^2,$$
where $`r_1\mathrm{}r_e0`$ is the dual partition to $`𝐬`$.
For $`𝐬`$ in $`𝒮^0(r,e,d)`$, this formula is compatible with the one of (3.14) via Theorem 4.1: we have
$$𝐬,2\varrho +r^2rd=r^2\underset{i=1}{\overset{e}{}}r_i^2,$$
comp. (3.15) (which is “dual” to (4.9)).
(iv) The equivariance of (4.4) can be rephrased (via descent theory for $`𝒢`$-torsors) by saying that the morphism (4.6) factors through a morphism of algebraic stacks
(4.10)
$$[M/𝒢][N/GL_r].$$
Note that both these stacks have only finitely many points. (The set of points in the special fiber of $`[M/𝒢]`$ resp. $`[N/GL_r]`$ is $`𝒮^0(r,e,d)`$ resp. $`𝒮(r,e)`$.)
###### Corollary 4.3
Let $`r_\phi 1,\phi `$. Then the standard model for $`GL_d`$ corresponding to $`𝐫`$ is flat over $`\mathrm{Spec}𝒪_E`$, with special fiber a normal complete intersection variety.
Proof. In this case the first condition $`A^e=0`$ in the definition (4.5) of $`\overline{N}`$ is a consequence of the second condition. Hence $`\overline{N}`$ is the variety of nilpotent matrices, which is a reduced and irreducible, normal and complete intersection variety. On the other hand, obviously $`𝐫=𝐫_{\mathrm{min}}`$, hence by Proposition 3.2 we have $`\mathrm{dim}N_{𝒪_E}E=\mathrm{dim}N_{𝒪_E}k`$. hence the generic point of $`\overline{N}`$ is the specialization of a point of $`N_{𝒪_E}E`$. Since $`N_{𝒪_E}k`$ is reduced, the flatness of $`N`$ follows from EGA IV 3.4.6.1. By Theorem 4.1 this implies the corresponding assertions for $`M`$.
###### Remark 4.4
Let us fix an isomorphism $`\mathrm{\Lambda }𝒪_F^d`$ and suppose that $`r=d=e`$. Write
$$\mathrm{\Lambda }_{𝒪_{F_0}}𝒪_S(𝒪_F_{𝒪_{F_0}}𝒪_S)^r=1𝒪_S^r\pi 𝒪_S^r\mathrm{}\pi ^{r1}𝒪_S^r.$$
Now let us consider $`𝒪_F_{𝒪_{F_0}}𝒪_S`$-submodules $`\mathrm{\Lambda }_{𝒪_{F_0}}𝒪_S`$ which are locally free $`𝒪_S`$-locally direct summands and are such that the composition
$$F:\mathrm{\Lambda }_{𝒪_{F_0}}𝒪_S\stackrel{pr}{}\pi ^{r1}𝒪_S^d$$
is a surjection (and therefore an isomorphism). The inverse of the isomorphism $`F`$ can be written
$$F^1(\pi ^{r1}v)=(1f_{r1}(v),\pi f_{r2}(v),\mathrm{},\pi ^{r1}f_0(v))$$
where $`f_i:𝒪_S^r𝒪_S^r`$, $`i=0,\mathrm{},r1`$, are $`𝒪_S`$-linear homomorphisms and $`f_0=\mathrm{Id}`$. Recall $`Q(T)=T^e+_{k=0}^{e1}b_kT^k`$. The condition that $``$ is stable under multiplication by $`\pi `$ translates to
$$f_{k+1}b_{rk1}=f_k(f_1b_{r1}f_0),k=1,\mathrm{},r2,$$
$$b_0=f_{r1}(f_1b_{r1}f_0)=0.$$
Therefore, all the $`f_k`$ are determined by $`f_1`$ and we have $`Q(f_1b_{r1}\mathrm{Id})=0`$. We also have $`F(\pi |)F^1=f_1b_{r1}\mathrm{Id}`$. Using these facts, we see that after choosing a basis of $`𝒪_S^r`$, the modules $``$ for which the composition $`F`$ is an isomorphism are in 1–1 correspondence with $`r\times r`$-matrices $`A`$ which satisfy $`Q(A)=0`$. Let $`U`$ denote the open subscheme of $`M`$ whose $`S`$-points correspond to modules $``$ for which the homomorphism $`F`$ above is an isomorphism. In terms of the diagram (4.4), the above implies that, in this case, there is a section $`s:N\stackrel{~}{M}`$ to the morphism $`\varphi :\stackrel{~}{M}N`$ such that the composition $`\pi s:NM`$ is an open immersion identifying $`N`$ with $`U`$. The smoothness of the morphism $`\varphi `$ then amounts to the smoothness of the conjugation action morphism $`N\times GL_rN`$. Also, since $`M`$ is projective, this shows that, in this case, the local model $`M`$ may be considered as a relative compactification over $`\mathrm{Spec}𝒪_E`$of the variety $`N`$. This last result for the special fibers is precisely the scheme-theoretic version of Lusztig’s result \[L\], section 2. Hence the special fiber of any standard model for $`GL_d`$, for which $`r=d`$ and $`e=d`$ (they are all identical), may be considered as a compactification of the nilpotent variety. We return in section 6 to the consequences of Theorem 4.1 for the special fibers.
## 5 The canonical flat model
We have seen in Corollary 3.3 that a standard model is rarely flat over $`\mathrm{Spec}𝒪_E`$. By Theorem 4.1 the same can be said of the scheme $`N`$. In this section we will first show that the flat scheme theoretic closure of the generic fiber $`N_{𝒪_E}E`$ in $`N`$ has good singularities. The idea is to use a variant of the Springer resolution of the nilpotent variety, as also in the work of Eisenbud and Saltman (\[E-S\]).
Recall our notations from the beginning of section 2. Let $`K`$ be the Galois hull of $`F`$ inside $`F_0^{\mathrm{sep}}`$. Let us order the different embeddings $`\varphi :FK`$. Then we can write
(5.1)
$$P(T)=\underset{\varphi }{}(T\varphi (\pi ))^{r_\varphi }=\underset{i=1}{\overset{e}{}}(Ta_i)^{r_i},Q(T)=\underset{i=1}{\overset{e}{}}(Ta_i).$$
Here $`\varphi (\pi )=a_i𝒪_K`$ are distinct roots.
Let us set $`n_k=_{i=1}^kr_i`$, for $`1ke`$. Let $`𝓕`$ be the scheme which classifies flags:
(5.2)
$$(0)=_e_{e1}\mathrm{}_0=𝒪_S^r$$
where $`_k`$ is locally on $`S`$ a direct summand of $`𝒪_S^r`$ of corank $`n_k`$. Following \[E-S\], we consider the subscheme $`𝒩`$ of $`((\mathrm{Mat}_{r\times r})\times 𝓕)_{𝒪_K}`$ classifying pairs $`(A,\{_{}\})`$ such that
(5.3)
$$(Aa_kI)_{k1}_k,1ke.$$
The scheme $`𝒩`$ supports an action $`GL_r`$ by $`g(A,\{_i\})=(gAg^1,\{g(_i)\})`$.
Obviously this is a variant of the Grothendieck-Springer construction. It differs from the original in two aspects: we consider partial flags instead of complete flags, and we fix the (generalized) eigenvalues $`a_i`$ of $`A`$.
###### Lemma 5.1
i) $`𝒩`$ is smooth over $`\mathrm{Spec}𝒪_K`$.
ii) There is a projective $`GL_r`$-equivariant morphism $`\mu :𝒩N_{𝒪_E}𝒪_K`$, given by $`(A,\{_i\})A`$.
iii) The morphism $`\mu _{\mathrm{Spec}K}`$ is an isomorphism between the generic fibers $`𝒩_{𝒪_K}K`$ and $`N_{𝒪_E}K`$.
Proof. Part (i) follows from the fact that the projection to the second factor is smooth
(5.4)
$$𝒩𝓕,$$
comp. \[E-S\], p. 190 (the fiber over $`\{_i\}`$ can be identified with the cotangent space of $`𝓕`$ at $`\{_i\}`$). Now suppose that $`A`$ is in $`M_{r\times r}(R)`$ for some $`𝒪_K`$-algebra $`R`$ and that locally on Spec $`R`$ there is a filtration $`\{_{}\}`$ of $`R^r`$ as described above. Then the characteristic polynomial of $`A`$ is equal to $`P(T)`$ and we have
(5.5)
$$\underset{k=1}{\overset{N}{}}(Aa_kI)=0M_{r\times r}(R).$$
This implies $`Q(A)=0M_{r\times r}(R)`$, cf. (5.1). This shows that the natural morphism $`𝒩(\mathrm{Mat}_{r\times r})_{𝒪_K}`$ factors through $`N_{𝒪_E}𝒪_K`$; the claim (ii) follows. Now $`N_{𝒪_E}K`$ consists of all semisimple matrices with eigenvalues $`a_i`$ with multiplicity $`r_i`$ (comp. remarks before (4.5)). Hence the fiber of $`\mu _{\mathrm{Spec}K}`$ over $`A`$ is the filtration $`_{}`$ associated to the eigenspace decomposition corresponding to $`A`$, i.e., is uniquely determined by $`A`$.
Now set $`N^{}=\mathrm{Spec}(\mu _{}(𝒪_𝒩))`$ and denote by $`\mu (𝒩)`$ the scheme-theoretic image of $`\mu :𝒩N_{𝒪_E}𝒪_K`$; this is a reduced closed subscheme of $`N_{𝒪_E}𝒪_K`$ since $`𝒩`$ is smooth (therefore reduced) and $`\mu `$ is proper. The scheme $`N^{}`$ supports an action of $`GL_r`$; there is a natural $`GL_r`$-equivariant morphism $`q^{}:N^{}N_{𝒪_E}𝒪_K`$. The morphism $`q^{}`$ is finite and factors as follows
(5.6)
$$q^{}:N^{}\mu (𝒩)N_{𝒪_E}𝒪_K.$$
Lemma 5.1 (iii) implies that $`q_{\mathrm{Spec}K}^{}`$ is an isomorphism.
Let $`k^{}`$ be the residue field of $`𝒪_K`$.
###### Proposition 5.2
a) The special fiber $`\overline{N}^{}=N^{}_{𝒪_K}k^{}`$ is normal and has rational singularities.
b) The scheme $`N^{}`$ is normal, Cohen-Macaulay and flat over $`\mathrm{Spec}𝒪_K`$.
c) $`N^{}=\mu (𝒩)`$.
d) $`N^{}`$ is the scheme theoretic closure of $`N_{𝒪_E}K`$ in $`N_{𝒪_E}𝒪_K`$. Its special fiber is the reduced closure of the orbit $`N_𝐭`$, where $`𝐭=𝐫^{}`$ is the dual partition to $`𝐫=(r_i)`$.
Proof. This follows closely the arguments of \[E-S\] (see p. 190-192, proof of Theorem 2.1) with new input the results of Mehta-van der Kallen (\[M-vdK\]). They show (using Frobenius splitting) that the closure of the orbit of a nilpotent matrix is normal and Cohen-Macaulay also in positive characteristic. For the duration of this proof, we will denote by $`\varpi `$ a uniformizer of $`𝒪_K`$. We will first consider the situation over $`\mathrm{Spec}k^{}`$ and use a bar to denote base change from $`\mathrm{Spec}𝒪_K`$ to $`\mathrm{Spec}k^{}`$. Consider the morphism $`\overline{\mu }:\overline{𝒩}\overline{N}`$ and the scheme $`\mathrm{Spec}(\overline{\mu }_{}(𝒪_{\overline{𝒩}}))\overline{N}`$. The morphism $`\overline{\mu }`$ factors as
(5.7)
$$\overline{\mu }:\overline{𝒩}\mathrm{Spec}(\overline{\mu }_{}(𝒪_{\overline{𝒩}}))\overline{\mu }(\overline{𝒩})\overline{N},$$
where again $`\overline{\mu }(\overline{𝒩})`$ denotes the scheme theoretic image of $`\overline{\mu }`$. The morphism $`\overline{𝒩}\overline{\mu }(\overline{𝒩})`$ is one to one on the open subset of $`\overline{𝒩}`$ of those $`(A,\{_{}\})`$ such that $`A_{k1}=_k`$, $`k=1,\mathrm{},e`$ and so it is birational; therefore the morphism $`\mathrm{Spec}(\overline{\mu }_{}(𝒪_{\overline{𝒩}}))\overline{\mu }(\overline{𝒩})`$ is finite and birational. Since $`\overline{𝒩}`$ is reduced, $`\overline{\mu }(\overline{𝒩})`$ is also reduced. As in \[E-S\], we see that $`\overline{\mu }(\overline{𝒩})(\mathrm{Mat}_{r\times r})_k^{}`$ is the reduced closure of the conjugation orbit $`N_𝐭`$ of the Jordan form for the dual partition $`𝐭=𝐫^{}`$. By \[M-vdK\], the closure of $`N_𝐭`$ is normal and has rational singularities. We conclude that $`\mathrm{Spec}(\overline{\mu }_{}(𝒪_{\overline{𝒩}}))\overline{\mu }(\overline{𝒩})`$ is an isomorphism which we use to identify these two schemes. We will first show that $`\mathrm{Spec}(\overline{\mu }_{}(𝒪_{\overline{𝒩}}))=\overline{\mu }(\overline{𝒩})`$ actually gives the special fiber of $`N^{}`$. The statement (a) then follows from the results of Mehta-van der Kallen. The cohomology exact sequence obtained by applying $`\mu _{}`$ to
$$0𝒪_𝒩\stackrel{\varpi }{}𝒪_𝒩𝒪_{\overline{𝒩}}0$$
gives an injective homomorphism
$$𝒪_N^{}/\varpi 𝒪_N^{}=\mu _{}(𝒪_𝒩)/\varpi \mu _{}(𝒪_𝒩)\overline{\mu }_{}(𝒪_{\overline{𝒩}})$$
and it is enough to show that this is an isomorphism. There is a commutative diagram
$$\begin{array}{ccc}𝒪_{\mu (𝒩)}/\varpi 𝒪_{\mu (𝒩)}& & 𝒪_{\overline{\mu }(\overline{𝒩})}\\ & & \\ 𝒪_N^{}/\varpi 𝒪_N^{}& & \overline{\mu }_{}(𝒪_{\overline{𝒩}}).\end{array}$$
From the definition of the scheme theoretic image, the upper horizontal homomorphism is an isomorphism. On the other hand, we have seen above that the right vertical homomorphism is an isomorphism, whence the claim. We will now show that $`N^{}\mu (𝒩)`$ is an isomorphism; this will establish (c). From the above it follows that the homomorphism
$$𝒪_{\mu (𝒩)}/\varpi 𝒪_{\mu (𝒩)}𝒪_N^{}/\varpi 𝒪_N^{}$$
is surjective. Since $`\mu `$ is proper, $`𝒪_N^{}=\mu _{}(𝒪_𝒩)`$ is finite over $`𝒪_{\mu (𝒩)}`$. Therefore, using Nakayama’s lemma, we conclude that $`𝒪_{\mu (𝒩)}𝒪_N^{}`$ is surjective locally over all points at the special fiber. Since $`N^{}_{𝒪_K}K\mu (𝒩)_{𝒪_K}K`$ is an isomorphism it follows that $`𝒪_{\mu (𝒩)}𝒪_N^{}`$ is surjective; it now follows from the definition of the scheme-theoretic image that $`𝒪_{\mu (𝒩)}𝒪_N^{}`$ is an isomorphism. This shows (c).
Now let us show part (b). By (a) the special fiber $`\overline{N}^{}`$ is normal and Cohen-Macaulay; in fact, it has dimension $`r^2_ir_i^2`$. This is equal to the dimension of the generic fiber $`N_{𝒪_E}K=N^{}_{𝒪_K}K`$. As a result the special fiber is reduced and its unique generic point lifts to the generic fiber; this implies that $`N^{}`$ is flat over $`\mathrm{Spec}𝒪_K`$ (EGA IV 3.4.6.1). Since $`\overline{N}^{}`$ is Cohen-Macaulay and $`N^{}\mathrm{Spec}𝒪_K`$ is flat, $`N^{}`$ is Cohen-Macaulay. Now $`N^{}_{𝒪_K}K=N_{𝒪_E}K`$ is smooth and $`\overline{N}^{}`$ generically smooth; this shows that $`N^{}`$ is regular in codimensions $`0`$ and $`1`$ and therefore, by Serre’s criterion, normal.
Finally, since $`N^{}=\mu (𝒩)N_{𝒪_E}𝒪_K`$ with identical generic fibers, and since $`N^{}`$ is flat over $`\mathrm{Spec}𝒪_K`$, $`N^{}`$ is the (flat) scheme theoretic closure of $`N_{𝒪_E}K`$ in $`N_{𝒪_E}𝒪_K`$. This shows (d).
###### Proposition 5.3
Let $`N^{\mathrm{loc}}`$ be the (flat) scheme theoretic closure of $`N_{𝒪_E}E`$ in $`N`$. Then the scheme $`N^{\mathrm{loc}}`$ is normal and Cohen-Macaulay. Its special fiber is the reduced closure of the orbit $`N_𝐭`$ with $`𝐭=𝐫^{}`$ the dual partition to $`𝐫=(r_\phi )`$. The special fiber is normal with rational singularities.
Proof. Denote by $`N^{\prime \prime }N`$ the scheme theoretic image of the finite composite morphism $`N^{}N_{𝒪_E}𝒪_KN`$. This is a $`GL_r`$-equivariant closed subscheme of $`N`$. We have $`𝒪_{N^{\prime \prime }}𝒪_N^{}`$ and so since $`N^{}`$ is flat over $`\mathrm{Spec}𝒪_K`$, we conclude that $`N^{\prime \prime }`$ is also flat over $`\mathrm{Spec}𝒪_E`$. We have $`N^{\prime \prime }_{𝒪_E}E=N_{𝒪_E}E`$; hence $`N^{\prime \prime }`$ is the flat scheme theoretic closure of $`N_{𝒪_E}E`$ in $`N`$, that is $`N^{\prime \prime }=N^{\mathrm{loc}}`$. The base change $`N^{\mathrm{loc}}_{𝒪_E}𝒪_K`$ is a closed subscheme of $`N_{𝒪_E}𝒪_K`$ which is flat over $`\mathrm{Spec}𝒪_K`$. Since we have $`N^{\mathrm{loc}}_{𝒪_E}K=N_{𝒪_E}K`$ by Proposition 5.2 (d) we have
(5.8)
$$N^{\mathrm{loc}}_{𝒪_E}𝒪_K=N^{}N_{𝒪_E}𝒪_K.$$
By Proposition 5.2 (b) $`N^{}`$ is normal. Therefore, $`N^{\mathrm{loc}}=N^{\prime \prime }/\mathrm{Gal}(K/E)`$ is also normal. Denoting by $`\overline{N}^{\mathrm{loc}}`$ the special fiber of $`N^{\mathrm{loc}}`$, we have
$$\overline{N}^{\mathrm{loc}}_kk^{}=\overline{N}^{}.$$
Hence the remaining assertions also follow from Proposition 5.2.
We note that we can define analogues of $`𝒩`$ for $`M`$ and $`\stackrel{~}{M}`$. Namely, we consider the $`𝒪_K`$-scheme $``$, which for an $`𝒪_K`$-scheme $`S`$ classifies the filtrations of $`𝒪_S`$-submodules
(5.9)
$$\{(0)=_e_{e1}\mathrm{}_0=\mathrm{\Lambda }_{𝒪_{F_0}}𝒪_S\},$$
where $`M^{\mathrm{loc}}(S)`$ and where $`_k`$ is locally on $`S`$ a direct summand of $``$ of corank $`n_k=_{i=1}^kr_i`$ such that
(5.10)
$$(\pi a_k\mathrm{Id})_{k1}_k,k=1,\mathrm{},e.$$
Similarly we define $`\stackrel{~}{}`$ by fixing in addition an isomorphism $`\psi :𝒪_S^r`$. We thus obtain a diagram with cartesian squares in which the vertical morphisms are projective with source a smooth $`𝒪_K`$-scheme,
(5.11)
$$\begin{array}{ccccc}& & \stackrel{~}{}& & 𝒩\\ & & & & \\ M_{𝒪_E}𝒪_K& & \stackrel{~}{M}_{𝒪_E}𝒪_K& & N_{𝒪_E}𝒪_K& .\end{array}$$
We now use Theorem 4.1 to transfer the previous results to the standard models.
###### Theorem 5.4
Let $`M^{\mathrm{loc}}`$ be the (flat) scheme theoretic closure of $`M_{𝒪_E}E`$ in $`M`$. Then
(i) $`M^{\mathrm{loc}}`$ is normal and Cohen-Macaulay.
(ii) The special fiber $`\overline{M}^{\mathrm{loc}}`$ is reduced, normal with rational singularities, and we have
$$\overline{M}^{\mathrm{loc}}=\underset{𝐬𝐭}{}M_𝐬^{\mathrm{loc}}.$$
(iii) There is a diagram with cartesian squares of $`𝒢\times GL_r`$-equivariant morphisms in which the horizontal morphisms are smooth and the vertical morphisms are closed embeddings,
$$\begin{array}{ccccc}M^{\mathrm{loc}}& \stackrel{\pi ^{\mathrm{loc}}}{}& \stackrel{~}{M}^{\mathrm{loc}}& \stackrel{\varphi ^{\mathrm{loc}}}{}& N^{\mathrm{loc}}\\ & & & & \\ M& \stackrel{\pi }{}& \stackrel{~}{M}& \stackrel{\varphi }{}& N& .& & \multicolumn{-6}{c}{}\end{array}$$
The $`𝒪_E`$-scheme $`M^{\mathrm{loc}}`$ is called the canonical model associated to the standard model for $`GL_d`$ corresponding to $`𝐫=(r_i)`$.
###### Remark 5.5
The use of Theorem 4.1 in proving parts (i) and (ii) of the above theorem is to enable us to appeal to the results of Mehta-van der Kallen on nilpotent orbit closures. An alternative approach which would not appeal to Theorem 4.1 might be obtained by studying directly the composite morphism
$$M_{𝒪_E}𝒪_KM$$
and using the theory of generalized Schubert varieties in the affine Grassmannian for $`GL_d`$. One can show directly that $``$ is smooth over $`\mathrm{Spec}𝒪_K`$ (comp. the proof of Lemma 5.1); the main point is that the special fiber $`\overline{}`$ can be written as a composite of smooth fibrations with fibers Grassmannian varieties (indeed, we may think of $`\overline{}`$ as a generalized Demazure-Bott-Samelson variety). Then, the arguments in the proofs of Propositions 5.2 and 5.3 can be repeated to obtain a direct proof of parts (i) and (ii) of Theorem 5.4, provided we know that the reduced closure of the stratum $`M_𝐭`$ in the special fiber $`\overline{M}`$ is normal and has rational singularities. By the discussion in §3, this reduced closure is isomorphic to the reduced closure $`X_𝐭:=\overline{𝒪}_𝐭`$ of the Schubert cell $`𝒪_𝐭`$ in the affine Grassmannian (defined as in \[BL\], comp. §3). It remains to show that the generalized Schubert variety $`X_𝐭`$ is normal and has rational singularities. Results of this type have been shown (in positive characteristic) by Mathieu \[Mat\]. Unfortunately, it is not clear that in positive characteristic the Schubert varieties he considers have the same scheme structure as the $`X_𝐭`$ (see the remarks on p. 410 of \[BL\]). Since this point is not cleared up, we use Theorem 4.1 to deduce results on the affine Grassmannian from results on nilpotent orbit closures, comp. section 6 below.
Suppose now that $`(|\mathrm{\Gamma }|,\mathrm{char}k)=1`$. Under this hypothesis we will show that a conjecture of de Concini and Procesi implies a rather explicit description of the scheme $`M^{\mathrm{loc}}`$.
We follow \[deC-P\] §1. Let $`x_1`$, $`\mathrm{}`$, $`x_r`$ be a set of variables; for every pair of integers $`t`$, $`h`$ with $`h0`$, $`1tr`$, we can consider the total symmetric function of degree $`h`$ in the first $`t`$ variables; this is defined to be the sum of all monomials in $`x_1`$, $`\mathrm{}`$, $`x_t`$ of degree $`h`$ and will be denoted by $`S_h^t(x_i)`$. We also indicate by the symbols $`\sigma _h`$ the elementary symmetric function of degree $`h`$ in the variables $`x_1`$, $`\mathrm{}`$, $`x_r`$ with the convention that $`\sigma _h=0`$ if $`h>r`$. Write
(5.12)
$$S_h^t(x_i)=a_{(h_1,\mathrm{},h_t)}x_1^{h_1}\mathrm{}x_t^{h_t}$$
For $`A\mathrm{Mat}_{r\times r}(R)=\mathrm{End}(R^r)`$, we set
(5.13)
$$S_h^t(A)(e_{i_1}\mathrm{}e_{i_t})=a_{(h_1,\mathrm{},h_t)}A^{h_1}e_{i_1}\mathrm{}A^{h_t}e_{i_t}$$
Since $`S_h^t(x_i)`$ is symmetric this defines a $`R`$-linear operator
(5.14)
$$S_h^t(A):^t(R^r)^t(R^r)$$
Now let us indicate by
(5.15)
$$T^r\sigma _1(A)T^{r1}+\sigma _2(A)T^{r2}\mathrm{}+(1)^r\sigma _r(A)$$
the characteristic polynomial $`\mathrm{det}(TIA)`$ of $`A`$. For $`t`$, $`h`$ as above, we now define the following element of $`\mathrm{End}(^t(R^r))`$:
$$F_h^t(A):=S_h^t(A)\sigma _1(A)S_{h1}^t(A)+\mathrm{}+(1)^h\sigma _h(A)$$
For each function $`f:\{1,2,\mathrm{},e\}𝐍`$ with $`0f(i)r_i`$, for all $`1ie`$, consider
(5.16)
$$Q_f(t)=\underset{i=1}{\overset{e}{}}(Ta_i)^{f(i)}$$
(a divisor of the polynomial $`P(T)`$).
Let us consider the subscheme $`N_0`$ of $`\mathrm{Mat}_{r\times r}`$ over $`\mathrm{Spec}𝒪_E`$ which is defined by the equations given by the conditions
(5.17)
$$\mathrm{det}(TIA)P(T),\mathrm{and}\underset{\sigma \mathrm{\Gamma }}{}F_h^t(A)^t(\sigma Q_f(A))=0,$$
$$\mathrm{for}\mathrm{all}f,\mathrm{and}\mathrm{for}t+h=r\underset{i,f(i)0}{}r_i+1,t1,h0.$$
(it is obvious that the generators of the ideal defining $`N_0`$ have all coefficients in $`𝒪_E`$). This is actually a closed subscheme of $`N`$; indeed, consider the second set of equations for $`f1`$, $`t=1`$, $`h=0`$. Since $`F_0^t(A)=I`$, we obtain
(5.18)
$$|\mathrm{\Gamma }|\underset{i=1}{\overset{e}{}}(Aa_iI)=0.$$
Since $`(|\mathrm{\Gamma }|,\mathrm{char}k)=1`$, this equation implies $`Q(A)=0`$. It is straightforward to see that $`N_0`$ is a $`GL_r`$-invariant subscheme of $`N`$.
###### Proposition 5.6
i) The schemes $`N`$ and $`N_0`$ over $`\mathrm{Spec}𝒪_E`$ have the same generic fiber.
ii) The reduced special fiber of $`N_0`$ is equal to the reduced closure of the conjugation orbit $`N_𝐭`$ of the Jordan form with partition $`𝐭=𝐫^{}`$.
Proof. We first show (i). By descent, we can check this after base changing to $`K`$. Consider the diagonal $`r\times r`$ matrix $`A_0`$ in which the element $`a_i`$ appears with multiplicity $`r_i`$. We have
(5.19)
$$\mathrm{rank}(Q_f(A_0))=r\underset{i,f(i)0}{}r_i.$$
It now follows from \[deC-P\], Proposition on p. 206, that
(5.20)
$$F_h^t(A_0)^t(Q_f(A_0))=0$$
if $`t+hr_{i,f(i)0}r_i+1`$ (comp. the proof of the Theorem on p. 207 loc. cit). This shows that $`A_0`$ satisfies the conditions defining $`N_0_{𝒪_E}K`$. Since $`N_{𝒪_E}K`$ is the reduced $`GL_r`$-orbit of such a matrix and $`N_0`$ is $`GL_r`$-equivariant, the result follows.
Now we prove (ii). Reducing the equations defining $`N_0`$ modulo the maximal ideal of $`𝒪_E`$ gives the following equations:
(5.21)
$$\mathrm{det}(TIA)T^r,\mathrm{and}|\mathrm{\Gamma }|F_h^t(A)^t(A^{_if(i)})=0$$
$$\mathrm{for}\mathrm{all}f,\mathrm{and}\mathrm{for}t+h=r\underset{i,f(i)0}{}r_i+1,t1,h0.$$
By the definition of $`F_h^t(A)`$, the difference $`F_h^t(A)S_h^t(A)`$ is in the ideal given by $`\mathrm{det}(TIA)T^r`$. Using this and taking $`J=\mathrm{supp}(f)I=\{1,\mathrm{},e\}`$ we see that the above set of equations generates the same ideal as the following one:
(5.22)
$$\mathrm{det}(TIA)T^r,\mathrm{and}S_h^t(A)^t(A^{|J|})=0$$
$$\text{for all subsets}JI,t+h=r\underset{iJ}{}r_i+1,t1,h0.$$
If the $`r_i`$ are arranged in decreasing order, it is the same to consider
(5.23)
$$\mathrm{det}(TIA)T^r,\mathrm{and}S_h^t(A)^t(A^k)=0$$
$$k=0,\mathrm{},e,t+h=rn_k+1,t1,h0.$$
(recall our notation $`n_k=_{i=1}^kr_i`$, $`n_0=0`$). In fact, we could also omit the first equation and write this as
(5.24)
$$S_h^t(A)^t(A^k)=0$$
$$\text{ for }k=0,\mathrm{},e,t+h=rn_k+1,t1,h0.$$
Indeed, by \[deC-P\], Lemma p. 206, the ideal generated by $`S_h^t(A)`$ for $`t+h=r+1`$ is the same as the one generated by the coefficients of the characteristic polynomial $`\sigma _1(A)`$, $`\mathrm{}`$, $`\sigma _r(A)`$.
Let us denote by $``$ the ideal of $`k[a_{ij}]_{1i,jr}`$ generated by the equations (5.24). Then it follows as in \[deC-P\], Theorem on p. 207, that its radical $`\mathrm{rad}()`$ defines the (reduced) closure of the orbit $`N_𝐭`$ of the Jordan form for the dual partition $`𝐭=𝐫^{}`$. This shows (ii).
De Concini and Procesi conjecture (loc. cit.) that when char $`k=0`$, we have $`\mathrm{rad}()=`$. Assume that this is true even if char $`k>0`$. Then
(5.25)
$$N_0_{𝒪_E}k=\overline{N}_𝐭,$$
(i.e, the special fiber of $`N_0`$ is already reduced). Since by \[M-vdK\] the reduced closure of $`N_𝐭`$ is normal, Cohen-Macaulay and has dimension $`r^2_ir_i^2`$, an argument as in the proof of Proposition 5.2 (ii) shows that $`N_0`$ is normal and flat over $`\mathrm{Spec}𝒪_E`$. It follows that $`N_0`$ is the scheme theoretic closure of $`N_{𝒪_E}E`$ in $`N`$, i.e. $`N_0=N^{\mathrm{loc}}`$. Transferring this result to the standard model we obtain the following result.
###### Theorem 5.7
Assume that $`(|\mathrm{\Gamma }|,\mathrm{char}k)=1`$ and that the conjecture of De Concini-Procesi holds over $`k`$. Then the canonical flat model $`M^{\mathrm{loc}}`$ can be described as the moduli scheme for the following moduli functor on $`(\mathrm{Sch}/𝒪_E)`$: The $`S`$-valued points are given by $`𝒪_F_{𝒪_{F_0}}𝒪_S`$-submodules $``$ of $`\mathrm{\Lambda }_{𝒪_{F_0}}𝒪_S`$ which are locally direct summands as $`𝒪_S`$-modules, with
$$\mathrm{det}(T\mathrm{Id}\pi |)\underset{\phi }{}(T\phi (\pi ))^{r_\phi }$$
and
$$\underset{\sigma \mathrm{\Gamma }}{}F_h^t(\pi |)^t(\sigma Q_f(\pi |))=0,$$
for all $`f`$ as before, and for $`t+h=r_{\phi ,f(\phi )0}r_\phi +1`$, $`t1,h0`$.
Let us fix integers $`d,r`$ and $`e`$ and let us state the following conjecture.
###### Conjecture 5.8
Consider the closed subscheme $`\overline{N}`$ of $`\mathrm{Mat}_{r\times r}`$ over $`k`$ defined in (4.5),
$$\overline{N}=\{A\mathrm{Mat}_{r\times r};A^e=0,\mathrm{det}(TIA)T^r\}.$$
Then this scheme is reduced.
Note that when $`re`$, the first condition describing $`\overline{N}`$ follows from the second (Cayley-Hamilton) and in this case the statement is a classical theorem of Kostant on the nilpotent variety. In a companion paper to ours, J. Weyman proves this conjecture when $`e=2`$, and for arbitrary $`e`$ when char $`k=0`$. We note that the reduced subscheme of $`\overline{N}`$ is simply the orbit closure corresponding to the partition $`(e^c,f)`$ of $`r`$, where we have written as usual $`r=ce+f`$, $`0f<e`$. Now De Concini and Procesi \[deC-P\] have determined the length of the intersection of any nilpotent orbit closure with the diagonal matrices (in arbitrary characteristic). Therefore, if Conjecture 5.8 holds we would obtain from their formula that
(5.26)
$$\mathrm{dim}_kk[X_1,\mathrm{},X_r]/(e_1,\mathrm{},e_r,X_1^e,\mathrm{},X_r^e)=\frac{r!}{((c+1)!)^f(c!)^{ef}}.$$
Here $`e_1,\mathrm{},e_r`$ are the elementary symmetric functions in the indeterminates $`X_1,\mathrm{},X_r`$.
###### Corollary 5.9
Let $`𝐫=𝐫_{\mathrm{min}}(r,e)`$. The corresponding standard model for $`GL_d`$ is flat over $`\mathrm{Spec}𝒪_E`$ if Conjecture 5.8 holds true.
Proof. If the conjecture holds true, we conclude from Theorem 4.1 that $`M_{𝒪_E}k`$ is reduced. On the other hand, by Proposition 3.2 the generic and the special fiber of $`M`$ are irreducible of the same dimension. The flatness of $`M`$ follows as in the proof of Proposition 5.2 from EGA IV.3.4.6.1.
###### Corollary 5.10
Assume that $`e=2`$ and order the embeddings so that $`r_1r_2`$. Then the canonical flat model $`M^{\mathrm{loc}}`$ represents the following moduli problem on $`(Sch/𝒪_E)`$: The $`S`$-valued points are given by $`𝒪_F_{F_0}𝒪_S`$-submodules $``$ of $`\mathrm{\Lambda }_{𝒪_{F_0}}𝒪_S`$ which are locally direct summands as $`𝒪_S`$-modules with
$$\mathrm{det}(T\mathrm{Id}\pi |)(T\phi _1(\pi ))^{r_1}(T\phi _2(\pi ))^{r_2}$$
and
$$^{r_2+1}(\pi \phi _1(\pi )\mathrm{Id}|)=0\text{if}r_1>r_2.$$
Here in the last case we used $`\phi _1`$ to identify $`F`$ with $`E`$.
Proof. In this case $`E`$ is a Galois extension, namely $`E=F_0`$ if $`r_1=r_2`$ and $`E=F`$ via $`\phi _1`$ if $`r_1>r_2`$. Let us discuss the case when $`r_1>r_2`$. In this case $`\mathrm{\Gamma }`$ is trivial, and the second identity above is the one in Theorem 5.7 corresponding to $`f`$ with $`f(1)=1`$, $`f(2)=0`$ and to $`t=r_2+1`$, $`h=0`$. It follows that the conditions above define a closed subscheme $`M^{}`$ of $`M`$ and it is easy to see that the generic fibers coincide. The special fiber of $`M^{}`$ is defined by the conditions
(5.27)
$$\mathrm{det}(T\mathrm{Id}\pi |)T^r,^{r_2+1}(\pi |)=0.$$
Consider the closed subscheme $`N^{}`$ of $`N`$ defined by the condition $`^{r_2+1}A=0`$. Then the special fiber of $`N^{}`$ is given as
(5.28)
$$\overline{N}^{}=\{A\mathrm{Mat}_{r\times r};A^2=0,^{r_2+1}A=0,\mathrm{det}(TIA)T^r\}.$$
We then obtain a diagram of morphisms
(5.29)
$$M^{}\stackrel{\pi ^{}}{}\stackrel{~}{M}^{}\stackrel{\varphi ^{}}{}N^{}$$
in which $`\pi ^{}`$ is a $`GL_r`$-torsor and $`\varphi ^{}`$ is smooth of relative dimension $`rd`$. According to Strickland (\[St\]), the special fiber of $`N^{}`$ is irreducible and reduced and in fact equal to the reduced closure of the nilpotent orbit corresponding to the partition $`𝐬=(2^{r_2},1^{r_1r_2})`$, hence of dimension $`r^2(r_1^2+r_2^2)`$. It follows that the special fiber of $`M^{}`$ is irreducible and reduced of dimension $`dr(r_1^2+r_2^2)=\mathrm{dim}M^{}_{𝒪_E}E`$. Hence $`M^{}`$ is $`𝒪_E`$-flat and therefore coincides with $`M^{\mathrm{loc}}`$.
Now assume $`r_1=r_2`$. In this case we are asserting that the standard model is flat over $`𝒪_{F_0}`$. This follows via Corollary 5.9 from Weyman’s result on Conjecture 5.8 concerning $`e=2`$.
Remark: The second identity above first appeared in \[P\] with the purpose of defining a flat local model for the unitary group corresponding to a ramified quadratic extension of $`𝐐_p`$. It was the starting point of the present paper.
###### Remark 5.11
Corollary 5.10 illustrates the fact that the identities in Theorem 5.7 are extremely redundant. It is an open problem to find a shorter list of identities, with coefficients in $`𝒪_E`$, in suitable tensor powers of $`\pi |`$ which describe the canonical flat model.
## 6 Applications to the affine Grassmannians
In this section we spell out some consequences of the results of the previous sections, as they pertain to the special fibers. For the first application we take up again the notation of section 3; in particular, we denote for a dominant coweight $`𝐬`$ of $`GL_d`$, by $`𝒪_𝐬`$ the corresponding orbit of $`\stackrel{~}{𝒢}`$ on $`\stackrel{~}{\mathrm{Grass}}_k`$.
###### Theorem 6.1
The reduced closure $`\overline{𝒪}_𝐬`$ is normal with rational singularities. Its singularities are in fact smoothly equivalent to singularities occurring in nilpotent orbit closures for a general linear group.
Proof. Let $`𝐬=(s_1,\mathrm{},s_d)`$. After translation by a scalar matrix we may assume $`s_d0`$. Let $`e`$ be any integer $`s_1`$, and put $`r=s_i`$. Then $`𝐬𝒮^0(r,e,d)`$. Hence $`𝒪_𝐬`$ may be identified with the corresponding stratum $`M_𝐬`$ of the special fiber $`\overline{M}`$ of any standard model $`M(\mathrm{\Lambda },𝐫)`$, where $`[F:F_0]=e`$ and $`\mathrm{\Lambda }=𝒪_F^d`$ and $`_{j=1}^er_j=r`$ (all of them have identical special fibers). And the closure $`\overline{𝒪}_𝐬`$ of $`𝒪_𝐬`$ can be identified with the closure $`\overline{M}_𝐬`$ of $`M_𝐬`$ in $`\overline{M}`$. From Theorem 4.1 it follows that $`\overline{M}_𝐬`$ is smoothly equivalent to the closure of the corresponding nilpotent orbit $`N_𝐬`$ of $`\overline{N}`$, which by Mehta - van der Kallen is normal with rational singularities.
###### Remark 6.2
Results of this type (normality of Schubert varieties) have been shown in positive characteristic in the context of Kac-Moody algebras by Mathieu \[Mat\]. However, it is not clear whether the Schubert varieties he considers have the same scheme structure as the $`\overline{𝒪}_𝐬`$ considered here (the corresponding statement is known in characteristic zero, by the integrability result of Faltings, comp. \[BL\], app. to section 7). It follows from the methods of Görtz \[G\] that, once a statement of the kind of Theorem 6.1 is known, it follows that any Schubert variety in any parahoric flag variety for $`GL_d(k((\mathrm{\Pi })))`$ is normal (and much more, comp. \[G1\]). For a completely different approach see Faltings’s paper \[F\].
In fact, we can be more precise than in Theorem 6.1. Let $`(V_e,\mathrm{\Pi }_e)`$ be the standard vector space of dimension $`e`$ over $`k`$, with the standard regular nilpotent endomorphism. For $`d1`$ let
(6.1)
$$(W,\mathrm{\Pi })=(V_e,\mathrm{\Pi }_e)^d.$$
Let $`0red`$ and consider the projective scheme $`X=X(r,e,d)`$ over $`k`$ which represents the following functor on the category of $`k`$-algebras. It associates to a $`k`$-algebra $`R`$ the set of $`R`$-submodules,
$`\{W_kR;`$ $`\text{is locally on }\mathrm{Spec}R\text{ a direct summand,}`$
$`\text{is }\mathrm{\Pi }\text{-stable and}`$
$`\mathrm{det}(T\mathrm{\Pi }|)T^r\}.`$
By the end of section 3, $`X`$ is the special fiber of any standard model $`M(\mathrm{\Lambda },𝐫)`$ where $`[F:F_0]=e`$, where $`\mathrm{\Lambda }=𝒪_F^d`$ and $`r=r_\phi `$. We obtain as a special fiber of the diagram (4.4) the diagram
(6.2)
$$X\stackrel{\pi }{}\stackrel{~}{X}\stackrel{\varphi }{}\overline{N},$$
in which $`\pi `$ is a $`GL_r`$-torsor and $`\varphi `$ is smooth of relative dimension $`rd`$. It is equivariant with respect to the action of the product group $`\overline{𝒢}\times GL_r`$, where $`\overline{𝒢}=_{k[\mathrm{\Pi }]/(\mathrm{\Pi }^e)/k}(GL_d)`$, cf. (3). We therefore obtain an extremely close relationship between affine Grassmannians and nilpotent varieties. Indeed, $`X`$ is the Schubert variety in the affine Grassmannian for $`GL_d`$ corresponding to the coweight $`(e^c,f,0,\mathrm{},0)`$ (where we have written as usual $`r=ce+f`$, $`0f<e`$), and the image of $`\stackrel{~}{X}`$ is the closure of the nilpotent orbit corresponding to the partition $`(e^c,f)`$ of $`r`$.
###### Remark 6.3
For this result the fact that the nilpotent endomorphism $`\mathrm{\Pi }`$ in (6.1) is “homogeneous” is essential. Indeed, let $`r=e2`$ and change momentarily notations to consider the following inhomogeneous example. Let
(6.3)
$$(W,\mathrm{\Pi })=k[\mathrm{\Pi }]/(\mathrm{\Pi }^e)k[\mathrm{\Pi }]/(\mathrm{\Pi }),$$
and define $`X`$ as above. Then it is easy to see that any $`X`$ satisfies
$$\mathrm{span}\{\mathrm{\Pi }v_1,\mathrm{},\mathrm{\Pi }^{e1}v_1\}\underset{}{}\underset{}{}W,$$
where $`v_1`$ resp. $`v_2`$ denotes the generator as $`k[\mathrm{\Pi }]`$-module of the first resp. second summand of (6.3). Hence $`X𝐏^1`$. Let $`oX`$ be the special point corresponding to
$$_o=\mathrm{span}\{v_2,\mathrm{\Pi }v_1,\mathrm{},\mathrm{\Pi }^{e1}v_1\}.$$
For $`_o`$, the Jordan type of $`\mathrm{\Pi }|_o`$ is $`(e1,1)`$ whereas for $`_0`$, the Jordan type of $`\mathrm{\Pi }|`$ is $`(e)`$. It follows that the fiber of $`\varphi `$ through a point of $`\pi ^1(_o)`$ has dimension equal to $`\mathrm{dim}GL_r\mathrm{dim}N_{(e1,1)}=e+2`$. On the other hand, the fiber of $`\varphi `$ through a point of $`\pi ^1(X\{o\})`$ has dimension equal to $`(\mathrm{dim}GL_r+1)\mathrm{dim}N_{(e)}=e+1`$. Hence $`\varphi `$ is not smooth in this case.
###### Remark 6.4
Let us fix $`𝐫`$ with $`r=r_\phi `$, and let us consider a standard model $`M(\mathrm{\Lambda },𝐫)`$ with special fiber $`X=X(r,e,d)`$. Let us order $`𝐫=(r_1,\mathrm{},r_e)`$. After extension of scalars from $`k`$ to $`k^{}`$ we obtain as special fiber of the diagram (5.11) the following diagram with cartesian squares,
(6.4)
$$\begin{array}{ccccc}\overline{}& & \overline{\stackrel{~}{}}& & \overline{𝒩}\\ & & & & \\ X_kk^{}& & \stackrel{~}{X}_kk^{}& & \overline{N}_kk^{}& .\end{array}$$
Here $`\overline{𝒩}`$ is the Springer resolution of the nilpotent orbit closure corresponding to $`𝐭=𝐫^{}`$, comp. beginning of section 5. On the other hand, as Ngô pointed out to us, the variety $`\overline{}`$ is an object which is well-known in the theory of the affine Grassmannians, comp. \[N-P\]. Namely, let us introduce the $`e`$ minuscule coweights $`\mu _i=(1^{r_i},0^{dr_i})`$ of $`GL_d`$. Corresponding to $`\mu _i`$ we have the Schubert variety $`\overline{𝒪}_{\mu _i}`$ in the affine Grassmannian for $`GL_d`$ over $`k^{}`$. Then we may identify the variety $`\overline{}`$ with the convolution in the sense of Lusztig, Ginzburg, Mirkovic and Vilonen
(6.5)
$$\overline{𝒪}_{(\mu _1,\mathrm{},\mu _e)}:=\overline{𝒪}_{\mu _1}\stackrel{~}{\times }\mathrm{}\stackrel{~}{\times }\overline{𝒪}_{\mu _e}$$
and the morphism from $`\overline{}`$ to $`X_kk^{}`$ factors through a proper surjective morphism which may be identified with the natural morphism (\[N-P\], §9)
(6.6)
$$m_{(\mu _1,\mathrm{},\mu _e)}:\overline{𝒪}_{(\mu _1,\mathrm{},\mu _e)}\overline{𝒪}_{\mu _1+\mathrm{}+\mu _e}.$$
Note that $`\overline{𝒪}_{\mu _1+\mathrm{}+\mu _e}`$ is just the Schubert variety corresponding to the coweight $`𝐭=𝐫^{}`$.
We note the following consequence of (6.2).
###### Proposition 6.5
If $`re`$, the scheme $`X(r,e,d)`$ is reduced and locally a complete intersection. If $`r>e`$ and Conjecture 5.8 is true, then $`X`$ is still reduced.
Proof. We argue with the special fiber $`\overline{M}`$ of a standard model as described before. If $`re`$, then we may take $`𝐫`$ such that $`r_j1`$, $`j`$. The result then follows from Corollary 4.3 If $`r>e`$, then $`X`$ is smoothly equivalent to (an open subscheme of) $`\overline{N}`$, and the assertion follows from a positive answer to Conjecture 5.8.
It is well-known that the Grassmannian over $`k`$ associated to $`GL_d`$ is not reduced (this happens already for $`d=1`$). Recall (\[BL\]) that
(6.7)
$$\stackrel{~}{\mathrm{Grass}}_k=GL_d(k((\mathrm{\Pi })))/GL_d(k[[\mathrm{\Pi }]]),$$
where $`GL_d(k((\mathrm{\Pi })))`$ resp. $`GL_d(k[[\mathrm{\Pi }]])`$ is the ind-group scheme resp. group scheme which to a $`k`$-algebra $`R`$ associates $`GL_d(R((\mathrm{\Pi })))`$ resp. $`GL_d(R[[\mathrm{\Pi }]])`$, and where the quotient is taken in the category of $`k`$-spaces and turns out to be an ind-scheme, \[BL\], 2.2. On the other hand, the analogous quotient for $`SL_d`$ instead of $`GL_d`$ is an ind-scheme which is reduced and even integral (\[BL\], 6.4.),
(6.8)
$$\stackrel{~}{\mathrm{Grass}}_k^{(0)}=SL_d(k((\mathrm{\Pi })))/SL_d(k[[\mathrm{\Pi }]]).$$
(At this point the blanket assumption in loc. cit. that char $`k=0`$ is not used.) This means that $`\stackrel{~}{\mathrm{Grass}}_k^{(0)}`$ can be obtained as an increasing union of integral $`k`$-schemes. However, the rather indirect proof of this fact in loc.cit. does not give an explicit presentation of $`\stackrel{~}{\mathrm{Grass}}_k^{(0)}`$ as such an increasing union. Based on Proposition 6.5 we are able to give such a presentation, provided Conjecture 5.8 holds true. For this recall (\[BL\], 2.3) that $`\stackrel{~}{\mathrm{Grass}}_k^{(0)}`$ represents the functor on $`k`$-algebras which to a $`k`$-algebra $`R`$ associates the set of special lattices in $`R((\mathrm{\Pi }))^d`$. (A lattice is a $`R[[\mathrm{\Pi }]]`$-submodule $`W`$ of $`R((\mathrm{\Pi }))^d`$ with $`\mathrm{\Pi }^fR[[\mathrm{\Pi }]]^dW\mathrm{\Pi }^fR[[\mathrm{\Pi }]]^d`$ for some $`f`$ and such that the $`R`$-module $`\mathrm{\Pi }^fR[[\mathrm{\Pi }]]^d/W`$ is projective. A lattice $`W`$ is special if the lattice $`^dW`$ in $`^dR((\mathrm{\Pi }))^d=R((\mathrm{\Pi }))`$ is trivial, i.e. equal to $`R[[\mathrm{\Pi }]]`$). Let $`\stackrel{~}{\mathrm{Grass}}_k[0]`$ be the space of lattices of total degree 0 (i.e. the rank of $`\mathrm{\Pi }^fR[[\mathrm{\Pi }]]^d/W`$ is equal to $`fd`$. Then $`\stackrel{~}{\mathrm{Grass}}_k[0]`$ is a connected component of $`\stackrel{~}{\mathrm{Grass}}_k`$ and
(6.9)
$$\stackrel{~}{\mathrm{Grass}}_k^{(0)}=(\stackrel{~}{\mathrm{Grass}}_k[0])_{\mathrm{red}},$$
(cf. \[BL\], 2.2. and 6.4.).
Let $`\stackrel{~}{X}_f`$ be the subscheme of $`\stackrel{~}{\mathrm{Grass}}_k[0]`$ which parametrizes the lattices $`W`$ with $`\mathrm{\Pi }^fR[[\mathrm{\Pi }]]^dW\mathrm{\Pi }^fR[[\mathrm{\Pi }]]^d`$ and let $`X_f`$ be the closed subscheme of $`W`$ in $`\stackrel{~}{X}_f`$ such that $`\mathrm{det}(T\mathrm{\Pi }|(W/\mathrm{\Pi }^fR[[\mathrm{\Pi }]]^d))T^{fd}`$. We have an exact sequence
$$0\mathrm{\Pi }^fR[[\mathrm{\Pi }]]^d/\mathrm{\Pi }^{f+1}R[[\mathrm{\Pi }]]^dW/\mathrm{\Pi }^{f+1}R[[\mathrm{\Pi }]]^dW/\mathrm{\Pi }^fR[[\mathrm{\Pi }]]^d0,$$
which implies that
$$\mathrm{det}(T\mathrm{\Pi }|(W/\mathrm{\Pi }^{f+1}R[[\mathrm{\Pi }]]^d))=\mathrm{det}(T\mathrm{\Pi }|(W/\mathrm{\Pi }^fR[[\mathrm{\Pi }]]^d))T^d.$$
We therefore obtain a chain of closed embeddings of $`k`$-schemes
(6.10)
$$X_0X_1\mathrm{}$$
###### Proposition 6.6
For $`d2`$, the chain (6.10) presents $`\stackrel{~}{\mathrm{Grass}}_k^{(0)}`$ as an increasing union of integral $`k`$-schemes. The same is true for arbitrary $`d`$, if Conjecture 5.8 holds true. In particular, this holds (by the theorem of Weyman) if char $`k=0`$.
Proof. We note that $`X_f=X(fd,2f,d)`$, hence is integral if $`d2`$ and for arbitrary $`d`$, if Conjecture 5.8 holds true. Hence $`X_f=(\stackrel{~}{X}_f)_{\mathrm{red}}`$. The claim follows from $`\stackrel{~}{\mathrm{Grass}}_k[0]=\underset{}{lim}\stackrel{~}{X}_f`$ and
$$\stackrel{~}{\mathrm{Grass}}_k^{(0)}=\underset{}{lim}(\stackrel{~}{X}_f)_{\mathrm{red}}=\underset{}{lim}X_f.$$
## 7 The complex of nearby cycles
In this section we suppose that the residue field $`k`$ of $`𝒪_E`$ is finite and shall aim for an expression for the complex of nearby cycles of a standard model for $`GL_d`$ corresponding to $`𝐫`$, by exploiting in more depth the resolution of singularities given by the scheme $`𝒩`$ in (5.3).
Recall that $`K`$ denotes the Galois hull of $`F`$ in $`F_0^{\mathrm{sep}}`$ and that $`k^{}`$ is the residue field of $`𝒪_K`$. Also, $`\mathrm{\Gamma }`$ denotes the Galois group of $`K`$ over $`E`$. We fix a prime number $`\mathrm{}`$ which is invertible in $`𝒪_E`$ and denote by $`R\psi =R\psi \overline{𝐐}_{\mathrm{}}`$ the complex of nearby cycles of the $`𝒪_E`$-scheme $`M`$. Since $`M^{\mathrm{loc}}`$ is the scheme-theoretic closure of the generic fiber in $`M`$, this complex has support in $`\overline{M}^{\mathrm{loc}}:=M^{\mathrm{loc}}_{𝒪_E}k`$. The complex is equipped with an action of $`\mathrm{Gal}(F_0^{\mathrm{sep}}/E)`$ which lifts the action on $`\overline{M}^{\mathrm{loc}}`$. We also fix a square root of the cardinality $`|k|`$ in $`\overline{𝐐}_{\mathrm{}}`$.
###### Theorem 7.1
There is an isomorphism between perverse sheaves pure of weight zero,
$$R\psi [\mathrm{dim}\overline{M}^{\mathrm{loc}}]\left(\frac{1}{2}\mathrm{dim}\overline{M}^{\mathrm{loc}}\right)=\underset{𝐬𝐭}{}_𝐬IC_{M_𝐬},$$
where $`IC_{M_𝐬}`$ denotes the intermediate extension of the constant sheaf$`\overline{𝐐}_{\mathrm{}}[\mathrm{dim}M_𝐬]\left(\frac{1}{2}\mathrm{dim}M_𝐬\right)`$, equipped with the action of $`\mathrm{Gal}(F_0^{\mathrm{sep}}/E)`$ which factors through $`\mathrm{Gal}(\overline{k}/k)`$, and $`_𝐬`$ is a $`\overline{𝐐}_{\mathrm{}}`$-vector space equipped with an action of $`\mathrm{Gal}(F_0^{\mathrm{sep}}/E)`$. The action of $`\mathrm{Gal}(F_0^{\mathrm{sep}}/E)`$ on $`_𝐬`$ factors through $`\mathrm{\Gamma }`$, and the degree of this representation is given as a Kostka number,
$$m_𝐬=\mathrm{dim}_𝐬=K_{𝐬^{},𝐫}$$
(cf. \[Mc\], p. 115).
By making use of naturality properties of the complex of nearby cycles, the diagram in Theorem 5.4 (iii) with smooth horizontal morphisms reduces us to proving the corresponding statements for $`N^{\mathrm{loc}}`$ instead of $`M^{\mathrm{loc}}`$ (or $`M`$). To be more precise, let us change notations and denote now by $`R\psi `$ the complex of nearby cycles of the $`𝒪_E`$-scheme $`N^{\mathrm{loc}}`$. We wish to prove the formula
(7.1)
$$R\psi [\mathrm{dim}\overline{N}^{\mathrm{loc}}]\left(\frac{1}{2}\mathrm{dim}\overline{N}^{\mathrm{loc}}\right)=\underset{𝐬𝐭}{}_𝐬IC_{N_𝐬},$$
with $`_𝐬`$ as above.
The left hand side is a perverse sheaf of weight zero on $`\overline{N}^{\mathrm{loc}}`$ (\[I\] Thm. 4.2 and Cor. 4.5), which is $`GL_r`$-equivariant. By \[BBD\], 5.3.8 its inverse image $`R\overline{\psi }`$ on $`\overline{N}^{\mathrm{loc}}_k\overline{k}`$ is semisimple and its simple constituents are all of the form $`IC_{N_𝐬_k\overline{k}}`$, for some $`𝐬𝐭`$, since $`N_𝐬_k\overline{k}`$ admits no non-trivial $`GL_r`$-equivariant irreducible local system. Therefore, the isotypical decomposition of $`R\overline{\psi }`$ has the form
(7.2)
$$R\overline{\psi }=\underset{𝐬𝐭}{}K_𝐬,$$
where $`K_𝐬`$ is a multiple of $`IC_{N_𝐬_k\overline{k}}`$. Obviously $`K_𝐬`$ is of the form
(7.3)
$$K_𝐬=_𝐬IC_{N_𝐬_k\overline{k}},$$
where $`_𝐬`$ is a $`\overline{𝐐}_{\mathrm{}}`$-vector space with an action of the inertia group $`I\mathrm{Gal}(F_0^{\mathrm{sep}}/E)`$. The fact that both sides of (7.3) come by extension of scalars from $`k`$, implies now that $`\mathrm{Gal}(F_0^{\mathrm{sep}}/E)`$ acts on $`_𝐬`$ and that we have a decomposition of the form (7.1).
It remains to show that the restriction of $`_𝐬`$ to $`\mathrm{Gal}(F_0^{\mathrm{sep}}/K)`$ is trivial and to determine its degree. By Deligne \[D\], Prop.3.7., we have
(7.4)
$$R\psi _kk^{}=R\psi ^{}$$
with $`R\psi ^{}`$ the complex of nearby cycles of the $`𝒪_K`$-scheme $`N^{\mathrm{loc}}_{𝒪_E}𝒪_K`$. Now $`N^{\mathrm{loc}}_{𝒪_E}K`$ has the smooth model $`𝒩`$ which maps via $`\mu `$ to $`N^{\mathrm{loc}}_{𝒪_E}𝒪_K`$. Since the complex of nearby cycles of a smooth scheme is the constant sheaf placed in degree 0, the functoriality with respect to push-forward under a proper morphism gives a natural identification of complexes on $`\overline{N}^{\mathrm{loc}}_kk^{}`$,
(7.5)
$$R\psi ^{}=R\overline{\mu }_{}\overline{𝐐}_{\mathrm{}}.$$
In particular, the action of $`\mathrm{Gal}(F_0^{\mathrm{sep}}/K)`$ on $`R\psi ^{}`$ is through $`\mathrm{Gal}(\overline{k}/k^{})`$. Now the morphism $`\overline{\mu }:\overline{𝒩}_kk^{}\overline{N}^{\mathrm{loc}}_kk^{}`$ comes by base change from the moment map $`\overline{\overline{\mu }}`$ for the variety $`\overline{𝓕}`$ of partial flags over $`k`$, cf. (5.2). This map is semi-small with all strata $`N_𝐬`$ of $`\overline{N}^{\mathrm{loc}}`$ relevant, comp. \[BM\]. Base changing the decomposition (7.1) from $`k`$ to $`k^{}`$ and identifying the left hand side with $`R\overline{\mu }_{}\overline{𝐐}_{\mathrm{}}[\mathrm{dim}\overline{N}^{\mathrm{loc}}](\frac{1}{2}\mathrm{dim}\overline{N}^{\mathrm{loc}})`$, we see that
(7.6)
$$_𝐬=R^{2d_𝐬}\overline{\mu }_{}\overline{𝐐}_{\mathrm{}}(d_𝐬)_{|N_𝐬_k\overline{k}},$$
as representations of $`\mathrm{Gal}(F_0^{\mathrm{sep}}/K)`$, i.e. of $`\mathrm{Gal}(\overline{k}/k^{})`$. Here $`d_𝐬=\frac{1}{2}\mathrm{codim}N_𝐬`$ is the relative dimension of $`\overline{\mu }`$ over $`N_𝐬`$. But by Spaltenstein \[Sp\] all irreducible components of $`\overline{\overline{\mu }}^1(N_𝐬_k\overline{k})`$ are defined over $`k`$, hence $`\mathrm{Gal}(\overline{k}/k^{})`$ acts trivially on the right hand side of (7.6). This shows that $`\mathrm{Gal}(F_0^{\mathrm{sep}}/K)`$ acts trivially on $`_𝐬`$. (Spaltenstein works over an algebraically closed field, but his results are valid over an arbitrary field, comp. \[HS\], §2.) A different argument, pointed out by G.Laumon, is to transpose the result of Braverman and Gaitsgory \[BG\] from $`𝐂`$ to a finite field and appeal to \[HS\], Corollary 2.3.
For the degree $`m_𝐬`$ of $`_𝐬`$ Borho and MacPherson give the formula \[BM\], 3.5.,
(7.7)
$$m_𝐬=\mathrm{dim}\mathrm{Hom}_{S_r}(\chi ^𝐬,\mathrm{Ind}_{S_𝐫}^{S_r}(\mathrm{sgn})).$$
Here $`S_𝐫=S_{r_1}\times \mathrm{}\times S_{r_e}`$ is a subgroup of the symmetric group $`S_r`$ and $`\mathrm{sgn}`$ is the sign character on all factors. Furthermore, $`\chi ^𝐬`$ denotes the unique irreducible representation of $`S_r`$ which occurs both in $`\mathrm{Ind}_{S_𝐬}^{S_r}(1)`$ and in $`\mathrm{Ind}_{S_𝐬}^{S_r}(\mathrm{sgn})`$, cf. \[Mc\], p. 115. By loc.cit. (7.7) can be identified with the Kostka number occurring in Theorem 7.1. Note that this agrees with the formula of Braverman and Gaitsgory \[BG\], Cor. 1.5.,
(7.8)
$$m_𝐬=\mathrm{dim}V(𝐬^{})_𝐫,$$
(where, however, in their formula $`V(𝐏)_𝐝`$ should be replaced by $`V(𝐏^{})_𝐝`$).
Here $`V(𝐬^{})`$ denotes the rational representation of $`GL_e`$,
$$V(𝐬^{})=\mathrm{Hom}_{S_r}(\chi ^𝐬^{},(\mathrm{nat}_{GL_e})^r)$$
where $`\mathrm{nat}`$ is the natural representation of $`GL_e`$ and $`V(𝐬^{})_𝐫`$ the weight space corresponding to $`𝐫`$. By \[Mc\], p. 163 the character of $`V(𝐬^{})`$ is given by the Schur function $`s_𝐬^{}`$ and hence (7.8) is indeed equal to $`K_{𝐬^{},𝐫}`$ by \[Mc\], p. 101. (Both sources \[BM\] and \[BG\] work over $`𝐂`$, but can be transposed to the present context.)
###### Remark 7.2
We have used here Theorem 4.1 in order to deduce Theorem 7.1 from the smoothness of $`𝒩`$ and the fact that the Springer resolution of the nilpotent orbit closure corresponding to $`𝐭=𝐫^{}`$ is semi-small. Of course, Theorem 4.1 implies also that the scheme $``$ is smooth and that the special fiber of $``$ is a semi-small resolution of $`\overline{M}_𝐭`$. On the other hand, this last fact has a direct proof, cf. \[N-P\], Lemma 9.3. In fact, the Remark 6.4 establishes an equivalence between the semi-smallness properties.
###### Remark 7.3
(B. C. Ngô) Denote by $`R\psi ^{}`$ the complex of nearby cycles of the $`𝒪_K`$-scheme $`M^{\mathrm{loc}}_{𝒪_E}𝒪_K`$. Let us fix an $`𝒪_F`$-basis of the module $`\mathrm{\Lambda }`$. Recall that then by (6.5), the special fiber $`\overline{}`$ can be identified with the convolution
$$\overline{𝒪}_{\mu _1}\stackrel{~}{\times }\mathrm{}\stackrel{~}{\times }\overline{𝒪}_{\mu _e}$$
and the morphism $`\overline{}\overline{M}^{\mathrm{loc}}\times _kk^{}`$ with the natural morphism
$$m_{(\mu _1,\mathrm{},\mu _e)}:\overline{𝒪}_{\mu _1}\stackrel{~}{\times }\mathrm{}\stackrel{~}{\times }\overline{𝒪}_{\mu _e}\overline{𝒪}_{\mu _1+\mathrm{}+\mu _e}$$
of (6.6); here $`\overline{𝒪}_\mu `$ denotes the reduced closure of the orbit which corresponds to the dominant coweight $`\mu `$ in the affine Grassmanian for $`GL_d`$ over $`k^{}`$. The same arguments as in the proof of Theorem 7.1, applied to the morphism $`\pi :M^{\mathrm{loc}}_{𝒪_E}𝒪_K`$ in place of $`𝒩N^{\mathrm{loc}}_{𝒪_E}𝒪_K`$, show now that
(7.9)
$$R\psi ^{}[\mathrm{dim}\overline{M}^{\mathrm{loc}}]\left(\frac{1}{2}\mathrm{dim}\overline{M}^{\mathrm{loc}}\right)=R\overline{\pi }_{}\overline{𝐐}_l[\mathrm{dim}\overline{}]\left(\frac{1}{2}\mathrm{dim}\overline{}\right).$$
By the above discussion, the right hand side of (7.9) is the convolution $`IC_{𝒪_{\mu _1}}\mathrm{}IC_{𝒪_{\mu _e}}`$. Hence, we obtain a relation between nearby cycles and convolution:
(7.10)
$$R\psi ^{}[\mathrm{dim}\overline{M}^{\mathrm{loc}}]\left(\frac{1}{2}\mathrm{dim}\overline{M}^{\mathrm{loc}}\right)=IC_{𝒪_{\mu _1}}\mathrm{}IC_{𝒪_{\mu _e}}.$$
Recall that there is an equivalence of tensor categories between the category of $`\stackrel{~}{𝒢}_k^{}`$-equivariant pure perverse $`\overline{𝐐}_l`$-sheaves of weight $`0`$ on the affine Grassmanian $`\stackrel{~}{\mathrm{Grass}}_k^{}`$ (with tensor structure given by the convolution product) and the category of finite dimensional $`\overline{𝐐}_l`$-representations of the Langlands dual group $`GL_d=\widehat{GL}_d`$ (see \[Gi\], and especially \[M-V\] §7). Under this equivalence the perverse sheaf $`IC_{𝒪_\mu }`$ corresponds to the representation $`V(\mu )`$ of $`GL_d`$ of highest weight $`\mu `$, and the convolution in (7.10) to the tensor product
$$V(\mu _1)\mathrm{}V(\mu _e).$$
This decomposes
$$V(\mu _1)\mathrm{}V(\mu _e)=\underset{\lambda \mu _1+\mathrm{}+\mu _e}{}_\lambda V(\lambda ),$$
where $`_\lambda `$ is a finite dimensional $`\overline{𝐐}_l`$-vector space. Using again the above equivalence of categories and (7.10), we obtain
(7.11)
$$R\psi ^{}[\mathrm{dim}\overline{M}^{\mathrm{loc}}]\left(\frac{1}{2}\mathrm{dim}\overline{M}^{\mathrm{loc}}\right)=\underset{\lambda \mu _1+\mathrm{}+\mu _e}{}_\lambda IC_{𝒪_\lambda }.$$
The right hand side of (7.11) corresponds to the expression in Theorem 7.1. Indeed, $`\{\lambda |\lambda \mu _1+\mathrm{}+\mu _e\}`$ corresponds to $`\{𝐬|𝐬𝐫^{}\}`$, and we can see directly that the Littlewood-Richardson number $`dim(_\lambda )`$ is equal to the Kostka number $`K_{𝐬^{},𝐫}`$ of 7.1.
###### Remark 7.4
As before, let us choose an ordering of the set of embeddings $`\varphi :FF_0^{\mathrm{sep}}`$. The Galois group $`\mathrm{\Gamma }=\mathrm{Gal}(K/E)`$ can be identified with the subgroup of elements $`\sigma `$ of the symmetric group $`S_e`$ which satisfy $`r_{\sigma (i)}=r_i`$. The group $`\mathrm{\Gamma }`$ acts by permutation of the factors on the tensor product $`V(\mu _1)\mathrm{}V(\mu _e)`$. Let us denote by $`\rho `$ the corresponding representation
$$\rho :\mathrm{\Gamma }GL(V(\mu _1)\mathrm{}V(\mu _e)).$$
Since the permutation action commutes with the action of $`GL_d`$ on the tensor product, the representation $`\rho `$ decomposes as
$$\rho =\underset{\lambda \mu _1+\mathrm{}+\mu _e}{}\rho _\lambda \mathrm{id}_{V(\lambda )}$$
where $`\rho _\lambda `$ is a representation of $`\mathrm{\Gamma }`$ on the vector space $`_\lambda `$. We conjecture that the representation of $`\mathrm{Gal}(K/E)`$ on $`_𝐬`$ (see Theorem 7.1) is isomorphic to $`\rho _\lambda `$ (with $`\lambda `$ corresponding to $`𝐬`$).
After our choice of an $`𝒪_F`$-basis of $`\mathrm{\Lambda }`$, the generic fiber of $`M^{\mathrm{loc}}`$ can be identified with the product of Grassmanians
$$M^{\mathrm{loc}}_{𝒪_K}K=\underset{i=1}{\overset{e}{}}\mathrm{Grass}_{r_i}(K^d).$$
For $`\sigma \mathrm{\Gamma }`$, permutation of the factors gives an isomorphism of $`K`$-schemes
$$\kappa (\sigma ):M^{\mathrm{loc}}_{𝒪_E}KM^{\mathrm{loc}}_{𝒪_E}K.$$
The isomorphism $`\kappa (\sigma )`$ induces an automorphism of the sheaf of vanishing cycles $`R\psi ^{}`$. Therefore, by (7.10) we obtain a “commutativity” isomorphism
$$\kappa (\sigma ):IC_{𝒪_{\mu _1}}\mathrm{}IC_{𝒪_{\mu _e}}IC_{𝒪_{\mu _1}}\mathrm{}IC_{𝒪_{\mu _e}}=IC_{𝒪_{\mu _{\sigma (1)}}}\mathrm{}IC_{𝒪_{\mu _{\sigma (e)}}}.$$
(notice here that, since $`\sigma \mathrm{\Gamma }`$, $`\mu _{\sigma (i)}=\mu _i`$). As was pointed out by Ngô, the conjecture follows, if $`\kappa (\sigma )`$ coincides with the isomorphism given using the permutation $`\sigma `$ and the “commutativity constraint” (for the tensor category of perverse sheaves of Remark 7.3) of \[M-V\].
###### Remark 7.5
Let $`x`$ be a point of $`M^{\mathrm{loc}}`$ with values in a finite extension $`𝐅_q`$ of $`k`$. Let $`\mathrm{Tr}^{ss}(\mathrm{Fr}_q,R\psi _x^M)`$ be the semi-simple trace of the geometric Frobenius on the stalk at $`x`$ of the complex of nearby cycles of $`M^{\mathrm{loc}}`$. It should be possible to transfer the conjecture of Kottwitz \[HN\] to this case to obtain a group-theoretic expression for this (the group $`R_{F/F_0}(GL_d)`$ relevant for this conjecture here is not split so that the original formulation does not apply directly). Note that
(7.12)
$$\mathrm{Tr}^{ss}(\mathrm{Fr}_q,R\psi _x^M)=\mathrm{Tr}^{ss}(\mathrm{Fr}_q,R\psi _y^N)$$
where $`y=\varphi (\stackrel{~}{x})`$ for an arbitrary point $`\stackrel{~}{x}\stackrel{~}{M}^{\mathrm{loc}}(𝐅_q)`$ mapping to $`x`$, and where $`R\psi ^N`$ is the complex of nearby cycles of $`N^{\mathrm{loc}}`$. The expression (7.12) only depends on the stratum $`M_𝐬`$ containing $`x`$ resp. $`N_𝐬`$ containing $`y`$ and may therefore be denoted by $`\mathrm{Tr}^{ss}(Fr_q,R\psi _𝐬^N)`$.
We also consider the analogous semi-simple traces for points with values in a finite extension $`𝐅_q`$ of $`k^{}`$,
(7.13) $`\mathrm{Tr}^{ss}(Fr_q,R\psi _x^{MK}),\mathrm{Tr}^{ss}(Fr_q,R\psi _x^{NK}),`$
$`\mathrm{Tr}^{ss}(Fr_q,R\psi _𝐬^{MK}),\mathrm{Tr}^{ss}(Fr_q,R\psi _𝐬^{NK}).`$
(these semi-simple traces may differ from the preceding ones when $`K/E`$ is ramified).
¿From Theorem 7.1 we obtain
(7.14)
$$\mathrm{Tr}^{ss}(\mathrm{Fr}_q;R\psi _𝐬^{NK})=q^{\frac{1}{2}\mathrm{dim}\overline{N}^{\mathrm{loc}}}\underset{𝐬𝐬^{}𝐭}{}m_𝐬^{}\mathrm{Tr}^{ss}(\mathrm{Fr}_q,(IC_{N_𝐬^{}})_𝐬),$$
provided $`𝐅_qk^{}`$. By Lusztig \[L\], the entity $`\mathrm{Tr}^{ss}(\mathrm{Fr}_q,(IC_{N_𝐬^{}})_𝐬)`$ can be expressed in terms of Kostka polynomials $`K_{𝐬^{},𝐬}(q^1)`$, comp. also \[Mc\], p. 245.
Consider the spectral sequence of nearby cycles,
(7.15)
$$E_2^{pq}=H^p(M^{\mathrm{loc}}_k\overline{k},R^q\psi ^M)H^{p+q}(M^{\mathrm{loc}}_{𝒪_E}F_0^{\mathrm{sep}},\overline{𝐐}_{\mathrm{}}).$$
We obtain an identity of semi-simple traces (cf. \[HN\])
(7.16)
$$\mathrm{Tr}^{ss}(\mathrm{Fr}_q,H^{}(M^{\mathrm{loc}}_{𝒪_E}F_0^{\mathrm{sep}},\overline{𝐐}_{\mathrm{}}))=\underset{xM^{\mathrm{loc}}(𝐅_q)}{}\mathrm{Tr}^{ss}(\mathrm{Fr}_q,R\psi _x^M).$$
Similarly we may consider the spectral sequence of nearby cycles for $`M^{\mathrm{loc}}_{𝒪_E}𝒪_K`$. We then obtain an identity of semi-simple traces (as $`\mathrm{Gal}(F_0^{\mathrm{sep}}/K)`$-modules), provided that $`𝐅_qk^{}`$,
(7.17)
$$\mathrm{Tr}^{ss}(Fr_q;H^{}((M_{𝒪_E}𝒪_K)_{𝒪_K}F_0^{\mathrm{sep}},\overline{𝐐}_{\mathrm{}}))=\underset{xM(𝐅_q)}{}\mathrm{Tr}^{ss}(\mathrm{Fr}_q,R\psi _x^{MK}).$$
However,
$$M^{\mathrm{loc}}_{𝒪_E}K=\underset{\phi }{}\mathrm{Grass}_{r_\phi }(V_\phi ),$$
cf. (2.7), and hence the left hand side of (7.17) is known. Taking into account (7.12) we obtain therefore from (7.16) a combinatorial identity involving Kostka polynomials, Kostka numbers etc. It might be interesting to identify this combinatorial identity explicitly.
## 8 Local models of $`EL`$-type
In this section we use the following notation (following \[RZ\]).
$`F_0`$ a complete discretely valued field with perfect residue field
$`F`$ a finite separable field extension of $`F_0`$,
$`B`$ a simple algebra with center $`F`$,
$`V`$ a finite $`B`$-module,
$`G=GL_B(V)`$, as algebraic group over $`F_0`$,
$`\mu :𝐆_{mK}G_K`$ a one parameter subgroup, defined over some sufficiently big extension $`K`$ contained in a fixed separable closure $`F_0^{\mathrm{sep}}`$ of $`F_0`$, given up to conjugation. We assume that the eigenspace decomposition of $`V_{𝐐_p}K`$ is given by
(8.1)
$$V_{F_0}K=V_0V_1.$$
$`E`$ the field of definition of the conjugacy class of $`\mu `$.
$``$ a periodic $`O_B`$-lattice chain in $`V`$.
In the sequel we denote by $`𝒪_{F_0},𝒪_F,𝒪_B,𝒪_E`$ the respective rings of integers.
To these data we associate the following functor on $`(Sch/𝒪_E)`$: The $`S`$-valued points are given by
1. a functor $`\mathrm{\Lambda }t_\mathrm{\Lambda }`$ to the category of $`𝒪_B_{𝒪_{F_0}}𝒪_S`$-modules on $`S`$
2. a morphism of functors $`\phi _\mathrm{\Lambda }:\mathrm{\Lambda }_{𝒪_{F_0}}𝒪_St_\mathrm{\Lambda }`$.
The requirements on these data are:
a) $`t_\mathrm{\Lambda }`$ is locally on $`S`$ a free $`𝒪_S`$-module of finite rank, and we have the following identity of polynomial functions on $`𝒪_B`$:
(8.2)
$$\mathrm{det}_{𝒪_S}(a|t_\mathrm{\Lambda })=\mathrm{det}_K(a|V_1).$$
b) $`\phi _\mathrm{\Lambda }`$ is surjective, for all $`\mathrm{\Lambda }`$.
We remark that the standard models considered in sections 27 are a special case. Indeed, let $`B=F`$ and let $``$ consist of the $`F^\times `$-multiples of a fixed $`𝒪_F`$-lattice $`\mathrm{\Lambda }_0`$ in the $`d`$-dimensional $`F`$-vector space $`V`$ and assume that under the decomposition (2.5) of $`F_{F_0}F_0^{\mathrm{sep}}`$ we have
(8.3)
$$V_{F_0}F_0^{\mathrm{sep}}=\underset{\phi }{}V_\phi ,V_0_KF_0^{\mathrm{sep}}=\underset{\phi }{}V_{0,\phi },V_1_KF_0^{\mathrm{sep}}=\underset{\phi }{}V_{1,\phi },$$
with $`\mathrm{dim}_{F_0^{\mathrm{sep}}}V_{0,\phi }=r_\phi `$. If $`(t_\mathrm{\Lambda },\phi _\mathrm{\Lambda })_\mathrm{\Lambda }`$ is an $`S`$-valued point of the moduli problem above, then $`=\mathrm{Ker}\phi _{\mathrm{\Lambda }_0}`$ is an $`S`$-valued point of the standard model for $`GL_d`$ corresponding to $`𝐫=(r_\phi )`$ and this establishes an isomorphism of moduli problems. Furthermore, $`E`$ has indeed the description given in section 2 in terms of $`𝐫`$.
In general, the above moduli problem is representable by a projective scheme over $`\mathrm{Spec}𝒪_E`$ which is a closed subscheme of a form over $`𝒪_E`$ of a product of Grassmannians. Let us make this more precise.
Let us first consider the case where $`B=F`$ is a totally ramified extension of degree $`e`$ of $`F_0`$. Let $`v_1,\mathrm{},v_d`$ be a basis of $`V`$. For $`i=0,\mathrm{},d1`$ consider the $`𝒪_F`$-lattice $`\mathrm{\Lambda }_i`$ in $`V`$ spanned by
(8.4)
$$\mathrm{\Lambda }_i=\mathrm{span}\{\pi ^1v_1,\mathrm{},\pi ^1v_i,v_{i+1},\mathrm{},v_d\}.$$
Here $`\pi `$ denotes a uniformizer in $`F`$. This yields a complete periodic lattice chain
(8.5)
$$\mathrm{}\mathrm{\Lambda }_0\mathrm{\Lambda }_1\mathrm{}\mathrm{\Lambda }_{d1}\pi ^1\mathrm{\Lambda }_0\mathrm{}$$
Choose $`I=\{i_0<i_1<\mathrm{}<i_{\mathrm{}}\}\{0,\mathrm{},d1\}`$. Then the periodic lattice $``$ is isomorphic to the subchain of (8.5) where only $`\mathrm{\Lambda }_i`$ with $`iI`$ are kept, for suitable $`I`$. The points of the above functor with values in a $`𝒪_E`$-scheme $`S`$ can now be interpreted as the isomorphism classes of commutative diagrams
$$\begin{array}{ccccccccc}\mathrm{\Lambda }_{i_0,S}& & \mathrm{\Lambda }_{i_1,S}& & \mathrm{}& & \mathrm{\Lambda }_{i_{\mathrm{}},S}& & \pi ^1\mathrm{\Lambda }_{i_0,S}\\ & & & & & & & & \\ _0& & _1& & \mathrm{}& & _{\mathrm{}}& & \pi ^1_0& ,\end{array}$$
where $`\mathrm{\Lambda }_{i_j,S}`$ is $`\mathrm{\Lambda }_{i_j}_{𝒪_{F_0}}𝒪_S`$ and where the $`_j`$ are $`𝒪_F_{𝒪_{F_0}}𝒪_S`$-submodules which locally on $`S`$ are direct summands of $`\mathrm{\Lambda }_{i_j,S}`$ as $`𝒪_S`$-modules and where we have an identity of polynomial functions on $`𝒪_F`$,
(8.6)
$$\mathrm{det}(a|_j)=\mathrm{det}_K(a|V_0),j=0,\mathrm{},\mathrm{}.$$
Suppose that $`K`$ is a sufficiently big Galois extension so that
(8.7)
$$F_{F_0}K=\underset{\phi :FK}{}K.$$
We obtain corresponding decompositions
(8.8)
$$V_{F_0}K=\underset{\phi }{}V_\phi ,V_0_{F_0}K=\underset{\phi }{}V_{0,\phi }.$$
Then $`\mathrm{dim}_KV_\phi =d`$ for all $`\phi `$. Let $`r_\phi =\mathrm{dim}_KV_{0,\phi }`$. Then the determinant condition (8.6) can be rewritten as
(8.9)
$$\mathrm{det}(a|_j)=\underset{\phi }{}\phi (a)^{r_\phi }.$$
In other words, each $`_i`$ is an $`S`$-valued point of the standard model for $`GL_d`$ corresponding to $`𝐫=(r_\phi )`$ (and $`\mathrm{\Lambda }_i`$). Our functor is representable by a closed subscheme of the product over $`𝒪_E`$ of these standard models, one for each $`i=0,\mathrm{},\mathrm{}`$.
Now let us consider the general case. Let $`\stackrel{˘}{F}_0`$ be the completion of the maximal unramified extension of $`F_0`$ in $`F_0^{\mathrm{sep}}`$. If $`F_1`$ is the maximal unramified extension of $`F_0`$ contained in $`F`$, we obtain a decomposition
(8.10)
$$F_1_{F_0}\stackrel{˘}{F}_0=\stackrel{˘}{F}_0\mathrm{}\stackrel{˘}{F}_0,$$
corresponding to the $`f=[F_1:F_0]`$ different embeddings $`\alpha :F_1\stackrel{˘}{F}_0`$ of $`F_1`$ into $`\stackrel{˘}{F}_0`$. For fixed $`\alpha `$, the extension $`\stackrel{˘}{F}_\alpha =F_{F_{1,\alpha }}\stackrel{˘}{F}_0`$ is a totally ramified field extension of degree $`e=[F:F_0]/f`$ of $`\stackrel{˘}{F}_0`$. The simple central algebra $`B_F\stackrel{˘}{F}_\alpha =B_{F_{1,\alpha }}\stackrel{˘}{F}_0`$ splits, i.e. is isomorphic to a matrix algebra over $`\stackrel{˘}{F}_\alpha `$. Similarly, we obtain a decomposition
(8.11)
$$V_{F_0}\stackrel{˘}{F}_0=\underset{\alpha }{}V_\alpha ,$$
where $`V_\alpha `$ is a $`\stackrel{˘}{F}_\alpha `$-vector space, all of the same dimension $`d`$.
Since $`F_1`$ is an unramified extension we obtain similar decompositions for the rings of integers, e.g.
(8.12)
$$𝒪_B_{𝒪_{F_0}}𝒪_{\stackrel{˘}{F}_0}=\underset{\alpha }{}𝒪_B_{𝒪_{F_1},\alpha }𝒪_{\stackrel{˘}{F}_0},$$
where for each $`\alpha `$ the summand $`𝒪_B_{𝒪_{F_1},\alpha }𝒪_{\stackrel{˘}{F}_0}=𝒪_B_{𝒪_F}𝒪_{\stackrel{˘}{F}_\alpha }`$ is a parahoric order in the matrix algebra $`B_F\stackrel{˘}{F}_\alpha `$. Similarly, the periodic lattice chain $``$ corresponds to periodic $`𝒪_{\stackrel{˘}{F}_\alpha }`$-lattice chains $`_\alpha `$ in each $`\stackrel{˘}{F}_\alpha `$-vector space $`V_\alpha `$ in (8.11). Let $`K`$ be a sufficiently big Galois extension of $`F_0`$. Then
(8.13)
$$V_{F_0}K=\underset{\alpha }{}V_\alpha _{\stackrel{˘}{F}_0}\stackrel{˘}{K},$$
where $`\stackrel{˘}{K}=\stackrel{˘}{F}_0.K`$, and for each $`\alpha `$
(8.14)
$$V_\alpha _{\stackrel{˘}{F}_0}\stackrel{˘}{K}=V_{\alpha ,0}V_{\alpha ,1}.$$
Let $`\stackrel{˘}{E}=E.\stackrel{˘}{F}_0`$. It now follows that our functor is representable by a projective scheme over $`𝒪_E`$ which after base change from $`𝒪_E`$ to $`𝒪_{\stackrel{˘}{F}}`$ becomes isomorphic to the product over all $`\alpha `$ of schemes considered before for the data $`(\stackrel{˘}{F}_\alpha /\stackrel{˘}{F}_0,V_\alpha ,\mu _\alpha ,_\alpha )`$.
Let us denote by $`M^{\mathrm{naive}}`$ the scheme over $`\mathrm{Spec}𝒪_E`$ associated in this way to the data fixed in the beginning of this section. This is what we call a naive local model of $`EL`$-type. We know from the special case of a standard model for $`GL_d`$ that $`M^{\mathrm{naive}}`$ is rarely flat over $`\mathrm{Spec}𝒪_E`$. We now define a closed subscheme $`M^{\mathrm{loc}}`$ of $`M^{\mathrm{naive}}`$ which stands a better chance of being flat over the base scheme.
Assume first that $`B=F`$ is a totally ramified extension. In the notation introduced after (8.5) let $`=_I`$. For every $`iI`$ we obtain a morphism
(8.15)
$$\pi _i:M^{\mathrm{naive}}M^{\mathrm{naive}}(\mathrm{\Lambda }_i)iI.$$
Here $`M^{\mathrm{naive}}(\mathrm{\Lambda }_i)=M(\mathrm{\Lambda }_i,\mu )`$ denotes the standard model associated to $`\mathrm{\Lambda }_i`$ (and $`(F,V,\mu )`$). We then define in this case
(8.16)
$$M^{\mathrm{loc}}=\underset{iI}{}\pi _i^1(M^{\mathrm{loc}}(\mathrm{\Lambda }_i))$$
(scheme-theoretic intersection inside $`M^{\mathrm{naive}}`$).
In the general case we have, with the notation used above,
(8.17)
$$M^{\mathrm{naive}}_{𝒪_E}𝒪_{\stackrel{˘}{E}}=\underset{\alpha }{}M^{\mathrm{naive}}(\stackrel{˘}{F}_\alpha /\stackrel{˘}{F}_0,V_\alpha ,\mu _\alpha ,_\alpha ).$$
Let
(8.18)
$$\stackrel{˘}{M}^{\mathrm{loc}}=\underset{\alpha }{}M^{\mathrm{loc}}(\stackrel{˘}{F}_\alpha /\stackrel{˘}{F}_0,V_\alpha ,\mu _\alpha ,_\alpha ).$$
The descent datum on $`M^{\mathrm{naive}}_{𝒪_E}𝒪_{\stackrel{˘}{E}}`$ respects the closed subscheme $`\stackrel{˘}{M}^{\mathrm{loc}}`$ and hence defines a closed subscheme $`M^{\mathrm{loc}}`$ of $`M^{\mathrm{naive}}`$.
We conjecture that $`M^{\mathrm{loc}}`$ is flat over $`\mathrm{Spec}𝒪_E`$. This would constitute the analogue of the result of Görtz \[G\] which confirms the conjecture in the case when $`F/F_0`$ is unramified.
George Pappas Michael Rapoport
Dept. of Mathematics Mathematisches Institut
Michigan State University der Universität zu Köln
E. Lansing Weyertal 86 – 90
MI 48824-1027 50931 Köln
USA Germany
email: pappas@math.msu.edu email: rapoport@mi.uni-koeln.de
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# Higher Derivative CP(𝑁) Model and Quantization of the Induced Chern-Simons Term
## Abstract
We consider higher derivative CP($`N`$) model in $`2+1`$ dimensions with the Wess-Zumino-Witten term and the topological current density squared term. We quantize the theory by using the auxiliary gauge field formulation in the path integral method and prove that the extended model remains renormalizable in the large $`N`$ limit. We find that the Maxwell-Chern-Simons theory is dynamically induced in the large $`N`$ effective action at a nontrivial UV fixed point. The quantization of the Chern-Simons term is also discussed.
The CP($`N`$) model in $`2+1`$ dimensions has many remarkable properties. In contrast to the perturbative nonrenormalizability, the theory is renormalizable in the large $`N`$ limit, in spite of the appearance of linear divergence . It also exhibits a nontrivial fixed point structure which divides the symmetric and broken phases with a second order phase transition, and dynamical generation of a gauge boson .
The purpose of this article is to extend the CP($`N`$) model with higher order derivative terms, and discuss its possible consequence in the large $`N`$ limit. In general, adding higher derivative terms in the nonlinear sigma models would make the theory less renormalizable and the path integral more involved. However, with a specific choice of higher derivative terms to be described below, the theory admits the auxiliary gauge field formulation and an exact evaluation of the path integral can be performed in the large $`N`$ limit. We find that (I) the extended theory remains renormalizable and (II) the auxiliary gauge field becomes dynamical with induced Maxwell-Chern-Simons terms. In order to achieve these results, two types of higher derivative terms are required. One is the third order derivative Wess-Zumino-Witten (WZW) term which lives in $`M_4`$ whose boundary is our $`2+1`$ dimensional space-time. It is well known that the coefficient of this term must be quantized in the case of Grassmann coset space $`\mathrm{Gr}(N,n)\mathrm{SU}(N)/\mathrm{SU}(Nn)\times \mathrm{U}(n)`$ with $`n>1`$ . The path integral evaluation of this Grassmann model is very involved even in the large $`N`$ limit due to the non-Abelian $`\mathrm{U}(n)`$ structure, but one can truncate the theory to the Abelian $`\mathrm{U}(1)`$ sector which yields the CP($`N`$) model. In the following, we consider the CP($`N`$) model as an Abelian truncated version of the Grassmann model, and a priori quantization of WZW coefficient survives. This leads to the interesting consequence that the coefficient of the induced Chern-Simons term is also quantized. The other is the topological current density squared term which was considered recently in the higher derivative extension of CP($`N`$) model in $`1+1`$ dimensions .
Let us start directly from the auxiliary gauge field formulation of the extended CP($`N`$) model for economic presentation. We will shortly show that this theory is equivalent to the aforementioned higher derivative CP($`N`$) model. The Lagrangian is given by
$$_A=\frac{N}{G}\left[(D_\mu z)^{}(D^\mu z)i\theta GA_\mu ϵ^{\mu \nu \rho }(_\nu z)^{}(_\rho z)\lambda (z^{}z1)\right],$$
(1)
where $`D_\mu _\mu iA_\mu `$ and $`z`$ is an $`N`$ components complex scalar field which obeys a constraint $`z^{}z=1`$. The first term is the usual $`\mathrm{CP}(N)\mathrm{SU}(N)/\mathrm{SU}(N1)\times \mathrm{U}(1)`$ model in the auxiliary gauge field formulation. The second term will be responsible for higher derivative terms when the auxiliary gauge field $`A_\mu `$ are eliminated through the equations of motion. Note that $`z`$ with the constraint $`z^{}z=1`$ contains $`2N1`$ real scalars, whereas the coset space is a ($`2N2`$)-dimensional manifold. This mismatch is due to the local U(1) symmetry of the model. More specifically, the U(1) gauge transformation is given by
$$z(x)e^{i\alpha (x)}z(x),A_\mu (x)A_\mu (x)+_\mu \alpha (x).$$
(2)
Note that the first term in the Lagrangian (1) is manifestly gauge invariant, whereas the second term changes by a total derivative, hence the action is gauge invariant.
The field $`z(z_1,\mathrm{},z_N)^T`$ is separated into $`2N2`$ Nambu-Goldstone bosons $`\psi (z_1,\mathrm{},z_{N1})^T`$ associated with the spontaneously broken SU($`N`$) symmetry and Higgs bosons $`z_N(\sigma +i\chi )/\sqrt{2}`$. In general there are two possible phases : (I) $`\sigma 0`$, $`\lambda =0`$ and (II) $`\sigma =0`$, $`\lambda 0`$. In phase (I) both global SU($`N`$) and local U(1) symmetries are broken simultaneously and $`\psi `$ arise as massless Goldstone bosons. Through the Higgs mechanism $`\chi `$ turns to a longitudinal mode of massive gauge boson $`A_\mu `$. On the other hand in phase (II) both global SU($`N`$) and local U(1) symmetries are not spontaneously broken. Instead $`\psi `$ and $`z_N`$ are combined into $`z`$ with a universal mass $`\lambda ^{1/2}`$. We will see later that the dimensionless coupling $`uG\mathrm{\Lambda }`$ shows a nontrivial ultraviolet (UV) fixed point $`u^{}`$ which arises as a zero of the Callan-Symanzik $`\beta `$-function and separates the weak coupling broken phase (I) from the strong coupling symmetric phase (II). Since we are interested in dynamical generation of gauge bosons, we confine our computation to the symmetric phase alone.
Solving the equations of motion and eliminating the $`A_\mu `$ fields, we see that the Lagrangian becomes
$$_J=\frac{N}{G}\left[(_\mu z)^{}(^\mu z)J_\mu J^\mu +\theta Gϵ^{\mu \nu \rho }J_\mu _\nu J_\rho +\frac{1}{4}\theta ^2G^2J_\mu (^2g^{\mu \nu }^\mu ^\nu )J_\nu \right],$$
(3)
where $`J^\mu (1/2i)[z^{}^\mu z(^\mu z^{})z]`$ with the constraint $`z^{}z=1`$. We see that the extra third and fourth terms as well as the original second term are not renormalizable in perturbative expansion. The geometrical implication of the above Lagrangian can be seen in the coadjoint orbit approach for the nonlinear sigma model . In terms of the coadjoint orbit variable
$$Q=izz^{}+i\frac{I}{N},z^{}z=1,$$
(4)
and the topological current density $`t^\mu =ϵ^{\mu \nu \rho }<Q_\nu Q_\rho Q>=ϵ^{\mu \nu \rho }_\nu J_\rho `$, we find that the Lagrangian (3) is equivalent to
$`S_Q={\displaystyle \frac{N}{G}}{\displaystyle d^3x\left[\frac{1}{2}<_\mu Q^\mu Q>+\frac{1}{4}\theta ^2G^2t^\mu t_\mu \right]}+S_{WZW}.`$ (5)
Here the symbol $`<\mathrm{}>`$ stands for a trace over an $`N\times N`$ matrix. And $`S_{WZW}`$ descends from $`M_4`$ whose boundary is our $`2+1`$ dimensional space-time, and is given by
$$S_{WZW}=iN\theta _{M_4}d^4xϵ^{\mu \nu \rho \sigma }<Q_\mu Q_\nu Q_\rho Q_\sigma Q>,$$
(6)
where $`N\theta `$ has to be $`k/4\pi `$ for integer $`k`$ from the aforementioned quantization condition . Note that the gauge formulation requires the coefficients of the topological current squared term to be fixed in terms of $`\theta `$.
The equivalence between (3) and (5) can be conveniently checked in terms of differential forms . Let us introduce 1-forms $`dQ=(_\mu Q)dx^\mu `$ and $`JJ_\mu dx^\mu =iz^{}dz=idz^{}z`$ with $`z^{}z=1`$. Then, by the Hodge dual $``$, we have the 1-form of topological current density given by $`tt_\mu dx^\mu =dJ`$. Now, using $`idQ=dzz^{}+zdz^{}`$, the first term in (5) can be rewritten as
$`{\displaystyle \frac{1}{2}}{\displaystyle }d^3x<_\mu Q^\mu Q>={\displaystyle \frac{1}{2}}{\displaystyle }<dQdQ>={\displaystyle }(dz^{}dzJJ),`$ (7)
which reproduces the first two terms in (3). Similarly the second term in (5) can be shown to be equal to
$`{\displaystyle }d^3xt_\mu t^\mu ={\displaystyle }(tt)={\displaystyle }(dJdJ)={\displaystyle }d^3xϵ_{\alpha \beta \gamma }ϵ^{\mu \nu \gamma }(^\alpha J^\beta )(_\mu J_\nu ),`$ (8)
which corresponds to the fourth term in (3) after a partial integration. Finally, WZW action (6), which is an integration over the $`M_4`$ of the 4-form $`<Q(dQ)^4>`$, yields after a straightforward computation
$`{\displaystyle _{M_4}}<Q(dQ)^4>=i{\displaystyle _{M_4}}\left(dz^{}dzdz^{}dz\right)=i{\displaystyle _{M_4}}d(JdJ),`$ (9)
where we have used $`(z^{}dz)^4=0=(dz^{}dz)(z^{}dz)^2`$ repeatedly. It yields the third term in (3) using the Stokes’ theorem.
The original CP($`N`$) model is not renormalizable in larger than two dimensions. However the theory may become renormalizable through a resummation of Feynman diagrams in a different way from the coupling perturbation expansion. In fact the $`1/N`$ expansion provides such a resummation technique and it makes the CP($`N`$) model in less than four dimensions renormalizable.
We can rewrite the Lagrangian (1) up to total derivative terms as
$$_A^{}=\frac{N}{G}z^{}\left[^2m^2\mathrm{\Gamma }\right]z+\frac{N\lambda }{G},$$
(10)
where we separate the Goldstone boson mass $`m^2`$ from $`\lambda m^2+\stackrel{~}{\lambda }`$ and $`\mathrm{\Gamma }`$ stands for the interaction vertices:
$$\mathrm{\Gamma }\stackrel{~}{\lambda }iA_\mu (^\mu \stackrel{}{^\mu }\theta Gϵ^{\mu \nu \rho }\stackrel{}{_\nu }_\rho )A_\mu A^\mu ,$$
(11)
where $`\stackrel{}{^\mu }`$ and $`\stackrel{}{_\nu }`$ do not operate on $`A_\mu `$. Path integrating $`z`$ and $`z^{}`$ provides the large $`N`$ effective action:
$$S_{\mathrm{eff}}=d^3x_A^{}+iN\mathrm{Tr}\mathrm{Ln}\left[^2m^2\right]iN\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n}\mathrm{Tr}\left[\frac{1}{^2m^2}\mathrm{\Gamma }\right]^n.$$
(12)
We divide the effective action up to quadratic terms ($`n=1,2`$) into two parts, $`S_{\mathrm{eff}}=S^I+S^{II}`$ where $`S^I`$ denotes the large $`N`$ effective action in the original CP($`N`$) model and $`S^{II}`$ stands for the extra induced terms. After some straightforward calculations, we obtain
$`S^I`$ $`=`$ $`N{\displaystyle }d^3x[{\displaystyle \frac{1}{G}}z^{}[^2m^2\mathrm{\Gamma }^I]z+({\displaystyle \frac{1}{u}}{\displaystyle \frac{1}{u^{}}})\mathrm{\Lambda }(m^2+\stackrel{~}{\lambda })`$ (14)
$`+{\displaystyle \frac{1}{4\pi }}m\stackrel{~}{\lambda }+{\displaystyle \frac{1}{6\pi }}m^3+{\displaystyle \frac{1}{2}}\stackrel{~}{\lambda }\mathrm{\Pi }_\lambda (i)\stackrel{~}{\lambda }{\displaystyle \frac{1}{4}}F_{\mu \nu }\mathrm{\Pi }_1(i)F^{\mu \nu }],`$
$`S^{II}`$ $`=`$ $`N{\displaystyle }d^3x[iA_\mu ϵ^{\mu \nu \rho }(_\nu z)^{}(_\rho z){\displaystyle \frac{1}{12}}F_{\mu \nu }\theta ^2G^2\mathrm{\Pi }_2(i)F^{\mu \nu }`$ (16)
$`{\displaystyle \frac{2}{3}}A_\mu \theta G\mathrm{\Pi }_2(i)ϵ^{\mu \nu \rho }_\nu A_\rho ],`$
where we have introduced the dimensionless coupling $`uG\mathrm{\Lambda }`$ and $`u^{}2\pi ^2`$. $`\mathrm{\Gamma }^I`$ is the interaction vertices in the original CP($`N`$) model without higher derivative terms, and the vacuum polarization functions are given by
$`\mathrm{\Pi }_\lambda (p)`$ $`=`$ $`{\displaystyle \frac{1}{4\pi \sqrt{p^2}}}\mathrm{arctan}{\displaystyle \frac{\sqrt{p^2}}{2m}},`$ (17)
$`\mathrm{\Pi }_1(p)`$ $`=`$ $`{\displaystyle \frac{1}{4\pi p^2}}\left[m{\displaystyle \frac{4m^2p^2}{2\sqrt{p^2}}}\mathrm{arctan}{\displaystyle \frac{\sqrt{p^2}}{2m}}\right],`$ (18)
$`\mathrm{\Pi }_2(p)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi ^2}}\mathrm{\Lambda }{\displaystyle \frac{3}{8\pi }}m+{\displaystyle \frac{3}{4}}p^2\mathrm{\Pi }_1(p).`$ (19)
We realize that there arise linear divergences in induced Chern-Simons and Maxwell terms which have no counter terms in the original Lagrangian.
Renormalization of the coupling $`G`$ can be worked out in the same manner as in the original CP($`N`$) model. The large $`N`$ effective potential is defined as the effective action divided by $`\mathrm{\Omega }d^3x`$ with $`\stackrel{~}{\lambda }0`$, $`A_\mu 0`$ and $`z(z^{})0`$. It is given by
$$\frac{1}{N}V_{\mathrm{eff}}=(\frac{1}{u}\frac{1}{u^{}})\mathrm{\Lambda }m^2\frac{m^3}{6\pi }.$$
(20)
The Goldstone boson mass $`m`$ is determined as a nontrivial solution to the gap equation $`dV_{\mathrm{eff}}/dm^2=0`$ and reads
$$m=4\pi \mathrm{\Lambda }\left(\frac{1}{u^{}}\frac{1}{u}\right).$$
(21)
We notice that $`m`$ can be independent of the ultraviolet cutoff $`\mathrm{\Lambda }`$ by imposing $`\mathrm{\Lambda }`$ dependence on the coupling $`u`$. In fact the scale invariance condition $`\mathrm{\Lambda }dm/d\mathrm{\Lambda }=0`$ leads us to the Callan-Symanzik $`\beta `$-function
$$\beta (u)\mathrm{\Lambda }\frac{du}{d\mathrm{\Lambda }}=u\left(1\frac{u}{u^{}}\right),$$
(22)
which shows a nontrivial UV fixed point at $`u=u^{}`$. In the original CP($`N`$) model the only divergence is the one which arises in the large $`N`$ effective action through a tadpole diagram coupled with $`\stackrel{~}{\lambda }`$ so that the scale invariance condition $`\mathrm{\Lambda }dm/d\mathrm{\Lambda }=0`$ is enough to achieve the cutoff independent theory. Since $`m`$ is scale invariant, the solution to the gap equation (21) suggests that the renormalization of coupling is given by
$$\left(\frac{1}{u}\frac{1}{u^{}}\right)\mathrm{\Lambda }=\left(\frac{1}{u_R}\frac{1}{u_R^{}}\right)\mu ,$$
(23)
where $`u_R`$ is the renormalized coupling at a reference energy scale $`\mu `$.
In the extended model, however, linear divergences arise in the induced Chern-Simons and Maxwell terms which do not have their counter terms in the classical action. Therefore the higher derivative theory seems to be nonrenormalizable, although the coupling $`u`$ can be renormalized in the same way as in the original CP($`N`$) model. However, since the extra linear divergences are always accompanied by the coupling $`Gu/\mathrm{\Lambda }`$ which cancels the linear divergences, we expect the large $`N`$ effective action (16) to be scale invariant in the continuum limit $`\mathrm{\Lambda }\mathrm{}`$.
To study this point in more detail, we first look at how $`S^I`$ can be scale invariant through the renormalization procedure. The induced kinetic terms of $`A_\mu `$ and $`\stackrel{~}{\lambda }`$ are UV finite in themselves so that we do not need wave function renormalization for them. Then the second term in the right hand side of Eq. (14) becomes UV finite through Eq. (23) from which the Z factor for the coupling $`G`$ can be read
$$Z^1\frac{G_R}{G}=1\frac{u_R}{u_R^{}}+\frac{u_R}{u^{}}\left(\frac{\mathrm{\Lambda }}{\mu }\right).$$
(24)
Here $`G_R`$ is connected to the dimensionless coupling $`u_R`$ through $`u_RG_R\mu `$. The kinetic term of $`z`$ has to be UV finite in itself so that we see
$$\frac{1}{G}z^{}\left[^2m^2\right]z=\frac{1}{G_R}z_R^{}\left[^2m^2\right]z_R,$$
(25)
where $`z_R`$ has been introduced through $`z=Z_z^{1/2}z_R`$ and $`Z_z`$ is thereby determined as $`Z_zZ`$ in order to cancel the Z factor from the coupling renormalization. Thus we realize that $`\mathrm{\Gamma }^I`$ in $`S^I`$ has to remain invariant through the renormalization procedure. This forces both of $`A_\mu `$, $`\stackrel{~}{\lambda }`$ to be unchanged through renormalization. This is consistent with the UV finiteness of kinetic terms for $`A_\mu `$ and $`\stackrel{~}{\lambda }`$ in $`S^I`$.
Now that we have explicitly shown that the original CP($`N`$) model can be renormalized through the $`1/N`$ resummation, let us look at what happens in $`S^{II}`$ if we take the continuum limit $`\mathrm{\Lambda }\mathrm{}`$. The Callan-Symanzik $`\beta `$-function (22) tells us that the dimensionless coupling $`uG\mathrm{\Lambda }`$ goes to the UV fixed point $`u=u^{}`$ as the cutoff $`\mathrm{\Lambda }`$ goes to infinity, whereas the bare coupling $`Gu/\mathrm{\Lambda }`$ reduces to zero. On the other hand, the extra vacuum polarization effect $`\mathrm{\Pi }_2`$ in Eq. (19) can be rewritten as
$$G\mathrm{\Pi }_2(p)=\frac{u}{u^{}}G\left[\frac{3}{8\pi }m\frac{3}{4}p^2\mathrm{\Pi }_1(p)\right].$$
(26)
Therefore, we can obtain the UV finite result $`G\mathrm{\Pi }_2(p)1`$ as $`\mathrm{\Lambda }\mathrm{}`$. In the same reasoning the extra contribution to the Maxwell term which contains $`G^2\mathrm{\Pi }_2(p)`$ vanishes in the continuum limit. Moreover, the first term in $`S^{II}`$ has an extra $`G`$ after the renormalization (25) and is thereby suppressed by a factor $`1/\mathrm{\Lambda }`$ as $`\mathrm{\Lambda }\mathrm{}`$. Thus in the continuum limit $`S^{II}`$ becomes a UV finite Chern-Simons action $`(2N\theta /3)d^3xϵ^{\mu \nu \rho }A_\mu _\nu A_\rho `$ without any ambiguity.
According to power counting of superficial degrees of UV divergences, the three-points function of $`A_\mu `$ shows a linear divergence. However the extra gauge coupling in $`\mathrm{\Gamma }`$ is invariant under the charge conjugation so that the three points function and its charge conjugation cancel each other in the same way as in the original CP($`N`$) model . Moreover all $`n`$-points functions with $`n4`$ are UV finite and the contributions from the extra gauge interaction are accompanied by $`G`$. Therefore in the continuum limit such extra effects are suppressed by a factor $`1/\mathrm{\Lambda }`$ and become irrelevant. Hence we can conclude that the large $`N`$ effective action in our extended model is renormalizable and the gauge sector is equivalent to the Maxwell-Chern-Simons theory which couples minimally to $`z`$ field. The effective Lagrangian at the lowest order in the derivative expansion can be rewritten as
$$_{\mathrm{eff}}=\frac{1}{4}_{\mu \nu }^{\mu \nu }\frac{\kappa }{2}ϵ^{\mu \nu \rho }𝒜_\mu _\nu 𝒜_\rho ,$$
(27)
where we introduced $`e^224\pi m/N`$ and redefined the gauge field by $`A_\mu =e𝒜_\mu `$ so that the covariant derivative is $`D_\mu _\mu ie𝒜_\mu `$. In this Lagrangian, the Chern-Simons coefficient becomes a dimensionful parameter $`\kappa 32\pi m\theta `$, and we supposed that the Goldstone bosons became very massive and decoupled after the symmetry restoration of both global $`\mathrm{SU}(N)`$ and local $`\mathrm{U}(1)`$.
Note that the ratio of the topological gauge boson mass $`\kappa `$ to the effective gauge coupling square $`e^2`$ is proportional to $`N\theta `$ and is quantized such as $`\kappa /e^2=k/3\pi `$. Recently, some authors showed that in the Maxwell-Chern-Simons theory coupled to fermion fields, the Lorentz symmetry is spontaneously broken through a dynamically induced magnetic field when $`\kappa /e^2`$ is quantized in a unit of $`1/4\pi `$ . Specifically, Ref. introduced $`N_f`$-flavored four-components fermion fields and showed that only the Lorentz symmetry is broken in $`\kappa =N_fe^2/2\pi `$, whereas both flavor U$`(2N_f)`$ and Lorentz symmetries are broken at the same time in $`\kappa =N_fe^2/4\pi `$ through a dynamically generated fermion mass and a magnetic field. Our current results provides a geometrical origin of quantization of the Chern-Simons coefficient due to the WZW term. In fact the quantization condition $`\kappa =N_fe^2/2\pi `$ and $`N_fe^2/4\pi `$ correspond to $`2k=3N_f`$ and $`4k=3N_f`$, respectively. We may also possibly prove that the induced Chern-Simons term may not receive any higher order $`1/N`$ corrections .
We have investigated the gauge formulation of higher derivative CP($`N`$) model in $`2+1`$ dimensions with WZW term and topological current density squared, and have proved its renormalizability in the large $`N`$ limit. We also have found that the Maxwell-Chern-Simons theory is dynamically generated in the effective action and the coefficient of the induced Chern-Simons term must be quantized which is a direct consequence of the quantization of the WZW term. If we couple the theory to fermions with U$`(2N_f)`$ flavors, the low energy effective theory shows spontaneous break down of the Lorentz symmetry associated with an induced magnetic field when $`2k=3N_f`$, $`4k=3N_f`$. It would be interesting to check whether this has some physical implications, especially in condensed matter phenomena such as anyon physics and the fractional quantum Hall effect .
There exist related subjects to be studied further. One of them is to compute the current algebra associated with (5) and (6), and to check whether it has some special properties at the fixed point $`u=u^{}`$. Another interesting problem is to consider possible extension to higher dimensions. The theory in $`2+1`$ dimensions was special in the sense that the topological current $`t^\mu `$ interacts with original gauge field of CP($`N`$) model. In $`D`$ dimensional extension, however, one would need extra $`D2`$ rank antisymmetric tensor fields $`B_{\mu _3\mathrm{}\mu _D}`$ with interaction
$`iϵ^{\mu _1\mu _2\mu _3\mathrm{}\mu _D}(_{\mu _1}z)^{}(_{\mu _2}z)B_{\mu _3\mathrm{}\mu _D}+{\displaystyle \frac{1}{2M^{D4}}}B_{\mu _3\mathrm{}\mu _D}B^{\mu _3\mathrm{}\mu _D}.`$ (28)
This theory which corresponds to the auxiliary field formulation of the dual variable description of SU(2) Yang-Mills theory in the infrared limit for CP(2) case, is not renormalizable in $`D4`$ even in $`1/N`$ expansion. So we have to consider it as an effective theory which describes a massless gauge field interacting with massive $`H=dB`$ field in the $`B^{}F`$-type interaction. It remains to investigate whether this type of Maxwell-Kalb-Ramond theory has some relevance with quark confinement in $`3+1`$ dimensions .
T.I. was supported by KOSEF Postdoctoral Fellowship and Korea Research Center for Theoretical Physics and Chemistry. P.O. was supported by the Samsung Research Fund, Sungkyunkwan University, 1999, and in part by the KOSEF through the project number 2000-1-11200-001-3. This work was also partially supported by BK21 Physics Research Program.
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# Actions and symmetries of NSR superstrings and D-strings
## a Field content and gauge symmetries.
The results reported here were obtained by a BRST-cohomological analysis whose details will be presented elsewhere. The input of this analysis is the field content and a set of gauge transformations of the fields. The fields are those of the $`2d`$ supergravity multiplet $`\{e_m^a,\chi _m^\alpha ,S\}`$ ($`m,a,\alpha `$ are $`2d`$ world-sheet, Lorentz and spinor indices, respectively), an arbitrary number of “matter multiplets” $`\{X^M,\psi _\alpha ^M,F^M\}`$ and an arbitrary number of gauge multiplets $`\{A_m^i,\lambda _\alpha ^i,\varphi ^i\}`$. $`e_m^a`$ and $`\chi _m^\alpha `$ are the world-sheet zweibein and gravitino, respectively, the $`A_m^i`$ are abelian gauge fields, and the $`X^M`$ may be regarded as coordinates of an enlarged target space. All fields, including the auxiliary fields are real (we use a Majorana-Weyl basis for the $`\gamma `$-matrices and the signature of the world-sheet metric is $`\eta _{ab}=\mathrm{diag}(1,1)`$; furthermore $`\epsilon ^{01}=\epsilon _{10}=1`$ and $`\gamma _{}`$ is defined through $`\gamma ^a\gamma ^b=\eta ^{ab}1\mathrm{l}+\epsilon ^{ab}\gamma _{}`$; our formulae apply with appropriate redefinitions also to Euclidean signature). We impose the following gauge symmetries: world-sheet diffeomorphisms, local (1,1) world-sheet supersymmetry, local $`2d`$ Lorentz invariance, invariance under abelian gauge transformations of the $`A_m^i`$, Weyl and super-Weyl invariance, and invariance under arbitrary local shifts of the field $`S`$. The corresponding gauge transformations of the fields are, written as BRST transformations,
$`se_m^a`$ $`=`$ $`\xi ^n_ne_m^a+(_m\xi ^n)e_n^a+C_b^ae_m^b+C^We_m^a`$ (2)
$`2\mathrm{i}\xi ^\alpha \chi _m^\beta (\gamma ^aC)_{\alpha \beta }`$
$`s\chi _m^\alpha `$ $`=`$ $`\xi ^n_n\chi _m^\alpha +(_m\xi ^n)\chi _n^\alpha \frac{1}{4}C^{ab}\chi _m^\beta \epsilon _{ab}(\gamma _{})_\beta ^\alpha `$ (5)
$`+\frac{1}{2}C^W\chi _m^\alpha +\mathrm{i}\eta ^\beta (\gamma _m)_\beta ^\alpha `$
$`+_m\xi ^\alpha +\frac{1}{4}\xi ^\beta (\omega _m^{ab}\epsilon _{ab}(\gamma _{})_\beta ^\alpha (\gamma _m)_\beta ^\alpha S)`$
$`sS`$ $`=`$ $`\xi ^n_nSC^WS+W+\mathrm{i}\xi ^\gamma (\gamma ^mC)_{\gamma \alpha }\chi _m^\alpha S`$ (7)
$`4\xi ^\gamma (\gamma _{}C)_{\gamma \alpha }\epsilon ^{nm}(_n\chi _m^\alpha +\frac{1}{4}\omega _n^{ab}\epsilon _{ab}\chi _m^\beta (\gamma _{})_\beta ^\alpha )`$
$`sX^M`$ $`=`$ $`\xi ^m_mX^M+\xi ^\alpha \psi _\alpha ^M`$ (8)
$`s\psi _\alpha ^M`$ $`=`$ $`\xi ^m_m\psi _\alpha ^M+\frac{1}{4}C^{ab}\epsilon _{ab}(\gamma _{})_\alpha ^\beta \psi _\beta ^M\frac{1}{2}C^W\psi _\alpha ^M`$ (10)
$`+\xi ^\beta C_{\beta \alpha }F^M\mathrm{i}\xi ^\beta (\gamma ^mC)_{\beta \alpha }(_mX^M\chi _m^\gamma \psi _\gamma ^M)`$
$`sF^M`$ $`=`$ $`\xi ^m_mF^MC^WF^M+\xi ^\alpha (\gamma ^m)_\alpha ^\beta \{_m\psi _\beta ^M`$ (13)
$`\frac{1}{4}\omega _m^{ab}\epsilon _{ab}(\gamma _{})_\beta ^\gamma \psi _\gamma ^M\chi _m^\gamma C_{\gamma \beta }F^M`$
$`+\mathrm{i}\chi _m^\gamma (\gamma ^nC)_{\gamma \beta }(_nX^M\chi _n^\delta \psi _\delta ^M)\}`$
$`sA_m^i`$ $`=`$ $`\xi ^n_nA_m^i+(_m\xi ^n)A_n^i+_mC^i`$ (15)
$`2\mathrm{i}\xi ^\alpha \chi _m^\beta (\gamma _{}C)_{\beta \alpha }\varphi ^i\xi ^\alpha (\gamma _m)_\alpha ^\beta \lambda _\beta ^i`$
$`s\varphi ^i`$ $`=`$ $`\xi ^n_n\varphi ^iC^W\varphi ^i+\xi ^\alpha (\gamma _{})_\alpha ^\beta \lambda _\beta ^i`$ (16)
$`s\lambda _\beta ^i`$ $`=`$ $`\xi ^n_n\lambda _\beta ^i+\frac{1}{4}C^{ab}\epsilon _{ab}(\gamma _{})_\beta ^\gamma \lambda _\gamma ^i\frac{3}{2}C^W\lambda _\beta ^i`$ (20)
$`+\mathrm{i}\xi ^\alpha \{(\gamma _{}\gamma ^mC)_{\alpha \beta }(_m\varphi ^i\chi _m\gamma _{}\lambda ^i)`$
$`+(\gamma _{}C)_{\alpha \beta }\epsilon ^{mn}(_mA_n^i+\chi _m\gamma _n\lambda ^i\mathrm{i}\chi _n\gamma _{}C\chi _m\varphi ^i)`$
$`+(\gamma _{}C)_{\alpha \beta }S\varphi ^i\}+2\eta _{SW}^\alpha (\gamma _{}C)_{\alpha \beta }\varphi ^i`$
where $`\xi ^m`$ are the ghosts of world-sheet diffeomorphisms, $`\xi ^\alpha `$ are the supersymmetry ghosts, $`C^{ab}`$ is the Lorentz ghost, $`C^W`$ and $`\eta _{SW}^\alpha `$ are the Weyl and super-Weyl ghosts, respectively, $`C^i`$ are the ghosts associated with the gauge fields $`A_m^i`$, and $`W`$ is the ghost corresponding to the local shifts of the auxiliary field $`S`$. A $`C`$ without any index denotes the charge conjugation matrix. These gauge transformations were obtained from an analysis of the $`2d`$ supergravity algebra in presence of the matter and gauge multiplets (the analysis is analogous to the superspace analysis in ). The corresponding BRST transformations of the ghosts are such that $`s^2=0`$.
## b Gauge invariant actions.
Owing to the use of auxiliary fields, the algebra of the gauge transformations closes off-shell. As a consequence, neither the BRST transformations (20), nor the BRST transformations of the ghosts contain antifields. This allows one to determine the action functionals which are invariant under the gauge transformations (20) by computing the cohomology of $`s`$ in the space of antifield independent local functionals with ghost number 0. Our result of this computation is the following. The most general Lagrangian $`L`$ which transforms under the above gauge transformations into a total derivative and which is polynomial in the derivatives of the fields is, modulo total derivatives,
$`L/e`$ $`=`$ $`\frac{1}{2}_mX^M_nX^N(h^{mn}G_{MN}+\epsilon ^{mn}B_{MN})`$ (29)
$`+\frac{\mathrm{i}}{2}\overline{\psi }^M\gamma ^m_m\psi ^NG_{MN}+\frac{1}{2}F^MF^NG_{MN}`$
$`+\chi _k\gamma ^n\gamma ^k(\psi ^N_nX^M\frac{1}{4}C\chi _n\overline{\psi }^M\psi ^N)G_{MN}`$
$`+(\frac{1}{2}F^M\overline{\psi }^K\psi ^N\mathrm{i}\overline{\psi }^N\gamma ^m\psi ^M_mX^K)\mathrm{\Gamma }_{NKM}`$
$`+\frac{1}{4}(F^M\overline{\psi }^K\gamma _{}\psi ^N\mathrm{i}\overline{\psi }^N\gamma ^m\gamma _{}\psi ^M_mX^K)H_{NKM}`$
$`\frac{\mathrm{i}}{12}\chi _m\gamma ^n\gamma ^m\psi ^M\overline{\psi }^N\gamma _n\gamma _{}\psi ^KH_{MNK}`$
$`+\frac{1}{16}\overline{\psi }^M(1\mathrm{l}+\gamma _{})\psi ^N\overline{\psi }^K(1\mathrm{l}+\gamma _{})\psi ^LR_{KMLN}`$
$`+\epsilon ^{mn}D_i_mA_n^i+\frac{\mathrm{i}}{4}\overline{\psi }^M\psi ^N\varphi ^i_N_MD_i`$
$`+\frac{1}{2}(\mathrm{i}\overline{\psi }^N\gamma _{}\lambda ^i\mathrm{i}F^N\varphi ^i+\chi _m\gamma ^m\psi ^N\varphi ^i)_ND_i,`$
where $`e=det(e_m^a)`$, and $`h_{mn}`$ is the world-sheet metric built from $`e_m^a`$. $`G_{MN}`$, $`B_{MN}`$ and $`D_i`$ are arbitrary functions of the $`X^M`$, except that $`G_{MN}`$ and $`B_{MN}`$ are symmetric and antisymmetric in $`M`$ and $`N`$, respectively, and $`\mathrm{\Gamma }_{MNK},H_{MNK}`$ and $`R_{MNKL}`$ are given by
$`\mathrm{\Gamma }_{MNK}`$ $`=`$ $`\frac{1}{2}(_MG_{NK}+_NG_{MK}_KG_{MN})`$
$`H_{MNK}`$ $`=`$ $`_MB_{NK}+_NB_{KM}+_KB_{MN}`$
$`R_{MNKL}`$ $`=`$ $`_K_{[M}(G+B)_{N]L}_L_{[M}(G+B)_{N]K}`$
where $`_M`$ denotes differentiation with respect to $`X^M`$. The redefinitions $`B_{MN}(X)B_{MN}(X)+_{[M}f_{N]}(X)`$ and $`D_i(X)D_i(X)+k_i`$ modify the Lagrangian only by total derivatives for any functions $`f_N(X)`$ and constants $`k_i`$.
Apart from the terms containing fields of the gauge multiplets, the Lagrangian (29) agrees with the one derived in when one eliminates the auxiliary fields $`F^M`$. Hence, the cohomological analysis shows that the Lagrangian derived in is in fact unique in absence of gauge multiplets (modulo total derivatives, and up to the choice of $`G_{MN}`$ and $`B_{MN}`$). It should be noted, however, that this uniqueness is tied to the gauge transformations (20) and may get lost when one allows that the gauge transformations get consistently deformed. E.g., one would expect that the world-sheet diffeomorphisms and supersymmetry transformations can be nontrivially deformed such that the deformed action is invariant under the deformed transformations if the background has special isometries, by analogy with the purely bosonic case . Furthermore, there are certainly actions invariant under nonabelian gauge transformations of the $`A_m^i`$ (leading to nonabelian Born-Infeld actions, among others). The general deformation problem is currently under study.
## c Simplified action.
For further discussion we shall assume in the following that the functions $`D_i`$ coincide with a subset of the fields $`X^M`$. We denote this subset by $`\{y^i\}`$ and the remaining $`X`$’s by $`x^\mu `$,
$$\{X^M\}=\{x^\mu ,y^i\},D_i=y^i.$$
(30)
In fact, this assumption is a very mild one because, except at stationary points of $`D_i(X)`$, (30) can be achieved by a field redefinition $`X^M\stackrel{~}{X}^M=\stackrel{~}{X}^M(X)`$, where this “coordinate transformation” is such that each nonconstant $`D_i(X)`$ becomes one of the $`\stackrel{~}{X}`$’s. Indeed, constant $`D_i`$ give only constributions to the Lagrangian which are total derivatives and can thus be neglected, at least classically; nonconstant $`D_i`$ can be assumed to be independent by a suitable choice of basis for the gauge fields and may thus be taken as $`\stackrel{~}{X}`$’s, at least locally (e.g., if $`D_1=D_2`$, the Lagrangian depends only on the combination $`A_m^1+A_m^2`$ which can be introduced as a new gauge field).
It is now easy to see that the Lagrangian (29) can actually be simplified by setting the fields $`\psi ^i,F^i,\lambda ^i,\varphi ^i`$ to zero. Indeed, owing to (30), the classical equations of motion for $`\lambda ^i`$ and $`\varphi ^i`$ yield $`\psi ^i=0`$ and $`F^i=0`$. The latter equations are algebraic and can be used in the Lagrangian. Then the Lagrangian does not contain $`\lambda ^i`$ and $`\varphi ^i`$ anymore and the only remnant of the gauge multiplets are the terms $`e\epsilon ^{mn}y^i_mA_n^i`$. This reflects that the gauge multiplets carry no dynamical degrees of freedom since the world-sheet is 2-dimensional. Of course, the transformations (20) must be adapted in order to provide the gauge symmetries of the simplified Lagrangian: those fields that are eliminated from the action must also be eliminated from the transformations of the remaining fields using the equations of motion of the eliminated fields. This only affects the supersymmetry transformations of $`y^i`$ and $`A_m^i`$. The new supersymmetry transformation of $`y^i`$ is then simply zero. This is not in contradiction with the supersymmetry algebra because the equations of motion for the $`A_m^i`$ give $`_my^i=0`$ (of course, after eliminating the fields $`\psi ^i,F^i,\lambda ^i,\varphi ^i`$, the supersymmetry algebra holds only on-shell). The $`y^i`$ are thus constant on-shell, their values being integration constants fixed only by initial conditions. This leads to the interpretation of the $`y^i`$ as coordinates of “frozen extra dimensions” mentioned in the beginning.
## d Born-Infeld actions.
Locally supersymmetric Born-Infeld actions arise from (29) for particular choices of $`G_{MN}`$ and $`B_{MN}`$, in complete analogy to the purely bosonic case . For instance, consider the case with only one gauge field ($`\{A_m^i\}=\{A_m\}`$, $`\{y^i\}=\{y\}`$) and the following particular choice of $`G_{MN}`$ and $`B_{MN}`$,
$`G_{\mu \nu }=\sqrt{1+y^2}g_{\mu \nu }(x),`$ $`B_{\mu \nu }=yb_{\mu \nu }(x)`$ (31)
$`G_{yy}=G_{y\mu }=B_{y\mu }=0.`$ (32)
$`y`$ can be eliminated algebraically. Eliminating also the world-sheet zweibein $`e_m^a`$, the Lagrangian becomes
$`L=\pm \sqrt{det(g_{mn}+_{mn})}+\mathrm{}`$ (33)
$`g_{mn}=g_{\mu \nu }(x)_mx^\mu _nx^\nu `$ (34)
$`_{mn}=_mA_n_nA_mb_{\mu \nu }(x)_mx^\mu _nx^\nu `$ (35)
where we have assumed $`det(g_{mn})<0`$ and $`det(g_{mn}+_{mn})<0`$, and the nonwritten terms involve fermions. The Born-Infeld Lagrangian $`L=\pm [+det(g_{mn}+_{mn})]^{1/2}+\mathrm{}`$ for Euclidean signature ($`det(g_{mn})>0`$) corresponds to $`G_{\mu \nu }=(1y^2)^{1/2}g_{\mu \nu }(x)`$.
## e Global symmetries.
Our second result concerns the global symmetries of the action (29). These can be obtained from the BRST cohomology in the space of antifield dependent local functionals with ghost number $`1`$ . We have computed this cohomology completely and present now the resulting global symmetries for the simplified form of the action described above (without the fields $`\psi ^i,F^i,\lambda ^i,\varphi ^i`$ and assuming (30)). The nontrivial global symmetries (a global symmetry is called trivial when it is equal to a gauge transformation on-shell) are generated by the following transformations,
$`\mathrm{\Delta }e_m^a`$ $`=`$ $`0,\mathrm{\Delta }\chi _m^\alpha =0`$ (36)
$`\mathrm{\Delta }X^M`$ $`=`$ $`^M,^i=K^i(y),^\mu =V^\mu (X)`$ (37)
$`\mathrm{\Delta }\psi _\alpha ^\mu `$ $`=`$ $`\psi _\alpha ^\nu _\nu V^\mu (X)`$ (38)
$`\mathrm{\Delta }F^\mu `$ $`=`$ $`F^\nu _\nu V^\mu (X)+\frac{1}{2}\overline{\psi }^\nu \psi ^\lambda _\nu _\lambda V^\mu (X)`$ (39)
$`\mathrm{\Delta }A_m^i`$ $`=`$ $`b_M^i(X)_mX^M+a_M^i(X)\epsilon _{m}^{}{}_{}{}^{n}_nX^M`$ (42)
$`\delta _{jk}A_m^j_iK^k(y)\chi _n\gamma _m\gamma ^n\gamma _{}\psi ^\mu a_\mu ^i(X)`$
$`+\frac{\mathrm{i}}{2}\overline{\psi }^\mu \gamma _m\{\gamma _{}\psi ^\nu _{[\nu }a_{\mu ]}^i(X)\psi ^\nu _{[\nu }b_{\mu ]}^i(X)\}`$
where $`^M`$, $`a_M^i`$ and $`b_M^i`$ have to solve the following generalized Killing vector equations,
$`_{}G_{MN}`$ $`=`$ $`2\delta _{i(M}a_{N)}^i,`$ (43)
$`_{}B_{MN}`$ $`=`$ $`2_{[M}p_{N]}+2\delta _{i[M}b_{N]}^i`$ (44)
for some functions $`p_M(X)`$ ($`_{}`$ is the Lie derivative along $`^M`$ and $`\delta _{iM}`$ is the Kronecker symbol). Note that the $`p_M`$ do not occur in the $`\mathrm{\Delta }`$-transformations; however, they do contribute to the corresponding Noether currents.
The equations (44) are actually the same as the equations which also determine the symmetries of bosonic string and D-string actions with the specification (30). In the absence of gauge fields (no $`A_m^i`$, $`y^i`$, $`K^i`$; $`\{^M\}\{V^\mu \}`$), they read
$`_VG_{\mu \nu }=0,_VB_{\mu \nu }=2_{[\mu }p_{\nu ]}.`$ (45)
These equations had been already discussed in . The first equation (45) is just the standard Killing vector equation for $`G_{\mu \nu }`$. Hence, the solutions of equations (45) are those Killing vector fields of $`G_{\mu \nu }`$ which solve the second equation (45) (for some $`p_\mu `$).
The situation changes when gauge fields are present. Then equations (44) read for $`M,N=\mu ,\nu `$:
$$(_V+K^i_i)G_{\mu \nu }=0,(_V+K^i_i)B_{\mu \nu }=2_{[\mu }p_{\nu ]},$$
(46)
where $`_V`$ is the Lie derivative along the vector field $`V^M`$ given by $`V^i=0`$, $`V^\mu =V^\mu (X)`$. The remaining equations (44) just determine the functions $`a_M^i`$ and $`b_M^i`$,
$`a_\mu ^i=_{}G_{\mu i},a_i^j=\frac{1}{2}_{}G_{ij}`$ (47)
$`b_\mu ^i=_{}B_{\mu i}+_ip_\mu ,b_i^j=\frac{1}{2}_{}B_{ij}.`$ (48)
Here we have used that $`p_i`$ and the parts of $`a_i^j`$ resp. $`b_i^j`$ which are antisymmetric resp. symmetric in $`i,j`$ can be set to zero without loss of generality (the corresponding contributions to $`\mathrm{\Delta }`$ can be removed by subtracting trivial global symmetries from $`\mathrm{\Delta }`$).
The global symmetries are thus completely determined by equations (46). Note that these equations reproduce (45) for $`K^i=0`$, except that now $`G_{\mu \nu }`$ and $`B_{\mu \nu }`$ depend in general not only on the $`x^\mu `$ but also on the $`y^i`$. Hence, in general $`V^\mu `$ and $`p_\mu `$ also depend on the $`y^i`$. For the discussion of equations (46), the $`y^i`$ may be viewed as parameters of $`G_{\mu \nu }`$ and $`B_{\mu \nu }`$ rather than as coordinates of extra dimensions. Solutions to equations (46) with $`K^i=0`$ can thus be regarded as solutions to equations (45) for parameter-dependent $`G_{\mu \nu }`$ and $`B_{\mu \nu }`$. In contrast, solutions to (46) with $`K^i0`$ have no counterparts among the solutions of (45). Such solutions may be called “dilatational” solutions, because in special cases they are true dilatations, as we will see in the example below (further examples can be found in ).
Finally we note that the solutions to equations (46) come in infinitely big families and that, as a consequence, the corresponding commutator algebra of the global symmetries is an infinite dimensional loop-like algebra. This has been observed already in and is a consequence of the fact that the action depends on the $`A_m^i`$ only via their field strengths . All members of a family of solutions arise from one solution $`V^\mu (X)`$, $`K^i(y)`$, $`p_\mu (X)`$ by multiplying that solution with arbitrary functions of the $`y^i`$. As the $`y^i`$ are constant on-shell, this infinite dimensionality of the space of global symmetries has no practical importance, i.e., in order to discuss the global symmetries it is sufficient to consider just one representative of each family.
## f Example.
To illustrate the results presented above, we specify them for a simple class of models characterized by Lagrangians containing only one $`U(1)`$ gauge field $`A_m`$ and the following choices for the background
$`G_{yM}=B_{y\mu }=0,G_{\mu \nu }=f(y)\eta _{\mu \nu },B_{\mu \nu }=B_{\mu \nu }(y),`$
leading to
$`L/e`$ $`=`$ $`\frac{1}{2}h^{mn}_mx^\mu _nx^\nu G_{\mu \nu }+\frac{1}{2}\epsilon ^{mn}_mx^\mu _nx^\nu B_{\mu \nu }`$
$`+\chi _k\gamma ^n\gamma ^k\psi ^\nu _nx^\mu G_{\mu \nu }\frac{1}{4}\chi _k\gamma ^n\gamma ^kC\chi _n\overline{\psi }^\mu \psi ^\nu G_{\mu \nu }`$
$`+\frac{\mathrm{i}}{2}\overline{\psi }^\mu \gamma ^m_m\psi ^\nu G_{\mu \nu }\frac{\mathrm{i}}{4}\overline{\psi }^\nu \gamma ^m\gamma _{}\psi ^\mu _my_yB_{\mu \nu }`$
$`+\frac{1}{2}\epsilon ^{mn}(_mA_n_nA_m)y`$
where the auxiliary fields $`F^\mu `$ have been eliminated. As shown in , in this case the general solution of equations (44) is (modulo trivial global symmetries)
$`K`$ $`=`$ $`2r(y)`$
$`V^\mu `$ $`=`$ $`r(y)[\mathrm{ln}f(y)]^{}x^\mu +r^\mu (y)+r^{[\mu \lambda ]}(y)\eta _{\lambda \nu }x^\nu `$
$`a_\mu `$ $`=`$ $`V^\lambda {}_{}{}^{}f(y)\eta _{\mu \lambda },a_y=0`$
$`b_\mu `$ $`=`$ $`(r(y)B_{\mu \nu }^{})^{}x^\nu +B_{\mu \nu }^{}V^\nu ,b_y=0`$
$`p_\mu `$ $`=`$ $`r(y)B_{\mu \nu }^{}x^\nu +B_{\mu \nu }V^\nu ,p_y=0`$
where a prime denotes differentiation with respect to $`y`$ and $`r(y)`$, $`r^\mu (y)`$ and $`r^{[\mu \lambda ]}(y)`$ are arbitrary functions of $`y`$ and correspond to families of dilatations, translations and Lorentz-transformations in target space, respectively. For three reasons the dilatations are special: (i) as discussed already above, they have no counterpart among the global symmetries of the ordinary superstring on a flat background; (ii) they can change the value of the “extra coordinate” $`y`$; (iii) they can map solutions to the classical equations of motion with vanishing field strength $`_mA_n_nA_m`$ to solutions with non-vanishing field strength, in contrast to the translations and Lorentz-transformations. Property (ii) holds because of $`\mathrm{\Delta }y=2r(y)`$ and is intimately related to property (iii) because the field strength is related to $`y`$ by $`f^{}(y)\epsilon ^{mn}_mA_n+\mathrm{}`$ where $``$ is equality on-shell. We stress that the presence of the dilatations is not an artefact of our formulation. Rather, they are of course present even after elimination of $`y`$. Finally we note that properties (ii) and (iii) extend to more complicated backgrounds for which solutions to (46) with $`K^i0`$ exist.
## g Acknowledgements.
FB thanks the Erwin-Schrödinger-Institute for hospitality and financial support during the time when this work was completed, and was supported by a DFG habilitation grant at earlier stages of the work. AK was supported by ÖNB under project grant number 7731 and by “Fonds zur Förderung der Wissenschaftlichen Forschung” under project grant number P13125-TPH.
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# Critical fields on the M5–brane and noncommutative open strings
## I Introduction
Noncommutative geometry has been shown to play a fascinating role in string theory. The original interest was sparked by the appearance of noncommutativity on D-branes in the presence of constant background Neveu–Schwarz two–form potentials $`B_{\mathrm{NS}}`$ . On the D–brane itself these potentials appear as a two–form adapted field strength $`=dA+B_{NS}`$. In the magnetic case (i.e. when $``$ has only spatial parts) a spatially noncommutative Yang–Mills theory (NCYM) on the D–brane can be decoupled from the bulk gravity . The natural, though conceptually challenging generalisation, was to include field strengths with non-zero electric components thus inducing spatio-temporal noncommutativity on the D-brane. This was examined in . The somewhat remarkable result was that by examining the D-brane in a decoupling limit near critical electric field strengthThe concept of a critical field only works for the electric case as it relies crucially on the Lorentz signature. one naturally constructed a unitary decoupled spatio-temporally noncommutative open string theory (NCOS). The thermodynamics of such theories was examined from the dual supergravity point of view in .
The crucial property of the NCOS limit described in is to keep fixed both the open string two–point function and the effective open string coupling constant
$$<X^AX^B>=2\pi \alpha ^{}G^{AB}+\mathrm{\Theta }^{AB}=2\pi \alpha ^{}\left(\frac{1}{g+2\pi \alpha ^{}}\right)^{AB},$$
(1)
$$G_0=g_s\frac{det^{1/2}(g+2\pi \alpha ^{})}{det^{1/2}(g)},$$
(2)
where $`G^{AB}`$ is the symmetric part and $`\mathrm{\Theta }^{AB}`$ the antisymmetric part of the two–point function. In this limit the leading divergent parts of $``$ cancel the contribution from $`g_{\mu \nu }/\alpha ^{}`$, leaving the two–point function and coupling governed by the finite subleading termsHere and in the rest of this paper we will ignore factors of $`2\pi `$ etc.. As discussed in , this NCOS limit introduces a fixed metric $`G_{\mathrm{OS}}^{AB}`$ and a new effective length scale $`\alpha _{\mathrm{eff}}^{}`$ defined as follows:
$$\alpha ^{}G^{AB}=\alpha _{\mathrm{eff}}^{}G_{\mathrm{OS}}^{AB}.$$
(3)
In this paper we will consider a decoupling limit in the M5–brane analogous to the noncommutative open string limit and examine its properties through an open membrane probe. Even though the analog of the string two–point function is not available for the open membrane, it was conjectured in that the decoupled five–brane theory should be formulated in terms of a so called open membrane metric with properties analogous to those of the open string metric. Thus we demand that the six–dimensional proper lengths measured by the open membrane metric are fixed in units of the 11D Planck length $`\mathrm{}_p`$. This defines a finite effective length scale of the decoupled five–brane theory that we shall denote $`\mathrm{}_g`$. This is suggestive of an open membrane theory underlying the decoupled spatio-temporal noncommutative M5–brane.
Evidence for the decoupled non–commutative five–brane theory can be obtained by comparing various reductions of the five–brane limit to NCOS and NCYM limits in string theory. Conversely, this provides a direct interpretation of the strong coupling behavior of the decoupled NCOS and NCYM theories.
The structure of the paper is as follows. We begin in Section II by describing the decoupling limit on the M5–brane. In section III we show that this limit reduces to the NCOS limit on the D4–brane. In Section IV we show that the single limit on the M5–brane reduces to both the NCOS and NCYM limits on the D3–brane. These two limits are related by a modular transformation on the two–torus. Here the role of the open membrane metric is shown to play a crucial role. We end with some conclusions and discussion.
## II The M5–brane and the decoupling limit
The five–brane may be effectively described by a six–dimensional self-dual two form field theory (this is neglecting the superpartners in the (2,0) supermultiplet). The adapted field strength is
$$=db+C,$$
(4)
where $`C`$ is the pull–back to the five–brane of the three–form potential in eleven–dimensional supergravity and $`b`$ is the two form potential on the five–brane worldvolume. The self–duality condition provides a nonlinear algebraic constraint involving the components of the field strength and the induced metric $`g_{\mu \nu }`$ on the brane as follows :
$$\frac{\sqrt{detg}}{6}ϵ_{\mu \nu \rho \sigma \lambda \tau }^{\sigma \lambda \tau }=\frac{1+K}{2}(G^1)_\mu {}_{}{}^{\lambda }_{\nu \rho \lambda }^{},$$
(5)
where $`ϵ^{012345}=1`$ and the scalar $`K`$ and the tensor $`G_{\mu \nu }`$ are given by
$`K`$ $`=`$ $`\sqrt{1+{\displaystyle \frac{\mathrm{}_p^6}{24}}^2},`$ (6)
$`G_{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{1+K}{2K}}\left(g_{\mu \nu }+{\displaystyle \frac{\mathrm{}_p^6}{4}}_{\mu \nu }^2\right),`$ (8)
where $`\mathrm{}_p`$ is the 11D Planck scale.
The relation (5) involves the metric, the field strength components and the plank length $`\mathrm{}_p`$. As we want to carry out a scaling in these quanities we must make sure that any scaling obeys the above relation (5). This is amply discussed in . To facilitate this we introduce a parametrisation of constant flux solutions to (5) as follows:
$$_{\mu \nu \rho }=\frac{h}{\sqrt{1+\mathrm{}_p^6h^2}}ϵ_{\alpha \beta \gamma }v_\mu ^\alpha v_\nu ^\beta v_\rho ^\gamma +hϵ_{abc}u_\mu ^au_\nu ^bu_\rho ^c,$$
(9)
$$G_{\mu \nu }=\frac{\left(1+\sqrt{1+h^2\mathrm{}_p^6}\right)^2}{4}\left(\frac{1}{1+h^2\mathrm{}_p^6}\eta _{\alpha \beta }v_\mu ^\alpha v_\nu ^\beta +\delta _{ab}u_\mu ^au_\nu ^b\right).$$
(10)
Here $`h`$ is a real field of dimension (mass)<sup>3</sup> and $`(v_\mu ^\alpha ,u_\mu ^a)`$, $`\alpha =0,1,2`$, $`a=3,4,5`$, are sechsbein fields in the nine–dimensional coset $`SO(5,1)/SO(2,1)\times SO(3)`$ satisfying
$$g^{\mu \nu }v_\mu ^\alpha v_\nu ^\beta =\eta ^{\alpha \beta },g^{\mu \nu }u_\mu ^av_\nu ^\beta =0,g^{\mu \nu }u_\mu ^au_\nu ^b=\delta ^{ab},$$
(11)
$$g_{\mu \nu }=\eta _{\alpha \beta }v_\mu ^\alpha v_\nu ^\beta +\delta _{ab}u_\mu ^au_\nu ^b.$$
(12)
A derivation of this parametrisation is given in .
The relation between the tensor $`G_{\mu \nu }`$ and the open membrane metric for the five–brane in analogy with the open string metric that occurs on D-branes was discussed in . It should be noted that the overall conformal scale of the open membrane metric was not determined. In this paper such an overall scale will play a role. We therefore define the open membrane metric as follows
$$\widehat{G}_{\mu \nu }\varphi (x)\left(g_{\mu \nu }+\frac{\mathrm{}_p^6}{4}_{\mu \nu }^2\right),$$
(13)
where the function $`\varphi (x)0`$ and $`x`$ is given by the dimensionless combination $`x=\mathrm{}_p^6^2`$. Using the parametrisation (9), this metric can be written as
$$\widehat{G}_{\mu \nu }=(1+\frac{1}{2}h^2\mathrm{}_p^6)\varphi (h^2\mathrm{}_p^6)\left(\frac{1}{1+h^2\mathrm{}_p^6}\eta _{\alpha \beta }v_\mu ^\alpha v_\nu ^\beta +\delta _{ab}u_\mu ^au_\nu ^b\right).$$
(14)
Below we shall determine the asymptotic behavior of $`\varphi (x)`$ for large $`x`$ from the requirements of the decoupling limit.
We now proceed with the definition of the decoupling limit. The properties that we demand for the decoupling limit we wish to take are as follows:
1. The Planck length $`\mathrm{}_p0`$, so that the gravitational interactions can be decoupled.
2. The proper six–dimensional lengths $`ds^2(\widehat{G})`$ of the open membrane metric are fixed in eleven–dimensional Planck units in the limit, i.e. $`\mathrm{}_p^2ds^2(\widehat{G})`$ is fixed, so that the limit describes a genuine six–dimensional theory with a finite length scale $`\mathrm{}_g`$.
3. The electric components contain a divergent piece and a constant piece, in analogy with the limit discussed in for open strings.
The first condition we satisfy by scaling $`\mathrm{}_pϵ^{\frac{1}{3}}`$ ($`ϵ0`$). In order to satisfy the second and third condition we impose that $`h\mathrm{}_p^3`$ diverges<sup>\**</sup><sup>\**</sup>\**This is in contrast with the limit in where $`h\mathrm{}_p^3`$ does not diverge..
We are therefore led to consider the following limit:
$`g_{\alpha \beta }ϵ^0vϵ^0`$ ; $`g_{ab}ϵ^1uϵ^{\frac{1}{2}},`$ (15)
$`\mathrm{}_pϵ^{\frac{1}{3}},h`$ $``$ $`ϵ^{\frac{3}{2}},ϵ0.`$ (16)
such that the components of $``$ given by (9) behave as follows:
$`_{012}`$ $``$ $`\mathrm{}_p^3(1{\displaystyle \frac{1}{2}}\mathrm{}_p^6h^2)ϵ^1+ϵ^0,`$ (17)
$`_{345}`$ $``$ $`hu^3ϵ^0.`$ (18)
The physics on the five–brane in the decoupling limit is uniquely defined by the two fixed noncommutativity parameters of dimension \[length\]<sup>2</sup> constructed from the the finite parts of $`_{012}`$ and $`_{345}`$ as follows:
$$\mathrm{\Theta }_\mathrm{T}(h^2\mathrm{}_p^9)^{2/3},\mathrm{\Theta }_\mathrm{S}(hu^3)^{2/3},$$
(19)
where we have set $`v=1`$. In order to satisfy requirement (ii) we demand in analogy with
$$\mathrm{}_p^2(\widehat{G}^1)^{\mu \nu }\mathrm{}_g^2G_{\mathrm{OM}}^{\mu \nu }\mathrm{is}\mathrm{fixed}.$$
(20)
This allows us to fix the conformal factor as follows:
$$\varphi (x)x^{\frac{2}{3}}\mathrm{as}x\mathrm{}.$$
(21)
With the conformal factor now fixed we find
$$\mathrm{}_p^2(\widehat{G}^1)^{\mu \nu }=(\mathrm{\Theta }_\mathrm{T}\eta ^{\alpha \beta }\mathrm{\Theta }_\mathrm{S}\delta ^{ab})\mathrm{}_g^2G_{\mathrm{OM}}^{\mu \nu }.$$
(22)
This defines a noncommutative M5–brane length scale $`\mathrm{}_g`$, a fixed metric $`G_{\mathrm{OM}}^{\mu \nu }`$ and a dimensionless parameter $`\lambda `$ as follows
$$\mathrm{}_g\sqrt{\mathrm{\Theta }_\mathrm{T}},\lambda \frac{\mathrm{\Theta }_\mathrm{S}}{\mathrm{\Theta }_\mathrm{T}},G_{\mathrm{OM}}^{\mu \nu }=(\eta ^{\alpha \beta }\lambda \delta ^{ab}).$$
(23)
## III The NCOS limit on the D4–brane
In this section we show that the decoupling limit (16) on the M–theory five–brane reduces to the NCOS limit on the D4–brane. This provides an interpretation of the spatio-temporal noncommutative five–brane as the strong coupling dual of the NCOS on the D4–brane.
In order to show this we wrap the five–brane, for finite $`ϵ`$, on a circle of fixed radius $`R`$ in the direction $`x^2`$ and identify
$$x^2=X^{11}X^{11}+R,_{AB}=R_{AB2},A,B=0,1,3,4,5.$$
(24)
Clearly this means that only $`_{01}`$ is nonzero on the D4–brane. We also use the following standard relations between M–theory and IIA string theory parameters:
$$g_s=\left(\frac{R}{\mathrm{}_p}\right)^{\frac{3}{2}},\alpha ^{}=\frac{\mathrm{}_p^3}{R}.$$
(25)
The scaling of the metric components in $`D=11`$ induces the same scaling for the ten–dimensional metric components and together with the requirement of fixed radius $`R`$ we find the following limit on the D4–brane (we reset our conventions such that $`\alpha ,\beta =0,1`$)
$`g_{\alpha \beta }ϵ^0,g_{ab}`$ $``$ $`ϵ^1,_{01}ϵ^1+ϵ^0,`$ (26)
$`\alpha ^{}ϵ^1,g_s`$ $``$ $`ϵ^{\frac{1}{2}},ϵ0.`$ (27)
As a result we find that length scales on the D4–brane, as measured by the open string metric $`G^{AB}`$, are kept fixed in the limit. This also holds for the noncommutativity parameters $`\mathrm{\Theta }^{AB}`$ appearing in the two–point function (1) and the open string coupling $`G_\mathrm{O}`$ given by (2). We identify the NCOS limit on the D4–brane with electric field strength $`_{01}=_c\frac{1}{2}\theta ^1`$, where the diverging critical electric field $`_c`$ and the fixed noncommutativity parameter $`\theta `$ are given by
$$_c=R\mathrm{}_p^3,\theta =\frac{h^2\mathrm{}_p^9}{R}.$$
(28)
Hence, using (18) and (19), we can write fixed D4–brane quantities in terms of the fixed five–brane data $`\mathrm{\Theta }_\mathrm{T}`$, $`\mathrm{\Theta }_\mathrm{S}`$ and $`R`$, or equivalently $`\mathrm{}_g`$, $`\lambda `$ and $`R`$:
$$\alpha ^{}G^{AB}=(\frac{\mathrm{\Theta }_\mathrm{T}^{\frac{3}{2}}}{R}\eta ^{\alpha \beta }\frac{\mathrm{\Theta }_\mathrm{T}^{\frac{1}{2}}\mathrm{\Theta }_\mathrm{S}}{R}\delta ^{ab})=\frac{\mathrm{}_g^3}{R}(\eta ^{\alpha \beta }\lambda \delta ^{ab})=\frac{\mathrm{}_g^3}{R}G_{\mathrm{OM}}^{AB},$$
(29)
$$\theta ^{AB}=(\frac{\mathrm{\Theta }_T^{\frac{3}{2}}}{R}ϵ^{\alpha \beta }0)=\frac{\mathrm{}_g^3}{R}(ϵ^{\alpha \beta }0),$$
(30)
$$G_\mathrm{O}=\left(\frac{R}{\mathrm{}_g}\right)^{\frac{3}{2}}.$$
(31)
Therefore the NCOS limit on the D4–brane has an effective open string scale $`\alpha _{\mathrm{eff}}^{}`$ and noncommutativity parameter $`\theta `$ given by
$$\alpha _{\mathrm{eff}}^{}\theta =\frac{\mathrm{}_g^3}{R},G_{\mathrm{OM}}^{AB}=G_{\mathrm{OS}}^{AB}.$$
(32)
Thus we find the following relations between open string moduli and M–theory open membrane moduli:
$`R`$ $`=`$ $`G_O\sqrt{\alpha _{\mathrm{eff}}^{}}`$ (33)
$`l_g`$ $`=`$ $`G_O^{\frac{1}{3}}\sqrt{\alpha _{\mathrm{eff}}^{}}.`$ (34)
These are formally equivalent to the standard relations between the moduli of M-theory and IIA superstring theory provided that we give $`\mathrm{}_g`$ a six–dimensional interpretation analogous to that of the eleven–dimensional Planck scale $`\mathrm{}_p`$ in M–theory. This suggests that the NCOS theory on the D4–brane generates an extra dimension when we increase the open string coupling and in the limit $`R\mathrm{}`$ we end up with a noncommutative (in all directions!) six–dimensional theory governed by the scale $`\mathrm{}_g`$, as displayed in Figure 1. Note that $`\alpha _{\mathrm{eff}}^{}=\theta `$ implies that a field theory limit taking $`\alpha _{\mathrm{eff}}^{}0`$ will at the same time also result in vanishing spatio-temporal noncommutativity.
## IV The noncommutative limits on the D3–brane
Next we carry out the double–dimensional reduction of the M–theory five–brane on a fixed two–torus in order to compare with the limits given in for the D3–brane directly reduced on a circle. We drop all Kaluza–Klein modes and identify the wrapped five–brane with the directly reduced D3–brane in nine dimensions. This gives the following relations between M–theory five brane and IIB three brane quantities
$$(x^2,x^5)=(X^{11},X^9)(X^{11}+R_2,X^9+R_5),$$
(35)
$`_{AB}=R_2_{AB2},{\displaystyle \frac{g_{AB}^{(E)}}{\alpha ^{}}}={\displaystyle \frac{\sqrt{R_2R_5u}}{\mathrm{}_p^3}}g_{AB},A,B=0,1,3,4,`$ (36)
$$g_s=\frac{R_2}{R_5u},\frac{R_B^{(S)}}{\sqrt{\alpha ^{}}}=\frac{\mathrm{}_p^{\frac{3}{2}}}{R_2^{\frac{1}{2}}R_5u}.$$
(37)
where the two–torus coordinate periodicities $`R_2`$ and $`R_5`$ are fixed quantities (we set the real part of the complex structure equal to zero) and where $`u`$ is the scale of the induced dreibein on the five–brane in the $`3,4,5`$ directions. The quantity $`R_B^{(S)}`$ is the IIB radius in string frame and $`g_{MN}^{(E)}`$ is the IIB Einstein metric which is related to the IIB string metric by $`g_{MN}^{(E)}=g_s^{\frac{1}{2}}g_{MN}^{(S)}`$.
The S-dual description of the D3–brane, giving a magnetic field strength, can be obtained by performing a modular transformation on the two–torus, which gives the following relations between the quantities in M–theory and the S-dual picture:
$`\stackrel{~}{}_{AB}=R_5_{AB5},\stackrel{~}{g}_s={\displaystyle \frac{R_5u}{R_2}},{\displaystyle \frac{\stackrel{~}{g}_{AB}^{(E)}}{\stackrel{~}{\alpha }^{}}}={\displaystyle \frac{g_{AB}^{(E)}}{\alpha ^{}}},{\displaystyle \frac{\stackrel{~}{R}_B^{(S)}}{\sqrt{\stackrel{~}{\alpha }^{}}}}=g_s^{\frac{1}{2}}{\displaystyle \frac{R_B^{(S)}}{\sqrt{\alpha ^{}}}},`$ (38)
where the tilde denotes quantities in the the IIB S-dual picture.
Inserting the scaling limit for the M5–brane (16) in the first set of relations above we obtain the following scaling limit for the D3–brane:
$`_{01}`$ $``$ $`ϵ^1+ϵ^0,_{34}=0,`$ (39)
$`{\displaystyle \frac{g_{AB}^{(S)}}{\alpha ^{}}}`$ $``$ $`\mathrm{diag}(ϵ^1,ϵ^1,ϵ^0,ϵ^0),g_sϵ^{\frac{1}{2}}.`$ (40)
We identify (40) as the open string limit where the critical field strength and the noncommutativity parameter $`\theta _{\mathrm{NCOS}}`$ are given in terms of M–theory quantities by
$$_c=R_2\mathrm{}_p^3,\theta _{\mathrm{NCOS}}=\frac{h^2\mathrm{}_p^9}{R_2}.$$
(41)
Identifying the finite quantities on the D3–brane with the finite quantities on the wrapped M five–brane we obtain the following relations (we reset our conventions such that $`\alpha =0,1`$ and $`a=3,4`$):
$$\alpha ^{}G^{AB}=(\frac{\mathrm{\Theta }_\mathrm{T}^{\frac{3}{2}}}{R_2}\eta ^{\alpha \beta }\frac{\mathrm{\Theta }_\mathrm{T}^{\frac{1}{2}}\mathrm{\Theta }_\mathrm{S}}{R_2}\delta ^{ab})=\frac{\mathrm{}_g^3}{R_2}\left(\eta ^{\alpha \beta }\lambda \delta ^{ab}\right),$$
(42)
$$\theta ^{AB}=(\frac{\mathrm{\Theta }_\mathrm{T}^{\frac{3}{2}}}{R_2}ϵ^{\alpha \beta }0)=\frac{\mathrm{}_g^3}{R_2}(ϵ^{\alpha \beta }0),$$
(43)
$$G_\mathrm{O}=\frac{R_2}{R_5}\sqrt{\frac{\mathrm{\Theta }_\mathrm{S}}{\mathrm{\Theta }_\mathrm{T}}}=\frac{R_2}{R_5}\sqrt{\lambda },$$
(44)
$$r_\mathrm{B}\frac{R_\mathrm{B}^{(S)}}{\sqrt{\alpha ^{}}}=\frac{\mathrm{\Theta }_\mathrm{S}^{\frac{1}{2}}\mathrm{\Theta }_\mathrm{T}^{\frac{1}{4}}}{R_2^{\frac{1}{2}}R_5}.$$
(45)
Hence the NCOS limit on the D3–brane has effective open string scale $`\alpha _{\mathrm{eff}}^{}`$ and noncommutativity parameter $`\theta _{\mathrm{NCOS}}`$ both given by
$$\alpha _{\mathrm{eff}}^{}\theta _{\mathrm{NCOS}}=\frac{\mathrm{\Theta }_\mathrm{T}^{\frac{3}{2}}}{R_2},$$
(46)
Importantly, the worldsheet sigma model coupling constant $`r_\mathrm{B}`$ given by (45) remains finite in the limit.
In the S-dual/modular transformed description we find (inserting the five brane limit into the second set of relations):
$`\stackrel{~}{}_{01}`$ $`=`$ $`0,\stackrel{~}{}_{34}=ϵ^0`$ (47)
$`{\displaystyle \frac{\stackrel{~}{g}_{AB}}{\stackrel{~}{\alpha }^{}}}`$ $``$ $`\mathrm{diag}(ϵ^{\frac{1}{2}},ϵ^{\frac{1}{2}},ϵ^{\frac{1}{2}},ϵ^{\frac{1}{2}}),\stackrel{~}{g}_sϵ^{\frac{1}{2}}`$ (48)
This we identify with the noncommutative field theory limit as described in . In this case the mass scales on the D3–brane are sent to zero, i.e. $`\stackrel{~}{\alpha }^{}G^{AB}0`$, decoupling the massive string modes.
The relations between the fixed M–theory quantities and the fixed quantities on the D3–brane are as follows:
$`g_{YM}^2`$ $`=`$ $`{\displaystyle \frac{R_5}{R_2\sqrt{\lambda }}},\theta _{\mathrm{NCYM}}={\displaystyle \frac{\mathrm{\Theta }_\mathrm{S}^{\frac{3}{2}}}{R_5}},`$ (49)
$`m`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Theta }_\mathrm{S}^{\frac{1}{4}}}{\sqrt{R_5}R_2}},`$ (50)
where $`m`$ is the periodicity in mass units of the compact scalar in the four–dimensional noncommutative action. Note that for the noncommutative field theory, instead of a fixed worldsheet sigma model radius $`r_\mathrm{B}`$ we now find a fixed kinetic term for the compact transverse scalar $`\mathrm{\Phi }\frac{X^9}{R_B^{(S)}}`$ in the limit when $`ϵ0`$.
The natural fixed moduli for the noncommutative M-five brane are $`l_g`$, the complex structure of the torus, $`\tau _{OM}`$ and the area of the torus, $`A_{OM}`$ as measured by the open membrane metric:
$$\tau _{OM}=\frac{R_2}{R_5}\sqrt{\lambda },A_{OM}=R_2R_5\frac{1}{\sqrt{\lambda }}$$
(51)
We now recover the standard relations between M theory and the S-dual descriptions of IIB for the noncommutative open string/membrane moduli. For the noncommutative open string,
$`G_O`$ $`=`$ $`\tau _{OM}`$ (52)
$`r_B`$ $`=`$ $`A_{OM}^{\frac{3}{4}}\tau _{OM}^{\frac{1}{4}}l_g^{\frac{3}{2}}.`$ (53)
For the S-dual, noncommutative field theory,
$`g_{YM}^2`$ $`=`$ $`{\displaystyle \frac{1}{\tau _{OM}}}`$ (54)
$`m`$ $`=`$ $`A_{OM}^{\frac{3}{4}}\tau _{OM}^{\frac{1}{4}}l_g^{\frac{1}{2}}.`$ (55)
The duality transformation is now obtained by a modular transformation on the torus as seen by the open membrane metric so that the two theories are related by:
$$\tau _{OM}\frac{1}{\tau _{OM}}$$
(56)
with the appropriate interpretation of duality related quantities. We remark that the couplings are independent of our choice of conformal factor for the open membrane metric.
The duality between NCOS and NCYM is possible for the D3–brane because both open string coupling and Yang-Mills coupling are independent of the closed string scale $`\alpha ^{}`$.
Finally, we consider the following limits of the NCOS on the D3–brane (we set $`\lambda =1`$ below):
1. $`T`$: Taking $`r_\mathrm{B}0`$ while keeping $`G_{\mathrm{O},\mathrm{A}}r_\mathrm{B}^1G_\mathrm{O}`$ fixed leads to the NCOS on the T-dual D4–brane with open string coupling $`G_{\mathrm{O},\mathrm{A}}`$. This is analogous to how the usual closed string coupling transforms under T-duality. It is interesting that this noncommutative open string theory exhibits a sort of T-duality.
2. $`S`$: Taking $`G_\mathrm{O}\mathrm{}`$ while keeping $`\theta _{\mathrm{NCYM}}G_\mathrm{O}\sqrt{\alpha _{\mathrm{eff}}^{}}`$ fixed leads to the S-dual NCYM on the D3–brane with $`g_{YM}0`$, as expected.
## V Discussion
We have argued for the existence of a decoupled noncommutative theory on the five–brane defined by the limit (16) by showing its relation to various well–defined limits of IIA and IIB string theory. Ultimately we are of course interested in finding an intrinsically six–dimensional definition of the decoupled theory. One may wonder to what extent the open membrane action underlies such a formulation and in particular whether there is an analog of the subtle cancellations between the diverging electric field and tension that occur in the string case. In the critical limit we expect the finite parts of the Wess–Zumino term to yield the non–commutative structure of the five–brane loop–space via the definition of the functional Moyal product given in . The role of the kinetic part of the action is more unclear, however, due to the usual membrane instability. Interestingly one may construct a stable, non–degenerate open membrane solution in the critical limit which is not a limit of any solution for finite $`ϵ`$. This is the analogue of the the critical string solution discussed in . One would hope that the quantisation of the membrane in this near critical background will provide the open membrane metric with the appropriate conformal factor given in this paper. This is ongoing work. It is not yet clear whether the critical field will cure the usual membrane sicknesses.
Finally we wish to make the following observation, given our choice of conformal factor for the open membrane metric we see that the line element for the self–dual string solution is exactly $`AdS_3\times S^3`$ in the near horizon limit. Given that the proceedure in this paper has been to reproduce the usual bulk relations for the decoupled theories on the brane one wonders whether it might be possible to formulate an AdS/CFT corresspondence for the self–dual string in the five brane.
In summary, the NCOS on the D4–brane with noncommutativity parameter $`\theta =\alpha _{\mathrm{eff}}^{}`$ has a dual description in the limit of strong coupling $`G_O>>1`$ as a noncommutative five–brane with fundamental length $`\mathrm{}_g=G_O^{\frac{1}{3}}\sqrt{\alpha _{\mathrm{eff}}^{}}`$ reduced on a circle of radius $`R=G_O\sqrt{\alpha _{\mathrm{eff}}^{}}`$. The S–duality of the NCOS and NCYM theories on a directly reduced D3–brane follows from the modular invariance of the noncommutative five–brane wrapped on a torus. The couplings on the D3–brane are identified with the complex structure of the torus in the open membrane metric.
Note Added
During the completion of this paper we received the preprint that also discusses the noncommutative open membrane limit and its relation to the noncommutative open string limit on the D4–brane. The preprint also contains an interesting discussion of NCOS theories at strong coupling for the other D–branes. Our paper emphasizes the role played by the open membrane metric and the relation between M/IIB moduli.
Acknowledgements
P.S. is grateful to M. Cederwall, H. Larsson and B.E.W. Nilsson for discussions. D.S.B. thanks R. Gopakumar for discussions during the Fradkin Memorial conference on Quantization, Gauge Theory and Strings. This work is supported by the European Commission TMR program ERBFMRX-CT96-0045, in which E.B., D.S.B. and J.P.v.d.S. are associated to the University of Utrecht. The work of J.P.v.d.S. and P.S. is part of the research program of the “Stichting voor Fundamenteel Onderzoek der Materie” (FOM).
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# Statistical Mechanics of Recurrent Neural Networks II. Dynamics
## 1 Introduction
This paper, on solving the dynamics of recurrent neural networks using non-equilibrium statistical mechanical techniques, is the sequel of , which was devoted to solving the statics using equilibrium techniques. I refer to for a general introduction to recurrent neural networks and their properties.
Equilibrium statistical mechanical techniques can provide much detailed quantitative information on the behaviour of recurrent neural networks, but they obviously have serious restrictions. The first one is that, by definition, they will only provide information on network properties in the stationary state. For associative memories, for instance, it is not clear how one can calculate quantities like sizes of domains of attraction without solving the dynamics. The second, and more serious, restriction is that for equilibrium statistical mechanics to apply the dynamics of the network under study must obey detailed balance, i.e. absence of microscopic probability currents in the stationary state. As we have seen in , for recurrent networks in which the dynamics take the form of a stochastic alignment of neuronal firing rates to post-synaptic potentials which, in turn, depend linearly on the firing rates, this requirement of detailed balance usually implies symmetry of the synaptic matrix. From a physiological point of view this requirement is clearly unacceptable, since it is violated in any network that obeys Dale’s law as soon as an excitatory neuron is connected to an inhibitory one. Worse still, we saw in that in any network of graded-response neurons detailed balance will always be violated, even when the synapses are symmetric. The situation will become even worse when we turn to networks of yet more realistic (spike-based) neurons, such as integrate-and-fire ones. In contrast to this, non-equilibrium statistical mechanical techniques, it will turn out, do not impose such biologically non-realistic restrictions on neuron types and synaptic symmetry, and they are consequently the more appropriate avenue for future theoretical research aimed at solving biologically more realistic models.
The common strategy of all non-equilibrium statistical mechanical studies is to derive and solve dynamical laws for a suitable small set of relevant macroscopic quantities from the dynamical laws of the underlying microscopic neuronal system. In order to make progress, as in equilibrium studies, one is initially forced to pay the price of having relatively simple model neurons, and of not having a very complicated spatial wiring structure in the network under study; the networks described and analysed in this paper will consequently be either fully connected, or randomly diluted. When attempting to obtain exact dynamical solutions within this class, one then soon finds a clear separation of network models into two distinct complexity classes, reflecting in the dynamics a separation which we also found in the statics. In statics one could get away with relatively simple mathematical techniques as long as the number of attractors of the dynamics was small compared to the number $`N`$ of neurons. As soon as the number of attractors became of the order of $`N`$, on the other hand, one entered the complex regime, requiring the more complicated formalism of replica theory. In dynamics we will again find that we can get away with relatively simple mathematical techniques as long as the number of attractors remains small, and find closed deterministic differential equations for macroscopic quantities with just a single time argument. As soon as we enter the complex regime, however, we will no longer find closed equations for one-time macroscopic objects: we will now have to work with correlation and response functions, which have two time arguments, and turn to the less trivial generating functional techniques<sup>1</sup><sup>1</sup>1A brief note about terminology: strictly speaking, in this paper we will apply these techniques only to models in which time is measured in discrete units, so that we should speak about generating functions rather than generating functionals. However, since these techniques can and have also been applied intensively to models with continuous time, they are in literature often referred to as generating functional techniques, for both discrete and continuous time..
In contrast to the situation in statics , I cannot in this paper give many references to textbooks on the dynamics, since these are more or less non-existent. There would appear to be two reasons for this. Firstly, in most physics departments non-equilibrium statistical mechanics (as a subject) is generally taught and applied far less intensively than equilibrium statistical mechanics, and thus the non-equilibrium studies of recurrent neural networks have been considerably less in number and later in appearance in literature than their equilibrium counterparts. Secondly, many of the popular textbooks on the statistical mechanics of neural networks were written around 1989, roughly at the point in time where non-equilibrium statistical mechanical studies just started being taken up. When reading such textbooks one could be forgiven for thinking that solving the dynamics of recurrent neural networks is generally ruled out, whereas, in fact, nothing could be further from the truth. Thus the references in this paper will, out of necessity, be mainly to research papers. I regret that, given constraints on page numbers and given my aim to explain ideas and techniques in a lecture notes style (rather than display encyclopedic skills), I will inevitably have left out relevant references. Another consequence of the scarce and scattered nature of the literature on the non-equilibrium statistical mechanics of recurrent neural networks is that a situation has developed where many mathematical procedures, properties and solutions are more or less known by the research community, but without there being a clear reference in literature where these were first formally derived (if at all). Examples of this are the fluctuation-dissipation theorems for parallel dynamics and the non-equilibrium analysis of networks with graded response neurons; often the separating boundary between accepted general knowledge and published accepted general knowledge is somewhat fuzzy.
The structure of this paper mirrors more or less the structure of . Again I will start with relatively simple networks, with a small number of attractors (such as systems with uniform synapses, or with a small number of patterns stored with Hebbian-type rules), which can be solved with relatively simple mathematical techniques. These will now also include networks that do not evolve to a stationary state, and networks of graded response neurons, which could not be studied within equilibrium statistical mechanics at all. Next follows a detour on correlation- and response functions and their relations (i.e. fluctuation-dissipation theorems), which serves as a prerequisite for the last section on generating functional methods, which are indeed formulated in the language of correlation- and response functions. In this last, more mathematically involved, section I study symmetric and non-symmetric attractor neural networks close to saturation, i.e. in the complex regime. I will show how to solve the dynamics of fully connected as well as extremely diluted networks, emphasising the (again) crucial issue of presence (or absence) of synaptic symmetry, and compare the predictions of the (exact) generating functional formalism to both numerical simulations and simple approximate theories.
## 2 Attractor Neural Networks with Binary Neurons
The simplest non-trivial recurrent neural networks consist of $`N`$ binary neurons $`\sigma _i\{1,1\}`$ (see ) which respond stochastically to post-synaptic potentials (or local fields) $`h_i(𝝈)`$, with $`𝝈=(\sigma _1,\mathrm{},\sigma _N)`$. The fields depend linearly on the instantaneous neuron states, $`h_i(𝝈)=_jJ_{ij}\sigma _j+\theta _i`$, with the $`J_{ij}`$ representing synaptic efficacies, and the $`\theta _i`$ representing external stimuli and/or neural thresholds.
### 2.1 Closed Macroscopic Laws for Sequential Dynamics
First I show how for sequential dynamics (where neurons are updated one after the other) one can calculate, from the microscopic stochastic laws, differential equations for the probability distribution of suitably defined macroscopic observables. For mathematical convenience our starting point will be the continuous-time master equation for the microscopic probability distribution $`p_t(𝝈)`$
$$\frac{d}{dt}p_t(𝝈)=\underset{i}{}\left\{w_i(F_i𝝈)p_t(F_i𝝈)w_i(𝝈)p_t(𝝈)\right\}w_i(𝝈)=\frac{1}{2}[1\sigma _i\mathrm{tanh}[\beta h_i(𝝈)]]$$
(1)
with $`F_i\mathrm{\Phi }(𝝈)=\mathrm{\Phi }(\sigma _1,\mathrm{},\sigma _{i1},\sigma _i,\sigma _{i+1},\mathrm{},\sigma _N)`$ (see ). I will discuss the conditions for the evolution of these macroscopic state variables to become deterministic in the limit of infinitely large networks and, in addition, be governed by a closed set of equations. I then turn to specific models, with and without detailed balance, and show how the macroscopic equations can be used to illuminate and understand the dynamics of attractor neural networks away from saturation.
A Toy Model. Let me illustrate the basic ideas with the help of a simple (infinite range) toy model: $`J_{ij}=(J/N)\eta _i\xi _j`$ and $`\theta _i=0`$ (the variables $`\eta _i`$ and $`\xi _i`$ are arbitrary, but may not depend on $`N`$). For $`\eta _i=\xi _i=1`$ we get a network with uniform synapses. For $`\eta _i=\xi _i\{1,1\}`$ and $`J>0`$ we recover the Hopfield model with one stored pattern. Note: the synaptic matrix is non-symmetric as soon as a pair $`(ij)`$ exists such that $`\eta _i\xi _j\eta _j\xi _i`$, so in general equilibrium statistical mechanics will not apply. The local fields become $`h_i(𝝈)=J\eta _im(𝝈)`$ with $`m(𝝈)=\frac{1}{N}_k\xi _k\sigma _k`$. Since they depend on the microscopic state $`𝝈`$ only through the value of $`m`$, the latter quantity appears to constitute a natural macroscopic level of description. The probability density of finding the macroscopic state $`m(𝝈)=m`$ is given by $`𝒫_t[m]=_𝝈p_t(𝝈)\delta [mm(𝝈)]`$. Its time derivative follows upon inserting (1):
$$\frac{d}{dt}𝒫_t[m]=\underset{𝝈}{}\underset{k=1}{\overset{N}{}}p_t(𝝈)w_k(𝝈)\left\{\delta [mm(𝝈)+\frac{2}{N}\xi _k\sigma _k]\delta \left[mm(𝝈)\right]\right\}$$
$$=\frac{d}{dm}\left\{\underset{𝝈}{}p_t(𝝈)\delta \left[mm(𝝈)\right]\frac{2}{N}\underset{k=1}{\overset{N}{}}\xi _k\sigma _kw_k(𝝈)\right\}+𝒪(\frac{1}{N})$$
Inserting our expressions for the transition rates $`w_i(𝝈)`$ and the local fields $`h_i(𝝈)`$ gives:
$$\frac{d}{dt}𝒫_t[m]=\frac{d}{dm}\left\{𝒫_t[m]\left[m\frac{1}{N}\underset{k=1}{\overset{N}{}}\xi _k\mathrm{tanh}[\eta _k\beta Jm]\right]\right\}+𝒪(N^1)$$
In the limit $`N\mathrm{}`$ only the first term survives. The general solution of the resulting Liouville equation is $`𝒫_t[m]=𝑑m_0𝒫_0[m_0]\delta \left[mm(t|m_0)\right]`$, where $`m(t|m_0)`$ is the solution of
$$\frac{d}{dt}m=\underset{N\mathrm{}}{lim}\frac{1}{N}\underset{k=1}{\overset{N}{}}\xi _k\mathrm{tanh}[\eta _k\beta Jm]mm(0)=m_0$$
(2)
This describes deterministic evolution; the only uncertainty in the value of $`m`$ is due to uncertainty in initial conditions. If at $`t=0`$ the quantity $`m`$ is known exactly, this will remain the case for finite time-scales; $`m`$ turns out to evolve in time according to (2).
Arbitrary Synapses. Let us now allow for less trivial choices of the synaptic matrix $`\{J_{ij}\}`$ and try to calculate the evolution in time of a given set of macroscopic observables $`𝛀(𝝈)=(\mathrm{\Omega }_1(𝝈),\mathrm{},\mathrm{\Omega }_n(𝝈))`$ in the limit $`N\mathrm{}`$. There are no restrictions yet on the form or the number $`n`$ of these state variables; these will, however, arise naturally if we require the observables $`𝛀`$ to obey a closed set of deterministic laws, as we will see. The probability density of finding the system in macroscopic state $`𝛀`$ is given by:
$$𝒫_t\left[𝛀\right]=\underset{𝝈}{}p_t(𝝈)\delta \left[𝛀𝛀(𝝈)\right]$$
(3)
Its time derivative is obtained by inserting (1). If in those parts of the resulting expression which contain the operators $`F_i`$ we perform the transformations $`𝝈F_i𝝈`$, we arrive at
$$\frac{d}{dt}𝒫_t\left[𝛀\right]=\underset{i}{}\underset{𝝈}{}p_t(𝝈)w_i(𝝈)\left\{\delta \left[𝛀𝛀(F_i𝝈)\right]\delta \left[𝛀𝛀(𝝈)\right]\right\}$$
Upon writing $`\mathrm{\Omega }_\mu (F_i𝝈)=\mathrm{\Omega }_\mu (𝝈)+\mathrm{\Delta }_{i\mu }(𝝈)`$ and making a Taylor expansion in powers of $`\{\mathrm{\Delta }_{i\mu }(𝝈)\}`$, we finally obtain the so-called Kramers-Moyal expansion:
$$\frac{d}{dt}𝒫_t\left[𝛀\right]=\underset{\mathrm{}1}{}\frac{(1)^{\mathrm{}}}{\mathrm{}!}\underset{\mu _1=1}{\overset{n}{}}\mathrm{}\underset{\mu _{\mathrm{}}=1}{\overset{n}{}}\frac{^{\mathrm{}}}{\mathrm{\Omega }_{\mu _1}\mathrm{}\mathrm{\Omega }_\mu _{\mathrm{}}}\left\{𝒫_t\left[𝛀\right]F_{\mu _1\mathrm{}\mu _{\mathrm{}}}^{(\mathrm{})}[𝛀;t]\right\}$$
(4)
It involves conditional averages $`f(𝝈)_{𝛀;t}`$ and the ‘discrete derivatives’ $`\mathrm{\Delta }_{j\mu }(𝝈)=\mathrm{\Omega }_\mu (F_j𝝈)\mathrm{\Omega }_\mu (𝝈)`$ <sup>2</sup><sup>2</sup>2 Expansion (4) is to be interpreted in a distributional sense, i.e. only to be used in expressions of the form $`𝑑𝛀𝒫_t(𝛀)G(𝛀)`$ with smooth functions $`G(𝛀)`$, so that all derivatives are well-defined and finite. Furthermore, (4) will only be useful if the $`\mathrm{\Delta }_{j\mu }`$, which measure the sensitivity of the macroscopic quantities to single neuron state changes, are sufficiently small. This is to be expected: for finite $`N`$ any observable can only assume a finite number of possible values; only for $`N\mathrm{}`$ may we expect smooth probability distributions for our macroscopic quantities.:
$$F_{\mu _1\mathrm{}\mu _l}^{(l)}[𝛀;t]=\underset{j=1}{\overset{N}{}}w_j(𝝈)\mathrm{\Delta }_{j\mu _1}(𝝈)\mathrm{}\mathrm{\Delta }_{j\mu _{\mathrm{}}}(𝝈)_{𝛀;t}f(𝝈)_{𝛀;t}=\frac{_𝝈p_t(𝝈)\delta \left[𝛀𝛀(𝝈)\right]f(𝝈)}{_𝝈p_t(𝝈)\delta \left[𝛀𝛀(𝝈)\right]}$$
(5)
Retaining only the $`\mathrm{}=1`$ term in (4) would lead us to a Liouville equation, which describes deterministic flow in $`𝛀`$ space. Including also the $`\mathrm{}=2`$ term leads us to a Fokker-Planck equation which, in addition to flow, describes diffusion of the macroscopic probability density. Thus a sufficient condition for the observables $`𝛀(𝝈)`$ to evolve in time deterministically in the limit $`N\mathrm{}`$ is:
$$\underset{N\mathrm{}}{lim}\underset{\mathrm{}2}{}\frac{1}{\mathrm{}!}\underset{\mu _1=1}{\overset{n}{}}\mathrm{}\underset{\mu _{\mathrm{}}=1}{\overset{n}{}}\underset{j=1}{\overset{N}{}}|\mathrm{\Delta }_{j\mu _1}(𝝈)\mathrm{}\mathrm{\Delta }_{j\mu _{\mathrm{}}}(𝝈)|_{𝛀;t}=0$$
(6)
In the simple case where all observables $`\mathrm{\Omega }_\mu `$ scale similarly in the sense that all ‘derivatives’ $`\mathrm{\Delta }_{j\mu }=\mathrm{\Omega }_\mu (F_i𝝈)\mathrm{\Omega }_\mu (𝝈)`$ are of the same order in $`N`$ (i.e. there is a monotonic function $`\stackrel{~}{\mathrm{\Delta }}_N`$ such that $`\mathrm{\Delta }_{j\mu }=𝒪(\stackrel{~}{\mathrm{\Delta }}_N)`$ for all $`j\mu `$), for instance, criterion (6) becomes:
$$\underset{N\mathrm{}}{lim}n\stackrel{~}{\mathrm{\Delta }}_N\sqrt{N}=0$$
(7)
If for a given set of observables condition (6) is satisfied, we can for large $`N`$ describe the evolution of the macroscopic probability density by a Liouville equation:
$$\frac{d}{dt}𝒫_t\left[𝛀\right]=\underset{\mu =1}{\overset{n}{}}\frac{}{\mathrm{\Omega }_\mu }\left\{𝒫_t\left[𝛀\right]F_\mu ^{(1)}[𝛀;t]\right\}$$
whose solution describes deterministic flow: $`𝒫_t[𝛀]=𝑑𝛀_0𝒫_0[𝛀_0]\delta [𝛀𝛀(t|𝛀_0)]`$ with $`𝛀(t|𝛀_0)`$ given, in turn, as the solution of
$$\frac{d}{dt}𝛀(t)=𝑭^{(1)}[𝛀(t);t]𝛀(0)=𝛀_0$$
(8)
In taking the limit $`N\mathrm{}`$, however, we have to keep in mind that the resulting deterministic theory is obtained by taking this limit for finite $`t`$. According to (4) the $`\mathrm{}>1`$ terms do come into play for sufficiently large times $`t`$; for $`N\mathrm{}`$, however, these times diverge by virtue of (6).
The Issue of Closure. Equation (8) will in general not be autonomous; tracing back the origin of the explicit time dependence in the right-hand side of (8) one finds that to calculate $`𝑭^{(1)}`$ one needs to know the microscopic probability density $`p_t(𝝈)`$. This, in turn, requires solving equation (1) (which is exactly what one tries to avoid). We will now discuss a mechanism via which to eliminate the offending explicit time dependence, and to turn the observables $`𝛀(𝝈)`$ into an autonomous level of description, governed by closed dynamic laws. The idea is to choose the observables $`𝛀(𝝈)`$ in such a way that there is no explicit time dependence in the flow field $`𝑭^{(1)}[𝛀;t]`$ (if possible). According to (5) this implies making sure that there exist functions $`\mathrm{\Phi }_\mu \left[𝛀\right]`$ such that
$$\underset{N\mathrm{}}{lim}\underset{j=1}{\overset{N}{}}w_j(𝝈)\mathrm{\Delta }_{j\mu }(𝝈)=\mathrm{\Phi }_\mu \left[𝛀(𝝈)\right]$$
(9)
in which case the time dependence of $`𝑭^{(1)}`$ indeed drops out and the macroscopic state vector simply evolves in time according to:
$$\frac{d}{dt}𝛀=𝚽\left[𝛀\right],𝚽[𝛀]=(\mathrm{\Phi }_1[𝛀],\mathrm{},\mathrm{\Phi }_n[𝛀])$$
Clearly, for this closure method to apply, a suitable separable structure of the synaptic matrix is required. If, for instance, the macroscopic observables $`\mathrm{\Omega }_\mu `$ depend linearly on the microscopic state variables $`𝝈`$ (i.e. $`\mathrm{\Omega }_\mu (𝝈)=\frac{1}{N}_{j=1}^N\omega _{\mu j}\sigma _j`$), we obtain with the transition rates defined in (1):
$$\frac{d}{dt}\mathrm{\Omega }_\mu =\underset{N\mathrm{}}{lim}\frac{1}{N}\underset{j=1}{\overset{N}{}}\omega _{\mu j}\mathrm{tanh}(\beta h_j(𝝈))\mathrm{\Omega }_\mu $$
(10)
in which case the only further condition for (9) to hold is that all local fields $`h_k(𝝈)`$ must (in leading order in $`N`$) depend on the microscopic state $`𝝈`$ only through the values of the observables $`𝛀`$; since the local fields depend linearly on $`𝝈`$ this, in turn, implies that the synaptic matrix must be separable: if $`J_{ij}=_\mu K_{i\mu }\omega _{\mu j}`$ then indeed $`h_i(𝝈)=_\mu K_{i\mu }\mathrm{\Omega }_\mu (𝝈)+\theta _i`$. Next I will show how this approach can be applied to networks for which the matrix of synapses has a separable form (which includes most symmetric and non-symmetric Hebbian type attractor models). I will restrict myself to models with $`\theta _i=0`$; introducing non-zero thresholds is straightforward and does not pose new problems.
### 2.2 Application to Separable Attractor Networks
Separable models: Description at the Level of Sublattice Activities. We consider the following class of models, in which the interaction matrices have the form
$$J_{ij}=\frac{1}{N}Q(𝝃_i;𝝃_j)𝝃_i=(\xi _i^1,\mathrm{},\xi _i^p)$$
(11)
The components $`\xi _i^\mu `$, representing the information (’patterns’) to be stored or processed, are assumed to be drawn from a finite discrete set $`\mathrm{\Lambda }`$, containing $`n_\mathrm{\Lambda }`$ elements (they are not allowed to depend on $`N`$). The Hopfield model corresponds to choosing $`Q(𝒙;𝒚)=𝒙𝒚`$ and $`\mathrm{\Lambda }\{1,1\}`$. One now introduces a partition of the system $`\{1,\mathrm{},N\}`$ into $`n_\mathrm{\Lambda }^p`$ so-called sublattices $`I_𝜼`$:
$$I_𝜼=\{i|𝝃_i=𝜼\}\{1,\mathrm{},N\}=\underset{𝜼}{}I_𝜼𝜼\mathrm{\Lambda }^p$$
(12)
The number of neurons in sublattice $`I_𝜼`$ is denoted by $`|I_𝜼|`$ (this number will have to be large). If we choose as our macroscopic observables the average activities (‘magnetisations’) within these sublattices, we are able to express the local fields $`h_k`$ solely in terms of macroscopic quantities:
$$m_𝜼(𝝈)=\frac{1}{|I_𝜼|}\underset{iI_𝜼}{}\sigma _i,h_k(𝝈)=\underset{𝜼}{}p_𝜼Q(𝝃_k;𝜼)m_𝜼$$
(13)
with the relative sublattice sizes $`p_𝜼=|I_𝜼|/N`$. If all $`p_𝜼`$ are of the same order in $`N`$ (which, for example, is the case if the vectors $`𝝃_i`$ have been drawn at random from the set $`\mathrm{\Lambda }^p`$) we may write $`\mathrm{\Delta }_{j𝜼}=𝒪(n_\mathrm{\Lambda }^pN^1)`$ and use (7). The evolution in time of the sublattice activities is then found to be deterministic in the $`N\mathrm{}`$ limit if $`lim_N\mathrm{}p/\mathrm{log}N=0`$. Furthermore, condition (9) holds, since
$$\underset{j=1}{\overset{N}{}}w_j(𝝈)\mathrm{\Delta }_{j𝜼}(𝝈)=\mathrm{tanh}[\beta \underset{𝜼^{}}{}p_𝜼^{}Q(𝜼;𝜼^{})m_𝜼^{}]m_𝜼$$
We may conclude that the situation is that described by (10), and that the evolution in time of the sublattice activities is governed by the following autonomous set of differential equations :
$$\frac{d}{dt}m_𝜼=\mathrm{tanh}[\beta \underset{𝜼^{}}{}p_𝜼^{}Q(𝜼;𝜼^{})m_𝜼^{}]m_𝜼$$
(14)
We see that, in contrast to the equilibrium techniques as described in , here there is no need at all to require symmetry of the interaction matrix or absence of self-interactions. In the symmetric case $`Q(𝒙;𝒚)=Q(𝒚;𝒙)`$ the system will approach equilibrium; if the kernel $`Q`$ is positive definite this can be shown, for instance, by inspection of the Lyapunov function<sup>3</sup><sup>3</sup>3A function of the state variables which is bounded from below and whose value decreases monotonically during the dynamics, see e.g. . Its existence guarantees evolution towards a stationary state (under some weak conditions). $`\{m_𝜼\}`$:
$$\{m_𝜼\}=\frac{1}{2}\underset{𝜼𝜼^{}}{}p_𝜼m_𝜼Q(𝜼;𝜼^{})m_𝜼^{}p_𝜼^{}\frac{1}{\beta }\underset{𝜼}{}p_𝜼\mathrm{log}\mathrm{cosh}[\beta \underset{𝜼^{}}{}Q(𝜼;𝜼^{})m_𝜼^{}p_𝜼^{}]$$
which is bounded from below and obeys:
$$\frac{d}{dt}=\underset{𝜼𝜼^{}}{}\left[p_𝜼\frac{d}{dt}m_𝜼\right]Q(𝜼;𝜼^{})\left[p_𝜼^{}\frac{d}{dt}m_𝜼^{}\right]0$$
(15)
Note that from the sublattice activities, in turn, follow the ‘overlaps’ $`m_\mu (𝝈)`$ (see ):
$$m_\mu (𝝈)=\frac{1}{N}\underset{i=1}{\overset{N}{}}\xi _i^\mu \sigma _i=\underset{𝜼}{}p_𝜼\eta _\mu m_𝜼$$
(16)
Simple examples of relevant models of the type (11), the dynamics of which are for large $`N`$ described by equation (14), are for instance the ones where one applies a non-linear operation $`\mathrm{\Phi }`$ to the standard Hopfield-type (or Hebbian-type) interactions . This non-linearity could result from e.g. a clipping procedure or from retaining only the sign of the Hebbian values:
$$J_{ij}=\frac{1}{N}\mathrm{\Phi }(\underset{\mu p}{}\xi _i^\mu \xi _j^\mu ):\mathrm{e}.\mathrm{g}.\mathrm{\Phi }(x)=\{\begin{array}{ccc}K& \mathrm{for}& xK\\ x& \mathrm{for}& K<x<K\\ K& \mathrm{for}& xK\end{array}\mathrm{or}\mathrm{\Phi }(x)=\mathrm{sgn}(x)$$
The effect of introducing such non-linearities is found to be of a quantitative nature, giving rise to little more than a re-scaling of critical noise levels and storage capacities. I will not go into full details, these can be found in e.g. , but illustrate this statement by working out the $`p=2`$ equations for randomly drawn pattern bits $`\xi _i^\mu \{1,1\}`$, where there are only four sub-lattices, and where $`p_𝜼=\frac{1}{4}`$ for all $`𝜼`$. Using $`\mathrm{\Phi }(0)=0`$ and $`\mathrm{\Phi }(x)=\mathrm{\Phi }(x)`$ (as with the above examples) we obtain from (14):
$$\frac{d}{dt}m_𝜼=\mathrm{tanh}[\frac{1}{4}\beta \mathrm{\Phi }(2)(m_𝜼m_𝜼)]m_𝜼$$
(17)
Here the choice made for $`\mathrm{\Phi }(x)`$ shows up only as a re-scaling of the temperature. From (17) we further obtain $`\frac{d}{dt}(m_𝜼+m_𝜼)=(m_𝜼+m_𝜼)`$. The system decays exponentially towards a state where, according to (16), $`m_𝜼=m_𝜼`$ for all $`𝜼`$. If at $`t=0`$ this is already the case, we find (at least for $`p=2`$) decoupled equations for the sub-lattice activities.
Separable Models: Description at the Level of Overlaps. Equations (14,16) suggest that at the level of overlaps there will be, in turn, closed laws if the kernel $`Q`$ is bi-linear:<sup>4</sup><sup>4</sup>4Strictly speaking, it is already sufficient to have a kernel which is linear in $`𝒚`$ only, i.e. $`Q(𝒙;𝒚)=_\nu f_\nu (𝒙)y_\nu `$, $`Q(𝒙;𝒚)=_{\mu \nu }x_\mu A_{\mu \nu }y_\nu `$, or:
$$J_{ij}=\frac{1}{N}\underset{\mu \nu =1}{\overset{p}{}}\xi _i^\mu A_{\mu \nu }\xi _j^\nu 𝝃_i=(\xi _i^1,\mathrm{},\xi _i^p)$$
(18)
We will see that now the $`\xi _i^\mu `$ need not be drawn from a finite discrete set (as long as they do not depend on $`N`$). The Hopfield model corresponds to $`A_{\mu \nu }=\delta _{\mu \nu }`$ and $`\xi _i^\mu \{1,1\}`$. The fields $`h_k`$ can now be written in terms of the overlaps $`m_\mu `$:
$$h_k(𝝈)=𝝃_kA𝒎(𝝈)𝒎=(m_1,\mathrm{},m_p)m_\mu (𝝈)=\frac{1}{N}\underset{i=1}{\overset{N}{}}\xi _i^\mu \sigma _i$$
(19)
For this choice of macroscopic variables we find $`\mathrm{\Delta }_{j\mu }=𝒪(N^1)`$, so the evolution of the vector $`𝒎`$ becomes deterministic for $`N\mathrm{}`$ if, according to (7), $`lim_N\mathrm{}p/\sqrt{N}=0`$. Again (9) holds, since
$$\underset{j=1}{\overset{N}{}}w_j(𝝈)\mathrm{\Delta }_{j\mu }(𝝈)=\frac{1}{N}\underset{k=1}{\overset{N}{}}𝝃_k\mathrm{tanh}\left[\beta 𝝃_kA𝒎\right]𝒎$$
Thus the evolution in time of the overlap vector $`𝒎`$ is governed by a closed set of differential equations:
$$\frac{d}{dt}𝒎=𝝃\mathrm{tanh}\left[\beta 𝝃A𝒎\right]_𝝃𝒎\mathrm{\Phi }(𝝃)_𝝃=𝑑𝝃\rho (𝝃)\mathrm{\Phi }(𝝃)$$
(20)
with $`\rho (𝝃)=lim_N\mathrm{}N^1_i\delta [𝝃𝝃_i]`$. Symmetry of the synapses is not required. For certain non-symmetric matrices $`A`$ one finds stable limit-cycle solutions of (20). In the symmetric case $`A_{\mu \nu }=A_{\nu \mu }`$ the system will approach equilibrium; the Lyapunov function (15) for positive definite matrices $`A`$ now becomes:
$$\{𝒎\}=\frac{1}{2}𝒎A𝒎\frac{1}{\beta }\mathrm{log}\mathrm{cosh}\left[\beta 𝝃A𝒎\right]_𝝃$$
Figure 1 shows in the $`m_1,m_2`$-plane the result of solving the macroscopic laws (20) numerically for $`p=2`$, randomly drawn pattern bits $`\xi _i^\mu \{1,1\}`$, and two choices of the matrix $`A`$. The first choice (upper row) corresponds to the Hopfield model; as the noise level $`T=\beta ^1`$ increases the amplitudes of the four attractors (corresponding to the two patterns $`𝝃^\mu `$ and their mirror images $`𝝃^\mu `$) continuously decrease, until at the critical noise level $`T_c=1`$ (see also ) they merge into the trivial attractor $`𝒎=(0,0)`$. The second choice corresponds to a non-symmetric model (i.e. without detailed balance); at the macroscopic level of description (at finite time-scales) the system clearly does not approach equilibrium; macroscopic order now manifests itself in the form of a limit-cycle (provided the noise level $`T`$ is below the critical value $`T_c=1`$ where this limit-cycle is destroyed). To what extent the laws (20) are in agreement with the result of performing the actual simulations in finite systems is illustrated in figure 2. Other examples can be found in .
As a second simple application of the flow equations (20) we turn to the relaxation times corresponding to the attractors of the Hopfield model (where $`A_{\mu \nu }=\delta _{\mu \nu }`$). Expanding (20) near a stable fixed-point $`𝒎^{}`$, i.e. $`𝒎(t)=𝒎^{}+𝒙(t)`$ with $`|𝒙(t)|1`$, gives the linearised equation
$$\frac{d}{dt}x_\mu =[\beta \underset{\nu }{}\xi _\mu \xi _\nu \mathrm{tanh}[\beta 𝝃𝒎^{}]_𝝃\delta _{\mu \nu }]x_\nu +𝒪(𝒙^2)$$
(21)
The Jacobian of (20), which determines the linearised equation (21), turns out to be minus the curvature matrix of the free energy surface at the fixed-point (c.f. the derivations in ). The asymptotic relaxation towards any stable attractor is generally exponential, with a characteristic time $`\tau `$ given by the inverse of the smallest eigenvalue of the curvature matrix. If, in particular, for the fixed point $`𝒎^{}`$ we substitute an $`n`$-mixture state, i.e. $`m_\mu =m_n(\mu n)`$ and $`m_\mu =0(\mu >n)`$, and transform (21) to the basis where the corresponding curvature matrix $`𝑫^{(n)}`$ (with eigenvalues $`D_\lambda ^n`$) is diagonal, $`𝒙\stackrel{~}{𝒙}`$, we obtain
$$\stackrel{~}{x}_\lambda (t)=\stackrel{~}{x}_\lambda (0)e^{tD_\lambda ^n}+\mathrm{}$$
so $`\tau ^1=\mathrm{min}_\lambda D_\lambda ^n`$, which we have already calculated (see ) in determining the character of the saddle-points of the free-energy surface. The result is shown in figure 3. The relaxation time for the $`n`$-mixture attractors decreases monotonically with the degree of mixing $`n`$, for any noise level. At the transition where a macroscopic state $`𝒎^{}`$ ceases to correspond to a local minimum of the free energy surface, it also de-stabilises in terms of the linearised dynamic equation (21) (as it should). The Jacobian develops a zero eigenvalue, the relaxation time diverges, and the long-time behaviour is no longer obtained from the linearised equation. This gives rise to critical slowing down (power law relaxation as opposed to exponential relaxation). For instance, at the transition temperature $`T_c=1`$ for the $`n=1`$ (pure) state, we find by expanding (20):
$$\frac{d}{dt}m_\mu =m_\mu [\frac{2}{3}m_\mu ^2𝒎^2]+𝒪(𝒎^5)$$
which gives rise to a relaxation towards the trivial fixed-point of the form $`𝒎t^{\frac{1}{2}}`$.
If one is willing to restrict oneself to the limited class of models (18) (as opposed to the more general class (11)) and to the more global level of description in terms of $`p`$ overlap parameters $`m_\mu `$ instead of $`n_\mathrm{\Lambda }^p`$ sublattice activities $`m_𝜼`$, then there are two rewards. Firstly there will be no restrictions on the stored pattern components $`\xi _i^\mu `$ (for instance, they are allowed to be real-valued); secondly the number $`p`$ of patterns stored can be much larger for the deterministic autonomous dynamical laws to hold ($`p\sqrt{N}`$ instead of $`p\mathrm{log}N`$, which from a biological point of view is not impressive.
### 2.3 Closed Macroscopic Laws for Parallel Dynamics
We now turn to the parallel dynamics counterpart of (1), i.e. the Markov chain
$$p_{\mathrm{}+1}(𝝈)=\underset{𝝈^{}}{}W[𝝈;𝝈^{}]p_{\mathrm{}}(𝝈^{})W[𝝈;𝝈^{}]=\underset{i=1}{\overset{N}{}}\frac{1}{2}\left[1+\sigma _i\mathrm{tanh}[\beta h_i(𝝈^{})]\right]$$
(22)
(with $`\sigma _i\{1,1\}`$, and with local fields $`h_i(𝝈)`$ defined in the usual way). The evolution of macroscopic probability densities will here be described by discrete mappings, in stead of differential equations.
The Toy Model. Let us first see what happens to our previous toy model: $`J_{ij}=(J/N)\eta _i\xi _j`$ and $`\theta _i=0`$. As before we try to describe the dynamics at the (macroscopic) level of the quantity $`m(𝝈)=\frac{1}{N}_k\xi _k\sigma _k`$. The evolution of the macroscopic probability density $`𝒫_t[m]`$ is obtained by inserting (22):
$$𝒫_{t+1}[m]=\underset{𝝈𝝈^{}}{}\delta \left[mm(𝝈)\right]W[𝝈;𝝈^{}]p_t(𝝈^{})=𝑑m^{}\stackrel{~}{W}_t[m,m^{}]𝒫_t[m^{}]$$
(23)
with
$$\stackrel{~}{W}_t[m,m^{}]=\frac{_{𝝈𝝈^{}}\delta \left[mm(𝝈)\right]\delta \left[m^{}m(𝝈^{})\right]W[𝝈;𝝈^{}]p_t(𝝈^{})}{_𝝈^{}\delta \left[m^{}m(𝝈^{})\right]p_t(𝝈^{})}$$
We now insert our expression for the transition probabilities $`W[𝝈;𝝈^{}]`$ and for the local fields. Since the fields depend on the microscopic state $`𝝈`$ only through $`m(𝝈)`$, the distribution $`p_t(𝝈)`$ drops out of the above expression for $`\stackrel{~}{W}_t`$ which thereby loses its explicit time-dependence, $`\stackrel{~}{W}_t[m,m^{}]\stackrel{~}{W}[m,m^{}]`$:
$$\stackrel{~}{W}[m,m^{}]=e^{_i\mathrm{log}\mathrm{cosh}(\beta Jm^{}\eta _i)}\delta \left[mm(𝝈)\right]e^{\beta Jm^{}_i\eta _i\sigma _i}_𝝈\text{}\mathrm{with}\mathrm{}_𝝈=2^N\underset{𝝈}{}\mathrm{}$$
Inserting the integral representation for the $`\delta `$-function allows us to perform the average:
$$\stackrel{~}{W}[m,m^{}]=\left[\frac{\beta N}{2\pi }\right]𝑑ke^{N\mathrm{\Psi }(m,m^{},k)}$$
$$\mathrm{\Psi }=i\beta km+\mathrm{log}\mathrm{cosh}\beta [J\eta m^{}ik\xi ]_{\eta ,\xi }\mathrm{log}\mathrm{cosh}\beta [J\eta m^{}]_\eta \text{}$$
Since $`\stackrel{~}{W}[m,m^{}]`$ is (by construction) normalised, $`𝑑m\stackrel{~}{W}[m,m^{}]=1`$, we find that for $`N\mathrm{}`$ the expectation value with respect to $`\stackrel{~}{W}[m,m^{}]`$ of any sufficiently smooth function $`f(m)`$ will be determined only by the value $`m^{}(m^{})`$ of $`m`$ in the relevant saddle-point of $`\mathrm{\Psi }`$:
$$𝑑mf(m)\stackrel{~}{W}[m,m^{}]=\frac{𝑑m𝑑kf(m)e^{N\mathrm{\Psi }(m,m^{},k)}}{𝑑m𝑑ke^{N\mathrm{\Psi }(m,m^{},k)}}f(m^{}(m^{}))(N\mathrm{})$$
Variation of $`\mathrm{\Psi }`$ with respect to $`k`$ and $`m`$ gives the two saddle-point equations:
$$m=\xi \mathrm{tanh}\beta [J\eta m^{}\xi k]_{\eta ,\xi },k=0$$
We may now conclude that $`lim_N\mathrm{}\stackrel{~}{W}[m,m^{}]=\delta \left[mm^{}(m^{})\right]`$ with $`m^{}(m^{})=\xi \mathrm{tanh}(\beta J\eta m^{})_{\eta ,\xi }`$, and that the macroscopic equation (23) becomes:
$$𝒫_{t+1}[m]=𝑑m^{}\delta \left[m\xi \mathrm{tanh}(\beta J\eta m^{})_{\eta \xi }\right]𝒫_t[m^{}](N\mathrm{})$$
This describes deterministic evolution. If at $`t=0`$ we know $`m`$ exactly, this will remain the case for finite time-scales, and $`m`$ will evolve according to a discrete version of the sequential dynamics law (2):
$$m_{t+1}=\xi \mathrm{tanh}[\beta J\eta m_t]_{\eta ,\xi }\text{}$$
(24)
Arbitrary Synapses. We now try to generalise the above approach to less trivial classes of models. As for the sequential case we will find in the limit $`N\mathrm{}`$ closed deterministic evolution equations for a more general set of intensive macroscopic state variables $`𝛀(𝝈)=(\mathrm{\Omega }_1(𝝈),\mathrm{},\mathrm{\Omega }_n(𝝈)`$ if the local fields $`h_i(𝝈)`$ depend on the microscopic state $`𝝈`$ only through the values of $`𝛀(𝝈)`$, and if the number $`n`$ of these state variables necessary to do so is not too large. The evolution of the ensemble probability density (3) is now obtained by inserting the Markov equation (22):
$$𝒫_{t+1}\left[𝛀\right]=𝑑𝛀^{}\stackrel{~}{W}_t[𝛀,𝛀^{}]𝒫_t\left[𝛀^{}\right]$$
(25)
$$\stackrel{~}{W}_t[𝛀,𝛀^{}]=\frac{_{𝝈𝝈^{}}\delta \left[𝛀𝛀(𝝈)\right]\delta \left[𝛀^{}𝛀(𝝈^{})\right]W[𝝈;𝝈^{}]p_t(𝝈^{})}{_𝝈^{}\delta \left[𝛀^{}𝛀(𝝈^{})\right]p_t(𝝈^{})}$$
$$=\delta \left[𝛀𝛀(𝝈)\right]e^{_i\left[\beta \sigma _ih_i(𝝈^{})\mathrm{log}\mathrm{cosh}(\beta h_i(𝝈^{}))\right]}_{𝛀^{};t}_𝝈\text{}$$
(26)
with $`\mathrm{}_𝝈=2^N_𝝈\mathrm{}`$, and with the conditional (or sub-shell) average defined as in (5). It is clear from (26) that in order to find autonomous macroscopic laws, i.e. for the distribution $`p_t(𝝈)`$ to drop out, the local fields must depend on the microscopic state $`𝝈`$ only through the macroscopic quantities $`𝛀(𝝈)`$: $`h_i(𝝈)=h_i[𝛀(𝝈)]`$. In this case $`\stackrel{~}{W}_t`$ loses its explicit time-dependence, $`\stackrel{~}{W}_t[𝛀,𝛀^{}]\stackrel{~}{W}[𝛀,𝛀^{}]`$. Inserting integral representations for the $`\delta `$-functions leads to:
$$\stackrel{~}{W}[𝛀,𝛀^{}]=\left[\frac{\beta N}{2\pi }\right]^n𝑑𝑲e^{N\mathrm{\Psi }(𝛀,𝛀^{},𝑲)}$$
$$\mathrm{\Psi }=i\beta 𝑲𝛀+\frac{1}{N}\mathrm{log}e^{\beta \left[_i\sigma _ih_i[𝛀^{}]iN𝑲𝛀(𝝈)\right]}_𝝈\frac{1}{N}\underset{i}{}\mathrm{log}\mathrm{cosh}[\beta h_i[𝛀^{}]]\text{}$$
Using the normalisation $`𝑑𝛀\stackrel{~}{W}[𝛀,𝛀^{}]=1`$, we can write expectation values with respect to $`\stackrel{~}{W}[𝛀,𝛀^{}]`$ of macroscopic quantities $`f[𝛀]`$ as
$$𝑑𝛀f[𝛀]\stackrel{~}{W}[𝛀,𝛀^{}]=\frac{𝑑𝛀𝑑𝑲f[𝛀]e^{N\mathrm{\Psi }(𝛀,𝛀^{},𝑲)}}{𝑑𝛀𝑑𝑲e^{N\mathrm{\Psi }(𝛀,𝛀^{},𝑲)}}$$
(27)
For saddle-point arguments to apply in determining the leading order in $`N`$ of (27), we encounter restrictions on the number $`n`$ of our macroscopic quantities (as expected), since $`n`$ determines the dimension of the integrations in (27). The restrictions can be found by expanding $`\mathrm{\Psi }`$ around its maximum $`\mathrm{\Psi }^{}`$. After defining $`𝒙=(𝛀,𝑲)`$, of dimension $`2n`$, and after translating the location of the maximum to the origin, one has
$$\mathrm{\Psi }(𝒙)=\mathrm{\Psi }^{}\frac{1}{2}\underset{\mu \nu }{}x_\mu x_\nu H_{\mu \nu }+\underset{\mu \nu \rho }{}x_\mu x_\nu x_\rho L_{\mu \nu \rho }+𝒪(𝒙^4)$$
giving
$$\frac{𝑑𝒙g(𝒙)e^{N\mathrm{\Psi }(𝒙)}}{𝑑𝒙g(𝒙)e^{N\mathrm{\Psi }(𝒙)}}g(\mathrm{𝟎})=\frac{𝑑𝒙[g(𝒙)g(\mathrm{𝟎})]e^{\frac{1}{2}N𝒙𝑯𝒙+N_{\mu \nu \rho }x_\mu x_\nu x_\rho L_{\mu \nu \rho }+𝒪(N𝒙^4)}}{𝑑𝒙e^{\frac{1}{2}N𝒙𝑯𝒙+N_{\mu \nu \rho }x_\mu x_\nu x_\rho L_{\mu \nu \rho }+𝒪(N𝒙^4)}}$$
$$=\frac{𝑑𝒚[g(𝒚/\sqrt{N})g(\mathrm{𝟎})]e^{\frac{1}{2}𝒚𝑯𝒚+_{\mu \nu \rho }y_\mu y_\nu y_\rho L_{\mu \nu \rho }/\sqrt{N}+𝒪(𝒚^4/N)}}{𝑑𝒚e^{\frac{1}{2}𝒚𝑯𝒚+_{\mu \nu \rho }y_\mu y_\nu y_\rho L_{\mu \nu \rho }/\sqrt{N}+𝒪(𝒚^4/N)}}$$
$$=\frac{𝑑𝒚\left[N^{\frac{1}{2}}𝒚\mathbf{}g(\mathrm{𝟎})+𝒪(𝒚^2/N)\right]e^{\frac{1}{2}𝒚𝑯𝒚}\left[1+_{\mu \nu \rho }y_\mu y_\nu y_\rho L_{\mu \nu \rho }/\sqrt{N}+𝒪(𝒚^6/N)\right]}{𝑑𝒚e^{\frac{1}{2}𝒚𝑯𝒚}\left[1+_{\mu \nu \rho }y_\mu y_\nu y_\rho L_{\mu \nu \rho }/\sqrt{N}+𝒪(𝒚^6/N)\right]}$$
$$=𝒪(n^2/N)+𝒪(n^4/N^2)+\mathrm{non}\mathrm{dominant}\mathrm{terms}(N,n\mathrm{})$$
with $`𝑯`$ denoting the Hessian (curvature) matrix of the surface $`\mathrm{\Psi }`$ at the minimum $`\mathrm{\Psi }^{}`$. We thus find
$$\underset{N\mathrm{}}{lim}n/\sqrt{N}=0:\underset{N\mathrm{}}{lim}d𝛀f[𝛀]\stackrel{~}{W}[𝛀,𝛀^{}]=f\left[𝛀^{}(𝛀^{})\right]$$
where $`𝛀^{}(𝛀^{})`$ denotes the value of $`𝛀`$ in the saddle-point where $`\mathrm{\Psi }`$ is minimised. Variation of $`\mathrm{\Psi }`$ with respect to $`𝛀`$ and $`𝑲`$ gives the saddle-point equations:
$$𝛀=\frac{𝛀(𝝈)e^{\beta \left[_i\sigma _ih_i[𝛀^{}]iN𝑲𝛀(𝝈)\right]}_𝝈}{e^{\beta \left[_i\sigma _ih_i[𝛀^{}]iN𝑲𝛀(𝝈)\right]}_𝝈},𝑲=0$$
We may now conclude that $`lim_N\mathrm{}\stackrel{~}{W}[𝛀,𝛀^{}]=\delta \left[𝛀𝛀^{}(𝛀^{})\right]`$, with
$$𝛀^{}(𝛀^{})=\frac{𝛀(𝝈)e^{\beta _i\sigma _ih_i[𝛀^{}]}_𝝈}{e^{\beta _i\sigma _ih_i[𝛀^{}]}_𝝈}$$
and that for $`N\mathrm{}`$ the macroscopic equation (25) becomes $`𝒫_{t+1}[𝛀]=𝑑𝛀^{}\delta [𝛀𝛀^{}(𝛀^{})]𝒫_t[𝛀^{}]`$. This relation again describes deterministic evolution. If at $`t=0`$ we know $`𝛀`$ exactly, this will remain the case for finite time-scales and $`𝛀`$ will evolve according to
$$𝛀(t+1)=\frac{𝛀(𝝈)e^{\beta _i\sigma _ih_i[𝛀(t)]}_𝝈}{e^{\beta _i\sigma _ih_i[𝛀(t)]}_𝝈}$$
(28)
As with the sequential case, in taking the limit $`N\mathrm{}`$ we have to keep in mind that the resulting laws apply to finite $`t`$, and that for sufficiently large times terms of higher order in $`N`$ do come into play. As for the sequential case, a more rigorous and tedious analysis shows that the restriction $`n/\sqrt{N}0`$ can in fact be weakened to $`n/N0`$. Finally, for macroscopic quantities $`𝛀(𝝈)`$ which are linear in $`𝝈`$, the remaining $`𝝈`$-averages become trivial, so that :
$$\mathrm{\Omega }_\mu (𝝈)=\frac{1}{N}\underset{i}{}\omega _{\mu i}\sigma _i:\mathrm{\Omega }_\mu (t+1)=\underset{N\mathrm{}}{lim}\frac{1}{N}\underset{i}{}\omega _{\mu i}\mathrm{tanh}\left[\beta h_i[𝛀(t)]\right]$$
(29)
(to be compared with (10), as derived for sequential dynamics).
### 2.4 Application to Separable Attractor Networks
Separable models: Sublattice Activities and Overlaps. The separable attractor models (11), described at the level of sublattice activities (13), indeed have the property that all local fields can be written in terms of the macroscopic observables. What remains to ensure deterministic evolution is meeting the condition on the number of sublattices. If all relative sublattice sizes $`p_𝜼`$ are of the same order in $`N`$ (as for randomly drawn patterns) this condition again translates into $`lim_N\mathrm{}p/\mathrm{log}N=0`$ (as for sequential dynamics). Since the sublattice activities are linear functions of the $`\sigma _i`$, their evolution in time is governed by equation (29), which acquires the form:
$$m_𝜼(t+1)=\mathrm{tanh}[\beta \underset{𝜼^{}}{}p_𝜼^{}Q(𝜼;𝜼^{})m_𝜼^{}(t)]$$
(30)
As for sequential dynamics, symmetry of the interaction matrix does not play a role.
At the more global level of overlaps $`m_\mu (𝝈)=N^1_i\xi _i^\mu \sigma _i`$ we, in turn, obtain autonomous deterministic laws if the local fields $`h_i(𝝈)`$ can be expressed in terms if $`𝒎(𝝈)`$ only, as for the models (18) (or, more generally, for all models in which the interactions are of the form $`J_{ij}=_{\mu p}f_{i\mu }\xi _j^\mu `$), and with the following restriction on the number $`p`$ of embedded patterns: $`lim_N\mathrm{}p/\sqrt{N}=0`$ (as with sequential dynamics). For the bi-linear models (18), the evolution in time of the overlap vector $`𝒎`$ (which depends linearly on the $`\sigma _i`$) is governed by (29), which now translates into the iterative map:
$$𝒎(t+1)=𝝃\mathrm{tanh}[\beta 𝝃A𝒎(t)]_𝝃$$
(31)
with $`\rho (𝝃)`$ as defined in (20). Again symmetry of the synapses is not required. For parallel dynamics it is far more difficult than for sequential dynamics to construct Lyapunov functions, and prove that the macroscopic laws (31) for symmetric systems evolve towards a stable fixed-point (as one would expect), but it can still be done. For non-symmetric systems the macroscopic laws (31) can in principle display all the interesting, but complicated, phenomena of non-conservative non-linear systems. Nevertheless, it is also not uncommon that the equations (31) for non-symmetric systems can be mapped by a time-dependent transformation onto the equations for related symmetric systems (mostly variants of the original Hopfield model).
As an example we show in figure 4 as functions of time the values of the overlaps $`\{m_\mu \}`$ for $`p=10`$ and $`T=0.5`$, resulting from numerical iteration of the macroscopic laws (31) for the model
$$J_{ij}=\frac{\nu }{N}\underset{\mu }{}\xi _i^\mu \xi _j^\mu +\frac{1\nu }{N}\underset{\mu }{}\xi _i^{\mu +1}\xi _j^\mu (\mu :\mathrm{mod}p)$$
i.e. $`A_{\lambda \rho }=\nu \delta _{\lambda \rho }+(1\nu )\delta _{\lambda ,\rho +1}(\lambda ,\rho :\mathrm{mod}p)`$, with randomly drawn pattern bits $`\xi _i^\mu \{1,1\}`$ The initial state is chosen to be the pure state $`m_\mu =\delta _{\mu ,1}`$. At intervals of $`\mathrm{\Delta }t=20`$ iterations the parameter $`\nu `$ is reduced in $`\mathrm{\Delta }\nu =0.25`$ steps from $`\nu =1`$ (where one recovers the symmetric Hopfield model) to $`\nu =0`$ (where one obtains a non-symmetric model which processes the $`p`$ embedded patterns in strict sequential order as a period-$`p`$ limit-cycle). The analysis of the equations (31) for the pure sequence processing case $`\nu =0`$ is greatly simplified by mapping the model onto the ordinary ($`\nu =1`$) Hopfield model, using the index permutation symmetries of the present pattern distribution, as follows (all pattern indices are periodic, mod $`p`$). Define $`m_\mu (t)=M_{\mu t}(t)`$, now
$$M_\mu (t+1)=\xi _{\mu +t+1}\mathrm{tanh}[\beta \underset{\rho }{}\xi _{\rho +1}M_{\rho t}(t)]_𝝃=\xi _\mu \mathrm{tanh}[\beta 𝝃𝑴(t)]_𝝃$$
We can now immediately infer, in particular, that to each stable macroscopic fixed-point attractor of the original Hopfield model corresponds a stable period-$`p`$ macroscopic limit-cycle attractor in the $`\nu =1`$ sequence processing model (e.g. pure states $``$ pure sequences, mixture states $``$ mixture sequences), with identical amplitude as a function of the noise level. Figure 4 shows for $`\nu =0`$ (i.e. $`t>80`$) a relaxation towards such a pure sequence.
Finally we note that the fixed-points of the macroscopic equations (14) and (20) (derived for sequential dynamics) are identical to those of (30) and (31) (derived for parallel dynamics). The stability properties of these fixed points, however, need not be the same, and have to be assessed on a case-by-case basis. For the Hopfield model, i.e. equations (20,31) with $`A_{\mu \nu }=\delta _{\mu \nu }`$, they are found to be the same, but already for $`A_{\mu \nu }=\delta _{\mu \nu }`$ the two types of dynamics would behave differently.
## 3 Attractor Neural Networks with Continuous Neurons
### 3.1 Closed Macroscopic Laws
General Derivation. We have seen in ) that models of recurrent neural networks with continuous neural variables (e.g. graded response neurons or coupled oscillators) can often be described by a Fokker-Planck equation for the microscopic state probability density $`p_t(𝝈)`$:
$$\frac{d}{dt}p_t(𝝈)=\underset{i}{}\frac{}{\sigma _i}\left[p_t(𝝈)f_i(𝝈)\right]+T\underset{i}{}\frac{^2}{\sigma _i^2}p_t(𝝈)$$
(32)
Averages over $`p_t(𝝈)`$ are denoted by $`G=𝑑𝝈p_t(𝝈)G(𝝈,t)`$. From (32) one obtains directly (through integration by parts) an equation for the time derivative of averages:
$$\frac{d}{dt}G=\frac{G}{t}+\underset{i}{}\left[f_i(𝝈)+T\frac{}{\sigma _i}\right]\frac{G}{\sigma _i}$$
(33)
In particular, if we apply (33) to $`G(𝝈,t)=\delta [𝛀𝛀(𝝈)]`$, for any set of macroscopic observables $`𝛀(𝝈)=(\mathrm{\Omega }_1(𝝈),\mathrm{},\mathrm{\Omega }_n(𝝈))`$ (in the spirit of the previous section), we obtain a dynamic equation for the macroscopic probability density $`P_t(𝛀)=\delta [𝛀𝛀(𝝈)]`$, which is again of the Fokker-Planck form:
$$\frac{d}{dt}P_t(𝛀)=\underset{\mu }{}\frac{}{\mathrm{\Omega }_\mu }\left\{P_t(𝛀)\underset{i}{}\left[f_i(𝝈)+T\frac{}{\sigma _i}\right]\frac{}{\sigma _i}\mathrm{\Omega }_\mu (𝝈)_{𝛀;t}\right\}$$
$$+T\underset{\mu \nu }{}\frac{^2}{\mathrm{\Omega }_\mu \mathrm{\Omega }_\nu }\left\{P_t(𝛀)\underset{i}{}\left[\frac{}{\sigma _i}\mathrm{\Omega }_\mu (𝝈)\right]\left[\frac{}{\sigma _i}\mathrm{\Omega }_\nu (𝝈)\right]_{𝛀;t}\right\}$$
(34)
with the conditional (or sub-shell) averages:
$$G(𝝈)_{𝛀,t}=\frac{𝑑𝝈p_t(𝝈)\delta [𝛀𝛀(𝝈)]G(𝝈)}{𝑑𝝈p_t(𝝈)\delta [𝛀𝛀(𝝈)]}$$
(35)
From (34) we infer that a sufficient condition for the observables $`𝛀(𝝈)`$ to evolve in time deterministically (i.e. for having vanishing diffusion matrix elements in (34)) in the limit $`N\mathrm{}`$ is
$$\underset{N\mathrm{}}{lim}\underset{i}{}\left[\underset{\mu }{}\left|\frac{}{\sigma _i}\mathrm{\Omega }_\mu (𝝈)\right|\right]^2_{𝛀;t}=0$$
(36)
If (36) holds, the macroscopic Fokker-Planck equation (34) reduces for $`N\mathrm{}`$ to a Liouville equation, and the observables $`𝛀(𝝈)`$ will evolve in time according to the coupled deterministic equations:
$$\frac{d}{dt}\mathrm{\Omega }_\mu =\underset{N\mathrm{}}{lim}\underset{i}{}\left[f_i(𝝈)+T\frac{}{\sigma _i}\right]\frac{}{\sigma _i}\mathrm{\Omega }_\mu (𝝈)_{𝛀;t}$$
(37)
The deterministic macroscopic equation (37), together with its associated condition for validity (36) will form the basis for the subsequent analysis.
Closure: A Toy Model Again. The general derivation given above went smoothly. However, the equations (37) are not yet closed. It turns out that to achieve closure even for simple continuous networks we can no longer get away with just a finite (small) number of macroscopic observables (as with binary neurons). This I will now illustrate with a simple toy network of graded response neurons:
$$\frac{d}{dt}u_i(t)=\underset{j}{}J_{ij}g[u_j(t)]u_i(t)+\eta _i(t)$$
(38)
with $`g[z]=\frac{1}{2}[\mathrm{tanh}(\gamma z)+1]`$ and with the standard Gaussian white noise $`\eta _i(t)`$ (see ). In the language of (32) this means $`f_i(𝒖)=_jJ_{ij}g[u_j]u_i`$. We choose uniform synapses $`J_{ij}=J/N`$, so $`f_i(𝒖)(J/N)_jg[u_j]u_i`$. If (36) were to hold, we would find the deterministic macroscopic laws
$$\frac{d}{dt}\mathrm{\Omega }_\mu =\underset{N\mathrm{}}{lim}\underset{i}{}[\frac{J}{N}\underset{j}{}g[u_j]u_i+T\frac{}{u_i}]\frac{}{u_i}\mathrm{\Omega }_\mu (𝒖)_{𝛀;t}$$
(39)
In contrast to similar models with binary neurons, choosing as our macroscopic level of description $`𝛀(𝒖)`$ again simply the average $`m(𝒖)=N^1_iu_i`$ now leads to an equation which fails to close:
$$\frac{d}{dt}m=\underset{N\mathrm{}}{lim}J\frac{1}{N}\underset{j}{}g[u_j]_{m;t}m$$
The term $`N^1_jg[u_j]`$ cannot be written as a function of $`N^1_iu_i`$. We might be tempted to try dealing with this problem by just including the offending term in our macroscopic set, and choose $`𝛀(𝒖)=(N^1_iu_i,N^1_ig[u_i])`$. This would indeed solve our closure problem for the $`m`$-equation, but we would now find a new closure problem in the equation for the newly introduced observable. The only way out is to choose an observable function, namely the distribution of potentials
$$\rho (u;𝒖)=\frac{1}{N}\underset{i}{}\delta [uu_i],\rho (u)=\rho (u;𝒖)=\frac{1}{N}\underset{i}{}\delta [uu_i]$$
(40)
This is to be done with care, in view of our restriction on the number of observables: we evaluate (40) at first only for $`n`$ specific values $`u_\mu `$ and take the limit $`n\mathrm{}`$ only after the limit $`N\mathrm{}`$. Thus we define $`\mathrm{\Omega }_\mu (𝒖)=\frac{1}{N}_i\delta [u_\mu u_i]`$, condition (36) reduces to the familiar expression $`lim_N\mathrm{}n/\sqrt{N}=0`$, and we get for $`N\mathrm{}`$ and $`n\mathrm{}`$ (taken in that order) from (39) a diffusion equation for the distribution of membrane potentials (describing a so-called ‘time-dependent Ornstein-Uhlenbeck process’ ):
$$\frac{d}{dt}\rho (u)=\frac{}{u}\left\{\rho (u)\left[J𝑑u^{}\rho (u^{})g[u^{}]u\right]\right\}+T\frac{^2}{u^2}\rho (u)$$
(41)
The natural<sup>5</sup><sup>5</sup>5For non-Gaussian initial conditions $`\rho _0(u)`$ the solution of (41) would in time converge towards the Gaussian solution. solution of (41) is the Gaussian distribution
$$\rho _t(u)=[2\pi \mathrm{\Sigma }^2(t)]^{\frac{1}{2}}e^{\frac{1}{2}[u\overline{u}(t)]^2/\mathrm{\Sigma }^2(t)}$$
(42)
in which $`\mathrm{\Sigma }=[T+(\mathrm{\Sigma }_0^2T)e^{2t}]^{\frac{1}{2}}`$, and $`\overline{u}`$ evolves in time according to
$$\frac{d}{dt}\overline{u}=JDzg[\overline{u}+\mathrm{\Sigma }z]\overline{u}$$
(43)
(with $`Dz=(2\pi )^{\frac{1}{2}}e^{\frac{1}{2}z^2}dz`$). We can now also calculate the distribution $`p(s)`$ of neuronal firing activities $`s_i=g[u_i]`$ at any time:
$$p(s)=𝑑u\rho (u)\delta [sg[u]]=\frac{\rho (g^{\mathrm{inv}}[s])}{_0^1𝑑s^{}\rho (g^{\mathrm{inv}}[s^{}])}$$
For our choice $`g[z]=\frac{1}{2}+\frac{1}{2}\mathrm{tanh}[\gamma z]`$ we have $`g^{\mathrm{inv}}[s]=\frac{1}{2\gamma }\mathrm{log}[s/(1s)]`$, so in combination with (42):
$$0<s<1:p(s)=\frac{e^{\frac{1}{2}[(2\gamma )^1\mathrm{log}[s/(1s)]\overline{u}]^2/\mathrm{\Sigma }^2}}{_0^1𝑑s^{}e^{\frac{1}{2}[(2\gamma )^1\mathrm{log}[s^{}/(1s^{})]\overline{u}]^2/\mathrm{\Sigma }^2}}$$
(44)
The results of solving and integrating numerically (43) and (44) is shown in figure 5, for Gaussian initial conditions (42) with $`\overline{u}_0=0`$ and $`\mathrm{\Sigma }_0=1`$, and with parameters $`\gamma =J=1`$ and different noise levels $`T`$. For low noise levels we find high average membrane potentials, low membrane potential variance, and high firing rates; for high noise levels the picture changes to lower average membrane potentials, higher potential variance, and uniformly distributed (noise-dominated) firing activities. The extreme cases $`T=0`$ and $`T=\mathrm{}`$ are easily extracted from our equations. For $`T=0`$ one finds $`\mathrm{\Sigma }(t)=\mathrm{\Sigma }_0e^t`$ and $`\frac{d}{dt}\overline{u}=Jg[\overline{u}]\overline{u}`$. This leads to a final state where $`\overline{u}=\frac{1}{2}J+\frac{1}{2}J\mathrm{tanh}[\gamma \overline{u}]`$ and where $`p(s)=\delta [s\overline{u}/J]`$. For $`T=\mathrm{}`$ one finds $`\mathrm{\Sigma }=\mathrm{}`$ (for any $`t>0`$) and $`\frac{d}{dt}\overline{u}=\frac{1}{2}J\overline{u}`$. This leads to an final state where $`\overline{u}=\frac{1}{2}J`$ and where $`p(s)=1`$ for all $`0<s<1`$.
None of the above results (not even those on the stationary state) could have been obtained within equilibrium statistical mechanics, since any network of connected graded response neurons will violate detailed balance . Secondly, there appears to be a qualitative difference between simple networks (e.g. $`J_{ij}=J/N`$) of binary neurons versus those of continuous neurons, in terms of the types of macroscopic observables needed for deriving closed deterministic laws: a single number $`m=N^1_i\sigma _i`$ versus a distribution $`\rho (\sigma )=N^1_i\delta [\sigma \sigma _i]`$. Note, however, that in the binary case the latter distribution would in fact have been been characterised fully by a single number: the average $`m`$, since $`\rho (\sigma )=\frac{1}{2}[1+m]\delta [\sigma 1]+\frac{1}{2}[1m]\delta [\sigma +1]`$. In other words: there we were just lucky.
### 3.2 Application to Graded Response Attractor Networks
Derivation of Closed Macroscopic Laws. I will now turn to attractor networks with graded response neurons of the type (38), in which $`p`$ binary patterns $`𝝃^\mu =(\xi _1^\mu ,\mathrm{},\xi _N^\mu )\{1,1\}^N`$ have been stored via separable Hebbian-type synapses (18): $`J_{ij}=(2/N)_{\mu \nu =1}^p\xi _i^\mu A_{\mu \nu }\xi _j^\nu `$ (the extra factor 2 is inserted for future convenience). Adding suitable thresholds $`\theta _i=\frac{1}{2}_jJ_{ij}`$ to the right-hand sides of (38), and choosing the non-linearity $`g[z]=\frac{1}{2}(1+\mathrm{tanh}[\gamma z])`$ would then give us
$$\frac{d}{dt}u_i(t)=\underset{\mu \nu }{}\xi _i^\mu A_{\mu \nu }\frac{1}{N}\underset{j}{}\xi _j^\nu \mathrm{tanh}[\gamma u_j(t)]u_i(t)+\eta _i(t)$$
so the deterministic forces are $`f_i(𝒖)=N^1_{\mu \nu }\xi _i^\mu A_{\mu \nu }_j\xi _j^\nu \mathrm{tanh}[\gamma u_j]u_i`$. Choosing our macroscopic observables $`𝛀(𝒖)`$ such that (36) holds, would lead to the deterministic macroscopic laws
$$\frac{d}{dt}\mathrm{\Omega }_\mu =\underset{N\mathrm{}}{lim}\underset{\mu \nu }{}A_{\mu \nu }\left[\frac{1}{N}\underset{j}{}\xi _j^\nu \mathrm{tanh}[\gamma u_j]\right]\left[\underset{i}{}\xi _i^\mu \frac{}{u_i}\mathrm{\Omega }_\mu (𝒖)\right]_{𝛀;t}+\underset{N\mathrm{}}{lim}\underset{i}{}\left[T\frac{}{u_i}u_i\right]\frac{}{u_i}\mathrm{\Omega }_\mu (𝒖)_{𝛀;t}$$
(45)
As with the uniform synapses case, the main problem to be dealt with is how to choose the $`\mathrm{\Omega }_\mu (𝒖)`$ such that (45) closes. It turns out that the canonical choice is to turn to the distributions of membrane potentials within each of the $`2^p`$ sub-lattices, as introduced in (12):
$$I_𝜼=\{i|𝝃_i=𝜼\}:\rho _𝜼(u;𝒖)=\frac{1}{|I_𝜼|}\underset{iI_𝜼}{}\delta [uu_i],\rho _𝜼(u)=\rho _𝜼(u;𝒖)$$
(46)
with $`𝜼\{1,1\}^p`$ and $`lim_N\mathrm{}|I_𝜼|/N=p_𝜼`$. Again we evaluate the distributions in (46) at first only for $`n`$ specific values $`u_\mu `$ and send $`n\mathrm{}`$ after $`N\mathrm{}`$. Now condition (36) reduces to $`lim_N\mathrm{}2^p/\sqrt{N}=0`$. We will keep $`p`$ finite, for simplicity. Using identities such as $`_i\mathrm{}=_𝜼_{iI_𝜼}\mathrm{}`$ and
$$iI_𝜼:\frac{}{u_i}\rho _𝜼(u;𝒖)=|I_𝜼|^1\frac{}{u}\delta [uu_i],\frac{^2}{u_i^2}\rho _𝜼(u;𝒖)=|I_𝜼|^1\frac{^2}{u^2}\delta [uu_i],$$
we then obtain for $`N\mathrm{}`$ and $`n\mathrm{}`$ (taken in that order) from equation (45) $`2^p`$ coupled diffusion equations for the distributions $`\rho _𝜼(u)`$ of membrane potentials in each of the $`2^p`$ sub-lattices $`I_𝜼`$:
$$\frac{d}{dt}\rho _𝜼(u)=\frac{}{u}\left\{\rho _𝜼(u)\left[\underset{\mu \nu =1}{\overset{p}{}}\eta _\mu A_{\mu \nu }\underset{𝜼^{}}{}p_𝜼^{}\eta _\nu ^{}𝑑u^{}\rho _𝜼^{}(u^{})\mathrm{tanh}[\gamma u^{}]u\right]\right\}+T\frac{^2}{u^2}\rho _𝜼(u)$$
(47)
Equation (47) is the basis for our further analysis. It can be simplified only if we make additional assumptions on the system’s initial conditions, such as $`\delta `$-distributed or Gaussian distributed $`\rho _𝜼(u)`$ at $`t=0`$ (see below); otherwise it will have to be solved numerically.
Reduction to the Level of Pattern Overlaps. It is clear that (47) is again of the time-dependent Ornstein-Uhlenbeck form, and will thus again have Gaussian solutions as the natural ones:
$$\rho _{t,𝜼}(u)=[2\pi \mathrm{\Sigma }_𝜼^2(t)]^{\frac{1}{2}}e^{\frac{1}{2}[u\overline{u}_𝜼(t)]^2/\mathrm{\Sigma }_𝜼^2(t)}$$
(48)
in which $`\mathrm{\Sigma }_𝜼(t)=[T+(\mathrm{\Sigma }_𝜼^2(0)T)e^{2t}]^{\frac{1}{2}}`$, and with the $`\overline{u}_𝜼(t)`$ evolving in time according to
$$\frac{d}{dt}\overline{u}_𝜼=\underset{𝜼^{}}{}p_𝜼^{}(𝜼𝑨𝜼^{})Dz\mathrm{tanh}[\gamma (\overline{u}_𝜼^{}+\mathrm{\Sigma }_𝜼^{}z)]\overline{u}_𝜼$$
(49)
Our problem has thus been reduced successfully to the study of the $`2^p`$ coupled scalar equations (49). We can also measure the correlation between the firing activities $`s_i(u_i)=\frac{1}{2}[1+\mathrm{tanh}(\gamma u_i)]`$ and the pattern components (similar to the overlaps in the case of binary neurons). If the pattern bits are drawn at random, i.e. $`lim_N\mathrm{}|I_𝜼|/N=p_𝜼=2^p`$ for all $`𝜼`$, we can define a ‘graded response’ equivalent $`m_\mu (𝒖)=2N^1_i\xi _i^\mu s_i(u_i)[1,1]`$ of the pattern overlaps:
$$m_\mu (𝒖)=\frac{2}{N}\underset{i}{}\xi _i^\mu s_i(𝒖)=\frac{1}{N}\underset{i}{}\xi _i^\mu \mathrm{tanh}(\gamma u_i)+𝒪(N^{\frac{1}{2}})$$
$$=\underset{𝜼}{}p_𝜼\eta _\mu 𝑑u\rho _𝜼(u;𝒖)\mathrm{tanh}(\gamma u)+𝒪(N^{\frac{1}{2}})$$
(50)
Full recall of pattern $`\mu `$ implies $`s_i(u_i)=\frac{1}{2}[\xi _i^\mu +1]`$, giving $`m_\mu (𝒖)=1`$. Since the distributions $`\rho _𝜼(u)`$ obey deterministic laws for $`N\mathrm{}`$, the same will be true for the overlaps $`𝒎=(m_1,\mathrm{},m_p)`$. For the Gaussian solutions (49) of (47) we can now proceed to replace the $`2^p`$ macroscopic laws (49), which reduce to $`\frac{d}{dt}\overline{u}_𝜼=𝜼𝑨𝒎\overline{u}_𝜼`$ and give $`\overline{u}_𝜼=\overline{u}_𝜼(0)e^t+𝜼𝑨_0^t𝑑se^{st}𝒎(s)`$, by $`p`$ integral equations in terms of overlaps only:
$$m_\mu (t)=\underset{𝜼}{}p_𝜼\eta _\mu Dz\mathrm{tanh}\left[\gamma \left(\overline{u}_𝜼(0)e^t+𝜼𝑨_0^t𝑑se^{st}𝒎(s)+z\sqrt{T+(\mathrm{\Sigma }_𝜼^2(0)T)e^{2t}}\right)\right]$$
(51)
with $`Dz=(2\pi )^{\frac{1}{2}}e^{\frac{1}{2}z^2}dz`$. Here the sub-lattices only come in via the initial conditions.
Extracting the Physics from the Macroscopic Laws. The equations describing the asymptotic (stationary) state can be written entirely without sub-lattices, by taking the $`t\mathrm{}`$ limit in (51), using $`\overline{u}_𝜼𝜼𝑨𝒎`$, $`\mathrm{\Sigma }_𝜼\sqrt{T}`$, and the familiar notation $`g(𝝃)_𝝃=lim_N\mathrm{}\frac{1}{N}_ig(𝝃_i)=2^p_{𝝃\{1,1\}^p}g(𝝃)`$:
$$m_\mu =\xi _\mu Dz\mathrm{tanh}[\gamma (𝝃𝑨𝒎+z\sqrt{T})]_𝝃\rho _𝜼(u)=[2\pi T]^{\frac{1}{2}}e^{\frac{1}{2}[u𝜼𝑨𝒎]^2/T}$$
(52)
Note the appealing similarity with previous results on networks with binary neurons in equilibrium . For $`T=0`$ the overlap equations (52) become identical to those found for attractor networks with binary neurons and finite $`p`$ (hence our choice to insert an extra factor 2 in defining the synapses), with $`\gamma `$ replacing the inverse noise level $`\beta `$ in the former.
For the simplest non-trivial choice $`A_{\mu \nu }=\delta _{\mu \nu }`$ (i.e. $`J_{ij}=(2/N)_\mu \xi _i^\mu \xi _j^\mu `$, as in the Hopfield model) equation (52) yields the familiar pure and mixture state solutions. For $`T=0`$ we find a continuous phase transition from non-recall to pure states of the form $`m_\mu =m\delta _{\mu \nu }`$ (for some $`\nu `$) at $`\gamma _c=1`$. For $`T>0`$ we have in (52) an additional Gaussian noise, absent in the models with binary neurons. Again the pure states are the first non-trivial solutions to enter the stage. Substituting $`m_\mu =m\delta _{\mu \nu }`$ into (52) gives
$$m=Dz\mathrm{tanh}[\gamma (m+z\sqrt{T})]$$
(53)
Writing (53) as $`m^2=\gamma m_0^m𝑑k[1Dz\mathrm{tanh}^2[\gamma (k+z\sqrt{T})]]\gamma m^2`$, reveals that $`m=0`$ as soon as $`\gamma <1`$. A continuous transition to an $`m>0`$ state occurs when $`\gamma ^1=1Dz\mathrm{tanh}^2[\gamma z\sqrt{T}]`$. A parametrisation of this transition line in the $`(\gamma ,T)`$-plane is given by
$$\gamma ^1(x)=1Dz\mathrm{tanh}^2(zx),T(x)=x^2/\gamma ^2(x),x0$$
(54)
Discontinuous transitions away from $`m=0`$ (for which there is no evidence) would have to be calculated numerically. For $`\gamma =\mathrm{}`$ we get the equation $`m=\mathrm{erf}[m/\sqrt{2T}]`$, giving a continuous transition to $`m>0`$ at $`T_c=2/\pi 0.637`$. Alternatively the latter number can also be found by taking $`lim_x\mathrm{}T(x)`$ in the above parametrisation:
$$T_c(\gamma =\mathrm{})=\underset{x\mathrm{}}{lim}x^2[1Dz\mathrm{tanh}^2(zx)]^2=\underset{x\mathrm{}}{lim}[Dz\frac{d}{dz}\mathrm{tanh}(zx)]^2=[2Dz\delta (z)]^2=2/\pi $$
The resulting picture of the network’s stationary state properties is illustrated in figure 6, which shows the phase diagram and the stationary recall overlaps of the pure states, obtained by numerical calculation and solution of equations (54) and (53).
Let us now turn to dynamics. It follows from (52) that the ‘natural’ initial conditions for $`\overline{u}_𝜼`$ and $`\mathrm{\Sigma }_𝜼`$ are of the form: $`\overline{u}_𝜼(0)=𝜼𝒌_0`$ and $`\mathrm{\Sigma }_𝜼(0)=\mathrm{\Sigma }_0`$ for all $`𝜼`$. Equivalently:
$$t=0:\rho _𝜼(u)=[2\pi \mathrm{\Sigma }_0^2]^{\frac{1}{2}}e^{\frac{1}{2}[u𝜼𝒌_0]^2/\mathrm{\Sigma }_0^2},𝒌_0\mathrm{}^p,\mathrm{\Sigma }_0\mathrm{}$$
These would also be the typical and natural statistics if we were to prepare an initial firing state $`\{s_i\}`$ by hand, via manipulation of the potentials $`\{u_i\}`$. For such initial conditions we can simplify the dynamical equation (51) to
$$m_\mu (t)=\xi _\mu Dz\mathrm{tanh}\left[\gamma \left(𝝃[𝒌_0e^t+𝑨_0^t𝑑se^{st}𝒎(s)]+z\sqrt{T+(\mathrm{\Sigma }_0^2T)e^{2t}}\right)\right]_𝝃$$
(55)
For the special case of the Hopfield synapses, i.e. $`A_{\mu \nu }=\delta _{\mu \nu }`$, it follows from (55) that recall of a given pattern $`\nu `$ is triggered upon choosing $`k_{0,\mu }=k_0\delta _{\mu \nu }`$ (with $`k_0>0`$), since then equation (55) generates $`m_\mu (t)=m(t)\delta _{\mu \nu }`$ at any time, with the amplitude $`m(t)`$ following from
$$m(t)=Dz\mathrm{tanh}\left[\gamma [k_0e^t+_0^t𝑑se^{st}m(s)+z\sqrt{T+(\mathrm{\Sigma }_0^2T)e^{2t}}]\right]$$
(56)
which is the dynamical counterpart of equation (53) (to which indeed it reduces for $`t\mathrm{}`$).
We finally specialize further to the case where our Gaussian initial conditions are not only chosen to trigger recall of a single pattern $`𝝃^\nu `$, but in addition describe uniform membrane potentials within the sub-lattices, i.e. $`k_{0,\mu }=k_0\delta _{\mu \nu }`$ and $`\mathrm{\Sigma }_0=0`$, so $`\rho _𝜼(u)=\delta [uk_0\eta _\nu ]`$. Here we can derive from (56) at $`t=0`$ the identity $`m_0=\mathrm{tanh}[\gamma k_0]`$, which enables us to express $`k_0`$ as $`k_0=(2\gamma )^1\mathrm{log}[(1+m_0)/(1m_0)]`$, and find (56) reducing to
$$m(t)=Dz\mathrm{tanh}[e^t\mathrm{log}[\frac{1+m_0}{1m_0}]^{\frac{1}{2}}+\gamma [_0^tdse^{st}m(s)+z\sqrt{T(1e^{2t})}]]$$
(57)
Solving this equation numerically leads to graphs such as those shown in figure 7 for the choice $`\gamma =4`$ and $`T\{0.25,0.5,0.75\}`$. Compared to the overlap evolution in large networks of binary networks (away from saturation) one immediately observes richer behaviour, e.g. non-monotonicity.
The analysis and results described in this section, which can be done and derived in a similar fashion for other networks with continuous units (such as coupled oscillators), are somewhat difficult to find in research papers. There are two reasons for this. Firstly, non-equilibrium statistical mechanical studies only started being carried out around 1988, and obviously concentrated at first on the (simpler) networks with binary variables. Secondly, due to the absence of detailed balance in networks of graded response networks, the latter appear to have been suspected of consequently having highly complicated dynamics, and analysis terminated with pseudo-equilibrium studies . In retrospect that turns out to have been too pessimistic a view on the power of non-equilibrium statistical mechanics: one finds that dynamical tools can be applied without serious technical problems (although the calculations are somewhat more involved), and again yield interesting and explicit results in the form of phase diagrams and dynamical curves for macroscopic observables, with sensible physical interpretations.
## 4 Correlation- and Response-Functions
We now turn to correlation functions $`C_{ij}(t,t^{})`$ and response functions $`G_{ij}(t,t^{})`$. These will become the language in which the generating functional methods are formulated, which will enable us to solve the dynamics of recurrent networks in the (complex) regime near saturation (we take $`t>t^{}`$):
$$C_{ij}(t,t^{})=\sigma _i(t)\sigma _j(t^{})G_{ij}(t,t^{})=\sigma _i(t)/\theta _j(t^{})$$
(58)
The $`\{\sigma _i\}`$ evolve in time according to equations of the form (1) (binary neurons, sequential updates), (22) (binary neurons, parallel updates) or (32) (continuous neurons). The $`\theta _i`$ represent thresholds and/or external stimuli, which are added to the local fields in the cases (1,22), or added to the deterministic forces in the case of a Fokker-Planck equation (32). We retain $`\theta _i(t)=\theta _i`$, except for a perturbation $`\delta \theta _j(t^{})`$ applied at time $`t^{}`$ in defining the response function. Calculating averages such as (58) requires determining joint probability distributions involving neuron states at different times.
### 4.1 Fluctuation-Dissipation Theorems
Networks of Binary Neurons. For networks of binary neurons with discrete time dynamics of the form $`p_{\mathrm{}+1}(𝝈)=_𝝈^{}W[𝝈;𝝈^{}]p_{\mathrm{}}(𝝈^{})`$, the probability of observing a given ‘path’ $`𝝈(\mathrm{}^{})𝝈(\mathrm{}^{}+1)\mathrm{}𝝈(\mathrm{}1)𝝈(\mathrm{})`$ of successive configurations between step $`\mathrm{}^{}`$ and step $`\mathrm{}`$ is given by the product of the corresponding transition matrix elements (without summation):
$$\mathrm{Prob}[𝝈(\mathrm{}^{}),\mathrm{},𝝈(\mathrm{})]=W[𝝈(\mathrm{});𝝈(\mathrm{}1)]W[𝝈(\mathrm{}1);𝝈(\mathrm{}2)]\mathrm{}W[𝝈(\mathrm{}^{}+1);𝝈(\mathrm{}^{})]p_{\mathrm{}^{}}(𝝈(\mathrm{}^{}))$$
This allows us to write
$$C_{ij}(\mathrm{},\mathrm{}^{})=\underset{𝝈(\mathrm{}^{})}{}\mathrm{}\underset{𝝈(\mathrm{})}{}\mathrm{Prob}[𝝈(\mathrm{}^{}),\mathrm{},𝝈(\mathrm{})]\sigma _i(\mathrm{})\sigma _j(\mathrm{}^{})=\underset{𝝈𝝈^{}}{}\sigma _i\sigma _j^{}W^{\mathrm{}\mathrm{}^{}}[𝝈;𝝈^{}]p_{\mathrm{}^{}}(𝝈^{})$$
(59)
$$G_{ij}(\mathrm{},\mathrm{}^{})=\underset{𝝈𝝈^{}𝝈^{\prime \prime }}{}\sigma _iW^\mathrm{}\mathrm{}^{}1[𝝈;𝝈^{\prime \prime }]\left[\frac{}{\theta _j}W[𝝈^{\prime \prime };𝝈^{}]\right]p_{\mathrm{}^{}}(𝝈^{})$$
(60)
From (59) and (60) it follows that both $`C_{ij}(\mathrm{},\mathrm{}^{})`$ and $`G_{ij}(\mathrm{},\mathrm{}^{})`$ will in the stationary state, i.e. upon substituting $`p_{\mathrm{}^{}}(𝝈^{})=p_{\mathrm{}}(𝝈^{})`$, only depend on $`\mathrm{}\mathrm{}^{}`$: $`C_{ij}(\mathrm{},\mathrm{}^{})C_{ij}(\mathrm{}\mathrm{}^{})`$ and $`G_{ij}(\mathrm{},\mathrm{}^{})G_{ij}(\mathrm{}\mathrm{}^{})`$. For this we do not require detailed balance. Detailed balance, however, leads to a simple relation between the response function $`G_{ij}(\tau )`$ and the temporal derivative of the correlation function $`C_{ij}(\tau )`$.
We now turn to equilibrium systems, i.e. networks with symmetric synapses (and with all $`J_{ii}=0`$ in the case of sequential dynamics). We calculate the derivative of the transition matrix that occurs in (60) by differentiating the equilibrium condition $`p_{\mathrm{eq}}(𝝈)=_𝝈^{}W[𝝈;𝝈^{}]p_{\mathrm{eq}}(𝝈^{})`$ with respect to external fields:
$$\frac{}{\theta _j}p_{\mathrm{eq}}(𝝈)=\underset{𝝈^{}}{}\left\{\frac{W[𝝈;𝝈^{}]}{\theta _j}p_{\mathrm{eq}}(𝝈^{})+W[𝝈;𝝈^{}]\frac{}{\theta _j}p_{\mathrm{eq}}(𝝈^{})\right\}$$
Detailed balance implies $`p_{\mathrm{eq}}(𝝈)=Z^1e^{\beta H(𝝈)}`$ (in the parallel case we simply substitute the appropriate Hamiltonian $`H\stackrel{~}{H}`$), giving $`p_{\mathrm{eq}}(𝝈)/\theta _j=[Z^1Z/\theta _j+\beta H(𝝈)/\theta _j]p_{\mathrm{eq}}(𝝈)`$, so that
$$\underset{𝝈^{}}{}\frac{W[𝝈;𝝈^{}]}{\theta _j}p_{\mathrm{eq}}(𝝈^{})=\beta \left\{\underset{𝝈^{}}{}W[𝝈;𝝈^{}]\frac{H(𝝈^{})}{\theta _j}p_{\mathrm{eq}}(𝝈^{})\frac{H(𝝈)}{\theta _j}p_{\mathrm{eq}}(𝝈)\right\}$$
(the term containing $`Z`$ drops out). We now obtain for the response function (60) in equilibrium :
$$G_{ij}(\mathrm{})=\beta \underset{𝝈𝝈^{}}{}\sigma _iW^\mathrm{}1[𝝈;𝝈^{}]\left\{\underset{𝝈^{\prime \prime }}{}W[𝝈^{};𝝈^{\prime \prime }]\frac{H(𝝈^{\prime \prime })}{\theta _j}p_{\mathrm{eq}}(𝝈^{\prime \prime })\frac{H(𝝈^{})}{\theta _j}p_{\mathrm{eq}}(𝝈^{})\right\}$$
(61)
The structure of (61) is similar to what follows upon calculating the evolution of the equilibrium correlation function (59) in a single iteration step:
$$C_{ij}(\mathrm{})C_{ij}(\mathrm{}1)=\underset{𝝈𝝈^{}}{}\sigma _iW^\mathrm{}1[𝝈;𝝈^{}]\left\{\underset{𝝈^{\prime \prime }}{}W[𝝈^{};𝝈^{\prime \prime }]\sigma _j^{\prime \prime }p_{\mathrm{eq}}(𝝈^{\prime \prime })\sigma _j^{}p_{\mathrm{eq}}(𝝈^{})\right\}$$
(62)
Finally we calculate the relevant derivatives of the two Hamiltonians $`H(𝝈)=_{i<j}J_{ij}\sigma _i\sigma _j+_i\theta _i\sigma _i`$ and $`\stackrel{~}{H}(𝝈)=_i\theta _i\sigma _i\beta ^1_i\mathrm{log}2\mathrm{cosh}[\beta h_i(𝝈)]`$ (with $`h_i(𝝈)=_jJ_{ij}\sigma _j+\theta _i`$), see :
$$H(𝝈)/\theta _j=\sigma _j\stackrel{~}{H}(𝝈)/\theta _j=\sigma _j\mathrm{tanh}[\beta h_j(𝝈)]$$
For sequential dynamics we hereby arrive directly at a fluctuation-dissipation theorem. For parallel dynamics we need one more identity (which follows from the definition of the transition matrix in (22) and the detailed balance property) to transform the $`tanh`$ occurring in the derivative of $`\stackrel{~}{H}`$:
$$\mathrm{tanh}[\beta h_j(𝝈^{})]p_{\mathrm{eq}}(𝝈^{})=\underset{𝝈^{\prime \prime }}{}\sigma _j^{\prime \prime }W[𝝈^{\prime \prime };𝝈^{}]p_{\mathrm{eq}}(𝝈^{})=\underset{𝝈^{\prime \prime }}{}W[𝝈^{};𝝈^{\prime \prime }]\sigma _j^{\prime \prime }p_{\mathrm{eq}}(𝝈^{\prime \prime })$$
For parallel dynamics $`\mathrm{}`$ and $`\mathrm{}^{}`$ are the real time labels $`t`$ and $`t^{}`$, and we obtain, with $`\tau =tt^{}`$:
$$\mathrm{𝐵𝑖𝑛𝑎𝑟𝑦}\&\mathrm{𝑃𝑎𝑟𝑎𝑙𝑙𝑒𝑙}:G_{ij}(\tau >0)=\beta [C_{ij}(\tau +1)C_{ij}(\tau 1)],G_{ij}(\tau 0)=0$$
(63)
For the continuous-time version (1) of sequential dynamics the time $`t`$ is defined as $`t=\mathrm{}/N`$, and the difference equation (62) becomes a differential equation. For perturbations at time $`t^{}`$ in the definition of the response function (60) to retain a non-vanishing effect at (re-scaled) time $`t`$ in the limit $`N\mathrm{}`$, they will have to be re-scaled as well: $`\delta \theta _j(t^{})N\delta \theta _j(t^{})`$. As a result:
$$\mathrm{𝐵𝑖𝑛𝑎𝑟𝑦}\&\mathrm{𝑆𝑒𝑞𝑢𝑒𝑛𝑡𝑖𝑎𝑙}:G_{ij}(\tau )=\beta \theta (\tau )\frac{d}{d\tau }C_{ij}(\tau )$$
(64)
The need to re-scale perturbations in making the transition from discrete to continuous times has the same origin as the need to re-scale the random forces in the derivation of the continuous-time Langevin equation from a discrete-time process. Going from ordinary derivatives to functional derivatives (which is what happens in the continuous-time limit), implies replacing Kronecker delta’s $`\delta _{t,t^{}}`$ by Dirac delta-functions according to $`\delta _{t,t^{}}\mathrm{\Delta }\delta (tt^{})`$, where $`\mathrm{\Delta }`$ is the average duration of an iteration step. Equations (63) and (64) are examples of so-called fluctuation-dissipation theorems (FDT).
Networks with Continuous Neurons. For systems described by a Fokker-Planck equation (32) the simplest way to calculate correlation- and response-functions is by first returning to the underlying discrete-time system and leaving the continuous time limit $`\mathrm{\Delta }0`$ until the end. In we saw that for small but finite time-steps $`\mathrm{\Delta }`$ the underlying discrete-time process is described by
$$t=\mathrm{}\mathrm{\Delta },p_{\mathrm{}\mathrm{\Delta }+\mathrm{\Delta }}(𝝈)=[1+\mathrm{\Delta }_𝝈+𝒪(\mathrm{\Delta }^{\frac{3}{2}})]p_\mathrm{}\mathrm{\Delta }(𝝈)$$
with $`\mathrm{}=0,1,2,\mathrm{}`$ and with the differential operator
$$_𝝈=\underset{i}{}\frac{}{\sigma _i}[f_i(𝝈)T\frac{}{\sigma _i}]$$
(65)
From this it follows that the conditional probability density $`p_\mathrm{}\mathrm{\Delta }(𝝈|𝝈^{},\mathrm{}^{}\mathrm{\Delta })`$ for finding state $`\sigma `$ at time $`\mathrm{}\mathrm{\Delta }`$, given the system was in state $`𝝈^{}`$ at time $`\mathrm{}^{}\mathrm{\Delta }`$, must be
$$p_\mathrm{}\mathrm{\Delta }(𝝈|𝝈^{},\mathrm{}^{}\mathrm{\Delta })=[1+\mathrm{\Delta }_𝝈+𝒪(\mathrm{\Delta }^{\frac{3}{2}})]^{\mathrm{}\mathrm{}^{}}\delta [𝝈𝝈^{}]$$
(66)
Equation (66) will be our main building block. Firstly, we will calculate the correlations:
$$C_{ij}(\mathrm{}\mathrm{\Delta },\mathrm{}^{}\mathrm{\Delta })=\sigma _i(\mathrm{}\mathrm{\Delta })\sigma _j(\mathrm{}^{}\mathrm{\Delta })=𝑑𝝈𝑑𝝈^{}\sigma _i\sigma _j^{}p_\mathrm{}\mathrm{\Delta }(𝝈|𝝈^{},\mathrm{}^{}\mathrm{\Delta })p_\mathrm{}^{}\mathrm{\Delta }(𝝈^{})$$
$$=𝑑𝝈\sigma _i[1+\mathrm{\Delta }_𝝈+𝒪(\mathrm{\Delta }^{\frac{3}{2}})]^{\mathrm{}\mathrm{}^{}}𝑑𝝈^{}\sigma _j^{}\delta [𝝈𝝈^{}]p_\mathrm{}^{}\mathrm{\Delta }(𝝈^{})$$
$$=𝑑𝝈\sigma _i[1+\mathrm{\Delta }_𝝈+𝒪(\mathrm{\Delta }^{\frac{3}{2}})]^{\mathrm{}\mathrm{}^{}}\left[\sigma _jp_\mathrm{}^{}\mathrm{\Delta }(𝝈)\right]$$
At this stage we can take the limits $`\mathrm{\Delta }0`$ and $`\mathrm{},\mathrm{}^{}\mathrm{}`$, with $`t=\mathrm{}\mathrm{\Delta }`$ and $`t^{}=\mathrm{}^{}\mathrm{\Delta }`$ finite, using $`lim_{\mathrm{\Delta }0}[1+\mathrm{\Delta }A]^{k/\mathrm{\Delta }}=e^{kA}`$:
$$C_{ij}(t,t^{})=𝑑𝝈\sigma _ie^{(tt^{})_𝝈}\left[\sigma _jp_t^{}(𝝈)\right]$$
(67)
Next we turn to the response function. A perturbation applied at time $`t^{}=\mathrm{}^{}\mathrm{\Delta }`$ to the Langevin forces $`f_i(𝝈)`$ comes in at the transition $`𝝈(\mathrm{}^{}\mathrm{\Delta })𝝈(\mathrm{}^{}\mathrm{\Delta }+\mathrm{\Delta })`$. As with sequential dynamics binary networks, the perturbation is re-scaled with the step size $`\mathrm{\Delta }`$ to retain significance as $`\mathrm{\Delta }0`$:
$$G_{ij}(\mathrm{}\mathrm{\Delta },\mathrm{}^{}\mathrm{\Delta })=\frac{\sigma _i(\mathrm{}\mathrm{\Delta })}{\mathrm{\Delta }\theta _j(\mathrm{}^{}\mathrm{\Delta })}=\frac{}{\mathrm{\Delta }\theta _j(\mathrm{}^{}\mathrm{\Delta })}𝑑𝝈𝑑𝝈^{}\sigma _ip_\mathrm{}\mathrm{\Delta }(𝝈|𝝈^{},\mathrm{}^{}\mathrm{\Delta })p_\mathrm{}^{}\mathrm{\Delta }(𝝈^{})$$
$$=𝑑𝝈𝑑𝝈^{}𝑑𝝈^{\prime \prime }\sigma _ip_\mathrm{}\mathrm{\Delta }(𝝈|𝝈^{\prime \prime },\mathrm{}^{}\mathrm{\Delta }+\mathrm{\Delta })\left[\frac{p_{\mathrm{}^{\prime \prime }\mathrm{\Delta }+\mathrm{\Delta }}(𝝈|𝝈^{},\mathrm{}^{}\mathrm{\Delta })}{\mathrm{\Delta }\theta _j}\right]p_\mathrm{}^{}\mathrm{\Delta }(𝝈^{})$$
$$=𝑑𝝈𝑑𝝈^{}𝑑𝝈^{\prime \prime }\sigma _i[1+\mathrm{\Delta }_𝝈+𝒪(\mathrm{\Delta }^{\frac{3}{2}})]^\mathrm{}\mathrm{}^{}1\delta [𝝈𝝈^{\prime \prime }]\left[\frac{1}{\mathrm{\Delta }}\frac{}{\theta _j}[1+\mathrm{\Delta }_{𝝈^{\prime \prime }}+𝒪(\mathrm{\Delta }^{\frac{3}{2}})]\delta [𝝈^{\prime \prime }𝝈^{}]\right]p_\mathrm{}^{}\mathrm{\Delta }(𝝈^{})$$
$$=𝑑𝝈𝑑𝝈^{}𝑑𝝈^{\prime \prime }\sigma _i[1+\mathrm{\Delta }_𝝈+𝒪(\mathrm{\Delta }^{\frac{3}{2}})]^\mathrm{}\mathrm{}^{}1\delta [𝝈𝝈^{\prime \prime }]\delta [𝝈^{\prime \prime }𝝈^{}][\frac{}{\sigma _j^{}}+𝒪(\mathrm{\Delta }^{\frac{1}{2}})]p_\mathrm{}^{}\mathrm{\Delta }(𝝈^{})$$
$$=𝑑𝝈\sigma _i[1+\mathrm{\Delta }_𝝈+𝒪(\mathrm{\Delta }^{\frac{3}{2}})]^\mathrm{}\mathrm{}^{}1[\frac{}{\sigma _j}+𝒪(\mathrm{\Delta }^{\frac{1}{2}})]p_\mathrm{}^{}\mathrm{\Delta }(𝝈)$$
We take the limits $`\mathrm{\Delta }0`$ and $`\mathrm{},\mathrm{}^{}\mathrm{}`$, with $`t=\mathrm{}\mathrm{\Delta }`$ and $`t^{}=\mathrm{}^{}\mathrm{\Delta }`$ finite:
$$G_{ij}(t,t^{})=𝑑𝝈\sigma _ie^{(tt^{})_𝝈}\frac{}{\sigma _j}p_t^{}(𝝈)$$
(68)
Equations (67) and (68) apply to arbitrary systems described by Fokker-Planck equations. In the case of conservative forces, i.e. $`f_i(𝝈)=H(𝝈)/\sigma _i`$, and when the system is in an equilibrium state at time $`t^{}`$ so that $`C_{ij}(t,t^{})=C_{ij}(tt^{})`$ and $`G_{ij}(t,t^{})=G_{ij}(tt^{})`$, we can take a further step using $`p_t^{}(𝝈)=p_{\mathrm{eq}}(𝝈)=Z^1e^{\beta H(𝝈)}`$. In that case, taking the time derivative of expression (67) gives
$$\frac{}{\tau }C_{ij}(\tau )=𝑑𝝈\sigma _ie^{\tau _𝝈}_𝝈\left[\sigma _jp_{\mathrm{eq}}(𝝈)\right]$$
Working out the key term in this expression gives
$$_𝝈[\sigma _jp_{\mathrm{eq}}(𝝈)]=\underset{i}{}\frac{}{\sigma _i}[f_i(𝝈)T\frac{}{\sigma _i}][\sigma _jp_{\mathrm{eq}}(𝝈)]=T\frac{}{\sigma _j}p_{\mathrm{eq}}(𝝈)\underset{i}{}\frac{}{\sigma _i}[\sigma _jJ_i(𝝈)]$$
with the components of the probability current density $`J_i(𝝈)=[f_i(𝝈)T\frac{}{\sigma _i}]p_{\mathrm{eq}}(𝝈)`$. In equilibrium, however, the current is zero by definition, so only the first term in the above expression survives. Insertion into our previous equation for $`C_{ij}(\tau )/\tau `$, and comparison with (68) leads to the FDT for continuous systems:
$$\mathrm{𝐶𝑜𝑛𝑡𝑖𝑛𝑢𝑜𝑢𝑠}:G_{ij}(\tau )=\beta \theta (\tau )\frac{d}{d\tau }C_{ij}(\tau )$$
(69)
We will now calculate the correlation and response functions explicitly, and verify the validity or otherwise of the FDT relations, for attractor networks away from saturation.
### 4.2 Example: Simple Attractor Networks with Binary Neurons
Correlation- and Response Functions for Sequential Dynamics. We will consider the continuous time version (1) of the sequential dynamics, with the local fields $`h_i(𝝈)=_jJ_{ij}\sigma _j+\theta _i`$, and the separable interaction matrix (18). We already solved the dynamics of this model for the case with zero external fields and away from saturation (i.e. $`p\sqrt{N}`$). Having non-zero, or even time-dependent, external fields does not affect the calculation much; one adds the external fields to the internal ones and finds the macroscopic laws (20) for the overlaps with the stored patterns being replaced by
$$\frac{d}{dt}𝒎(t)=\underset{N\mathrm{}}{lim}\frac{1}{N}\underset{i}{}𝝃_i\mathrm{tanh}\left[\beta 𝝃_i𝑨𝒎(t)+\theta _i(t)\right]𝒎(t)$$
(70)
Fluctuations in the local fields are of vanishing order in $`N`$ (since the fluctuations in $`𝒎`$ are), so that one can easily derive from the master equation (1) the following expressions for spin averages:
$$\frac{d}{dt}\sigma _i(t)=\mathrm{tanh}\beta [𝝃_i𝑨𝒎(t)+\theta _i(t)]\sigma _i(t)$$
(71)
$$ij:\frac{d}{dt}\sigma _i(t)\sigma _j(t)=\mathrm{tanh}\beta [𝝃_i𝑨𝒎(t)+\theta _i(t)]\sigma _j(t)+\mathrm{tanh}\beta [𝝃_j𝑨𝒎(t)+\theta _j(t)]\sigma _i(t)2\sigma _i(t)\sigma _j(t)$$
(72)
Correlations at different times are calculated by applying (71) to situations where the microscopic state at time $`t^{}`$ is known exactly, i.e. where $`p_t^{}(𝝈)=\delta _{𝝈,𝝈^{}}`$ for some $`𝝈^{}`$:
$$\sigma _i(t)|_{𝝈(t^{})=𝝈^{}}=\sigma _i^{}e^{(tt^{})}+_t^{}^t𝑑se^{st}\mathrm{tanh}\beta [𝝃_i𝑨𝒎(s;𝝈^{},t^{})+\theta _i(s)]$$
(73)
with $`𝒎(s;𝝈^{},t^{})`$ denoting the solution of (70) following initial condition $`𝒎(t^{})=\frac{1}{N}_i\sigma _i^{}𝝃_i`$. If we multiply both sides of (73) by $`\sigma _j^{}`$ and average over all possible states $`𝝈^{}`$ at time $`t^{}`$ we obtain in leading order in $`N`$:
$$\sigma _i(t)\sigma _j(t^{})=\sigma _i(t^{})\sigma _j(t^{})e^{(tt^{})}+_t^{}^t𝑑se^{st}\mathrm{tanh}\beta [𝝃_i𝑨𝒎(s;𝝈(t^{}),t^{})+\theta _i(s)]\sigma _j(t^{})$$
Because of the existence of deterministic laws for the overlaps $`𝒎`$ in the $`N\mathrm{}`$ limit, we know with probability one that during the stochastic process the actual value $`𝒎(𝝈(t^{}))`$ must be given by the solution of (70), evaluated at time $`t^{}`$. As a result we obtain, with $`C_{ij}(t,t^{})=\sigma _i(t)\sigma _j(t^{})`$:
$$C_{ij}(t,t^{})=C_{ij}(t^{},t^{})e^{(tt^{})}+_t^{}^t𝑑se^{st}\mathrm{tanh}\beta [𝝃_i𝑨𝒎(s)+\theta _i(s)]\sigma _j(t^{})$$
(74)
Similarly we obtain from the solution of (71) an equation for the leading order in $`N`$ of the response functions, by derivation with respect to external fields:
$$\frac{\sigma _i(t)}{\theta _j(t^{})}=\beta \theta (tt^{})_{\mathrm{}}^t𝑑se^{st}\left[1\mathrm{tanh}^2\beta [𝝃_i𝑨𝒎(s)+\theta _i(s)]\right]\left[\frac{1}{N}\underset{k}{}(𝝃_i𝑨𝝃_k)\frac{\sigma _k(s)}{\theta _j(t^{})}+\delta _{ij}\delta (st^{})\right]$$
or
$$G_{ij}(t,t^{})=\beta \delta _{ij}\theta (tt^{})e^{(tt^{})}\left[1\mathrm{tanh}^2\beta [𝝃_i𝑨𝒎(t^{})+\theta _i(t^{})]\right]$$
$$+\beta \theta (tt^{})_t^{}^t𝑑se^{st}\left[1\mathrm{tanh}^2\beta [𝝃_i𝑨𝒎(s)+\theta _i(s)]\right]\frac{1}{N}\underset{k}{}(𝝃_i𝑨𝝃_k)G_{kj}(s,t^{})$$
(75)
For $`t=t^{}`$ we retain in leading order in $`N`$ only the instantaneous single-site contribution
$$\underset{t^{}t}{lim}G_{ij}(t,t^{})=\beta \delta _{ij}\left[1\mathrm{tanh}^2\beta [𝝃_i𝑨𝒎(t)+\theta _i(t)]\right]$$
(76)
This leads to the following ansatz for the scaling with $`N`$ of the $`G_{ij}(t,t^{})`$, which can be shown to be correct by insertion into (75), in combination with the correctness at $`t=t^{}`$ following from (76):
$$i=j:G_{ii}(t,t^{})=𝒪(1),ij:G_{ij}(t,t^{})=𝒪(N^1)$$
Note that this implies $`\frac{1}{N}_k(𝝃_i𝑨𝝃_k)G_{kj}(s,t^{})=𝒪(\frac{1}{N})`$. In leading order in $`N`$ we now find
$$G_{ij}(t,t^{})=\beta \delta _{ij}\theta (tt^{})e^{(tt^{})}\left[1\mathrm{tanh}^2\beta [𝝃_i𝑨𝒎(t^{})+\theta _i(t^{})]\right]$$
(77)
For those cases where the macroscopic laws (70) describe evolution to a stationary state $`𝒎`$, obviously requiring stationary external fields $`\theta _i(t)=\theta _i`$, we can take the limit $`t\mathrm{}`$, with $`tt^{}=\tau `$ fixed, in the two results (74,77). Using the $`t\mathrm{}`$ limits of (71,72) we subsequently find time translation invariant expressions: $`lim_t\mathrm{}C_{ij}(t,t\tau )=C_{ij}(\tau )`$ and $`lim_t\mathrm{}G_{ij}(t,t\tau )=G_{ij}(\tau )`$, with in leading order in $`N`$
$$C_{ij}(\tau )=\mathrm{tanh}\beta [𝝃_i𝑨𝒎+\theta _i]\mathrm{tanh}\beta [𝝃_jA𝒎+\theta _j]+\delta _{ij}e^\tau \left[1\mathrm{tanh}^2\beta [𝝃_i𝑨𝒎+\theta _i]\right]$$
(78)
$$G_{ij}(\tau )=\beta \delta _{ij}\theta (\tau )e^\tau \left[1\mathrm{tanh}^2\beta [𝝃_i𝑨𝒎+\theta _i]\right]$$
(79)
for which indeed the Fluctuation-Dissipation Theorem (64) holds: $`G_{ij}(\tau )=\beta \theta (\tau )\frac{d}{d\tau }C_{ij}(\tau )`$.
Correlation- and Response Functions for Parallel Dynamics. We now turn to the parallel dynamical rules (22), with the local fields $`h_i(𝝈)=_jJ_{ij}\sigma _j+\theta _i`$, and the interaction matrix (18). As before, having time-dependent external fields amounts simply to adding these fields to the internal ones, and the dynamic laws (31) are found to be replaced by
$$𝒎(t+1)=\underset{N\mathrm{}}{lim}\frac{1}{N}\underset{i}{}𝝃_i\mathrm{tanh}\left[\beta 𝝃_i𝑨𝒎(t)+\theta _i(t)\right]$$
(80)
Fluctuations in the local fields are again of vanishing order in $`N`$, and the parallel dynamics versions of equations (71,72), to be derived from (22), are found to be
$$\sigma _i(t+1)=\mathrm{tanh}\beta [𝝃_iA𝒎(t)+\theta _i(t)]$$
(81)
$$ij:\sigma _i(t+1)\sigma _j(t+1)=\mathrm{tanh}\beta [𝝃_i𝑨𝒎(t)+\theta _i(t)]\mathrm{tanh}\beta [𝝃_j𝑨𝒎(t)+\theta _j(t)]$$
(82)
With $`𝒎(t;𝝈^{},t^{})`$ denoting the solution of the map (80) following initial condition $`𝒎(t^{})=\frac{1}{N}_i\sigma _i^{}𝝃_i`$, we immediately obtain from equations (81,82) the correlation functions:
$$C_{ij}(t,t)=\delta _{ij}+[1\delta _{ij}]\mathrm{tanh}\beta [𝝃_i𝑨𝒎(t1)+\theta _i(t1)]\mathrm{tanh}\beta [𝝃_j𝑨𝒎(t1)+\theta _j(t1)]$$
(83)
$$t>t^{}:C_{ij}(t,t^{})=\mathrm{tanh}\beta [𝝃_i𝑨𝒎(t1;𝝈(t^{}),t^{})+\theta _i(t1)]\sigma _j(t^{})$$
$$=\mathrm{tanh}\beta [𝝃_i𝑨𝒎(t1)+\theta _i(t1)]\mathrm{tanh}\beta [𝝃_j𝑨𝒎(t^{}1)+\theta _j(t^{}1)]$$
(84)
From (81) also follow equations determining the leading order in $`N`$ of the response functions $`G_{ij}(t,t^{})`$, by derivation with respect to the external fields $`\theta _j(t^{})`$:
$$\begin{array}{ccc}t^{}>t1:\hfill & & G_{ij}(t,t^{})=0\hfill \\ t^{}=t1:\hfill & & G_{ij}(t,t^{})=\beta \delta _{ij}\left[1\mathrm{tanh}^2\beta [𝝃_i𝑨𝒎(t1)+\theta _i(t1)]\right]\hfill \\ t^{}<t1:\hfill & & G_{ij}(t,t^{})=\beta \left[1\mathrm{tanh}^2\beta [𝝃_i𝑨𝒎(t1)+\theta _i(t1)]\right]\frac{1}{N}_k(𝝃_i𝑨𝝃_k)G_{kj}(t1,t^{})\hfill \end{array}$$
(85)
It now follows iteratively that all off-diagonal elements must be of vanishing order in $`N`$: $`G_{ij}(t,t1)=\delta _{ij}G_{ii}(t,t1)G_{ij}(t,t2)=\delta _{ij}G_{ii}(t,t2)\mathrm{}`$, so that in leading order
$$G_{ij}(t,t^{})=\beta \delta _{ij}\delta _{t,t^{}+1}\left[1\mathrm{tanh}^2\beta [𝝃_i𝑨𝒎(t^{})+\theta _i(t^{})]\right]$$
(86)
For those cases where the macroscopic laws (80) describe evolution to a stationary state $`𝒎`$, with stationary external fields, we can take the limit $`t\mathrm{}`$, with $`tt^{}=\tau `$ fixed, in (83,84,86). We find time translation invariant expressions: $`lim_t\mathrm{}C_{ij}(t,t\tau )=C_{ij}(\tau )`$ and $`lim_t\mathrm{}G_{ij}(t,t\tau )=G_{ij}(\tau )`$, with in leading order in $`N`$:
$$C_{ij}(\tau )=\mathrm{tanh}\beta [𝝃_i𝑨𝒎+\theta _i]\mathrm{tanh}\beta [𝝃_j𝑨𝒎+\theta _j]+\delta _{ij}\delta _{\tau ,0}\left[1\mathrm{tanh}^2\beta [𝝃_i𝑨𝒎+\theta _i]\right]$$
(87)
$$G_{ij}(\tau )=\beta \delta _{ij}\delta _{\tau ,1}\left[1\mathrm{tanh}^2\beta [𝝃_i𝑨𝒎+\theta _i]\right]$$
(88)
obeying the Fluctuation-Dissipation Theorem (63): $`G_{ij}(\tau >0)=\beta [C_{ij}(\tau +1)C_{ij}(\tau 1)]`$.
### 4.3 Example: Graded Response Neurons with Uniform Synapses
Let us finally find out how to calculate correlation and response function for the simple network (38) of graded response neurons, with (possibly time-dependent) external forces $`\theta _i(t)`$, and with uniform synapses $`J_{ij}=J/N`$:
$$\frac{d}{dt}u_i(t)=\frac{J}{N}\underset{j}{}g[\gamma u_j(t)]u_i(t)+\theta _i(t)+\eta _i(t)$$
(89)
For a given realisation of the external forces and the Gaussian noise variables $`\{\eta _i(t)\}`$ we can formally integrate (89) and find
$$u_i(t)=u_i(0)e^t+_0^t𝑑se^{st}\left[J𝑑u\rho (u;𝒖(s))g[\gamma u]+\theta _i(s)+\eta _i(s)\right]$$
(90)
with the distribution of membrane potentials $`\rho (u;𝒖)=N^1_i\delta [uu_i]`$. The correlation function $`C_{ij}(t,t^{})=u_i(t)u_j(t^{})`$ immediately follows from (90). Without loss of generality we can define $`tt^{}`$. For absent external forces (which were only needed in order to define the response function), and upon using $`\eta _i(s)=0`$ and $`\eta _i(s)\eta _j(s^{})=2T\delta _{ij}\delta (ss^{})`$, we arrive at
$$C_{ij}(t,t^{})=T\delta _{ij}(e^{t^{}t}e^{t^{}t})$$
$$+\left[u_i(0)e^t+J𝑑ug[\gamma u]_0^t𝑑se^{st}\rho (u;𝒖(s))\right]\left[u_j(0)e^t^{}+J𝑑ug[\gamma u]_0^t^{}𝑑s^{}e^{s^{}t^{}}\rho (u;𝒖(s^{}))\right]$$
For $`N\mathrm{}`$, however, we know the distribution of potentials to evolve deterministically: $`\rho (u;𝒖(s))\rho _s(u)`$ where $`\rho _s(u)`$ is the solution of (41). This allows us to simplify the above expression to
$$N\mathrm{}:C_{ij}(t,t^{})=T\delta _{ij}(e^{t^{}t}e^{t^{}t})$$
$$+\left[u_i(0)e^t+J𝑑ug[\gamma u]_0^t𝑑se^{st}\rho _s(u)\right]\left[u_j(0)e^t^{}+J𝑑ug[\gamma u]_0^t^{}𝑑s^{}e^{s^{}t^{}}\rho _s^{}(u)\right]$$
(91)
Next we turn to the response function $`G_{ij}(t,t^{})=\delta u_i(t)/\delta \xi _j(t^{})`$ (its definition involves functional rather than scalar differentiation, since time is continuous). After this differentiation the forces $`\{\theta _i(s)\}`$ can be put to zero. Functional differentiation of (90), followed by averaging, then leads us to
$$G_{ij}(t,t^{})=\theta (tt^{})\delta _{ij}e^{t^{}t}J𝑑ug[\gamma u]\frac{}{u}_0^t𝑑se^{st}\frac{1}{N}\underset{k}{}\underset{𝜽\mathrm{𝟎}}{lim}\delta [uu_k(s)]\frac{\delta u_k(s)}{\delta \theta _j(t^{})}$$
In view of (90) we make the self-consistent ansatz $`\delta u_k(s)/\delta \xi _j(s^{})=𝒪(N^1)`$ for $`kj`$. This produces
$$N\mathrm{}:G_{ij}(t,t^{})=\theta (tt^{})\delta _{ij}e^{t^{}t}$$
(92)
Since equation (41) evolves towards a stationary state, we can also take the limit $`t\mathrm{}`$, with $`tt^{}=\tau `$ fixed, in (91). Assuming non-pathological decay of the distribution of potentials allows us to put $`lim_t\mathrm{}_0^t𝑑se^{st}\rho _s(u)=\rho (u)`$ (the stationary solution of (41)), with which we find not only (92) but also (91) reducing to time translation invariant expressions for $`N\mathrm{}`$, $`lim_t\mathrm{}C_{ij}(t,t\tau )=C_{ij}(\tau )`$ and $`lim_t\mathrm{}G_{ij}(t,t\tau )=G_{ij}(\tau )`$, in which
$$C_{ij}(\tau )=T\delta _{ij}e^\tau +J^2\left\{𝑑u\rho (u)g[\gamma u]\right\}^2G_{ij}(\tau )=\theta (\tau )\delta _{ij}e^\tau $$
(93)
Clearly the leading orders in $`N`$ of these two functions obey the fluctuation-dissipation theorem (69): $`G_{ij}(\tau )=\beta \theta (\tau )\frac{d}{d\tau }C_{ij}(\tau )`$. As with the binary neuron attractor networks for which we calculated the correlation and response functions earlier, the impact of detailed balance violation (occurring when $`A_{\mu \nu }A_{\nu \mu }`$ in networks with binary neurons and synapses (18), and in all networks with graded response neurons ) on the validity of the fluctuation-dissipation theorems, vanishes for $`N\mathrm{}`$, provided our networks are relatively simple and evolve to a stationary state in terms of the macroscopic observables (the latter need not necessarily happen, see e.g. figures 1 and 4). Detailed balance violation, however, would be noticed in the finite size effects .
## 5 Dynamics in the Complex Regime
The approach we followed so far to derive closed macroscopic laws from the microscopic equations fails when the number of attractors is no longer small compared to the number $`N`$ of microscopic neuronal variables. In statics we have seen that, at the work floor level, the fingerprint of complexity is the need to use replica theory, rather than the relatively simple and straightforward methods based on (or equivalent to) calculating the density of states for given realisations of the macroscopic observables. This is caused by the presence of a number of ‘disorder’ variables per degree of freedom which is proportional to $`N`$, over which we are forced to average the macroscopic laws. One finds that in dynamics this situation is reflected in the inability to find an exact set of closed equations for a finite number of observables (or densities). We will see that the natural dynamical counterpart of equilibrium replica theory is generating functional analysis.
### 5.1 Overview of Methods and Theories
Let us return to the simplest setting in which to study the problem: single pattern recall in an attractor neural network with $`N`$ binary neurons and $`p=\alpha N`$ stored patterns in the non-trivial regime, where $`\alpha >0`$. We choose parallel dynamics, i.e. (22), with Hebbian-type synapses of the form (18) with $`A_{\mu \nu }=\delta _{\mu \nu }`$, i.e. $`J_{ij}=N^1_\mu ^p\xi _i^\mu \xi _j^\mu `$, giving us the parallel dynamics version of the Hopfield model . Our interest is in the recall overlap $`m(𝝈)=N^1_i\sigma _i\xi _i^1`$ between system state and pattern one. We saw in that for $`N\mathrm{}`$ the fluctuations in the values of the recall overlap $`m`$ will vanish, and that for initial states where all $`\sigma _i(0)`$ are drawn independently the overlap $`m`$ will obey
$$m(t+1)=𝑑zP_t(z)\mathrm{tanh}[\beta (m(t)+z)]P_t(z)=\underset{N\mathrm{}}{lim}\frac{1}{N}\underset{i}{}\delta [z\frac{1}{N}\underset{\mu >1}{}\xi _i^1\xi _i^\mu \underset{ji}{}\xi _j^\mu \sigma _j(t)]$$
(94)
and that all complications in a dynamical analysis of the $`\alpha >0`$ regime are concentrated in the calculation of the distribution $`P_t(z)`$ of the (generally non-trivial) interference noise.
Gaussian Approximations. As a simple approximation one could just assume that the $`\sigma _i`$ remain uncorrelated at all times, i.e. $`\mathrm{Prob}[\sigma _i(t)=\pm \xi _i^1]=\frac{1}{2}[1\pm m(t)]`$ for all $`t0`$, such that the argument given in for $`t=0`$ (leading to a Gaussian $`P(z)`$) would hold generally, and where the mapping (94) would describe the overlap evolution at all times:
$$P_t(z)=[2\pi \alpha ]^{\frac{1}{2}}e^{\frac{1}{2}z^2/\alpha }:m(t+1)=Dz\mathrm{tanh}[\beta (m(t)+z\sqrt{\alpha })]$$
(95)
with the Gaussian measure $`Dz=(2\pi )^{\frac{1}{2}}e^{\frac{1}{2}z^2}dz`$. This equation, however, must be generally incorrect. Firstly, figure 5 in shows that knowledge of $`m(t)`$ only does not permit prediction of $`m(t+1)`$. Secondly, expansion of the right-hand side of (95) for small $`m(t)`$ shows that (95) predicts a critical noise level (at $`\alpha =0`$) of $`T_c=\beta _c^1=1`$, and a storage capacity (at $`T=0`$) of $`\alpha _c=2/\pi 0.637`$, whereas both numerical simulations and equilibrium statistical mechanical calculations point to $`\alpha _c0.139`$. Rather than taking all $`\sigma _i`$ to be independent, a weaker assumption would be to just assume the interference noise distribution $`P_t(z)`$ to be a zero-average Gaussian one, at any time, with statistically independent noise variables $`z`$ at different times. One can then derive (for $`N\mathrm{}`$ and fully connected networks) an evolution equation for the width $`\mathrm{\Sigma }(t)`$, giving :
$$P_t(z)=[2\pi \mathrm{\Sigma }^2(t)]^{\frac{1}{2}}e^{\frac{1}{2}z^2/\mathrm{\Sigma }^2(t)}:m(t+1)=Dz\mathrm{tanh}[\beta (m(t)+z\mathrm{\Sigma }(t))]$$
$$\mathrm{\Sigma }^2(t+1)=\alpha +2\alpha m(t+1)m(t)h[m(t),\mathrm{\Sigma }(t)]+\mathrm{\Sigma }^2(t)h^2[m(t),\mathrm{\Sigma }(t)]$$
with $`h[m,\mathrm{\Sigma }]=\beta \left[1Dz\mathrm{tanh}^2[\beta (m+z\mathrm{\Sigma })]\right]`$. These equations describe correctly the qualitative features of recall dynamics, and are found to work well when retrieval actually occurs. For non-retrieval trajectories, however, they appear to underestimate the impact of interference noise: they predict $`T_c=1`$ (at $`\alpha =0`$) and a storage capacity (at $`T=0`$) of $`\alpha _c0.1597`$ (which should have been about 0.139). A final refinement of the Gaussian approach consisted in allowing for correlations between the noise variables $`z`$ at different times (while still describing them by Gaussian distributions). This results in a hierarchy of macroscopic equations, which improve upon the previous Gaussian theories and even predict the correct stationary state and phase diagrams, but still fail to be correct at intermediate times. The fundamental problem with all Gaussian theories, however sophisticated, is clearly illustrated in figure 6 of : the interference noise distribution is generally not of a Gaussian shape. $`P_t(z)`$ is only approximately Gaussian when pattern recall occurs. Hence the successes of Gaussian theories in describing recall trajectories, and their perpetual problems in describing the non-recall ones.
Non-Gaussian Approximations. In view of the non-Gaussian shape of the interference noise distribution, several attempts have been made at constructing non-Gaussian approximations. In all cases the aim is to arrive at a theory involving only macroscopic observables with a single time-argument. Figure 6 of suggests that for a fully connected network with binary neurons and parallel dynamics a more accurate ansatz for $`P_t(z)`$ would be the sum of two Gaussians. In the following choice was proposed, guided by the structure of the exact formalism to be described later:
$$P_t(z)=P_t^+(z)+P_t^{}(z),P_t^\pm (z)=\underset{N\mathrm{}}{lim}\frac{1}{N}\underset{i}{}\delta _{\sigma _i(t),\pm \xi _i^1}\delta [z\frac{1}{N}\underset{\mu >1}{}\xi _i^1\xi _i^\mu \underset{ji}{}\xi _j^\mu \sigma _j(t)]$$
$$P_t^\pm (z)=\frac{1\pm m(t)}{2\mathrm{\Sigma }(t)\sqrt{2\pi }}e^{\frac{1}{2}[zd(t)]^2/\mathrm{\Sigma }^2(t)}$$
followed by a self-consistent calculation of $`d(t)`$ (representing an effective ‘retarded self-interaction’, since it has an effect equivalent to adding $`h_i(𝝈(t))h_i(𝝈(t))+d(t)\sigma _i(t)`$), and of the width $`\mathrm{\Sigma }(t)`$ of the two distributions $`P_t^\pm (z)`$, together with
$$m(t+1)=\frac{1}{2}[1+m(t)]Dz\mathrm{tanh}[\beta (m(t)+d(t)+z\mathrm{\Sigma }(t))]+\frac{1}{2}[1m(t)]Dz\mathrm{tanh}[\beta (m(t)d(t)+z\mathrm{\Sigma }(t))]$$
The resulting three-parameter theory, in the form of closed dynamic equations for $`\{m,d,\mathrm{\Sigma }\}`$, is found to give a nice (but not perfect) agreement with numerical simulations.
A different philosophy was followed in (for sequential dynamics). First (as yet exact) equations are derived for the evolution of the two macroscopic observables $`m(𝝈)=m_1(𝝈)`$ and $`r(𝝈)=\alpha ^1_{\mu >1}m_\mu ^2(𝝈)`$, with $`m_\mu (𝝈)=N^1_i\xi _i^1\sigma _i`$, which are both found to involve $`P_t(z)`$:
$$\frac{d}{dt}m=𝑑zP_t(z)\mathrm{tanh}[\beta (m+z)]\frac{d}{dt}r=\frac{1}{\alpha }𝑑zP_t(z)z\mathrm{tanh}[\beta (m+z)]+1r$$
Next one closes these equations by hand, using a maximum-entropy (or ‘Occam’s Razor’) argument: instead of calculating $`P_t(z)`$ from (94) with the real (unknown) microscopic distribution $`p_t(𝝈)`$, it is calculated upon assigning equal probabilities to all states $`𝝈`$ with $`m(𝝈)=m`$ and $`r(𝝈)=r`$, followed by averaging over all realisations of the stored patterns with $`\mu >1`$. In order words: one assumes (i) that the microscopic states visited by the system are ‘typical’ within the appropriate $`(m,r)`$ sub-shells of state space, and (ii) that one can average over the disorder. Assumption (ii) is harmless, the most important step is (i). This procedure results in an explicit (non-Gaussian) expression for the noise distribution in terms of $`(m,r)`$ only, a closed two-parameter theory which is exact for short times and in equilibrium, accurate predictions of the macroscopic flow in the $`(m,r)`$-plane (such as that shown in figure 5 of ), but (again) deviations in predicted time-dependencies at intermediate times. This theory, and its performance, was later improved by applying the same ideas to a derivation of a dynamic equation for the function $`P_t(z)`$ itself (rather than for $`m`$ and $`r`$ only) ; research is still under way with the aim to construct a theory along these lines which is fully exact.
Exact Results: Generating Functional Analysis. The only fully exact procedure available at present is known under various names, such as ‘generating functional analysis’, ‘path integral formalism’ or ‘dynamic mean-field theory’, and is based on a philosophy different from those described so far. Rather than working with the probability $`p_t(𝝈)`$ of finding a microscopic state $`𝝈`$ at time $`t`$ in order to calculate the statistics of a set of macroscopic observables $`𝛀(𝝈)`$ at time $`t`$, one here turns to the probability $`\mathrm{Prob}[𝝈(0),\mathrm{},𝝈(t_m)]`$ of finding a microscopic path $`𝝈(0)𝝈(1)\mathrm{}𝝈(t_m)`$. One also adds time-dependent external sources to the local fields, $`h_i(𝝈)h_i(𝝈)+\theta _i(t)`$, in order to probe the networks via perturbations and define a response function. The idea is to concentrate on the moment generating function $`Z[𝝍]`$, which, like $`\mathrm{Prob}[𝝈(0),\mathrm{},𝝈(t_m)]`$, fully captures the statistics of paths:
$$Z[𝝍]=e^{i_i_{t=0}^{t_m}\psi _i(t)\sigma _i(t)}$$
(96)
It generates averages of the relevant observables, including those involving neuron states at different times, such as correlation functions $`C_{ij}(t,t^{})=\sigma _i(t)\sigma _j(t^{})`$ and response functions $`G_{ij}(t,t^{})=\sigma _i(t)/\theta _j(t^{})`$, upon differentiation with respect to the dummy variables $`\{\psi _i(t)\}`$:
$$\sigma _i(t)=i\underset{𝝍\mathrm{𝟎}}{lim}\frac{Z[𝝍]}{\psi _i(t)}C_{ij}(t,t^{})=\underset{𝝍\mathrm{𝟎}}{lim}\frac{^2Z[𝝍]}{\psi _i(t)\psi _j(t^{})}G_{ij}(t,t^{})=i\underset{𝝍\mathrm{𝟎}}{lim}\frac{^2Z[𝝍]}{\psi _i(t)\theta _j(t^{})}$$
(97)
Next one assumes (correctly) that for $`N\mathrm{}`$ only the statistical properties of the stored patterns will influence the macroscopic quantities, so that the generating function $`Z[𝝍]`$ can be averaged over all pattern realisations, i.e. $`Z[𝝍]\overline{Z[𝝍]}`$. As in replica theories (the canonical tool to deal with complexity in equilibrium) one carries out the disorder average before the average over the statistics of the neuron states, resulting for $`N\mathrm{}`$ in what can be interpreted as a theory describing a single ‘effective’ binary neuron $`\sigma (t)`$, with an effective local field $`h(t)`$ and the dynamics $`Prob[\sigma (t+1)=\pm 1]=\frac{1}{2}[1\pm \mathrm{tanh}[\beta h(t)]]`$. However, this effective local field is found to generally depend on past states of the neuron, and on zero-average but temporally correlated Gaussian noise contributions $`\varphi (t)`$:
$$h(t|\{\sigma \},\{\varphi \})=m(t)+\theta (t)+\alpha \underset{t^{}<t}{}R(t,t^{})\sigma (t^{})+\sqrt{\alpha }\varphi (t)$$
(98)
The first comprehensive neural network studies along these lines, dealing with fully connected networks, were carried out in , followed by applications to a-symmetrically and symmetrically extremely diluted networks (we will come back to those later). More recent applications include sequence processing networks <sup>6</sup><sup>6</sup>6In the case of sequence recall the overlap $`m`$ is defined with respect to the ‘moving’ target, i.e. $`m(t)=\frac{1}{N}_i\sigma _i(t)\xi _i^t`$. For $`N\mathrm{}`$ the differences between different models are found to show up only in the actual form taken by the effective local field (98), i.e. in the dependence of the ‘retarded self-interaction’ kernel $`R(t,t^{})`$ and the covariance matrix $`\varphi (t)\varphi (t^{})`$ of the interference-induced Gaussian noise on the macroscopic objects $`𝑪=\{C(s,s^{})=lim_N\mathrm{}\frac{1}{N}_iC_{ii}(s,s^{})\}`$ and $`𝑮=\{G(s,s^{})=lim_N\mathrm{}\frac{1}{N}_iG_{ii}(s,s^{})\}`$. For instance<sup>7</sup><sup>7</sup>7In the case of extremely diluted models the structure variables are also treated as disorder, and thus averaged out.:
$$\begin{array}{ccccccc}\mathrm{𝑚𝑜𝑑𝑒𝑙}& & \mathrm{𝑠𝑦𝑛𝑎𝑝𝑠𝑒𝑠}J_{ij}& & R(t,t^{})& & \varphi (t)\varphi (t^{})\\ & & & & & & \\ & & & & & & \\ \mathrm{𝑓𝑢𝑙𝑙𝑦}\mathrm{𝑐𝑜𝑛𝑛𝑒𝑐𝑡𝑒𝑑},& & \frac{1}{N}_{\mu =1}^{\alpha N}\xi _i^\mu \xi _j^\mu & & [(1\mathrm{I}𝑮)^1𝑮](t,t^{})& & [(1\mathrm{I}𝑮)^1𝑪(1\mathrm{I}𝑮^{})^1](t,t^{})\\ \mathrm{𝑠𝑡𝑎𝑡𝑖𝑐}\mathrm{𝑝𝑎𝑡𝑡𝑒𝑟𝑛𝑠}& & & & & & \\ \mathrm{𝑓𝑢𝑙𝑙𝑦}\mathrm{𝑐𝑜𝑛𝑛𝑒𝑐𝑡𝑒𝑑},& & \frac{1}{N}_{\mu =1}^{\alpha N}\xi _i^{\mu +1}\xi _j^\mu & & 0& & _{n0}[(𝑮^{})^n𝑪𝑮^n](t,t^{})\\ \mathrm{𝑝𝑎𝑡𝑡𝑒𝑟𝑛}\mathrm{𝑠𝑒𝑞𝑢𝑒𝑛𝑐𝑒}& & & & & & \\ \mathrm{𝑠𝑦𝑚𝑚}\mathrm{𝑒𝑥𝑡𝑟}\mathrm{𝑑𝑖𝑙𝑢𝑡𝑒𝑑},& & \frac{c_{ij}}{c}_{\mu =1}^{\alpha c}\xi _i^\mu \xi _j^\mu & & G(t,t^{})& & C(t,t^{})\\ \mathrm{𝑠𝑡𝑎𝑡𝑖𝑐}\mathrm{𝑝𝑎𝑡𝑡𝑒𝑟𝑛𝑠}& & & & & & \\ \mathrm{𝑎𝑠𝑦𝑚𝑚}\mathrm{𝑒𝑥𝑡𝑟}\mathrm{𝑑𝑖𝑙𝑢𝑡𝑒𝑑},& & \frac{c_{ij}}{c}_{\mu =1}^{\alpha c}\xi _i^\mu \xi _j^\mu & & 0& & C(t,t^{})\\ \mathrm{𝑠𝑡𝑎𝑡𝑖𝑐}\mathrm{𝑝𝑎𝑡𝑡𝑒𝑟𝑛𝑠}& & & & & & \end{array}$$
with the $`c_{ij}`$ drawn at random according to $`P(c_{ij})=\frac{c}{N}\delta _{c_{ij},1}+(1\frac{c}{N})\delta _{c_{ij},0}`$ (either symmetrically, i.e. $`c_{ij}=c_{ji}`$, or independently) and where $`c_{ii}=0`$, $`lim_N\mathrm{}c/N=0`$, and $`c\mathrm{}`$. In all cases the observables (overlaps and correlation- and response-functions) are to be solved from the following closed equations, involving the statistics of the single effective neuron experiencing the field (98):
$$m(t)=\sigma (t)C(t,t^{})=\sigma (t)\sigma (t^{})G(t,t^{})=\sigma (t)/\theta (t^{})$$
(99)
It is now clear that Gaussian theories can at most produce exact results for asymmetric networks. Any degree of symmetry in the synapses is found to induce a non-zero retarded self-interaction, via the kernel $`K(t,t^{})`$, which constitutes a non-Gaussian contribution to the local fields. Exact closed macroscopic theories apparently require a number of macroscopic observables which grows as $`𝒪(t^2)`$ in order to predict the dynamics up to time $`t`$. In the case of sequential dynamics the picture is found to be very similar to the one above; instead of discrete time labels $`t\{0,1,\mathrm{},t_m\}`$, path summations and matrices, there one has a real time variable $`t[0,t_m]`$, path-integrals and integral operators. The remainder of this paper is devoted to the derivation of the above results and their implications.
### 5.2 Generating Functional Analysis for Binary Neurons
General Definitions. I will now show more explicitly how the generating functional formalism works for networks of binary neurons. We define parallel dynamics, i.e. (22), driven as usual by local fields of the form $`h_i(𝝈;t)=_jJ_{ij}\sigma _j+\theta _i(t)`$, but with a more general choice of Hebbian-type synapses, in which we allow for a possible random dilution (to reduce repetition in our subsequent derivations):
$$J_{ij}=\frac{c_{ij}}{c}\underset{\mu =1}{\overset{p}{}}\xi _i^\mu \xi _j^\mu p=\alpha c$$
(100)
Architectural properties are reflected in the variables $`c_{ij}\{0,1\}`$, whereas information storage is to be effected by the remainder in (100), involving $`p`$ randomly and independently drawn patterns $`𝝃^\mu =(\xi _1^\mu ,\mathrm{},\xi _N^\mu )\{1,1\}^N`$. I will deal both with symmetric and with asymmetric architectures (always putting $`c_{ii}=0`$), in which the variables $`c_{ij}`$ are drawn randomly according to
$$\mathrm{𝑠𝑦𝑚𝑚𝑒𝑡𝑟𝑖𝑐}:c_{ij}=c_{ji},i<jP(c_{ij})=\frac{c}{N}\delta _{c_{ij},1}+(1\frac{c}{N})\delta _{c_{ij},0}$$
(101)
$$\mathrm{𝑎𝑠𝑦𝑚𝑚𝑒𝑡𝑟𝑖𝑐}:ijP(c_{ij})=\frac{c}{N}\delta _{c_{ij},1}+(1\frac{c}{N})\delta _{c_{ij},0}$$
(102)
(one could also study intermediate degrees of symmetry; this would involve only simple adaptations). Thus $`c_{kl}`$ is statistically independent of $`c_{ij}`$ as soon as $`(k,l)\{(i,j),(j,i)\}`$. In leading order in $`N`$ one has $`_jc_{ij}=c`$ for all $`i`$, so $`c`$ gives the average number of neurons contributing to the field of any given neuron. In view of this, the number $`p`$ of patterns to be stored can be expected to scale as $`p=\alpha c`$. The connectivity parameter $`c`$ is chosen to diverge with $`N`$, i.e. $`lim_N\mathrm{}c^1=0`$. If $`c=N`$ we obtain the fully connected (parallel dynamics) Hopfield model. Extremely diluted networks are obtained when $`lim_N\mathrm{}c/N=0`$.
For simplicity we make the so-called ‘condensed ansatz’: we assume that the system state has an $`𝒪(N^0)`$ overlap only with a single pattern, say $`\mu =1`$. This situation is induced by initial conditions: we take a randomly drawn $`𝝈(0)`$, generated by
$$p(𝝈(0))=\underset{i}{}\left\{\frac{1}{2}[1+m_0]\delta _{\sigma _i(0),\xi _i^1}+\frac{1}{2}[1m_0]\delta _{\sigma _i(0),\xi _i^1}\right\}\mathrm{so}\frac{1}{N}\underset{i}{}\xi _i^1\sigma _i(0)=m_0$$
(103)
The patterns $`\mu >1`$, as well as the architecture variables $`c_{ij}`$, are viewed as disorder. One assumes that for $`N\mathrm{}`$ the macroscopic behaviour of the system is ‘self-averaging’, i.e. only dependent on the statistical properties of the disorder (rather than on its microscopic realisation). Averages over the disorder are written as $`\overline{\mathrm{}}`$. We next define the disorder-averaged generating function:
$$\overline{Z[𝝍]}=\overline{e^{i_i_t\psi _i(t)\sigma _i(t)}}$$
(104)
in which the time $`t`$ runs from $`t=0`$ to some (finite) upper limit $`t_m`$. Note that $`\overline{Z[\mathrm{𝟎}]}=1`$. With a modest amount of foresight we define the macroscopic site-averaged and disorder-averaged objects $`m(t)=N^1_i\xi _i^1\overline{\sigma _i(t)}`$, $`C(t,t^{})=N^1_i\overline{\sigma _i(t)\sigma _i(t^{})}`$ and $`G(t,t^{})=N^1_i\overline{\sigma _i(t)}/\theta _i(t^{})`$. According to (97) they can be obtained from (104) as follows:
$$m(t)=\underset{𝝍\mathrm{𝟎}}{lim}\frac{i}{N}\underset{j}{}\xi _j^1\frac{\overline{Z[𝝍]}}{\psi _j(t)}$$
(105)
$$C(t,t^{})=\underset{𝝍\mathrm{𝟎}}{lim}\frac{1}{N}\underset{j}{}\frac{^2\overline{Z[𝝍]}}{\psi _j(t)\psi _j(t^{})}G(t,t^{})=\underset{𝝍\mathrm{𝟎}}{lim}\frac{i}{N}\underset{j}{}\frac{^2\overline{Z[𝝍]}}{\psi _j(t)\theta _j(t^{})}$$
(106)
So far we have only reduced our problem to the calculation of the function $`\overline{Z[𝝍]}`$ in (104), which will play a part similar to that of the disorder-averaged free energy in equilibrium calculations (see ).
Evaluation of the Disorder-Averaged Generating Function. As in equilibrium replica calculations, the hope is that progress can be made by carrying out the disorder averages first. In equilibrium calculations we use the replica trick to convert our disorder averages into feasible ones; here the idea is to isolate the local fields at different times and different sites by inserting appropriate $`\delta `$-distributions:
$$1=\underset{it}{}𝑑h_i(t)\delta [h_i(t)\underset{j}{}J_{ij}\sigma _j(t)\theta _i(t)]=\{d𝒉d\widehat{𝒉}\}e^{i_{it}\widehat{h}_i(t)[h_i(t)_jJ_{ij}\sigma _j(t)\theta _i(t)]}$$
with $`\{d𝒉d\widehat{𝒉}\}=_{it}[d\widehat{h}_i(t)dh_i(t)/2\pi ]`$, giving
$$\overline{Z[𝝍]}=\{d𝒉d\widehat{𝒉}\}e^{i_{it}\widehat{h}_i(t)[h_i(t)\theta _i(t)]}e^{i_{it}\psi _i(t)\sigma _i(t)}\overline{\left[e^{i_{it}\widehat{h}_i(t)_jJ_{ij}\sigma _j(t)}\right]}_{\mathrm{pf}}$$
in which $`\mathrm{}_{\mathrm{pf}}`$ refers to averages over a constrained stochastic process of the type (22), but with prescribed fields $`\{h_i(t)\}`$ at all sites and at all times. Note that with such prescribed fields the probability of generating a path $`\{𝝈(0),\mathrm{},𝝈(t_m)\}`$ is given by
$$\mathrm{Prob}[𝝈(0),\mathrm{},𝝈(t_m)|\{h_i(t)\}]=p(𝝈(0))e^{_{it}[\beta \sigma _i(t+1)h_i(t)\mathrm{log}2\mathrm{cosh}[\beta h_i(t)]]}$$
so
$$\overline{Z[𝝍]}=\{d𝒉d\widehat{𝒉}\}\underset{𝝈(0)}{}\mathrm{}\underset{𝝈(t_m)}{}p(𝝈(0))e^{N[\{𝝈\},\{\widehat{𝒉}\}]}\underset{it}{}e^{i\widehat{h}_i(t)[h_i(t)\theta _i(t)]i\psi _i(t)\sigma _i(t)+\beta \sigma _i(t+1)h_i(t)\mathrm{log}2\mathrm{cosh}[\beta h_i(t)]}$$
(107)
with
$$[\{𝝈\},\{\widehat{𝒉}\}]=\frac{1}{N}\mathrm{log}\overline{\left[e^{i_{it}\widehat{h}_i(t)_jJ_{ij}\sigma _j(t)}\right]}$$
(108)
We concentrate on the term $`[\mathrm{}]`$ (with the disorder), of which we need only know the limit $`N\mathrm{}`$, since only terms inside $`\overline{Z[𝝍]}`$ which are exponential in $`N`$ will retain statistical relevance. In the disorder-average of (108) every site $`i`$ plays an equivalent role, so the leading order in $`N`$ of (108) should depend only on site-averaged functions of the $`\{\sigma _i(t),\widehat{h}_i(t)\}`$, with no reference to any special direction except the one defined by pattern $`𝝃^1`$. The simplest such functions with a single time variable are
$$a(t;\{𝝈\})=\frac{1}{N}\underset{i}{}\xi _i^1\sigma _i(t)k(t;\{\widehat{𝒉}\})=\frac{1}{N}\underset{i}{}\xi _i^1\widehat{h}_i(t)$$
(109)
whereas the simplest ones with two time variables would appear to be
$$q(t,t^{};\{𝝈\})=\frac{1}{N}\underset{i}{}\sigma _i(t)\sigma _i(t^{})Q(t,t^{};\{\widehat{𝒉}\})=\frac{1}{N}\underset{i}{}\widehat{h}_i(t)\widehat{h}_i(t^{})$$
(110)
$$K(t,t^{};\{𝝈,\widehat{𝒉}\})=\frac{1}{N}\underset{i}{}\widehat{h}_i(t)\sigma _i(t^{})$$
(111)
It will turn out that all models of the type (100), with either (101) or (102), have the crucial property that (109,110,111) are in fact the only functions to appear in the leading order of (108):
$$[\mathrm{}]=\mathrm{\Phi }[\{a(t;\mathrm{}),k(t;\mathrm{}),q(t,t^{};\mathrm{}),Q(t,t^{};\mathrm{}),K(t,t^{};\mathrm{})\}]+\mathrm{}(N\mathrm{})$$
(112)
for some as yet unknown function $`\mathrm{\Phi }[\mathrm{}]`$. This allows us to proceed with the evaluation of (107). We can achieve site factorisation in (107) if we isolate the macroscopic objects (109,110,111) by introducing suitable $`\delta `$-distributions (taking care that all exponents scale linearly with $`N`$, to secure statistical relevance). Thus we insert
$$1=\underset{t=0}{\overset{t_m}{}}𝑑a(t)\delta [a(t)a(t;\{𝝈\})]=\left[\frac{N}{2\pi }\right]^{t_m+1}𝑑𝒂𝑑\widehat{𝒂}e^{iN_t\widehat{a}(t)[a(t)\frac{1}{N}_j\xi _j^1\sigma _j(t)]}$$
$$1=\underset{t=0}{\overset{t_m}{}}𝑑k(t)\delta [k(t)k(t;\{\widehat{𝒉}\})]=\left[\frac{N}{2\pi }\right]^{t_m+1}𝑑𝒌𝑑\widehat{𝒌}e^{iN_t\widehat{k}(t)[k(t)\frac{1}{N}_j\xi _j^1\widehat{h}_j(t)]}$$
$$1=\underset{t,t^{}=0}{\overset{t_m}{}}𝑑q(t,t^{})\delta [q(t,t^{})q(t,t^{};\{𝝈\})]=\left[\frac{N}{2\pi }\right]^{(t_m+1)^2}𝑑𝒒𝑑\widehat{𝒒}e^{iN_{t,t^{}}\widehat{q}(t,t^{})[q(t,t^{})\frac{1}{N}_j\sigma _j(t)\sigma _j(t^{})]}$$
$$1=\underset{t,t^{}=0}{\overset{t_m}{}}𝑑Q(t,t^{})\delta [Q(t,t^{})Q(t,t^{};\{\widehat{𝒉}\})]=\left[\frac{N}{2\pi }\right]^{(t_m+1)^2}𝑑𝑸𝑑\widehat{𝑸}e^{iN_{t,t^{}}\widehat{Q}(t,t^{})[Q(t,t^{})\frac{1}{N}_j\widehat{h}_j(t)\widehat{h}_j(t^{})]}$$
$$1=\underset{t,t^{}=0}{\overset{t_m}{}}𝑑K(t,t^{})\delta [K(t,t^{})K(t,t^{};\{𝝈,\widehat{𝒉}\})]=\left[\frac{N}{2\pi }\right]^{(t_m+1)^2}𝑑𝑲𝑑\widehat{𝑲}e^{iN_{t,t^{}}\widehat{K}(t,t^{})[K(t,t^{})\frac{1}{N}_j\widehat{h}_j(t)\sigma _j(t^{})]}$$
Insertion of these integrals into (107), followed by insertion of (112) and usage of the short-hand
$$\mathrm{\Psi }[𝒂,\widehat{𝒂},𝒌,\widehat{𝒌},𝒒,\widehat{𝒒},𝑸,\widehat{𝑸},𝑲,\widehat{𝑲}]=i\underset{t}{}[\widehat{a}(t)a(t)+\widehat{k}(t)k(t)]$$
$$+i\underset{t,t^{}}{}[\widehat{q}(t,t^{})q(t,t^{})+\widehat{Q}(t,t^{})Q(t,t^{})+\widehat{K}(t,t^{})K(t,t^{})]$$
(113)
then leads us to
$$\overline{Z[𝝍]}=𝑑𝒂𝑑\widehat{𝒂}𝑑𝒌𝑑\widehat{𝒌}𝑑𝒒𝑑\widehat{𝒒}𝑑𝑸𝑑\widehat{𝑸}𝑑𝑲𝑑\widehat{𝑲}e^{N\mathrm{\Psi }[𝒂,\widehat{𝒂},𝒌,\widehat{𝒌},𝒒,\widehat{𝒒},𝑸,\widehat{𝑸},𝑲,\widehat{𝑲}]+N\mathrm{\Phi }[𝒂,𝒌,𝒒,𝑸,𝑲]+𝒪(\mathrm{})}$$
$$\times \{d𝒉d\widehat{𝒉}\}\underset{𝝈(0)}{}\mathrm{}\underset{𝝈(t_m)}{}p(𝝈(0))\underset{it}{}e^{i\widehat{h}_i(t)[h_i(t)\theta _i(t)]i\psi _i(t)\sigma _i(t)+\beta \sigma _i(t+1)h_i(t)\mathrm{log}2\mathrm{cosh}[\beta h_i(t)]}$$
$$\times \underset{i}{}e^{i\xi _i^1_t[\widehat{a}(t)\sigma _i(t)+\widehat{k}(t)\widehat{h}_i(t)]i_{t,t^{}}[\widehat{q}(t,t^{})\sigma _i(t)\sigma _i(t^{})+\widehat{Q}(t,t^{})\widehat{h}_i(t)\widehat{h}_i(t^{})+\widehat{K}(t,t^{})\widehat{h}_i(t)\sigma _i(t^{})]}$$
(114)
in which the term denoted as $`𝒪(\mathrm{})`$ covers both the non-dominant orders in (108) and the $`𝒪(\mathrm{log}N)`$ relics of the various pre-factors $`[N/2\pi ]`$ in the above integral representations of the $`\delta `$-distributions (note: $`t_m`$ was assumed fixed). We now see explicitly in (114) that the summations and integrations over neuron states and local fields fully factorise over the $`N`$ sites. A simple transformation $`\{\sigma _i(t),h_i(t),\widehat{h}_i(t)\}\{\xi _i^1\sigma _i(t),\xi _i^1h_i(t),\xi _i^1\widehat{h}_i(t)\}`$ brings the result into the form
$$\{d𝒉d\widehat{𝒉}\}\underset{𝝈(0)}{}\mathrm{}\underset{𝝈(t_m)}{}p(𝝈(0))\underset{it}{}e^{i\widehat{h}_i(t)[h_i(t)\xi _i^1\theta _i(t)]i\xi _i^1\psi _i(t)\sigma _i(t)+\beta \sigma _i(t+1)h_i(t)\mathrm{log}2\mathrm{cosh}[\beta h_i(t)]}$$
$$\times \underset{i}{}e^{i\xi _i^1_t[\widehat{a}(t)\sigma _i(t)+\widehat{k}(t)\widehat{h}_i(t)]i_{t,t^{}}[\widehat{q}(t,t^{})\sigma _i(t)\sigma _i(t^{})+\widehat{Q}(t,t^{})\widehat{h}_i(t)\widehat{h}_i(t^{})+\widehat{K}(t,t^{})\widehat{h}_i(t)\sigma _i(t^{})]}$$
$$=e^{N\mathrm{\Xi }[\widehat{𝒂},\widehat{𝒌},\widehat{𝒒},\widehat{𝑸},\widehat{𝑲}]}$$
with
$$\mathrm{\Xi }[\widehat{𝒂},\widehat{𝒌},\widehat{𝒒},\widehat{𝑸},\widehat{𝑲}]=\frac{1}{N}\underset{i}{}\mathrm{log}\{dhd\widehat{h}\}\underset{\sigma (0)\mathrm{}\sigma (t_m)}{}\pi _0(\sigma (0))e^{_t\{i\widehat{h}(t)[h(t)\xi _i^1\theta _i(t)]i\xi _i^1\psi _i(t)\sigma (t)\}}$$
$$\times e^{_t\{\beta \sigma (t+1)h(t)\mathrm{log}2\mathrm{cosh}[\beta h(t)]\}i_t[\widehat{a}(t)\sigma (t)+\widehat{k}(t)\widehat{h}(t)]i_{t,t^{}}[\widehat{q}(t,t^{})\sigma (t)\sigma (t^{})+\widehat{Q}(t,t^{})\widehat{h}(t)\widehat{h}(t^{})+\widehat{K}(t,t^{})\widehat{h}(t)\sigma (t^{})]}$$
(115)
in which $`\{dhd\widehat{h}\}=_t[dh(t)d\widehat{h}(t)/2\pi ]`$ and $`\pi _0(\sigma )=\frac{1}{2}[1+m_0]\delta _{\sigma ,1}+\frac{1}{2}[1m_0]\delta _{\sigma ,1}`$. At this stage (114) acquires the form of an integral to be evaluated via the saddle-point (or ‘steepest descent’) method:
$$\overline{Z[\{\psi (t)\}]}=𝑑𝒂𝑑\widehat{𝒂}𝑑𝒌𝑑\widehat{𝒌}𝑑𝒒𝑑\widehat{𝒒}𝑑𝑸𝑑\widehat{𝑸}𝑑𝑲𝑑\widehat{𝑲}e^{N\left\{\mathrm{\Psi }[\mathrm{}]+\mathrm{\Phi }[\mathrm{}]+\mathrm{\Xi }[\mathrm{}]\right\}+𝒪(\mathrm{})}$$
(116)
in which the functions $`\mathrm{\Psi }[\mathrm{}]`$, $`\mathrm{\Phi }[\mathrm{}]`$ and $`\mathrm{\Xi }[\mathrm{}]`$ are defined by (112,113,115).
The Saddle-Point Problem. The disorder-averaged generating function (116) is for $`N\mathrm{}`$ dominated by the physical saddle-point of the macroscopic surface
$$\mathrm{\Psi }[𝒂,\widehat{𝒂},𝒌,\widehat{𝒌},𝒒,\widehat{𝒒},𝑸,\widehat{𝑸},𝑲,\widehat{𝑲}]+\mathrm{\Phi }[𝒂,𝒌,𝒒,𝑸,𝑲]+\mathrm{\Xi }[\widehat{𝒂},\widehat{𝒌},\widehat{𝒒},\widehat{𝑸},\widehat{𝑲}]$$
(117)
with the three contributions defined in (112,113,115). It will be advantageous at this stage to define the following effective measure (which will be further simplified later):
$$f[\{\sigma \},\{h\},\{\widehat{h}\}]_{}=\frac{1}{N}\underset{i}{}\left\{\frac{\{dhd\widehat{h}\}_{\sigma (0)\mathrm{}\sigma (t_m)}M_i[\{\sigma \},\{h\},\{\widehat{h}\}]f[\{\sigma \},\{h\},\{\widehat{h}\}]}{\{dhd\widehat{h}\}_{\sigma (0)\mathrm{}\sigma (t_m)}M_i[\{\sigma \},\{h\},\{\widehat{h}\}]}\right\}$$
(118)
with
$$M_i[\{\sigma \},\{h\},\{\widehat{h}\}]=\pi _0(\sigma (0))e^{_t\{i\widehat{h}(t)[h(t)\xi _i^1\theta _i(t)]i\xi _i^1\psi _i(t)\sigma (t)+\beta \sigma (t+1)h(t)\mathrm{log}2\mathrm{cosh}[\beta h(t)]\}}$$
$$\times e^{i_t[\widehat{a}(t)\sigma (t)+\widehat{k}(t)\widehat{h}(t)]i_{t,t^{}}[\widehat{q}(t,t^{})\sigma (t)\sigma (t^{})+\widehat{Q}(t,t^{})\widehat{h}(t)\widehat{h}(t^{})+\widehat{K}(t,t^{})\widehat{h}(t)\sigma (t^{})]}$$
in which the values to be inserted for $`\{\widehat{m}(t),\widehat{k}(t),\widehat{q}(t,t^{}),\widehat{Q}(t,t^{}),\widehat{K}(t,t^{})\}`$ are given by the saddle-point of (117). Variation of (117) with respect to all the original macroscopic objects occurring as arguments (those without the ‘hats’) gives the following set of saddle-point equations:
$$\widehat{a}(t)=i\mathrm{\Phi }/a(t)\widehat{k}(t)=i\mathrm{\Phi }/k(t)$$
(119)
$$\widehat{q}(t,t^{})=i\mathrm{\Phi }/q(t,t^{})\widehat{Q}(t,t^{})=i\mathrm{\Phi }/Q(t,t^{})\widehat{K}(t,t^{})=i\mathrm{\Phi }/K(t,t^{})$$
(120)
Variation of (117) with respect to the conjugate macroscopic objects (those with the ‘hats’), in turn, and usage of our newly introduced short-hand notation $`\mathrm{}_{}`$, gives:
$$a(t)=\sigma (t)_{}k(t)=\widehat{h}(t)_{}$$
(121)
$$q(t,t^{})=\sigma (t)\sigma (t^{})_{}Q(t,t^{})=\widehat{h}(t)\widehat{h}(t^{})_{}K(t,t^{})=\widehat{h}(t)\sigma (t^{})_{}$$
(122)
The coupled equations (119,120,121,122) are to be solved simultaneously, once we have calculated the term $`\mathrm{\Phi }[\mathrm{}]`$ (112) which depends on the synapses. This appears to be a formidable task; it can, however, be simplified considerably upon first deriving the physical meaning of the above macroscopic quantities. We apply (105,106) to (116), using identities such as
$$\frac{\mathrm{\Xi }[\mathrm{}]}{\psi _j(t)}=\frac{i}{N}\xi _j^1\left[\frac{\{dhd\widehat{h}\}_{\sigma (0)\mathrm{}\sigma (t_m)}M_j[\{\sigma \},\{h\},\{\widehat{h}\}]\sigma (t)}{\{dhd\widehat{h}\}_{\sigma (0)\mathrm{}\sigma (t_m)}M_j[\{\sigma \},\{h\},\{\widehat{h}\}]}\right]$$
$$\frac{\mathrm{\Xi }[\mathrm{}]}{\theta _j(t)}=\frac{i}{N}\xi _j^1\left[\frac{\{dhd\widehat{h}\}_{\sigma (0)\mathrm{}\sigma (t_m)}M_j[\{\sigma \},\{h\},\{\widehat{h}\}]\widehat{h}(t)}{\{dhd\widehat{h}\}_{\sigma (0)\mathrm{}\sigma (t_m)}M_j[\{\sigma \},\{h\},\{\widehat{h}\}]}\right]$$
$$\frac{^2\mathrm{\Xi }[\mathrm{}]}{\psi _j(t)\psi _j(t^{})}=\frac{1}{N}\left[\frac{\{dhd\widehat{h}\}_{\sigma (0)\mathrm{}\sigma (t_m)}M_j[\{\sigma \},\{h\},\{\widehat{h}\}]\sigma (t)\sigma (t^{})}{\{dhd\widehat{h}\}_{\sigma (0)\mathrm{}\sigma (t_m)}M_j[\{\sigma \},\{h\},\{\widehat{h}\}]}\right]N\left[\frac{\mathrm{\Xi }[\mathrm{}]}{\psi _j(t)}\right]\left[\frac{\mathrm{\Xi }[\mathrm{}]}{\psi _j(t^{})}\right]$$
$$\frac{^2\mathrm{\Xi }[\mathrm{}]}{\theta _j(t)\theta _j(t^{})}=\frac{1}{N}\left[\frac{\{dhd\widehat{h}\}_{\sigma (0)\mathrm{}\sigma (t_m)}M_j[\{\sigma \},\{h\},\{\widehat{h}\}]\widehat{h}(t)\widehat{h}(t^{})}{\{dhd\widehat{h}\}_{\sigma (0)\mathrm{}\sigma (t_m)}M_j[\{\sigma \},\{h\},\{\widehat{h}\}]}\right]N\left[\frac{\mathrm{\Xi }[\mathrm{}]}{\theta _j(t)}\right]\left[\frac{\mathrm{\Xi }[\mathrm{}]}{\theta _j(t^{})}\right]$$
$$\frac{^2\mathrm{\Xi }[\mathrm{}]}{\psi _j(t)\theta _j(t^{})}=\frac{i}{N}\left[\frac{\{dhd\widehat{h}\}_{\sigma (0)\mathrm{}\sigma (t_m)}M_j[\{\sigma \},\{h\},\{\widehat{h}\}]\sigma (t)\widehat{h}(t^{})}{\{dhd\widehat{h}\}_{\sigma (0)\mathrm{}\sigma (t_m)}M_j[\{\sigma \},\{h\},\{\widehat{h}\}]}\right]N\left[\frac{\mathrm{\Xi }[\mathrm{}]}{\psi _j(t)}\right]\left[\frac{\mathrm{\Xi }[\mathrm{}]}{\theta _j(t^{})}\right]$$
and using the short-hand notation (118) wherever possible. Note that the external fields $`\{\psi _i(t),\theta _i(t)\}`$ occur only in the function $`\mathrm{\Xi }[\mathrm{}]`$, not in $`\mathrm{\Psi }[\mathrm{}]`$ or $`\mathrm{\Phi }[\mathrm{}]`$, and that overall constants in $`\overline{Z[𝝍]}`$ can always be recovered a posteriori, using $`\overline{Z[\mathrm{𝟎}]}=1`$:
$$m(t)=\underset{𝝍\mathrm{𝟎}}{lim}\frac{i}{N}\underset{i}{}\xi _i^1\frac{𝑑𝒂\mathrm{}𝑑\widehat{𝑲}\left[\frac{N\mathrm{\Xi }}{\psi _i(t)}\right]e^{N[\mathrm{\Psi }+\mathrm{\Phi }+\mathrm{\Xi }]+𝒪(\mathrm{})}}{𝑑𝒂\mathrm{}𝑑\widehat{𝑲}e^{N[\mathrm{\Psi }+\mathrm{\Phi }+\mathrm{\Xi }]+𝒪(\mathrm{})}}=\underset{𝝍\mathrm{𝟎}}{lim}\sigma (t)_{}$$
$$C(t,t^{})=\underset{𝝍\mathrm{𝟎}}{lim}\frac{1}{N}\underset{i}{}\frac{𝑑𝒂\mathrm{}𝑑\widehat{𝑲}\left[\frac{N^2\mathrm{\Xi }}{\psi _i(t)\psi _i(t^{})}+\frac{N\mathrm{\Xi }}{\psi _i(t)}\frac{N\mathrm{\Xi }}{\psi _i(t^{})}\right]e^{N[\mathrm{\Psi }+\mathrm{\Phi }+\mathrm{\Xi }]+𝒪(\mathrm{})}}{𝑑𝒂\mathrm{}𝑑\widehat{𝑲}e^{N[\mathrm{\Psi }+\mathrm{\Phi }+\mathrm{\Xi }]+𝒪(\mathrm{})}}=\underset{𝝍\mathrm{𝟎}}{lim}\sigma (t)\sigma (t^{})_{}$$
$$iG(t,t^{})=\underset{𝝍\mathrm{𝟎}}{lim}\frac{1}{N}\underset{i}{}\frac{𝑑𝒂\mathrm{}𝑑\widehat{𝑲}\left[\frac{N^2\mathrm{\Xi }}{\psi _i(t)\theta _i(t^{})}+\frac{N\mathrm{\Xi }}{\psi _i(t)}\frac{N\mathrm{\Xi }}{\theta _i(t^{})}\right]e^{N[\mathrm{\Psi }+\mathrm{\Phi }+\mathrm{\Xi }]+𝒪(\mathrm{})}}{𝑑𝒂\mathrm{}𝑑\widehat{𝑲}e^{N[\mathrm{\Psi }+\mathrm{\Phi }+\mathrm{\Xi }]+𝒪(\mathrm{})}}=\underset{𝝍\mathrm{𝟎}}{lim}\sigma (t)\widehat{h}(t^{})_{}$$
Finally we obtain useful identities from the seemingly trivial statements $`N^1_i\xi _i^1\overline{Z[\mathrm{𝟎}]}/\theta _i(t)=0`$ and $`N^1_i^2\overline{Z[\mathrm{𝟎}]}/\theta _i(t)\theta _i(t^{})=0`$:
$$0=\underset{𝝍\mathrm{𝟎}}{lim}\frac{i}{N}\underset{i}{}\xi _i^1\frac{𝑑𝒂\mathrm{}𝑑\widehat{𝑲}\left[\frac{N\mathrm{\Xi }}{\theta _i(t)}\right]e^{N[\mathrm{\Psi }+\mathrm{\Phi }+\mathrm{\Xi }]+𝒪(\mathrm{})}}{𝑑𝒂\mathrm{}𝑑\widehat{𝑲}e^{N[\mathrm{\Psi }+\mathrm{\Phi }+\mathrm{\Xi }]+𝒪(\mathrm{})}}=\underset{𝝍\mathrm{𝟎}}{lim}\widehat{h}(t)_{}$$
$$0=\underset{𝝍\mathrm{𝟎}}{lim}\frac{1}{N}\underset{i}{}\frac{𝑑𝒂\mathrm{}𝑑\widehat{𝑲}\left[\frac{N^2\mathrm{\Xi }}{\theta _i(t)\theta _i(t^{})}+\frac{N\mathrm{\Xi }}{\theta _i(t)}\frac{N\mathrm{\Xi }}{\theta _i(t^{})}\right]e^{N[\mathrm{\Psi }+\mathrm{\Phi }+\mathrm{\Xi }]+𝒪(\mathrm{})}}{𝑑𝒂\mathrm{}𝑑\widehat{𝑲}e^{N[\mathrm{\Psi }+\mathrm{\Phi }+\mathrm{\Xi }]+𝒪(\mathrm{})}}=\underset{𝝍\mathrm{𝟎}}{lim}\widehat{h}(t)\widehat{h}(t^{})_{}$$
In combination with (121,122), the above five identities simplify our problem considerably. The dummy fields $`\psi _i(t)`$ have served their purpose and will now be put to zero, as a result we can now identify our macroscopic observables at the relevant saddle-point as:
$$a(t)=m(t)k(t)=0q(t,t^{})=C(t,t^{})Q(t,t^{})=0K(t,t^{})=iG(t^{},t)$$
(123)
Finally we make a convenient choice for the external fields, $`\theta _i(t)=\xi _i^1\theta (t)`$, with which the effective measure $`\mathrm{}_{}`$ of (124) simplifies to
$$f[\{\sigma \},\{h\},\{\widehat{h}\}]_{}=\frac{\{dhd\widehat{h}\}_{\sigma (0)\mathrm{}\sigma (t_m)}M[\{\sigma \},\{h\},\{\widehat{h}\}]f[\{\sigma \},\{h\},\{\widehat{h}\}]}{\{dhd\widehat{h}\}_{\sigma (0)\mathrm{}\sigma (t_m)}M[\{\sigma \},\{h\},\{\widehat{h}\}]}$$
(124)
with
$$M[\{\sigma \},\{h\},\{\widehat{h}\}]=\pi _0(\sigma (0))e^{_t\{i\widehat{h}(t)[h(t)\theta (t)]+\beta \sigma (t+1)h(t)\mathrm{log}2\mathrm{cosh}[\beta h(t)]\}i_t[\widehat{a}(t)\sigma (t)+\widehat{k}(t)\widehat{h}(t)]}$$
$$\times e^{i_{t,t^{}}[\widehat{q}(t,t^{})\sigma (t)\sigma (t^{})+\widehat{Q}(t,t^{})\widehat{h}(t)\widehat{h}(t^{})+\widehat{K}(t,t^{})\widehat{h}(t)\sigma (t^{})]}$$
In summary: our saddle-point equations are given by (119,120,121,122), and the physical meaning of the macroscopic quantities is given by (123) (apparently many of them must be zero). Our final task is finding (112), i.e. calculating the leading order of
$$[\{𝝈\},\{\widehat{𝒉}\}]=\frac{1}{N}\mathrm{log}\overline{\left[e^{i_{it}\widehat{h}_i(t)_jJ_{ij}\sigma _j(t)}\right]}$$
(125)
which is where the properties of the synapses (100) come in.
### 5.3 Parallel Dynamics Hopfield Model Near Saturation
The Disorder Average. The fully connected Hopfield network (here with parallel dynamics) is obtained upon choosing $`c=N`$ in the recipe (100), i.e. $`c_{ij}=1\delta _{ij}`$ and $`p=\alpha N`$. The disorder average thus involves only the patterns with $`\mu >1`$. In view of our objective to write (125) in the form (112), we will substitute the observables defined in (109,110,111) whenever possible. Now (125) gives
$$[\mathrm{}]=\frac{1}{N}\mathrm{log}\overline{\left[e^{iN^1_t_\mu _{ij}\xi _i^\mu \xi _j^\mu \widehat{h}_i(t)\sigma _j(t)}\right]}$$
$$=i\alpha \underset{t}{}K(t,t;\{𝝈,\widehat{𝒉}\})i\underset{t}{}a(t)k(t)+\alpha \mathrm{log}\overline{\left[e^{i_t[_i\xi _i\widehat{h}_i(t)/\sqrt{N}][_i\xi _i\sigma _i(t)/\sqrt{N}]}\right]}+𝒪(N^1)$$
(126)
We concentrate on the last term:
$$\overline{\left[e^{i_t[_i\xi _i\widehat{h}_i(t)/\sqrt{N}][_i\xi _i\sigma _i(t)/\sqrt{N}]}\right]}=𝑑𝒙𝑑𝒚e^{i𝒙𝒚}\overline{\underset{t}{}\left\{\delta [x(t)\frac{_i\xi _i\sigma _i(t)}{\sqrt{N}}]\delta [y(t)\frac{_i\xi _i\widehat{h}_i(t)}{\sqrt{N}}]\right\}}$$
$$=\frac{d𝒙d𝒚d\widehat{𝒙}d\widehat{𝒚}}{(2\pi )^{2(t_m+1)}}e^{i[\widehat{𝒙}𝒙+\widehat{𝒚}𝒚𝒙𝒚]}\overline{\left[e^{\frac{i}{\sqrt{N}}_i\xi _i_t[\widehat{x}(t)\sigma _i(t)+\widehat{y}(t)\widehat{h}_i(t)]}\right]}$$
$$=\frac{d𝒙d𝒚d\widehat{𝒙}d\widehat{𝒚}}{(2\pi )^{2(t_m+1)}}e^{i[\widehat{𝒙}𝒙+\widehat{𝒚}𝒚𝒙𝒚]+_i\mathrm{log}\mathrm{cos}\left[\frac{1}{\sqrt{N}}_t[\widehat{x}(t)\sigma _i(t)+\widehat{y}(t)\widehat{h}_i(t)]\right]}$$
$$=\frac{d𝒙d𝒚d\widehat{𝒙}d\widehat{𝒚}}{(2\pi )^{2(t_m+1)}}e^{i[\widehat{𝒙}𝒙+\widehat{𝒚}𝒚𝒙𝒚]\frac{1}{2N}_i\{_t[\widehat{x}(t)\sigma _i(t)+\widehat{y}(t)\widehat{h}_i(t)]\}^2+𝒪(N^1)}$$
$$=\frac{d𝒙d𝒚d\widehat{𝒙}d\widehat{𝒚}}{(2\pi )^{2(t_m+1)}}e^{i[\widehat{𝒙}𝒙+\widehat{𝒚}𝒚𝒙𝒚]\frac{1}{2}_{t,t^{}}\left[\widehat{x}(t)\widehat{x}(t^{})q(t,t^{})+2\widehat{x}(t)\widehat{y}(t^{})K(t^{},t)+\widehat{y}(t)\widehat{y}(t^{})Q(t,t^{})\right]+𝒪(N^1)}$$
Together with (126) we have now shown that the disorder average (125) is indeed, in leading order in $`N`$, of the form (112) (as claimed), with
$$\mathrm{\Phi }[𝒂,𝒌,𝒒,𝑸,𝑲]=i\alpha \underset{t}{}K(t,t)i𝒂𝒌+\alpha \mathrm{log}\frac{d𝒙d𝒚d\widehat{𝒙}d\widehat{𝒚}}{(2\pi )^{2(t_m+1)}}e^{i[\widehat{𝒙}𝒙+\widehat{𝒚}𝒚𝒙𝒚]\frac{1}{2}[\widehat{𝒙}𝒒\widehat{𝒙}+2\widehat{𝒚}𝑲\widehat{𝒙}+\widehat{𝒚}𝑸\widehat{𝒚}]}$$
$$=i\alpha \underset{t}{}K(t,t)i𝒂𝒌+\alpha \mathrm{log}\frac{d𝒖d𝒗}{(2\pi )^{t_m+1}}e^{\frac{1}{2}[𝒖𝒒𝒖+2𝒗𝑲𝒖2i𝒖𝒗+𝒗𝑸𝒗]}$$
(127)
(which, of course, can be simplified further).
Simplification of the Saddle-Point Equations. We are now in a position to work out equations (119,120). For the single-time observables this gives $`\widehat{a}(t)=k(t)`$ and $`\widehat{k}(t)=a(t)`$, and for the two-time ones:
$$\widehat{q}(t,t^{})=\frac{1}{2}\alpha i\frac{𝑑𝒖𝑑𝒗u(t)u(t^{})e^{\frac{1}{2}[𝒖𝒒𝒖+2𝒗𝑲𝒖2i𝒖𝒗+𝒗𝑸𝒗]}}{𝑑𝒖𝑑𝒗e^{\frac{1}{2}[𝒖𝒒𝒖+2𝒗𝑲𝒖2i𝒖𝒗+𝒗𝑸𝒗]}}$$
$$\widehat{Q}(t,t^{})=\frac{1}{2}\alpha i\frac{𝑑𝒖𝑑𝒗v(t)v(t^{})e^{\frac{1}{2}[𝒖𝒒𝒖+2𝒗𝑲𝒖2i𝒖𝒗+𝒗𝑸𝒗]}}{𝑑𝒖𝑑𝒗e^{\frac{1}{2}[𝒖𝒒𝒖+2𝒗𝑲𝒖2i𝒖𝒗+𝒗𝑸𝒗]}}$$
$$\widehat{K}(t,t^{})=\alpha i\frac{𝑑𝒖𝑑𝒗v(t)u(t^{})e^{\frac{1}{2}[𝒖𝒒𝒖+2𝒗𝑲𝒖2i𝒖𝒗+𝒗𝑸𝒗]}}{𝑑𝒖𝑑𝒗e^{\frac{1}{2}[𝒖𝒒𝒖+2𝒗𝑲𝒖2i𝒖𝒗+𝒗𝑸𝒗]}}\alpha \delta _{t,t^{}}$$
At the physical saddle-point we can use (123) to express all non-zero objects in terms of the observables $`m(t)`$, $`C(t,t^{})`$ and $`G(t,t^{})`$, with a clear physical meaning. Thus we find $`\widehat{a}(t)=0`$, $`\widehat{k}(t)=m(t)`$, and
$$\widehat{q}(t,t^{})=\frac{1}{2}\alpha i\frac{𝑑𝒖𝑑𝒗u(t)u(t^{})e^{\frac{1}{2}[𝒖𝑪𝒖2i𝒖[1\mathrm{I}𝑮]𝒗]}}{𝑑𝒖𝑑𝒗e^{\frac{1}{2}[𝒖𝑪𝒖2i𝒖[1\mathrm{I}𝑮]𝒗]}}=0$$
(128)
$$\widehat{Q}(t,t^{})=\frac{1}{2}\alpha i\frac{𝑑𝒖𝑑𝒗v(t)v(t^{})e^{\frac{1}{2}[𝒖𝑪𝒖2i𝒖[1\mathrm{I}𝑮]𝒗]}}{𝑑𝒖𝑑𝒗e^{\frac{1}{2}[𝒖𝑪𝒖2i𝒖[1\mathrm{I}𝑮]𝒗]}}=\frac{1}{2}\alpha i\left[(1\mathrm{I}𝑮)^1𝑪(1\mathrm{I}𝑮^{})^1\right](t,t^{})$$
(129)
$$\widehat{K}(t,t^{})+\alpha \delta _{t,t^{}}=\alpha i\frac{𝑑𝒖𝑑𝒗v(t)u(t^{})e^{\frac{1}{2}[𝒖𝑪𝒖2i𝒖[1\mathrm{I}𝑮]𝒗]}}{𝑑𝒖𝑑𝒗e^{\frac{1}{2}[𝒖𝑪𝒖2i𝒖[1\mathrm{I}𝑮]𝒗]}}=\alpha (1\mathrm{I}𝑮)^1(t,t^{})$$
(130)
(with $`G^{}(t,t^{})=G(t^{},t)`$, and using standard manipulations of Gaussian integrals). Note that we can use the identity $`(1\mathrm{I}𝑮)^11\mathrm{I}=_\mathrm{}0𝑮^{\mathrm{}}1\mathrm{I}=_{\mathrm{}>0}𝑮^{\mathrm{}}=𝑮(1\mathrm{I}𝑮)^1`$ to compactify (130) to
$$\widehat{K}(t,t^{})=\alpha [𝑮(1\mathrm{I}𝑮)^1](t,t^{})$$
(131)
We have now expressed all our objects in terms of the disorder-averaged recall overlap $`𝒎=\{m(t)\}`$ and the disorder-averaged single-site correlation- and response functions $`𝑪=\{C(t,t^{})\}`$ and $`𝑮=\{G(t,t^{})\}`$. We can next simplify the effective measure (124), which plays a crucial role in the remaining saddle-point equations. Inserting $`\widehat{a}(t)=\widehat{q}(t,t^{})=0`$ and $`\widehat{k}(t)=m(t)`$ into (124), first of all, gives us
$$M[\{\sigma \},\{h\},\{\widehat{h}\}]=\pi _0(\sigma (0))\times $$
$$e^{_t\{i\widehat{h}(t)[h(t)m(t)\theta (t)_t^{}\widehat{K}(t,t^{})\sigma (t^{})]+\beta \sigma (t+1)h(t)\mathrm{log}2\mathrm{cosh}[\beta h(t)]\}i_{t,t^{}}\widehat{Q}(t,t^{})\widehat{h}(t)\widehat{h}(t^{})}$$
(132)
Secondly, causality ensures that $`G(t,t^{})=0`$ for $`tt^{}`$, from which, in combination with (131), it follows that the same must be true for the kernel $`\widehat{K}(t,t^{})`$, since
$$\widehat{K}(t,t^{})=\alpha [𝑮(1\mathrm{I}𝑮)^1](t,t^{})=\alpha \left\{𝑮+𝑮^2+𝑮^3+\mathrm{}\right\}(t,t^{})$$
This, in turn, guarantees that the function $`M[\mathrm{}]`$ in (132) is already normalised:
$$\{dhd\widehat{h}\}\underset{\sigma (0)\mathrm{}\sigma (t_m)}{}M[\{\sigma \},\{h\},\{\widehat{h}\}]=1$$
One can prove this iteratively. After summation over $`\sigma (t_m)`$ (which due to causality cannot occur in the term with the kernel $`\widehat{K}(t,t^{})`$) one is left with just a single occurrence of the field $`h(t_m)`$ in the exponent, integration over which reduces to $`\delta [\widehat{h}(t_m)]`$, which then eliminates the conjugate field $`\widehat{h}(t_m)`$. This cycle of operations is next applied to the variables at time $`t_m1`$, etc. The effective measure (124) can now be written simply as
$$f[\{\sigma \},\{h\},\{\widehat{h}\}]_{}=\underset{\sigma (0)\mathrm{}\sigma (t_m)}{}\{dhd\widehat{h}\}M[\{\sigma \},\{h\},\{\widehat{h}\}]f[\{\sigma \},\{h\},\{\widehat{h}\}]$$
with $`M[\mathrm{}]`$ as given in (132). The remaining saddle-point equations to be solved, which can be slightly simplified by using the identity $`\sigma (t)\widehat{h}(t^{})_{}=i\sigma (t)_{}/\theta (t^{})`$, are
$$m(t)=\sigma (t)_{}C(t,t^{})=\sigma (t)\sigma (t^{})_{}G(t,t^{})=\sigma (t)_{}/\theta (t^{})$$
(133)
Extracting the Physics from the Saddle-Point Equations. At this stage we observe in (133) that we only need to insert functions of spin states into the effective measure $`\mathrm{}_{}`$ (rather than fields or conjugate fields), so the effective measure can again be simplified. Upon inserting (129,131) into the function (132) we obtain $`f[\{\sigma \}]_{}=_{\sigma (0)\mathrm{}\sigma (t_m)}\mathrm{Prob}[\{\sigma \}]f[\{\sigma \}]`$, with
$$\mathrm{Prob}[\{\sigma \}]=\pi _0(\sigma (0))\{d\varphi \}P[\{\varphi \}]\underset{t}{}[\frac{1}{2}[1+\sigma (t+1)\mathrm{tanh}[\beta h(t|\{\sigma \},\{\varphi \})]]$$
(134)
in which $`\pi _0(\sigma (0))=\frac{1}{2}[1+\sigma (0)m_0]`$, and
$$h(t|\{\sigma \},\{\varphi \})=m(t)+\theta (t)+\alpha \underset{t^{}<t}{}[𝑮(1\mathrm{I}𝑮)^1](t,t^{})\sigma (t^{})+\alpha ^{\frac{1}{2}}\varphi (t)$$
(135)
$$P[\{\varphi \}]=\frac{e^{\frac{1}{2}_{t,t^{}}\varphi (t)\left[(1\mathrm{I}𝑮^{})𝑪^1(1\mathrm{I}𝑮)\right](t,t^{})\varphi (t^{})}}{(2\pi )^{(t_m+1)/2}\mathrm{det}^{\frac{1}{2}}\left[(1\mathrm{I}𝑮^{})𝑪^1(1\mathrm{I}𝑮)\right]}$$
(136)
(note: to predict neuron states up until time $`t_m`$ we only need the fields up until time $`t_m1`$). We recognise (134) as describing an effective single neuron, with the usual dynamics $`\mathrm{Prob}[\sigma (t+1)=\pm 1]=\frac{1}{2}[1\pm \mathrm{tanh}[\beta h(t)]]`$, but with the fields (135). This result is indeed of the form (98), with a retarded self-interaction kernel $`R(t,t^{})`$ and covariance matrix $`\varphi (t)\varphi (t^{})`$ of the Gaussian $`\varphi (t)`$ given by
$$R(t,t^{})=[𝑮(1\mathrm{I}𝑮)^1](t,t^{})\varphi (t)\varphi (t^{})=[(1\mathrm{I}𝑮)^1𝑪(1\mathrm{I}𝑮^{})^1](t,t^{})$$
(137)
For $`\alpha 0`$ we loose all the complicated terms in the local fields, and recover the type of simple expression we found earlier for finite $`p`$: $`m(t+1)=\mathrm{tanh}[\beta (m(t)+\theta (t))]`$.
It can be shown (space limitations prevent a demonstration in this paper) that the equilibrium solutions obtained via replica theory in replica-symmetric ansatz can be recovered as those time-translation invariant solutions<sup>8</sup><sup>8</sup>8i.e. $`m(t)=m`$, $`C(t,t^{})=C(tt^{})`$ and $`G(t,t^{})=G(tt^{})`$ of the above dynamic equations which (i) obey the parallel dynamics fluctuation-dissipation theorem, and (ii) obey $`lim_\tau \mathrm{}G(\tau )=0`$. It can also be shown that the AT instability, where replica symmetry ceases to hold, corresponds to a dynamical instability in the present formalism, where so-called anomalous response sets in: $`lim_\tau \mathrm{}G(\tau )0`$.
Before we calculate the solution explicitly for the first few time-steps, we first work out the relevant averages using (134). Note that always $`C(t,t)=\sigma ^2(t)_{}=1`$ and $`G(t,t^{})=R(t,t^{})=0`$ for $`tt^{}`$. As a result the covariance matrix of the Gaussian fields can be written as
$$\varphi (t)\varphi (t^{})=[(1\mathrm{I}𝑮)^1𝑪(1\mathrm{I}𝑮^{})^1](t,t^{})=\underset{s,s^{}0}{}[\delta _{t,s}+R(t,s)]C(s,s^{})[\delta _{s^{},t^{}}+R(t^{},s^{})]$$
$$=\underset{s=0}{\overset{t}{}}\underset{s^{}=0}{\overset{t^{}}{}}[\delta _{t,s}+R(t,s)]C(s,s^{})[\delta _{s^{},t^{}}+R(t^{},s^{})]$$
(138)
Considering arbitrary positive integer powers of the response function immediately shows that
$$(𝑮^{\mathrm{}})(t,t^{})=0\mathrm{if}t^{}>t\mathrm{}$$
(139)
which, in turn, gives
$$R(t,t^{})=\underset{\mathrm{}>0}{}(𝑮^{\mathrm{}})(t,t^{})=\underset{\mathrm{}=1}{\overset{tt^{}}{}}(𝑮^{\mathrm{}})(t,t^{})$$
(140)
Similarly we obtain from $`(1\mathrm{I}𝑮)^1=1\mathrm{I}+𝑹`$ that for $`t^{}t`$: $`(1\mathrm{I}𝑮)^1(t,t^{})=\delta _{t,t^{}}`$. To suppress notation we will simply put $`h(t|..)`$ instead of $`h(t|\{\sigma \},\{\varphi \})`$; this need not cause any ambiguity. We notice that summation over neuron variables $`\sigma (s)`$ and integration over Gaussian variables $`\varphi (s)`$ with time arguments $`s`$ higher than than those occurring in the function to be averaged can always be carried out immediately, giving (for $`t>0`$ and $`t^{}<t`$):
$$m(t)=\underset{\sigma (0)\mathrm{}\sigma (t1)}{}\pi _0(\sigma (0))\{d\varphi \}P[\{\varphi \}]\mathrm{tanh}[\beta h(t1|..)]\underset{s=0}{\overset{t2}{}}\frac{1}{2}[1+\sigma (s+1)\mathrm{tanh}[\beta h(s|..)]]$$
(141)
$$G(t,t^{})=\beta \{C(t,t^{}+1)\underset{\sigma (0)\mathrm{}\sigma (t1)}{}\pi _0(\sigma (0))\{d\varphi \}P[\{\varphi \}]\mathrm{tanh}[\beta h(t1|..)]\mathrm{tanh}[\beta h(t^{}|..)]$$
$$\times \underset{s=0}{\overset{t2}{}}\frac{1}{2}[1+\sigma (s+1)\mathrm{tanh}[\beta h(s|..)]]\}$$
(142)
(which we obtain directly for $`t^{}=t1`$, and which follows for times $`t^{}<t1`$ upon using the identity $`\sigma [1\mathrm{tanh}^2(x)]=[1+\sigma \mathrm{tanh}(x)][\sigma \mathrm{tanh}(x)]`$). For the correlations we distinguish between $`t^{}=t1`$ and $`t^{}<t1`$:
$$C(t,t1)=\underset{\sigma (0)\mathrm{}\sigma (t2)}{}\pi _0(\sigma (0))\{d\varphi \}P[\{\varphi \}]\mathrm{tanh}[\beta h(t1|..)]\mathrm{tanh}[\beta h(t2|..)]\underset{s=0}{\overset{t3}{}}\frac{1}{2}[1+\sigma (s+1)\mathrm{tanh}[\beta h(s|..)]]$$
(143)
whereas for $`t^{}<t1`$ we have
$$C(t,t^{})=\underset{\sigma (0)\mathrm{}\sigma (t1)}{}\pi _0(\sigma (0))\{d\varphi \}P[\{\varphi \}]\mathrm{tanh}[\beta h(t1|..)]\sigma (t^{})\underset{s=0}{\overset{t2}{}}\frac{1}{2}[1+\sigma (s+1)\mathrm{tanh}[\beta h(s|..)]]$$
(144)
Let us finally work out explicitly the final macroscopic laws (141,142,143,144), with (135,136), for the first few time-steps. For arbitrary times our equations will have to be evaluated numerically; we will see below, however, that this can be done in an iterative (i.e. easy) manner. At $`t=0`$ we just have the two observables $`m(0)=m_0`$ and $`C(0,0)=1`$.
The First Few Time-Steps. The field at $`t=0`$ is $`h(0|..)=m_0+\theta (0)+\alpha ^{\frac{1}{2}}\varphi (0)`$, since the retarded self-interaction does not yet come into play. The distribution of $`\varphi (0)`$ is fully characterised by its variance, which (138) claims to be
$$\varphi ^2(0)=C(0,0)=1$$
Therefore, with $`Dz=(2\pi )^{\frac{1}{2}}e^{\frac{1}{2}z^2}dz`$, we immediately find (141,142,143,144) reducing to
$$m(1)=Dz\mathrm{tanh}[\beta (m_0+\theta (0)+z\sqrt{\alpha })]C(1,0)=m_0m(1)$$
(145)
$$G(1,0)=\beta \left\{1Dz\mathrm{tanh}^2[\beta (m_0+\theta (0)+z\sqrt{\alpha })]\right\}$$
(146)
For the self-interaction kernel this implies, using (140), that $`R(1,0)=G(1,0)`$. We now move on to $`t=2`$. Here equations (141,142,143,144) give us
$$m(2)=\frac{1}{2}\underset{\sigma (0)}{}d\varphi (0)d\varphi (1)P[\varphi (0),\varphi (1)]\mathrm{tanh}[\beta h(1|..)][1+\sigma (0)m_0]$$
$$C(2,1)=\frac{1}{2}\underset{\sigma (0)}{}d\varphi (1)d\varphi (0)P[\varphi (0),\varphi (1)]\mathrm{tanh}[\beta h(1|..)]\mathrm{tanh}[\beta h(0|..)][1+\sigma (0)m_0]$$
$$C(2,0)=\frac{1}{2}\underset{\sigma (0)\sigma (1)}{}\{d\varphi \}P[\{\varphi \}]\mathrm{tanh}[\beta h(1|..)]\sigma (0)\frac{1}{2}[1+\sigma (1)\mathrm{tanh}[\beta h(0|..)]][1+\sigma (0)m_0]$$
$$G(2,1)=\beta \{1\frac{1}{2}\underset{\sigma (0)}{}d\varphi (0)d\varphi (1)P[\varphi (0),\varphi (1)]\mathrm{tanh}^2[\beta h(1|..)][1+\sigma (0)m_0]\}$$
$$G(2,0)=\beta \{C(2,1)\frac{1}{2}\underset{\sigma (0)}{}d\varphi (0)d\varphi (1)P[\varphi (0),\varphi (1)]\mathrm{tanh}[\beta h(1|..)]\mathrm{tanh}[\beta h(0|..)][1+\sigma (0)m_0]\}=0$$
We already know that $`\varphi ^2(0)=1`$; the remaining two moments we need in order to determine $`P[\varphi (0),\varphi (1)]`$ follow again from (138):
$$\varphi (1)\varphi (0)=\underset{s=0}{\overset{1}{}}[\delta _{1,s}+\delta _{0,s}R(1,0)]C(s,0)=C(1,0)+G(1,0)$$
$$\varphi ^2(1)=\underset{s=0}{\overset{1}{}}\underset{s^{}=1}{\overset{1}{}}[\delta _{1,s}+\delta _{0,s}R(1,0)]C(s,s^{})[\delta _{s^{},1}+\delta _{s^{},0}R(1,0)]=G^2(1,0)+2C(0,1)G(1,0)+1$$
We now know $`P[\varphi (0),\varphi (1)]`$ and can work out all macroscopic objects with $`t=2`$ explicitly, if we wish. I will not do this here in full, but only point at the emerging pattern of all calculations at a given time $`t`$ depending only on macroscopic quantities that have been calculated at times $`t^{}<t`$, which allows for iterative solution. Let us just work out $`m(2)`$ explicitly, in order to compare the first two recall overlaps $`m(1)`$ and $`m(2)`$ with the values found in simulations and in approximate theories. We note that calculating $`m(2)`$ only requires the field $`\varphi (1)`$, for which we found $`\varphi ^2(1)=G^2(1,0)+2C(0,1)G(1,0)+1`$:
$$m(2)=\frac{1}{2}\underset{\sigma (0)}{}𝑑\varphi (1)P[\varphi (1)]\mathrm{tanh}[\beta (m(1)+\theta (1)+\alpha G(1,0)\sigma (0)+\alpha ^{\frac{1}{2}}\varphi (1))][1+\sigma (0)m_0]$$
$$=\frac{1}{2}[1+m_0]Dz\mathrm{tanh}[\beta (m(1)+\theta (1)+\alpha G(1,0)+z\sqrt{\alpha [G^2(1,0)+2m_0m(1)G(1,0)+1]})]$$
$$+\frac{1}{2}[1m_0]Dz\mathrm{tanh}[\beta (m(1)+\theta (1)\alpha G(1,0)+z\sqrt{\alpha [G^2(1,0)+2m_0m(1)G(1,0)+1]})]$$
Exact Results Versus Simulations and Gaussian Approximations. I close this section on the fully connected networks with a comparison of some of the approximate theories, the (exact) generating functional formalism, and numerical simulations, for the case $`\theta (t)=0`$ (no external stimuli at any time). The evolution of the recall overlap in the first two time-steps has been described as follows:
$$\begin{array}{cccc}\mathrm{𝑁𝑎𝑖𝑣𝑒}\mathrm{𝐺𝑎𝑢𝑠𝑠𝑖𝑎𝑛}\mathrm{𝐴𝑝𝑝𝑟𝑜𝑥𝑖𝑚𝑎𝑡𝑖𝑜𝑛}:\hfill & & m(1)=\hfill & Dz\mathrm{tanh}[\beta (m(0)+z\sqrt{\alpha })]\hfill \\ & & m(2)=\hfill & Dz\mathrm{tanh}[\beta (m(1)+z\sqrt{\alpha })]\hfill \\ \mathrm{𝐴𝑚𝑎𝑟𝑖}\mathrm{𝑀𝑎𝑔𝑖𝑛𝑢}\mathrm{𝑇ℎ𝑒𝑜𝑟𝑦}:\hfill & & m(1)=\hfill & Dz\mathrm{tanh}[\beta (m(0)+z\sqrt{\alpha })]\hfill \\ & & m(2)=\hfill & Dz\mathrm{tanh}[\beta (m(1)+z\mathrm{\Sigma }\sqrt{\alpha })]\hfill \\ & & \mathrm{\Sigma }^2=\hfill & 1+2m(0)m(1)G+G^2\hfill \\ & & G=\hfill & \beta \left[1Dz\mathrm{tanh}^2[\beta (m(0)+z\sqrt{\alpha })]\right]\hfill \\ \mathrm{𝐸𝑥𝑎𝑐𝑡}\mathrm{𝑆𝑜𝑙𝑢𝑡𝑖𝑜𝑛}:\hfill & & m(1)=\hfill & Dz\mathrm{tanh}[\beta (m(0)+z\sqrt{\alpha })]\hfill \\ & & m(2)=\hfill & \frac{1}{2}[1+m_0]Dz\mathrm{tanh}[\beta (m(1)+\alpha G+z\mathrm{\Sigma }\sqrt{\alpha })]\hfill \\ & & & +\frac{1}{2}[1m_0]Dz\mathrm{tanh}[\beta (m(1)\alpha G+z\mathrm{\Sigma }\sqrt{\alpha })]\hfill \\ & & \mathrm{\Sigma }^2=\hfill & 1+2m(0)m(1)G+G^2\hfill \\ & & G=\hfill & \beta \left[1Dz\mathrm{tanh}^2[\beta (m(0)+z\sqrt{\alpha })]\right]\hfill \end{array}$$
We can now appreciate why the more advanced Gaussian approximation (Amari-Maginu theory, ) works well when the system state is close to the target attractor. This theory gets the moments of the Gaussian part of the interference noise distribution at $`t=1`$ exactly right, but not the discrete part, whereas close to the attractor both the response function $`G(1,0)`$ and one of the two pre-factors $`\frac{1}{2}[1\pm m_0]`$ in the exact expression for $`m(2)`$ will be very small, and the latter will therefore indeed approach a Gaussian shape. One can also see why the non-Gaussian approximation of made sense: in the calculation of $`m(2)`$ the interference noise distribution can indeed be written as the sum of two Gaussian ones (although for $`t>2`$ this will cease to be true). Numerical evaluation of these expressions result in explicit predictions which can be tested against numerical simulations. This is done in figure 8, which confirms the picture sketched above, and hints that the performance of the Gaussian approximations is indeed worse for those initial conditions which fail to trigger pattern recall.
### 5.4 Extremely Diluted Attractor Networks Near Saturation
Extremely diluted attractor networks are obtained upon choosing $`lim_N\mathrm{}c/N=0`$ (while still $`c\mathrm{}`$) in definition (100) of the Hebbian-type synapses. The disorder average now involves both the patterns with $`\mu >1`$ and the realisation of the ‘wiring’ variables $`c_{ij}\{0,1\}`$. Again, in working out the key function (125) we will show that for $`N\mathrm{}`$ the outcome can be written in terms of the macroscopic quantities (109,110,111). We carry out the average over the spatial structure variables $`\{c_{ij}\}`$ first:
$$[\mathrm{}]=\frac{1}{N}\mathrm{log}\overline{\left[e^{\frac{i}{c}_{ij}c_{ij}_\mu \xi _i^\mu \xi _j^\mu _t\widehat{h}_i(t)\sigma _j(t)}\right]}=\frac{1}{N}\mathrm{log}\overline{\underset{i<j}{}e^{\frac{i}{c}_\mu \xi _i^\mu \xi _j^\mu [c_{ij}_t\widehat{h}_i(t)\sigma _j(t)+c_{ji}_t\widehat{h}_j(t)\sigma _i(t)]}}$$
At this stage we have to distinguish between symmetric and asymmetric dilution.
The Disorder Average. First we deal with the case of symmetric dilution: $`c_{ij}=c_{ji}`$ for all $`ij`$. The average over the $`c_{ij}`$, with the distribution (101), is trivial:
$$\overline{\underset{i<j}{}e^{\frac{i}{c}c_{ij}_\mu \xi _i^\mu \xi _j^\mu _t[\widehat{h}_i(t)\sigma _j(t)+\widehat{h}_j(t)\sigma _i(t)]}}=\overline{\underset{i<j}{}\left\{1+\frac{c}{N}[e^{\frac{i}{c}_\mu \xi _i^\mu \xi _j^\mu _t[\widehat{h}_i(t)\sigma _j(t)+\widehat{h}_j(t)\sigma _i(t)]}1]\right\}}$$
$$=\overline{\underset{i<j}{}\left\{1\frac{c}{N}\left[\frac{i}{c}\underset{\mu }{}\xi _i^\mu \xi _j^\mu \underset{t}{}[\widehat{h}_i(t)\sigma _j(t)+\widehat{h}_j(t)\sigma _i(t)]+\frac{1}{2c^2}[\underset{\mu }{}\xi _i^\mu \xi _j^\mu \underset{t}{}[\widehat{h}_i(t)\sigma _j(t)+\widehat{h}_j(t)\sigma _i(t)]]^2+𝒪(c^{\frac{3}{2}})\right]\right\}}$$
$$=\overline{\underset{i<j}{}e^{\frac{i}{N}_\mu \xi _i^\mu \xi _j^\mu _t[\widehat{h}_i(t)\sigma _j(t)+\widehat{h}_j(t)\sigma _i(t)]\frac{1}{2cN}[_\mu \xi _i^\mu \xi _j^\mu _t[\widehat{h}_i(t)\sigma _j(t)+\widehat{h}_j(t)\sigma _i(t)]]^2+𝒪(\frac{1}{N\sqrt{c}})+𝒪(\frac{c}{N^2})}}$$
We separate in the exponent the terms where $`\mu =\nu `$ in the quadratic term (being of the form $`_{\mu \nu }\mathrm{}`$), and the terms with $`\mu =1`$. Note: $`p=\alpha c`$. We also use the definitions (109,110,111) wherever we can:
$$[\mathrm{}]=i\underset{t}{}a(t)k(t)\frac{1}{2}\alpha \underset{st}{}[q(s,t)Q(s,t)+K(s,t)K(t,s)]+𝒪(c^{\frac{1}{2}})+𝒪(c/N)+$$
$$\frac{1}{N}\mathrm{log}\left\{\overline{e^{\frac{i}{N}_{\mu >1}_t[_i\xi _i^\mu \widehat{h}_i(t)][_j\xi _j^\mu \sigma _j(t)]\frac{1}{4cN}_{ij}_{\mu \nu }_{st}\xi _i^\mu \xi _j^\mu \xi _i^\nu \xi _j^\nu [\widehat{h}_i(s)\sigma _j(s)+\widehat{h}_j(s)\sigma _i(s)][\widehat{h}_i(t)\sigma _j(t)+\widehat{h}_j(t)\sigma _i(t)]}}\right\}$$
Our ‘condensed ansatz’ implies that for $`\mu >1`$: $`N^{\frac{1}{2}}_i\xi _i^\mu \sigma _i(t)=𝒪(1)`$ and $`N^{\frac{1}{2}}_i\xi _i^\mu \widehat{h}_i(t)=𝒪(1)`$. Thus the first term in the exponent containing the disorder is $`𝒪(c)`$, contributing $`𝒪(c/N)`$ to $`[\mathrm{}]`$. We therefore retain only the second term in the exponent. However, the same argument applies to the second term. There all contributions can be seen as uncorrelated in leading order, so that $`_{ij}_{\mu \nu }\mathrm{}=𝒪(Np)`$, giving a non-leading $`𝒪(N^1)`$ cumulative contribution to $`[\mathrm{}]`$. Thus, provided $`lim_N\mathrm{}c^1=lim_N\mathrm{}c/N=0`$ (which we assumed), we have shown that the disorder average (125) is again, in leading order in $`N`$, of the form (112) (as claimed), with
$$\mathrm{𝑆𝑦𝑚𝑚𝑒𝑡𝑟𝑖𝑐}:\mathrm{\Phi }[𝒂,𝒌,𝒒,𝑸,𝑲]=i𝒂𝒌\frac{1}{2}\alpha \underset{st}{}[q(s,t)Q(s,t)+K(s,t)K(t,s)]$$
(147)
Next we deal with the asymmetric case (102), where $`c_{ij}`$ and $`c_{ji}`$ are independent. Again the average over the $`c_{ij}`$ is trivial; here it gives
$$\overline{\underset{i<j}{}\left\{e^{\frac{i}{c}c_{ij}_\mu \xi _i^\mu \xi _j^\mu _t\widehat{h}_i(t)\sigma _j(t)}e^{\frac{i}{c}c_{ji}_\mu \xi _i^\mu \xi _j^\mu _t\widehat{h}_j(t)\sigma _i(t)}\right\}}$$
$$=\overline{\underset{i<j}{}\left\{1+\frac{c}{N}[e^{\frac{i}{c}_\mu \xi _i^\mu \xi _j^\mu _t\widehat{h}_i(t)\sigma _j(t)}1]\right\}\left\{1+\frac{c}{N}[e^{\frac{i}{c}_\mu \xi _i^\mu \xi _j^\mu _t\widehat{h}_j(t)\sigma _i(t)}1]\right\}}$$
$$=\overline{\underset{i<j}{}\left\{1\frac{c}{N}[\frac{i}{c}\underset{\mu }{}\xi _i^\mu \xi _j^\mu \underset{t}{}\widehat{h}_i(t)\sigma _j(t)+\frac{1}{2c^2}[\underset{\mu }{}\xi _i^\mu \xi _j^\mu \underset{t}{}\widehat{h}_i(t)\sigma _j(t)]^2+𝒪(c^{\frac{3}{2}})]\right\}}$$
$$\overline{\times \left\{1\frac{c}{N}[\frac{i}{c}\underset{\mu }{}\xi _i^\mu \xi _j^\mu \underset{t}{}\widehat{h}_j(t)\sigma _i(t)+\frac{1}{2c^2}[\underset{\mu }{}\xi _i^\mu \xi _j^\mu \underset{t}{}\widehat{h}_j(t)\sigma _i(t)]^2+𝒪(c^{\frac{3}{2}})]\right\}}$$
(in which the horizontal bars of the two constituent lines are to be read as connected)
$$=\overline{\underset{i<j}{}e^{\frac{i}{N}_\mu \xi _i^\mu \xi _j^\mu _t[\widehat{h}_i(t)\sigma _j(t)+\widehat{h}_j(t)\sigma _i(t)]\frac{1}{2cN}[_\mu \xi _i^\mu \xi _j^\mu _t\widehat{h}_i(t)\sigma _j(t)]^2\frac{1}{2cN}[_\mu \xi _i^\mu \xi _j^\mu _t\widehat{h}_j(t)\sigma _i(t)]^2+𝒪(\frac{1}{N\sqrt{c}})+𝒪(\frac{c}{N^2})}}$$
Again we separate in the exponent the terms where $`\mu =\nu `$ in the quadratic term (being of the form $`_{\mu \nu }\mathrm{}`$), and the terms with $`\mu =1`$, and use the definitions (109,110,111):
$$[\mathrm{}]=i\underset{t}{}a(t)k(t)\frac{1}{2}\alpha \underset{st}{}q(s,t)Q(s,t)+𝒪(c^{\frac{1}{2}})+𝒪(c/n)$$
$$+\frac{1}{N}\mathrm{log}\left\{\overline{e^{\frac{i}{N}_{\mu >1}_t[_i\xi _i^\mu \widehat{h}_i(t)][_j\xi _j^\mu \sigma _j(t)]\frac{1}{2cN}_{ij}_{\mu \nu }\xi _i^\mu \xi _j^\mu \xi _i^\nu \xi _j^\nu _{st}\widehat{h}_i(s)\sigma _j(s)\widehat{h}_i(t)\sigma _j(t)}}\right\}$$
The scaling arguments given in the symmetric case, based on our ‘condensed ansatz’, apply again, and tell us that the remaining terms with the disorder are of vanishing order in $`N`$. We have again shown that the disorder average (125) is, in leading order in $`N`$, of the form (112), with
$$\mathrm{𝐴𝑠𝑦𝑚𝑚𝑒𝑡𝑟𝑖𝑐}:\mathrm{\Phi }[𝒂,𝒌,𝒒,𝑸,𝑲]=i𝒂𝒌\frac{1}{2}\alpha \underset{st}{}q(s,t)Q(s,t)$$
(148)
Extracting the Physics from the Saddle-Point Equations. First we combine the above two results (147,148) in the following way (with $`\mathrm{\Delta }=1`$ for symmetric dilution and $`\mathrm{\Delta }=0`$ for asymmetric dilution):
$$\mathrm{\Phi }[𝒂,𝒌,𝒒,𝑸,𝑲]=i𝒂𝒌\frac{1}{2}\alpha \underset{st}{}[q(s,t)Q(s,t)+\mathrm{\Delta }K(s,t)K(t,s)]$$
(149)
We can now work out equations (119,120), and use (123) to express the result at the physical saddle-point in terms of the trio $`\{m(t),C(t,t^{}),G(t,t^{})\}`$. For the single-time observables this gives (as with the fully connected system) $`\widehat{a}(t)=k(t)`$ and $`\widehat{k}(t)=a(t)`$; for the two-time ones we find:
$$\widehat{Q}(t,t^{})=\frac{1}{2}i\alpha C(t,t^{})\widehat{q}(t,t^{})=0\widehat{K}(t,t^{})=\alpha \mathrm{\Delta }G(t,t^{})$$
We now observe that the remainder of the derivation followed for the fully connected network can be followed with only two minor adjustments to the terms generated by $`\widehat{K}(t,t^{})`$ and by $`\widehat{Q}(t,t^{})`$: $`\alpha 𝑮(1\mathrm{I}𝑮)^1\alpha \mathrm{\Delta }𝑮`$ in the retarded self-interaction, and $`(1\mathrm{I}𝑮)^1𝑪(1\mathrm{I}𝑮^{})^1𝑪`$ in the covariance of the Gaussian noise in the effective single neuron problem. This results in the familiar saddle-point equations (133) for an effective single neuron problem, with state probabilities (134) equivalent to the dynamics $`\mathrm{Prob}[\sigma (t+1)=\pm 1]=\frac{1}{2}[1\pm \mathrm{tanh}[\beta h(t)]]`$, and in which $`\pi _0(\sigma (0))=\frac{1}{2}[1+\sigma (0)m_0]`$ and
$$h(t|\{\sigma \},\{\varphi \})=m(t)+\theta (t)+\alpha \mathrm{\Delta }\underset{t^{}<t}{}G(t,t^{})\sigma (t^{})+\alpha ^{\frac{1}{2}}\varphi (t)P[\{\varphi \}]=\frac{e^{\frac{1}{2}_{t,t^{}}\varphi (t)𝑪^1(t,t^{})\varphi (t^{})}}{(2\pi )^{(t_m+1)/2}\mathrm{det}^{\frac{1}{2}}𝑪}$$
(150)
Physics of Networks with Asymmetric Dilution. Asymmetric dilution corresponds to $`\mathrm{\Delta }=0`$, i.e. there is no retarded self-interaction, and the response function no longer plays a role. In (150) we now only retain $`h(t|\mathrm{})=m(t)+\theta (t)+\alpha ^{\frac{1}{2}}\varphi (t)`$, with $`\varphi ^2(t)=C(1,1)=1`$. We now find (141) simply giving
$$m(t+1)=\underset{\sigma (0)\mathrm{}\sigma (t)}{}\pi _0(\sigma (0))\{d\varphi \}P[\{\varphi \}]\mathrm{tanh}[\beta h(t|\mathrm{})]\underset{s=0}{\overset{t1}{}}\frac{1}{2}\left[1+\sigma (s+1)\mathrm{tanh}[\beta h(s|\mathrm{})]\right]$$
$$=Dz\mathrm{tanh}[\beta (m(t)+\theta (t)+z\sqrt{\alpha })]$$
(151)
Apparently this is the one case where the simple Gaussian dynamical law (95) is exact at all times. Similarly, for $`t>t^{}`$ equations (142,143,144) for correlation and response functions reduce to
$$C(t,t^{})=$$
$$\frac{d\varphi _ad\varphi _be^{\frac{1}{2}\frac{\varphi _a^2+\varphi _b^22C(t1,t^{}1)\varphi _a\varphi _b}{1C^2(t1,t^{}1)}}}{2\pi \sqrt{1C^2(t1,t^{}1)}}\mathrm{tanh}[\beta (m(t1)+\theta (t1)+\varphi _a\sqrt{\alpha })]\mathrm{tanh}[\beta (m(t^{}1)+\theta (t^{}1)+\varphi _b\sqrt{\alpha })]$$
(152)
$$G(t,t^{})=\beta \delta _{t,t^{}+1}\left\{1Dz\mathrm{tanh}^2[\beta (m(t1)+\theta (t1)+z\sqrt{\alpha })]\right\}$$
(153)
Let us also inspect the stationary state $`m(t)=m`$, for $`\theta (t)=0`$. One easily proves that $`m=0`$ as soon as $`T>1`$, using $`m^2=\beta m_0^m𝑑k[1Dz\mathrm{tanh}^2[\beta (k+z\sqrt{\alpha })]]\beta m^2`$. A continuous bifurcation occurs from the $`m=0`$ state to an $`m>0`$ state when $`T=1Dz\mathrm{tanh}^2[\beta z\sqrt{\alpha }]`$. A parametrisation of this transition line in the $`(\alpha ,T)`$-plane is given by
$$T(x)=1Dz\mathrm{tanh}^2(zx),\alpha (x)=x^2T^2(x),x0$$
For $`\alpha =0`$ we just jet $`m=\mathrm{tanh}(\beta m)`$ so $`T_c=1`$. For $`T=0`$ we obtain the equation $`m=\mathrm{erf}[m/\sqrt{2\alpha }]`$, giving a continuous transition to $`m>0`$ solutions at $`\alpha _c=2/\pi 0.637`$. The remaining question concerns the nature of the $`m=0`$ state. Inserting $`m(t)=\theta (t)=0`$ (for all $`t`$) into (152) tells us that $`C(t,t^{})=f[C(t1,t^{}1)]`$ for $`t>t^{}>0`$, with ‘initial conditions’ $`C(t,0)=m(t)m_0`$, where
$$f[C]=\frac{d\varphi _ad\varphi _b}{2\pi \sqrt{1C^2}}e^{\frac{1}{2}\frac{\varphi _a^2+\varphi _b^22C\varphi _a\varphi _b}{1C^2}}\mathrm{tanh}[\beta \sqrt{\alpha }\varphi _a]\mathrm{tanh}[\beta \sqrt{\alpha }\varphi _b]$$
In the $`m=0`$ regime we have $`C(t,0)=0`$ for any $`t>0`$, inducing $`C(t,t^{})=0`$ for any $`t>t^{}`$, due to $`f[0]=0`$. Thus we conclude that $`C(t,t^{})=\delta _{t,t^{}}`$ in the $`m=0`$ phase, i.e. this phase is para-magnetic rather than of a spin-glass type. The resulting phase diagram is given in figure 9, together with that of symmetric dilution (for comparison).
Physics of Networks with Symmetric Dilution. This is the more complicated situation. In spite of the extreme dilution, the interaction symmetry makes sure that the spins still have a sufficient number of common ancestors for complicated correlations to build up in finite time. We have
$$h(t|\{\sigma \},\{\varphi \})=m(t)+\theta (t)+\alpha \underset{t^{}<t}{}G(t,t^{})\sigma (t^{})+\alpha ^{\frac{1}{2}}\varphi (t)P[\{\varphi \}]=\frac{e^{\frac{1}{2}_{t,t^{}}\varphi (t)𝑪^1(t,t^{})\varphi (t^{})}}{(2\pi )^{(t_m+1)/2}\mathrm{det}^{\frac{1}{2}}𝑪}$$
(154)
The effective single neuron problem (134,154) is found to be exactly of the form found also for the Gaussian model in (which, in turn, maps onto the parallel dynamics SK model ) with the synapses $`J_{ij}=J_0\xi _i\xi _j/N+Jz_{ij}/\sqrt{N}`$ (in which the $`z_{ij}`$ are symmetric zero-average and unit-variance Gaussian variables, and $`J_{ii}=0`$ for all $`i`$), with the identification:
$$J\sqrt{\alpha }J_01$$
(this becomes clear upon applying the generating functional analysis to the Gaussian model, page limitations prevent me from explicit demonstration here). Since one can show that for $`J_0>0`$ the parallel dynamics SK model gives the same equilibrium state as the sequential one, we can now immediately write down the stationary solution of our dynamic equations which corresponds to the FDT regime, with $`q=lim_\tau \mathrm{}lim_t\mathrm{}C(t,t+\tau )`$:
$$q=Dz\mathrm{tanh}^2[\beta (m+z\sqrt{\alpha q})]m=Dz\mathrm{tanh}[\beta (m+z\sqrt{\alpha q})]$$
(155)
These are neither identical to the equations for the fully connected Hopfield model, nor to those of the asymmetrically diluted model. Using the equivalence with the (sequential and parallel) SK model we can immediately translate the phase transition lines as well, giving:
$$\begin{array}{cccc}& \mathrm{𝑆𝐾}\mathrm{𝑚𝑜𝑑𝑒𝑙}& & \mathrm{𝑆𝑦𝑚𝑚𝑒𝑡𝑟𝑖𝑐𝑎𝑙𝑙𝑦}\mathrm{𝐷𝑖𝑙𝑢𝑡𝑒𝑑}\mathrm{𝑀𝑜𝑑𝑒𝑙}\\ PF:\hfill & T=J_0\mathrm{for}J_0>J& & T=1\mathrm{for}\alpha <1\\ PSG:\hfill & T=J\mathrm{for}J_0<J& & T=\sqrt{\alpha }\mathrm{for}\alpha >1\\ FSG(\mathrm{in}\mathrm{RS}):\hfill & T=J_0(1q)\mathrm{for}T<J_0& & T=1q\mathrm{for}T<1\\ FSG(\mathrm{in}\mathrm{RSB}):\hfill & J_0=J\mathrm{for}T<J& & \alpha =1\mathrm{for}T<\sqrt{\alpha }\\ AT\mathrm{line}:\hfill & T^2=J^2Dz\mathrm{cosh}^4\beta [J_0m+Jz\sqrt{q}]& & T^2=\alpha Dz\mathrm{cosh}^4\beta [m+z\sqrt{\alpha q}]\end{array}$$
where $`q=Dz\mathrm{tanh}^2\beta [m+z\sqrt{\alpha q}]`$. Note that for $`T=0`$ we have $`q=1`$, so that the equation for $`m`$ reduces to the one found for asymmetric dilution: $`m=\mathrm{erf}[m/\sqrt{2\alpha }]`$. However, the phase diagram shows that the line $`FSG`$ is entirely in the RSB region and describes physically unrealistic re-entrance (as in the SK model), so that the true transition must be calculated using Parisi’s replica-symmetry breaking (RSB) formalism (see e.g. ), giving here $`\alpha _c=1`$.
The extremely diluted models analysed here were first studied in (asymmetric dilution) and (symmetric dilution). We note that it is not extreme dilution which is responsible for a drastic simplification in the macroscopic dynamics in the complex regime (i.e. close to saturation), but rather the absence of synaptic symmetry. Any finite degree of synaptic symmetry, whether in a fully connected or in an extremely diluted attractor network, immediately generates an effective retarded self-interaction in the dynamics, which is ultimately responsible for highly non-trivial ‘glassy’ dynamics.
## 6 Epilogue
In this paper I have tried to explain how the techniques from non-equilibrium statistical mechanics can be used to solve the dynamics of recurrent neural networks. As in the companion paper on statics in this volume, I have restricted myself to relatively simple models, where one can most clearly see the potential and restrictions of these techniques, without being distracted by details. I have dealt with binary neurons and graded response neurons, and with fully connected and extremely diluted networks, with symmetric but also with non-symmetric synapses. Similar calculations could have been done for neuron models which are not based on firing rates, such as coupled oscillators or integrate-and-fire type ones, see e.g. . My hope is that bringing together methods and results that have so far been mostly scattered over research papers, and by presenting these in a uniform language to simplify comparison, I will have made the area somewhat more accessible to the interested outsider.
At another level I hope to have compensated somewhat for the incorrect view that has sometimes surfaced in the past that statistical mechanics applies only to recurrent networks with symmetric synapses, and is therefore not likely to have a lasting impact on neuro-biological modeling. This was indeed true for equilibrium statistical mechanics, but it is not true for non-equilibrium statistical mechanics. This does not mean that there are no practical restrictions in the latter; the golden rule of there not being any free lunches is obviously also valid here. Whenever we wish to incorporate more biological details in our models, we will have to reduce our ambition to obtain exact solutions, work much harder, and turn to our computer at an earlier stage. However, the practical restrictions in dynamics are of a quantitative nature (equations tend to become more lengthy and messy), rather than of a qualitative one (in statics the issue of detailed balance decides whether or not we can at all start a calculation). The main stumbling block that remains is the issue of spatial structure. Short-range models are extremely difficult to handle, and this is likely to remain so for a long time. In statistical mechanics the state of the art in short-range models is to be able to identify phase transitions, and calculate critical exponents, but this is generally not the type of information one is interested in when studying the operation of recurrent neural networks.
Yet, since dynamical techniques are still far less hampered by the need to impose biologically dubious (or even unacceptable) model constraints than equilibrium techniques, and since there are now well-established and efficient methods and techniques to obtain model solutions in the form of macroscopic laws for large systems (some are exact, some are useful approximations), the future in the statistical mechanical analysis of biologically more realistic recurrent neural networks is clearly in the non-equilibrium half of the statistical mechanics playing field.
### Acknowledgements
I is my pleasure to thank Heinz Horner, David Sherrington and Nikos Skantzos for their direct and indirect contributions to this review.
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# PROBING THE QCD VACUUM WITH FLAVOUR SINGLET OBJECTS: 𝜂' ON THE LATTICE
## 1 Introduction
Flavour symmetric hadronic states are of particular importance for the understanding of nonperturbative aspects of quantum chromodynamics (QCD). In fact they are expected to reveal important insight into the topological vacuum structure induced by gluonic self interactions, as borne out in the famous Witten-Veneziano formula $`^\mathrm{?}`$ that relates the flavour octet/singlet mass splitting, $`M_0^2`$ <sup>a</sup><sup>a</sup>aUpper (lower) case letters refer to masses in physical (lattice) units., to the topological susceptibility, $`\chi _q`$, in the large-$`N_c`$ limit of the theory:
$$M_0^2:=M_\eta ^{}^2M_8^2=2N_f\chi _q/f_\pi ^2,$$
(1)
with $`N_f`$ being the number of active flavours and $`f_\pi `$ the pion decay constant.
While glueballs have been the target of many ab initio lattice investigations over the past fifteen years $`^\mathrm{?}`$ much less attention has been paid to the direct calculation of flavour singlet states with fermionic content such as the $`\eta ^{}`$ meson. The reason is obvious: the extraction of flavour singlet mesons from the lattice is even more challenging than the one of glueball states, due to the substantial cost of computing quark (rather than Wilson) loop correlators. In fact such Zweig-rule forbidden diagrams (in general misleadingly called ‘disconnected diagrams’) consist of two quark loops connected via gluon lines only and need to be calculated in the momentum zero state, which amounts to the computationally very expensive evaluation of the trace of the inverse Dirac operator. The two pioneering studies in the field used quenched vacuum configurations, with Wilson fermions and wall sources $`^\mathrm{?}`$ and staggered fermions with stochastic sources $`^\mathrm{?}`$.
This situation will change with the imminent advent of near-Teraflops computers that promise both to provide (a) sufficient sampling rates for full QCD vacuum configurations and (b) the analysing power to deal with the above mentioned loop-loop correlators. Albeit present day simulations are still restricted to the case of two active flavours, $`N_f=2`$, it is of great interest to fathom state-of-the art techniques to deal with Zweig rule forbidden diagrams. First results on flavour singlet masses were presented recently by CP-PACS $`^\mathrm{?}`$ and UKQCD $`^\mathrm{?}`$.
In the present paper we present a study on the $`\eta ^{}`$ mass using the full set of QCD configurations generated by SESAM ($`16^3\times 32`$, ‘small lattice’) and T$`\chi `$$`^\mathrm{?}`$ ($`24^3\times 40`$, ‘large lattice’), both at $`\beta =5.6`$ with standard Wilson action (see Table 1).
## 2 Disconnected diagrams
### 2.1 The problem
We consider the pseudoscalar flavour singlet operator in a flavour symmetric theory
$$\eta ^{}(x)=\underset{i=1}{\overset{N_f}{}}\overline{q}_i(x)\gamma _5q_i(x),$$
(2)
with $`N_f`$ flavours. By the usual Wick contraction it leads to the flavour singlet propagator in terms of the inverse Dirac operator, $`M^1`$:
$`C_\eta ^{}(0|x)`$ $``$ $`N_f\text{tr}((M^1(0|x))^{}M^1(0|x))`$ (3)
$`N_f^2\text{tr}(\gamma _5M^1(0|0))^{}\text{tr}(\gamma _5M^1(x|x)),`$
which is a sum of fermionic connected and disconnected terms with traces to be taken in the spin and colour spaces. In the rest of the paper we shall refer to them as ‘one-loop’ and ‘two-loop’ contributions, respectively. The momentum zero projection
$$C_\eta ^{}(t)\eta ^{}(t)\eta ^{}(0)_{conn}\eta ^{}(t)\eta ^{}(0)_{disc}$$
(4)
is expected to decay exponentially, $`\mathrm{exp}(m_\eta ^{}t)`$ and reveal the flavour singlet mass, $`m_\eta ^{}`$. On an antiperiodic (of length $`T`$ in time) lattice one should encounter the usual $`\mathrm{cosh}`$ behaviour at large values of $`t`$ and $`Tt`$
$$G(t)\mathrm{exp}(m_\eta ^{}t)+\mathrm{exp}(m_\eta ^{}(Tt)).$$
(5)
From this parametrization local effective masses $`m_\eta ^{}(t)`$ can be retrieved by solving the implicit equations
$$\frac{G(t+1)}{G(t)}=\frac{\mathrm{exp}(m_\eta ^{}(t+1))+\mathrm{exp}(m_\eta ^{}(Tt1))}{\mathrm{exp}(m_\eta ^{}t)+\mathrm{exp}(m_\eta ^{}(Tt))},$$
(6)
where the l.h.s. ratios are taken from the lattice.
### 2.2 Diagonal improved stochastic estimator
Obviously the momentum zero projection embodied in eq. 3 requires the knowledge of the traces, $`\text{tr}(\gamma _5M^1(x|x))`$, on each time slice. But their exact evaluation is far beyond the scope of present computers.
One therefore takes resort to a stochastic estimate of the disconnected diagrams by introducing an ensemble of $`N_e`$ nonlocal sources, $`\varphi (x)^a,a=1,\mathrm{}N_e`$ with uncorrelated stochastic entries on the lattice sites $`x\text{colour}\times \text{spin}\times R_4`$ space such that in the limit $`N_e\mathrm{}`$ the ensemble averages show the following limiting behaviours
$`\varphi (x){\displaystyle \frac{1}{N_e}}{\displaystyle \underset{a=1}{\overset{N_e}{}}}\varphi ^a(x)`$ $``$ $`0`$ (7)
$`\varphi (x)^{}\varphi (y)`$ $``$ $`\delta _{x,y}.`$ (8)
In our case, we chose $`Z_2`$ noise for the components of the source vectors. By solving the Dirac equation
$$M(z,x)\xi (x)=\varphi (z)$$
(9)
on the stochastic sources $`\varphi (z)`$ one retrieves a bias free estimate for the inverse Dirac operator in the large $`N_e`$ limit
$$\varphi (y)^{}\xi (x)=\underset{z}{}M^1(z,x)\varphi (y)^{}\varphi (z)M^1(y,x),$$
(10)
and hence for the one-loop term, $`\text{tr}(\gamma _5M^1)`$.
But beware of the subleading terms when computing components of $`M^1`$ as we do here (actually in spin space); for they might be affected by the dominant diagonal ones in so far the latter are not sufficiently rejected by our approximation to the Kronecker $`\delta `$ in eq. 8. So we improve the estimate on nondiagonals, $`M^1(y,x)`$ ($`xy`$), by subtracting out, from the expression on the l.h.s. of eq. 10, those leading error terms by redefining the estimator to be $`^\mathrm{?}`$:
$`\varphi (y)^{}\xi (x)M^1(x,x)\varphi (y)^{}\varphi (x)`$ $`=`$ $`M^1(y,x)`$ (11)
$`+{\displaystyle \underset{zx,y}{}}M^1(z,x)\varphi (y)^{}\varphi (z).`$
Note that in this improvement step, the subleading term on the l.h.s. again is gained by SET. Throughout this work we shall apply eq. 11 in order to achieve improved estimates to $`\text{tr}(\gamma _5M^1)`$. We mention in passing that an alternative route would be to work in a spin explicite mode by using spin projected sources $`^\mathrm{?}`$ $`^\mathrm{?}`$.
### 2.3 Quality of the signal
In Table 1 we collected the characteristic parameters of our present analysis. We have used five different sea quark masses and two different lattice sizes to have some control on finite-size effects. While the number of vacuum configurations varies from 156 to 195, the number of independent stochastic sources has been chosen to be 400 (100) for the small (large) lattices. The pseudoscalar and vector mass ratios quoted refer to our final spectrum analysis $`^\mathrm{?}`$. The results presented in this talk are based on local sources only. The errors quoted are statistical and have been obtained by jackknifing.
As a result we obtain two-loop correlators, $`C_{disc}(t)`$, with statistical quality illustrated for our lightest quark mass on the small lattice in Fig. 1. The plot also contains the connected correlator, $`C_8(t)`$, corresponding to the nonsinglet pseudoscalar meson, in order to demonstrate the reduced accuracy of the (Zweig-rule forbidden) two-loop contribution. At the time separation $`t10`$, we are faced with a statistical error of order $`25`$ % on the latter.
Given this situation it is practically impossible to establish an effective mass plateau in the $`\eta ^{}`$ channel, $`C_\eta ^{}(t)`$, (eq. 3). For further analysis we made single exponential fits to $`C_\eta ^{}(t)`$ and $`C_8(t)`$ over the t-ranges listed in Table 1. The quality of our fits is illustrated for the lightest quark mass on the large lattice in Fig. 2.
## 3 Physics results
### 3.1 Chiral extrapolations
Because of the well-known technical limitations of the hybrid Monte Carlo algorithm $`^\mathrm{?}`$, the SESAM and T$`\chi `$L configurations correspond to two mass degenerate light sea quark flavours ($`N_f=2`$), with the unrenormalized mass value (in lattice units)
$$m_q=1/2(\kappa ^1\kappa _c^1).$$
(12)
From our previous light spectrum analysis $`^\mathrm{?}`$, we quote the lattice spacing
$$a_\rho ^1(\kappa _l)=2.302(64)\text{GeV}$$
(13)
and the critical and light quark $`\kappa `$ values:
$$\kappa _c=.158507(44),\kappa _{light}=.158462(42).$$
(14)
From our data we cannot decide whether it is $`m_\eta ^{}^2`$ or $`m_\eta ^{}`$ that follows a linear behaviour in the quark mass: the two fits shown in Fig. 3 work equally well, with $`\chi ^2/d.o.f.𝒪(1)`$. Note that we make no distinction between sea and valence quarks, as we choose the quark masses in the fermion loops to equal the sea quark masses (symmetric extrapolation in the sense of ref. $`^\mathrm{?}`$).
In our $`N_f=2`$ world, according to eq. 3, we would not expect to encounter the full effect of Zweig rule fordbidden diagrams, and hence we anticipate to underestimate the real world $`\eta ^{}`$ mass. From the experimental mass splitting
$$M_0^2=M_\eta ^{}^2M_8^2,$$
(15)
we therefore compute, in the spirit of the Witten-Veneziano formula eq. 1, the ‘pseudoexperimental’ value in the $`N_f=2`$ world:
$$M_\eta ^{}^2(N_F=2)=2/3M_0^2+M_\pi ^2=(716\text{MeV})^2.$$
(16)
This value corresponds in lattice units to the full squares in the two alternative chiral extrapolations shown in Fig. 3. To compare: our lattice analysis yields at the light quark mass
$$M_\eta ^{}^2=(551(85)\text{MeV})^2\text{and}M_\eta ^{}=615(53)\text{MeV},$$
(17)
by use of a linear ansatz in $`m_\eta ^{}^2(m_q)`$ and $`m_\eta ^{}(m_q)`$, respectively.
### 3.2 Impact of topology
We expect two-loop diagrams to be sensitive to the topology of the vacuum. A simple check is to look for a dependency of the ratio of disconnected and connected correlators, $`R_Q(t)`$, on the size of the topological charge $`|Q|`$, as determined in ref. $`^\mathrm{?}`$. In Fig. 4, we have plotted, for $`\kappa _{sea}=.1575`$ and the small lattice, the data with cuts according to $`|Q|1.5`$ (top01) and $`|Q|>1.5`$ (top24). We definitely find a dependency of $`R_Q`$ on $`|Q|`$. Note that the disconnected piece vanishes in the vacuum sector with small values of $`|Q|`$!
## 4 Discussion and outlook
The analysis presented here looks rather promising, but leaves room for further investigations:
Smearing. The $`\eta ^{}`$ correlator being a difference of one-loop and two-loop terms it is of great importance to attain early asymptotics of the flavour singlet correlator. But with local sources the connected (flavour octet) correlator does require analysis at $`t`$-values beyond $`12`$, in view of appreciable excited state contributions. Yet for the flavour singlet situation, we had to analyse in the t-range 6 to 9, for lack of statistics. Thus it appears very promising to improve the ground state overlap by use of source and sink smearing; for this will help to establish mass plateaus in the flavour singlet channel and hence reduce errors. Indeed, we have observed mass plateaus appearing at $`t4`$ to $`5`$, as illustrated in Fig. 5. In a forthcoming paper, we shall elaborate on this point $`^\mathrm{?}`$.
Spectral methods. Another possible direction to go is to construct the two-loop correlator from the low lying eigenfunctions of the Dirac operator. For illustraton we show in Fig. 6 the result of an attempt to saturate the spectral representation of the two-loop correlator by the 300 lowest eigenmodes of $`(\gamma _5M)`$ $`^\mathrm{?}`$. We find very nice agreement between the results from this eigenmode approach (EVA) from SET. This is very encouraging, since in the deeper chiral regime EVA will become even more efficient whereas the performance of SET will deteriorate.
## Acknowledgements
T.S., H.N., and B.O. thank the DFG-Graduiertenkolleg “Feldtheoretische und Numerische Methoden in der Statistischen und Elementarteilchenphysik” for support. The HMC productions were run on an APE100 at NIC Zeuthen and INFN Roma while the loop computations were carried out on APE100 systems at DESY Zeuthen and at the University of Bielefeld. We are grateful to our colleagues F. Rapuano and G. Martinelli for the fruitful T$`\chi `$L-collaboration. Analysis was also performed on the CRAY T3E system of ZAM at Research Center Jülich. K.S. thanks the organizers of Confinement2000, in particular Prof. H. Suganuma, for inviting him to their very stimulating symposion.
## References
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# A Note on Closed Geodesics for a Class of non–compact Riemannian Manifolds
## 1 Introduction
This paper is concerned with the existence of closed geodesics on a non–compact manifold $`M`$. There are very few papers on such a problem, see . In particular, Tanaka deals with the manifod $`M=\times S^N`$, endowed with a metric $`g(s,\xi )=g_0(\xi )+h(s,\xi )`$, where $`g_0`$ is the standard product metric on $`\times S^N`$. Under the assumption that $`h(s,\xi )0`$ as $`|s|\mathrm{}`$, he proves the existence of a closed geodesic, found as a critical point of the energy functional
$$E(u)=\frac{1}{2}_0^1g(u)[\dot{u},\dot{u}]𝑑t,$$
(1)
defined on the loop space $`\mathrm{\Lambda }=\mathrm{\Lambda }(M)=H^1(S^1,M)`$ <sup>2</sup><sup>2</sup>2We will identify $`S^1`$ with $`[0,1]/\{0,1\}`$.. The lack of compactness due to the unboundedness of $`M`$ is overcome by a suitable use of the concentration–compactness principle. To carry out the proof, the fact that $`M`$ has the specific form $`M=\times S^N`$ is fundamental, because this permits to compare $`E`$ with a functional at infinity whose behavior is explicitely known.
In the present paper, we consider a perturbed metric $`g_\epsilon =g_0+\epsilon h`$, and extend Tanaka’s result in two directions. First, we show the existence of at least $`N`$, in some cases $`2N`$, closed geodesics on $`M=\times S^N`$, see Theorem 2.6. Such a theorem can also be seen as an extension to cilyndrical domains of the result by Carminati . Next, we deal with the case in which $`M=\times M_0`$ for a general compact $`N`$–dimensional manifold $`M_0`$ <sup>3</sup><sup>3</sup>3By manifold we mean a smooth, connected manifold.. The existence result we are able to prove requires that either $`M_0`$ possesses a non–degenerate closed geodesic, see Theorem 3.5, or that $`\pi _1(M_0)\{0\}`$ and the geodesics on $`M_0`$ are isolated, see Theorem 4.3.
The approach we use is different that Tanaka’s one, and relies on a perturbation result discussed in that leads to rather simple proofs. Roughly, the main advantages of using this abstract perturbation method are that
* we can obtain sharper results, like the multeplicity ones;
* we can deal with a general manifold like $`M=\times M_0`$, not only $`M=\times S^N`$, when the results – for the reasons indicated before – cannot be easily obtained by using Tanaka’s approach.
The author wishes to thank professor Ambrosetti, for suggesting the problem, teaching him the perturbative technique, and much more.
## 2 Spheres
In this section we assume $`M=\times S^N`$, where $`S^N=\{\xi ^{N+1}:|\xi |=1\}`$ <sup>1</sup><sup>1</sup>1Hereafter, we use the notation $`\xi \eta =_i\xi _i\eta _i`$ for the scalar product in $`^{N+1}`$, and $`|\xi |^2=\xi \xi `$.. For $`s`$, $`rT_s`$, $`\xi S^N`$, $`\eta T_\xi S^N`$, let
$$g_0(s,\xi )((r,\eta ),(r,\eta ))=|r|^2+|\eta |^2$$
(2)
be the standard product metric on $`M=\times S^N`$. We consider a perturbed metric
$$g_\epsilon (s,\xi )((r,\eta ),(r,\eta ))=|r|^2+|\eta |^2+\epsilon h(s,\xi )((r,\eta ),(r,\eta )),$$
(3)
where $`h(s,\xi )`$ is a bilinear form, not necessarily positive definite.
Define the space of closed loops
$$\mathrm{\Lambda }=\{u=(r,x)H^1(S^1,)\times H^1(S^1,S^N)\}$$
(4)
Closed geodesics on $`(M,g_\epsilon )`$ are the critical points of $`E_\epsilon :\mathrm{\Lambda }`$ given by
$$E_\epsilon (u)=\frac{1}{2}_0^1g_\epsilon (u)[\dot{u},\dot{u}]𝑑t.$$
(5)
One has that
$$E_\epsilon (u)=E_\epsilon (r,x)=E_0(r,x)+\epsilon G(r,x),$$
(6)
where
$$E_0(r,x)=\frac{1}{2}_0^1\left(|\dot{r}|^2+|\dot{x}|^2\right)𝑑t$$
and
$$G(r,x)=\frac{1}{2}_0^1h(r,x)[(\dot{r},\dot{x}),(\dot{r},\dot{x})]𝑑t.$$
(7)
In particular, we split $`E_0`$ into two parts, namely
$$E_0(r,x)=L_0(r)+E_{M_0}(x),$$
(8)
where
$$L_0(r)=\frac{1}{2}_0^1|\dot{r}|^2𝑑t,E_{M_0}(x)=\frac{1}{2}_0^1|\dot{x}|^2𝑑t.$$
The form of $`E_\epsilon `$ suggests to apply the perturbative results of that we recall below for the reader’s convenience.
###### Theorem 2.1 ().
Let $`H`$ be a real Hilbert space, $`E_\epsilon C^2(H)`$ be of the form
$$E_\epsilon (u)=E_0(u)+\epsilon G(u),$$
(9)
where $`GC^2(E)`$.. Suppose that there exists a finite dimensional manifold $`Z`$ such that
* $`E_0^{}(z)=0`$ for all $`zZ`$;
* $`E_0^{\prime \prime }(z)`$ is a compact perturbation of the identity, for all $`zZ`$;
* $`T_zZ=\mathrm{ker}E_0^{\prime \prime }(z)`$ for all $`zZ`$.
There exist a positive number $`\epsilon _0`$ and a smooth function $`w:Z\times (\epsilon _0,\epsilon _0)H`$ such that the critical points of
$$\mathrm{\Phi }_\epsilon (z)=E_\epsilon (z+w(z,\epsilon )),zZ,$$
(10)
are critical points of $`E_\epsilon `$.
Moreover, it is possible to show that
$$\mathrm{\Phi }_\epsilon (z)=b+\epsilon \mathrm{\Gamma }(z)+o(\epsilon ),$$
(11)
where $`b=E_0(z)`$ and $`\mathrm{\Gamma }=G_{|Z}`$. From this “first order” expansion, one infers
###### Theorem 2.2 ().
Let $`H`$ be a real Hilbert space, $`E_\epsilon C^2(H)`$ be of the form (9). Suppose that (AS1)–(AS3) hold. Then any strict local extremum of $`G_{|Z}`$ gives rise to a critical point of $`E_\epsilon `$, for $`|\epsilon |`$ sufficiently small.
In the present situation, the critical points of $`E_0`$ are nothing but the great circles of $`S^N`$, namely
$$z_{p,q}=p\mathrm{cos}2\pi t+q\mathrm{sin}2\pi t,$$
(12)
where $`p,q^{N+1}`$, $`pq=0`$, $`|p|=|q|=1`$. Hence $`E_0`$ has a “critical manifold” given by
$$Z=\{z(r,p,q)=(r,z_{p,q}())r,z_{p,q}\text{ as in (}\text{12}\text{)}\}.$$
###### Lemma 2.3.
$`Z`$ satisfies (AS2)–(AS3).
###### Proof.
The first assertion is known, see for instance .
For the second statement, we closely follow .
For $`zZ`$, of the form $`z(t)=(r,z_{p,q}(t))`$, it turns out that
$$E_0^{\prime \prime }(z)[h,k]=_0^1\left[\dot{h}\dot{k}|\dot{z}|^2hk\right]𝑑t$$
for any $`h,kT_zZ`$.
Let $`e_i^{N+1}`$, $`i=2,\mathrm{},N+1`$, be orthonormal vectors such that $`\{\frac{1}{2\pi }\dot{z}_{p,q},e_2,\mathrm{}e_{N+1}\}`$ is a basis of $`T_zZ`$, and set
$$e_i(t)=\{\begin{array}{cc}\dot{z}_{u^1,u^2}(t)/2\pi \hfill & \text{ if }i=1\hfill \\ e_i\hfill & \text{ if }i>1,\hfill \end{array}$$
Then, for $`h,k`$ as before, we can write a “Fourier–type” expansion
$$h(t)=h_0(t)\frac{d}{dt}+\underset{i=1}{\overset{N1}{}}h_i(t)e_i(t),k(t)=k_0(t)\frac{d}{dt}+\underset{i=1}{\overset{N1}{}}k_i(t)e_i(t).$$
(13)
Assume now that $`h\mathrm{ker}E_0^{\prime \prime }(z_{p,q})`$, i.e.
$$_0^1\dot{h}\dot{k}𝑑t=_0^1|\dot{z}|^2hk𝑑tkT_{z_{p,q}}Z.$$
We plug (13) into this relations, and we get the system
$$\{\begin{array}{cc}\stackrel{..}{h_1}=0\hfill & \\ \stackrel{..}{h_j}+4\pi ^2h_j=0j=2,\mathrm{},N1\hfill & \\ \stackrel{..}{h_0}=0.\hfill & \end{array}$$
(14)
Recalling that $`h_0`$ and $`h_1`$ are periodic, we find
$$\{\begin{array}{cc}h_0=\lambda _0,h_1=\lambda _1\hfill & \\ h_j=\lambda _j\mathrm{cos}2\pi t+\mu _j\mathrm{sin}2\pi tj=2,\mathrm{}N1.\hfill & \end{array}$$
(15)
Therefore, $`hT_zZ`$. This shows that $`\mathrm{ker}E_0^{\prime \prime }(z_{p,q})T_{z_{p,q}}Z`$. Since the converse inclusion is is always true, the lemma follows. ∎
###### Lemma 2.4.
Suppose
* $`h(r,)0`$ pointwise on $`S^N`$, as $`|r|\mathrm{}`$,
then
$$\mathrm{\Phi }_\epsilon bE_0(z).$$
Recall that $`\mathrm{\Phi }_\epsilon `$ was defined in (10).
###### Proof.
This is proved as in . We just sketch the argument. The idea is to use the contraction mapping principle to characterize the function $`w(\epsilon ,z)`$ (see Theorem 1). Indeed, define
$$H(\alpha ,w,z_r,\epsilon )=\left(\begin{array}{c}E_\epsilon ^{}(z_r+w)\alpha \dot{z}\\ w\dot{z}.\end{array}\right)$$
So $`H=0`$ if and only if $`w(T_{z_r}Z)^{}`$ and $`E_\epsilon ^{}(z_r+w)T_{z_r}Z`$. Now,
$$H(\alpha ,w,z_r,\epsilon )=0H(0,0,z_r,0)+\frac{H}{(\alpha ,w)}(0,0,z_r,0)[\alpha ,w]+R(\alpha ,w,z_r,\epsilon )=0,$$
where $`R(\alpha ,w,z_r,\epsilon )=H(\alpha ,w,z_r,\epsilon )\frac{H}{(\alpha ,w)}(0,0,z_r,0)[\alpha ,w]`$.
Setting
$$R_{z_r,\epsilon }(\alpha ,w)=\left[\frac{H}{(\alpha ,w)}(0,0,z_r,0)\right]^1R(\alpha ,w,z_r,\epsilon ),$$
one finds that
$$H(\alpha ,w,z_r,\epsilon )=0(\alpha ,w)=R_{z_r,\epsilon }(\alpha ,w).$$
By the Cauchy–Schwarz inequality, it turns out that $`R_{zr,\epsilon }`$ is a contraction mapping from some ball $`B_{\rho (\epsilon )}`$ into itself. If $`|\epsilon |`$ is sufficiently small, we have proved the existence of $`(\alpha ,w)`$ uniformly for $`z_rZ`$. We want to study the asymptotic behavior of $`w=w(\epsilon ,z_r)`$ as $`|r|+\mathrm{}`$. We denote by $`R_\epsilon ^0`$ the functions $`R_{z_r,\epsilon }`$ corresponding to the unperturbed energy functional $`E_0=E_{M_0}`$. It is easy to see (, Lemma 3) that the function $`w^0`$ found with the same argument as before satisfies $`w^0(z_r)0`$ as $`|r|+\mathrm{}`$. Thus, by the continuous dependence of $`w(\epsilon ,z_r)`$ on $`\epsilon `$ and the characterization of $`w(\epsilon ,z_r)`$ and $`w^0`$ as fixed points of contractive mappings, we deduce as in , proof of Lemma 3.2, that $`lim_r\mathrm{}w(\epsilon ,z_r)=0`$. In conclusion, we have that $`lim_{|r|+\mathrm{}}\mathrm{\Phi }_\epsilon (z_r+w(\epsilon ,z_r))=E_{M_0}(z_0)`$.
###### Remark 2.5.
There is a natural action of the group $`O(2)`$ on the space $`\mathrm{\Lambda }`$, given by
$`\{\pm 1\}\times S^1\times \mathrm{\Lambda }`$ $`\mathrm{\Lambda }`$
$`(\pm 1,\theta ,u)`$ $`u(\pm t+\theta ),`$
under which the energy $`E_\epsilon `$ is invariant. Since this is an isometric action under which $`Z`$ is left unchanged, it easily follows that the function $`w`$ constructed in Theorem 2.1 is invariant, too.
###### Theorem 2.6.
Assume that the functions $`h_{ij}=h_{ji}`$’s are smooth, bounded, and (h1) holds Then $`M=\times M_0`$ has at least $`N`$ non-trivial closed geodesics, distinct modulo the action of the group $`O(2)`$. Furthermore, if
* the matrix $`[h_{ij}(p,)]`$ representing the bilinear form $`h`$ is positive definite for $`p+\mathrm{}`$, and negative definite for $`p\mathrm{}`$,
then $`M`$ possesses at least $`2N`$ non–trivial closed geodesics, geometrically distinct.
###### Proof.
Observe that $`Z=\times Z_0`$, where $`Z_0=\{z_{p,q}|p|=|q|=1,pq=0\}`$. According to Theorem 2.1, it suffices to look for critical points of $`\mathrm{\Phi }_\epsilon `$. From Lemma 2.4, it follows that either $`\mathrm{\Phi }_\epsilon =b`$ everywhere, or has a critical point $`(\overline{r},\overline{p},\overline{q})`$. In any case such a critical point gives rise to a (non–trivial) closed geodesic of $`(M,g_\epsilon )`$.
From Remark 2.5, we know that $`\mathrm{\Phi }_\epsilon `$ is $`O(2)`$–invariant. This allows us to introduce the $`O(2)`$–category $`\mathrm{cat}_{O(2)}`$. One has
$$\mathrm{cat}_{O(2)}(Z)\mathrm{cat}(Z/O(2))\mathrm{cuplength}(Z/O(2))+1.$$
Since $`\mathrm{cuplength}(Z/O(2))N1`$, (see ), then $`\mathrm{cat}_{O(2)}(Z)N`$. Finally, by the Lusternik–Schnirel’man theory, $`M`$ carries at least $`N`$ closed geodesics, distinct modulo the action $`O(2)`$. This proves the first statement.
Next, let
$$\mathrm{\Gamma }(r,p,q)=G((r,z_{p,q}))=\frac{1}{2}_0^1h(r,z_{p,q}(t))[\dot{z}_{p,q},\dot{z}_{p,q}]𝑑t$$
(16)
Then (h) immediately implies that
$$\mathrm{\Gamma }(r,p,q)0\text{ as }|r|\mathrm{},$$
(17)
Moreover, if (h2) holds, then $`\mathrm{\Gamma }(r,p,q)>0`$ for $`r>r_0`$, and $`\mathrm{\Gamma }(r,p,q)<0`$ for $`r<r_0`$. Since (recall equation (11))
$$\mathrm{\Phi }_\epsilon (r,p,q)=b+\epsilon \mathrm{\Gamma }(r,p,q)+o(\epsilon ),$$
(18)
it follows that
$$\{\begin{array}{cc}\mathrm{\Phi }_\epsilon (r,p,q)>b\hfill & \text{ for }r>r_0\hfill \\ \mathrm{\Phi }_\epsilon (r,p,q)<b\hfill & \text{ for }r<r_0.\hfill \end{array}$$
We can now exploit again the $`O(2)`$ invariance.
By assumption, and a simple continuity argument, $`\{\mathrm{\Phi }_\epsilon >b\}[R_0,\mathrm{})\times Z_0`$, and similarly $`\{\mathrm{\Phi }_\epsilon <b\}[\mathrm{},R_0)\times Z_0`$, for a suitably large $`R_0>0`$. Hence $`\mathrm{cat}_{O(2)}(\{\mathrm{\Phi }_\epsilon >b\})\mathrm{cat}_{O(2)}(Z_0)=N`$. The same argument applies to $`\{\mathrm{\Phi }_\epsilon <b\}`$. This proves the existence of at least $`2N`$ closed geodesics. ∎
###### Remark 2.7.
* In , the existence of $`N`$ closed geodesics on $`S^N`$ endowed with a metric close to the standard one is proved. Such a result does not need any study of $`\mathrm{\Phi }_\epsilon `$ and its behavior. The existence of $`2N`$ geodesics is, as far as we know, new. We emphasize that it strongly depends on the form of $`M=\times M_0`$.
* In , the metric $`g`$ on $`M`$ is possibly not perturbative. No multiplicity result is given.
## 3 The general case
In this section we consider a compact riemannian manifold $`(M_0,g_0)`$, and in analogy to the previous section, we put
$$g_\epsilon (s,\xi )((r,\eta ),(r,\eta ))=|r|^2+g_0(\xi )(\eta ,\eta )+\epsilon h(s,\xi )((r,\eta ),(r,\eta )).$$
(19)
Again, we define $`\mathrm{\Lambda }=\{u=(r,x)rH^1(S^1,),xH^1(S^1,M_0)\}`$,
$$E_{M_0}(x)=\frac{1}{2}_0^1g_0(x)(\dot{x},\dot{x})𝑑t,E_0(r,x)=\frac{1}{2}_0^1|\dot{r}|^2𝑑t+E_{M_0}(x),$$
and finally
$$E_\epsilon (r,x)=E_0(r,x)+\epsilon G(r,x),$$
with $`G`$ as in (7). It is well known () that $`M_0`$ has a closed geodesic $`z_0`$. The functional $`E_{M_0}`$ has again a critical manifold $`Z`$ given by
$$Z=\{u()=(\rho ,z_0(+\tau ))\rho \text{ constant, }\tau S^1\}.$$
Let $`Z_0=\{z_0(+\tau )\tau S^1\}`$. It follows that $`Z\times Z_0`$. The counterpart of $`\mathrm{\Gamma }`$ in (11) is
$$\mathrm{\Gamma }(r,\tau )=\frac{1}{2}_0^1h(r,z_\tau )[\dot{z}_\tau ,\dot{z}_\tau ]𝑑t.$$
(20)
Let us recall some facts from .
###### Remark 3.1.
There is a linear operator $`A_z:T_z\mathrm{\Lambda }(M_0)T_z\mathrm{\Lambda }`$, which is a compact perturbation of the identity, such that
$$E_{M_0}^{\prime \prime }(z)[h,k]=A_zhk_1=_0^1\stackrel{}{\stackrel{}{A_zh}}\dot{k}𝑑t.$$
In particular, $`E_0`$ satisfies (AS2).
###### Definition 3.2.
Let
$$\mathrm{ker}E_{M_0}^{\prime \prime }(z_0)=\{hT_{z_0}\mathrm{\Lambda }(M_0)A_{z_0}hk_1=0kT_{z_0}\mathrm{\Lambda }(M_0)\}.$$
We say that a closed geodesic $`z_0`$ of $`M_0`$ is non–degenerate, if
$$dim\mathrm{ker}E_{M_0}^{\prime \prime }(z_0)=1.$$
###### Remark 3.3.
For example, it is known that when $`M_0`$ has negative sectional curvature, then all the geodesics of $`M_0`$ are non–degenerate. See . Moreover, it is easy to see that the existence of non–degenerate closed geodesics is a generic property.
###### Lemma 3.4.
If $`z_0`$ is a non–degenerate closed geodesic of $`M_0`$, then $`Z`$ satisfies (AS2).
###### Proof.
It is always true that $`T_{z_r}Z\mathrm{ker}E_0^{\prime \prime }(z_r)`$. By (26), we have that $`dimT_{z_r}Z=dim\mathrm{ker}E_0^{\prime \prime }(z_r)`$. This implies that $`T_{z_r}Z=\mathrm{ker}E_0^{\prime \prime }(z_r)`$. A generic element of $`Z`$ has the form $`(\rho ,z^\tau )`$ for $`\rho `$ and $`z^\tau =z(+\tau )`$; then
$$T_{(\rho ,z^\tau )}M=\times T_{z^\tau }M_0,$$
and any two vector fields $`Y`$ and $`W`$ along a curve on $`M=\times M_0`$ can be decomposed into
$$Y=h(t)\frac{d}{dt}+y(t)T_{z^\tau }Z_0,$$
(21)
$$W=k(t)\frac{d}{dt}+w(t)T_{z^\tau }Z_0.$$
(22)
In addition, there results (see )
$$E_{M_0}^{\prime \prime }(z_0)[y,w]=_0^1[g_0(D_ty,D_tw)g_0(R_{M_0}y(t),\dot{z}_0(t))\dot{z}_0(t)w(t))]dt,$$
(23)
and
$$R_M(r,z)=R_{}(r)+R_{M_0}(z)=R_{M_0}(z),$$
(24)
where $`R_M`$, $`R_{M_0}`$, etc. stand for the curvature tensors of $`M`$, $`M_0`$, etc. By (23), (21) and (22), as in the previous section, $`E_0^{\prime \prime }(\rho ,z_\tau )[Y,W]=0`$ is equivalent to the system
$$\{\begin{array}{cc}\ddot{h}=0\hfill & \\ _0^1g_0(z)[D_ty,D_tw]R_{M_0}(y(t),\dot{z}_r(t))\dot{z}_r(t)w(t)dt=0.\hfill & \end{array}$$
(25)
As in the case of the sphere, the first equation implies that $`h`$ is constant. The second equation in (25) implies that $`y\mathrm{ker}E_{M_0}^{\prime \prime }(z^\tau )=\mathrm{ker}E_{M_0}^{\prime \prime }(z_0)`$. Hence,
$$\mathrm{ker}E_0^{\prime \prime }(z_r)=\{(h,y)h\text{ is constant, and }y\mathrm{ker}E_{M_0}^{\prime \prime }(z_0)\}.$$
(26)
This completes the proof. ∎
###### Theorem 3.5.
Let $`M_0`$ be a compact, connected manifold of dimension $`N<\mathrm{}`$. Assume that $`M_0`$ admits a non–degenerate closed geodesic $`z`$, and that (in local coordinates) $`h_{ij}(p,)a_{}`$ as $`p\mathrm{}`$, and $`h_{ij}(p,)a_+`$ as $`p+\mathrm{}`$.
1. If $`a_{}=a_+`$ and $`h_{ij}(p,)`$ satisfies (h2), then $`M`$ has at least one closed geodesic.
2. If $`a_{}a_+`$ and $`h_{ij}(p,)[u,v]a(uv)`$ is negative definite for $`p\mathrm{}`$ and positive definite for $`p+\mathrm{}`$, then $`M`$ has at least two non-trivial closed geodesic.
###### Proof.
Lemma 3.4 allows us to repeat all the argument in Theorem 2.6, and the result follows immediately. ∎
## 4 Isolated geodesics
In this final section, we discuss one situaion where the critical manifold $`Z`$ may be degenerate. Here, the non–degeneracy condition (AS3) fails, and $`T_zZ\mathrm{ker}E_0^{\prime \prime }(z)`$ strictly. Fix a closed geodesic $`Z_0`$ for $`M_0`$, and put $`\stackrel{~}{W}=(T_{z_0}Z)^{}`$. Since $`T_zZ\mathrm{ker}E_0^{\prime \prime }(z)`$ strictly, there exists $`k>0`$ such that $`\stackrel{~}{W}=(\mathrm{ker}E_0^{\prime \prime }(z_0))^{}^k`$. Repeating the preceding finite dimensional reduction, one can find again a unique map $`\stackrel{~}{w}=\stackrel{~}{w}(z,\zeta )`$, where $`zZ`$ and $`\zeta ^k`$, in such a way that $`E_\epsilon ^{}=0`$ reduces to an equation like
$$A(z+\zeta +\stackrel{~}{w}(z,\zeta ))=0.$$
If $`z_0`$ is an isolated minimum of the energy $`E_{M_0}`$ over some connected component of $`\mathrm{\Lambda }(M_0)`$, then it is possible to show that there exists again a function $`\mathrm{\Gamma }:Z`$ such that
$$A(z+\zeta +\stackrel{~}{w}(z,\zeta ))=0\frac{\mathrm{\Gamma }}{r}(R,\tau )\frac{\mathrm{\Gamma }}{r}(R,\tau )0$$
for some $`R`$ and all $`\tau S^1`$. For more details, see . In particular, we will use the following result.
###### Theorem 4.1.
Let $`H`$ be a real Hilbert space, $`f_\epsilon :H`$ is a family of $`C^2`$–functionals of the form $`f_\epsilon =f_0+\epsilon G`$, and that:
* $`f_0`$ has a finite dimensional manifold $`Z`$ of critical points, each of them being a minimum of $`f_0`$;
* for all $`zZ`$, $`f_0^{\prime \prime }(z)`$ is a compact perturbation of the identity.
Fix $`z_0Z`$, put $`W=(T_{z_0}Z)^{}`$, and suppose that $`(f_0)_{|W}`$ has an isolated minimum at $`z_0`$. Then, for $`\epsilon `$ sufficiently small, $`f_\epsilon `$ has a critical point, provided $`\mathrm{deg}(\mathrm{\Gamma }^{},B_R,0)0`$.
###### Remark 4.2.
Theorem 4.1 has been presented in a linear setting. For Riemannian manifold, we can either reduce to a local situation and then apply the exponential map, or directly resort to the slightly more general degree theory on Banach manifold developed in .
###### Theorem 4.3.
Assume that $`\pi _1(M_0)\{0\}`$, and that all the critical points of $`E_0`$, the energy functional of $`M_0`$, are isolated. Suppose the bilinear form $`h`$ satisfies (h1), and
* $`{\displaystyle \frac{h}{r}}(R,\xi ){\displaystyle \frac{h}{r}}(R,\xi )0`$ for some $`R>0`$ and all $`\xi S^1`$.
Then, for $`\epsilon >0`$ sufficiently small, the manifold $`M=\times M_0`$ carries at least one closed geodesic.
###### Proof.
We wish to use Theorem 4.1. Since $`\pi _1(M_0)\{0\}`$, then $`E_0`$ has a geodesic $`z_0`$ such that $`E_0(z_0)=\mathrm{min}E_0`$ over some component $`C`$ of $`\mathrm{\Lambda }(M_0)`$. See .
We consider the manifold
$$Z=\{u\mathrm{\Lambda }u(t)=(\rho ,z_0(t+\tau )),\rho \text{ constant, }\tau S^1\}.$$
Here we do not know, a priori, if $`Z`$ is non–degenerate in the sense of condition (AS2). But of course $`(E_0)_W`$ has a minimum at the point $`(\rho ,z_0)`$, where $`W=(T_{\rho ,z^\tau )}Z)^{}`$. We now check that it is isolated for $`(E_0)_W`$. We still know that $`Z=\times Z_0`$. Take any point $`(\rho ,z_\tau )Z`$, and observe that $`T_{(\rho ,z^\tau )}Z=\{(r,y)r,yT_{z^\tau }Z_0\}`$. For all $`(r,y)W`$ sufficiently close to $`(\rho ,z_0)`$, it holds in particular that $`yz_\tau `$. Hence
$$E_0(r,y)=L_0(r)+E_{M_0}(y)E_{M_0}(y)>E_{M_0}(z^\tau )=E_{M_0}(z_0)=E_0(\rho ,z_0)$$
since $`L_00`$ and $`z_0`$ (and hence $`z^\tau `$, due to $`O(2)`$ invariance) is an isolated minimum of $`E_{M_0}`$ by assumption.
Finally, thanks to assumption (h3), $`\frac{\mathrm{\Gamma }}{r}(R,\tau )\frac{\mathrm{\Gamma }}{r}(R,\tau )0`$.
This concludes the proof. ∎
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# HARES: an efficient method for first-principles electronic structure calculations of complex systems
## I Introduction
Technological advances are pushing the size of device components and the demands on their performance to ever smaller sizes and higher standards. These trends, which are only expected to accelerate in the future, make it imperative that the structure and behavior of systems at the atomistic level are thoroughly understood from a fundamental perspective. While experimental methods are steadily improving in their ability to probe atomistic processes of materials, computational approaches provide a complementary technique for systematic and well-controlled studies. Tuning and optimization of the properties of materials, taking into account their chemical composition, require computational methods that are of general applicability, unbiased and accurate.
First-principles methods based on density functional theory (DFT) dft have proven to be an accurate and reliable tool in understanding and predicting a wide range of physical properties of finite (such as molecules and clusters) and extended structures (such as bulk crystalline solids, their defects and surfaces). Such methods must obtain the quantum mechanical ground state of the interacting electrons and ions, which makes them computationally very intensive. The computational cost scales as a rather high power (typically 3) of the number of atoms or electrons in the system, which limits the sizes that can be investigated in a reasonable time. Improvement of the efficiency of first-principles calculations is therefore an important goal. This goal can be reached either by improving the algorithms to obtain better scaling with system size, or by exploiting modern computational resources and in particular parallel architectures.
The scaling of the calculation of the quantum mechanical ground state of interacting electrons and ions can be improved by exploiting what W. Kohn has called “the short-sightedness” of quantum mechanics: due to screening, interactions are essentially short ranged. A natural way to express this property is through the density matrix of the system. Thus, implementations of linear scaling (referred to as $`O(N)`$ methods, with $`N`$ a number representing the system size, like the number of electrons) typically involve density-matrix expressions with localized orbitals linear . In terms of the computational time, $`O(N)`$ methods prove advantageous for insulating systems with $`N>10^3`$ or metallic systems with $`N>10^4`$ electrons (the number of atoms in the system is typically an order of magnitude smaller than the number of electrons). This difference in efficiency of the $`O(N)`$ methods between insulators and metals is related to the long-range behavior of the density matrix which has a different fall off, i.e. exponential in the former vs. power law in the latter. These methods are well suited for the calculation of the total energy of the system, which provides useful information about its optimal structure, dynamics, and response to mechanical loading. In addition to the total energy, it is often important to study the electronic structure of the system. This is necessary for understanding electronic, optical and magnetic properties, and is relevant for the study of both insulators (semiconductors) and metals. A DFT electronic structure calculation requires the calculation of the eigenvalue spectrum of the single-particle Hamiltonian, a problem which is not easily amenable to improvements in scaling since diagonalization of the Hamiltonian typically scales as $`N^3`$. Localization of the electronic orbitals can be detrimental to the accuracy of electronic properties though it helps improve the efficiency of $`O(N)`$ methods.
In DFT calculations, the single-particle Hamiltonian matrix itself depends on the eigenfunctions, so the complete solution must be obtained by iterating the solution to self-consistency. The size of the Hamiltonian matrix depends on the basis set used to represent the electronic wavefunctions and the electronic charge density. Ideally, one would like to work with a sparse Hamiltonian, which can be solved efficiently using iterative algorithms; the use such algorithms reduces both computer memory and time requirements. A natural way to generate a sparse Hamiltonian is to use a real space grid for the representation of the electronic eigenfunctions and charge density: each term in the Hamiltonian, evaluated at some point in space, acts only on the wavefunctions at the same point in space, except for the Laplacian in the kinetic energy operator which involves several points simultaneously. The number of points included in the evaluation of the Laplacian determines the few off-diagonal non-zero matrix elements in each row (or column) of the Hamiltonian matrix. The calculation can be made even more efficient by using an adaptive grid in real space for representation of the eigenfunctions, with points distributed according to the electronegativity of ions. The adaptive grid can be mapped onto a regular grid in curvilinear space through the proper definition of a metric. The regular grid in curvilinear space makes it possible to exploit fully the capabilities of modern computational platforms based on parallel processing. Thus, this formulation of the problem satisfies all requirements for very efficient electronic structure calculations: (a) Sparsity for fast iterative diagonalization of the Hamiltonian; (b) Adaptability for efficient distribution of grid points as demanded by the physical system; and (c) Efficient parallelization of the computation due to the natural distribution of the regular curvilinear space grid onto the processor grid.
Our original implementation of such a method acres demonstrated the feasibility of performing calculations within this framework. The data structures and operations involved in this method make it easily parallelizable, particularly using high performance Fortran (HPF). In the present paper we discuss several algorithmic issues that enhance the performance of the method and their implementation using HPF; we refer to the new implementation as HARES for HPF-Adaptive-grid Real-space Electronic Structure. The paper is organized as follows: In section II we briefly review the theory underlying the HARES method and present an analysis of the computational effort involved in the various parts of the calculation. In section III we discuss the recent algorithmic enhancements and their implementation. In section IV we illustrate the efficacy of these algorithmic enhancements through several applications of HARES to interesting systems. These include: (a) a few simple elemental crystals and a few molecules composed of atoms in the first row of the periodic table which typically present a computational challenge to plane-wave (PW) methods; (b) blue bronze, a quasi one-dimensional conductor; and (c) a zeolite, that is, a complex structure composed of Si-O tetrahedra and large pores, which represents a molecular sieve. Section V contains our conclusions.
## II Theoretical Framework
### II.1 Density Functional Theory
The problem of finding the quantum mechanical ground state of electrons in solids is a many body problem which, at present, can be solved only approximately. The computational framework of choice for a wide range of problems involving a system of ions and interacting electrons is DFT dft . The central theorem of DFT, proven by Hohenberg and Kohn, states that the ground state energy of an electronic system is a unique functional of its charge density $`\rho (𝐫)`$ and is an extremum (a minimum) with respect to variations in the charge density. Kohn and Sham dft expressed the charge density in terms of single particle wavefunctions $`\psi _\alpha (𝐫)`$ (referred to as Kohn-Sham orbitals) and occupation numbers $`f_\alpha `$
$$\rho (𝐫)=\underset{\alpha }{}f_\alpha |\psi _\alpha (𝐫)|^2.$$
The ground state energy functional is then given by
$$\begin{array}{c}\hfill \underset{\alpha }{}f_\alpha \psi _\alpha ^{}(𝐫)\left[\frac{1}{2}^2+V_{\mathrm{ext}}(𝐫)\right]\psi _\alpha (𝐫)𝑑𝐫+\\ \hfill E_\mathrm{H}[\rho (𝐫)]+E_{\mathrm{XC}}[\rho (𝐫)],\end{array}$$
(1)
where $`V_{\mathrm{ext}}(𝐫)`$ is the external potential experienced by the electrons due to the presence of the ions, $`E_\mathrm{H}`$ is the electrostatic (also known as Hartree) energy due to Coulomb repulsion of electrons and $`E_{\mathrm{XC}}`$ is the exchange-correlation (XC) contribution, which embodies the many-body properties of the interacting electron system. A variational argument in terms of the single-particle states $`\psi _\alpha (𝐫)`$ leads to a set of single-particle equations for fictitious non-interacting particles that produce the same density as the real electrons:
$$\left[\frac{1}{2}^2+V_{\mathrm{eff}}(\rho (𝐫),𝐫)\right]\psi _\alpha (𝐫)=ϵ_\alpha \psi _\alpha (𝐫).$$
(2)
The effective potential $`V_{\mathrm{eff}}`$ in these single-particle equations is:
$$V_{\mathrm{eff}}(\rho (𝐫),𝐫)=V_{\mathrm{ext}}(𝐫)+V_\mathrm{H}[\rho (𝐫)]+V_{\mathrm{xc}}[\rho (𝐫)]$$
(3)
where $`V_\mathrm{H}`$ is the electrostatic potential due to Coulomb repulsion between electrons (known as the Hartree potential) and $`V_{\mathrm{XC}}=\delta E_{\mathrm{XC}}/\delta \rho (𝐫)`$ is the exchange correlation potential. The system of Eqs. (2), referred to as Kohn-Sham equations, represents a set of nonlinear coupled equations due to the dependence of $`V_\mathrm{H}`$ and $`V_{\mathrm{XC}}`$ on the density (and hence the wave functions $`\psi _\alpha `$); these equations are solved iteratively, beginning with a guess for the $`\psi _\alpha `$’s, until self-consistency is achieved.
The only significant approximation in this set of equations is the form of $`E_{\mathrm{XC}}[\rho (𝐫)]`$, which is not analytically known. The standard choices involve expressions that depend locally on $`\rho (𝐫)`$ (known as the Local Density Approximation — LDA), or involve both $`\rho (𝐫)`$ and its gradients (known as the Generalized Gradient Approximation — GGA). Such expressions have been derived from analyzing the behavior of the uniform or non-uniform electron gas in certain limits, or by fitting the results of accurate calculations based on quantum Monte Carlo techniques for sampling the many-body wavefunction; they work well in reproducing the energetics of a wide variety of ground state structures of extended (crystalline) or finite (cluster or molecular) systems. We discuss next the basic features of the LDA and GGA approaches and their limitations and capabilities.
In the LDA, the exchange-correlation functional is expressed as:
$$E_{\mathrm{XC}}^{\mathrm{LDA}}[\rho (𝐫)]=\rho (𝐫)ϵ_{\mathrm{XC}}^0(\rho (𝐫))\mathrm{d}^3𝐫.$$
(4)
when the number of spin-up and spin-down states in the system are equal (we refer to this as the “spin compensated” case). In this expression $`ϵ_{\mathrm{XC}}^0(\rho )`$ is the exchange correlation energy of the uniform electron gas of density $`\rho `$, which can be obtained by more elaborate computational methods like quantum Monte Carlo, or even analytically in certain limits. The ground state of a number of physical systems, such as atoms, molecules, and magnetic crystals can exhibit nonzero spin-polarization. In the above description of DFT, the single particle states label $`\alpha `$ can be extended to include the spin (up or down) quantum number, and the energy functional can be readily generalized to take into account the spin-polarized electrons. The spin-polarized version of the exchange-correlation energy functional in the so called Local Spin Density Approximation (LSDA) is written as
$$E_{\mathrm{XC}}^{\mathrm{LSDA}}[\rho _{}(𝐫),\rho _{}(𝐫)]=\rho (𝐫)ϵ_{\mathrm{XC}}^0(\rho _{}(𝐫),\rho _{}(𝐫))\mathrm{d}^3𝐫,$$
(5)
where $`\rho _{}(𝐫)`$ and $`\rho _{}(𝐫)`$ are the electronic densities of spin-up and spin-down electrons, in terms of which the total electronic density $`\rho (𝐫)`$ is:
$$\rho (𝐫)=\rho _{}(𝐫)+\rho _{}(𝐫).$$
Among the various proposed XC functionals $`ϵ_{\mathrm{XC}}^0(\rho )`$, in our approach to the DFT method we have implemented two different parametrizations of the Ceperley-Alder (CA) functional CA for both spin-polarized and spin-compensated systems: the first is the Perdew-Zunger PZ-CA parametrization (PZ-CA) and the second is the Perdew-Wang PW-CA parametrization (PW-CA). The PW-CA functional uses a more accurate spin interpolation formula for the correlation, proposed by Vosko, Wilk, and Nusair VWN , which is based on the random-phase approximation; the PZ-CA functional uses the von Barth-Hedin BH spin-dependence for the correlation which is correct for the exchange part of the functional.
Although the L(S)DA has proved successful in a variety of chemical and physical applications, it suffers certain well known deficiencies. Among those, the most serious are:
1. The tendency to produce more bonding in solids than is observed experimentally; manifestations of this tendency include the underestimate of the lattice constant or bond length and the overestimate of the cohesive energy and the bulk modulus.
2. Poor representation of activation energies which are related to chemical reactions or transitions between structures.
3. Incorrect relative stability of different magnetic phases for some magnetic materials.
In order to correct these deficiencies, expressions for the XC functional which go beyond the density and include gradients of the density have been devised. In the so-called generalized gradient approximation (GGA), the XC energy functional is expressed as follows:
$$\begin{array}{c}\hfill E_{\mathrm{XC}}^{\mathrm{GGA}}[\rho _{}(𝐫),\rho _{}(𝐫)]=\\ \hfill \rho (𝐫)ϵ_{\mathrm{XC}}^{\mathrm{GGA}}(\rho _{}(𝐫),\rho _{}(𝐫),\rho _{}(𝐫),\rho _{}(𝐫))\mathrm{d}^3𝐫.\end{array}$$
(6)
Applications of GGA to real materials show a tendency to over-correct the deficiencies of L(S)DA. For instance, the lattice constants of common crystalline solids tend to be overestimated within GGA calculations, while bulk moduli are underestimated juan.1993.agc ; juan.1995.ugg . We have implemented in HARES the recently developed parameter-free GGA functionals (PW91 pw91 ; perdew96 and PBE96 perdew96 ; pbe96 ; perdew98 ). Present capabilities include the use of the GGA functional in two different modes : (a) fully self-consistent GGA calculations and (b) a posteriori correction of the total energy with the perturbative GGA XC correction applied at the end of a self-consistent LDA calculation as:
$$\mathrm{\Delta }E_{\mathrm{total}}=E_{\mathrm{XC}}^{\mathrm{GGA}}[\rho ^{\mathrm{LDA}}]E_{\mathrm{XC}}^{\mathrm{LDA}}[\rho ^{\mathrm{LDA}}].$$
(7)
It is important to point out that the core–valence XC interaction is significantly different between LDA and GGA as was noted by Fuchs et al. fuchs98 Therefore, in order to ensure the reliability of the GGA results, it is necessary to perform either the fully self-consistent GGA calculation using GGA-constructed pseudopotentials (referred to as mode (a) above) or the a posteriori GGA correction after the self-consistent LDA calculation using the LDA-constructed pseudopotentials with the partial-core electron density (referred to as mode (b) above).
For the spin-polarized systems, we consider two different modes of the computation: (1) The conventional unconstrained calculation, where the total electron density and the magnetic moment are determined simultaneously and self-consistently; and (2) the fixed spin moment (FSM) method FSM ; moruzzi86 ; singh , which constrains the magnetic moment to be constant, but allows the possibility of different Fermi energies for the spin-up and the spin-down electron densities. The latter method has certain advantages: A series of FSM calculations with different magnetic moments provide the total energy as a function of the magnetic moment, yielding detailed information about the magnetic phase. In addition, the FSM calculations rapidly achieve self-consistency and are numerically more stable compared with the unconstrained calculations.
### II.2 Computational Approach
There exist a variety of methods for solving the set of single-particle equations derived from DFT, Eqs. (2). In the broadest classification, these methods fall into two categories, depending on how they describe the single-particle wavefunctions and charge density: Methods in the first category use explicit basis sets to represent the wavefunctions and charge density, while those in the second category use finite, discrete grids (or meshes) of points on which these functions are represented.jrc94 ; briggs95 ; baroni92 ; iyer95 ; hoshi95 ; gygi95 A standard approach of the first type employs a PW basis, which is a natural basis for periodic systems.ihm ; jrc96 The plane waves needed in the expansion are determined by the reciprocal lattice of the crystal while the number of plane waves included in the basis is determined by the highest kinetic energy, a parameter referred to as the “energy cutoff”. HARES falls in the second category of methods, as it employs a discrete mesh for the calculation. An important difference in the two types of methods is that in the former all operators have a unique representation once the basis set is chosen, whereas in the latter operators involving differentials have many possible representations with different order of approximation. In this sense, the latter type of methods involve an additional degree of approximation. Both types of methods map the Kohn-Sham problem onto a matrix eigenvalue problem, denoted by $`H_{\mathrm{KS}}`$. One of the desirable features of grid-based methods is to produce a sparse matrix $`H_{\mathrm{KS}}`$; this makes it possible to employ iterative algorithms for its solution.
#### II.2.1 Adaptive Coordinate Transformation
HARES uses a uniform grid in curvilinear space which is analytically mapped onto a grid in real space with resolution (grid-spacing) adapted to natural inhomogeneities in the problem. With the use of the adaptive grid, one can use HARES for both all-electron and pseudopotential calculations. However, use of pseudopotentials proves effective in most practical calculations. In the Kohn-Sham problem, inhomogeneities arise fundamentally from $`V_{\mathrm{ext}}(𝐫)`$, which is the potential that each electron experiences due to the presence of the nuclei. The Cartesian coordinates $`x^i(\xi ^\alpha ;P^m)`$ depend on the curvilinear coordinates $`\xi ^\alpha `$ and a set of parameters $`P^m`$ that allows tuning of the coordinate representation to a particular physical problem. The Jacobian of the transformation is
$$J_\alpha ^i(\xi ;P)=x^i/\xi ^\alpha $$
(8)
and describes how derivatives transform between the coordinate systems; its determinant $`|J|=detJ_\alpha ^i`$ is a measure of how the volume element is changed by the coordinate transformation. The metric giving the elemental length associated with infinitesimal displacement is given by
$$g^{\alpha \beta }=(J^1)_i^\alpha \delta ^{ij}(J^1)_j^\beta .$$
(9)
Details of the coordinate transformation can be found in Ref. acres, .
The gist of the transformation is to enhance spatial resolution in the region where it is desirable to increase the accuracy of the finite-difference derivatives and the representation of charge density inhomogeneities. The equivalent enhancement of resolution in the PW approach is the increase of the energy cutoff. The connection between the effective energy cutoff and the local resolution of the HARES grid is given by the factor $`|J|^{2/3}`$. The differential equation of the Kohn-Sham problem in the adaptive grid representation becomes a finite matrix eigenvalue problem, with only the kinetic energy term (the Laplacian in the single-particle equations) having off-diagonal elements.
The uniform mesh in $`\xi `$coordinates is subsequently broken into blocks that are distributed over a number of processors on a parallel computer architecture. The wavefunctions, potentials, and charge density are represented on this mesh allowing for balanced distribution on processors. In the iterative solution of the eigenvalue problem, an operation that is performed frequently during the calculation is the product of the Hamiltonian matrix with a vector representing a single-particle wavefunction. In parallel execution of this operation, it is the kinetic energy term (the Laplacian) with off-diagonal elements that requires most of the communication and makes the solution of the eigenvalue problem nontrivial.
#### II.2.2 Boundary Conditions
In DFT calculations based on a real-space grid, boundary conditions enter in the way the Laplacian is applied to a function. There are only two aspects of a calculation where this is relevant: (i) the kinetic energy operator, that is, the Laplacian acting on wavefunctions and (ii) the calculation of the electrostatic potential, obtained by solving the Poisson equation, that is, the Laplacian acting on the potential. Since the calculation of the Laplacian (represented as a finite difference) of a function at a given grid point uses values of the function at adjacent grid points, imposition of the boundary conditions requires knowledge of the function at a few grid points outside the boundary. For calculations on infinite crystalline solids, we use periodic boundary conditions (PBC) demanding that the function is periodic in space with the period of a unit cell that models a physical system. Thus, application of the Laplacian at any point in the unit cell involves values of the function inside the unit cell making implementation of boundary conditions for this case straight-forward. For calculations on finite systems (atoms, molecules or clusters), we use open boundary conditions (OBC). In this case, the treatment of boundary conditions is more intricate since the grid points adjacent to the ones on the boundary fall both inside and outside the region containing the system. In Fig. 1, we show how this is treated in HARES. We choose a rectangular box with a spherical region inside, the interior of which is large enough to contain the physical system. The thickness of the buffer region surrounding this sphere depends on the order of the finite-difference Laplacian, and is equal to this order times the grid-spacing. Thus, the distance between the sphere and the box vanishes in the continuum limit and the ratio of the volumes of the sphere to the box becomes $`\pi /6`$. The wavefunctions and charge density vanish outside the spherical region; this takes care of the boundary conditions in the kinetic energy part. There is one more term that deserves special attention in handling OBC: it is the electrostatic potential which has long range, and therefore cannot be assumed to be zero in the buffer region outside the spherical region. We obtain values for the electrostatic potential in the buffer region through a multipole expansion up to order 4, using the density inside the spherical region; we use the values at those points as the boundary conditions for $`V_\mathrm{H}`$ in the solution of the Poisson equation.
#### II.2.3 Scaling of Computational Effort
It is useful at this point to analyze the computational effort involved in various aspects of a DFT calculation using HARES and compare it to that performed using the PW method.
In the PW method, there is one grid in real (direct) space and another in reciprocal space. The wavefunctions are stored in reciprocal space on part of the grid inside a spherical region with diameter equal to half of the side of the full grid. These are transformed into real space using the fast Fourier transform (FFT) whenever necessary. The calculations of the charge density, the local ionic potential, and the exchange correlation energy are carried out in real space, whereas the calculation of the kinetic energy and the Hartree energy are carried out in reciprocal space. The calculation of the energy related to the nonlocal pseudopotential can be done on either grid. For small system sizes, the most time-consuming part is often that of performing the FFTs, which scales as $`O(N\mathrm{log}N)`$. The number of FFTs scales linearly with system size giving an overall $`O(N^2\mathrm{log}N)`$ scaling for the entire calculation. For large system sizes, the orthonormalization of wavefunctions, or equivalent constraints imposed during minimization of the energy functional, dominate the computational time. These operations scale as $`O(N^3)`$. Calculation of the contribution from the nonlocal pseudopotentials normally scales as $`O(N^3)`$, but can in principle be improved to $`O(N^2)`$ scaling by exploiting the short-range character of the potentials in real space.
In HARES, the wavefunctions are stored on the full grid in real space and all operations are performed on the same grid, eliminating the need for FFTs. The calculation of the kinetic energy is carried out using finite-difference formulae for derivatives on a grid in real space. The Hartree energy and the long-range electrostatic potential due to periodic charge density are computed by solving the Poisson equation, which scales as $`O(N)`$. For large system sizes, orthonormalization of the wavefunctions dominates the computational time, which then scales as $`O(N^3)`$. The treatment of the nonlocal pseudopotentials also scales as $`O(N^3)`$, unless their short-range character is exploited. We summarize the comparison of scaling between HARES and a PW method in Table 1.
The current parallel implementation of HARES is in high performance Fortran (HPF), which involves single instruction multiple data (SIMD) coding. Since the wavefunctions and charge density are stored in real space and distributed across processors, communication between processors is necessary in calculating:
1. the finite-difference derivatives using the CSHIFT operation, which cyclically shifts the data in an array on a grid along the specified direction; and
2. the inner product of two functions using the the SUM operation.
The scaling of inter-processor communication with system size is presented in Table 2. For large enough system size, the SUM operations dominate the communications in a parallel calculation.
## III Algorithmic Enhancements and Implementation
### III.1 Non-orthogonal Unit Cell
Finite-difference formulae for derivatives of functions represented on a grid with finite spacing are designed to achieve high accuracy and are reasonably accurate for a polynomial function up to certain order. These formulae are typically derived for functions of one variable on grids of uniform spacing. Their generalization to higher dimensions is trivial through direct product if the grid is orthogonal in the various dimensions. For periodic crystals with non-orthogonal unit cell, it is often not possible to design an orthogonal grid with the same periodicity in all directions; in this case the implementation of finite-difference formulae is not trivial.
The coordinate transformation employed in HARES provides a very simple method to treat non-orthogonal unit cells and grids. In this case, the transformation is uniform throughout the unit cell and maps an orthogonal unit cell in $`\xi `$space onto a non-orthogonal unit cell in $`x`$space. If $`F`$ is a matrix that gives the deformation of the orthogonal unit cell into the one under study (i.e. its columns are the non-orthogonal unit cell vectors), one can always obtain a transformation that is symmetric by filtering out the rotational part of $`F`$ as follows: first obtain an auxiliary matrix $`M`$ defined as $`M=FF^T`$ and then diagonalize $`M`$ to obtain a diagonal matrix $`D`$. The Jacobian for a rotation-free transformation is then given by
$$J=T^1D^{\frac{1}{2}}T,$$
where $`T`$ is a matrix that diagonalizes $`M`$: $`D=TMT^1`$. Once the mapping onto an orthogonal $`\xi `$grid is obtained, the derivatives, length and volume elements can be obtained using the formalism described in section II.2.1.
### III.2 Preconditioned Conjugate Gradients Solver
The dominant part of a DFT calculation often consists of solving an eigenvalue problem, that is, obtaining the lowest few (compared with the full spectrum) eigenvalues and eigenvectors of a very large matrix. In the case of PW basis, the size of the matrix is determined by the number of PW components included in the basis. In the case of HARES, the size of the matrix is determined by the number of points that constitute the real space grid. For small enough matrices the standard techniques of linear algebra can be employed, which give the exact (within the numerical accuracy of the algorithm) eigenvalues and eigenvectors of the matrix. When the size of the matrix is large, the conventional methods are not practical and the only alternative is to employ iterative approaches which approximate the eigenvalues and eigenvectors in successively improving steps. We considered two iterative algorithms for the diagonalization task in HARES:
1. An Inverse Iteration (II) algorithm with multigrid preconditioning norm\_thesis ;
2. A Conjugate Gradient (CG) algorithm conj with suitable preconditioning in real space.
The implementation of the former has been presented earlier acres and we want to focus on the CG algorithm in this subsection.
In Fig. 2, we present a flowchart of the CG algorithm. It is similar to the one in Ref. conj, , presented for a PW basis. In real space, most steps in the algorithm remain unchanged except for the preconditioning. The main idea in preconditioning is to filter out high Fourier components in the wavefunctions and the charge density. We achieve this through multiple application of a coarsening transformation. For example, the coarsening applied to the charge density gives:
$$\begin{array}{c}\hfill \rho (k,l,m)\frac{1}{2}\rho (k,l,m)+\\ \hfill \frac{1}{12}[\rho (k\pm 1,l,m)+\rho (k,l\pm 1,m)+\rho (k,l,m\pm 1)],\end{array}$$
(10)
where $`k,l,m`$ are indices of the grid points at which the charge density $`\rho `$ is calculated. We find that application of this transformation twice on the function under consideration results in adequate preconditioning. Better preconditioning is possible in principle for selected cases but may not be worth the extra effort required; this method provides a preconditioner that works reasonably well in all cases we have considered.
Another aspect of a DFT calculation is that the eigenvalue problem needs to be solved repeatedly while updating the charge density to achieve self-consistency between the wavefunctions and the corresponding effective potential, which depends on the density. Specifically, a self-consistent DFT calculation starts with an initial guess for the density and subsequently follows an iterative procedure of alternating steps of diagonalization and improved estimate for the charge density. Each iteration $`i`$ starts with a charge density $`\rho _{in}^i`$ and obtains the part of the eigenspectrum of $`H_{\mathrm{KS}}`$ that corresponds to occupied single-particle states and some low-lying unoccupied states, which depends on $`\rho _{in}^i`$ through the effective potential. At the end of an iteration the eigenfunctions are used to obtain an output charge density $`\rho _{out}^i`$. An improved estimate for the charge density is obtained from $`\rho _{in}^i`$ and $`\rho _{out}^i`$; this is discussed in detail in the following subsection.
We observe that diagonalization with full convergence, that is, with the tolerance for the difference between the $`\rho _{in}^i`$ and $`\rho _{out}^i`$ set to be very low, can be computationally demanding and depends on the initial guess for the density. Since the charge density at the initial stage of self-consistency is far away from the self-consistent one, accurate diagonalization of the Hamiltonian matrix at this stage is not worthwhile. Accordingly, we limit the number of CG steps for diagonalization at the initial stages of the self-consistency loop to relatively few — we found that two CG iterations at this stage yield optimal efficiency. Accurate diagonalization is achieved in the course of self-consistency as the initial guess for the eigensolver is steadily improved. We have found that the performance of the II and CG algorithms for convergence to self-consistency is comparable. The differences are rather small and system-dependent. The advantage of the CG algorithm is that all operations take place on the same grid, whereas in the II algorithm multi-grid preconditioning requires the mesh sizes to be a power of two.
### III.3 Charge Density Mixing
We return now to the way in which the charge density is updated at the end of each iteration. In general, the new charge density at the end of step $`i`$ can be constructed from the charge densities of previous steps:
$$\rho _{new}^i=\underset{j=id}{\overset{i}{}}(\kappa _j^{in}\rho _{in}^j+\kappa _j^{out}\rho _{out}^j)=\rho _{in}^{i+1},$$
(11)
where $`\kappa _j`$’s are mixing coefficients and $`d`$ is called the “depth” of the mixing procedure, i.e. the number of previous iterations used in the improved estimate of the density. Various mixing schemes are available and have been discussed in Ref. kresse, .
In our work, we added another feature to mixing: we optimize the strength of mixing, that is, the values of the parameters $`\kappa _j`$, as a function of iteration. This feature can be used along with most of the mixing schemes employed in the literature. To illustrate the basic idea, we take a simple realization of a mixing scheme:
$$\rho _{new}^i=\rho _{in}^i+\kappa (\rho _{out}^i\rho _{in}^i).$$
(12)
$`\kappa =1`$ is an extreme case where no knowledge of the input density is used, and $`\kappa =0`$ corresponds to the opposite extreme where the density is not updated at all. For $`\kappa =1`$, there can be oscillations between input and output densities corresponding to underdamped mixing. On the other hand, for small $`\kappa `$ the oscillations are overdamped resulting in slow update of the density and hence the approach to self-consistency. We have devised a way to achieve the optimal critical damping in the mixing procedure. In our scheme, we calculate the root mean square change in density $`\delta \rho ^i`$ at each iteration $`i`$:
$$\delta \rho ^i=\left(\frac{1}{\mathrm{\Omega }}|\rho _{out}^i\rho _{in}^i|^2\mathrm{d}^3𝐫\right)^{\frac{1}{2}},$$
where $`\mathrm{\Omega }`$ is the unit cell volume, and define the rate of self-consistency as
$$R_i\frac{\mathrm{log}(\delta \rho ^i)}{t},$$
where $`t`$ is a fictitious time associated with iterations. $`R_i`$ is a rough measure of how well a calculation is evolving toward a self-consistent solution. We monitor both the rate and the mixing coefficient at each iteration and can estimate $`\lambda R/\kappa `$. If $`R_i`$ is too small ($`<0.2`$), the approach to self-consistency is too slow. In that case, the mixing coefficient is increased or decreased, depending on the sign of $`\lambda `$, by an amount $`\mathrm{\Delta }\kappa `$, — a positive $`\lambda `$ corresponds to smaller $`\kappa `$ and a negative $`\lambda `$ to larger $`\kappa `$ compared with the optimal value of this parameter. This has the effect of keeping the strength of mixing $`\kappa `$ near a value where $`R`$ is optimal i.e., in the neighborhood of critical damping. The magnitude of change $`\mathrm{\Delta }\kappa `$ is reduced over time to make sure that it converges to its optimal value.
### III.4 Nonlocal Pseudopotential Projector
The nonlocal part of the pseudopotentials $`V_{NL}(𝐫,𝐫^{})`$ is used in a separable form kleinman ; gonze to facilitate fast calculation of its product with the wavefunction $`\psi `$:
$$𝐫|V_{NL}|\psi =\underset{\alpha ,l,m}{}\frac{𝐫|\varphi _{lm}^\alpha \varphi _{lm}^\alpha |\psi }{\eta _{lm}^\alpha },$$
where $`\alpha `$ is an atomic index, $`l`$ and $`m`$ are angular quantum indices, and $`\eta _{lm}^\alpha `$ are constants related to the normalization of the nonlocal projectors $`\varphi _{lm}^\alpha `$. Straightforward application of $`V_{NL}(𝐫,𝐫^{})`$ on the wavefunctions scales as $`O(N^3)`$, since there are $`N_e`$ wavefunctions on $`N`$ grid points, and $`N_a`$ atoms (the number of grid points $`N`$ is proportional to the system size measured by either $`N_a`$ or $`N_e`$, which also scale with each other). This computation can be accelerated significantly, by noticing that $`𝐫|\varphi _{lm}^\alpha `$ is a localized function centered on atom $`\alpha `$. Thus, a calculation of $`\varphi _{lm}^\alpha |\psi `$ in real space involves only the grid points near atom $`\alpha `$ making it an $`O(N^0)`$ computation, which gives $`O(N^2)`$ for the overall calculation involving the non-local pseudopotential. We use a filtered pseudopotential approach bernholc:pseudopotentials with a filter that is smooth in $`k`$-space (as opposed to the theta-function used in plane-wave approaches). This minimizes the errors in the evaluation of $`\varphi _{lm}^\alpha |\psi `$ in addition to using the adaptive grid. The parallel implementation of such a calculation is not trivial, since it involves only those processors which store the grid points near the given atom $`\alpha `$. To address this issue, we represent the atomic projectors $`\varphi _{lm}^\alpha (𝐫)`$ in terms of the packed projectors $`\chi _{lm}^j(𝐫)`$ shown schematically as follows:
$$\begin{array}{c}\hfill \left\{\varphi _{lm}^\alpha (𝐫)\right|\alpha =1,\mathrm{},N\}\\ \hfill \left\{\chi _{lm}^j(𝐫),\beta _{lm}^j(𝐫)\right|j=1,\mathrm{},M_d\}\end{array}$$
(13)
where $`\chi _{lm}^j(𝐫)`$ is defined as the projector of the $`j`$-th largest magnitude at a given grid point $`𝐫`$ (e.g. $`\varphi _{lm}^\gamma (𝐫)`$), and $`\beta _{lm}^j(𝐫)`$ is the index of the atom from which $`\chi _{lm}^j(𝐫)`$ was generated (in the case of the above example $`\beta _{lm}^j(𝐫)=\gamma `$). We only keep a number of important pseudoprojectors $`M_\mathrm{d}`$, (thus letting $`j`$ vary from 1 to $`M_\mathrm{d}`$), which we call the depth of the packed projectors. We find that $`M_\mathrm{d}=3`$ is typically sufficient; for example, this is exact if the nonlocal projectors of at most three atoms are nonzero at any point $`𝐫`$. This changes the scaling of memory requirements for the nonlocal potential from $`O(N^2)`$ to $`O(N)`$. With this choice of packing, the expression for the inner product becomes:
$$\varphi _{lm}^\alpha |\psi =\underset{j}{}\underset{𝐫}{}\delta _{\beta _{lm}^j(𝐫),\alpha }\chi _{lm}^j|𝐫𝐫|\psi .$$
This is readily evaluated using an EXTRINSIC subroutine call in HPF, which involves execution of the whole routine on each processor, but on different data. Effectively, an inner product $`\varphi _{lm}^\alpha |\psi `$ with contribution from only the grid points inside a sphere centered at atom $`\alpha `$ is calculated by distributing the data with respect to $`\alpha `$ rather than grid points.
### III.5 Computation of Forces
Forces on the atoms are calculated using the Hellman-Feynman theorem, as is usual in DFT calculations:
$$𝐅_\alpha =\underset{i}{}\psi _i\left|\frac{V_{\mathrm{ext}}}{𝐑_\alpha }\right|\psi _i,$$
where $`𝐑_\alpha `$ is the position of atom $`\alpha `$. For reasons similar to those mentioned in the previous subsection, the contribution of the nonlocal pseudopotential to atomic forces is computationally demanding and scales as $`O(N^3)`$. With packed projectors $`𝐫|\chi _{lm}^j`$, the scaling of this computational cost is improved to $`O(N^2)`$. This is a significant improvement for problems involving structural relaxation of large systems.
The implementation of packed projectors in the calculation of forces deserves further elaboration. The force calculation with nonlocal pseudopotentials involves both the projectors $`\varphi _{lm}^\alpha (𝐫)`$ and their derivatives $`\varphi _{lm}^\alpha /𝐑_\alpha `$. Since the latter is needed only during the calculation of forces, it does not need to be packed and stored but can be obtained at the time when it is needed. The inner products of $`\varphi _{lm}^\alpha `$ with $`\psi _i`$ have to be computed for all atoms $`\alpha `$ at once at the beginning of a force calculation since they are packed. With these improvements, we obtain a factor of 7 speedup in the calculation of forces for systems containing about 30 atoms.
Finally, we should note that as in any method that involves a computational basis which changes with the positions of the atoms, the adaptive grid generates fictitious forces referred to as Pulay forces. We have implemented the correction related to the Pulay forces and found that it is not significant when the grid is refined to the point where it yields adequate accuracy. Thus, for all practical purposes, with an accurate grid the Pulay correction to the forces can be neglected (see also Ref. acres, ).
### III.6 Geometry Optimization
In modern ab initio total energy calculations, one of the objectives is to obtain minimum energy geometries (corresponding to local minima or, if possible, the global minimum of the energy), where the Hellmann-Feynman force on each atom is zero or, more precisely, smaller in magnitude than a prescribed value (typically $`0.5`$ mRy/a.u.). The process in which the initial ionic geometry is sequentially updated to relax to a neighboring local minimum is referred to as the ionic relaxation. Mathematically, the ionic relaxation is a nonlinear optimization problem, which is a subject of vast interest in applied mathematics.
Many optimization algorithms have been proposed so far, which can be broadly classified into three groups:
1. algorithms which require only the evaluation of the function;
2. algorithms requiring the function values and its gradients;
3. algorithms requiring the function values and its first and second derivatives (the gradients and the Hessians).
In electronic structure calculations, the function to be optimized is the total energy and its gradients are the forces on the atoms. Since the forces are not too computationally demanding compared to the self-consistency loop (especially after implementing the improvement discussed in the previous subsection), it is natural to use the methods of class (ii) for ionic relaxation.
Among the gradient algorithms, the quasi-Newton (also referred to as variable metric) method is known to be most efficient brodlie77 . The inverse Hessian is approximated and updated at each iteration. Suppose $`|R_0`$ is an initial estimate of the minimizer of the total energy $`E_{\mathrm{total}}`$, $`|g_0`$ is the corresponding gradient, and $`H_0`$ is the initial guess for the inverse Hessian. At the $`n`$-th relaxation step, the next approximate minimizer is given by
$$|R_{n+1}=|R_n\beta _nH_n|g_n$$
(14)
where the step size $`\beta _n`$ is determined by the line search walsh75 (or the line minimization). The inverse Hessian is updated by
$$H_{n+1}=H_n+\mathrm{\Delta }_n$$
(15)
where $`\mathrm{\Delta }_n`$ is the correction to $`H_n`$ and is determined by requiring that it satisfies the quasi-Newton condition,
$$H_{n+1}|h_n=|d_n$$
(16)
with $`|h_n|g_{n+1}|g_n`$ and $`|d_n|R_{n+1}|R_n`$.
Among different update formulae for the inverse Hessian brodlie77 , we have implemented the initially-scaled Broyden-Fletcher-Goldfarb-Shanno (IS-BFGS) expression shanno78 ; chetty95 :
$$\begin{array}{c}\hfill H_{n+1}=\stackrel{~}{H}_n\frac{\stackrel{~}{H}_n|h_nd_n|+|d_nh_n|\stackrel{~}{H}_n}{d_n|h_n}\\ \hfill +\left(1+\frac{h_n|\stackrel{~}{H}_n|h_n}{d_n|h_n}\right)\frac{d_n|d_n}{d_n|h_n}\end{array}$$
(17)
where $`\stackrel{~}{H}_n=(d_n|h_n/h_n|H_n|h_n)H_n`$, when $`n=0`$ and $`\stackrel{~}{H}_n=H_n`$, otherwise.
Implemented in combination with the approximate line search algorithm fletcher70 , the IS-BFGS provides an efficient ionic relaxation tool which assures the convergence of the approximate inverse Hessian to the correct inverse Hessian and is numerically stable. As noticed in other works powell77 ; shanno78b , restarting the update sequence \[Eq. (17)\] can be beneficiary in some cases, and we have used the following restart criterion: $`d_n|g_{n+1}0.5\sqrt{d_n|d_ng_{n+1}|g_{n+1}}`$. This implies that the displacement of an atom is not too different from the direction of the calculated force acting on it.
### III.7 Dual real-space grid calculations
Typically, the representation of the charge density and local potentials in a DFT calculation needs twice as much spatial resolution as that of the wavefunctions. To exploit this aspect of electronic structure calculations, we have developed a version of HARES which employs two separate real-space grids — a coarser one and a finer one. The wavefunctions are represented on the coarser grid and the charge density on the finer one with half the grid-spacing of the former. The finer grid corresponds to the FFT grid in a PW calculation. This version of the code reduces memory requirements substantially. The computational cost of the worst scaling part of the calculation (the wavefunction orthogonalization) is reduced by a factor of 8.
The transformation from the coarser to the finer grid is performed only when the charge density needs to be calculated from the wavefunctions. We use wavelet interpolants arias to achieve this. The Poisson equation is solved on the finer mesh to obtain the electrostatic potential. The exchange correlation potential and the local part of the pseudopotential are also calculated on the finer grid; these terms are convolutions in $`k`$-space and in a PW method are calculated on the FFT grid. The kinetic energy operator and the nonlocal pseudopotential act on the wavefunctions directly on the coarser grid. Naturally, this necessitates usage of a higher order Laplacian in the calculation of the kinetic energy, while the one with lower order is adequate in the solution of the Poisson equation.
To check the accuracy of the dual-grid code, we calculated the energy difference between two configurations of the O<sub>2</sub> molecule and compared it with the result of a PW calculation. Both methods were used at different energy cutoffs, or equivalently, of spatial resolution. We found that the energy difference as a function of the energy cutoff behaves the same way in both methods. In fact, at a low energy cutoff, the sign of the energy difference was inverted in both calculations and the magnitude was within 6 % of the correct value.
While the dual-grid approach to real-space electronic structure calculation enhances the performance significantly, we caution that the errors introduced by breaking of translational invariance, a feature inherent in real-space grids (see Ref. acres, ), are larger than those in the single finer grid calculations. This is due to the coarser grid used in the representation of wavefunctions and the higher order expression for the Laplacian on the coarser grid. Within the PW method, the same error enters in calculating the XC potential on the real-space FFT mesh. These errors enter into the calculation of forces and can be minimized if necessary by using Fourier-filtered pseudopotentials. We expect the use of the dual-grid approach to be very advantageous at the initial stages of ionic relaxation of a system with large number of atoms, when the accuracy in forces or wavefunctions is not crucial since the system is presumably far from its optimal structure.
### III.8 Performance
In Table 3, we present a comparison of the performance of HARES and that of a PW code, for DFT calculations on the O<sub>2</sub> molecule and a zeolite, Si<sub>24</sub>O<sub>48</sub>. We have selected for the comparison the academic version of CASTEP castep , a PW package which uses all the standard methods for such calculations and is freely available to academic researchers, as is also the case for HARES. We believe this provides the most meanigful comparison of the performance of the different approaches, for codes at equivalent levels of development and availability to the academic community. It is clear that the performance of HARES for the oxygen molecule is definitely better than the PW code. In general, we find that HARES performs better for metallic systems. For the zeolite, which is an insulator, CASTEP uses a variational method for direct minimization of the total energy, which is not applicable to the case of metals. This makes the performance of CASTEP very good for such systems, though the performance of HARES is not unacceptable (a factor of 1.66 slower than CASTEP). This advantage of CASTEP over HARES is lost when applied to metallic systems, in many of which the academic version of CASTEP available to us failed to converge in the self-constistency loop. We have been informed that in commercial versions of this package the problematic convergence to self-consistency in metallic systems has been solved and the performance has been improved castep-extra . We point out that a comparison of ab initio packages, is meaningful only for methods that employ the same type of pseudopotentials. In the comparison discussed here, both methods use conventional norm-conserving pseudopotentials. There exists another class of pseudopotentials, called ultra-soft pseudopotentials dhv , which reduce the size of the Hamiltonian matrix substantially, making the calculations more efficient. Codes that employ this class of pseudopotentials are naturally faster, whether they use a PW basis or a real-space grid. For instance, the VASP code based on the PW formulation uses ultra-soft pseudopotentials and has proven quite effectivevasp ; the commercial version of CASTEP also uses these pseudopotentials castep-extra . This class of pseudopotentials has not been yet implemented in HARES.
## IV Applications
As a test of the accuracy and the efficiency of the algorithmic improvements discussed above, we offer a range of example applications of HARES. These include representative elemental crystals, some molecules, and a couple of rather complex materials — blue molybdenum bronze and the TON zeolite. All the calculations were performed on a Silicon Graphics Origin 2000, using from 2 to 16 processors in parallel mode.
### IV.1 Study of elemental solids
The simplest test of the method is its application to elemental crystalline solids. We have calculated the basic structural and electronic properties for representative elemental solids, including alkali metals (Li, K), group II A metals (Be, Ca), sp-electron metals (Al, Ga), d-electron non-magnetic metals (V, Cu, Mo), d-electron magnetic metals (Fe, Ni), and semiconductors and insulators (Si, C). The properties of these solids are extracted from total energy calculations for a given crystal structure, using the LDA and applying the a posteriori GGA corrections. We have used norm-conserving pseudopotentials from Bachelet et al bachelet.1982.pwh for V, and pseudopotentials generated with the Troullier and Martins troullier91 scheme for the all the other elements. We perform the calculations as follows: we choose a sufficiently dense grid of k-points in the Monkhorst-Pack scheme monkhorst.1976.spb and fold it to the irreducible part of the Brillouin zone by applying the symmetry operations of the point group of the crystal including inversion which is always a symmetry operation in reciprocal space. We make sure the calculation is converged with respect to the real space grid spacing in the neighborhood of the anticipated equilibrium lattice constant, and keep the grid spacing approximately constant for a range of lattice constants up to about twice the equilibrium value. We then fit the resulting energies to powers of $`\mathrm{\Omega }^{2/3}`$, where $`\mathrm{\Omega }`$ is the volume of the unit cell. We thus obtain accurate values of the equilibrium lattice constant and minimum total energy. These values are used to fit the two-parameter Universal Binding Energy Relation rose.1984.ufe , which has a simple analytical form from which the bulk modulus and the cohesive energy are obtained. This procedure relies on the fact that the total energy differences used to fit the Universal Binding Energy Relation converge quicker than the total energy with respect to the grid parameters. For the nonmagnetic materials we use the non-spin-polarized code to perform these calculations for reasons of computational efficiency. The free atoms in most of the cases we calculated are polarized, the extreme case being Molybdenum were all six valence electrons have the same spin. We thus add to the cohesive energy as calculated previously the difference of a free atom calculation using the spin-polarized and spin-average codes. The free atom calculations can be converged to the same, high degree of accuracy by use of the adaptive grid and thus the relative energies of the spin polarized and unpolarized atoms are evaluated on an equal footing with the other energy differences.
In Table 4, we summarize the results of the calculations. The experimental values are a compilation of results from Kittelkittel.1996.iss ; for certain elements we considered a simpler lattice than the experimental ground state, and in these cases we compare our results to the all-electron calculations of Moruzzi et almoruzzi.1978.cep . Note that, as expected, the LDA results for the lattice constant are lower, and for the bulk modulus are higher than the experimental values. The gradient correction tends to improve the situation. Overall, the agreement with experiment is quite good, except for the cohesive energies, which is a well known deficiency of the approach we are using. Looking at the gradient corrected results, other than Ca which seems to have large offsets, the lattice constants are within 2.0% of the experimental values, bulk moduli within 22%, and cohesive energies are within 26%. In Fig. 3 we present the scaled results of the gradient corrected calculations and the corresponding universal curve. The magnetic moments of the ferromagnetic materials at their equilibrium lattice constant values as calculated with the LDA (GC) are 2.06 (2.30) $`\mu _\text{B}`$/atom for Fe and 0.57 (0.59) $`\mu _\text{B}`$/atom for Ni. These values compare well with the experimental values of 2.22 $`\mu _\text{B}`$/atom for Fe and 0.61 $`\mu _\text{B}`$/atom for Ni, as well as with the summary of the results of various theoretical methods that is presented in Ref. moroni:1997:upa, .
### IV.2 Study of small molecules
In Sec. II.2.2, we described how two types of boundary conditions — OBC and PBC — can be readily used in a HARES calculation. Many DFT methods designed to do calculations for solids use PBC and are constrained to use a large supercell to study an isolated molecule or a cluster of atoms. Here, we present results for four molecules N<sub>2</sub>, O<sub>2</sub>, H<sub>2</sub>O, and NH<sub>3</sub> obtained using HARES with the two types of boundary conditions keeping all other computational parameters fixed. We use a periodic box of dimensions 24 $`\times `$ 18 $`\times `$ 18 a.u.<sup>3</sup> for the N<sub>2</sub> and O<sub>2</sub> molecules and 20 $`\times `$ 20 $`\times `$ 20 a.u.<sup>3</sup> box for the H<sub>2</sub>O and NH<sub>3</sub> molecules, with a grid spacing of 0.25 a.u.
In Table 5, the various results for small molecules are summarized. In ammonia, we also calculated the inversion barrier of the potential energy surface, by relaxing the positions of the hydrogen atoms in the plane for selected heights of the Nitrogen atom. We find that the bond-lengths obtained with OBC tend to be smaller than those obtained with PBC, though the difference is quite small, in most cases smaller than the accuracy in the reported results. The energies, on the other hand, have significantly larger differences. We suggest that this is due to the electrostatic interaction of the field in a supercell calculation with PBC. This interaction between the molecule and its periodic images changes the energy of the molecule. In the examples considered here a cubic cell geometry is used, which results in zero dipole interaction (it can be shown analytically that the interaction energy of a dipole with a full shell of dipoles is zero). This indicates that the discrepancy is due to higher-order multipole terms. The computational time for the calculations with PBC and OBC is similar, so there is no particular advantage to either approach from the point of computational cost. It appears, however, that for truly isolated systems the OBC approach gives more realistic results due to the absence of any spurious long range fields.
### IV.3 Electronic Structure of Blue Bronze, K<sub>3</sub>Mo<sub>10</sub>O<sub>30</sub>
The material called blue bronze (BB), whose chemical composition is A<sub>0.3</sub>MoO<sub>3</sub> with A an alkali metal, exhibits a variety of interesting physical properties including a metal-to-semiconductor transition at $`T_c=180`$ K, quasi-one-dimensional electronic properties above $`T_c`$, and the existence of incommensurate and commensurate charge density wave (CDW) phases schlenker89 . Recently, a family of molybdenum bronzes has been extensively studied in experiments using angle-resolved photoemission spectroscopy (ARPES) to explore a possible realization of non-Fermi-liquid behavior due to its low-dimensional electronic properties claessen95 ; gweon96 ; denlinger99 . To our knowledge, the only published electronic band calculation of BB is based on a tight-binding (TB) method using some model structures whangbo86 . The dispersion of the TB bands around the Fermi level is qualitatively different from the ARPES experimental results gweon96 . It is of great importance to have an accurate and reliable ab initio calculation of the electronic structure of BB in order to interpret the ARPES measurement in terms of possibly interesting physics. With a large number of atoms in the unit cell, including 10 Mo atoms and 30 O atoms which are typically difficult to handle with PW approaches, BB provides a challenging system for performing state-of-the-art ab initio calculations of the electronic structure; this requires a highly efficient computational tool such as HARES.
The structure and the lattice parameters of BB is well documented in Ref. ghedira85, : the Bravais lattice is centered monoclinic (CM), the space group is C2/m, and the lattice constants of the simple monoclinic cell are $`a_2=16.2311`$ Å, $`a_3=7.5502`$ Å, and $`a_1=9.8614`$ Å with the angle $`\beta =94.895^{}`$ between $`𝐚_1`$ and $`𝐚_2`$. The basic building block of BB is the MoO<sub>6</sub> octahedron; ten octahedra form a rigid unit by edge-sharing. Within the simple monoclinic (SM) cell, two rigid units are arranged so that one of them is located at the apex and the other at the center of the cell. As a result, neighboring rigid units share a corner oxygen to form a slab spanned by $`𝐚_2𝐚_1`$ and $`𝐚_3`$ and four infinitely-connected MoO chains (per primitive unit cell) parallel to $`𝐚_3`$. The crystal structure is illustrated in Fig. 4 where the simple monoclinic cell is indicated by a box; the two different classes of octahedra are indicated by different colors: yellow for those that participate actively to conduction along the high-conduction direction ($`𝐚_3`$) and blue for those that are apparently inactive.
We have performed electronic structure calculations for BB with HARES. For these calculations we use the Ceperley-Alder XC functional as parametrized by Perdew and Zunger CA ; PZ-CA . The ions are represented by norm-conserving pseudopotentials generated by the Troullier-Martins scheme troullier91 in the fully separable form of Kleinman and Bylander kleinman ; gonze . The Brillouin zone (BZ) integrations are performed using a $`4\times 4\times 4`$ Monkhorst-Pack k-point meshmonkhorst.1976.spb in the BZ of the CM cell. We used a grid spacing of $`h0.33`$ a.u. which gives around 200,000 real-space grid points.
In Fig. 5, we show isosurfaces of the valence electron density obtained from the fully self-consistent calculation. A few interesting features can be observed from the green isosurfaces which correspond to high electron density:
1. the high density region has a cylindrical shape with its axes overlapping with MoO infinite chains;
2. the electron density is “marginally” connected along the \[$`\overline{1}10`$\] crystallographic direction, which is within the slab but perpendicular to the chain direction;
3. the electron density is “barely” connected along the slab normal.
From these observations, we expect that electronic conduction will be highly anisotropic and the chain direction is the most favored. On the other hand, the blue isosurfaces, corresponding to low electron density, indicate that the valence electrons are largely depleted around the potassium ion sites, suggesting that K atoms play the role of donors.
The full spectrum of the energy bands is shown in Fig. 6. The lower valence band manifold (30 bands) has O($`2s`$) orbital character, the upper valence band manifold (90 bands) has O($`2p`$) character mixed with Mo($`4d`$) orbitals near the top of the energy range. There are two bands crossing the Fermi level shown in red, well separated from both the valence band and the conduction band manifolds. These bands disperse primarily along the chain direction which is indicated by gray panels in Fig. 6.
The two partially filled bands of our ab initio calculation are qualitatively different from the TB bands whangbo86 as illustrated in Fig. 7. For instance, the two LDA bands cross each other near the BZ boundary whereas the TB bands do not. Within the SM BZ (from -$`z_0`$ to $`z_0`$), the LDA bands have occupied bandwidths of 1.3 and 0.3 eV while both of the TB bands disperse by 0.3 eV. The occupied bandwidth of 1.3 eV for the low-lying band is in good agreement with the APRES datagweon96 . These qualitative differences in the ab initio and TB bands originate from the correct and unbiased description of the interactions and the use of a realistic atomic structure in our calculations. Our analysis of the wavefunction character at the $`\mathrm{\Gamma }`$-point shows anisotropic hybridization between the Mo $`d`$-states and the O $`p`$-states. It would be rather difficult to describe this situation within the context of the TB approach.
The significant dispersion of the partially filled bands only along the chain direction results in planar Fermi surfaces, which are nested along the chain direction. In turn, the nested Fermi surfaces induce a CDW. The estimated CDW wave vector is $`0.75`$ $`𝐛_3`$ compared with the observed one in the range of $`0.72`$$`0.75`$ $`𝐛_3`$. A more detailed analysis of the physics of this material will be presented elsewhere bb2000 .
In summary, our application of HARES to the electronic structure of BB suggests that an accurate and reliable method with a realistic atomic structure is needed in order to investigate the behavior of such complex materials. The ab initio energy bands are in good agreement with the ARPES measurement and the nature of the electronic states relevant to conduction can thus be elucidated.
### IV.4 Zeolite: Na<sub>n</sub>Al<sub>n</sub> Si<sub>24-n</sub>O<sub>48</sub>
The word “zeolite” (of Greek origin) means “boiling stone” and derives from the visible loss of water when natural zeolite minerals are heated. Zeolites are materials with unique properties which make them useful in a variety of applications such as oil cracking, nuclear waste management, catalysis and animal feed supplements. They form a well-defined class of naturally occurring crystalline alumino-silicate minerals. They have elegant three-dimensional structures arising from a framework of \[SiO$`{}_{4}{}^{}]^4`$ and \[AlO$`{}_{4}{}^{}]^5`$ coordination tetrahedra linked at their corners. The frameworks are generally very open and contain cavities that enclose cations and water molecules. The presence of cavities make zeolites porous and gives rise to their low density and unique properties. Since the cations, water or other molecules that can be contained in these cavities, interact weakly with the cavity walls, these entities have high mobility in the solid zeolite. As a result a number of interesting physical and chemical properties arise: facile ion exchange, easy water loss upon heating, molecular sieve behavior, etc.
In this Section, we study a zeolite which is referred to by the code name “TON” zeo\_atlas and has the general chemical formula Na<sub>n</sub>Al<sub>n</sub> Si<sub>24-n</sub>O<sub>48</sub>; we will consider the structures for $`n=0`$ and $`n=1`$. The framework of tetrahedra in its crystal structure is displayed in Fig. 8. TON has an orthorhombic crystal structure with three different types of pores or channels parallel to the $`c`$-axis. For $`n=0`$, its unit cell has 24 formula units of SiO<sub>2</sub> with a volume of about 1320 Å<sup>3</sup>. Its space group is Cmc2<sub>1</sub> and four of the 24 formula units are symmetry-independent. For $`n=1`$ (called Theta-1), the structure has been determined from X-ray powder experiments acta\_zeo . Theta-1 is the first reported unidimensional medium-pore high-silica zeolite.
Starting with the experimental geometry acta\_zeo , we relaxed the atomic structure of Si<sub>24</sub>O<sub>48</sub> using HARES. All the bond lengths obtained from the calculation are within 2% of the experimental values. We next considered four independent Si sites where an Al atom can be substituted for a Si atom to obtain AlSi<sub>23</sub>O<sub>48</sub>, and relaxed its atomic structure. Interestingly enough, we find that all four possible structures have very similar energies. Within the accuracy of our calculations the four sites cannot be differentiated.
Addition of a sodium atom to the AlSi<sub>23</sub>O<sub>48</sub> structure introduces a variety of structural possibilities. We explored three possible structures, based on the three types of cavities in which the Na atom can be placed. In Fig. 9, we show a unit cell of a structure with Na added in the largest of the three cavities. Since AlSi<sub>23</sub>O<sub>48</sub> is missing one electron due to the substitution of a Si atom by an Al atom, the added Na atom will naturally prefer to stay in the vicinity of the Al atom to which is can donate its valence electron. In the relaxed structure, the Na atom sits closest to three O atoms that are bonded to the Al atom and as a result the Al-O bonds are elongated.
In Fig. 9, we also show an isosurface of the electron density for NaAlSi<sub>23</sub>O<sub>48</sub> and compare it with that of Si<sub>24</sub>O<sub>48</sub>. The two systems have very similar charge distribution except in the region localized near the Na and Al atoms. The isosurfaces clearly have the same topology as the geometrical structure in Fig. 9: O atoms are at the center of the bulging regions of the isosurface and Si atoms at the joints or vertices. This indicates that the bonding is primarily ionic, with negatively charged O and positively charged Si atoms. Partial covalent character is also evident from the fully connected isosurface.
In NaAlSi<sub>23</sub>O<sub>48</sub>, the valence electron donated by Na compensates for the one missing electron in the four bonds formed by the Al atom. The nature of bonding of Al with O atoms on the opposite side of the Na atom is clearly different from that with the three O atoms on the Na side. The latter is very similar to the bonding character between Si and O in Si<sub>24</sub>O<sub>48</sub>. The charge on both Al and Na is positive, which has the effect of displacing the Al atom slightly away from the Na atom resulting in longer Al-O bonds. This introduces small structural distortions and changes in charge distribution in the neighboring SiO<sub>4</sub> tetrahedra. Since the addition of Na results in compensating electrostatic and covalent interactions, we expect that the energy barrier in the process of attachment of Na (or in general a cation) to the walls of cavities in this zeolite should be very small. Further investigation of the chemical activity inside these pores and its effect in the electronic structure of the zeolite will be the subject of future studies zeo2000 .
## V Summary
In this paper, we provided a comprehensive review of the theory underlying HARES, which is a method for ab initio electronic structure calculations implementated using HPF on a shared memory parallel computer architecture. Several applications of the method to calculate the properties of simple and complex physical systems were presented to illustrate its capabilities. We obtained the bulk features of elemental solids such as equilibrium lattice constant, bulk modulus and cohesive energy, for elements from many different columns of the Periodic Table, and find good agreement with experiment within the limitations of DFT/LDA calculations. For the small molecules N<sub>2</sub>, O<sub>2</sub>, H<sub>2</sub>O, and NH<sub>3</sub>, we find that the structural features do not depend on boundary conditions (open or periodic) used in the calculation, while the energy is sensitive to the the choice of boundary conditions. Application of the method to blue molybdenum bronze and a zeolite demonstrate that it can be used effectively to study complex material systems. In the case of Blue Bronze the results help to clarify important issues of the electronic structure pertaining to recent experiments. In our study of the TON zeolite Si<sub>24</sub>O<sub>48</sub> and its variations containing Al and Na atoms, we demonstrated the ability of the method to capture the nature of bonding between a cation and the walls of cavities in the zeolite; such interactions are related to the mobility of ions and molecules inside pores of the zeolite framework and should give rise to interesting physical and chemical behavior.
## Acknowledgements
The original development of the adaptive-grid real-space method was supported by the Office of Naval Research through the Common High-performance Scientific Software Initiative (CHSSI) and the High Performance Computation Modernization Office. This set of codes is available upon request <sup>1</sup><sup>1</sup>1email:kaxiras@cmt.harvard.edu. The subsequent development of HARES was funded by Ryoka Systems Inc. The authors wish to acknowledge useful discussions and collaborations with Melvin Chen and Greg Smith, and useful comments from Nick Choly, Ioannis Remediakis, Jose Soler and G.-H. Gweon.
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# 1 Introduction
## 1 Introduction
We consider the $`A_{𝒩1}^{(1)}`$ open spin chain ($`𝒩>2`$) with N sites , . This is an integrable system which has been solved by means of the nested Bethe ansatz , . Also, this model, in the critical regime, can be thought as the lattice analogue of a certain two dimensional field theory i.e., the $`A_{𝒩1}^{(1)}`$ affine Toda field theory . We focus here in the special case that the chain has $`U_q\left(SU(𝒩)\right)`$ symmetry , , and we solve it using a simpler method compared to nesting, namely, the analytical Bethe ansatz , .
The analytical Bethe ansatz has been used for models with crossing symmetry e.g. $`A_1^{(1)}`$, $`A_2^{(2)}`$, $`A_{2n}^{(2)}`$, (see , ). This is the first time that this method has been used for a model without crossing symmetry i.e., the $`A_{𝒩1}^{(1)}`$ spin chain. Although this chain does not have such symmetry, the corresponding $`R`$ matrix satisfies a crossing property -, and consequently one can show that the transfer matrix satisfies an analogous property. One can also generalise the results of , and derive a fusion procedure for the corresponding open chain transfer matrix. The crossing property, the fusion of the transfer matrix, and the quantum group symmetry play an essential role in the derivation of the analytical Bethe ansatz.
The outline of the paper is as follows: In section 2 we describe the model and we derive the crossing property for the $`R`$ matrix and the transfer matrix. In the next section we deduce the fusion procedure for the open chain transfer matrix. In section 4 we derive the asymptotic behaviour of the transfer matrix and together with the results of the previous sections, periodicity and analyticity we find the spectrum of the transfer matrix and the Bethe ansatz equations. We illustrate the method using the $`A_2^{(1)}`$ chain, but the results can nevertheless be generalised for any $`𝒩`$. Finally, in the last section we give a brief discussion about the results of this work.
## 2 The model
There are two basic building blocks for constructing open spin chains:
1. The $`R`$ matrix, which is a solution of the Yang-Baxter equation
$`R_{12}(\lambda _1\lambda _2)R_{13}(\lambda _1)R_{23}(\lambda _2)=R_{23}(\lambda _2)R_{13}(\lambda _1)R_{12}(\lambda _1\lambda _2)`$ (2.1)
(see, e.g., ). We assume that the $`R`$ matrix has the unitarity property
$`R_{12}(\lambda )R_{21}(\lambda )=\zeta (\lambda ),`$ (2.2)
where $`R_{21}(\lambda )=𝒫_{12}R_{12}(\lambda )𝒫_{12}=R_{12}(\lambda )^{t_1t_2}`$, $`t`$ denotes transpose, and $`𝒫`$ is the permutation matrix, and also the property
$`R_{12}(\lambda )^{t_1}M_1R_{12}(\lambda 2\rho )^{t_2}M_1^1=\zeta ^{}(\lambda ),`$ (2.3)
with $`M^t=M`$,
$`[M_1M_2,R_{12}(\lambda )]=0.`$ (2.4)
and
$`\zeta (\lambda )=\mathrm{sinh}\mu (\lambda +i)\mathrm{sinh}\mu (\lambda +i),\zeta ^{}(\lambda )=\mathrm{sinh}\mu (\lambda +\rho )\mathrm{sinh}\mu (\lambda +\rho ).`$ (2.5)
For the purposes of this work we are going to need also the $`R`$ matrix that involves different representations of $`U_q(SU(𝒩))`$ , , in particular, $`𝒩`$ and $`\overline{}𝒩`$. This matrix is given by crossing
$`R_{\overline{1}2}(\lambda )=V_1R_{12}(\lambda \rho )^{t_2}V_1=V_2^{t_2}R_{12}(\lambda \rho )^{t_1}V_2^{t_2},`$ (2.6)
where $`V^2=1`$ and $`M=V^tV`$, for the $`A_1^{(1)}`$ case $`R_{\overline{1}2}(\lambda )=R_{12}(\lambda )`$. $`R_{\overline{1}2}(\lambda )`$ also satisfies the unitarity property,
$`R_{\overline{1}2}(\lambda )R_{2\overline{1}}(\lambda )=\zeta ^{}(\lambda ),`$ (2.7)
this equation is equivalent to (2.3), with $`R_{2\overline{1}}(\lambda )=R_{\overline{1}2}(\lambda )^{t_1t_2}`$. Moreover
$`R_{\overline{1}2}(\lambda )^{t_1}M_1R_{\overline{1}2}(\lambda 2\rho )^{t_2}M_1^1=\zeta (\lambda ),`$ (2.8)
which is equivalent to (2.2). $`R_{\overline{1}2}(\lambda )`$ is also a solution of the Yang-Baxter equation
$`R_{\overline{1}2}(\lambda _1\lambda _2)R_{\overline{1}3}(\lambda _1)R_{23}(\lambda _2)=R_{23}(\lambda _2)R_{\overline{1}3}(\lambda _1)R_{\overline{1}2}(\lambda _1\lambda _2).`$ (2.9)
2. The matrices $`K^{}`$, and $`K^+`$ which are solutions of the boundary Yang-Baxter equation
$`R_{12}(\lambda _1\lambda _2)K_1^{}(\lambda _1)R_{21}(\lambda _1+\lambda _2)K_2^{}(\lambda _2)`$
$`=K_2^{}(\lambda _2)R_{12}(\lambda _1+\lambda _2)K_1^{}(\lambda _1)R_{21}(\lambda _1\lambda _2),`$ (2.10)
and,
$`R_{12}(\lambda _1+\lambda _2)K_1^+(\lambda _1)^{t_1}M_1^1R_{21}(\lambda _1\lambda _22\rho )M_1K_2^+(\lambda _2)^{t_2}`$
$`=K_2^+(\lambda _2)^{t_2}M_1R_{12}(\lambda _1\lambda _22\rho )M_1^1K_1^+(\lambda _1)^{t_1}R_{21}(\lambda _1+\lambda _2),`$ (2.11)
there exist an automorphism between $`K^{}`$ and $`K^+`$ i.e.,
$`K^+(\lambda )=MK^{}(\lambda \rho )^t.`$ (2.12)
For the following we are going to need a reflection equation that involves $`R_{\overline{1}2}`$ as well, in particular,
$`R_{\overline{1}2}(\lambda _1\lambda _2)K_{\overline{1}}^{}(\lambda _1)R_{2\overline{1}}(\lambda _1+\lambda _2)K_2^{}(\lambda _2)`$
$`=K_2^{}(\lambda _2)R_{\overline{1}2}(\lambda _1+\lambda _2)K_{\overline{1}}^{}(\lambda _1)R_{2\overline{1}}(\lambda _1\lambda _2),`$ (2.13)
and,
$`R_{\overline{1}2}(\lambda _1+\lambda _2)K_{\overline{1}}^+(\lambda _1)^{t_1}M_1^1R_{2\overline{1}}(\lambda _1\lambda _22\rho )M_1K_2(\lambda _2)^{t_2}`$
$`=K_2^+(\lambda _2)^{t_2}M_1R_{\overline{1}2}(\lambda _1\lambda _22\rho )M_1^1K_{\overline{1}}^+(\lambda _1)^{t_1}R_{2\overline{1}}(\lambda _1+\lambda _2).`$ (2.14)
In the scattering language if we think that the $`K_i`$ matrix describes the scattering of a soliton with the boundary, then $`K_{\overline{i}}`$ describes the scattering of an anti-soliton with the boundary. Subsequently $`R_{\overline{1}2}`$ describes the scattering of a soliton with an anti-soliton.
The corresponding transfer matrix $`t(\lambda )`$ for an open chain of $`N`$ spins is given by ,
$`t(\lambda )=tr_0K_0^+(\lambda )T_0(\lambda )K_0^{}(\lambda )\widehat{T}_0(\lambda ),`$ (2.15)
where $`tr_0`$ denotes trace over the “auxiliary space” 0, $`T_0(\lambda )`$ is the monodromy matrix
$`T_0(\lambda )=R_{0N}(\lambda )\mathrm{}R_{01}(\lambda ),`$ (2.16)
and $`\widehat{T}_0(\lambda )`$ is given by
$`\widehat{T}_0(\lambda )=R_{10}(\lambda )\mathrm{}R_{N0}(\lambda ).`$ (2.17)
(We usually suppress the “quantum-space” subscripts $`1,\mathrm{},N`$.) Indeed, it can be shown that this transfer matrix has the commutativity property
$`[t(\lambda ),t(\lambda ^{})]=0.`$ (2.18)
In this paper, we consider the case of the $`A_{𝒩1}^{(1)}`$ $`R`$ matrix
$`R_{12}(\lambda )_{jj,jj}`$ $`=`$ $`\mathrm{sinh}\mu (\lambda +i),`$
$`R_{12}(\lambda )_{jk,jk}`$ $`=`$ $`\mathrm{sinh}(\mu \lambda ),jk,`$
$`R_{12}(\lambda )_{jk,kj}`$ $`=`$ $`\mathrm{sinh}(i\mu )e^{\mu \lambda sign(jk)},jk,`$ (2.19)
$`1j,k𝒩,`$
which depends on the anisotropy parameter $`\mu 0`$, and which becomes $`SU(𝒩)`$ invariant for $`\mu 0`$. This $`R`$ matrix has the properties (2.2) and (2.3), with
$`M_{jk}=\delta _{jk}e^{i\mu (𝒩2j+1)},\rho =i𝒩/2.`$ (2.20)
The corresponding open spin chain Hamiltonian $``$ for $`K^{}(\lambda )=1`$ and $`K^+(\lambda )=M`$ is:
$`={\displaystyle \underset{n=1}{\overset{N1}{}}}_{nn+1}+{\displaystyle \frac{tr_0M_0_{N0}}{trM}},`$ (2.21)
where the two-site Hamiltonian $`_{jk}`$ is given by
$`_{jk}={\displaystyle \frac{i}{2}}𝒫_{jk}{\displaystyle \frac{d}{d\lambda }}R_{jk}(\lambda )|_{\lambda =0}.`$ (2.22)
One can verify that the Hamiltonian is Hermitian.
We consider the case that $`K^{}(\lambda )=1`$ and $`K^+(\lambda )=M`$, and so the transfer matrix is $`U_q\left(SU(𝒩)\right)`$ invariant , . Following we show that the transfer matrix satisfies a crossing property. To prove the crossing property we need (2.1), (2.3) and the following identity
$`𝒫_{12}^{t_2}M_2R_{21}(\lambda )^{t_1}=R_{21}(\lambda )^{t_1}M_1^1𝒫_{12}^{t_2}.`$ (2.23)
It is important to mention that in order to show (2.23) we considered the “unusual” reflection equation (2.13) for $`\lambda _1\lambda _2=\rho `$. Then the crossing property for the transfer matrix is given by
$`t(\lambda )=\overline{t}(\lambda \rho ),`$ (2.24)
where
$`\overline{t}(\lambda )=tr_0M_0T_{\overline{0}}(\lambda )\widehat{T}_{\overline{0}}(\lambda ),`$ (2.25)
and
$`T_{\overline{0}}(\lambda )`$ $`=`$ $`R_{\overline{0}N}(\lambda )\mathrm{}R_{\overline{0}1}(\lambda ),`$
$`\widehat{T}_{\overline{0}}(\lambda )`$ $`=`$ $`R_{1\overline{0}}(\lambda )\mathrm{}R_{N\overline{0}}(\lambda ),`$ (2.26)
the proof of equation (2.24) follows exactly the proof in . The only difference is that in this case $`\overline{t}(\lambda )`$ is involved as well because of (2.6). The transfer matrix $`\overline{t}(\lambda )`$ satisfies the commutativity property
$`[\overline{t}(\lambda ),\overline{t}(\lambda ^{})]=0.`$ (2.27)
Relation (2.24) is one of the basic results of this paper and it plays an essential role in the derivation of the transfer matrix eigenvalues. The “new” transfer matrix (see also e.g., ) leads apparently to a non local Hamiltonian, however, this is not a problem since $`\overline{t}(\lambda )`$ has an auxiliary character in our calculations as we are going to see later.
## 3 Fusion
The fusion procedure for spin chains with crossing symmetry is known . We generalise this procedure for the case that the $`R`$ matrix does not have crossing symmetry. From now on the indices 1 and 2 refer to the auxiliary space. We consider the equation (2.6) for $`\lambda =\rho `$. Then the $`R_{\overline{1}2}(\lambda )`$ degenerates to a projector onto an one dimensional subspace
$`P_{\overline{1}2}^{}={\displaystyle \frac{1}{𝒩}}V_1𝒫_{12}^{t_2}V_1,`$ (3.1)
also
$`P_{\overline{1}2}^+=1P_{\overline{1}2}^{}`$ (3.2)
is a projector. We consider the Yang-Baxter equation (2.9) for $`\lambda _1\lambda _2=\rho `$, then the fused $`R`$ matrix is given by
$`R_{<\overline{1}2>3}(\lambda )`$ $`=`$ $`P_{\overline{1}2}^+R_{\overline{1}3}(\lambda )R_{23}(\lambda +\rho )P_{\overline{1}2}^+,`$
$`R_{<2\overline{1}>3}(\lambda )`$ $`=`$ $`P_{2\overline{1}}^+R_{23}(\lambda )R_{\overline{1}3}(\lambda +\rho )P_{2\overline{1}}^+,`$ (3.3)
also, one finds
$`R_{3<\overline{1}2>}(\lambda )`$ $`=`$ $`P_{\overline{1}2}^+R_{32}(\lambda \rho )R_{3\overline{1}}(\lambda )P_{\overline{1}2}^+,`$
$`R_{3<2\overline{1}>}(\lambda )`$ $`=`$ $`P_{2\overline{1}}^+R_{3\overline{1}}(\lambda \rho )R_{32}(\lambda )P_{2\overline{1}}^+,`$ (3.4)
similarly, one can fuse the spaces 1 and $`\overline{2}`$. The fused $`R`$ matrices obey the general Yang-Baxter equation
$`R_{j_1j_2}(\lambda _1)R_{j_1j_3}(\lambda _1+\lambda _2)R_{j_2j_3}(\lambda _2)=R_{j_2j_3}(\lambda _2)R_{j_1j_3}(\lambda _1+\lambda _2)R_{j_1j_2}(\lambda _1).`$ (3.5)
Consider the reflection equation (2.13) for $`\lambda _1\lambda _2=\rho `$, then the fused $`K`$ matrices are given by
$`K_{<\overline{1}2>}^{}(\lambda )=P_{\overline{1}2}^+K_{\overline{1}}^{}(\lambda )R_{2\overline{1}}(2\lambda +\rho )K_2^{}(\lambda +\rho )P_{2\overline{1}}^+,`$
$`K_{<\overline{1}2>}^+(\lambda )^{t_{12}}=P_{2\overline{1}}^+K_{\overline{1}}^+(\lambda )^{t_1}M_2R_{2\overline{1}}(2\lambda 3\rho )M_2^1K_2^+(\lambda +\rho )^{t_2}P_{\overline{1}2}^+.`$ (3.6)
The above $`K`$ matrices obey the reflection equations
$`R_{3<\overline{1}2>}(\lambda _1\lambda _2)K_3^{}(\lambda _1)R_{<\overline{1}2>3}(\lambda _1+\lambda _2)K_{<\overline{1}2>}^{}(\lambda _2)`$
$`=K_{<\overline{1}2>}^{}(\lambda _2)R_{3<2\overline{1}>}(\lambda _1+\lambda _2+\rho )K_3^{}(\lambda _1)R_{<2\overline{1}>3}(\lambda _1\lambda _2\rho ),`$ (3.7)
and,
$`R_{<\overline{1}2>3}(\lambda _1+\lambda _2)^{t_{123}}K_3^+(\lambda _1)^{t_3}M_3^1R_{3<\overline{1}2>}(\lambda _1\lambda _22\rho )^{t_{123}}M_3K_{<\overline{1}2>}^+(\lambda _2)^{t_{12}}`$
$`=K_{<\overline{1}2>}^+(\lambda _2)^{t_{12}}M_3R_{<2\overline{1}>3}(\lambda _1\lambda _23\rho )^{t_{123}}M_3^1K_3^+(\lambda _1)^{t_3}R_{3<2\overline{1}>}(\lambda _1+\lambda _2+\rho )^{t_{123}},`$ (3.8)
analogously we obtain the $`K_{<1\overline{2}>}(\lambda )`$ matrices.
Having fused the $`R`$ and $`K`$ matrices we can show that the corresponding fused transfer matrix is (for a detailed computation see e.g., )
$`\widehat{t}(\lambda )=\zeta (2\lambda +2\rho )\overline{t}(\lambda )t(\lambda +\rho )\mathrm{\Delta }[K^+(\lambda )]\delta [T(\lambda )]\mathrm{\Delta }[K^{}(\lambda )]\delta [\widehat{T}(\lambda )],`$ (3.9)
notice that $`\overline{t}(\lambda )`$ appears in the last equation as well as in (2.24). In this case we fuse the spaces $`\overline{1}2`$, and the quantum determinants are , ,
$`\delta [T(\lambda )]=tr_{12}[P_{\overline{1}2}^{}T_{\overline{1}}(\lambda )\widehat{T}_2(\lambda +\rho )],`$
$`\delta [\widehat{T}(\lambda )]=tr_{12}[P_{2\overline{1}}^{}T_{\overline{1}}(\lambda )\widehat{T}_2(\lambda +\rho )],`$
$`\mathrm{\Delta }[K^{}(\lambda )]=tr_{12}[P_{\overline{1}2}^{}K_1^{}(\lambda )R_{2\overline{1}}(2\lambda +\rho )K_2^{}(\lambda +\rho )V_1V_2],`$
$`\mathrm{\Delta }[K^+(\lambda )]=tr_{12}[P_{\overline{1}2}^{}V_1V_2K_2^+(\lambda +\rho )M_2^1R_{\overline{1}2}(2\lambda 3\rho )M_2K_1^+(\lambda )].`$ (3.10)
Similar relations to (3.9) and (3.10) we can derive for the $`1\overline{2}`$ fusion. Using unitarity and the crossing property we prove the identities
$`P_{\overline{1}2}^{}R_{\overline{1}m}(\lambda )R_{2m}(\lambda +\rho )P_{\overline{1}2}^{}`$ $`=`$ $`\zeta (\lambda +\rho )P_{\overline{1}2}^{},`$
$`P_{1\overline{2}}^{}R_{1m}(\lambda )R_{\overline{2}m}(\lambda +\rho )P_{1\overline{2}}^{}`$ $`=`$ $`\zeta ^{}(\lambda +\rho )P_{1\overline{2}}^{},`$
$`m=1,2,\mathrm{},N`$ (3.11)
and by computing the quantum determinants explicitly we find
$`\delta [T(\lambda )]=\delta [\widehat{T}(\lambda )]=\zeta (\lambda +\rho )^Nor\zeta ^{}(\lambda +\rho )^N,`$ (3.12)
depending on the spaces we fuse, i.e., $`\overline{1}2`$ or $`1\overline{2}`$ respectively. For the special case $`K^{}(\lambda )=1`$ and $`K^+(\lambda )=M`$, these are solutions of the reflection equations (2.13) and (2.14) correspondingly, one can show from (3.10) that,
$`\mathrm{\Delta }[K^{}(\lambda )]=g(2\lambda +\rho ),\mathrm{\Delta }[K^+(\lambda )]=g(2\lambda 3\rho ),`$ (3.13)
where,
$`g(\lambda )=\mathrm{sinh}\mu (\lambda +\rho ).`$ (3.14)
## 4 Analytical Bethe Ansatz
We focus here in the simplest case, namely, the $`A_2^{(1)}`$ open chain. The asymptotic behaviour of the $`R`$ matrix for $`\lambda \mathrm{}`$ follows from (2.19)
$`R_{0k}`$ $``$ $`{\displaystyle \frac{1}{2}}e^{\mu \lambda }\left(\begin{array}{ccc}e^{i\mu S_{1,k}}& pJ_{1,k}^{}& pJ_{3,k}^{}\\ 0& e^{i\mu S_{2,k}}& pJ_{2,k}^{}\\ 0& 0& e^{i\mu S_{3,k}}\end{array}\right),`$ (4.4)
$`R_{k0}`$ $``$ $`{\displaystyle \frac{1}{2}}e^{\mu \lambda }\left(\begin{array}{ccc}e^{i\mu S_{1,k}}& 0& 0\\ pJ_{1,k}^+& e^{i\mu S_{2,k}}& 0\\ pJ_{3,k}^+& pJ_{2,k}^+& e^{i\mu S_{3,k}}\end{array}\right),`$ (4.8)
where $`p=2\mathrm{sinh}(i\mu )`$, and the matrix elements are:
$`S_i`$ $`=`$ $`e_{i,i},i=1,2,3,`$
$`J_i^+`$ $`=`$ $`e_{i,i+1},J_i^{}=e_{i+1,i},i=1,2,`$
$`J_3^+`$ $`=`$ $`e_{1,3},J_3^{}=e_{3,1},`$ (4.9)
with,
$`(e_{i,j})_{kl}=\delta _{ik}\delta _{jl}.`$ (4.10)
The leading asymptotic behaviour of the monodromy matrix is given by
$`T_0(\lambda )`$ $``$ $`({\displaystyle \frac{1}{2}})^Ne^{\mu \lambda N}\left(\begin{array}{ccc}e^{i\mu 𝒮_1}& p𝒥_1^{}& p𝒥_3^{}\\ 0& e^{i\mu 𝒮_2}& p𝒥_2^{}\\ 0& 0& e^{i\mu 𝒮_3}\end{array}\right),`$ (4.14)
$`\widehat{T}_0(\lambda )`$ $``$ $`({\displaystyle \frac{1}{2}})^Ne^{\mu \lambda N}\left(\begin{array}{ccc}e^{i\mu 𝒮_1}& 0& 0\\ p𝒥_1^+& e^{i\mu 𝒮_2}& 0\\ p𝒥_3^+& p𝒥_2^+& e^{i\mu 𝒮_3}\end{array}\right),`$ (4.18)
where,
$`𝒮_i`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N}{}}}S_{i,k}i=1,2,3,`$
$`𝒥_1^\pm `$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N}{}}}e^{i\mu S_{1,N}}\mathrm{}e^{i\mu S_{1,k+1}}J_{1,k}^\pm e^{i\mu S_{2,k1}}\mathrm{}e^{i\mu S_{2,1}},`$
$`𝒥_2^\pm `$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N}{}}}e^{i\mu S_{2,N}}\mathrm{}e^{i\mu S_{2,k+1}}J_{2,k}^\pm e^{i\mu S_{3,k1}}\mathrm{}e^{i\mu S_{3,1}},`$
$`𝒥_3^\pm `$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N}{}}}(e^{i\mu S_{1,N}}\mathrm{}e^{i\mu S_{1,k+1}}J_{3,k}^\pm +p𝒥_{1,k+1}^\pm J_{2,k}^\pm )e^{i\mu S_{3,k1}}\mathrm{}e^{i\mu S_{3,1}}.`$ (4.19)
Then from equations (2.15) for $`K^{}(\lambda )=1`$ and $`K^+(\lambda )=M`$, (2.20), and (4.18) we conclude that the leading asymptotic behaviour of the transfer matrix has the following form
$`t(\lambda )`$ $``$ $`({\displaystyle \frac{1}{2}})^{2N}e^{2\mu \lambda N}(e^{2i\mu 2i\mu 𝒮_1}+e^{2i\mu }p^2𝒥_1^{}𝒥_1^++e^{2i\mu }p^2𝒥_3^{}𝒥_3^+`$ (4.20)
$`+`$ $`e^{2i\mu 2i\mu 𝒮_3}+p^2𝒥_2^{}𝒥_2^++e^{2i\mu 𝒮_2}).`$
We introduce the operators $`M_1`$ and $`M_2`$
$`𝒮_1=NM_1,𝒮_2=M_1M_2,𝒮_3=M_2,`$ (4.21)
we consider simultaneous eigenstates of $`M_i`$ and the transfer matrix i.e.,
$`M_i|\mathrm{\Lambda }^{(m)}=m_i|\mathrm{\Lambda }^{(m)},t(\lambda )|\mathrm{\Lambda }^{(m)}=\mathrm{\Lambda }^{(m)}(\lambda )|\mathrm{\Lambda }^{(m)},`$ (4.22)
(to simplify our notation we write $`(m)`$ instead of $`(m_1m_2)`$). We choose these states to be annihilated by $`𝒥_i^+`$
$`𝒥_i^+|\mathrm{\Lambda }^{(m)}=0.`$ (4.23)
We conclude, from (4.20), (4.22), and (4.23) that the asymptotic behaviour of the corresponding eigenvalue is given by
$`\mathrm{\Lambda }^{(m)}(\lambda )({\displaystyle \frac{1}{2}})^{2N}e^{2\mu \lambda N}(e^{2i\mu (1+Nm_1)}+e^{2i\mu (1m_2)}+e^{2i\mu (m_1m_2)}).`$ (4.24)
In order to determine the asymptotic behaviour of $`\overline{t}(\lambda )`$, we also need the asymptotic behaviour of $`R_{\overline{0}k}(\lambda )(R_{k\overline{0}}(\lambda ))`$ for $`\lambda \mathrm{}`$
$`R_{\overline{0}k}`$ $``$ $`{\displaystyle \frac{1}{2}}e^{\mu \lambda \frac{3i\mu }{2}}\left(\begin{array}{ccc}e^{i\mu S_{3,k}}& qpJ_{2,k}^{}& q^2pJ_{3,k}^{}\\ 0& e^{i\mu S_{2,k}}& qpJ_{1,k}^{}\\ 0& 0& e^{i\mu S_{1,k}}\end{array}\right),`$ (4.28)
$`R_{k\overline{0}}`$ $``$ $`{\displaystyle \frac{1}{2}}e^{\mu \lambda \frac{3i\mu }{2}}\left(\begin{array}{ccc}e^{i\mu S_{3,k}}& 0& 0\\ qpJ_{2,k}^+& e^{i\mu S_{2,k}}& 0\\ q^2pJ_{3,k}^+& qpJ_{1,k}^+& e^{i\mu S_{1,k}}\end{array}\right),`$ (4.32)
we define $`R_{\overline{0}k}(\lambda )`$ from (2.6) using
$`V=\left(\begin{array}{ccc}& & q\\ & 1& \\ q^1& & \end{array}\right)`$ (4.36)
where $`q=e^{i\mu }`$, and the asymptotic behaviour of $`\overline{t}(\lambda )`$ is given by
$`\overline{t}(\lambda )`$ $``$ $`({\displaystyle \frac{1}{2}})^{2N}e^{2\mu \lambda N3i\mu N}(e^{2i\mu +2i\mu 𝒮_1}+e^{2i\mu }p^2\overline{𝒥}_1^{}\overline{𝒥}_1^++e^{2i\mu }p^2\overline{𝒥}_3^{}\overline{𝒥}_3^+`$ (4.37)
$`+`$ $`e^{2i\mu +2i\mu 𝒮_3}+p^2\overline{𝒥}_2^{}\overline{𝒥}_2^++e^{2i\mu 𝒮_2}),`$
with
$`\overline{𝒥}_1^\pm `$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N}{}}}e^{i\mu S_{2,N}}\mathrm{}e^{i\mu S_{2,k+1}}J_{1,k}^\pm e^{i\mu S_{1,k1}}\mathrm{}e^{i\mu S_{1,1}},`$
$`\overline{𝒥}_2^\pm `$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N}{}}}e^{i\mu S_{3,N}}\mathrm{}e^{i\mu S_{3,k+1}}J_{2,k}^\pm e^{i\mu S_{2,k1}}\mathrm{}e^{i\mu S_{2,1}},`$
$`\overline{𝒥}_3^\pm `$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N}{}}}(e^{i\mu S_{3,N}}\mathrm{}e^{i\mu S_{3,k+1}}J_{3,k}^\pm p\overline{𝒥}_{2,k+1}^\pm J_{1,k}^\pm )e^{i\mu S_{1,k1}}\mathrm{}e^{i\mu S_{1,1}}.`$ (4.38)
The $`|\mathrm{\Lambda }^{(m)}`$ states are also annihilated by $`\overline{𝒥}_i^+`$ i.e.,
$`\overline{𝒥}_i^+|\mathrm{\Lambda }^{(m)}=0,`$ (4.39)
where the corresponding eigenvalue of $`\overline{t}(\lambda )`$ is
$`\overline{\mathrm{\Lambda }}^{(m)}(\lambda )({\displaystyle \frac{1}{2}})^{2N}e^{2\lambda \mu N3i\mu N}(e^{2i\mu (1+Nm_1)}+e^{2i\mu (1m_2)}+e^{2i\mu (m_1m_2)}).`$ (4.40)
We consider the state with all “spins” up i.e.,
$`|\mathrm{\Lambda }^{(0)}={\displaystyle \underset{k=1}{\overset{N}{}}}|+_{(k)},`$ (4.41)
this is annihilated by $`𝒥_i^+`$ and $`\overline{𝒥}_i^+`$, where (we suppress the $`(k)`$ index)
$`|+=\left(\begin{array}{c}1\\ 0\\ 0\end{array}\right).`$ (4.45)
We assume this is an eigenstate of the transfer matrix and it is also an eigenstate of $`\overline{t}(\lambda )`$. The action of the $`R`$ matrix on the $`|+`$ state gives lower and upper triangular matrices i.e.,
$`+|R_{0k}(\lambda )=+|\left(\begin{array}{ccc}A_k& 0& 0\\ C_{1,k}& D_{1,k}& 0\\ C_{2,k}& D_{3,k}& D_{4,k}\end{array}\right),R_{k0}(\lambda )|+=\left(\begin{array}{ccc}A_k& B_{1,k}& B_{2,k}\\ 0& D_{1,k}& D_{2,k}\\ 0& 0& D_{4,k}\end{array}\right)|+,`$ (4.52)
where the matrices $`A`$, $`B_i`$, $`C_i`$, $`D_i`$ act on the quantum space and they are determined by the form of the $`R`$ matrix (2.19). Then the action of the transfer matrix on the pseudo-vacuum is,
$`\mathrm{\Lambda }^{(0)}|T_0(\lambda )=\mathrm{\Lambda }^{(0)}|\left(\begin{array}{ccc}𝒜& 0& 0\\ 𝒞_1& 𝒟_1& 0\\ 𝒞_2& 𝒟_3& 𝒟_4\end{array}\right),\widehat{T}_0(\lambda )|\mathrm{\Lambda }^{(0)}=\left(\begin{array}{ccc}𝒜& _1& _2\\ 0& 𝒟_1& 𝒟_2\\ 0& 0& 𝒟_4\end{array}\right)|\mathrm{\Lambda }^{(0)},`$ (4.59)
and the matrix elements are given by
$`𝒜=A_N\mathrm{}A_1,𝒟_1=D_{1,N}\mathrm{}D_{1,1},𝒟_4=D_{4,N}\mathrm{}D_{4,1},`$ (4.60)
$`𝒞_1`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N}{}}}D_{1,N}\mathrm{}D_{1,k+1}C_{1,k}A_{k1}\mathrm{}A_1,𝒞_0=0,`$
$`𝒟_3`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N}{}}}D_{4,N}\mathrm{}D_{4,k+1}D_{3,k}D_{1,k1}\mathrm{}D_{1,1},𝒟_{3,0}=0,`$
$`𝒞_2`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N}{}}}D_{4,N}\mathrm{}D_{4,k+1}(C_{2,k}A_{k1}\mathrm{}A_1+D_{3,k}𝒞_{1,k1}).`$ (4.61)
Similarly, to find $`_1`$, $`_2`$ and $`𝒟_2`$ we replace $`C_1`$, $`D_3`$, $`C_2`$ and $`𝒞_1`$ with $`B_1`$, $`D_2`$, $`B_2`$ and $`_1`$ correspondingly. The transfer matrix eigenvalue for the pseudo-vacuum state is
$`\mathrm{\Lambda }^{(0)}(\lambda )`$ $`=`$ $`\mathrm{\Lambda }^{(0)}|(e^{2i\mu }𝒜^2+𝒟_1^2+e^{2i\mu }𝒟_4^2+𝒞_1_1+e^{2i\mu }𝒞_2_2+e^{2i\mu }𝒟_3𝒟_2)|\mathrm{\Lambda }^{(0)},`$ (4.62)
and after some tedious algebra we find
$`\mathrm{\Lambda }^{(0)}(\lambda )=f(\lambda )(a(\lambda )^{2N}{\displaystyle \frac{\mathrm{sinh}2\mu (\lambda +i)}{\mathrm{sinh}(2\mu \lambda )}}+b(\lambda )^{2N}(1+{\displaystyle \frac{\mathrm{sinh}\mu (2\lambda +i)}{\mathrm{sinh}\mu (2\lambda +3i)}})).`$ (4.63)
We make the assumption that a general eigenvalue has the form of a “dressed” pseudo-vacuum eigenvalue i.e.,
$`\mathrm{\Lambda }^{(m)}(\lambda )=f(\lambda )(a(\lambda )^{2N}{\displaystyle \frac{\mathrm{sinh}2\mu (\lambda +i)}{\mathrm{sinh}(2\mu \lambda )}}A_1(\lambda )+b(\lambda )^{2N}(A_2(\lambda )+{\displaystyle \frac{\mathrm{sinh}\mu (2\lambda +i)}{\mathrm{sinh}\mu (2\lambda +3i)}}A_3(\lambda ))).`$ (4.64)
Again, the action of $`R_{\overline{0}k}(\lambda )`$ on the $`|+`$ state gives upper and lower triangular matrices (see (4.52)), so we find an analogous equation for the $`\overline{\mathrm{\Lambda }}^{(m)}(\lambda )`$,
$`\overline{\mathrm{\Lambda }}^{(m)}(\lambda )=f(\lambda )(\overline{a}(\lambda )^{2N}{\displaystyle \frac{\mathrm{sinh}\mu (2\lambda +i)}{\mathrm{sinh}\mu (2\lambda +3i)}}\overline{A}_1(\lambda )+\overline{b}(\lambda )^{2N}(\overline{A}_2(\lambda )+{\displaystyle \frac{\mathrm{sinh}2\mu (\lambda +i)}{\mathrm{sinh}(2\mu \lambda )}}\overline{A}_3(\lambda ))),`$ (4.65)
(we suppress the $`(m)`$ index from the “dressing” functions), where
$`f(\lambda )={\displaystyle \frac{\mathrm{sinh}(2\mu \lambda )\mathrm{sinh}\mu (2\lambda +3i)}{\mathrm{sinh}2\mu (\lambda +i)\mathrm{sinh}\mu (2\lambda +i)}},`$ (4.66)
$`a(\lambda )`$ $`=`$ $`\mathrm{sinh}\mu (\lambda +i),b(\lambda )=\mathrm{sinh}(\mu \lambda ),`$ (4.67)
and $`\overline{a}(\lambda )`$, $`\overline{b}(\lambda )`$ are $`a(\lambda \rho )`$, $`b(\lambda \rho )`$, respectively. It is obvious that $`\mathrm{\Lambda }^{(0)}(\lambda )=\overline{\mathrm{\Lambda }}^{(0)}(\lambda \rho )`$. We conclude from the asymptotic behaviour of the transfer matrix that
$`A_1(\lambda )e^{2i\mu m_1},A_2(\lambda )e^{2i\mu (m_2m_1)},A_3(\lambda )e^{2i\mu m_2},`$ (4.68)
and
$`\overline{A}_1(\lambda )e^{2i\mu m_1},\overline{A}_2(\lambda )e^{2i\mu (m_2m_1)},\overline{A}_3(\lambda )e^{2i\mu m_2}.`$ (4.69)
We substitute the eigenvalues to the fusion equations (3.9) and we obtain conditions involving $`A_1(\lambda )`$, $`A_3(\lambda )`$, and $`\overline{A}_1(\lambda )`$, $`\overline{A}_3(\lambda )`$. It is clear that the $`\mathrm{\Lambda }^{(0)}(\lambda )`$ satisfies (3.9) a fact that further supports our assumption that $`|\mathrm{\Lambda }^{(0)}`$ is an eigenstate of the transfer matrices, $`t(\lambda )`$ and $`\overline{t}(\lambda )`$. The fusion equations give
$`A_1(\lambda +\rho )\overline{A}_1(\lambda )=1,\overline{A}_3(\lambda +\rho )A_3(\lambda )=1,`$ (4.70)
where notice that we obtain two conditions from (3.9), one from $`\overline{1}2`$ fusion and one from $`1\overline{2}`$, whereas e.g., for the $`A_1^{(1)}`$ case we obtain only one such condition. From the crossing property of the transfer matrix (2.24),
$`A_i(\lambda \rho )=\overline{A}_i(\lambda ),i=1,2,3.`$ (4.71)
Combining the last two conditions we find
$`A_1(\lambda )A_1(\lambda )=1,\overline{A}_3(\lambda )\overline{A}_3(\lambda )=1.`$ (4.72)
Note that the previous equations mix the “dressing” functions of $`\mathrm{\Lambda }^{(m)}(\lambda )`$ and $`\overline{\mathrm{\Lambda }}^{(m)}(\lambda )`$, which is expected because of the form of equations (2.24) and (3.9). In the case of a model with crossing symmetry e.g., $`A_1^{(1)}`$, the two eigenvalues become degenerate. From the periodicity of the transfer matrix we obtain
$`t(\lambda +{\displaystyle \frac{i\pi }{\mu }})=t(\lambda ),`$ (4.73)
we expect the eigenvalues to be periodic as well. We impose $`A_2(\lambda )`$ ($`\overline{A}_2(\lambda )`$) to have the same poles with $`A_1(\lambda )`$ and $`A_3(\lambda )`$ ($`\overline{A}_1(\lambda )`$ and $`\overline{A}_3(\lambda )`$). Also, the residue of $`\mathrm{\Lambda }^{(m)}(\lambda )`$ at $`\lambda =i`$ should vanish, thus we obtain the following condition,
$`A_2(i)=A_3(i).`$ (4.74)
We put all the above requirements together and we find that
$`A_1(\lambda )={\displaystyle \underset{j=1}{\overset{m_1}{}}}{\displaystyle \frac{\mathrm{sinh}\mu (\lambda +\lambda _j^{(1)})}{\mathrm{sinh}\mu (\lambda +\lambda _j^{(1)}+i)}}{\displaystyle \frac{\mathrm{sinh}\mu (\lambda \lambda _j^{(1)}i)}{\mathrm{sinh}\mu (\lambda \lambda _j^{(1)})}},`$ (4.75)
$`A_2(\lambda )`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{m_1}{}}}{\displaystyle \frac{\mathrm{sinh}\mu (\lambda +\lambda _j^{(1)}+2i)}{\mathrm{sinh}\mu (\lambda +\lambda _j^{(1)}+i)}}{\displaystyle \frac{\mathrm{sinh}\mu (\lambda \lambda _j^{(1)}+i)}{\mathrm{sinh}\mu (\lambda \lambda _j^{(1)})}}`$ (4.76)
$`{\displaystyle \underset{j=1}{\overset{m_2}{}}}{\displaystyle \frac{\mathrm{sinh}\mu (\lambda +\lambda _j^{(2)}+i)}{\mathrm{sinh}\mu (\lambda +\lambda _j^{(2)}+2i)}}{\displaystyle \frac{\mathrm{sinh}\mu (\lambda \lambda _j^{(2)}i)}{\mathrm{sinh}\mu (\lambda \lambda _j^{(2)})}},`$
$`A_3(\lambda )={\displaystyle \underset{j=1}{\overset{m_2}{}}}{\displaystyle \frac{\mathrm{sinh}\mu (\lambda +\lambda _j^{(2)}+3i)}{\mathrm{sinh}\mu (\lambda +\lambda _j^{(2)}+2i)}}{\displaystyle \frac{\mathrm{sinh}\mu (\lambda \lambda _j^{(2)}+i)}{\mathrm{sinh}\mu (\lambda \lambda _j^{(2)})}}.`$ (4.77)
We obtain $`\overline{A}_i(\lambda )`$ from (4.71). It is easy to check that the eigenvalues satisfy all the above conditions. Moreover we want the eigenvalues to be analytical, so the poles must vanish. This condition leads to the Bethe ansatz equations
$`e_1(\lambda _i^{(1)})^{2N}`$ $`=`$ $`{\displaystyle \underset{ij=1}{\overset{m_1}{}}}e_2(\lambda _i^{(1)}\lambda _j^{(1)})e_2(\lambda _i^{(1)}+\lambda _j^{(1)}){\displaystyle \underset{j=1}{\overset{m_2}{}}}e_1(\lambda _i^{(1)}\lambda _j^{(2)})e_1(\lambda _i^{(1)}+\lambda _j^{(2)})`$
$`1`$ $`=`$ $`{\displaystyle \underset{ij=1}{\overset{m_2}{}}}e_2(\lambda _i^{(2)}\lambda _j^{(2)})e_2(\lambda _i^{(2)}+\lambda _j^{(2)}){\displaystyle \underset{j=1}{\overset{m_1}{}}}e_1(\lambda _i^{(2)}\lambda _j^{(1)})e_1(\lambda _i^{(2)}+\lambda _j^{(1)})`$ (4.78)
where we have defined $`e_n(\lambda )`$ as
$`e_n(\lambda )={\displaystyle \frac{\mathrm{sinh}\mu (\lambda +\frac{in}{2})}{\mathrm{sinh}\mu (\lambda \frac{in}{2})}}.`$ (4.79)
The exact computation for the general case becomes complicated, however, one can “guess” the form of the general eigenvalue having in mind all the conditions that it must satisfy. The expression for the spectrum of the transfer matrix for any $`𝒩`$ is given by
$`\mathrm{\Lambda }^{(m)}(\lambda )b(\lambda )^{2N}{\displaystyle \underset{k=1}{\overset{𝒩}{}}}{\displaystyle \frac{\mathrm{sinh}(2\mu \lambda )}{\mathrm{sinh}\mu (2\lambda +(k1)i)}}{\displaystyle \frac{\mathrm{sinh}\mu (2\lambda +i)}{\mathrm{sinh}\mu (2\lambda +ki)}}A_k(\lambda ),`$ (4.80)
where
$`A_k(\lambda )`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{m_{k1}}{}}}{\displaystyle \frac{\mathrm{sinh}\mu (\lambda +\lambda _j^{(k1)}+ki)}{\mathrm{sinh}\mu (\lambda +\lambda _j^{(k1)}+(k1)i)}}{\displaystyle \frac{\mathrm{sinh}\mu (\lambda \lambda _j^{(k1)}+i)}{\mathrm{sinh}\mu (\lambda \lambda _j^{(k1)})}}`$ (4.81)
$`{\displaystyle \underset{j=1}{\overset{m_k}{}}}{\displaystyle \frac{\mathrm{sinh}\mu (\lambda +\lambda _j^{(k)}+(k1)i)}{\mathrm{sinh}\mu (\lambda +\lambda _j^{(k)}+ki)}}{\displaystyle \frac{\mathrm{sinh}\mu (\lambda \lambda _j^{(k)}i)}{\mathrm{sinh}\mu (\lambda \lambda _j^{(k)})}},k=1,\mathrm{},𝒩,`$
$`m_0=N`$ and $`m_𝒩=0`$. Also, $`\overline{\mathrm{\Lambda }}^{(m)}(\lambda )=\mathrm{\Lambda }^{(m)}(\lambda \rho )`$. The procedure we described uniquely fixes the “dressing” functions. By inspection we can verify that the above eigenvalues indeed satisfy all the requirements we derived previously i.e, analyticity, asymptotic behaviour, crossing, etc. e.g.,
$`\overline{A}_k(\lambda )=A_k(\lambda \rho ),k=1,\mathrm{},𝒩,`$ (4.82)
$`A_1(\lambda )A_1(\lambda )=1,\overline{A}_𝒩(\lambda )\overline{A}_𝒩(\lambda )=1.`$ (4.83)
We can see from (4.81) that every two terms have the same poles. From the analyticity of the eigenvalues we obtain the Bethe ansatz equations
$`1`$ $`=`$ $`{\displaystyle \underset{ij=1}{\overset{m_k}{}}}e_2(\lambda _i^{(k)}\lambda _j^{(k)})e_2(\lambda _i^{(k)}+\lambda _j^{(k)})`$ (4.84)
$`{\displaystyle \underset{j=1}{\overset{m_{k+1}}{}}}e_1(\lambda _i^{(k)}\lambda _j^{(k+1)})e_1(\lambda _i^{(k)}+\lambda _j^{(k+1)})`$
$`{\displaystyle \underset{j=1}{\overset{m_{k1}}{}}}e_1(\lambda _i^{(k)}\lambda _j^{(k1)})e_1(\lambda _i^{(k)}+\lambda _j^{(k1)}),`$
$`k=1,\mathrm{},𝒩,`$
for $`𝒩=3`$ we recover (4.78). The results, as expected, coincide with the known ones obtained by nesting , ( for $`\xi \pm i\mathrm{}`$).
## 5 Discussion
We generalised the fusion procedure for open spin chains without crossing symmetry. Furthermore, we showed that even though the $`R`$ matrix does not have crossing symmetry the transfer matrix satisfies a crossing property (2.24). We applied these results to diagonalise the transfer matrix via the analytical Bethe ansatz method. We found explicit expressions for the transfer matrix spectrum (4.80) and we deduced the Bethe ansatz equations (4.84), avoiding nesting. The main realization in this paper was the necessity of the transfer matrix $`\overline{t}(\lambda )`$ (2.25) in the derivation of the analytical Bethe ansatz. Indeed, it was necessary to consider $`\overline{t}(\lambda )`$ together with the usual transfer matrix in order to derive the conditions that the eigenvalues should satisfy.
Here, we considered the special case where the chain has a $`U_q\left(SU(𝒩)\right)`$ symmetry. However, we believe that the previous analysis can be extended even in the case of the reduced symmetry $`U_q\left(SU(l)\right)\times U_q\left(SU(𝒩l)\right)\times U(1)`$ . Moreover, the Bethe ansatz equations are known for open spin chains with “soliton preserving” boundary conditions. There is also the case of “soliton non-preserving” boundary conditions (see e.g., , ) for which the Bethe ansatz equations are not known. Using the analytical Bethe ansatz, one can presumably derive the corresponding transfer matrix spectrum and the Bethe ansatz equations, avoiding nesting. We hope to address these questions in a future work .
## 6 Acknowledgements
I am grateful to P. Bowcock, E. Corrigan and G.W. Delius for helpful discussions. I would also like to thank R.I. Nepomechie for valuable suggestions and for prior collaborations. This work was supported by the European Commission under the TMR Network “Integrability, non-perturbative effects, and symmetry in quantum field theory”, contract number FMRX-CT96-0012.
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# Influence of vortex–vortex interaction on critical currents across low–angle grain boundaries in YBa2Cu3O7-δ thin films
## Abstract
Low–angle grain boundaries with misorientation angles $`\theta <5^{}`$ in optimally doped thin films of YBCO are investigated by magnetooptical imaging. By using a numerical inversion scheme of Biot–Savart’s law the critical current density across the grain boundary can be determined with a spatial resolution of about 5 $`\mu `$m. Detailed investigation of the spatially resolved flux density and current density data shows that the current density across the boundary varies with varying local flux density. Combining the corresponding flux and current pattern it is found that there exists a universal dependency of the grain boundary current on the local flux density. A change in the local flux density means a variation in the flux line–flux line distance. With this knowledge a model is developped that explains the flux–current relation by means of magnetic vortex–vortex interaction.
The current limiting effect of grain boundaries in high–temperature superconductors (HTSC’s) is a topic of primary importance for the application of these materials. In the last 12 years many measurements have been carried out which show that the transport current across a grain boundary exhibits an exponential decay with increasing misorientation angle $`\varphi `$ . The reason for this exponential decay is a topic of ongoing discussions and several mechanisms to explain this behavior have been taken into account. At the grain boundary the local distortion of the crystal symmetry causes an array of dislocation cores in the superconducting film. The strain field of these dislocations creates regularly ordered normal conducting regions. This leads to a reduction of the effective superconducting interface at the grain boundary and additionally to a reduction of the order parameter in the superconducting regions . This reduction can be explained by a local bending of the electronic band structure . A further point of interest is the influence of oxygen deficiency or oxygen disorder which can lead to the appearance of localized states at the grain boundary . Considering the pinning scenario of the flux lines located inside the grain boundary it is necessary to remark that the vortices are anisotropic which leads to an enhanced coherence length in direction of the boundary and therefore to a reduction of the pinning force density. Most of this effects which lead to the observed exponential decay, however, play just a secondary role if one considers grain boundaries with low misorientation angles in thin films. It can be experimentally found that the exponential decrease of the current density starts above a certain threshold angle of about $`\varphi _0`$ = 5 in case of zero field . The regime of low–angle grain boundaries (LAGB’s) which means in this case grain boundaries with misorientation angles $`\varphi <5^{}`$ can no longer be described by a weak link behavior because the transport current densities which can occur across these grain boundaries can reach the values of the unperturbed film.
In this paper we present local measurements of the critical current density across LAGB’s in thin films performed by a magnetooptical technique. It can be shown that it is not feasible to characterize these grain boundaries only by a global transport current. The critical currents, however, are very sensitive to the local magnetic flux density conditions inside the grain boundary. By applying an appropriate external magnetic field it can be managed that the current limiting role of the grain boundary vanishes. This effect was also found in a similar form in YBCO bulk grain boundaries .
To investigate the role of LAGB’s on the critical current density in superconducting films the following sample geometry is used. YBCO thin films are grown epitaxially on SrTiO<sub>3</sub> (STO) bicrystalline substrates by pulsed laser deposition. The STO–substrates contain a symmetric tilt grain boundary with misorientation angles $`\varphi <5^{}`$. A sketch of the substrate geometry is shown in Fig. 1.
The geometry of this grain boundary is adapted from the superconducting film during the growth process. We want to focus in the following on two optimally doped YBCO films with the dimensions of 1 mm $`\times `$ 1 mm $`\times `$ 300 nm patterned by chemical etching, which contain a 2 and a 3 grain boundary, respectively.
The measurements that are presented in this paper are performed by applying a magnetooptical technique. As a field sensing layer a ferrimagnetic Lutetium doped Iron garnet film is used which is grown on a Gallium–Gadolinium garnet substrate by liquid–phase epitaxy. This field sensing layer allows the depiction of the magnetic flux density distribution with a spatial resolution of 3 – 5 $`\mu `$m . The garnet film is observed by a polarization light microscope and the images are obtained by a charge–coupled device camera with a resolution of 1000 $`\times `$ 1000 picture elements. Due to the fixed magnetic anisotropy of the indicator film a flux density range of about 2 to 150 mT can be observed with high quantitative precision.
In a first measurement a sample with a 3 grain boundary is examined by use of the magnetooptical technique. Fig. 2 shows a grayscale image of the sample after zero–field cooling (ZFC) to 5 K with an afterwards applied field of $`B_{ex}`$ = 48 mT. Bright parts refer to high magnetic flux densities, black indicates flux–free regions. The image shows $`B_z`$, the flux density component perpendicular to the film.
The gray square represents the region of the superconducting film. Magnetic flux has begun to penetrate the sample in a well–known cushion–like form along the sample’s borders . The influence of the grain boundary can be seen in the two horizontal bright lines in the upper half of the sample. These lines indicate a large penetration of the external flux along the LAGB . This is what can be expected and can be easily understood by a reduced critical current density across the grain boundary, which leads to an enlarged flux penetration in this region. Starting out from the greyscale image in Fig. 2, it can be pointed out that the penetration depth at the grain boundary is about twice as large as in the unperturbed film. That means in a first order approximation a reduced critical current density by a factor of 2 across the LAGB.
In a next step the sample is driven into the fully penetrated state by applying an external magnetic flux density of $`B_{ex}`$ 500 mT. Afterwards the applied external field is reduced gradually. Fig. 3 shows a series of snapshots at B<sub>ex</sub> = 112 mT, 96 mT, 80 mT and 48 mT.
The behavior of the grain boundary is no longer as trivial to understand as for the ZFC case shown in Fig. 2. For $`B_{ex}`$ =112 mT, which is shown in the top left image of Fig. 3 the square–shaped film shows a perfect four–fold symmetry of the flux density distribution. The white discontinuity lines which indicate the trapped flux inside the sample are exactly crossed, no perturbation by the LAGB can be detected. This means that there exists no current limiting effect across the grain boundary; grain boundary and unperturbed film exhibit the same critical current. With decreasing external magnetic field, the grain boundary reappears continuously, the influence of the grain boundary is clearly visible again at $`B_{ex}`$ = 48 mT. The black lines along the LAGB indicate an expulsion of the trapped magnetic flux due to the collapsing critical currents across the grain boundary. These images prove that the critical current density across a LAGB shows a strong dependence of the magnetic flux density. A similar behavior was found for twin boundaries in YBCO single crystals . The current density increases with increasing flux density up to the value in the unperturbed film, that means that the current limiting effect of an LAGB can be compensated by applying an appropriate magnetic field.
To obtain further information about the current limiting role of LAGB’s in thin films it is necessary to determine the critical currents across the grain boundary quantitatively. This is possible by a detailed examination of the magnetooptical data. From the measured perpendicular component of the magnetic flux density $`B_z`$ the corresponding current density distribution can be calculated by a numerical inversion of Biot–Savart’s law. The relation
$`B_z(x,y)=`$ (1)
$`\mu _0H_{ex}+\mu _0{\displaystyle _V}{\displaystyle \frac{j_x(𝐫^{})(yy^{})j_y(𝐫^{})(xx^{})}{4\pi |𝐫𝐫^{}|^3}}d^3r^{}`$ (2)
which is valid for a two–dimensional current density distribution $`𝐣=(j_x,j_y,0)`$ can be inverted unambigiously by using Fourier transformation and convolution theorem . The lateral resolution of the calculated current density distribution is about 5–7 $`\mu m`$ and is therefore slightly reduced compared to that of the magnetic field data. This reduction appears because of noise effects in the measurement .
Two different representations of the calculated currents are shown in Fig. 4.
The left image in Fig. 4 shows an overlay of the flux density distribution known from Fig. 2 and the from these data calculated corresponding current stream lines as solid black lines. The lines which appear outside the sample’s region are generated by numerical artefacts in the calculation. The influence of these virtual currents on the current pattern in the superconducting film is very small and thus can be neglected. An important feature of the current density distribution is the strong bending of the stream lines in the region of the grain boundary that can easily be identified by the two bright lines in the upper part of the image. In the right image the absolute value of the current density is plotted as a grayscale. The white color indicates a current density of about $`2.5\times 10^{11}`$ A/m<sup>2</sup>. This representation also shows clearly the perturbing influence of the grain boundary in the upper half of the sample. The small white spots in the center of the sample are artefacts of the numerical calculus. A profile of the absolute values of $`j`$ which is taken along the horizontal black line is plotted below. It shows clearly the critical current in the flux–penetrated regions and the screening currents in the center of the sample.
The measured flux density distribution and the corresponding calculated current density distribution are now used to investigate the local relation between field and current. This investigation should clarify the remarkable behavior of the LAGB in the decreasing field shown in Fig. 3. To obtain the local relation between flux and current density we take the spatially resolved data from Fig. 2 and Fig. 4, respectively, and note down the values of flux density and current density for every single picture element. That means for everyone of the about 1000 $`\times `$ 1000 picture elements we get a couple $`(B,j)`$ that can be plotted in a $`B`$$`j`$ diagram as shown in Fig. 5.
This technique was now applied on two different regions of the superconductor. First, of course, on the area of the grain boundary and second, for comparison, on an area of the unperturbed film . As a result we obtain the two curves plotted in Fig. 6.
The large difference between the two curves is obvious. The upper curve depicts the field dependence of the critical current density in the unperturbed area. It shows a constant value of $`j_c=2.5\times 10^{11}`$ A/m<sup>2</sup> over the considered flux density range from 20 to 60 mT. A totally different behavior occurs for the $`B`$$`j`$ relation in case of the currents across the grain boundary. A strong increase with increasing flux density can be detected. The experimental data is shown in the lower curve in Fig. 6. The curve has a nearly linear slope with a slight bending at $`B`$ = 40 mT. This bending is due to local variations in the microstructure of the film and will not be discussed any further. The increasing current density across the grain boundary does not reach the value of the unperturbed film in the considered flux density range, but meets the other curve at $`B80`$ mT. This behavior goes along with the non–perturbing influence of the LAGB in the first image of Fig. 3. Note, that the flux density values of Fig. 3 are valid for the applied external flux density whereas in this case $`B`$ is the local magnetic flux density.
Fig. 7 shows the $`B`$$`j`$ relation for two grain boundaries with different misorientation angles. The upper curve shows the experimental data for a symmetric 2 grain boundary, the lower curve is again the curve from Fig. 6 for comparison.
Both of the curves show nearly the same slope, they are just seperated by an offset of about 4 $`\times 10^{10}`$ A/m<sup>2</sup>, if the small hump of the 3 degree between 50 and 60 mT is neglected. This hump is probably related to a local variation in the microstructure of the sample. Note, that a quantitative comparison of the two measurement makes sense in this case because both films exhibit a field independent critical current density of $`j_c=2.5\times 10^{11}`$ A/m<sup>2</sup> in the unperturbed region.
The parallel shape for different angles suggests a universal field dependence of critical currents across LAGB’s which can be totally seperated from the microstructural properties. The uniform shape of the $`B`$$`j`$ relations is an evidence for a additional pinning mechanism of the flux lines which is independent of the microstructural pinning of the grain boundary. Only the local magnetic flux density and thus the flux line–flux line distance originates this effect. As a consequence the critical current can be written as
$$j_c(\varphi ,B)=j_{c1}(\varphi )+j_{c2}(B).$$
(3)
In this equation $`j_{c1}(\varphi )`$ represents the part of the critical current density which is caused by the intrinsic pinning of the grain boundary. This $`j_{c1}`$ shows the well known exponential decay with increasing misorientation angle $`\varphi `$. $`j_{c2}(B)`$ has its origin only in flux line–flux line interaction and is totally independent of the microstructure of the grain boundary. We focus now on the contribution $`j_{c2}(B)`$ and try to understand the magnetic field dependency that we observe in our measurements.
To explain the shape of the B–j curve we assume a single vortex located exactly on the LAGB and take a look at the interaction with an Abrikosov flux line lattice in the vicinity of the grain boundary in presence of a Lorentz force. A sketch of the chosen model geometry is shown in Fig. 8. The flux line on the grain boundary is represented as dark gray circle at the top in Fig. 8, the neighboring vortices are presented light gray and white. In the following the force per unit length shall be calculated that is required to drive the vortex along the grain boundary through the nearest neighbors.
For this calculation numerous assumptions are made which have to be discussed first. The complex flux line–flux line interaction is reduced to magnetic interaction. The distance between two flux lines is several hundred nanometers in the considered flux density range which is at least two orders of magnitude larger than the coherence length. Therefore the flux line core interaction can be neglected in this first order calculation. Also neglected is the anisotropy of the vortex which is located exactly on the LAGB. A flux line on a grain boundary shows a crossover from an isotropic Abrikosov vortex to an extremely anisotropic Josephson vortex with increasing misorientation angle . A certain degree of anisotropy definitely appears in case of the LAGB’s, but the fact that the misorientation angles are very small gives rise to use isotropic vortices in this model. In addition to this the flux lines in the unperturbed film are assumed to be immobile and only the interaction with the nearest neighbors is concerned. The model neglects any bending effects of the flux lines, e. g. only the two–dimensional projection of the vortices is taken into account.
With all these restrictions the pinning contribution of this model can be calculated. The magnetic interaction force (per unit length) between two vortices is given by deGennes
$$F_{IA}=\frac{\mathrm{\Phi }_0^2}{2\pi \lambda ^3\mu _0}K_1\left(\frac{a}{\lambda }\right).$$
Here $`\lambda `$ is the London penetration depth, which is $`\lambda 150`$ nm at $`T=5`$ K, $`\mathrm{\Phi }_0`$ is the flux quantum, $`a`$ the flux line distance and $`K_1`$ the modified Bessel function or MacDonald function of first order. The interaction with the two nearest neighbors in Fig. 8 compensates the Lorentz force $`f_L`$, that tries to move the flux line in the LAGB towards the two light gray flux lines. The pinning force of this geometry is now given by the maximum force that appears, if the flux line in the LAGB is forced to pass through the two nearest neighbors
$$F_{pin}=\begin{array}{c}\\ \mathrm{max}\\ \mathrm{a},\theta \end{array}\left[2F_{IA}(\frac{a}{\mathrm{sin}\theta })\mathrm{cos}\theta \right],$$
$`\theta `$ is defined in Fig. 8. To obtain the contribution to the critical current density across the LAGB it is necessary to calculate the current density from the force per unit length. Due to the fact that the magnetic interaction force is present over the whole length of the vortex, one obtains easily $`j_{c2}(a)=F_{pin}(a)/\mathrm{\Phi }_0`$, which is the well known definition of a Lorentz force. For better comparison to the experimental data, the dependency on the flux line–flux line distance $`j(a)`$ is transformed into a flux density dependence $`j(B)`$ using $`B=2\mathrm{\Phi }_0/\sqrt{3}a^2`$ for a triangular Abrikosov lattice.
The resulting relation for the interesting flux density range is plotted in Fig. 9.
The plot shows a similar increase of the critical current density with increasing flux density as found in the experimental data and the calculation yields current densities of the right order of magnitude.
For an optimal comparison to the measurement the field independent part $`j_{c1}`$ of the critcal current density of the grain boundary has to be estimated. This can be performed by comparing the data for very low magnetic flux densities where the contribution $`j_{c2}`$ is small. Using the data below a local flux density $`B`$ = 30 mT a value of $`j_{c1}=1.4\times 10^{11}`$ A/m<sup>2</sup> fits the data best in case of the 2 grain boundary. Plotting now $`j_c=j_{c1}+j_{c2}`$ versus the experimental data one obtains Fig. 10.
Both the experimental and the calculated curve show a similar shape. The measured data show a stronger increase than the model predicts but the slope of both curves is in the same range. The largest deviation from the model prediction is found for higher magnetic fields. A possible improvement of the very simple model might be the consideration of more than just nearest–neighbor interaction of the flux lines especially for the higher field range .
To summarize our results, the critical current density across low–angle grain boundaries in thin films of YBCO is investigated by a high–resolution method. The analysis of the local field dependence of the critical current shows a uniform behavior for different misorientation angles. This uniformity can be explained by seperating two parts that contribute to the critical current density. One part is correlated to the microstructural properties of the grain boundary and shows the typical drop for increasing misorientation angles. The second part is independent of the microstructure and can be described by vortex–vortex interaction in and in the vicinity of the grain boundary. A model which takes the deGennes magnetic interaction into account is able to reproduce the measured current densities.
The authors are grateful to Ch. Jooss, R. Warthmann and M. V. Indenbom for stimulating and helpful discussions and to G. Cristiani and H.–U. Habermeier for the preparation of the excellent samples.
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# Parametrization of Multiple Pathways in Proteins: Fast Folding versus Tight Transitions
## I Introduction
Globular proteins exists in two states the native, or folded, state and the denaturated, or unfolded, state . The folding of a protein is mostly driven by thermodynamics, and it typically occurs on a time scale of order seconds. The folding, and the resulting conformations result from an interplay between polar and nonpolar amino acids on the protein, their hydrogen bonds, and how all these interact with and constrain the surrounding water. The folding is therefore an immensely complex process and to date it is still a completely open problem how one may predict a protein’s folding properties given its structure.
On a qualitative level it is remarkable that the native state of a globular protein is solid with specified positions of all amino acids. This fact is thermodynamically reflected in a denaturation phase transition which is exceedingly sharp, in spite of the small number of molecules in the protein. Enhancing this qualitatively remarkable phenomenon is the fact that a typical polymer folds with a very weak transition . In fact single domain protein folding transition fairly accurately behaves as a two-state thermodynamical system where only denatured and native state can be populated and nearly no states with energies between these two extremal states are observed. Another remarkable feature of proteins is that they reach their unique ground state in a time that is much smaller than the time needed to search the extremely large configurational space. This is commonly known as the Levinthal paradox, and is noteworthy because it is difficult to merge with the two-state thermodynamics of proteins.
Research into the more specific aspects of the dynamics of protein folding is divided into several sub camps. For example, there is a huge effort to study simplified models numerically using ideas originating from spin glass theory. The basic idea here is that the different possible conformations of a protein define an energy landscape containing multiple minima. Such a picture is supported by experimental evidence at very low temperatures and structures confined by potentials less than about 2kcal/mol . However, if the dynamics of proteins at normal temperatures were that of spin glasses, its folding would be even slower than a random search in configuration space, resulting in time scales on an astronomical scale. This is not the case. Today, one imagines two possible ways to explain how proteins can fold on reasonable timescales: (1) The existence of folding funnels , or (2) the existence of folding pathways . The basic idea behind the folding funnel is that it is a deformation of the energy landscape in such a way that it resembles a funnel. The protein then “falls” into the global ground state without getting stuck.
The folding pathway is a concept related to a specific sequence of events. In contrast to a funnel, a pathway is characterized by subsequent formation of identifiable structures that must be formed in a specific order. Pathway studies of proteins that fold have been extensively carried out by Fersht . Occasionally one observes more than one possible pathway. The mechanism leading to the existence of guiding pathways may be hierarchical from small scales to larger scales in the sense of Baldwin and Rose , or it may be purely sequential as suggested in . It is the aim of this paper to demonstrate a third possibility, that encompass both of the above suggested scenarios but do it in a way that ensure the experimentally observed two state nature of protein folding-unfolding.
In the next section, we introduce and discuss our parametrization scheme for protein folding. This is based on our earlier work on one-pathway protein folding . In order for a protein to be biologically relevant, certain demands has to be put on its folding properties. One such demand is that the folding-unfolding transition is sharp. The sharpest possible transition is obtained when the system is two state. We analyse in Section III a class of Hamiltonians parametrized by the scheme outlined in Section II B with respect to sharpness using the van’t Hoff relation. We find that the typical system within this class behaves as a two-state system with maximal sharpness of transition. In Section IV, we finally demonstrate that the class of Hierarchical folders encoded by the parametrization scheme of Section II have folding times compatible with the ones that are experimentally observed. Further we demonstrate that the phase transition is sharpened by an entropic barrier and therefore that the folding time increase strongly with temperature $`T`$ in some $`T`$ interval below the equilibrium transition temperature $`T_c`$. In Section V we conclude that both demands, the sharpness of transition and the short folding times, are encoded by sequentially encoded hierarchical folders.
## II Pathways
We take the point of view that proteins follow certain pathways when folding from the denaturated to the native state. The “standard” theorist’s way of envisionin such a pathway is as a curve in some rugged energy landscape leading to the global minimum point — which corresponds to the native state. In the present work, we take a different point of view. We define a pathway as a sequence of partially folded structures which sequentially reduces both energy and entropy of the protein, and thus guide it towards the native state. By a partially folded structure we mean a protein conformation where parts of it has taken on the structure of the native state.
The simplest model of subsequent entropy pruning is a zipper , where one at each step $`k`$ of the folding gains one unit of energy, say $`E(n+1)=E(n)a`$, and simultaneously reduces the available number of conformational states with a factor $`g`$, i.e., state space $`S(k)S(k+1)=S(k)/g`$. If each step is associated with binding of one amino acid from a random coil then $`g6`$ corresponding to the possible orientations of one amino acid along the chain . If we deal with reduction of effective number of orientations form a compact molten globule, then $`g6/e2`$ where the factor $`e`$ comes from excluded volume effects . For individual amino acid contacts the corresponding energy gains may be estimated from the Miyazawa-Jernigan matrix , suggesting energies $`a`$ of order $`12`$ kcal/mol. In the first subsection we will explore the mathematical formulation of this simple zipper model, and demonstrate that it indeed predicts a first order phase transition, albeit a fairly weak one.
Folding of real proteins are presumably more complicated than a simple zipper. Protein structure data collected by Baldwin and Rose support that the native state of real proteins is reached from progressively larger parts of the protein which are stabilized by the smaller structures already in place. This indeed suggest a sequential ordering of events with some similarity to the zipper, but with the important difference that folding may initiate at different places along the amino acid chain simultaneously. If this was the case, the resulting independence of events would imply a factorization of the partition function into sub-parts, with the result of weakening or eventually destroying the phase transition altogether. As discussed in Section III, the melting phase transition of a single domain protein is remarkably sharp, implying that folding the protein is a very cooperative and correlated process.
We will in Section III see that the observed sharpness of the phase transition of proteins puts severe constraint on any protein folding model, e.g. ruling out both the simplest version of the zipper model and the simplest version of the hierarchy model. In addition, there is growing evidence that proteins may follow alternate pathways during the folding process , thus requiring a formalism that merge the conflicting interest between the need of dealing with more than one pathway, to fold in a reasonable time, and to maintain the native state sharply folded when it finally is reached.
The purpose of this section is to develop such a formalism, thus encompassing both the pathway formalism and the possibility that folding proceed hierarchically from smaller to larger structures. We first assume that there is only one pathway. This will allow us to introduce the central ideas in our approach in the next subsection. In subsection II B we subsequently generalize theses ideas into a scenario where several pathways are allowed.
### A Single Pathway
A folding pathway is a sequence of partially folded structures. Each structure along the pathway may contain the previously folded structures as a subset. If there is only one pathway available for the protein, this will be the case. We are therefore dealing with a hierarchy where structures use previously folded structures as building blocks. It is also clear that the folding process towards the native state proceeds from smaller to larger scales.
We number the partially folded structures as they appear in the folding process, $`i=1,\mathrm{},N`$. The binary variable $`\psi _i`$ equals 1 if partially folded structure $`i`$ exists — also if it forms part of an even larger folded structure. Otherwise, it is zero. Thus, the folding pathway may be parametrized by the sequence of inequalities
$$\psi _1\psi _2\psi _3\mathrm{}\psi _N.$$
(1)
We will in the following refer to the $`\psi _i`$ variables as folding variables.
There is energy and entropy associated with the partially folded structures. There are the direct binding energies that keeps the folded structures together. Perturbing the folded parts of the protein requires more energy than perturbing the unfolded parts. The latter may change conformation with very low or no energy cost. Thus, these changes appear as entropy on the thermodynamical level. In order to model these mechanisms, we introduce a second set of variables, the structural variables $`\xi _1`$ to $`\xi _N`$. These variables contain e.g. the angles $`\varphi `$ and $`\psi `$ associated with $`NC_\alpha `$ and the $`C_\alpha C^{}`$ bonds of the part of the protein that folded. However, they are simplifications — and to simplify even further, we assume that they only take two values ,
$$\xi _i=\{\begin{array}{cc}1,\hfill & \\ 1\mathrm{\Xi }.\hfill & \end{array}$$
(2)
This latter simplification is not essential to retain an analytically tractable model. However, it makes it easier to introduce the model and discuss it in general terms. The choice of two states corresponds to a degeneracy factor $`g=2`$, and one can easily generalize the results we will obtain to larger $`g`$ values by simply assuming that $`g1`$ states are associated with the higher energy level.
We now proceed to construct the simplest possible Hamiltonian that reflect the above discussion . It is
$$=\underset{i=1}{\overset{N}{}}a^{(i)}\psi _i\xi _i,$$
(3)
where $`a^{(i)}`$ are positive coupling constants.
When the protein is at stage $`j`$ in the folding process, the energy associated with it is
$$E_i=\underset{i=1}{\overset{j}{}}\xi _i.$$
(4)
The parts of the protein not yet folded provide entropy, which in this particular case consists of the variables $`\xi _{j+1},\mathrm{},\xi _N`$. The degeneracy of this protein configuration is $`2^{Nj}`$.
Working with the Hamiltonian Eq. (3) with the constraints (1) is somewhat cumbersome. We will therefore make a coordinate transformation that builds the constraints into the Hamiltonian directly. As we will see in the next section, this transformation allows us to generalize the one-pathway Hamiltonian presented in this section to multiple-pathway Hamiltonians: Specify a set of pathways, and the corresponding Hamiltonian can be written down immediately.
We assume in the following that the value $`\mathrm{\Xi }`$ defined in Eq. (2) is so large that the $`\xi _i`$’s for which $`\psi _i=1`$, will not take on this value for any realistic temperature ($`T<100`$ C). In this case, we may set
$`\psi _1`$ $`=`$ $`\varphi _1,`$ (5)
$`\mathrm{}`$ (6)
$`\psi _i`$ $`=`$ $`\varphi _1\varphi _2\mathrm{}\varphi _i,`$ (7)
$`\mathrm{}`$ (8)
$`\psi _N`$ $`=`$ $`\varphi _1\varphi _2\mathrm{}\varphi _i\mathrm{}\varphi _N,`$ (9)
where there are no constraints on the new variables $`\varphi _i`$. We see that the constraints Eq. (1) automatically are obeyed. The Hamiltonian (3) now becomes
$$=a^{(1)}\varphi _1a^{(2)}\varphi _1\varphi _2\mathrm{}a^{(N)}\varphi _1\mathrm{}\varphi _N,$$
(11)
We furthermore note that the degeneracy provided by the $`\xi _i`$ variables (as they may take on two possible values, 1 and $`1\mathrm{\Xi }`$, when the corresponding $`\psi _i`$ is zero), is also build into the Hamiltonian. Assume e.g. that all $`\varphi _1`$ to $`\varphi _j`$ are equal to one, this corresponds to $`\psi _1=\mathrm{}\psi _j=1`$. Furthermore assume that $`\varphi _{j+1}=\psi _{j+1}=0`$. The remaining $`\varphi _{j+2}`$ to $`\varphi _N`$ can take on any value without change in energy. Thus, to within a factor 2, the degeneracy is the same that provided by the $`\xi _i`$ variables in the original Hamiltonian (3).
The Hamiltonian (11) has a first order transition at a temperature $`T=1/\mathrm{log}2`$ when $`a^{(1)}=\mathrm{}=a^{(N)}=1`$ . This Hamiltonian is easily generalized to take into account the coupling between water and protein .
### B Multiple Pathways
We need to define what we mean by “multiple pathways.” We return to the single pathway defined using the folding variables by Eq. (II A). Suppose we break this sequence of inequalities by removing one of them, say $`\psi _i\psi _{i+1}`$, so that we have
$$\psi _1\mathrm{}\psi _i,$$
(12)
and
$$\psi _{i+1}\mathrm{}\psi _N.$$
(13)
Now, the status (0 or 1) of any $`\psi _j`$, where $`ji`$ is independent of the status of any $`\psi _k`$, where $`k>i`$. That is, we have created two independent folding “domains.”
In terms of the $`\varphi _i`$ variables, the two groups of inequalities (12) and (13), may be expressed as
$`\psi _1`$ $`=`$ $`\varphi _1,`$ (14)
$`\mathrm{}`$ (15)
$`\psi _i`$ $`=`$ $`\varphi _1\mathrm{}\varphi _i,`$ (16)
and
$`\psi _{i+1}`$ $`=`$ $`\varphi _{i+1},`$ (18)
$`\mathrm{}`$ (19)
$`\psi _N`$ $`=`$ $`\varphi _{i+1}\mathrm{}\varphi _N.`$ (20)
We may also construct the corresponding Hamiltonian,
$$=a^{(1)}\varphi _1a^{(2)}\varphi _1\varphi _2\mathrm{}a^{(i)}\varphi _1\mathrm{}\varphi _ia^{(i+1)}\varphi _{i+1}\mathrm{}a^{(N)}\varphi _{i+1}\mathrm{}\varphi _N.$$
(22)
When setting up the Hamiltonian (22), we have made the assumption that both pathways, $`12\mathrm{}i`$ and $`(i+1)\mathrm{}N`$ have the same degeneracies associated with them, fixing the relative entropic contribution from each branch. By introducing additional independent variables $`\varphi _{N+j}`$, we may tailor these contributions. For example, the Hamiltonian
$$=a^{(1)}\varphi _1a^{(2)}\varphi _1\varphi _2\mathrm{}a^{(i)}\varphi _1\mathrm{}\varphi _ia^{(i+1)}\varphi _{i+1}\varphi _{N+1}\mathrm{}a^{(N)}\varphi _{i+1}\mathrm{}\varphi _N\varphi _{N+1}$$
(23)
will increase the entropic contribution of the second branch by $`T\mathrm{log}2`$ compared to the first branch.
We may generalize the ideas presented above to systems with multiple folding pathways. This is done simply by generating telescoping groups of $`\varphi _i`$ variables, one for each pathway. The corresponding Hamiltonian may be written down directly,
$$=\underset{i}{}a_i^{(1)}\varphi _i\underset{ij}{}a_{ij}^{(2)}\varphi _i\varphi _j\mathrm{}\underset{ijk\mathrm{}}{}a_{ijk\mathrm{}}^{(N)}\varphi _i\varphi _j\varphi _k\mathrm{}.$$
(24)
We have here added subscripts to the $`a^{(i)}`$ coefficients, reflecting which $`\varphi _i`$ variables are involved in the corresponding term in the Hamiltonian.
In order to construct a Hamiltonian given a set of pathways, we start with organizing the pathways in a branching structure where each numbering refer to a folded substructure of the protein, as shown in Fig. 1. The interpretation of this figure is that folding of structures 1 and 5 are independent starting points of the folding process. In order to structure 4 to be formed, structure 5 must already be in place. However, this structure is independent of structure 1 having been formed. Structure 2, on the other hand, needs both structures 1 and 5 in place. Lastly, structure 3 can only form if both structures 2 and 4 have formed. In the end, all structures must always be formed. A concrete example of such a set of pathways has been reported for staphylococcal nuclease .
Assuming now that the reduction in degeneracy when either structure 1 or 5 form is the same, we may set up a “folding table” as shown in Table I. From this table, we read off the Hamiltonian
$$=a_1^{(1)}\varphi _1a_5^{(1)}\varphi _5a_{45}^{(2)}\varphi _4\varphi _5a_{125}^{(3)}\varphi _1\varphi _2\varphi _5a_{12345}^{(5)}\varphi _1\varphi _2\varphi _3\varphi _4\varphi _5.$$
(25)
In the following, we will study the thermodynamical and dynamical properties of Hamiltonians constructed in this way. There are two main questions to be addressed which are essential in connection with protein folding: (1) sharpness of folding transition and (2) shortness of folding time. In the next section, we will consider the sharpness of the transition.
## III Sharpness of Transition
An excellent measure of sharpness of a first-order transition is to study the van’t Hoff relation which relates enthalpy change through the transition $`\mathrm{\Delta }H`$, the peak of the heat capacity $`C_\mathrm{p}`$, the absorbed heat, $`Q`$ through the formula
$$\mathrm{\Delta }H=\alpha \frac{T_\mathrm{c}^2C_\mathrm{p}(T_\mathrm{c})}{Q}.$$
(26)
Here $`T_\mathrm{c}`$ is the transition temperature, and $`\alpha `$ is the van’t Hoff coefficient. The smaller this coefficient is, the sharper the transition.
We may demonstrate this in the following way. In the present case, $`\mathrm{\Delta }H=Q`$ (as there is no pressure in the system), which is the area of the bump in the heat capacity plot. The height of the hump, $`h`$, is $`C_\mathrm{p}`$, while the area $`Q`$ is essentially the width of the top $`b`$ multiplied by the height, $`h`$, $`Q=bh`$. We may thus write the van’t Hoff equation
$$b^2h=\alpha T_\mathrm{c}^2.$$
(27)
¿From this equation it might seem that the opposite of what we are claiming is true: The larger the $`\alpha `$, the wider the transition, as the height $`h`$ is proportional to $`\alpha `$. However, one must not forget the presence of the width $`b`$. We rewrite the left hand side of Eq. (27) as
$$(bh)^2\frac{1}{h}=\alpha T_\mathrm{c}^2.$$
(28)
Now keeping the area of the hump, $`bh=Q`$, constant, we see that the height $`h`$ is inversely proportional to $`\alpha `$. Hence, our statement that a smaller $`\alpha `$ implies a sharper transition follows, as sharpness is a relative concept.
In order to understand the the $`\alpha `$ coefficient properly, we examine its physical meaning in the following. Let us assume a system with an equally spaced energy spectrum, 1,…,$`N`$. We furthermore assume a degeneracy $`g_n`$ for the $`n`$th energy level. The partition function is then
$$Z=\underset{n=0}{\overset{N}{}}g_n\mathrm{e}^{n/T}.$$
(29)
For proteins there should be a unique ground state, so that $`g_N=1`$. The energy is given by
$$E(T)=T^2\frac{}{T}\mathrm{log}Z=n,$$
(30)
and the heat capacity is given by
$$C(T)=\frac{}{T}E(T)=\frac{n^2n^2}{T^2}.$$
(31)
These relations are of course valid at any temperature. The $`\alpha `$ factor can now be expressed in terms of the fluctuations of the energy levels at $`T=T_\mathrm{c}`$,
$$n^2n^2=\frac{N^2}{\alpha }.$$
(32)
For any probability distribution we must have $`n^2n^2N^2/4`$. We only have equality when $`P(n=0)=P(n=N)=1/2`$. Thus in general
$$\alpha 4.$$
(33)
Equality $`\alpha =4`$ is only found when the system exhibits two state behaviour, i.e., when $`g_0=g_N\mathrm{e}^{N/T}`$ and otherwise $`g_n=0`$. This is fulfilled with the Hamiltonian $`=N\varphi _1\mathrm{}\varphi _N`$, i.e., when all terms $`a^{(i)}=0`$ for all $`i<N`$ in Eq. (22). In general Eq. (32) can be reformulated as
$$s^2s^2=\frac{1}{\alpha }.$$
(34)
which express $`\alpha `$ as the variance of the distribution of normalized energy states ($`s=n/N`$) at the transition point. This distribution could be bimodal, ($`\alpha =4`$), be flat ($`\alpha =12`$) or be centered around a mean intermediate energy ($`\alpha >12`$). The flat distribution is obtained when all $`a^{(i)}=const`$ in the Hamiltonian (22). Of special interest is the class of Hamiltonians that correspond to the hierarchical folding suggested by Baldwin and Rose . We show in Fig. 2 an example of such a hierarchical folding network with four levels. In this scheme the folding of smaller units form a necessary but not sufficient condition for the folding of larger units. A corresponding Hamiltonian may be constructed as outlined in Section II B, and we find
$$=a_1^{(1)}\varphi _1a_2^{(1)}\varphi _2a_3^{(1)}\varphi _3a_4^{(1)}\varphi _4a_{125}^{(3)}\varphi _1\varphi _2\varphi _5a_{346}^{(3)}\varphi _3\varphi _4\varphi _6a_{1234567}^{(7)}\varphi _1\varphi _2\varphi _3\varphi _4\varphi _5\varphi _6\varphi _7.$$
(35)
At each level in the hierarchy a new variable is needed in order to make folding at level below a necessary but sufficient condition for folding at this level.
As the number of terms in this hierarchical Hamiltonian grows, $`\alpha \mathrm{}`$. The reason for this is the statistical broadening due to the exponentially growing number of independent branches when moving “backwards” in the hierarchy towards smaller units.
Privalov and Khechinasvili measured experimentally the value
$$\alpha =\mathrm{\hspace{0.33em}4.2}\pm 0.2.$$
(36)
for large group of small globular proteins. Thus real single domain proteins are close to two-state behaviour, i.e., far from a simply guided Hamiltonian, and very far from protein folding viewed as subsequent folding of independent subunits. Folding of single domain proteins, that in practical terms means proteins of about 100 amino acids is indeed close to a maximally cooperative process.
### A A Class of Hierarchical Hamiltonians with $`\alpha =4`$.
Based on spatial correlations of amino acids in real proteins, Baldwin and Rose argue convincingly that some sort of subsequent folding hierarchy must be implemented in real protein folding . In last section we saw that the simplest hierarchical scheme is excluded for single domain proteins. We now explore a more elaborate hierarchical scheme, where correlations between folding subunits naturally becomes implemented on an early level. The simplest of these is the single pathway model of Eq. (11). A general formulation is the Hamiltonian (24). It has up to $`N`$ groups of products of various length of $`\varphi _i`$ variables. We will in next subsection investigate the broad subclass of these Hamiltonians which are restricted by having maximum one term of each product number, and see that they typically have $`\alpha 4`$. This means that we investigate Hamiltonians as in Eq. (24) where the coefficient $`a_{ijk\mathrm{}}^{(m)}=1`$ for one combination of $`ijk\mathrm{}`$ and otherwise zero. Thus, there is one term with only one $`\varphi _i`$, one term with two $`\varphi _i`$ variables etc.
We note that the folding system in Fig. 1 as described by Eq. (25) does not include only one term at each level. However, the system in Fig. 1 can also be parametrized by introducing one additional auxiliary variable $`\varphi _6`$, thereby creating an entropy barrier in the thermodynamics. The corresponding Hamiltonian for Fig. 1 would then be
$$=a_1^{(1)}\varphi _1a_5^{(1)}\varphi _5\varphi _6a_{456}^{(3)}\varphi _4\varphi _5\varphi _6a_{1256}^{(4)}\varphi _1\varphi _2\varphi _5\varphi _6a_{123456}^{(6)}\varphi _1\varphi _2\varphi _3\varphi _4\varphi _5\varphi _6.$$
(37)
This is the type of Hamiltonians that we wish to study. This subclass deals with systems where indeed different regions may start independent folding, but where some foldings are more difficult to initiate and thus becomes rate limiting. In other words, that different events takes very different times implies that some folding events acts as effective nucleation centers. When these barriers are overcome, subsequent foldings spread across the system along some fairly well defined paths.
The nucleation barriers were associated to cases where several variables must fold simultaneously without previous guiding ($`\varphi _5`$ and $`\varphi _6`$ in the above Hamiltonian (37)). In addition, different folding paths typically become correlated e.g. through a common origin, as in the above Hamiltonian where the events $`\varphi _5`$ and $`\varphi _6`$ both are needed in order to proceed folding along the parallel ways parametrized by respectively $`\varphi _2`$ and $`\varphi _4`$. Such correlation could e.g. be due to steric interactions that facilitate folding by bringing far-apart amino acids closer to each other. These associated entropy barriers and correlations are responsible for the effective two state behaviour that this class of hierarchical models exhibit. We will examine this in detail in the next section.
### B Averaging over the Class of Hamiltonians
We will in following show that the type of Hamiltonians discussed in the previous section displays remarkable thermodynamic properties. We shall consider more specificly the class of Hamiltonians defined by Eq. (24), where there is only one term for each number of $`\varphi _i`$s. The partition function for a Hamiltonians belonging to this class is
$$Z=\underset{\{\varphi \}}{}\mathrm{e}^{/T}.$$
(38)
To gain insight into the thermodynamics of this class of Hamiltonians, we have calculated the mean of the partition function over all possible Hamiltonians. The result is
$$Z=1+\underset{s=1}{\overset{N}{}}\left(\begin{array}{c}N\\ s\end{array}\right)\underset{k=1}{\overset{s}{}}\left[1+\frac{s!(Nk)!}{N!(sk)!}(\mathrm{e}^{1/T}1)\right].$$
(39)
The calculation is outlined in Appendix A. This gives us the thermodynamics of an average Hamiltonian. In order to find out if this is a representative result we have to calculate the fluctuations around this mean. All thermodynamics properties of a given Hamiltonian is derived from $`\mathrm{log}Z`$, and the mean can be expressed as $`\mathrm{log}Z=\mathrm{log}z+\mathrm{log}Z`$, with $`z=Z/Z`$. We have in Appendix B transformed the expression for the fluctuations of the partition function into the function, Eq. (B20) and the result is plotted in Fig. 3. The fluctuations are extremely small, in fact $`z^2z^20`$ as the system size increases. For example for for $`N100`$ they contribute $`10^4k_\mathrm{B}T`$ to the free energy, and hence without any physical relevance. In other words two randomly chosen Hamiltonians will have exactly the same thermodynamic properties, the system is self-averaging and Eq. (39) is the exact solution of the partition function.
In order to study the folding transition, we define an order parameter, $`u`$, by the mean of the sum over all the $`\varphi `$ variables
$$u=\frac{2}{N}\underset{i=1}{\overset{N}{}}\varphi _i1.$$
(40)
When the protein is in the ground state $`\varphi _i=1`$ and the order parameter is $`u=1`$, at high temperatures $`\varphi _i=1/2`$ and $`u=0`$. Since the class of hamiltonians that we study are self-averaging for large $`N`$, we may calculate the thermal mean values by taking the average over all Hamiltonians. We may thus write
$`{\displaystyle \underset{i}{}}\varphi _i`$ $`=`$ $`{\displaystyle \frac{1}{M}}{\displaystyle \underset{\{a\}}{}}{\displaystyle \underset{\{\varphi \}}{}}{\displaystyle \underset{i}{}}\varphi _i\mathrm{e}^{/T}/Z`$ (41)
$`=`$ $`{\displaystyle \frac{1}{Z}}{\displaystyle \underset{s=1}{\overset{N}{}}}s\left(\begin{array}{c}N\\ s\end{array}\right){\displaystyle \underset{k=1}{\overset{s}{}}}\left[1+{\displaystyle \frac{s!(Nk)!}{N!(sk)!}}(\mathrm{e}^{1/T}1)\right],`$ (44)
where $`M`$ is the total number of Hamiltonians. Fig. 4 shows the order parameter at different temperatures. We clearly see that there is sharp first-order transition around the temperature $`T_\mathrm{c}1.44`$. The sharpness of the transition is characterized by the van’t Hoff $`\alpha `$ coefficient, which is calculated from the partition function Eq. (39). The results plotted in Fig. 5 show that it goes quickly to $`\alpha =4`$ for large $`N`$, which is the lowest possible value, i.e. the sharpest possible transition. We conclude hence that the folding scheme we have introduced gives a first-order phase transition with a van’t Hoff coefficient of $`\alpha 4`$, corresponding to a two-state thermodynamics.
We note that for the single-path Hamiltonian, Eq. (11) with all $`a^{(i)}=1`$, $`\alpha `$ equals 12. However, if the last term is set equal to $`a^{(N)}=ϵN`$, where $`ϵ0`$ is small (and goes to zero as $`N`$ in the limit $`N\mathrm{}`$), then $`\alpha 4`$ .
## IV Folding Time
Proteins in their native form have well defined structures that are separated from their denatured counterparts through some barrier, as quantified by value of $`\alpha `$. The denaturated state consists of an astronomical numbers of microscopic conformations, whereas the native state is unique. The fact that the unique native state anyway is reached fairly fast from this large number of unfolded configurations is known as the Levinthal paradox. Thus, the folding proteins must have some property encoded which lower configurational entropy systematically as folding proceeds. However, such guided folding has to be merged with the fact that the folding process exhibits effectively two-state behaviour, i.e. has few or no detectable intermediates , respectively $`\alpha 4`$ .
Minimalization of folding time and the van’t Hoff parameter $`\alpha `$ poses seemingly contradictory goals. This one might expect because absence of significant intermediates restricts the possibility of guided folding between the denatured states and the native state. As an example the following Hamiltonian has $`\alpha =4`$, i.e., a maximally sharp transition ,
$$=N\varphi _1\varphi _2\mathrm{}\varphi _N.$$
(45)
However, it has absolutely no guiding at all: When in the denaturated state, i.e., when at least one $`\varphi _i=0`$, there is no energy gain in turning this to $`\varphi _i=1`$, unless all other variables $`\varphi _j`$ already are equal to one. There is no energy gain unless such a protein moves into the native state. Searching for the native state in this system will grow exponentially with system size $`N`$. In general, it may seem that the more “two state” the transition is, the less guiding can be involved, and hence, the longer the folding time will be.
In order to discuss folding times, we need an operational definition of folding time. We base our definition on the one-step Monte Carlo method . The Monte Carlo method is a way to generate a biased random walk through configuration space such that the relative frequence of occupancy of any configuration is proportional to the Boltzmann factor. Thus Monte Carlo “time”, measuring the number of steps of the random walk, has strictly speaking nothing to do with “real” time. However, in a discrete system, such as our model, we believe that time defined in this way, essentially brings us as close as possible to time in a corresponding continuous system. Hence, we will base our definition of time in the model by one-step Monte Carlo. Hence, we define one unit of time as the times where all variables have been updated once in average. This definition of time is common when dealing with dynamical questions concerning discrete system .
The non-guided Hamiltonian (45) will fold in a time proportional to $`2^N/N`$ when starting from a random initial configuration. The fully guided (zipper) Hamiltonian (11) will similarly fold in a time proportional to $`N`$. The Baldwin-Rose picture of hierarchical folding will give a folding time proportional to $`\mathrm{log}N`$. These examples, as noted above, seem to suggest that Hamiltonians with larger $`\alpha `$ fold faster. However real single-domain proteins have $`\alpha `$ close to its minimum value 4, and are still able to fold in fairly short time.
We have in Section III A suggested a class of hierarchical folders, and seen that these typically have $`\alpha 4`$ that is compatible with two-state behaviour seen in experiments.
In order to examine the folding behavior, we restrict ourselves initially to the case of zero temperature, i.e., when a term has folded it never opens again. The fact that we at all can reach the ground state at zero temperature follows from all terms in the Hamiltonian being attractive (i.e., negative signs) and that the various terms in the Hamiltonian guide the system toward the uniquely defined ground state (where all $`\varphi _i=1`$). Thus, there are no energy barriers in our idealized model system, and folding is therefore only limited by degeneracy/entropy barriers.
The zero temperature limit makes it possible to characterize the folding as function of number of subsequent folding steps, where one step is defined as an event where at least one term in the Hamiltonian becomes non-zero. Fig. 6a shows the time to fold subsequent steps of the Hamiltonian for a system of size $`N=100`$. The ($`+`$) and ($`x`$) symbols denote trajectories of two different Hamiltonians of type (24), whereas the full line show the ensemble averaged folding time (over all Hamiltonians and trajectories). We note that although the trajectories typically fold subsequent levels of the Hamiltonian, this is not necessarily the case, thus opening for alternate pathways.
Entropy barriers are caused by the necessity to fold several $`\varphi _i`$ terms simultaneously, and the chance to do this decreases as the $`1/2^k`$, where $`k`$ is the number of terms. In Fig. 7a we display the entropy barrier $`dS/dk`$ defined as the number of simultaneous folding terms versus the folding step $`k`$ for the $`N=100`$ system. The histogram is $`dS/dk`$ for one particular Hamiltonian, whereas the dashed line shows the barriers for the ensemble averaged values. Fig. 7b displays a data collapse of $`dS/dk`$ for three different system sizes, demonstrating that both the position and the maximal barrier scale with system size approximately as $`\sqrt{N}`$, which implies that the folding time grows with systems size as $`2^\sqrt{N}`$. An argument for the obtained scaling is that after $`k`$ folding steps, then as long as $`k`$ is small, one typically has folded $`k^2/2`$ terms. When $`k^2/2`$ becomes comparable to $`N`$, overlap between subsequent steps becomes significant, and subsequent folding involves fewer new folding variables and therefore folding becomes easier.
In order to estimate whether these numbers are consistent with a fine grained behavior of real proteins then the degeneracy of each degree of freedom should be of the order of 6, corresponding to the orientational possibilities of a single amino acid on a random peptide chain. Thus exchanging $`26`$ in the above analysis, and noting that individual amino acids change conformations on a timescale of the order of nanoseconds , we very roughly estimate a folding time of $`6^\sqrt{N}\times 10^9\mathrm{s}10^1`$ s for $`N=100`$ which is of the order found for proteins.
Fig. 8a examines the zero temperature energy as function of folding step for a specific $`N=40`$ system, where we associate one unit of energy to each contact term in the Hamiltonian. Thus the coefficients $`a^{(i)}=1`$, implying that all Hamiltonians of type (24) have equilibrium transition temperature $`T_c=1/\mathrm{log}2`$. The figure shows that the first $`k12`$ steps is associated with a slow linear decrease in energy, hereafter a few subsequent steps is responsible for $`65`$% of the total energy gain due to folding. Combining this with the fact that nearly all entropy reduction has to take place for $`k<12`$ (see Fig. 7), we conclude that entropy is reduced before energy is gained.
We now examine the effect of finite temperature where the increased importance of entropy converts entropy barriers into free energy barriers. This makes the initial folding steps thermodynamically disfavoured at temperatures below the equilibrium melting temperature $`T_c`$. In Fig. 8b this is examined quantitatively, where free energy is shown versus an effective order parameter $`u^{}`$. defined as the number of $`\phi `$ variable which is equal to one and belongs to a term in the Hamiltonian which contributes a finite amount (i.e. $`E=1`$) to the energy. Thus $`u^{}`$ is linearly decreasing with the residual entropy of the system. The plot show the free energy profile for three temperatures, the $`T=0`$ case, the temperature where folding becomes dynamically suppressed, $`T=T^{}=\mathrm{\hspace{0.33em}0.27}T_c`$ (for $`N=40`$), and the equilibrium denaturation temperature ($`T_c=1/\mathrm{log}2`$) where the ground state just balances the denatured state. We conclude that in our simplified scenario, the sharpness of the transition at $`T_c`$ arises because of an entropy barrier that will make folding much slower in some limited interval below the equilibrium melting temperature. With the parameters of the figure one first obtains the fairly fast folding examined in Fig. 6 when the temperature becomes below $`T^{}0.27T_c`$.
We stress that one in principle could include energy barriers into our formalism by adding positive terms to the hamiltonian (i.e. $`a_{ijk\mathrm{}}^{(n)}>0`$ in Eq. (24)). Such terms could arise either if the denatured state has binding contacts which are not in the native state or if the transition state have non native contacts. Here we focussed on the possibility for a transition that is purely limited by entropy barriers, thereby defining a transition state which has little entropy and large enthalpy difference to the ground state, a feature observed in for $`\alpha `$spectrin SH3 by . As a consequence we predict a temperature dependence that is opposite to recent measurement on protein folding kinetics, where typically the folding time decrease with temperature $`T`$ for $`T<40^o`$C . Energy barriers therefore plays a role at least at these lower temperatures. Future measurement of folding time behaviour close to the denaturation temperature could teach us whether the two-state behaviour found for folding of small proteins is due primarily to enthalpic or entropic barriers.
## V Conclusion
In this paper we have explored the idea that entropy correlations and barriers may appear at any step during a sequential folding process. We have seen that this is natural to expect as consequence of a hierarchical folding process where different branches of the hierarchical folding structure are correlated. We have showed that this folding scheme can result in two-state thermodynamics even though the energy gain is microscopically small at each folding step. We have seen that this two state thermodynamics can be merged with a fairly fast folding process, with time increasing with $`N`$ as $`2^\sqrt{N}`$ instead of $`2^N`$ as one naively should expect for two-state folders. Further we have found the qualitatively same subdivision between initial and final folding as reported in , without the need to resorting to a special treatment of any terms. This subdivision into early entropically disfavoured guiding and last step easy match predicts a folding time that increases with temperature when denaturation is approached.
We thank A. Bakk and J. S. Høye for many fruitful discussions. P. G. D. thanks the NFR for financial support. P. G. D. and A. H. thank Nordita and the Niels Bohr Institute for friendly hospitality and support.
## A Calculating the average partition function
We can make $`N!/m!`$ different products of $`m`$ $`\varphi _i`$. Thus, there is in total
$$M=\frac{N!}{(N1)!}\frac{N!}{(N2)!}\mathrm{}\frac{N!}{(NN)!}$$
(A1)
different Hamiltonians (note, however, that most of them are equal up to permutations of $`\varphi _i`$).
We may calculate the mean value of the partition function
$`Z`$ $`=`$ $`{\displaystyle \frac{1}{M}}{\displaystyle \underset{\{a\}}{}}{\displaystyle \underset{\{\varphi \}}{}}\mathrm{e}^{/T}`$ (A2)
$`=`$ $`1+{\displaystyle \frac{1}{M}}{\displaystyle \underset{\{\varphi \}}{}}{\displaystyle \underset{i}{}}\mathrm{e}^{\varphi _i/T}{\displaystyle \underset{ij}{}}\mathrm{e}^{\varphi _i\varphi _j/T}\mathrm{}{\displaystyle \underset{ijk\mathrm{}}{}}\mathrm{e}^{\varphi _i\varphi _j\varphi _k\mathrm{}/T}`$ (A3)
The Hamiltonian is invariant to all permutations of $`\varphi `$, and this suggests the new variable $`s=_{i=1}^N\varphi _i`$. We then find
$$\underset{i}{}\mathrm{e}^{\varphi _i/T}=\frac{s!}{(s1)!}{}_{}{}^{1/T}+\frac{N!}{(N1)!}\frac{s!}{(s1)!},$$
(A4)
and
$$\underset{ij}{}\mathrm{e}^{\varphi _i\varphi _j/T}=\frac{s!}{(s2)!}{}_{}{}^{1/T}+\frac{N!}{(N2)!}\frac{s!}{(s2)!},$$
(A5)
and so on. This gives us an expression for the mean partition function
$$Z=1+\frac{1}{M}\underset{s=1}{\overset{N}{}}\left(\begin{array}{c}N\\ s\end{array}\right)\underset{k=1}{\overset{s}{}}\left[\frac{s!}{(sk)!}\mathrm{e}^{1/T}+\frac{N!}{(Nk)!}\frac{s!}{(sk)!}\right]\underset{k=s+1}{\overset{N}{}}\frac{N!}{(Nk)!}.$$
(A6)
This expression can be further simplified to
$$Z=1+\underset{s=1}{\overset{N}{}}\left(\begin{array}{c}N\\ s\end{array}\right)\underset{k=1}{\overset{s}{}}\left[1+\frac{s!(Nk)!}{N!(sk)!}(\mathrm{e}^{1/T}1)\right].$$
(A7)
## B Fluctuations of Partition Function
In this Appendix, we find the expression for $`Z^2`$ which was used to show that the relative fluctuations $`(Z^2Z^2)/Z^2`$ goes to zero as $`N\mathrm{}`$, as is seen in Fig. 3. We have that
$`Z^2`$ $`=`$ $`{\displaystyle \frac{1}{M}}{\displaystyle \underset{\{a\}}{}}{\displaystyle \underset{\{\varphi \}}{}}{\displaystyle \underset{\{\varphi ^{}\}}{}}\mathrm{e}^{H(\varphi )/TH(\varphi ^{})/T}`$ (B1)
$`=`$ $`{\displaystyle \underset{\{\varphi \}}{}}{\displaystyle \underset{\{\varphi ^{}\}}{}}G(\varphi ,\varphi ^{}),`$ (B2)
where
$$G(\varphi ,\varphi ^{})=\underset{\{\varphi \}}{}\underset{\{\varphi ^{}\}}{}\underset{i}{}\mathrm{e}^{(\varphi _i+\varphi _i^{})/T}\underset{ij}{}\mathrm{e}^{(\varphi _i\varphi _j+\varphi _i^{}\varphi _j^{})/T}\mathrm{}\underset{ijk\mathrm{}}{}\mathrm{e}^{(\varphi _i\varphi _j\varphi _k\mathrm{}+\varphi _i^{}\varphi _j^{}\varphi _k^{}\mathrm{})/T}.$$
(B3)
The sum over $`\varphi `$ and $`\varphi ^{}`$ is simplified by noting that the function $`G(\varphi ,\varphi ^{})`$ is symmetric to all permutations of pairs $`(\varphi _i,\varphi _i^{})`$. This means that $`G`$ can be expressed as
$$G(\varphi ,\varphi ^{})=G(s,s^{},p),$$
(B4)
with $`s=_{i=1}^N\varphi _i`$, $`s^{}=_{i=1}^N\varphi _i^{}`$ and $`p=_{i=1}^N\varphi _i\varphi _i^{}`$. The sum over all $`\varphi _i`$ and $`\varphi _i^{}`$ is thus reduced to a sum over $`s`$, $`s^{}`$ and $`p`$.
$`Z^2`$ $`=`$ $`{\displaystyle \frac{1}{M}}{\displaystyle \underset{s=0}{\overset{N}{}}}{\displaystyle \underset{p=p_{\mathrm{min}}}{\overset{s}{}}}\left(\begin{array}{c}N\\ p\end{array}\right)\left(\begin{array}{c}Np\\ sp\end{array}\right)\left(\begin{array}{c}Ns\\ sp\end{array}\right)G(s,s,p)`$ (B11)
$`+`$ $`{\displaystyle \frac{1}{M}}2{\displaystyle \underset{s=1}{\overset{N}{}}}{\displaystyle \underset{s^{}=0}{\overset{s1}{}}}{\displaystyle \underset{p=p_{\mathrm{min}}}{\overset{s^{}}{}}}\left(\begin{array}{c}N\\ p\end{array}\right)\left(\begin{array}{c}Np\\ sp\end{array}\right)\left(\begin{array}{c}Ns\\ s^{}p\end{array}\right)G(s,s^{},p),`$ (B18)
where $`p_{\mathrm{min}}=(s+s^{}N)\mathrm{\Theta }(s+s^{}N)`$. The function $`G(s,s^{},p)`$ is a product of $`N`$ sums. Each sum can be be evaluated. The results is
$$G(s,s^{},p)=\underset{m=1}{\overset{N}{}}A_m(s,s^{},p),$$
(B19)
where we have that
$`A_m(s,s^{},p)`$ $`=`$ $`\mathrm{\Theta }(pm){\displaystyle \frac{p!}{(pm)!}}\mathrm{e}^{2/T}`$ (B20)
$`+`$ $`\left[\mathrm{\Theta }(sm){\displaystyle \frac{s!}{(sm)!}}\mathrm{\Theta }(pm){\displaystyle \frac{p!}{(pm)!}}\right]\mathrm{e}^{1/T}`$ (B21)
$`+`$ $`\left[\mathrm{\Theta }(s^{}m){\displaystyle \frac{s^{}!}{(s^{}m)!}}\mathrm{\Theta }(pm){\displaystyle \frac{p!}{(pm)!}}\right]`$ (B22)
$`\mathrm{e}^{1/T}`$ (B23)
$`+`$ $`{\displaystyle \frac{N!}{(Nm)!}}+\mathrm{\Theta }(pm){\displaystyle \frac{p!}{(pm)!}}\mathrm{\Theta }(sm){\displaystyle \frac{s!}{(sm)!}}\mathrm{\Theta }(s^{}m){\displaystyle \frac{s^{}!}{(s^{}m)!}}.`$ (B24)
Here $`\mathrm{\Theta }(n)`$ is the step function ($`\mathrm{\Theta }(n)=1`$ for $`n0`$).
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# Can Prosody Aid the Automatic Classification of Dialog Acts in Conversational Speech?
## 1 Abstract
Identifying whether an utterance is a statement, question, greeting, and so forth is integral to effective automatic understanding of natural dialog. Little is known, however, about how such dialog acts (DAs) can be automatically classified in truly natural conversation. This study asks whether current approaches, which use mainly word information, could be improved by adding prosodic information.
The study is based on more than 1000 conversations from the Switchboard corpus. DAs were hand-annotated, and prosodic features (duration, pause, F0, energy, and speaking rate) were automatically extracted for each DA. In training, decision trees based on these features were inferred; trees were then applied to unseen test data to evaluate performance. Performance was evaluated for prosody models alone, and after combining the prosody models with word information—either from true words or from the output of an automatic speech recognizer.
For an overall classification task, as well as three subtasks, prosody made significant contributions to classification. Feature-specific analyses further revealed that although canonical features (such as F0 for questions) were important, less obvious features could compensate if canonical features were removed. Finally, in each task, integrating the prosodic model with a DA-specific statistical language model improved performance over that of the language model alone, especially for the case of recognized words. Results suggest that DAs are redundantly marked in natural conversation, and that a variety of automatically extractable prosodic features could aid dialog processing in speech applications.
Keywords: automatic dialog act classification, prosody, discourse modeling, speech understanding, spontaneous speech recognition.
## 2 Introduction
### 2.1 Why Model Dialog?
Identifying whether an utterance is a statement, question, greeting, and so forth is integral to understanding and producing natural dialog. Human listeners easily discriminate such dialog acts (DAs) in everyday conversation, responding in systematic ways to achieve the mutual goals of the participants ( ?, ?). Little is known, however, about how to build a fully automatic system that can successfully identify DAs occurring in natural conversation.
At first blush, such a goal may appear misguided, because most current computer dialog systems are designed for human-computer interactions in specific domains. Studying unconstrained human-human dialogs would seem to make the problem more difficult than necessary, since task-oriented dialog (whether human-human or human-computer) is by definition more constrained and hence easier to process. Nevertheless, for many other applications, as well as for basic research in dialog, developing DA classifiers for conversational speech is clearly an important goal. For example, optimal automatic summarization and segmentation of natural conversations (such as meetings or interviews) for archival and retrieval purposes requires not only knowing the string of words spoken, but also who asked questions, who answered them, whether answers were agreements or disagreements, and so forth. Another motivation for speech technology is to improve word recognition. Because dialog is highly conventional, different DAs tend to involve different word patterns or phrases. Knowledge about the likely DA of an utterance could therefore be applied to constrain word hypotheses in a speech recognizer. Modeling of DAs from human-human conversation can also guide the design of better and more natural human-computer interfaces. On the theoretical side, information about properties of natural utterances provides useful comparison data to check against descriptive models based on contrived examples or speech produced under laboratory settings. Automatic methods for classifying dialog acts could also be applied to the problem of labeling large databases when hand-annotation is not feasible, thereby providing data to further basic research.
### 2.2 Word-Based Approaches to Dialog Act Detection
Automatic modeling of dialog has gained interest in recent years, particularly in the domain of human-computer dialog applications. One line of work has focused on predicting the most probable next dialog act in a conversation, using mainly information about the DA history or context ( ?, ?, ?, ?, ?, ?, ?). A second, related line of research has focused on DA recognition and classification, taking into account both the DA history and features of the current DA itself ( ?, ?, ?, ?). In all of these previous approaches, DA classification has relied heavily on information that can be gleaned from words, such as cue phrases and N-grams, or information that can be derived from word sequences, such as syntactic form.
### 2.3 Why Use Prosody?
This work focuses on exploring another, relatively untapped potential knowledge source for automatic DA classification: prosody. By prosody we mean information about temporal, pitch, and energy characteristics of utterances that are independent of the words. We were interested in prosody for several reasons. First, some DAs are inherently ambiguous from word information alone. For example, declarative questions (e.g., “John is here?”) have the same word order as statements, and hence when lexical and syntactic cues are consistent with that of a statement, may be distinguishable as a question only via prosody. Second, in a real application, word recognition may not be perfect. Indeed, state-of-the-art recognizers still show over 30% word error rate for large-vocabulary conversational speech. Third, there are potential applications for which a full-fledged speech recognizer may not be available or practical, and a less computationally expensive, but somewhat less accurate method to track the structure of a dialog is acceptable. Fourth, an understanding of prosodic properties of different utterance types can lead to more natural output from speech synthesis systems. And finally, it is of basic theoretical interest to descriptive accounts in linguistics, as well as to psycholinguistic theories of sentence processing, to understand how different DAs are signaled prosodically.
### 2.4 Previous Studies of Prosody and Discourse
The main context in which prosody has been explored specifically for the purpose of dialog processing is in the area of discourse segmentation—both at the utterance level and at higher levels such as the organization of utterances into turns and topics. The segmentation studies span both descriptive and computational fields, and describe or attempt to detect utterance and topic boundaries using various acoustic-prosodic features, including pitch range, intonational contour, declination patterns, utterance duration, pre-boundary lengthening phenomena, pause patterns, speaking rate, and energy patterns. There has been increasing work in studying spontaneous speech, in both human-human and human-machine dialog. In most cases the features cuing the segments are coded by hand, but could potentially be estimated by automatic means for speech applications ( ?, ?, ?, ?, ?, ?, ?, ?, ?, ?). Although much of the work on prosody and segmentation has been descriptive, some recent studies have developed classifiers and tested performance using a fully automatic detection paradigm. For example, ? (?) found that features derived from a pitch tracker (F0, but also voicing and energy information) provide cues to intonational phrase boundaries; such a system could be used as a front end for audio browsing and playback. Similarly, in experiments on subsets of the German Verbmobil spontaneous speech corpus, prosodic features (including features reflecting duration, pause, F0, and energy) were found to improve segmentation performance (into DAs) over that given by a language model alone ( ?, ?). The Verbmobil work was in the context of an overall system for automatically classifying DAs, but the prosodic features were used only at the segmentation stage.
A second line of relevant previous work includes studies on the automatic detection of pitch accents, phrase accents, and boundary tones for speech technology. It has become increasingly clear that a transcribed word sequence does not provide enough information for speech understanding, since the same sequence of words can have different meanings depending, in part, on prosody. The location and type of accents and boundary tones can provide important cues for tasks such as lexical or syntactic disambiguation, and can be used to rescore word hypotheses and reduce syntactic or semantic search complexity ( ?, ?, ?, ?, ?). These and many related studies model F0, energy, and duration patterns to detect and classify accents and boundary tones; information on the location and type of prosodic events can then be used to assign or constrain meaning, typically at the level of the utterance. Such information is relevant to dialog processing, since the locations of major phrase boundaries delimit utterance units, and since tonal information can specify pragmatic meaning in certain contexts (e.g., a rising final boundary tone suggests questions). First developed for formal speech, such approaches have also been applied to spontaneous human-computer dialog, where the modeling problem becomes more difficult as a result of less constrained speech styles.
Beyond the detection of accents, boundary tones, and discourse-relevant segment boundaries, there has been only limited investigation into automatic processing specifically to identify DAs in conversational speech. In one approach, Taylor et al. (?, ?) used hidden Markov models (HMMs) to model accents and boundary tones in different conversational “moves” in the Maptask corpus ( ?), with the aim of applying move-specific language models to improve speech recognition. The event recognizer used “tilt” parameters ( ?), or F0, amplitude, duration, and a feature capturing the shape (rise, fall, or combination). As reported in many other studies of accent detection, performance degraded sharply from speaker-dependent formal styles to speaker-independent spontaneous speech (e.g., ?). The automatic detection of moves was thus limited by somewhat low accent detection accuracy (below 40%); however, overall results suggested that intonation can be a good predictor of move type.
In another study, ? (?) aimed to automatically identify utterances in human-machine dialog likely to contain emotional content such as exclamations of puzzlement, self-talk, or other types of paralinguistic information that the system would not be able to process. The approach involved clustering utterances based on vector-quantized F0 patterns and overall regression fits on the contours. Patterns deviating from a typically relatively flat overall slope were found to be likely to contain such paralinguistic content.
Finally, researchers on the Verbmobil project ( ?, ?), following ideas of ? (?), addressed an interesting case of ambiguity in human-machine interaction in the context of a train-scheduling system. Apparently, subjects often interrupt the announcement of train schedules to repeat a specific departure or arrival time. The repeat can serve one of three functional roles: confirmation of understanding, questioning of the time, or feedback that the user is still listening. The tendency of users to interrupt in this manner is even more pronounced when talking to an automatic system with synthesized speech output, since the synthesis can often be difficult to comprehend. To aid in automatically identifying responses, Gaussian classifiers were trained on F0 features similar to those mentioned in earlier work ( ?, ?), including the slope of the regression line of the whole contour and of the final portion, as well as utterance onset- and offset-related values. Similarly, ? (?) used F0 information to distinguish user queries from acknowledgments in a direction-giving system. To this end, the shape of pitch contours was classified either by a hand-written rule system, or a trained neural network.
### 2.5 Current Study
For the present work, we were interested in automatic methods that could be applied to spontaneous human-human dialog, which is notoriously more variable than read speech or most forms of human-computer dialog ( ?, ?, ?). We also wanted to cover the full set of dialog act labels observed, and thus needed to be able to define the extraction and computation of all proposed features for all utterances in the data. We took an exploratory approach, including a large set of features from the different categories of prosodic features used in the work on boundary and discourse described earlier. However, our constraints were somewhat different than in previous studies.
One important difference is that because we were interested in using prosodic features in combination with a language model in speech recognition, our features were designed to not rely on any word information; as explained later, this feature independence allows a probabilistic combination of prosodic and word-based models. A second major difference between our approach and work based on hand-labeled prosodic annotations is that our features needed to be automatically extractable from the signal. This constraint was practical rather than theoretical: it is currently not feasible to automatically detect abstract events such as accents and phrase boundaries reliably in spontaneous human-human dialog with variable channel quality (such as in telephone speech). Nevertheless, it is also the case that we do not yet fully understand how abstract categories characterize DAs in natural speech styles, and that an understanding could be augmented by information about correlations between DAs and other feature types. For example, even for DAs with presumed canonical boundary tone indicators (such as the rising intonation typical of questions), other features may additionally characterize the DA. For instance, descriptive analyses of Dutch question intonation have found that in addition to a final F0 rise, certain interrogatives differ from declaratives in features located elsewhere, such as in onset F0 and in overall pitch range ( ?, ?). Thus, we focussed on global and rather simple features, and assumed no landmarks in our utterances other than the start and end times.
Our investigation began as part of a larger project ( ?, ?, ?) on DA classification in human-human telephone conversations, using three knowledge sources: (1) a dialog grammar (a statistical model of the sequencing of DAs in a conversation), (2) DA-specific language models (statistical models of the word sequences associated with particular types of DAs), and (3) DA-specific prosodic models. Results revealed that the modeling was driven largely by DA priors (represented as unigram frequencies in the dialog grammar) because of an extreme skew in the distribution of DAs in the corpus—nearly 70% of the utterances in the corpus studied were either statements (declaratives) or brief backchannels (such as “uh-huh”). Because of the skew, it was difficult to assess the potential contribution of features of the DAs themselves, including the prosodic features. Thus, to better investigate whether prosody can contribute to DA classification in natural dialog, for this paper we eliminate additional knowledge sources that could confound our results. Analyses are conducted in a domain of uniform priors (all DAs are made equally likely). We also exclude contextual information from the dialog grammar (such as the DA of the previous utterance). In this way, we hope to gain a better understanding of the inherent prosodic properties of different DAs, which can in turn help in the building of better integrated models for natural speech corpora in general.
Our approach builds on a methodology previously developed for a different task involving conversational speech ( ?). The method is based on constructing a large database of automatically extracted acoustic-prosodic features. In training, decision tree classifiers are inferred from the features; the trees are then applied to unseen data to evaluate performance and to study feature usage.
The analyses examine decision tree performance in four DA-classification tasks. We begin with a task involving multiway classification of the DAs in our corpus. We then examine three binary classification tasks found to be problematic for word-based classification: Question detection, Agreement detection, and the detection of Incomplete Utterances. For each task, we train classifiers using various subsets of features to gain an understanding of the relative importance of different feature types. In addition, we integrate tree models with DA-specific language models to explore the role of prosody when word information is also available, from either a transcript or a speech recognizer.
## 3 Method
### 3.1 Speech Data
Our data were taken from the Switchboard corpus of human-human telephone conversations on various topics ( ?). The original release of this corpus contains roughly three million words from more than 2430 different conversations, each roughly 10 minutes in duration. The corpus was collected at Texas Instruments and is distributed by the Linguistics Data Consortium (LDC). A set of roughly 500 speakers representing all major dialects of American English participated in the task in exchange for a per-call remuneration. Speakers could participate as often as they desired; many speakers participated multiple times. Speakers were aware that their speech was being recorded, but were informed only generally that TI speech researchers were interested in the conversations. Speakers registered by choosing topics of interest (e.g., recycling, sports) from a predetermined set, and by indicating times that they would be available. They were automatically connected to another caller by a “robot operator” based on matching of registrants to topics and available times. An advantage of this procedure is the absence of experimenter bias. Conversations were therefore between strangers; however, transcribers rated the majority of conversations as sounding highly “natural”. There were some clear advantages to using this corpus for our work, including its size, the availability of transcriptions, and sentence-level segmentations. But most important, it was one of the only large English conversational-speech corpora available at the time, for which we could obtain N-best word recognition output from a state-of-the-art recognition system.
### 3.2 Dialog Act Labeling
#### 3.2.1 Labeling system
We developed a DA labeling system for Switchboard, taking as a starting point the DAMSL system ( ?) of DA labeling for task-oriented dialog. We adapted the DAMSL system to allow better coverage for Switchboard, and also to create labels that provide more information about the lexical and syntactic realization of DAs. Certain classes in DAMSL were never used, and conversely it was necessary to expand some of the DAMSL classes to provide a variety of labels. The adapted system, “SWBD-DAMSL”, is described in detail in ? (?).
SWBD-DAMSL defines approximately 60 unique tags, many of which represent orthogonal information about an utterance and hence can be combined. The labelers made use of 220 of these combined tags, which we clustered for our larger project into 42 classes ( ?). To simplify analyses, the 42 classes were further grouped into seven disjoint main classes, consisting of the frequently occurring classes plus an “Other” class containing DAs each occurring less than 2% of the time. The groups are shown in Table 1. The full set of DAs is listed in Appendix A, along with actual frequencies. The full list is useful for getting a feel for the heterogeneity of the “Other” class. Table 2 shows three typical exchanges found in the corpus, along with the kinds of annotations we had at our disposal.
For the Statement classes, independent analyses showed that the two SWBD-DAMSL types of Statements, Descriptions and Opinions, were similar in their lexical and their prosodic features, although they did show some differences in their distribution in the discourse, which warrants their continued distinction in the labeling system. Since, as explained in the Introduction, we do not use dialog grammar information in this work, there is no reason not to group the two types together for analysis. For the Question category we grouped together the main question types described by ? (?, ?), namely, Declarative Questions, Yes-No Questions, and Wh-Questions.
#### 3.2.2 Labeling procedure
Since there was a large set of data to label, and limited time and labor resources, we decided to have our main set of DA labels produced based on the text transcripts alone. Llabelers were given the transcriptions of the full conversations, and thus could use contextual information, as well as cues from standard punctuation (e.g., question marks), but did not listen to the soundfiles. A similar approach was used for the same reason in the work of ? (?). We were aware, however, that labeling without listening is not without problems. One concern is that certain DAs are inherently ambiguous from transcripts alone. A commonly noted example is the distinction between simple Backchannels, which acknowledge a contribution (e.g., “uh-huh”) and explicit Agreements (e.g., “that’s exactly it”). There is considerable lexical overlap between these two DAs, with emphatic intonation conveying an Agreement (e.g., “right” versus “right!”). Emphasis of this sort was not marked by punctuation in the transcriptions, and Backchannels were nearly four times as likely in our corpus; thus, labelers when in doubt were instructed to mark an ambiguous case as a Backchannel. We therefore expected that some percentage of our Backchannels were actually Agreements. In addition to the known problem of Backchannel/Agreement ambiguities, we were concerned about other possible mislabelings. For example, rising intonation could reveal that an utterance is a Declarative Question rather than a Statement. Similarly, hesitant-sounding prosody could indicate an Incomplete Utterance (from the point of view of the speaker’s intention), even if the utterance is potentially complete based on words alone.
Such ambiguities are of particular concern for the analyses at hand, which seek to determine the role of prosody in DA classification. If some DAs are identifiable only when prosody is made available, then a subset of our original labels will not only be incorrect, they will also be biased toward the label cued by a language model. This will make it difficult to determine the degree to which prosodic cues can contribute to DA classification above and beyond the language model cues. We took two steps toward addressing these concerns within the limits of our available resources. First, we instructed our labelers to flag any utterances that they felt were ambiguous from text alone. In future work such utterances could be labeled after listening. Given that this was not possible yet for all of the labeled data, we chose to simply remove all flagged utterances for the present analyses.
Second, we conducted experiments to assess the loss incurred by labeling with transcripts only. We asked one of the most experienced of our original DA labelers<sup>1</sup><sup>1</sup>1We thank Traci Curl for reannotating the data and for helpful discussions. to reannotate utterances after listening to the soundfiles. So that the factor of listening would not be confounded with that of inter-labeler agreement, all conversations to be relabeled were taken from the set of conversations that she had labeled originally. In the interest of time, the relabeling was done with the original labels available. Instructions were to listen to all of the utterances, and take the time needed to make any changes in which she felt the original labels were inconsistent with what she heard. This approach is not necessarily equivalent to relabeling from scratch, since the labeler may be biased toward retaining previous labels. Nevertheless, it should reveal the types of DAs for which listening is most important. This was the goal of a first round (Round I) of relabeling, in which we did not give any information about which DAs to pay attention to. The rate of changes for the individual DA types, however, was assumed to be conservative here, since the labeler had to divide her attention over all DA types. Results are shown in the left column of Table 3.
Only 114 changes were made in Round I, for an overall rate of change of under 2%. Given that attention was divided over all DAs in this round, the most meaningful information from Round I is not the overall rate of changes, which is expected to be conservative, but rather the distribution of types of changes. The most prominent change made after listening was the conversion of Backchannels (b) to Agreements (aa). Details on the prosodic cues associated with this change are described elsewhere ( ?). As the table shows for top changes, this change accounted for 43, or 37.7%, of the 114 changes made; the next most frequent change (within the two different original Statement labels) accounted for less than 20% of the changes.<sup>2</sup><sup>2</sup>2In addition, many of the sd$``$sv changes were in fact an indirect result of b$``$aa changes for the following utterance. The salience of the b$``$aa changes is further seen after normalizing the number of changes by the DA priors. On this measure, b$``$aa changes occur for over 4% of original b labels. In contrast, the normalized rates for the second and third most frequent types of changes in Round I were 22/989 (2.22%) for sv$``$sd and 17/2147 (0.79%) for sd$``$sv. For all changes not involving either b or aa, the rate was only about 1%. A complete list of recall and precision rates by DA type (where labels after listening are used as reference labels, and labels from transcripts alone are used as hypothesized labels), can be found in Appendix B.
To address the issue of attention to changing the original labels, we ran a second round of relabeling (Round II). Since b$``$aa changes were clearly the most salient from Round I, we discussed these changes with the labeler, and then asked her to relabel additional conversations with attention to these changes. Thus, we expected her to focus relatively more attention on b$``$aa in Round II (although she was instructed also to label any other glaring changes). We viewed Round II as a way to obtain an upper bound on the DA-specific change rate, since b$``$aa changes were the most frequently occurring changes after listening, and since the labeler was biased toward focusing attention on these changes. For Round II, we used a completely separate set of data from Round I, to avoid confounding the relabeling procedure. The overall distribution of DAs was similar to that in the set used in Round I.
As shown in Table 3, the number of changes made in Round II was the same (by coincidence) as in Round I. However, since there were fewer total utterances in Round I, the rate of change relative to total DAs increased from Round I to Round II. In Round II, b$``$aa changes greatly increased from Round I, both relative to total DAs and relative to DA-specific priors. At the same time, other types of changes decreased from Round I to Round II.
The most important result from Round II is the rate of b$``$aa changes relative to the prior for the b class. This value was about 10%, and is a reasonable estimate of the upper bound on DA changes for any particular class from listening, since it is unlikely that listening would affect other DAs more than it did Backchannels, given both the predominance of b$``$aa changes in Round I, and the fact that the labeler was biased to attend to b$``$aa changes in Round II. These results suggest that at least 90% of the utterances in any of our originally labeled DA classes are likely to be marked with the same DA label after listening, and that for most other DAs this value should be considerably higher. Therefore, although our transcript-only labels contained some errors, based on the results of the relabeling experiments we felt that it was reasonable to use the transcript-only labels as estimates of after-listening labels.
#### 3.2.3 Interlabeler reliability
Interlabeler reliability on our main (transcript-only) set of annotations was assessed using the Kappa statistic ( ?, ?, ?), or the ratio of the proportion of times that raters agree (corrected for chance agreement) to the maximum proportion of times that the rates could agree (corrected for chance agreement). Kappa computed for the rating of the original 42 classes was 0.81, which is considered high for this type of task. Post hoc grouping of the ratings using the seven main classes just described yielded a Kappa of 0.85.
### 3.3 Training and Test Sets
We partitioned the available data into three subsets for training and testing. The three subsets were not only disjoint but also shared no speakers. The training set (TRN) contained 1794 conversation sides; its acoustic waveforms were used to train decision trees, while the corresponding transcripts served as training data for the statistical language models used in word-based DA classification. The held-out set (HLD) contained 436 conversation sides; it was used to test tree performance as well as DA classification based on true words. A much smaller development test set (DEV) consisting of 38 matched conversation sides (19 conversations) was used to perform experiments involving automatic word recognition, as well as corresponding experiments based on prosody and true words.<sup>3</sup><sup>3</sup>3The DEV set was so called because of its role in the WS97 projects that focused on word recognition. The TRN and HLD sets contained single, unmatched conversation sides, but since no discourse context was required for the studies reported here this was not a problem. The three corpus subsets with their statistics are summarized in Table 4.
### 3.4 Dialog Act Segmentation
In a fully automated system, DA classification presupposes the ability to also find the boundaries between utterances. In spite of extensive work on this problem in recent years, to our knowledge there are currently no systems that reliably perform utterance segmentation for spontaneous conversational speech when the true words are not known. For this work we did not want to confound the issue of DA classification with DA segmentation; thus, we used utterance boundaries marked by human labelers according to the LDC annotation guidelines described in ? (?). To keep results using different knowledge sources comparable, these DA boundaries were also made explicit for purposes of speech recognition and language modeling.<sup>4</sup><sup>4</sup>4Note that the very notion of utterances and utterance boundaries is a matter of debate and subject to research ( ?). We adopted a pragmatic approach by choosing a pre-existing segmentation for this rather large corpus.
The utterance boundaries were marked between words. To estimate the locations of the boundaries in the speech waveforms, a forced alignment of the acoustic training data was merged with the training transcriptions containing the utterance boundary annotations marked by the LDC. This yielded word and pause times of the training data with respect to the acoustic segmentations. By using these word times along with the linguistic segmentation marks, the start and end times for linguistic segments were found.
This technique was not perfect, however. One problem is that many of the words included in the linguistic transcription had been excised from the acoustic training data. Some speech segments were considered not useful for acoustic training and thus had been excluded deliberately. In addition, the alignment program was allowed to skip words at the beginning and end of an acoustic segment if there was insufficient acoustic evidence for the word. This caused misalignments in the context of highly reduced pronunciations or for low-energy speech, both of which are frequent in Switchboard. Errors in the boundary times for DAs crucially affect the prosodic analyses, since prosodic features are extracted assuming that the boundaries are reasonably correct. Incorrect estimates affect the accuracy of global features (e.g., DA duration) and may render local features meaningless (e.g., F0 measured at the supposed end of the utterance). Since features for DAs with known problematic end estimates would be misleading in the prosodic analyses, they were omitted from all of our TRN and HLD data. The time boundaries of the DEV test set, however, were carefully handmarked for other purposes, so we were able to use exact values for this test set. Overall, we were missing 30% of the utterances in the TRN and HLD sets because of problems with time boundaries; this figure was higher for particular utterance types, especially for short utterances such as backchannels, for which as much as 45% of the utterances were affected. Thus, the DEV set was mismatched with respect to the TRN and HLD sets in terms of the percentage of utterances affected by problematic segmentations.
### 3.5 Prosodic Features
The prosodic database included a variety of features that could be computed automatically without reference to word information. In particular, we attempted to have good coverage of features and feature extraction regions that were expected to play a role in the three focused analyses mentioned in the Introduction: detection of Questions, Agreements, and Incomplete Utterances. Based on the literature on question intonation ( ?, ?, ?), we expected Questions to show rising F0 at the end of the utterance, particularly for Declarative and Yes-No Questions. Thus, F0 should be a helpful cue for distinguishing Questions from other long DAs such as Statements. Many Incomplete Utterances give the impression of being cut off prematurely, so the prosodic behavior at the end of such an utterance may be similar to that of the middle of a normal utterance. Specifically, energy can be expected to be higher at the end of an abandoned utterance compared to energy at the end of a completed one. In addition, unlike most completed utterances, the F0 contour at the end of an Incomplete Utterance is neither rising nor falling. We expected Backchannels to differ from Agreements by the amount of effort used in speaking. Backchannels function to acknowledge another speaker’s contributions without taking the floor, whereas Agreements assert an opinion. We therefore expected Agreements to have higher energy, greater F0 movement, and a higher likelihood of accents and boundary tones than Backchannels.
#### 3.5.1 Duration features
Duration was expected to be a good cue for discriminating Statements and Questions from DAs functioning to manage the dialog (e.g., Backchannels), although this difference is also encoded to some extent in the language model. In addition to the duration of the utterance in seconds, we included features correlated with utterance duration, but based on frame counts conditioned on the value of other feature types, as shown in Table 5.
The duration-pause set of features computes duration, ignoring pause regions. Such features may be useful if pauses are unrelated to DA classification. (If pauses are relevant, however, this should be captured by the pause features described in the next section.) The F0-based count features reflect either the number of frames or recognized intonational events (accents or boundaries) based on F0 information (see F0 features, below). The first four of these features capture time in speaking by using knowledge about the presence and location of voiced frames, which may be more robust for our data than relying on pause locations from the alignments. The last two features are intended to capture the amount of information in the utterance, by counting accents and phrase boundaries. Duration-normalized versions of many of these features are included under their respective feature type in the following sections.
#### 3.5.2 Pause features
To address the possibility that hesitation could provide a cue to the type of DA, we included features intended to reflect the degree of pausing, as shown in Table 6. To obtain pause locations we used information available from forced alignments; however, this was only for convenience (the alignment information was included in our database for other purposes). In principle, pause locations can be detected by current recognizers with high accuracy without knowledge of the words. Pauses with durations below 100 milliseconds (10 frames) were excluded since they are more likely to reflect segmental information than hesitation. Features were normalized to remove the inherent correlation with utterance duration. The last feature provides a more global measure of pause behavior, including pauses during which the other speaker was talking. The measure counts only those speech frames occurring in regions of at least 1 second of continuous speaking. The window was run over the conversation (by channel), writing out a binary value for each frame; the feature was then computed based on the frames within a particular DA.
#### 3.5.3 F0 features
F0 features, shown in Table 7, included both raw values (obtained from ESPS/Waves+) and values from a linear regression (least-squares fit) to the frame-level F0 values.
To capture overall pitch range, mean F0 values were calculated over all voiced frames in an utterance. To normalize differences in F0 range over speakers, particularly across genders, utterance-level values were normalized with respect to the mean and standard deviation of F0 values measured over the whole conversation side. F0 difference values were normalized on a log scale. The standard deviation in F0 over an utterance was computed as a possible measure of expressiveness over the utterance. Minimum and maximum F0 values, calculated after median smoothing to eliminate spurious values, were also included for this purpose.<sup>5</sup><sup>5</sup>5A more linguistically motivated measure of the maximum F0 would be to take the F0 value at the RMS maximum of the sonorant portion of the nuclear-accented syllable in the phrase (e.g., ?). However, our less sophisticated measure of pitch range was used as an approximation because we did not have information about the location of accents or phrase boundaries available.
We included parallel measures that used only “good” F0 values, or values above a threshold (f0\_min) estimated as the bottom of a speaker’s natural F0 range. The f0\_min can be calculated in two ways. For both methods, a smoothed histogram of all the calculated F0 values for a conversation side is used to find the F0 mode. The true f0\_min comes from the minimum F0 value to the left of this mode. Because the histogram can be flat or not sufficiently smoothed, the algorithm could be fooled into choosing a value greater than the true minimum. A simpler way to estimate the f0\_min takes advantage of the fact that values below the minimum typically result from pitch halving. Thus, a good estimate of f0\_min is to take the point at 0.75 times the F0 value at the mode of the histogram. This measure closely approximates the true f0\_min, and is more robust for use with the Switchboard data.<sup>6</sup><sup>6</sup>6We thank David Talkin for suggesting this method. The percentage of “good” F0 values was also included to measure (inversely) the degree of creaky voice or vocal fry.
The rising/falling behavior of pitch contours is a good cue to their utterance type. We investigated several ways to measure this behavior. To measure overall slope, we calculated the gradient of a least-squares fit regression line for the F0 contour. While this gives an adequate measure for the overall gradient of the utterance, it is not always a good indicator of the type of rising/falling behavior in which we are most interested. Rises at the end can be swamped by the declination of the preceding part of the contour, and hence the overall gradient for a contour can be falling. We therefore marked two special regions at the end of the contour, corresponding to the last 200 milliseconds (end region) and the 200 milliseconds previous to that (penultimate region). For each of these regions we measured the mean F0 and gradient, and used the differences between these as features. The starting value in the regression line was also included as a potential cue to F0 register (the actual first value is prone to F0 measurement error).
In addition to these F0 features, we included intonational-event features, or features intended to capture local pitch accents and phrase boundaries. The event features were obtained using the event recognizer described in ? (?). The event detector uses an HMM approach to provide an intonational segmentation of an utterance, which gives the locations of pitch accents and boundary tones. When compared to human intonation transcriptions of Switchboard,<sup>7</sup><sup>7</sup>7As labeled by the team of students at Edinburgh; see Acknowledgments. this system correctly identifies 64.9% of events, but has a high false alarm rate, resulting in an accuracy of 31.7%.
#### 3.5.4 Energy features
We included two types of energy features, as shown in Table 8. The first set of features was computed based on standard RMS energy. Because our data were recorded from telephone handsets with various noise sources (background noise as well as channel noise), we also included a signal-to-noise ratio (SNR) feature to try to capture the energy from the speaker. SNR values were calculated using the SRI recognizer with a Switchboard-adapted front end ( ?, ?). Values were calculated over the entire conversation side, and those extracted from regions of speech were used to find a cumulative distribution function (CDF) for the conversation. The frame-level SNR values were then represented by their CDF value to normalize the SNR values across speakers and conversations.
#### 3.5.5 Speaking rate (enrate) features
We were also interested in overall speaking rate. However, we needed a measure that could be run directly on the signal, since our features could not rely on word information. For this purpose, we experimented with a signal processing measure, “enrate” ( ?), which estimates a syllable-like rate by looking at the energy in the speech signal after preprocessing. Studies comparing enrate values to values based on hand-transcribed syllable rates for Switchboard show a correlation of about .46 for the version of the software used in the present work.<sup>8</sup><sup>8</sup>8We thank Nelson Morgan, Eric Fosler-Lussier, and Nikki Mirghafori for allowing us to use the software and note that the measure has since been improved (mrate), with correlations increasing to about .67 as described in ? (?).
The measure can be run over the entire signal, but because it uses a large window, values are less meaningful if significant pause time is included in the window. We calculated frame-level values over a 2-second speech interval. The enrate value was calculated for a 25-millisecond frame window with a window step size of 200 milliseconds. Output values were calculated every 10 milliseconds to correspond to other measurements. We included pauses of less than 1 second and ignored speech regions of less than 1 second, where pause locations were determined as described earlier.
If the end of a speech segment was approaching, meaning that the 2-second window could not be filled, no values were written out. The enrate values corresponding to particular utterances were then extracted from the conversation-side values. This way, if utterances were adjacent, information from surrounding speech regions could be used to get enrate values for the beginnings and ends of utterances that normally would not fill the 2-second speech window. Features computed for use in tree-building are listed in Table 9.
#### 3.5.6 Gender features
As a way to check the effectiveness of our F0 normalizations we included the gender of the speaker. It is also possible that features could be used differently by men and women, even after appropriate normalization for pitch range differences. We also included the gender of the listener to check for a possible sociolinguistic interaction between the conversational dyad and the ways in which speakers employ different prosodic features.
### 3.6 Decision Tree Classifiers
For our prosodic classifiers, we used CART-style decision trees ( ?). Decision trees can be trained to perform classification using a combination of discrete and continuous features, and can be inspected to gain an understanding of the role of different features and feature combinations.
We downsampled our data (in both training and testing) to obtain an equal number of datapoints in each class. Although an inherent drawback is a loss of power in the analyses due to fewer datapoints, downsampling was warranted for two reasons. First, as noted earlier, the distribution of frequencies for our DA classes was severely skewed. Because decision trees split according to an entropy criterion, large differences in class size wash out any effect of the features themselves, causing the tree not to split. By downsampling to equal class priors we assure maximum sensitivity to the features. A second motivation for downsampling was that by training our classifiers on a uniform distribution of DAs, we facilitated integration with other knowledge sources (see section on Integration). After expanding the tree with questions, the training algorithm used a tenfold cross-validation procedure to avoid overfitting the training data. Leaf nodes were successively pruned if they failed to reduce the entropy in the cross-validation procedure.
We report tree performance using two metrics, accuracy and efficiency. Accuracy is the number of correct classifications divided by the total number of samples. Accuracy is based on hard decisions; the classification is that class with the highest posterior probability. Because we downsampled to equal class priors, the chance performance for any tree with N classes is 100/N%. For any particular accuracy level, there is a trade-off between recall and false alarms. In the real world there may well be different costs to a false positive versus a false negative in detecting a particular utterance type. In the absence of any model of how such costs would be assigned for our data, we report results assuming equal costs to these errors.
Efficiency measures the relative reduction in entropy between the prior class distribution and the posterior distribution predicted by the tree. Two trees may have the same classification accuracy, but the tree that more closely approximates the probability distributions of the data (even if there is no effect on decisions) has higher efficiency (lower entropy). Although accuracy and efficiency are typically correlated, the relationship between the measures is not strictly monotonic since efficiency looks at probability distributions and accuracy looks only at decisions.
### 3.7 Dialog Act Classification from Word Sequences
Two methods were used for classification of DAs from word information. For experiments using the correct words $`W`$, we needed to compute the likelihoods $`P(W|U)`$ for each DA or utterance type $`U`$, i.e., the probability with which $`U`$ generates the word sequence $`W`$. The predicted DA type would then be the one with maximum likelihood. To estimate these probabilities, we grouped the transcripts of the training corpus by DA type, and trained a standard trigram language model using backoff smoothing ( ?) for each DA. This was done for the original 42 DA categories, yielding 42 DA-specific language models. Next, for experiments involving a DA class $`C`$ comprising several of the original DAs $`U_1`$, $`U_2`$, …, $`U_n`$, we combined the DA likelihoods in a weighted manner:
$$P(W|C)=P(W|U_1)P(U_1|C)+\mathrm{}+P(W|U_n)P(U_n|C)$$
Here, $`P(U_1|C)`$, …, $`P(U_n|C)`$ are the relative frequencies of the various DAs within class $`C`$.
For experiments involving (necessarily imperfect) automatic word recognition, we were given only the acoustic information $`A`$. We therefore needed to compute acoustic likelihoods $`P(A|U)`$, i.e., the probability that utterance type $`U`$ generates the acoustic manifestation $`A`$. In principle, this can be accomplished by considering all possible word sequences $`W`$ that might have generated the acoustics $`A`$, and summing over them:
$$P(A|U)=\underset{W}{}P(A|W)P(W|U)$$
Here $`P(W|U)`$ is estimated by the same DA-specific language models as before, and $`P(A|W)`$ is the acoustic score of a speech recognizer, expressing how well the acoustic observations match the word sequence $`W`$. In practice, however, we could only consider a finite number of potential word hypotheses $`W`$; in our experiments we generated the 2500 most likely word sequences for each utterance, and carried out the above summation over only those sequences. The recognizer used was a state-of-the-art HTK large-vocabulary recognizer, which nevertheless had a word error rate of 41% on the test corpus.<sup>9</sup><sup>9</sup>9Note that the summation over multiple word hypotheses is preferable to the more straightforward approach of looking at only the one best hypothesis and treating it as the actual words for the purpose of DA classification.
### 3.8 Integration of Knowledge Sources
To use multiple knowledge sources for DA classification, i.e., prosodic information as well as other, word-based evidence, we combined tree probabilities $`P(U|F)`$ and word-based likelihoods $`P(W|U)`$ multiplicatively. This approach can be justified as follows. The likelihood-based classifier approach dictates choosing the DA with the highest likelihood based on both the prosodic features $`F`$ and the words $`W`$, $`P(F,W|U)`$. To make the computation tractable, we assumed, similar to ? (?), that the prosodic features are independent of the words once conditioned on the DA. We recognize, however, that this assumption is a simplification.<sup>10</sup><sup>10</sup>10Utterance length is one feature for which this independence assumption is clearly violated. Utterance length is represented by a prosodic feature (utterance duration) as well as implicitly in the DA-specific language models. ? (?) suggest a way to deal with this particular problem by conditioning the language models on utterance length. Our prosodic model averages over all examples of a particular DA; it is “blind” to any differences in prosodic features that correlate with word information. For example, statements about a favorite sports team use different words than statements about personal finance, and the two different types of statements tend to differ prosodically (e.g., in animation level as reflected by overall pitch range). In future work, such differences could potentially be captured by using more sophisticated models designed to represent semantic or topic information. For practical reasons, however, we consider our prosodic models independent of the words once conditioned on the DA, i.e.:
$`P(F,W|U)`$ $`=`$ $`P(W|U)P(F|W,U)`$
$``$ $`P(W|U)P(F|U)`$
$``$ $`P(W|U)P(U|F)`$
The last line is justified because, as noted earlier, we trained the prosodic trees on downsampled data or a uniform distribution of DA classes. According to Bayes’ Law, the required likelihood $`P(F|U)`$ equals $`P(U|F)P(F)/P(U)`$. The second factor, $`P(F)`$, is the same for all DA types $`U`$, and $`P(U)`$ is equalized by the downsampling procedure. Hence, the probability estimated by the tree, $`P(U|F)`$, is proportional to the likelihood $`P(F|U)`$. Overall, this justifies multiplying $`P(W|U)`$ and $`P(U|F)`$.<sup>11</sup><sup>11</sup>11In practice we needed to adjust the dynamic ranges of the two probability estimates by finding a suitable exponential weight $`\lambda `$, to make
$$P(F,W|U)P(W|U)P(F|U)^\lambda .$$
## 4 Results and Discussion
We first examine results of the prosodic model for a seven-way classification involving all DAs. We then look at results from a words-only analysis, to discover potential subtasks for which prosody could be particularly helpful. The words-only analysis reveals that even if correct words are available, certain DAs tend to be misclassified. We examine the potential role of prosody for three such subtasks: (1) the detection of Questions, (2) the detection of Agreements, and (3) the detection of Incomplete Utterances. In all analyses we seek to understand the relative importance of different features and feature types, as well as to determine whether integrating prosodic information with a language model can improve classification performance overall.
### 4.1 Seven-Way Classification
We applied the prosodic model first to a seven-way classification task for the full set of DAs: Statements, Questions, Incomplete Utterances, Backchannels, Agreements, Appreciations, and Other. Note that “Other” is a catch-all class representing numerous heterogeneous DAs that occurred infrequently in our data. Therefore we do not expect this class to have consistent features. As described in the Method section, data were downsampled to equal class sizes to avoid confounding results with information from prior frequencies of each class. Because there are seven classes, chance accuracy for this task is 100/7% or 14.3%. For simplicity, we assumed equal cost to all decision errors, i.e., to all possible confusions among the seven classes.
A tree built using the full database of features described earlier yields a classification accuracy of 41.15%. This gain in accuracy is highly significant by a binomial test, $`p<.0001`$. If we are interested in probability distributions rather than decisions, we can look at the efficiency of the tree, or the relative reduction in entropy over the prior distribution. By using the all-features prosodic tree for this seven-way classification, we reduce the number of bits necessary to describe the class of each datapoint by 16.8%.
The all-features tree is large (52 leaves), making it difficult to interpret the tree directly. In such cases we found it useful to summarize the overall contribution of different features by using a measure of “feature usage”, which is proportional to the number of times a feature was queried in classifying the datapoints. The measure thus accounts for the position of the feature in the tree: features used higher in the tree have greater usage values than those lower in the tree since there are more datapoints at the higher nodes. The measure is normalized to sum to 1.0 for each tree. Table 10 lists usage by feature type.
Table 10 indicates that when all features are available, a duration-related feature is used in more than half of the queries. Notably, gender features are not used at all; this supports the earlier hypothesis that, as long as features are appropriately normalized, it is reasonable to create gender-independent prosodic models for these data. A summary of individual feature usage, as shown in Table 11, reveals that the raw duration feature (ling\_dur)—which is a “free” measure in our work since we assumed locations of utterance boundaries—accounted for only 14% of the queries in the tree; the remaining queries of the 55% accounted for by duration features were those associated with the computation of F0- and pause-related information. Thus, the power of duration for the seven-way classification comes largely from measures involving computation of other prosodic features. The most-queried feature, regr\_num\_frames (the number of frames used in computing the F0 regression line), may be better than other duration measures at capturing actual speech portions (as opposed to silence or nonspeech sounds), and may be better than other F0-constrained duration measures (e.g., f0\_num\_good\_utt) because of a more robust smoothing algorithm. We can also note that the high overall rate of F0 features given in Table 11 represents a summation over many different individual features.
Since we were also interested in feature importance, individual trees were built using the leave-one-out method, in which the feature list is systematically modified and a new tree is built for each subset of allowable features. It was not feasible to leave out individual features because of the large set of features used; we therefore left out groups of features corresponding to the feature types as defined in the Method section. We also included a matched set of “leave-one-in” trees for each of the feature types (i.e., trees for which all other feature types were removed) and a single leave-two-in tree, built post hoc, which made available the two feature types with highest accuracy from the leave-one-in analyses. Note that the defined feature lists specify the features available for use in building a particular prosodic model; whether or not features are actually used is determined by the tree learning algorithm and depends on the data. Figure 1 shows results for the set of leave-one-out and leave-one-in trees, with the all-features tree provided for comparison. The upper graph indicates accuracy values; the lower graph shows efficiency values. Each bar indicates a separate tree.
We first tested whether there was any significant loss in leaving out a feature type, by doing pairwise comparisons between the all-features tree and each of the leave-one-out trees.<sup>12</sup><sup>12</sup>12To test whether one tree (A) was significantly better than another (B), we counted the number of test instances on which A and B differed, and on how many instances A was correct but B was not; we then applied a Sign test to these counts. Although trees with more features to choose from typically perform better than those with fewer features, additional features can hurt performance. The greedy tree-growing algorithm does not look ahead to determine the overall best feature set, but rather seeks to maximize entropy reduction locally at each split. This limitation of decision trees is another motivation for conducting the leave-one-out analyses. Since we cannot predict the direction of change for different feature sets, comparisons on tree results were conducted using two-tailed tests.
Results showed that the only two feature types whose removal caused a significant reduction in accuracy were duration ($`p<0.0001`$) and enrate ($`p<0.05`$). The enrate-only tree, however, yields accuracies on par with other feature types whose removal did not affect overall performance; this suggests that the contribution of enrate in the overall tree may be through interactions with other features. All of the leave-one-in trees were significantly less accurate than the all-features tree. Although the tree using only duration achieved an accuracy close to that of the all-features tree, it was still significantly less accurate by a Sign test ($`p<0.01`$). Adding F0 features (the next-best feature set in the leave-one-in trees) did not significantly improve accuracy over the duration-only tree alone, suggesting that for this task the two feature types are highly correlated. Nevertheless, for each of the leave-one-in trees, all feature types except gender yielded accuracies significantly above chance by a binomial test ($`p<.0001`$ for the first five trees). The gender-only tree was slightly better than chance by either a one- or a two-tailed test.<sup>13</sup><sup>13</sup>13It is not clear here whether a one- or two-tailed test is more appropriate. Trees typically should not do worse than chance; however, because they minimize entropy and not accuracy, the accuracy can fall slightly below chance. However, this was most likely due to a difference in gender representation across classes.
Taken together, these results suggest that there is considerable redundancy in the features for DA classification, since removing one feature type at a time (other than duration) makes little difference to accuracy. Results also suggest, however, that features are not perfectly correlated; there must be considerable interaction among features in classifying DAs, because trees using only individual feature types are significantly less accurate than the all-features tree.
Finally, duration is clearly of primary importance to this classification. This is not surprising, as the task involves a seven-way classification including longer utterances (such as Statements) and very brief ones (such as Backchannels like “uh-huh”). Two questions of further interest regarding duration, however, are (1) will a prosody model that uses mostly duration add anything to a language model (in which duration is implicitly encoded), and (2) is duration useful for other tasks involving classification of DAs similar in length? Both questions are addressed in the following analyses.
As just discussed, the all-features tree (as well as others including only subsets of feature types) provides significant information for the seven-way classification task. Thus, if one were to use only prosodic information (no words or context), this is the level of performance resulting for the case of equal class frequencies. To explore whether the prosodic information could be of use when lexical information is also available, we integrated the tree probabilities with likelihoods from our DA-specific trigram language models built from the same data. For simplicity, integration results are reported only for the all-features tree in this and all further analyses, although as noted earlier this is not guaranteed to be the optimal tree.
Since our trees were trained with uniform class priors, we combined tree probabilities $`P(U|F)`$ with the word-based likelihoods $`P(W|U)`$ multiplicatively, as described in the Integration section.<sup>14</sup><sup>14</sup>14The relative weight assigned to the prosodic and the word likelihoods was optimized on the test set due to lack of an additional held-out data set. Although in principle this could bias results, we found empirically that similar performance is obtained using a range of weighting values; this is not surprising since only a single parameter is estimated.
The integration was performed separately for each of our two test sets (HLD and DEV), and within the DEV set for both transcribed and recognized words. Results are shown in Table 12. Classification performance is shown for each of the individual classifiers, as well as for the optimized combined classifier.
As shown, for all three analyses, adding information from the tree to the word-based model improved classification accuracy. Although the gain appears modest in absolute terms, for the HLD test set it was highly significant by a Sign test,<sup>15</sup><sup>15</sup>15One-tailed, because model integration assures no loss in accuracy. $`p<.001`$. For the smaller DEV test set, the improvements did not reach significance; however, the pattern of results suggests that this is likely to be due to the small sample size. It is also the case that the tree model does not perform as well for the DEV as the HLD set. This is not attributable to small sample size, but rather to a mismatch between the DEV set and the training data involving how data were segmented, as noted in the Method section. The mismatch in particular affects duration features, which were important in these analyses as discussed earlier. Nevertheless, word-model results are lower for N-best than for true words in the DEV data, while by definition the tree results stay the same. We see that accordingly, integration provides a larger win for the recognized than for the true words. Thus, we would expect that results for recognized words for the HLD set (if they could be obtained) should show an even larger win than the win observed for the true words in that set.
These results provide an answer to one of the questions posed earlier: does prosody provide an advantage over words if the prosody model uses mainly duration? The results indicate that the answer is yes. Although the number of words in an utterance is highly correlated with duration, and word counts are represented implicitly by the probability of the end-of-utterance marker in a language model, a duration-based tree model still provides added benefit over words alone. One reason may be that duration (reflected by the various features we included) is simply a better predictor of DA than is word count. Another independent possibility is that the benefit from the tree model comes from its ability to threshold feature values directly and iteratively.
### 4.2 Dialog Act Confusions Based on Word Information
Next we explored additional tasks for which prosody could aid DA classification. Since our trees allow N-ary classification, the logical search space of possible tasks was too large to explore systematically. We therefore looked to the language model to guide us in identifying particular tasks of interest. Specifically, we were interested in DAs that tended to be misclassified even given knowledge of the true words. We examined the pattern of confusions made when our seven DAs were classified using the language model alone. Results are shown in Figure 2. Each subplot represents the data for one actual DA.<sup>16</sup><sup>16</sup>16Because of the heterogeneous makeup of the “Other” DA class, we were not per se interested in its pattern of confusions, and hence the graph for that data is not shown. Bars reflect the normalized rate at which the actual DA was classified as each of the seven possible DAs, in each of the three test conditions (HLD, DEV/true, and DEV/N-best).
As shown, classification is excellent for the Statement class, with few misclassifications even when only the recognized words are used.<sup>17</sup><sup>17</sup>17The high classification rate for Statements by word information was a prime motivation for downsampling our data in order to examine the inherent contribution of prosody, since as noted in the Method section, Statements make up most of the data in this corpus. For the remaining DAs, however, misclassifications occur at considerable rates.<sup>18</sup><sup>18</sup>18Exact classification accuracy values for the various DAs shown in Figure 2 are provided in the text as needed for the subtasks examined, i.e. under “words” in Tables 15, 17, and 18. Classification of Questions is a case in point: even with true words, Questions are often misclassified as Statements (but not vice versa), and this pattern is exaggerated when testing on recognized as opposed to true words. The asymmetry is partially attributable to the presence of declarative Questions. The effect associated with recognized words appears to reflect a high rate of missed initial “do” in our recognition output, as discovered in independent error analyses ( ?). For both Statements and Questions, however, there is little misclassification involving the remaining classes. This probably reflects the length distinction as well as the fact that most of the propositional content in our corpus occurred in Statements and Questions, whereas other DAs generally served to manage the communication—a distinction likely to be reflected in the words. Thus, our first subtask was to examine the role of prosody in the classification of Statements and Questions. A second problem visible in Figure 2 is the detection of Incomplete Utterances. Even with true words, classification of these DAs is at only 75.0% accuracy. Knowing whether or not a DA is complete would be particularly useful for both language modeling and understanding. Since the misclassifications are distributed over the set of DAs, and since logically any DA can have an incomplete counterpart, our second subtask was to classify a DA as either incomplete or not-incomplete (all other DAs). A third notable pattern of confusions involves Backchannels and explicit Agreements. This was an expected difficult discrimination as discussed earlier, since the two classes share words such as “yeah” and “right”. In this case, the confusions are considerable in both directions.
### 4.3 Subtask 1: Detection of Questions
As illustrated in the previous section, Questions in our corpus were often misclassified as Statements based on words alone. Based on the literature, we hypothesized that prosodic features, particularly those capturing the final F0 rise typical of some Question types in English, could play a role in reducing the rate of misclassifications. To investigate the hypothesis, we built a series of classifiers using only Question and Statement data. We first examined results for an all-features tree, shown in Figure 3. Each node displays the name of the majority class, as well as the posterior probability of the classes (in the class order indicated in the top node). Branches are labeled with the name of the feature and threshold value determining the split. The tree yields an accuracy of 74.21%, which is significantly above the chance level of 50% by a binomial test, $`p<.0001`$; the tree reduces the number of bits necessary to describe the class of each datapoint by 20.9%.
#### 4.3.1 Feature importance
The feature usage of the tree is summarized in Table 13. As predicted, F0 features help differentiate Questions from Statements, and in the expected direction (Questions have higher F0 means and higher end gradients than Statements). What was not obvious at the outset is the extent to which other features also cue this distinction. In the all-features tree, F0 features comprise only about 28% of the total queries. Two other features, regr\_dur and cont\_speech\_frames, are each queried more often than the F0 features together. Questions are shorter in duration (from starting to ending voiced frame) than Statements. They also have a lower percentage of frames in continuous speech regions than Statements. Further inspection suggests that the pause feature in this case (and also most likely for the seven-way classification discussed earlier) indirectly captures information about turn boundaries surrounding the DA of interest. Since our speakers were recorded on different channels, the end of one speaker’s turn is often associated with the onset of a long pause (during which the other speaker is talking). Furthermore, long pauses reduce the frame count for the continuous-speech-frames feature enrate measure because of the windowing described earlier. Therefore, this measure reflects the timing of continuous speech spurts across speakers, and is thus different in nature from the other pause features that look only inside an utterance.
To further examine the role of features, we built additional trees using partial feature sets. Results are summarized in Figure 4. As suggested by the leave-one-out trees, there is no significant effect on accuracy when any one of the feature types is removed. Although we predicted that Questions should differ from Statements mainly by intonation, results indicate that a tree with no F0 features achieves the same accuracy as a tree with all features for the present task. Removal of all pause features, which resulted in the largest drop in accuracy, yields a tree with an accuracy of 73.43%, which is not significantly different from that of the all-features tree ($`p=.2111`$, n.s.). Thus, if any feature type is removed, other feature types compensate to provide roughly the same overall accuracy. However, it is not the case that the main features used are perfectly correlated, with one substituting for another that has been removed. Inspection of the leave-one-out tree reveals that upon removal of a feature type, new features (features, and feature types, that never appeared in the all-features tree) are used. Thus, there is a high degree of redundancy in the features that differentiate Questions and Statements, but the relationship among these features and the allowable feature sets for tree building is complex.
Inspection of the leave-one-in tree results in Figure 4 indicates, not surprisingly, that the feature types most useful in the all-features analyses (duration and pause) yield the highest accuracies for the leave-one-in analyses (all of which are significantly above chance, $`p<.0001`$). It is interesting, however, that enrate, which was used only minimally in the all-features tree, allows classification at 68.09%, which is better than that of the F0-only tree. Furthermore, the enrate-only classifier is a mere shrub: as shown in Figure 5, it splits only once, on an unnormalized feature that expresses simply the variability in enrate over the utterance. As noted in the Method section, enrate is expected to correlate with speaking rate, although for this work we were not able to investigate the nature of this relationship. However, the result has interesting potential implications. Theoretically, it suggests that absolute speaking rate may be less important for DA classification than variation in speaking rate over an utterance; a theory of conversation should be able to account for the lower variability in questions than in statements. For applications, results suggest that the inexpensive enrate measure could be used alone to help distinguish these two types of DAs in a system in which other feature types are not available.
We ran one further analysis on question classification. The aim was to determine the extent to which our grouping of different kinds of questions into one class affected the features used in question classification. As described in the Method section, our Question class included Yes-No Questions, Wh-questions, and Declarative Questions. These different types of questions are expected to differ in their intonational characteristics ( ?, ?, ?, ?). Yes-No Questions and Declarative Questions typically involve a final F0 rise; this is particularly true for Declarative Questions whose function is not conveyed syntactically. Wh-Questions, on the other hand, often fall in F0, as do Statements.
We broke down our Question class into the originally coded Yes-No Questions, Wh-Questions, and Declarative Questions, and ran a four-way classification along with Statements. The resulting all-features tree is shown in Figure 6, and a summary of the feature usage is provided in Table 14.
The tree achieves an accuracy of 47.15%, a highly significant increase over chance accuracy (25%) by a binomial test, $`p<.0001`$. Unlike the case for the grouped Question class, the most queried feature type is now F0. Inspection of the tree reveals that the pattern of results is consistent with the literature on question intonation. Final rises (end\_grad, norm\_f0\_diff, and utt\_grad) are associated with Yes-No and Declarative Questions, but not with Wh-Questions. Wh-Questions show a higher average F0 (f0\_mean\_zcv) than Statements.
To further assess feature importance, we again built trees after selectively removing feature types. Results are shown in Figure 7.
In contrast to Figure 4, in which accuracy was unchanged by removal of any single feature type, the data in Figure 7 show a sharp reduction in accuracy when F0 features are removed. This result is highly significant by a Sign test ($`p<.001`$, two-tailed) despite the small amount of data in the analyses, resulting from downsampling to the size of the least frequent question subclass. For all other feature types, there was no significant reduction in accuracy when the feature type was removed. Thus, F0 plays an important role in question detection, but because different kinds of questions are signaled in different ways intonationally, combining questions into a single class as in the earlier analysis smoothes over some of the distinctions. In particular, the grouping tends to conceal the features associated with the final F0 rise (probably because the rise is averaged in with final falls).
#### 4.3.2 Integration with language model
To answer the question of whether prosody can aid Question classification when word information is also available, tree probabilities were combined with likelihoods from our DA-specific trigram language models, using an optimal weighting factor. Results were computed for the two test sets (HLD and DEV) and within the DEV set for both transcribed and recognized words. The outcome is shown in Table 15.
The prosodic tree model yielded accuracies significantly better than chance for both test sets ($`p<.0001`$). The tree alone was also more accurate than the recognized words alone for this task. Integration yielded consistent improvement over the words alone. The larger HLD set showed a highly significant gain in accuracy for the combined model over the words-only model, $`p<.001`$ by a Sign test. Significance tests were not meaningful for the DEV set because of a lack of power given the small sample size; however, the pattern of results for the two sets is similar (the spread is greatest for the recognized words) and therefore suggestive.
### 4.4 Subtask 2: Detection of Incomplete Utterances
A second problem area in the words-only analyses was the classification of Incomplete Utterances. Utterances labeled as incomplete in our work included three different main phenomena:<sup>19</sup><sup>19</sup>19In addition, the class included a variety of utterance types deemed “uninterpretable” because of premature cut-off.
| Turn exits: | | (A) | We have young children. |
| --- | --- | --- | --- |
| | $``$ | (A) | So … |
| | | (B) | Yeah, that’s tough then. |
| Other-interruptions: | $``$ | (A) | We eventually — |
| | | (B) | Well you’ve got to start somewhere. |
| Self-interruptions: | $``$ | (A) | And they were definitely — |
| (repairs) | | (A) | At halftime they were up by four. |
Although the three cases represent different phenomena, they are similar in that in each case the utterance could have been completed (and coded as the relevant type) but was not. An all-features tree built for the classification of Incomplete Utterances and all other classes combined (Non-Incomplete) yielded an accuracy of 72.16% on the HLD test set, a highly significant improvement over chance, $`p<.0001`$.
#### 4.4.1 Feature analyses
The all-features tree is complex and thus not shown, but feature usage by feature type is summarized in Table 16.
As indicated, the most-queried feature for this analysis is duration. Not surprisingly, Incomplete Utterances are shorter overall than complete ones; certainly they are by definition shorter than their completed counterparts. However, duration cannot completely differentiate Incomplete from Non-Incomplete utterances, because inherently short DAs (e.g., Backchannels, Agreements) are also present in the data. For these cases, other features such as energy and enrate play a role.
Results for trees run after features were selectively left out are shown in Figure 8. Removal of duration features resulted in a significant loss in accuracy, to 68.63%, $`p<.0001`$. Removal of any of the other feature types, however, did not significantly affect performance. Furthermore, a tree built using only duration features yielded an accuracy of 71.28%, which was not significantly less accurate than the all-features tree. These results clearly indicate that duration features are primary for this task. Nevertheless, good accuracy could be achieved using other feature types alone; for all trees except the gender-only tree, accuracy was significantly above chance, $`p<.0001`$. Particularly noteworthy is the energy-only tree, which achieved an accuracy of 68.97%. Typically, utterances fall to a low energy value when close to completion. However, when speakers stop mid-stream, this fall has not yet occurred, and thus the energy stays unusually high. Inspection of the energy-only tree revealed that over 75% of the queries involved SNR rather than RMS features, suggesting that at least for telephone speech, it is crucial to use a feature that can capture the energy from the speaker over the noise floor.
#### 4.4.2 Integration with language model
We again integrated the all-features tree with a DA-specific language model to determine whether prosody could aid classification with word information present. Results are presented in Table 17. Like the earlier analyses, integration improves performance over the words-only model for all three test cases. Unlike earlier analyses, however, the relative improvement when true words are used is minimal, and the effect is not significant for either the HLD/true-words or the DEV/true-words data. However, the relative improvement for the DEV/N-best case is much larger. The effect is just below the significance threshold for this small dataset ($`p=.067`$), but would be expected, based on the pattern of results in the previous analyses, to easily reach significance for a set of data the size of the HLD set.
Results suggest that for this task, prosody is an important knowledge source when word recognition is not perfect. When true words are available, however, it is not clear whether adding prosody aids performance. One factor underlying this pattern of results may be that the tree information is already accounted for in the language model. Consistent with this possibility is the fact that the tree uses mainly duration features, which are indirectly represented in the language model by the end-of-sentence marker. On the other hand, typically the word lengths of true and N-best lists are similar, and our results differ for the two cases, so it is unlikely that this could be the only factor.
Another possibility is that when true words are available, certain canonical Incomplete Utterances can be detected with excellent accuracy. A likely candidate here is the turn exit. Turn exits typically contain one or two words from a small inventory of possibilities—mainly coordinating conjunctions (“and”, “but”) and fillers (“uh”, “um”). Similarly, because Switchboard consists mainly of first-person narratives, a typical self-interrupted utterance in this corpus is a noncommittal false start such as “I—” or “I think—”. Both the turn exits and the noncommittal false starts are lexically cued and are thus likely to be well captured by a language model that has true words available.
A third possible reason for the lack of improvement over true words is that the prosodic model loses sensitivity because it averages over phenomena with different characteristics. False starts in our data typically involved a sudden cut-off, whereas for turn exits the preceding speech was often drawn out as in a hesitation. As a preliminary means of investigating this possibility, we built a tree for Incomplete Utterances only, but breaking down the class into those ending at turn boundaries (mainly turn exits and interrupted utterances) versus those ending within a speaker’s turn (mainly false starts). The resulting tree achieved high accuracy (81.17%) and revealed that the two subclasses differed on several features. For example, false starts were longer in duration, higher in energy, and had faster speaking rates than the turn exit/other-interrupted class. Thus, as we also saw for the case of Question detection, a prosodic model for Incomplete Utterances is probably best built on data that have been broken down to isolate subsets of phenomena whose prosodic features pattern differently.
### 4.5 Subtask 3: Detection of Agreements
Our final subtask examined whether prosody could aid in the detection of explicit Agreements (e.g., “that’s exactly right”). As shown earlier, Agreements were most often misclassified as Backchannels (e.g., “uh-huh”, “yeah”). Thus, our experiments focused on the distinction by including only these two DAs in the trees. An all-features tree for this task classified the data with an accuracy of 68.77% (significantly above chance by a binomial test, $`p<.0001`$) and with an efficiency of 12.21%.
#### 4.5.1 Feature analyses
The all-features tree is shown in Figure 9. It uses duration, pause, and energy features. From inspection we see that Agreements are consistently longer in duration and have higher energy (as measured by mean SNR) than Backchannels. The pause feature in this case may play a role similar to that discussed for the question classification task. Although Agreements and Backchannels were about equally likely to occur turn-finally, Backchannels were more than three times as likely as Agreements to be the only DA in a turn. Thus, Backchannels were more often surrounded by nonspeech regions (pauses during which the other speaker was typically talking), causing the cont\_speech\_frames window to not be filled at the edges of the DA and thereby lowering the value of the feature.
Significance tests for the leave-one-out trees showed that removal of the main feature types used in the all-features tree—that is, duration, pause, and energy features—resulted in a significant reduction in classification accuracy: $`p<.001`$, $`p<.05`$, and $`p<.05`$, respectively. Although significant, the reduction was not large in absolute terms, as seen from the figure and the $`\alpha `$ levels for significance. For the leave-one-in trees, results were in all cases significantly lower than that of the all-features trees; however, duration and pause features alone each yielded accuracy rates near that of the all-features tree. Although neither F0 nor enrate was used in the all-features tree, each individually was able to distinguish the DAs at rates significantly better than chance ($`p<.0001`$).
#### 4.5.2 Integration with language model
Integration results are reported in Table 18. Several observations are noteworthy. First, integrating the tree with word models improves performance considerably for all three test sets. Sign tests run for the larger HLD set showed a highly significant gain in accuracy by adding prosody, $`p<.00001`$. The DEV set did not contain enough samples for sufficient power to reject the null hypothesis, but showed the same pattern of results as the HLD set for both true and recognized words, and thus would be expected to reach significance for a larger data set. Second, for this analysis, the prosodic tree has better accuracy than the true words for the HLD set. Third, comparison of the data for the different test sets reveals an unusual pattern of results. Typically (and in the previous analyses), accuracy results for tree and word models were better for the HLD than for the DEV set. As noted in the Method section, HLD waveforms were segmented into DAs in the same manner (automatically) as the training data, while DEV data were carefully segmented by hand. For this task, however, results for both tree and word models are considerably better for the DEV data, i.e., the mismatched case (see also Figure 2). This pattern can be understood as follows. In the automatically segmented training and HLD data, utterances with “bad” estimated start or end times were thrown out of the analysis, as described in the Method section. The DAs most affected by the bad time marks were very short DAs, many of which were brief, single-word Backchannels such as “yeah”. Thus, the data remaining in the training and HLD sets are biased toward longer DAs, while the data in the DEV set retain the very brief DAs. Since the present task pits Backchannels against the longer Agreements, an increase in the percentage of shorter Backchannels (from training to test, as occurs when testing on the DEV data) can only enhance discriminability for the prosodic trees as well as for the language model.
## 5 Summary and General Discussion
### 5.1 Feature Importance
Across analyses we found that a variety of features were useful for DA classification. Results from the leave-one-out and leave-one-in trees showed that there is considerable redundancy in the features; typically there is little loss when one feature type is removed. Interestingly, although canonical or predicted features such as F0 for questions are important, less predictable features (such as pause features for questions) show similar or even greater influence on results.
Duration was found to be important not only in the seven-way classification, which included both long and short utterance types, but also for subtasks within general length categories (e.g., Statements versus Questions, Backchannels versus Agreements). Duration was also found to be useful as an added knowledge source to language model information, even though the length in words of an utterance is indirectly captured by the language model. Across tasks, the most-queried duration features were not raw duration, but rather duration-related measures that relied on the computation of other feature types.
F0 information was found to be important, as expected, for the classification of Questions, particularly when questions were broken down by type. However, it was also of use in many other classification tasks. In general, the main contribution from F0 features for all but the Question task came from global features (such as overall mean or gradient) rather than local features (such as the penultimate and end features, or the intonational event features). An interesting issue to explore in future work is whether this is a robustness effect, or whether global features are inherently better predictors of DAs than local features such as accents and boundaries.
Energy features were particularly helpful for classifying Incomplete Utterances, but also for the classification of Agreements and Backchannels. Analysis of the usage of energy features over all tasks revealed that SNR-based features were queried more than 4.8 times as often as features based on the raw RMS energy. Similarly, when the individual leave-one-in analyses for energy features were computed using only RMS versus only SNR features, results were consistently better for the SNR experiments. This suggests that for telephone speech or speech data collected under noisy conditions, it is important to estimate the energy of the speaker above the noise floor.
Enrate, the experimental speaking-rate feature from ? (?), proved to be useful across analyses in the following way. Although no task was significantly affected when enrate features were removed, enrate systematically achieved good performance when used alone. It was always better alone than at least one of the other main prosodic feature types alone. Furthermore, it provided remarkable accuracy for the classification of Questions and Statements, without any conversation-level normalization. Thus, the measure could be a valuable feature to include in a system, particularly if other more costly features cannot be computed.
Finally, across analyses, gender was not used in the trees. This suggests that gender-dependent features such as F0 were sufficiently normalized to allow gender-independent modeling. Since many of the features were normalized with respect to all values from a conversation side, it is possible that men and women do differ in the degree to which they use different prosodic features (even after normalization for pitch range), but that we cannot discern these differences here because speakers have been normalized individually.
Overall, the high degree of feature compensation found across tasks suggests that automatic systems could be successful using only a subset of the feature types. However, we also found that different feature types are used to varying degrees in the different tasks, and it is not straightforward at this point to predict which features will be most important for a task. Therefore, for best coverage on a variety of classification tasks, it is desirable to have as many different feature types available as possible.
### 5.2 Integration of Trees with Language Models
Not only were the prosodic trees able to classify the data at rates significantly above chance, but they also provided a consistent advantage over word information alone. To summarize the integration experiments: all tasks with the exception of the Incomplete Utterance task showed a significant improvement over words alone for the HLD set. For the Incomplete Utterance task, results for the DEV set were marginally significant. In all cases, the DEV set lacked power because of small sample size, making it difficult to reach significance in the comparisons. However, the relative win on the DEV set was consistently larger for the experiments using recognized rather than true words. This pattern of results suggests that prosody can provide significant benefit over word information alone, particularly when word recognition is imperfect.
## 6 Future Work
### 6.1 Improved DA Classification
One aim of future work is to optimize the prosodic features, and better understand the correlations among them. In evaluating the contribution of features, it is important to take into account such factors as measurement robustness and inherent constraints leading to missing data in our trees. For example, duration is used frequently, but it is also (unlike, e.g., F0 information) available and fairly accurately extracted for all utterances. We would also like to better understand which of our features capture functional versus semantic or paralinguistic information, as well as the extent to which features are speaker-dependent.
A second goal is to explore additional features that do not depend on the words. For example, we found that whether or not an utterance is turn-initial and/or turn-final, and the rate of interruption (including overlaps) by the other speaker, can significantly improve tree performance for certain tasks. In our overall model, we consider turn-related features to be part of the dialog grammar. Nevertheless, if one wanted to design a system that did not use word information, turn features could be used along with the prosodic features to improve performance overall.
Third, although we chose to use decision trees for the reasons given earlier, we might have used any suitable probabilistic classifier, i.e., any model that estimates the posterior probabilities of DAs given the prosodic features. We have conducted preliminary experiments to assess how neural networks compare to decision trees for the type of data studied here. Neural networks are worth investigating since they offer potential advantages over decision trees. They can learn decision surfaces that lie at an angle to the axes of the input feature space, unlike standard CART trees, which always split continuous features on one dimension at a time. The response function of neural networks is continuous (smooth) at the decision boundaries, allowing them to avoid hard decisions and the complete fragmentation of data associated with decision tree questions. Most important, neural networks with hidden units can learn new features that combine multiple input features. Results from preliminary experiments on a single task showed that a softmax network ( ?) without hidden units resulted in a slight improvement over a decision tree on the same task. The fact that hidden units did not afford an advantage indicates that complex combinations of features (as far as the network could learn them) may not better predict DAs for the task than linear combinations of our input features.
Thus, whether or not substantial gains can be obtained using alternative classifier architectures remains an open question. One approach that looks promising given the redundancy among different feature types is a combination of parallel classifiers, each based on a subcategory of features, for example using the mixture-of-experts framework ( ?). We will also need to develop an effective way to combine specialized classifiers (such as those investigated for the subtasks in this study) into an overall classifier for the entire DA set.
Finally, many questions remain concerning the best way to integrate the various knowledge sources. Instead of treating words and prosody as independent knowledge sources, as done here for simplicity, we could provide both types of cues to a single classifier. This would allow the model to account for interactions between prosodic cues and words, such as word-specific prosodic patterns. The main problem with such an approach is the large number of potential input values that “word features” can take on. A related question is how to combine prosodic classifiers most effectively with dialog grammars and the contextual knowledge sources.
### 6.2 Automatic Dialog Act Classification and Segmentation
Perhaps the most important area for future work is the automatic segmentation of dialogs into utterance units. As explained earlier, we side-stepped the segmentation problem for the present study by using segmentations by human labelers. Eventually, however, a fully automatic dialog annotation system will have to perform both segmentation and DA classification. Not only is this combined task more difficult, it also raises methodological issues, such as how to evaluate the DA classification on incorrectly identified utterance units. One approach, taken by ? (?), is to evaluate recognized DA sequences in terms of substitution, deletion, and insertion errors, analogous to the scoring of speech recognition output.
As noted in the Introduction, a large body of work addresses segmentation into intonational units or prosodic phrases, and utterance segmentation can be considered as a special case of prosodic boundary detection. To our knowledge, there are no published results for performing utterance-level segmentation of spontaneous speech by using only acoustic evidence, i.e., without knowledge of the correct words. Studies have investigated segmentation assuming that some kind of word-level information is given. ? (?) and ? (?) investigate DA segmentation and classification in the (task-oriented) Verbmobil domain, combining neural-network prosodic models with N-gram models for segment boundary detection, as well as N-gram and decision tree DA models with N-gram discourse grammars for DA classification, in a mathematical framework very similar to the one used here. ? (?) and ? (?) both investigated segmentation of spontaneous, Switchboard-style conversations using word-level N-gram models. ? (?) observed that word-level N-gram segmentation models work best when using a combination of parts-of-speech and cue words, rather than words alone.
Both ? (?) and ? (?) propose an A search for integrated DA segmentation and labeling. However, the results of ? (?) show only a small improvement over a sequential (first segment, then label) approach, and ? (?) found that segmentation accuracy did not change significantly as a result of modeling DAs in the segment language model. These findings indicate that a DA-independent utterance segmentation, followed by DA labeling using the methods described here, will be a reasonable strategy for extending our approach to unsegmented speech. This is especially important since our prosodic features rely on known utterance boundaries for extraction and normalization.
### 6.3 Dialog Act Classification and Word Recognition
As mentioned in the Introduction, in addition to dialog modeling as a final goal, there are other practical reasons for developing methods for automatic DA classification. In particular, DA classification holds the potential to improve speech recognition accuracy, since language models constrained by the DA can be applied when the utterance type is known. There has been little work involving speech recognition output for large annotated natural speech corpora. One relevant experiment has been conducted as part of our larger WS97 discourse modeling project, described in detail elsewhere ( ?).
To put an upper bound on the potential benefit of the approach, it is most meaningful to consider the extent to which word recognition accuracy could be improved if one’s automatic DA classifier had perfect accuracy. We therefore conducted experiments in which our language models were conditioned on the correct (i.e., hand-labeled) DA type. From the perspective of overall word accuracy results, the outcome was somewhat discouraging. Overall, the word error rate dropped by only 0.9% absolute, from a baseline of 41.2% to 40.9%. On the other hand, if one considers the Switchboard corpus statistics, results are in line with what one would predict for this corpus. In Switchboard, roughly 83% of all test set words were contained in the Statement category. Statements are thus already well-represented in the baseline language model. It is not surprising, then, that the error rate for Statements was reduced by only 0.5%. The approach was successful, however, for reducing word error for other DA types. For example, for Backchannels and No-Answers, word error was reduced significantly (by 7% and 18%, respectively). But because these syntactically restricted categories tend to be both less frequent and shorter than Statements, they contributed too few words to have much of an impact on the overall word error rate.
The DA-specific error reduction results suggest that although overall word accuracy for Switchboard was little improved in our experiments, DA classification could substantially benefit word recognition results for other types of speech data, or when evaluating on specific DA types. This should be true particularly for domains with a less skewed distribution of DA types. Similarly, DA modeling could reduce word error for corpora with a more uniform distribution of utterance lengths, or for applications where it is important to correctly recognize words in a specific subset of DAs.
## 7 Conclusion
We have shown that in a large database of natural human-human conversations, assuming equal class prior probabilities, prosody is a useful knowledge source for a variety of DA classification tasks. The features that allow this classification are task-dependent. Although canonical features are used in predicted ways, other less obvious features also play important roles. Overall there is a high degree of correlation among features, such that if one feature type is not available, other features can compensate. Finally, integrating prosodic decision trees with DA-specific statistical language models improves performance over that of the language models alone, particularly in a realistic setting where word information is based on automatic recognition. We conclude that DAs are redundantly marked in free conversation, and that a variety of automatically extractable prosodic features could aid the processing of natural dialog in speech applications.
## Appendix A Appendix A: Table of Original Dialog Acts
The following table lists the 42 original (before grouping into classes) dialog acts. Counts and relative frequencies were obtained from the corpus of 197,000 utterances used in model training.
| Dialog Act | Tag | Example | Count | % |
| --- | --- | --- | --- | --- |
| Statement-non-opinion | sd | Me, I’m in the legal department. | 72,824 | 36 |
| Acknowledge (Backchannel) | b | Uh-huh. | 37,096 | 19 |
| Statement-opinion | sv | I think it’s great. | 25,197 | 13 |
| Agree/Accept | aa | That’s exactly it. | 10,820 | 5 |
| Abandoned or Turn-Exit | % …-/ | So, -/ | 10,569 | 5 |
| Appreciation | ba | I can imagine. | 4,633 | 2 |
| Yes-No-Question | qy | Do you have to have any special training? | 4,624 | 2 |
| Non-verbal | x | $`<`$Laughter$`>`$,$`<`$Throat\_clearing$`>`$ | 3,548 | 2 |
| Yes-Answer | ny | Yes. | 2,934 | 1 |
| Conventional-closing | fc | Well, it’s been nice talking to you. | 2,486 | 1 |
| Uninterpretable | % | But, uh, yeah. | 2,158 | 1 |
| Wh-Question | qw | Well, how old are you? | 1,911 | 1 |
| No-Answer | nn | No. | 1,340 | 1 |
| Acknowledge-Answer | bk | Oh, okay. | 1,277 | 1 |
| Hedge | h | I don’t know if I’m making any sense or not. | 1,182 | 1 |
| Declarative Yes-No-Question | qy^d | So you can afford to get a house? | 1,174 | 1 |
| Other | o,fo | Well give me a break, you know. | 1,074 | 1 |
| Backchannel-Question | bh | Is that right? | 1,019 | 1 |
| Quotation | ^q | He’s always saying “why do they have to be here?” | 934 | .5 |
| Summarize/Reformulate | bf | Oh, you mean you switched schools for the kids. | 919 | .5 |
| Affirmative Non-Yes Answers | na | It is. | 836 | .4 |
| Action-directive | ad | Why don’t you go first | 719 | .4 |
| Collaborative Completion | ^2 | Who aren’t contributing. | 699 | .4 |
| Repeat-phrase | b^m | Oh, fajitas. | 660 | .3 |
| Open-Question | qo | How about you? | 632 | .3 |
| Rhetorical-Questions | qh | Who would steal a newspaper? | 557 | .2 |
| Hold before Answer/Agreement | ^h | I’m drawing a blank. | 540 | .3 |
| Reject | ar | Well, no. | 338 | .2 |
| Negative Non-No Answers | ng | Uh, not a whole lot. | 292 | .1 |
| Signal-non-understanding | br | Excuse me? | 288 | .1 |
| Other Answers | no | I don’t know. | 279 | .1 |
| Conventional-opening | fp | How are you? | 220 | .1 |
| Or-Clause | qrr | or is it more of a company? | 207 | .1 |
| Dispreferred Answers | arp,nd | Well, not so much that. | 205 | .1 |
| Third-party-talk | t3 | My goodness, Diane, get down from there. | 115 | .1 |
| Offers, Options & Commits | oo,cc,co | I’ll have to check that out. | 109 | .1 |
| Self-talk | t1 | What’s the word I’m looking for? | 102 | .1 |
| Downplayer | bd | That’s all right. | 100 | .1 |
| Maybe/Accept-part | aap/am | Something like that. | 98 | $`<`$.1 |
| Tag-Question | ^ g | Right? | 93 | $`<`$.1 |
| Declarative Wh-Question | qw^d | You are what kind of buff? | 80 | $`<`$.1 |
| Apology | fa | I’m sorry. | 76 | $`<`$.1 |
| Thanking | ft | Hey thanks a lot. | 67 | $`<`$.1 |
## Appendix B Appendix B: Estimated accuracy of transcript-based labeling
The table below shows the estimated recall and precision of hand-labeling utterances using only the transcribed words.
The estimates are computed using the results of “Round I” relabeling with listening to speech (see the Method section) as reference labels. DA types are sorted by their occurrence count in the relabeled subcorpus of 44 conversations.
For a given DA type, let $`a`$ be the number of original (labeled from text only) DA tokens of that type, $`b`$ the number of DA tokens after relabeling with listening, and $`c`$ the number of tokens that remained unchanged in the relabeling. Recall is estimated as $`\frac{b}{a}`$ and precision as $`\frac{c}{a}`$.
| Dialog Act | Tag | Recall (%) | Precision (%) | Count |
| --- | --- | --- | --- | --- |
| Statement-non-opinion | sd | 98.8 | 98.9 | 2147 |
| Statement-opinion | sv | 97.9 | 97.7 | 989 |
| Acknowledge (Backchannel) | b | 99.1 | 95.4 | 986 |
| Abandoned/Uninterpretable | % | 99.8 | 99.4 | 466 |
| Agree/Accept | aa | 86.5 | 99.3 | 327 |
| Yes-No-Question | qy | 100.0 | 98.0 | 144 |
| Non-verbal | x | 100.0 | 100.0 | 99 |
| Appreciation | ba | 100.0 | 94.6 | 70 |
| Yes-Answer | ny | 95.7 | 98.5 | 70 |
| Wh-Question | qw | 98.3 | 100.0 | 59 |
| Summarize/Reformulate | bf | 100.0 | 97.8 | 44 |
| Hedge | h | 93.0 | 97.6 | 43 |
| Quotation | ^q | 100.0 | 100.0 | 38 |
| Declarative Yes-No-Question | qy^d | 92.1 | 97.2 | 38 |
| Acknowledge-Answer | bk | 100.0 | 100.0 | 34 |
| No-Answer | nn | 100.0 | 100.0 | 33 |
| Other | o,fo | 100.0 | 100.0 | 33 |
| Open-Question | qo | 100.0 | 100.0 | 27 |
| Backchannel-Question | bh | 95.5 | 100.0 | 22 |
| Action-directive | ad | 100.0 | 95.5 | 21 |
| Collaborative Completion | ^2 | 100.0 | 94.7 | 18 |
| Hold before Answer/Agreement | ^h | 100.0 | 100.0 | 18 |
| Affirmative Non-Yes Answers | na | 100.0 | 100.0 | 18 |
| Repeat-phrase | b^m | 100.0 | 100.0 | 17 |
| Conventional-closing | fc | 100.0 | 100.0 | 16 |
| Reject | ar | 100.0 | 100.0 | 13 |
| Or-Clause | qrr | 100.0 | 100.0 | 11 |
| Other Answers | no | 100.0 | 100.0 | 10 |
| Rhetorical-Questions | qh | 80.0 | 100.0 | 10 |
| Signal-non-understanding | br | 100.0 | 87.5 | 7 |
| Negative Non-No Answers | ng | 100.0 | 100.0 | 6 |
| Maybe/Accept-part | aap/am | 100.0 | 100.0 | 5 |
| Conventional-opening | fp | 100.0 | 100.0 | 5 |
| Tag-Question | ^g | 100.0 | 100.0 | 4 |
| Offers, Options & Commits | oo,cc,co | 100.0 | 100.0 | 3 |
| Thanking | ft | 100.0 | 100.0 | 2 |
| Downplayer | bd | 100.0 | 100.0 | 1 |
| Declarative Wh-Question | qw^d | 100.0 | 100.0 | 1 |
| Self-talk | t1 | 100.0 | 50.0 | 1 |
| Third-party-talk | t3 | 100.0 | 100.0 | 1 |
| Dispreferred Answers | arp,nd | - | - | 0 |
| Apology | fa | - | - | 0 |
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# 1 Introduction
## 1 Introduction
Conformally invariant boundary conditions of two-dimensional conformal field theories arise in the study of defects in systems of condensed matter physics, of percolation probabilities and of (open) string perturbation theory in the background of D-branes. They are presently under active investigation. Boundary conditions that preserve all bulk symmetries for theories with charge conjugation modular invariant have been treated by Cardy . The two basic results are the following: Boundary conditions are labelled by the primary fields of the theory, and the annulus multiplicities are given by the fusion rules. Together with the information that all bulk symmetries are preserved, these two results allow in particular to recover the so-called boundary states, which encode all one-point amplitudes on the disk.
More recently, these results have been generalized in several directions . In particular, those boundary conditions have been classified which preserve an abelian orbifold subalgebra of the algebra $`𝔄`$ of bulk symmetries, i.e. for which the preserved symmetries can be characterized as the subalgebra $`\overline{𝔄}=𝔄^G`$ of symmetries that are fixed by some abelian group $`G`$ of automorphisms of $`𝔄`$. Boundary states for such boundary conditions have been given explicitly, and the integrality of the annulus coefficients was proven. It was also shown that correlation functions in the presence of such boundary conditions can be written as linear combinations of twisted conformal blocks. As a special case, the boundary states can be expressed in terms of twisted Ishibashi states $`|\lambda ,\omega `$, which are characterized by the identity
$$(Y_n\mathrm{𝟏}+(1)^{\mathrm{\Delta }_Y1}\mathrm{𝟏}\omega (Y_n))|\lambda ,\omega =0$$
(1)
for every primary field $`Y`$ (of conformal weight $`\mathrm{\Delta }_Y`$) in the chiral symmetry algebra $`𝔄`$. Here $`\omega G`$ is an automorphism of the bulk symmetries that leaves the Virasoro algebra invariant, $`\omega (L_n)=L_n`$. In a sense, the relations (1) express the fact that at the boundary left movers and right movers are connected by the automorphism $`\omega `$. The automorphism $`\omega `$ has been called the automorphism type, or gluing automorphism, of the boundary condition. We will say that such boundary conditions possess a definite automorphism type, in this case $`\omega `$.
In the present letter, we study more general patterns of symmetry breaking by boundaries, in which left movers and right movers are not necessarily related any more by automorphisms. In more precise terms, this means that the boundary conditions preserve a subalgebra $`\overline{𝔄}`$ of the algebra $`𝔄`$ of bulk symmetries that cannot any longer be characterized as a fixed algebra under some group of automorphisms. We refer to such boundary conditions as boundary conditions without automorphism type, or without gluing automorphism. Examples of such boundary conditions appear already for the $`_2`$-orbifold of a compactified free boson. Other examples are provided by various conformal embeddings; boundary conditions associated to conformal embeddings have been studied in , in particular in their relation with certain graph algebras .
Boundary conditions without automorphism type are of direct relevance in string theory: they correspond to non-BPS branes. Indeed, every chiral algebra automorphism $`\omega `$ maps the vertex operator of a space-time supercharge to the vertex operator of another supercharge. Therefore validity of (1) immediately implies that boundary conditions with automorphism type preserve half of the space-time supersymmetries, and hence are BPS.
The purpose of the present note is to generalize Cardy’s results and those of once more. We consider conformally invariant boundary conditions of a rational conformal field theory with chiral algebra $`𝔄`$ that preserve a subalgebra $`\overline{𝔄}`$ of $`𝔄`$ that is still rational, but otherwise arbitrary. We give a natural labelling of the boundary conditions and compute the annulus coefficients. By a modular transformation, this allows to determine the boundary states.
For certain conformal field theories, one can construct (nets of) factors; their irreducible local sectors (inner unitary equivalence class of representations) are in one-to-one correspondence to the primary fields of the conformal field theory. Our main tool is an adaptation of a certain form of induction for sectors, the so-called $`\alpha `$-induction, which was developped in the framework of subfactor theory. Applying $`\alpha `$-induction to a sector of the subfactor, it produces a sector of the ambient factor. Among the sectors obtained this way there are ordinary, local, sectors as well as solitonic sectors.
For the purposes of this letter, we do not have to know the relevant nets of subfactors and their sectors in any detail. Rather, we simply postulate that the process of $`\alpha `$-induction works at the level of the representation category of the conformal field theory under investigation. Thus we regard the irreducible sectors as the primary fields, or rather as the associated basis elements of the fusion ring; general sectors correspond to arbitrary elements of the fusion ring, and the composition of sectors is simply the fusion product. In fact, all we need to know is the action of $`\alpha `$-induction on primary fields. It will provide us with solitonic sectors which precisely label symmetry breaking boundary conditions. Moreover, the fusion of these sectors will provide us with the annulus multiplicities. As in the Cardy case, these data, together with the preserved symmetries $`\overline{𝔄}`$, allow to construct the boundary states.
In section 2 we discuss the labelling problem for bulk and boundary fields. Motivated by constructions from topological field theory, we are led to the concepts of bulk and boundary categories. Our prescription for the boundary category is presented in section 3. It is based on imposing a version of $`\alpha `$-induction at the level of the fusion rules. In section 4 two illustrative examples are analyzed.
## 2 Symmetry breaking boundary conditions
Before we explain the case of boundary conditions without automorphism type, we briefly rephrase some of the results of on boundary conditions that do possess an automorphism type. As was shown there, boundary conditions leaving $`\overline{𝔄}=𝔄^G`$ invariant correspond to orbits $`[\overline{\mu },\psi ]`$ of primary fields $`\overline{\mu }`$ of the $`\overline{𝔄}`$-theory with respect to a group $`𝒢G^{}`$ of simple current fields (the degeneracy label $`\psi `$ is a character of a suitable subgroup of $`𝒢`$). The monodromy charge of $`\overline{\mu }`$ with respect to $`𝒢`$ is not restricted. These labels $`[\overline{\mu },\psi ]`$, in turn, can be seen to correspond to representations of $`𝔄`$; these are twisted representations when the monodromy charge of $`\overline{\mu }`$ is not zero. <sup>1</sup><sup>1</sup>1 When the action of the twists is only projective, additional subtleties arise. Twisted representations of vertex operator algebras have been investigated e.g. in . The notion of fusion of such representations has also been studied , but little is known about the resulting fusion ring. However, one can re-write the annulus amplitudes derived in as sums of characters of the twisted representations. It is therefore reasonable to expect that the annulus coefficients determined in precisely coincide with the fusion rules of the twisted representations.
A second ingredient we will need is the description of correlation functions in the presence of boundaries through three-dimensional topological field theory that was developped in for symmetry preserving boundary conditions. In that context, a three-dimensional manifold $`M_X`$, the connecting manifold, was constructed to compute correlators on a two-dimensional world sheet $`X`$. $`M_X`$ has a boundary, and this boundary is isomorphic to the so-called double $`\widehat{X}`$ of $`X`$. The connecting manifold is universal in the sense that it is the same for all rational conformal field theories. One also needs to prescribe a Wilson graph in $`M_X`$. Bulk points on the world sheet $`X`$ possess two pre-images on its double $`\widehat{X}=M_X`$, and these two points are connected by a natural interval in $`M_X`$. A Wilson line carrying the bulk label is placed in that interval.
As for the boundary data, a circular Wilson line must be placed parallel to each boundary component. Insertions of boundary fields are linked with short Wilson lines to the corresponding circular Wilson line. This is summarized in the following picture for the case when $`X`$ is the disk (then $`\widehat{X}`$ is the sphere and $`M_X`$ is a solid three-ball):
(2)
We will assume that the Wilson graph is universal as well, i.e. that the same graph is still to be used for symmetry breaking boundary conditions. The boundary graph must be labelled, too. In the Cardy case, we use the same type of labels as for the bulk components, i.e. primary fields of $`𝔄`$. This nicely fits together with Cardy’s result that symmetry preserving boundary conditions are in one-to-one correspondence with primary fields. Moreover, in this description trivalent vertices, couplings, appear naturally in the boundary graph. They involve two boundary conditions and the chiral label of a boundary insertion.
We remark that typically the Wilson graph is not connected. What is crucial in the present context is that there is never a Wilson line that connects a bulk insertion with a boundary insertion or with a segment of the Wilson graph that encodes a boundary condition. We are thereby led to the following general picture. Boundary conditions as well as field insertions on the boundary should be characterized by basis elements of the same ‘fusion ring’, and the structure constants of this fusion ring should coincide with the annulus coefficients. In more technical terms, the labels should correspond to the (isomorphism classes of) simple objects of a suitable tensor category. We call this structure the boundary category. Bulk fields, on the other hand, will have to be described by a different structure, namely as the simple objects of a bulk category; their treatment is beyond the scope of the present note. Notice that the interpretation in terms of boundary insertions requires that associated to each boundary label there comes a natural state space, which is a module over (at least) the Virasoro algebra.
In Cardy’s case – torus partition function of charge conjugation type, and only symmetry preserving boundary conditions – both the boundary and the bulk category are just the one associated to the fusion ring of the full bulk symmetry $`𝔄`$, and the state space associated to a label $`\lambda `$ is nothing but the corresponding irreducible $`𝔄`$-module $`_\lambda `$. For a general boundary insertion, the state space must be an appropriate generalization of a twisted representation, or in other words, a solitonic sector. In Cardy’s case, the tensor category has actually a lot more structure. In particular it is modular , i.e., roughly speaking, there is a braiding , leading to the unitarity of the modular matrix $`S`$, as well as a twist of Wilson lines, which corresponds to the modular matrix $`T`$.
This is too much structure than may be expected to be present for a boundary category in the general case. As examples of twisted representations show, the twist of a Wilson line will, in general, no longer be a well-defined operation on the boundary. Correspondingly, the boundary sectors will not have a unique conformal weight up to integers. (The bulk category, on the other hand, will still possess a $`T`$-matrix.) Similarly, neither the bulk nor the boundary category will possess a braiding any more, in general. However, it is still possible to braid an object of the bulk category with an object of the boundary category. This braiding allows us to establish a generalization of the diagonalizing matrix $`\stackrel{~}{S}`$ that was introduced in . We expect that this matrix $`\stackrel{~}{S}`$ is square; nevertheless its two indices take values in two different sets: the rows are labelled by the bulk fields, whereas the columns are labelled by the boundary conditions. Up to normalization, $`\stackrel{~}{S}`$ provides the coefficients for the expansion of the boundary states with respect to the boundary conformal blocks (generalized Ishibashi states). We conjecture that the matrix $`\stackrel{~}{S}`$ obtained this way is invertible and hence in particular indeed a square matrix (and can be chosen unitary). This implies that the bulk and the boundary category have the same number of simple objects; as a consequence, the number of boundary conditions can be read off from the bulk modular invariant. The equality between the number of bulk fields and the number of boundary conditions has been derived in for boundary conditions with definite automorphism type; we will see in the examples below that it holds for more general symmetry breakings as well.
## 3 Solitonic sectors from $`\alpha `$-induction
According to the reasoning above, we would like to view the boundary labels $`\mu `$ as simple objects of a suitable category, and accordingly the associated intertwiners as morphisms of that category. It will, however, be important that we can regard the sectors also as (isomorphism classes of) representations – including both ordinary and solitonic representations – of the chiral algebra or vertex operator algebra $`𝔄`$. These representation spaces provide the spaces of open string states whose partition function is the annulus amplitude.
How can we then obtain the boundary category? We wish to find ‘solitonic’ $`𝔄`$-representations. Fortunately, there is one situation in which this can be done totally explicitly. Namely, chiral WZW theories can also be analyzed in the framework of nets of operator algebras. In that context, one can employ the notion of $`\alpha `$-induction to arrive at solitonic sectors of the chiral theories .
The concrete construction of operator algebras for more general chiral conformal field theories is a difficult problem; see , as well as for the case of coset and orbifold theories. For the purposes of the present note, we need not address these questions which are definitely quite important. Rather, we will only abstract from $`\alpha `$-induction and its adjoint operation, $`\sigma `$-restriction, a few properties at the level of fusion rings. We present these properties in the form of a recipe. (But we expect that they are indeed realized in any decent conformal field theory, and actually that the existence of such an induction procedure can be entirely established in the context of the relevant tensor categories.)
We start by prescribing those bulk symmetries that are preserved by the boundary conditions. These symmetries must form a consistent rational subalgebra $`\overline{𝔄}`$ of the chiral algebra $`𝔄`$. We assume that the collection of all (isomorphism classes of) ordinary – i.e. non-solitonic – irreducible modules $`\overline{}_{\overline{\lambda }}`$ over $`\overline{𝔄}`$ gives rise to a modular fusion ring with basis $`\{\overline{\mathrm{\Phi }}_{\overline{\lambda }}\}`$. From the embedding of $`\overline{𝔄}`$ into $`𝔄`$ we determine how the vacuum module $`_\mathrm{\Omega }`$ of $`𝔄`$ decomposes into a direct sum of irreducible $`\overline{𝔄}`$-modules:
$$_\mathrm{\Omega }=\underset{\overline{\mu }}{}b_{\overline{\mu }}\overline{}_{\overline{\mu }}.$$
(3)
(Thus $`b_{\overline{\lambda }}_0`$ is the multiplicity with which the $`\overline{𝔄}`$-module $`\overline{}_{\overline{\lambda }}`$ appears in the vacuum module of the $`𝔄`$-theory.) This allows us to introduce the element
$$\overline{\theta }:=\underset{\overline{\mu }}{}b_{\overline{\mu }}\overline{\mathrm{\Phi }}_{\overline{\mu }}$$
(4)
of the fusion ring of $`\overline{𝔄}`$; we refer to $`\overline{\theta }`$ as the extending sector of the $`\overline{𝔄}`$-theory.
We now construct a new fusion ring as follows. First, it is generated by objects $`\alpha _{\overline{\lambda }}`$ for each basis element $`\overline{\mathrm{\Phi }}_{\overline{\lambda }}`$ of the fusion ring of $`\overline{𝔄}`$. The fusion product is defined by
$$\alpha _{\overline{\lambda }}\alpha _{\overline{\mu }}:=\alpha _{\overline{\lambda }\overline{\mu }},$$
(5)
and we also require that $`\alpha _{\overline{\lambda }\overline{\mu }}=\alpha _{\overline{\lambda }}+\alpha _{\overline{\mu }}`$ and $`\alpha _{\overline{\lambda }_{}^+}=(\alpha _{\overline{\lambda }})^+`$. This would not constitute anything new beyond what is encoded in the fusion ring of $`\overline{𝔄}`$, were it not for another piece of information. Namely, the fusion ring element $`\alpha _{\overline{\lambda }}`$ is also supposed to represent a – possibly twisted or solitonic – representation of $`𝔄`$, which for brevity we denote by the same symbol. An important point is that even for irreducible $`\overline{\lambda }`$ the $`𝔄`$-representation $`\alpha _{\overline{\lambda }}`$ need not necessarily be irreducible, and that the $`\alpha _{\overline{\mu }}`$, respectively their irreducible subrepresentations, for different values of $`\overline{\mu }`$ are allowed to be isomorphic. (Thus in particular the $`\alpha _{\overline{\lambda }}`$ generically do not form a basis of the fusion ring.) In view of Schur’s lemma, this information is conveniently encoded in the intertwiner spaces $`\mathrm{Hom}_𝔄(\alpha _{\overline{\lambda }},\alpha _{\overline{\mu }})`$. For instance, $`\alpha _{\overline{\lambda }}`$ is irreducible if and only if $`\mathrm{Hom}_𝔄(\alpha _{\overline{\lambda }},\alpha _{\overline{\lambda }})`$ is one-dimensional. Also, when $`\alpha `$ is a simple and $`\beta `$ any arbitrary object of the fusion category, then the dimension of $`\mathrm{Hom}_𝔄(\alpha ,\beta )`$ tells us how many times $`\alpha `$ appears in the decomposition of $`\beta `$. The $`\mathrm{Hom}`$ spaces are defined in terms of the intertwiner spaces in the fusion category of $`\overline{𝔄}`$ as follows:
$$\mathrm{Hom}_𝔄(\alpha _{\overline{\lambda }},\alpha _{\overline{\mu }}):=\mathrm{Hom}_{\overline{𝔄}}(\overline{\mathrm{\Phi }}_{\overline{\lambda }},\overline{\theta }\overline{\mathrm{\Phi }}_{\overline{\mu }}).$$
(6)
This system of $`\mathrm{Hom}`$ spaces obeys tight consistency constraints. For example, from $`\mathrm{Hom}_𝔄(\alpha _{\overline{\lambda }},\alpha _{\overline{\lambda }})`$ we compute the number $`n_{\overline{\lambda }}`$ of irreducible subsectors of $`\alpha _{\overline{\lambda }}`$. If one would just prescribe an extending sector $`\overline{\theta }`$ at random, one might find contradictions of the type that more than $`n_{\overline{\lambda }}`$ irreducible sectors have non-trivial intertwiners with $`\alpha _{\overline{\lambda }}`$. The existence of a system of $`\mathrm{Hom}`$ spaces that is free of contradiction is therefore highly non-trivial and requires special properties of $`\overline{\theta }`$. A necessary condition is of course that all irreducible $`\overline{𝔄}`$-subsectors of $`\overline{\theta }`$ are mutually local, but this condition is typically far from being sufficient. It would be rewarding to find a characterization of consistent extending sectors purely at the level of fusion rings. It will then be particularly interesting to compare the problem of classifying consistent extending sectors with the problem of classifying modular invariant partition functions of extension type.
It is also important that along with $`\alpha `$-induction there comes an “adjoint” operation, known as $`\sigma `$-restriction. Namely, every sector $`\beta `$ of $`𝔄`$, whether solitonic or not, may be seen as a (typically reducible) sector of $`\overline{𝔄}`$, which we denote as $`\sigma (\beta )`$. Induction and restriction are related by the reciprocity relation
$$\mathrm{Hom}_𝔄(\alpha _{\overline{\lambda }},\beta )\mathrm{Hom}_{\overline{𝔄}}(\overline{\lambda },\sigma (\beta )).$$
(7)
This implies that
$$\sigma (\alpha _{\overline{\mu }})=\overline{\theta }\overline{\mu }$$
(8)
and allows us to decompose induced solitonic sectors into irreducible $`\overline{𝔄}`$-sectors.
Let us pause and compare these ideas to the situation studied in . In that case, $`\overline{\theta }`$ can be written as a sum over so-called simple current sectors $`\overline{J}`$ which form a finite abelian group $`𝒢`$ under fusion, and each such simple current appears with multiplicity one:
$$\overline{\theta }=\underset{\overline{J}𝒢}{}\overline{J}.$$
(9)
Formula (6) then just summarizes how the fusion rules of a simple current extension are related to those of the original theory (in the category theoretical setting, this is discussed in ). However, it only allows for a direct determination of the extended fusion rules as long as no fixed points – that is, sectors $`\overline{\lambda }`$ with $`\overline{J}\overline{\lambda }=\overline{\lambda }`$ for some $`\overline{J}𝒢`$ – are involved. Indeed, in the simple current situation the induced sector $`\alpha _{\overline{\lambda }}`$ is reducible if and only if $`\overline{\lambda }`$ is a fixed point. The decomposition of $`\alpha _{\overline{\lambda }}`$ for a fixed point is precisely what is known as fixed point resolution in the theory of simple current extensions.
In the general case there is the following analogue of the problem caused by simple current fixed points. It can happen that the relations (6) do not provide enough information for decomposing all $`\alpha _{\overline{\lambda }}`$ into irreducible sectors. In that case, the category must be enlarged: sufficiently many additional irreducibles have to be introduced to provide subobjects. There exists a general procedure for doing so . But unfortunately fully explicit formulae, in particular for the modular $`S`$-matrix of the enlarged theory, are only known in the simple current case , where it leads in particular to the group character $`\psi `$ that appears in the description of boundary conditions with definite automorphism type, see above.
A more explicit understanding of these new irreducibles in the general case and in particular what their braiding properties with bulk fields are, might be called the generalized fixed point problem. To be precise, the task is to express the fusion products of the new irreducibles in terms of chiral data of the $`\overline{𝔄}`$-theory, like e.g. the modular matrices for one-point conformal blocks on the torus. We consider this to be a central problem in the study of solitonic sectors, and hence of conformally invariant boundary conditions. <sup>2</sup><sup>2</sup>2 A special version of the fixed point problem arises already when one aims to express the modular $`S`$-matrix of the $`𝔄`$-theory through chiral data of the $`\overline{𝔄}`$-theory. For exceptional extensions, no general solution to this problem is known. In the present letter, we restrict ourselves to examples where either this problem does not occur at all or where it can be resolved by using the knowledge about the simple current case.
Our prescription provides us explicitly with labels for the boundary conditions and the boundary insertions. The annulus multiplicities are just the tensor product multiplicities in the boundary category, and the open string states are organized in terms of the induced sectors. The induced sectors come in two classes: ordinary, non-solitonic sectors correspond to symmetry preserving boundary conditions, while the solitonic sectors are in correspondence with symmetry breaking boundary conditions. In the case of boundary conditions with automorphism type, the latter are just the orbits with non-vanishing monodromy charge. In the subfactor framework, ordinary and solitonic sectors can be distinguished by their localization properties.
Before we support our findings by examples, we wish to add a speculative comment. In the operator-algebraic definition of $`\alpha `$-induction, a braiding among $`\overline{𝔄}`$-sectors enters. In two dimensions there are two independent braidings – ‘over’- and ‘under’-braiding – which are each others’ inverse. As a consequence, there are in fact two $`\alpha `$-inductions, called $`\alpha ^\pm `$. It has been shown in that $`\alpha _{\overline{\lambda }}^\pm `$ is not solitonic if and only if $`\alpha _{\overline{\lambda }}^+`$ and $`\alpha _{\overline{\lambda }}^{}`$ are isomorphic. This suggests that solitonic representations, and thus symmetry breaking boundary conditions, actually come in pairs. However, only one version of $`\alpha `$-induction may be used at a time; so there is a twofold choice on which set of (symmetry breaking) boundary conditions one must take. It will be interesting to see whether this can explain the observations in , where two distinct sets of symmetry breaking boundary conditions were found; any two boundary conditions of the same set are compatible, while two boundary conditions belonging to distinct sets are mutually incompatible.
## 4 Examples
Our general ideas are easily illustrated by examples; we present two of them. The first example is the $`E_6`$-type modular invariant of $`A_1`$ at level $`10`$, which has already been discussed extensively elsewhere . We will show how the structures developed above allow to rederive and systematize the results of on the boundary conditions of this theory. The second example deals with the exceptional modular invariant of $`G_2`$ at level 3 and is, to the best of our knowledge, new.
The fusion ring of $`A_1`$ at level $`10`$ has eleven simple sectors, which we label by $`\overline{\mu }=\mathrm{\hspace{0.17em}0},1,\mathrm{},10`$. In this notation, the $`E_6`$-type modular invariant of $`A_1`$ reads
$$Z=|\chi _0+\chi _6|^2+|\chi _4+\chi _{10}|^2+|\chi _3+\chi _7|^2.$$
(10)
It corresponds to the conformal embedding into $`B_2`$ at level $`1`$. The first block comes from the vacuum $`o`$, the second from the vector $`v`$ and the third block from the spinor $`s`$ of $`B_2`$. The relevant aspects of $`\alpha `$-induction for this example can be found in \[13, Sec. 2.2 of II\]; here we summarize the most important features.
From the modular invariant (10) we read off the extending sector as $`\overline{\theta }=\overline{\mathrm{\Phi }}_0+\overline{\mathrm{\Phi }}_6`$. The dimensions of the $`\mathrm{Hom}`$ spaces are thus given by
$$\begin{array}{c}dim\mathrm{Hom}_𝔄(\alpha _{\overline{\mu }_1},\alpha _{\overline{\mu }_2})=dim\mathrm{Hom}_{\overline{𝔄}}(\overline{\mathrm{\Phi }}_{\overline{\mu }_1},\overline{\mathrm{\Phi }}_{\overline{\mu }_2})+dim\mathrm{Hom}_{\overline{𝔄}}(\overline{\mathrm{\Phi }}_{\overline{\mu }_1},\overline{\mathrm{\Phi }}_6\overline{\mathrm{\Phi }}_{\overline{\mu }_2})\hfill \\ \text{ }=\delta _{\overline{\mu }_1,\overline{\mu }_2}+\overline{𝒩}_{6,\overline{\mu }_2}^{\overline{\mu }_1},\hfill \end{array}$$
(11)
where $`\overline{𝒩}`$ are the fusion rules of $`A_1`$ at level $`10`$. Applying this to the case $`\overline{\lambda }_1=\overline{\lambda }_2`$, we find that the sectors $`\alpha _{\overline{\lambda }}`$ are irreducible for $`\overline{\lambda }=\mathrm{\hspace{0.17em}0},1,2,8,9,10`$, and contain two irreducible subsectors else.
Computing the $`\mathrm{Hom}`$ spaces between the irreducible $`\alpha _{\overline{\lambda }}`$ shows that they all vanish, except for $`\mathrm{Hom}(\alpha _2,\alpha _8)`$, which is one-dimensional. Hence the two irreducible sectors $`\alpha _2`$ and $`\alpha _8`$ are isomorphic, $`\alpha _2\alpha _8`$. Furthermore,
$$dim\mathrm{Hom}(\alpha _2,\alpha _4)=1=dim\mathrm{Hom}(\alpha _{10},\alpha _4),$$
(12)
so that $`\alpha _4\alpha _2+\alpha _{10}`$. Similarly, one finds $`\alpha _5\alpha _1+\alpha _9`$ and $`\alpha _6\alpha _0+\alpha _2`$. Thus these sectors do not give rise to new irreducible sectors. According to our general conjecture, and in accordance with the results of , we expect in total 6 boundary conditions and thus one additional simple object, which we call $`\alpha _3^{\left(1\right)}`$. Indeed we have $`dim\mathrm{Hom}_𝔄(\alpha _3,\alpha _{\overline{\mu }})=\delta _{\overline{\mu },3}+dim\mathrm{Hom}_{\overline{𝔄}}(\overline{\mathrm{\Phi }}_{\overline{\mu }},\overline{\mathrm{\Phi }}_6\overline{\mathrm{\Phi }}_3)=\delta _{\overline{\mu },3}+dim\mathrm{Hom}_{\overline{𝔄}}(\overline{\mathrm{\Phi }}_{\overline{\mu }},\overline{\mathrm{\Phi }}_3+\overline{\mathrm{\Phi }}_5+\overline{\mathrm{\Phi }}_7+\overline{\mathrm{\Phi }}_9)`$. So $`\alpha _3^{\left(1\right)}`$ appears in the decompositions
$$\alpha _3\alpha _9+\alpha _3^{\left(1\right)}\alpha _7\alpha _1+\alpha _3^{\left(1\right)}.$$
(13)
The $`\sigma `$-restriction is found from formula (7). First,
$$\sigma (\alpha _0)06,\sigma (\alpha _{10})410,\sigma (\alpha _3^{\left(1\right)})37,$$
(14)
showing that these sectors are the three non-solitonic sectors of $`B_2`$ that can already be inferred from the partition function (10), namely
$$o=\alpha _0,v=\alpha _{10},s=\alpha _3^{\left(1\right)}.$$
(15)
It is convenient to introduce a similar notation $`\stackrel{ˇ}{o},\stackrel{ˇ}{v},\stackrel{ˇ}{s}`$ for the three solitonic $`B_2`$-sectors; they restrict as follows:
$$\begin{array}{c}\sigma (\stackrel{ˇ}{o})\sigma (\alpha _1)157,\sigma (\stackrel{ˇ}{v})\sigma (\alpha _9)359,\hfill \\ \sigma (\stackrel{ˇ}{s})\sigma (\alpha _2)2468.\hfill \end{array}$$
(16)
It is readily checked that all annulus amplitudes reported in can indeed be written as linear combinations of the corresponding six specific sums of $`A_1`$-characters.
The fusion products of the sectors $`\alpha _{\overline{\mu }}`$ are computed with formula (5). For $`o,v`$ and $`s`$ we get the usual Ising fusion rules; they indeed provide the annuli of the Cardy boundary conditions. $`\alpha _0`$ acts generally as the identity under fusion. The remaining fusion rules between ordinary and solitonic sectors turn out to be
$$\begin{array}{ccc}v\stackrel{ˇ}{o}=\stackrel{ˇ}{v},\hfill & v\stackrel{ˇ}{v}=\stackrel{ˇ}{o},\hfill & v\stackrel{ˇ}{s}=\stackrel{ˇ}{s},\hfill \\ s\stackrel{ˇ}{o}=\stackrel{ˇ}{s},\hfill & s\stackrel{ˇ}{v}=\stackrel{ˇ}{s},\hfill & s\stackrel{ˇ}{s}=\stackrel{ˇ}{o}+\stackrel{ˇ}{v}.\hfill \end{array}$$
(17)
The fusion between two solitonic sectors produces ordinary as well as solitonic sectors; we find
$$\begin{array}{cc}\stackrel{ˇ}{o}\stackrel{ˇ}{o}=o+\stackrel{ˇ}{s},\hfill & \stackrel{ˇ}{v}\stackrel{ˇ}{v}=o+\stackrel{ˇ}{s},\hfill \\ \stackrel{ˇ}{o}\stackrel{ˇ}{v}=v+\stackrel{ˇ}{s},\hfill & \stackrel{ˇ}{v}\stackrel{ˇ}{s}=s+\stackrel{ˇ}{o}+\stackrel{ˇ}{v},\hfill \\ \stackrel{ˇ}{o}\stackrel{ˇ}{s}=s+\stackrel{ˇ}{o}+\stackrel{ˇ}{v},\hfill & \stackrel{ˇ}{s}\stackrel{ˇ}{s}=o+v+2\stackrel{ˇ}{s}.\hfill \end{array}$$
(18)
These fusion products exactly give the annulus multiplicities that have been found by different arguments in . Also the $`\stackrel{~}{S}`$ matrix can be computed. It reads
$$\stackrel{~}{S}=\frac{1}{d}\left(\begin{array}{cccccc}1& \sqrt{2}& 1& \frac{1}{\sqrt{2}}(1+\sqrt{3})& 1+\sqrt{3}& \frac{1}{\sqrt{2}}(1+\sqrt{3})\\ \frac{1}{2}d& 0& \frac{1}{2}d& \frac{1}{2}d& 0& \frac{1}{2}d\\ 1& \sqrt{2}& 1& \frac{1}{\sqrt{2}}(1+\sqrt{3})& 1+\sqrt{3}& \frac{1}{\sqrt{2}}(1+\sqrt{3})\\ \frac{1}{\sqrt{2}}(1+\sqrt{3})& 1+\sqrt{3}& \frac{1}{\sqrt{2}}(1+\sqrt{3})& 1& \sqrt{2}& 1\\ \frac{1}{2}d& 0& \frac{1}{2}d& \frac{1}{2}d& 0& \frac{1}{2}d\\ \frac{1}{\sqrt{2}}(1+\sqrt{3})& 1\sqrt{3}& \frac{1}{\sqrt{2}}(1+\sqrt{3})& 1& \sqrt{2}& 1\end{array}\right)$$
(19)
with $`d:=1/2\sqrt{3+\sqrt{3}}`$.
Our second example is the exceptional modular invariant of $`G_2`$ at level 3. It reads
$$Z=|\chi _{00}+\chi _{11}|^2+2|\chi _{02}|^2$$
(20)
and describes the conformal embedding into $`E_6`$ at level 1. Here $`G_2`$-sectors are characterized by their highest weights. The multiplicity two in the second term of $`Z`$ indicates that the 27-dimensional irreducible representation of $`E_6`$ and its conjugate restrict to the same irreducible $`G_2`$-representation. $`G_2`$ at level $`3`$ has six irreducible sectors. A careful analysis of the $`\mathrm{Hom}`$ spaces shows that $`\alpha _{01},\alpha _{03}`$ and $`\alpha _{10}`$ are irreducible and all isomorphic. $`\alpha _{00}`$ is, as always, irreducible, and indeed not isomorphic to $`\alpha _{01}`$. $`\alpha _{11}`$ contains two irreducibles, and one finds $`\alpha _{11}\alpha _{00}+\alpha _{01}`$ so that it does not give rise to any new irreducibles. Finally, $`dim\mathrm{Hom}(\alpha _{02},\alpha _{02})=\mathrm{\hspace{0.17em}3}`$, and $`\alpha _{02}`$ contains $`\alpha _{01}`$ as a subobject. We choose the notation $`\alpha _{02}^{(\pm )}`$ for its two other subobjects:
$$\alpha _{02}=\alpha _{01}+\alpha _{02}^{(+)}+\alpha _{02}^{()}.$$
(21)
The computation of the $`\sigma `$-restriction is straightforward, too. We get
$$\begin{array}{c}\sigma (\alpha _{01})=\overline{\mathrm{\Phi }}_{01}+\overline{\mathrm{\Phi }}_{02}+\overline{\mathrm{\Phi }}_{03}+\overline{\mathrm{\Phi }}_{10}+\overline{\mathrm{\Phi }}_{11},\hfill \\ \sigma (\alpha _{00})=\overline{\mathrm{\Phi }}_{00}+\overline{\mathrm{\Phi }}_{11},\sigma (\alpha _{02})=\sigma (\alpha _{01})+2\overline{\mathrm{\Phi }}_{02},\hfill \end{array}$$
(22)
from which we also learn that $`\sigma (\alpha _{02}^{(\pm )})=\overline{\mathrm{\Phi }}_{02}`$. We can therefore identify $`\alpha _{00}`$ as the vacuum sector of the $`E_6`$-theory and $`\alpha _{02}^{(\pm )}`$ as the sectors corresponding to the two 27-dimensional irreducible representations of $`E_6`$. In addition there is a single solitonic sector, given by $`\alpha _{01}`$. Notice that we obtain again the same number of simple objects in the bulk and in the boundary category.
It is readily checked that the fusion rules of $`\alpha _{00}`$ and $`\alpha _{02}^{(\pm )}`$ are indeed the $`_3`$ fusion rules of $`E_6`$ at level 1. The fusion rules involving $`\alpha _{01}`$ turn out to be
$$\begin{array}{c}\alpha _{01}\alpha _{00}=\alpha _{01},\alpha _{01}\alpha _{02}^{(+)}=\alpha _{01}=\alpha _{01}\alpha _{02}^{()},\hfill \\ \alpha _{01}\alpha _{01}=\alpha _{00}+3\alpha _{01}+\alpha _{02}^{(+)}+\alpha _{02}^{()}.\hfill \end{array}$$
(23)
Thus the fusion graph of $`\alpha _{01}`$ looks like
(24)
It has already been displayed in , where also the $`_3`$ fusion rules of the ordinary sectors were established by a different method..
These fusion rules provide the annulus multiplicities. Combining them with the modular $`S`$-matrix of the $`G_2`$-theory yields three symmetry preserving boundary states,
$$\begin{array}{cc}|00\hfill & =3^{1/4}(|00+|11+|02,++|02,),\hfill \\ |02,+\hfill & =3^{1/4}(|00+|11+\mathrm{e}^{2\pi \mathrm{i}/3}|02,++\mathrm{e}^{2\pi \mathrm{i}/3}|02,),\hfill \\ |02,\hfill & =3^{1/4}(|00+|11+\mathrm{e}^{2\pi \mathrm{i}/3}|02,++\mathrm{e}^{2\pi \mathrm{i}/3}|02,)\hfill \end{array}$$
(25)
as well as the single symmetry breaking boundary state
$$|01=\frac{1}{2}\mathrm{\hspace{0.17em}3}^{1/4}(\sqrt{3}+\sqrt{7})|00+\frac{1}{2}\mathrm{\hspace{0.17em}3}^{1/4}(\sqrt{3}\sqrt{7})|11.$$
(26)
We finally remark that the system of equations for the coefficients of the Ishibashi states is highly over-determined. We regard it as a non-trivial check of our ideas that a solution exists at all.
Acknowledgement:
We are grateful to J. Böckenhauer, D.E. Evans, G. Felder, J. Fröhlich and J.-B. Zuber for helpful discussions, and to P. Bantay and B. Schellekens for a careful reading of the manuscript. C.S. would like to thank the Schrödinger Institute for hospitality.
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# Deformation cohomology for Yetter-Drinfel’d modules and Hopf (bi)modules
## 1 Introduction
If $`A`$ is a bialgebra over a field $`k`$, a left-right Yetter-Drinfel’d module over $`A`$ is a $`k`$-linear space $`M`$ which is a left $`A`$-module, a right $`A`$-comodule and such that a certain compatibility condition between these two structures holds. Yetter-Drinfel’d modules were introduced by D. Yetter in under the name of “crossed bimodules” (they are called “quantum Yang-Baxter modules” in ; the present name is taken from ). If $`A`$ is a finite dimensional Hopf algebra then the category of left-right Yetter-Drinfel’d modules is equivalent to the category of left modules over $`D(A)`$, the Drinfel’d double of $`A`$ (see , ), even as braided tensor categories, and also to the category of Hopf bimodules over $`A`$ (see , , , ). An important class of examples occurs as follows: if $`M`$ is a finite dimensional vector space and $`REnd(MM)`$ is a solution to the quantum Yang-Baxter equation, then the so-called “FRT construction” (see for instance ) associates to $`R`$ a certain bialgebra $`A(R)`$, and $`M`$ becomes a left-right Yetter-Drinfel’d module over $`A(R)`$ (see , ).
In this paper we introduce a cohomology theory for left-right Yetter-Drinfel’d modules. If $`A`$ is a bialgebra and $`M,N`$ are left-right Yetter-Drinfel’d modules over $`A`$, we construct a double complex $`\{Y^{n,p}(M,N)\}`$ whose total cohomology is the desired cohomology $`H^{}(M,N)`$. For $`M=N=k`$, this cohomology is just the Gerstenhaber-Schack cohomology of the bialgebra $`A`$. In general, we prove that $`H^0(M,N)`$ is $`Hom(M,N)`$ in the category of Yetter-Drinfel’d modules, and $`H^1(M,N)`$ is isomorphic to the group $`Ext^1(M,N)`$ of extensions of $`M`$ by $`N`$ in the category of Yetter-Drinfel’d modules; in particular, if $`A`$ is a finite dimensional Hopf algebra, this implies that $`H^1(M,N)Ext_{D(A)}^1(M,N)`$, where $`D(A)`$ is the Drinfel’d double of $`A`$, and we raise the problem whether $`H^n(M,N)Ext_{D(A)}^n(M,N)`$ for any $`n2`$.
Similarly, we construct a cohomology theory $`H^{}(M,N)`$ for $`M,N`$ being this time left-right Hopf modules over a bialgebra $`A`$, as the total cohomology of a certain double complex $`\{C^{n,p}(M,N)\}`$. We prove that this cohomology vanishes if $`M`$ and $`N`$ are of the form $`M=VA`$ and $`N=WA`$ for some linear spaces $`V,W`$ (in particular, this cohomology vanishes if the bialgebra $`A`$ has a skew antipode, since in this case any left-right Hopf module over $`A`$ is of this form). Finally, motivated by the recent work of R. Taillefer, we consider the case when $`M`$ and $`N`$ are not only left-right Hopf modules, but even Hopf bimodules, and we construct a subbicomplex $`\{T^{n,p}(M,N)\}`$ of the above $`\{C^{n,p}(M,N)\}`$, yielding a cohomology theory for Hopf bimodules, similar to the one introduced in . It is likely that these two cohomologies are isomorphic, at least when $`A`$ is a finite dimensional Hopf algebra.
Let us finally mention that our cohomology theories classify deformations of the corresponding structures, in the sense of Gerstenhaber’s deformation theory (see ).
Throughout, $`k`$ will be a fixed field and all linear spaces, algebras etc. will be over $`k`$. Unadorned $``$ and $`Hom`$ are also over $`k`$. If $`V`$ is a $`k`$-linear space and $`n`$ is a natural number, we denote $`V^n`$ by $`V^n`$. For bialgebras and Hopf algebras we refer to , ; we shall use Sweedler’s sigma notation $`\mathrm{\Delta }(a)=a_1a_2`$, $`\mathrm{\Delta }_2(a)=a_1a_2a_3`$ etc.
## 2 Cohomology for Yetter-Drinfel’d modules
Let $`A`$ be a bialgebra with multiplication $`\mu `$ and comultiplication $`\mathrm{\Delta }`$ and $`(M,\omega _M,\rho _M)`$ a left-right Yetter-Drinfel’d module over $`A`$, that is $`(M,\omega _M)`$ is a left $`A`$-module (we denote by $`\omega _M(am)=am`$ the left $`A`$-module structure of $`M`$), $`(M,\rho _M)`$ is a right $`A`$-comodule (we denote by $`\rho _M:MMA`$, $`\rho _M(m)=m_0m_1`$ the comodule structure of $`M`$), and the following compatibility condition holds:
$$(a_2m)_0(a_2m)_1a_1=a_1m_0a_2m_1$$
(1)
for all $`aA,mM`$. Let also $`(N,\omega _N,\rho _N)`$ be another left-right Yetter-Drinfel’d module, with the same kind of notation.
For any natural numbers $`n,p0`$, we denote
$$Y^{n,p}(M,N)=Hom(A^nM,NA^p)$$
If $`fY^{n,p}(M,N)`$, $`a^1,a^2,\mathrm{},a^nA`$, $`mM`$, we shall denote
$$f(a^1\mathrm{}a^nm)=f(a^1\mathrm{}a^nm)^0$$
$$f(a^1\mathrm{}a^nm)^1\mathrm{}f(a^1\mathrm{}a^nm)^p$$
For any $`n,p0`$ and for any $`i=0,1,\mathrm{},n+1`$, define $`b_i^{n,p}:Y^{n,p}(M,N)Y^{n+1,p}(M,N)`$, by:
$`b_0^{n,p}(f)(a^1\mathrm{}a^{n+1}m)=(a^1)_1f(a^2\mathrm{}a^{n+1}m)^0`$
$$(a^1)_2f(a^2\mathrm{}a^{n+1}m)^1\mathrm{}(a^1)_{p+1}f(a^2\mathrm{}a^{n+1}m)^p$$
$`b_i^{n,p}(f)(a^1\mathrm{}a^{n+1}m)=f(a^1\mathrm{}a^ia^{i+1}\mathrm{}a^{n+1}m)`$ for all $`1in`$
$`b_{n+1}^{n,p}(f)(a^1\mathrm{}a^{n+1}m)=f(a^1\mathrm{}a^n(a^{n+1})_{p+1}m)^0`$
$$f(a^1\mathrm{}a^n(a^{n+1})_{p+1}m)^1(a^{n+1})_1\mathrm{}f(a^1\mathrm{}a^n(a^{n+1})_{p+1}m)^p(a^{n+1})_p$$
Define now
$$d_m^{n,p}:Y^{n,p}(M,N)Y^{n+1,p}(M,N),d_m^{n,p}=\underset{i=0}{\overset{n+1}{}}(1)^ib_i^{n,p}$$
Then one can prove, case by case, that for any $`0i<jn+2`$ the following relation holds:
$$b_j^{n+1,p}b_i^{n,p}=b_i^{n+1,p}b_{j1}^{n,p}$$
and using this relation it follows that
$$d_m^{n+1,p}d_m^{n,p}=0$$
for all $`n,p0`$.
Now, for any $`n,p0`$ and for any $`i=0,1,\mathrm{},p+1`$, define $`c_i^{n,p}:Y^{n,p}(M,N)Y^{n,p+1}(M,N)`$, by:
$`c_0^{n,p}(f)(a^1\mathrm{}a^nm)=(f((a^1)_2\mathrm{}(a^n)_2m)^0)_0`$
$$(f((a^1)_2\mathrm{}(a^n)_2m)^0)_1(a^1)_1\mathrm{}(a^n)_1$$
$$f((a^1)_2\mathrm{}(a^n)_2m)^1\mathrm{}f((a^1)_2\mathrm{}(a^n)_2m)^p$$
$`c_i^{n,p}(f)(a^1\mathrm{}a^nm)=f(a^1\mathrm{}a^nm)^0f(a^1\mathrm{}a^nm)^1`$
$$\mathrm{}(f(a^1\mathrm{}a^nm)^i)_1(f(a^1\mathrm{}a^nm)^i)_2\mathrm{}$$
$$f(a^1\mathrm{}a^nm)^p$$
for all $`1ip`$
$`c_{p+1}^{n,p}(f)(a^1\mathrm{}a^nm)=f((a^1)_1\mathrm{}(a^n)_1m_0)(a^1)_2\mathrm{}(a^n)_2m_1`$
Define
$$d_c^{n,p}:Y^{n,p}(M,N)Y^{n,p+1}(M,N),d_c^{n,p}=\underset{i=0}{\overset{p+1}{}}(1)^ic_i^{n,p}$$
Then one can prove, case by case, that for any $`0i<jp+2`$ the following relation holds:
$$c_j^{n,p+1}c_i^{n,p}=c_i^{n,p+1}c_{j1}^{n,p}$$
and from this relation it follows that
$$d_c^{n,p+1}d_c^{n,p}=0$$
for all $`n,p0`$.
Also, one can prove, case by case, that for any $`0in+1`$ and $`0jp+1`$, the following relation holds:
$$c_j^{n+1,p}b_i^{n,p}=b_i^{n,p+1}c_j^{n,p}$$
Note that for the cases $`j=p+1,i=n+1`$ and $`j=0,i=0`$ one has to use the Yetter-Drinfel’d module condition (1), and these are the only two places where this condition is used.
Using this relation, it follows immediately that
$$d_c^{n+1,p}d_m^{n,p}=d_m^{n,p+1}d_c^{n,p}$$
for all $`n,p0`$.
In conclusion, $`(Y^{n,p}(M,N),d_m^{n,p},d_c^{n,p})`$ is a double complex. We shall denote by $`H^n(M,N)`$, for $`n0`$, the cohomology of the total complex associated to this double complex.
It is easy to see that $`H^0(M,N)`$ is the set of morphisms from $`M`$ to $`N`$ in the category of Yetter-Drinfel’d modules. Also, it is easy to see that $`H^1(M,N)=Z^1(M,N)/B^1(M,N)`$, where $`Z^1(M,N)`$ is the set of all pairs $`(\omega ^{},\rho ^{})Hom(AM,N)Hom(M,NA)`$ that satisfy the relations:
$``$ $`\omega ^{}(id\omega _M)+\omega _N(id\omega ^{})=\omega ^{}(\mu id)`$
$``$ $`(\rho ^{}id)\rho _M+(\rho _Nid)\rho ^{}=(id\mathrm{\Delta })\rho ^{}`$
$``$ $`(\omega ^{}\mu )(id\tau _Mid)(\mathrm{\Delta }\rho _M)+(\omega _N\mu )(id\tau _Nid)(\mathrm{\Delta }\rho ^{})=(id\mu )(\rho ^{}id)\tau _M(id\omega _M)(\mathrm{\Delta }id)+(id\mu )(\rho _Nid)\tau _N(id\omega ^{})(\mathrm{\Delta }id)`$
where $`\tau _M:AMMA`$, $`\tau _M(am)=ma`$ (and similarly for $`\tau _N`$), and $`B^1(M,N)`$ is the set of all pairs $`(d_m(f),d_c(f))Hom(AM,N)Hom(M,NA)`$, with $`fHom(M,N)`$, and
$``$ $`d_m(f)=\omega _N(idf)f\omega _M`$
$``$ $`d_c(f)=\rho _Nf(fid)\rho _M`$
Now, the category of Yetter-Drinfel’d modules is abelian, so we can consider the abelian group $`Ext^1(M,N)`$, which is the set of equivalence classes of extensions of $`M`$ by $`N`$ in the category of Yetter-Drinfel’d modules, the group law being the Baer sum (see , pp. 78-79).
If $`(\omega ^{},\rho ^{})Hom(AM,N)Hom(M,NA)`$, denote by $`N_{(\omega ^{},\rho ^{})}M`$ the $`k`$-linear space $`NM`$, endowed with a left multiplication and a right comultiplication, as follows:
$``$ $`a(n,m)=(an+\omega ^{}(am),am)`$
$``$ $`\lambda :NM(NM)A(NA)(MA)`$
$$\lambda ((n,m))=\rho _N(n)+\rho ^{}(m)+\rho _M(m)$$
Then one can check, by a direct computation, that $`N_{(\omega ^{},\rho ^{})}M`$ with these structures is a left-right Yetter-Drinfel’d module if and only if the pair $`(\omega ^{},\rho ^{})`$ belongs to $`Z^1(M,N)`$. Moreover, the sequence
$$0NN_{(\omega ^{},\rho ^{})}MM0$$
is an extension of $`M`$ by $`N`$ in the category of Yetter-Drinfel’d modules, and any extension of $`M`$ by $`N`$ is equivalent to one of this form. If $`(\omega ^{},\rho ^{}),(\omega ^{\prime \prime },\rho ^{\prime \prime })Z^1(M,N)`$, then one can also check that the two extensions determined by these pairs are equivalent if and only if $`(\omega ^{},\rho ^{})(\omega ^{\prime \prime },\rho ^{\prime \prime })B^1(M,N)`$. So, we have a bijection
$$H^1(M,N)Ext^1(M,N)$$
and one can prove that this is actually a group isomorphism. In conclusion, we have the following
###### Proposition 1
The groups $`H^1(M,N)`$ and $`Ext^1(M,N)`$ are isomorphic.
In particular, if $`A`$ is a finite dimensional Hopf algebra, it is well-known that the category of left-right Yetter-Drinfel’d modules over $`A`$ is isomorphic to the category of left modules over the Drinfel’d double of $`A`$ (see , ), so in this case we have a group isomorphism $`H^1(M,N)Ext_{D(A)}^1(M,N)`$. It is natural to ask the following
Question: Is it true that $`H^n(M,N)Ext_{D(A)}^n(M,N)`$ for any $`n2`$?
###### Example 2
Let $`M=N=k`$ with trivial Yetter-Drinfel’d module structure over the bialgebra $`A`$. In this case the double complex $`(Y^{n,p}(M,N),d_m^{n,p},d_c^{n,p})`$ becomes:
$``$ $`Y^{n,p}(k,k)=Hom(A^n,A^p)`$ for all $`n,p0`$
$``$ $`b_0^{n,p}(f)(a^1\mathrm{}a^{n+1})=a^1f(a^2\mathrm{}a^{n+1})`$
$``$ $`b_i^{n,p}(f)(a^1\mathrm{}a^{n+1})=f(a^1\mathrm{}a^ia^{i+1}\mathrm{}a^{n+1})`$
for all $`1in`$
$``$ $`b_{n+1}^{n,p}(f)(a^1\mathrm{}a^{n+1})=f(a^1\mathrm{}a^n)a^{n+1}`$
where the dots represent the canonical (diagonal) left and right $`A`$-module structures on $`A^p`$
$``$ $`c_0^{n,p}(f)(a^1\mathrm{}a^n)=(a^1)_1\mathrm{}(a^n)_1f((a^1)_2\mathrm{}(a^n)_2)`$
$``$ $`c_i^{n,p}(f)(a^1\mathrm{}a^n)=(idid\mathrm{}\mathrm{\Delta }id\mathrm{}id)(f(a^1\mathrm{}a^n))`$
for all $`1ip`$, where $`\mathrm{\Delta }`$ is applied on the $`i^{th}`$ position
$``$ $`c_{p+1}^{n,p}(f)(a^1\mathrm{}a^n)=f((a^1)_1\mathrm{}(a^n)_1)(a^1)_2\mathrm{}(a^n)_2`$
and this is the double complex that gives the Gerstenhaber-Schack cohomology $`H_b^{}(A,A)`$ of the bialgebra $`A`$ (see , ).
###### Remark 3
A positive answer to the above question would imply that for a finite dimensional Hopf algebra $`A`$ we have $`H_b^{}(A,A)Ext_{D(A)}^{}(k,k)`$, and this would give another proof for the vanishing of the “hat” Gerstenhaber-Schack cohomology of a semisimple cosemisimple Hopf algebra (let us note that the original proof in uses also the Drinfel’d double of $`A`$).
###### Remark 4
If $`A`$ is finite dimensional and the field $`k`$ is algebraically closed, there exists a geometric interpretation of $`H^1(M,M)`$ for any finite dimensional Yetter-Drinfel’d module $`M`$ (actually, it is this geometric approach who suggested how to define $`H^1(M,M)`$), similar to the one for the “hat” Gerstenhaber-Schack cohomology given in , that we shall now briefly describe (the proofs are similar to the ones in ). Let $`M`$ be a finite dimensional $`k`$-linear space, consider the affine algebraic variety $`𝒜`$$`=Hom(AM,M)\times Hom(M,MA)`$, and define $`𝒴𝒟`$$`(M)`$ to be the set of all pairs $`(\omega ,\rho )𝒜`$ such that $`(M,\omega ,\rho )`$ is a left-right Yetter-Drinfel’d module over $`A`$. Since the Yetter-Drinfel’d conditions are polynomial, $`𝒴𝒟`$$`(M)`$ is a subvariety of $`𝒜`$. Then, if we take a Yetter-Drinfel’d module $`(M,\omega ,\rho )`$ (that is, a point $`(\omega ,\rho )𝒴𝒟`$$`(M)`$), one can prove that the tangent space $`T_{(\omega ,\rho )}(𝒴𝒟`$$`(M))`$ is contained in $`Z^1(M,M)`$.
Now, the algebraic group $`G=GL(M)`$ acts on $`𝒴𝒟`$$`(M)`$ by transport of structures, that is $`g(\omega ,\rho )=(\omega ^g,\rho ^g)`$, where $`\omega ^g=g\omega (id_Ag^1)`$ and $`\rho ^g=(gid_A)\rho g^1`$ (and obviously there is a bijection between orbits and isomorphism classes of Yetter-Drinfel’d module structures on $`M`$). If we fix $`(\omega ,\rho )𝒴𝒟`$$`(M)`$ and we denote by $`\overline{Orb_{(\omega ,\rho )}}`$ the closure of the orbit through $`(\omega ,\rho )`$, then one can prove that $`B^1(M,M)`$ is contained in the tangent space $`T_{(\omega ,\rho )}(\overline{Orb_{(\omega ,\rho )}})`$.
## 3 Cohomology for Hopf modules
Let $`A`$ be a bialgebra with multiplication $`\mu `$ and comultiplication $`\mathrm{\Delta }`$ and $`(M,\omega _M,\rho _M)`$ a left-right Hopf module over $`A`$, that is $`(M,\omega _M)`$ is a left $`A`$-module (with notation $`\omega _M(am)=am`$), $`(M,\rho _M)`$ is a right $`A`$-comodule (with notation $`\rho _M:MMA`$, $`\rho _M(m)=m_0m_1`$) and the following compatibility condition holds:
$$(am)_0(am)_1=a_1m_0a_2m_1$$
(2)
for all $`aA,mM`$. Let $`(N,\omega _N,\rho _N)`$ be another left-right Hopf module over $`A`$. We denote by $`C^{n,p}(M,N)=Hom(A^nM,NA^p)`$ and for $`fC^{n,p}(M,N)`$ we use the same notation as in the previous section for $`f(a^1\mathrm{}a^nm)`$. For any $`i=0,1,\mathrm{},n+1`$, define $`b_i^{n,p}:C^{n,p}(M,N)C^{n+1,p}(M,N)`$, by:
$`b_0^{n,p}(f)(a^1\mathrm{}a^{n+1}m)=(a^1)_1f(a^2\mathrm{}a^{n+1}m)^0`$
$$(a^1)_2f(a^2\mathrm{}a^{n+1}m)^1\mathrm{}(a^1)_{p+1}f(a^2\mathrm{}a^{n+1}m)^p$$
$`b_i^{n,p}(f)(a^1\mathrm{}a^{n+1}m)=f(a^1\mathrm{}a^ia^{i+1}\mathrm{}a^{n+1}m)`$
for all $`1in`$
$`b_{n+1}^{n,p}(f)(a^1\mathrm{}a^{n+1}m)=f(a^1\mathrm{}a^na^{n+1}m)`$
For any $`i=0,1,\mathrm{},p+1`$, define $`c_i^{n,p}:C^{n,p}(M,N)C^{n,p+1}(M,N)`$, by:
$`c_0^{n,p}(f)(a^1\mathrm{}a^nm)=(f(a^1\mathrm{}a^nm)^0)_0(f(a^1\mathrm{}a^nm)^0)_1`$
$$f(a^1\mathrm{}a^nm)^1\mathrm{}f(a^1\mathrm{}a^nm)^p$$
$`c_i^{n,p}(f)(a^1\mathrm{}a^nm)=f(a^1\mathrm{}a^nm)^0f(a^1\mathrm{}a^nm)^1`$
$$\mathrm{}(f(a^1\mathrm{}a^nm)^i)_1(f(a^1\mathrm{}a^nm)^i)_2\mathrm{}f(a^1\mathrm{}a^nm)^p$$
for all $`1ip`$
$`c_{p+1}^{n,p}(f)(a^1\mathrm{}a^nm)=f((a^1)_1\mathrm{}(a^n)_1m_0)(a^1)_2(a^2)_2\mathrm{}(a^n)_2m_1`$
Then one can prove, as in the previous section, that $`(C^{n,p}(M,N),d_m^{n,p},d_c^{n,p})`$ is a double complex, where $`d_m^{n,p}`$ and $`d_c^{n,p}`$ are defined by the same formulae as in the previous section. We shall denote by $`H^n(M,N)`$, for $`n0`$, the cohomology of the total complex associated to this double complex.
It is easy to see that $`H^0(M,N)`$ is the set of morphisms from $`M`$ to $`N`$ in the category of left-right Hopf modules over $`A`$. Also, it is easy to see that $`H^1(M,N)=Z^1(M,N)/B^1(M,N)`$, where $`Z^1(M,N)`$ is the set of all pairs $`(\omega ^{},\rho ^{})Hom(AM,N)Hom(M,NA)`$ that satisfy the relations:
$`\omega ^{}(id\omega _M)+\omega _N(id\omega ^{})=\omega ^{}(\mu id)`$
$`(\rho ^{}id)\rho _M+(\rho _Nid)\rho ^{}=(id\mathrm{\Delta })\rho ^{}`$
$`(\omega ^{}\mu )(id\tau _Mid)(\mathrm{\Delta }\rho _M)+(\omega _N\mu )(id\tau _Nid)(\mathrm{\Delta }\rho ^{})=\rho ^{}\omega _M+\rho _N\omega ^{}`$
where $`\tau _M:AMMA`$, $`\tau _M(am)=ma`$ (and similarly for $`\tau _N`$), and $`B^1(M,N)`$ is the set of all pairs $`(d_m(f),d_c(f))Hom(AM,N)Hom(M,NA)`$, with $`fHom(M,N)`$ and
$`d_m(f)=\omega _N(idf)f\omega _M`$
$`d_c(f)=\rho _Nf(fid)\rho _M`$
As in the previous section, one can prove that $`H^1(M,N)`$ is isomorphic to the group $`Ext^1(M,N)`$ of extensions of $`M`$ by $`N`$ in the category of left-right Hopf modules.
Let $`A`$ be a bialgebra and $`V`$ a $`k`$-linear space. Then $`M=VA`$ becomes a left-right Hopf module over $`A`$, with module structure $`a(vb)=vab`$ for all $`a,bA,vV`$, and with comodule structure $`\rho :VAVAA`$, $`\rho (va)=va_1a_2`$. Let also $`W`$ be a $`k`$-linear space and consider the left-right Hopf module $`N=WA`$. One can check that in this case all the rows and the columns of the double complex corresponding to $`M`$ and $`N`$ are acyclic. Indeed, if $`gKer(d_m^{n+1,p})`$, then $`g=d_m^{n,p}(f)`$, where $`f:A^nVAWAA^p`$,
$$f(a^1\mathrm{}a^nva)=(1)^{n+1}g(a^1\mathrm{}a^nav1)$$
and if $`gKer(d_c^{n,p+1})`$, then $`g=d_c^{n,p}(f)`$, where $`f:A^nVAWAA^p`$,
$$f(a^1\mathrm{}a^nva)=(id_W\epsilon id_A^{p+1})(g(a^1\mathrm{}a^nva))$$
Since the rows of the double complex $`\{C^{n,p}(M,N)\}`$ are acyclic, if we consider the double complex $`\{D^{n,p}\}`$ obtained by adding to $`\{C^{n,p}(M,N)\}`$ one more column consisting of $`\{Ker(d_m^{0,p})\}`$ for all $`p0`$, then by the “Acyclic Assembly Lemma” (see , p. 59) the total complex associated to $`\{D^{n,p}\}`$ is acyclic. Then, a long exact sequence argument shows that the total cohomology of $`\{C^{n,p}(M,N)\}`$ may be computed as the cohomology of the added column, that is
$$H^{p+1}(VA,WA)=Ker(d_m^{0,p+1})Ker(d_c^{0,p+1})/d_c^{0,p}(Ker(d_m^{0,p}))$$
###### Proposition 5
$`H^{p+1}(VA,WA)=0`$ for all $`p0`$.
Proof: Let $`gKer(d_m^{0,p+1})Ker(d_c^{0,p+1})`$. Since $`gKer(d_c^{0,p+1})`$, there exists $`fC^{0,p}(VA,WA)`$ such that $`d_c^{0,p}(f)=g`$, namely $`f=(id_W\epsilon id_A^{p+1})g`$. It will be enough to prove that $`fKer(d_m^{0,p})`$. Since $`gKer(d_m^{0,p+1})`$, it follows that
$$g(vab)=g_W(vb)a_1g(vb)^0a_2g(vb)^1\mathrm{}a_{p+2}g(vb)^{p+1}$$
for all $`a,bA`$ and $`vV`$, where we denoted
$$g(vb)=g_W(vb)g(vb)^0\mathrm{}g(vb)^{p+1}WAA^{p+1}$$
By applying $`id_W\epsilon id_A^{p+1}`$ we obtain
$$(id_W\epsilon id_A^{p+1})(g(vab))=$$
$$\epsilon (g(vb)^0)g_W(vb)a_1g(vb)^1\mathrm{}a_{p+1}g(vb)^{p+1}$$
that is
$$f(vab)=a_1f(vb)^0a_2f(vb)^1\mathrm{}a_{p+1}f(vb)^p$$
which means that $`d_m^{0,p}(f)=0`$, q.e.d.
Suppose that the bialgebra $`A`$ has a skew antipode. In this case, it is very well-known that any left-right Hopf module over $`A`$ is of the form $`M=VA`$, for some $`k`$-linear space $`V`$ (see, for instance, , p.16). Hence we have obtained the following
###### Proposition 6
If $`A`$ is a bialgebra with a skew antipode (for example, a Hopf algebra with bijective antipode) then for any left-right Hopf modules $`M`$ and $`N`$ over $`A`$ and for any natural number $`n1`$ we have $`H^n(M,N)=0`$.
Let us introduce some notation. If $`A`$ is a bialgebra, we denote by $`modA`$ the category of right $`A`$-modules, by $`Acomod`$ the category of left $`A`$-comodules, by $`A_{lr}^r`$ the category whose objects are $`k`$-linear spaces $`M`$ which are bimodules, left-right Hopf modules and right-right Hopf modules over $`A`$, by $`A_l^{lr}`$ the category whose objects are $`k`$-linear spaces $`M`$ which are bicomodules, left-left Hopf modules and left-right Hopf modules over $`A`$, and finally by $`A_{lr}^{lr}`$ the category of Hopf bimodules (or two-sided two-cosided Hopf modules in the terminology of ) over $`A`$. We shall see now how the above double complex $`\{C^{n,p}\}`$ yields very naturally some cohomology theories for the categories $`A_{lr}^r`$, $`A_l^{lr}`$ and $`A_{lr}^{lr}`$.
Let $`M,NA_{lr}^r`$ and $`n,p`$ some natural numbers. Then $`A^nM`$ becomes a right $`A`$-module with structure
$$(a^1\mathrm{}a^nm)b=a^1\mathrm{}a^nmb$$
and $`NA^p`$ becomes a right $`A`$-module with structure
$$(xa^1\mathrm{}a^p)b=xb_1a^1b_2\mathrm{}a^pb_{p+1}$$
Now, if we define $`R^{n,p}(M,N)=Hom_{modA}(A^nM,NA^p)`$, then one can check, by a direct computation, that $`d_m^{n,p}(R^{n,p}(M,N))R^{n+1,p}(M,N)`$ and $`d_c^{n,p}(R^{n,p}(M,N))R^{n,p+1}(M,N)`$, so that $`(R^{n,p}(M,N),d_m^{n,p},d_c^{n,p})`$ is a double complex, giving a cohomology theory for objects in $`A_{lr}^r`$.
Similarly, if $`M,NA_l^{lr}`$ and $`n,p`$ are natural numbers, then $`A^nM`$ becomes a left $`A`$-comodule, with structure
$$A^nMAA^nM$$
$$a^1\mathrm{}a^nm(a^1)_1\mathrm{}(a^n)_1m_{(1)}(a^1)_2\mathrm{}(a^n)_2m_{(0)}$$
where we denoted by $`mm_{(1)}m_{(0)}`$ the left $`A`$-comodule structure of $`M`$, and $`NA^p`$ becomes a left $`A`$-comodule, with structure
$$NA^pANA^p$$
$$xa^1\mathrm{}a^px_{(1)}x_{(0)}a^1\mathrm{}a^p$$
If we denote by $`L^{n,p}(M,N)=Hom^{Acomod}(A^nM,NA^p)`$, then one can check also by a direct computation that $`d_m^{n,p}(L^{n,p}(M,N))L^{n+1,p}(M,N)`$ and $`d_c^{n,p}(L^{n,p}(M,N))L^{n,p+1}(M,N)`$, so that $`(L^{n,p}(M,N),d_m^{n,p},d_c^{n,p})`$ is a double complex, yielding a cohomology theory for objects in $`A_l^{lr}`$.
Finally, if $`M,NA_{lr}^{lr}`$, then on $`A^nM`$ and $`NA^p`$ we can introduce all the above right $`A`$-module and left $`A`$-comodule structures, so if we denote by
$$T^{n,p}(M,N)=Hom_{modA}^{Acomod}(A^nM,NA^p)$$
then $`d_m^{n,p}(T^{n,p}(M,N))T^{n+1,p}(M,N)`$ and $`d_c^{n,p}(T^{n,p}(M,N))T^{n,p+1}(M,N)`$, hence $`(T^{n,p}(M,N),d_m^{n,p},d_c^{n,p})`$ is a double complex, which yields a cohomology theory for Hopf bimodules. Note that $`\{T^{n,p}(M,N)\}`$ is a sort of “mirror” version of the double complex for Hopf bimodules introduced by R. Taillefer in . It is likely that the cohomologies given by these double complexes are isomorphic, at least when $`A`$ is a finite dimensional Hopf algebra.
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# Creation and Decay of 𝜂-Mesic Nuclei
## Abstract
First experimental results on photoproduction of $`\eta `$-mesic nuclei are analyzed. In an experiment performed at the 1 GeV electron synchrotron of the Lebedev Physical Institute, correlated $`\pi ^+n`$ pairs arising from the reaction
$$\gamma +{}_{}{}^{12}\mathrm{C}N+{}_{\eta }{}^{}(A1)N+\pi ^++n+(A2)$$
and flying transversely to the photon beam have been observed. When the photon energy exceeds the $`\eta `$-meson production threshold, a distribution of the $`\pi ^+n`$ pairs over their total energy is found to have a peak in the subthreshold region of the internal-conversion process $`\eta p\pi ^+n`$ which signals about formation of $`\eta `$-mesic nuclei.
The idea that a bound state of the $`\eta `$-meson and a nucleus (the so-called $`\eta `$-mesic nucleus) can exist in Nature was put forward long ago by Peng pen85 who relied on the first estimates of the $`\eta N`$ scattering length $`a_{\eta N}`$ obtained by Bhalerao and Liu bha85 . Owing to $`\mathrm{Re}a_{\eta N}>0`$, an average attractive potential exists between slow $`\eta `$ and nucleons. This can result in binding $`\eta A`$ systems, provided the life time of $`\eta `$ in nuclei is long enough liu86 . Modern calculations bat98 ; gre97 predict a rather strong $`\eta N`$ attraction which is sufficient for binding $`\eta `$ in all nuclei with $`A4`$.
The very first experiments on searching for the $`\eta `$-mesic nuclei performed at BNL chr88 and LAMPF lie88 gave negative results. Meantime, studies of the reactions $`p(d,{}_{}{}^{3}\mathrm{He})\eta `$ ber88 ; may96 , $`{}_{}{}^{18}\mathrm{O}(\pi ^+,\pi ^{}){}_{}{}^{18}\mathrm{Ne}`$ joh93 , and $`\stackrel{}{d}(d,{}_{}{}^{4}\mathrm{He})\eta `$ wil97 suggest that a quasi-bound $`\eta A`$ state is formed in these reactions wil93 ; kon94 .
In the present work we report on first results concerning formation of $`\eta `$-mesic nuclei in photoreactions. A very efficient trigger for searching for $`\eta `$-mesic nuclei sok91 consists in detecting decay products of the $`\eta `$-mesic nucleus, viz. a $`\pi N`$ pair produced in the reaction $`\eta N\pi N`$ inside the nucleus. Here $`\eta `$ itself is produced at an earlier stage, in the reaction $`\gamma N\eta N`$ in our case. Both these reactions are mediated by the $`S_{11}(1535)`$ resonance which affects also a propagation of the intermediate $`\eta `$ in the medium (via multiple $`\eta N`$ rescattering) and leads to capturing slow $`\eta `$ into a bound state (Fig. 1). Formation of the bound state of the $`\eta `$ and the nucleus becomes possible when the momentum of the produced $`\eta `$ is small (typically less than 150 MeV/c). This requirement suggests photon energies $`E_\gamma =650850`$ MeV as most suitable for creating $`\eta `$-nuclei.
$`\pi N`$ pairs emerging from $`\eta `$-mesic nucleus decays have an opening angle $`\theta _{\pi N}=180^{}`$ and specific kinetic energies of their components (though smeared by the Fermi motion), $`E_\pi 300`$ MeV, $`E_n100`$ MeV. Among four possible isotopic combinations $`\pi ^+n`$, $`\pi ^{}p`$, $`\pi ^0n`$, $`\pi ^0p`$ the first one is quite suitable for measuring energies of the particles.
Accordingly, in an experiment performed at the 1 GeV electron synchrotron of the Lebedev Physical Institute, correlated $`\pi ^+n`$ pairs arising from the reaction
$$\gamma +{}_{}{}^{12}\mathrm{C}N+{}_{\eta }{}^{}(A1)N+\pi ^++n+(A2)$$
(1)
have been searched for. An experimental setup (Fig. 2) consisted of a carbon target $`\mathrm{}\mathrm{\hspace{0.17em}4}\mathrm{cm}\times 4\mathrm{cm}`$ and two time-of-flight scintillator spectrometers having a time resolution of $`\delta \tau 0.1`$ ns. A plastic anticounter A of charged particles (of the 90% efficiency), placed in front of the neutron detectors, and $`dE/dx`$ layers, placed between start and stop detectors in the pion spectrometer, were used for a better identification of particles.
Strategy of measurements was as follows. Two bremsstrahlung-beam energies were used, $`E_{\gamma \mathrm{max}}=650`$ MeV and 850 MeV, i.e., well below and well above $`\eta `$ production threshold on the free nucleon which is 707 MeV. The first, “calibration” run was performed at 650 MeV with the spectrometers positioned at angles $`\theta _n=\theta _\pi =50^{}`$ around the beam. In that run, this was a quasi-free photoproduction $`\gamma p\pi ^+n`$ which dominated the observed yield of the $`\pi ^+n`$ pairs. Then, at the same “low” energy 650 MeV, the spectrometers were positioned at $`\theta _n=\theta _\pi =90^{}`$ (the “background” run). In such a kinematics, the quasi-free production did not contribute and the observed count was presumably dominated by double-pion production. At last, the third run (the “effect$`+`$background” run) was performed at the same $`90^{}/90^{}`$ position, however with the higher photon beam energy of 850 MeV, at which $`\eta `$ mesons are produced too.
In accordance with measured velocities of particles detected by the spectrometers, all candidates to the $`\pi ^+n`$ events were separated into three classes: fast-fast (FF), fast-slow (FS), and slow-slow (SS) events. The FF events mostly correspond to $`\pi ^0\pi ^0`$ production which results in hitting detectors by photons or $`e^+/e^{}`$. The FS events mostly emerge from the $`\pi ^+n`$ pairs. Comparing yields and time spectra in these runs (and, in particular, using the SS events for extrapolating and subtracting a background), we have found a clear excess of the FS events which appeared when the photon energy exceeded $`\eta `$ production threshold. The total cross section of photoproduction of such excess pairs, averaged over the photon energy range $`650850`$ MeV, was found to be about $`\sigma _{\mathrm{tot}}(\pi ^+n)10`$ $`\mu `$b. See Ref. sok98 for more details. In the present work a further analysis of the excess FS events is done and their energy characteristics are determined.
In order to find kinetic energies of the neutron and pion, the velocities $`\beta _i=L_i/ct_i`$ of both the particles have to be determined. They are subject to fluctuations stemming from errors $`\delta t_i`$ and $`\delta L_i`$ in the time-of-flight $`t_i`$ and the flight base $`L_i`$. Such fluctuations are clearly seen in the case of the ultra-relativistic FF events which have experimentally observed velocities close but not equal to 1 (see Fig. 3). Therefore, an experimental $`\beta `$-resolution of the setup can be directly inferred from the FF events. Then, using this information and applying an inverse-problem statistical method described in Ref. pav83 , one can unfold the experimental spectrum, obtain a smooth velocity distribution in the physical region $`\beta _i<1`$ (Fig. 4), and eventually find a distribution of the particle’s kinetic energies $`E_i=M_i[(1\beta _i^2)^{1/2}1]`$. Finding $`E_i`$, we introduced corrections related with average energy losses of particles in absorbers and in the detector matter. It is worth to say that the number of the $`\pi ^+n`$ FS events visibly increases when the photon beam energy becomes sufficient for producing $`\eta `$ mesons.
Of the most interest is the distribution of the $`\pi ^+n`$ events over their total energy $`E_{\mathrm{tot}}=E_n+E_\pi `$, because creation and decay of $`\eta `$-mesic nuclei is expected to produce a relatively narrow peak in $`E_{\mathrm{tot}}`$ of the width $`5070`$ MeV (see, e.g., sok98 ; lvo98 ). Such a peak was indeed observed: see Fig. 5, in which an excess of the FS events appears when the photon energy exceeds the $`\eta `$-production threshold. Subtracting a smooth background, we have found a 1-dimensional energy distribution of the $`\pi ^+n`$ events presumably coming from (bound) $`\eta `$ decaying in the nucleus, see Fig. 6. The experimental width of this distribution is about 100 MeV, including the apparatus resolution. Its center lies by $`\mathrm{\Delta }E=40`$ MeV below the energy excess $`m_\eta m_\pi =408`$ MeV in the reaction $`\eta N\pi N`$, and it is well below the position of the $`S_{11}(1535)`$ resonance too. Up to effects of binding of protons annihilated in the decay subprocess $`\eta p\pi ^+n`$, the value $`\mathrm{\Delta }E`$ characterizes the binding energy of $`\eta `$ in the nucleus. The width of that peak is determined both by the width of the $`\eta `$-bound state and by the Fermi motion.
Whereas the fixed opening angle $`\theta _{\pi n}=180^{}`$ chosen in the kinematics with $`\theta _n=\theta _\pi =90^{}`$ selects $`\pi ^+n`$ pairs carrying a low total momentum in the direction of the photon beam, an independent check of the transverse momentum $`p_{}=p_\pi p_n`$ is meaningful. The corresponding distribution is shown in Fig. 7. On the top of a background, there is a narrower peak in $`p_{}`$ having a width compatible with the Fermi momentum of nucleons in the nucleus.
In conclusion, an excess of correlated $`\pi ^+n`$ pairs with the opening angle $`\theta _{\pi N}=180^{}`$ has been experimentally observed when the energy of photons exceeded the $`\eta `$-production threshold. A distribution of the pairs over their total kinetic energy was found to have a peak lying below threshold of the elementary process $`\pi N\eta N`$. A narrow peak is also found in the pair’s distribution over their total transverse momentum. All that suggests that these $`\pi ^+n`$ pairs arise from creation and decay of captured bound $`\eta `$ in the nucleus, i.e., they arise through the stage of formation of an $`\eta `$-mesic nucleus.
The work was supported by RFBR grant 99-02-18224.
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# Antiferromagnetism, Stripes, and Superconductivity in the t-J Model with Coulomb Interaction
## Abstract
We study mean-field phases of the t-J model with long-range Coulomb interaction. In the order of increasing doping density we find a classical antiferromagnet, charge and spin stripes, and a uniform $`d`$-wave superconductor, at the realistic doping parameters. Both in-phase and anti-phase stripes exist as metastable configurations, but the in-phase stripes have a slightly lower energy. The dependence of the stripe width and the inter-stripe spacing on the doping is examined. Effects of fluctuations around the mean-field states are discussed.
The cuprate materials which exhibit high-T<sub>c</sub> superconductivity show many different ordering tendencies as the hole doping concentration ($`x`$) in the system is varied. At zero doping the cuprates are antiferromagnetic insulators below the Néel temperature (Fig. 1(a)). The ordered antiferromagnetic phase ceases to exist with more than about $`2\%`$ of holes. Around $`5\%`$ doping the system begins to show superconductivity at low temperatures. In the intervening $`25`$% the system is lacking either antiferromagnetic order or superconductivity. In the doping range between $`5\%`$ and $`15\%`$, superconductivity co-exists with the (dynamic) stripe order. The stripes are most readily seen in neutron scattering experiments at non-zero energy transfer, hence “dynamic”, where one finds evidence of one-dimensional periodic modulation of the antiferromagnetic order as well as the charge density.
As noted from the early days of high-T<sub>c</sub>, the cuprates are characterized by strong electron-electron repulsion, which gives rise to a Mott-insulating state at half-filling. The t-J model, which incorporates such electron repulson, has been extensively studied in connection with high-T<sub>c</sub> superconductivity. Although it is quite likely that the t-J model correctly captures the short-distance correlation of electrons in the cuprates, one must not a priori overlook the long-range part of Coulomb interaction. In particular for the t-J model without Coulomb interaction, phase separation occurs for a wide range of doping concentration, which pre-empts the possibility of a high-pairing-scale superconducting state.
Various mean-field theories which assumes a uniform ground state have been in existence. These theories have varying degrees of success in understanding the phases of the cuprates. Such mean-field approaches have however been subject to the skepticism that the no-double-occupancy constraint is treated only approximately. Attempts to improve the treatment of the constraint result in a strongly fluctuating gauge field. Recently one of us (D.-H.L.) was able to integrate out the gauge field exactly at low energies in the uniform superconducting phase. It was shown that despite the drastic modification of the excitations, mean-field vacua serves as a good starting point in characterizing the zero-temperature state.
This work is performed under the assumption that mean-field theory will capture the short-distance/high-energy ordering tendency of the t-J + Coulomb model, and that the long-distance/low-energy properties can be understood by studying the soft fluctuations of the mean-field order parameters.
The Hamiltonian we consider is the following:
$`H`$ $`=`$ $`t{\displaystyle \underset{ij}{}}(b_jb_i^{}f_{j\alpha }^{}f_{i\alpha }+h.c.)+J{\displaystyle \underset{ij}{}}(𝐒_i𝐒_j{\displaystyle \frac{1}{4}}n_in_j)`$ (1)
$`+`$ $`{\displaystyle \frac{V_c}{2}}{\displaystyle \underset{ij}{}}{\displaystyle \frac{1}{r_{ij}}}(n_i\overline{n})(n_j\overline{n}).`$ (2)
The no-double-occupancy constraint is expressed as $`f_{i\alpha }^{}f_{i\alpha }+b_i^{}b_i1=0`$ in terms of the spinon ($`f_{i\sigma }`$) and holon ($`b_i`$) operators. Other notations include $`n_i=1b_i^{}b_i`$ (site electron density), and $`\overline{n}=1x`$ (average electron density). Our first goal is to understand the mean-field phases sustained by this model. Unlike other mean-field theories in the past, we include the magnetic order parameter on an equal footing with all other order parameters. Among other things this gives us the possibility of obtaining the long-range ordered antiferromagnetic state observed in experiments.
Mean-field theory
Our mean-field theory, in essence, is a variational approach: a trial wavefunction is constructed and parameters are varied to obtain the minimum energy. We construct a wavefunction which allows local magnetic moments (but not limited to antiferromagnetism), superconducting pairing (but not limited to $`d`$-wave symmetry), and modulations in the charge density. The trial wavefunction $`|\mathrm{\Psi }`$ is given by $`|\mathrm{\Psi }=|\mathrm{\Psi }_b|\mathrm{\Psi }_f`$ where the bosonic $`(|\mathrm{\Psi }_b)`$ and fermionic ($`|\mathrm{\Psi }_f)`$ states are independently constructed from their respective vacua as follows:
$`|\mathrm{\Psi }_b=(\chi _jb_j^{})^{N_b}|0_b`$ (3)
$`|\mathrm{\Psi }_f={\displaystyle \underset{a,b}{}}\left(u_j^af_j^{}+v_j^af_j\right)\left(w_l^bf_l^{}+z_l^bf_l\right)|0_f.`$ (4)
Repeated indices $`j`$ and $`l`$ implies summation over lattice sites. The bosons are assumed to be condensed, and the fermion ground state is constructed by occupying the (yet undetermined) quasiparticle orbitals labeled by $`a,b`$. The mean-field single-particle wavefunctions $`u_j^a,v_j^a,w_l^b,z_l^b`$ and $`\chi _j`$ are varied to minimize
$`\mathrm{\Psi }|H|\mathrm{\Psi }`$ $``$ $`{\displaystyle \underset{i}{}}\lambda _i\mathrm{\Psi }|b_i^{}b_i+f_{i\alpha }^{}f_{i\alpha }1|\mathrm{\Psi }`$ (5)
$``$ $`\mu {\displaystyle \underset{i}{}}\mathrm{\Psi }|n_i\overline{n}|\mathrm{\Psi }.`$ (6)
Lagrange multipliers $`\lambda _i`$, and $`\mu `$ are introduced to guarantee that the average occupation obeys the constraints locally as well as globally.
The calculation is carried out numerically on a $`N_x\times N_y`$ lattice with a periodic boundary condition. Not assuming any translational invariance, we first perform totally unrestricted minimization for $`N_x,N_y16`$. After the nature of the solution is established, we perform a more restricted search for $`N_x,N_y`$ up to $`120`$. The results reported below are for $`t/J=3`$, $`V_c/J=5`$. Other choices of $`t/J`$ and $`V_c/J`$ are also studied, with results that are not qualitatively different from those presented below.
We find three prominent types of order. In the order of increasing doping, they are antiferromagnetic insulator, charge/spin stripes, and uniform $`d`$-wave superconductor.
Antiferromagnet (Fig. 1(a)): At zero doping the mean-field ground state shows antiferromagnetic long-range order. Each electron is surrounded by four neighbors with opposite spins. It is an insulator because the occupation constraint forbids the electrons to hop. This mean-field state is a caricature of the insulating antiferromagnet observed in the undoped cuprates. For $`x<x_{c1}0.02`$, the doped holes are localized, often in the form of elongated puddles (or finite-length stripes). However, these localized puddles do not disrupt the overall antiferromagnetic order.
Stripes (Fig. 1(b)-(c)): For $`x_{c1}<x<x_c0.14`$ the mean-field ground state shows charge corrugation in the form of stripes. A stripe is a region which is extended in one direction (say $`\widehat{y}`$) and localized in the other, with partially occupied sites. There are two types of stripes, anti-phase and in-phase. The antiferromagnetic order parameter goes through a $`\pi `$ phase shift across the anti-phase stripe, whereas it remains in-phase across the in-phase stripe. As a result the anti-phase stripe modulates the antiferromagnetic order with a period twice that for the charge density.
We find that the in-phase stripes are the ground states for $`x_{c1}<x<x_c`$. At $`x=x_c`$ a first order phase transition occurs, after which the system becomes uniform. In all cases we have studied the in-phase stripes are bond-centered and have a width of two lattice constants. The stripe spacing, on the other hand, depends on doping and increases as the doping level decreases. The smallest distance between the in-phase stripes we observe is four lattice constants, and it occurs near $`x_c`$. This spacing is maintained in the doping range $`0.095x0.14`$. For this range the site hole density in each stripe varies from $`0.19`$ per site (corresponding to a line density of $`0.38`$) at $`x=0.095`$ to $`0.28`$ (line density $`0.56`$) at $`x=0.14`$. As the doping decreases below $`0.095`$ a new stripe configuration with the stripe spacing equal to five lattice constants emerges as the ground state. The stripe width is still two. The evolution of the stripe spacing with the doping concentration is shown in Fig. 2. We infer from this the existence of a series of integer stripe spacings as doping decreases. Each spacing has a non-zero range of stability giving rise to plateaus in the modulation period of the charge density. The step-wise evolution of the stripe spacing is clearly a lattice commensuration effect. In the presence of thermal fluctuation of the shape of stripes (quantum fluctuation is known not to roughen the stripe shape), we expect a smoother evolution of the modulation period. Inside each stripe there exists strong superconducting pairing as well as weak magnetism, as illustrated in Fig. 1(b). The pairing gap inside the in-phase stripe is comparable to the maximal pairing gap ($`0.05J`$) observed in the uniform $`d`$-wave phase, while the magnetic moments are a fraction of the full moment, $`S_z=\pm 1/2`$, of the insulator. An in-phase stripe is in some sense an optimum pairing state which is confined in the one-dimensional geometry. As in-phase stripes get closely spaced, it is likely that transverse fluctuation smears out the charge corrugation and results in a high-pairing-scale superconductor.
In the entire range of $`x_{c1}xx_c`$ the anti-phase stripes are metastable mean-field solutions. The energy difference between the most favorable anti-phase and in-phase stripes is shown as function of doping in Fig. 3. The largest difference ($`10\%`$) occurs at $`x=0.025`$ and the smallest ($`0.14\%`$) at $`x_c`$. Note that anti-phase stripes come very close in energy to the in-phase ones near $`x_c`$, and therefore, fluctuations that are omitted by the mean-field theory may stabilize the anti-phase stripes. (One such candidate is the transverse fluctuation of the stripes.) When that happens, the progression of the ground states vs. doping will be according to Figs. 1(a) through 1(e). We discuss the properties of the anti-phase stripe in the following.
The anti-phase stripes are also bond-centered and have a width equal to two lattice constants. The stripe spacing evolves in a step-wise fashion similar to Fig. 2. The range of hole density inside the anti-phase stripe is consequently similar to the in-phase stripe case. The anti-phase stripes also have non-zero pairing and magnetic moments (Fig. 2(c)). However, the pairing scale is considerably smaller (by a factor of three) than that in the in-phase stripes. In this sense, anti-phase stripes are unfavorable as far as pairing is concerned.
Uniform $`d`$-wave superconductor (Fig. 1(d)-(e)): The homogeneous phase, $`x>x_c`$, is characterized by $`d`$-wave pairing and, for $`x`$ near $`x_c`$, some residual antiferromagnetism. The pairing scale is maximum at $`x=x_c`$ where $`\mathrm{\Delta }_{max}0.05J`$ and decreases monotonically as $`x`$ increases. The antiferromagnetic moments disappear completely for $`x0.2`$.
We find it gratifying that our mean-field theory produces states which extrapolate between extreme classical (antiferromagnet at $`x=0`$) and extreme quantum (uniform superconducting) limits.
Coulomb interaction and high pairing scale: In light of the above mean-field results, we argue that high-$`T_c`$ superconductivity is a cooperative effect due to the short-range magnetic exchange and the long-range Coulomb interaction. Larger antiferromagnetic exchange favors higher pairing scale, however it also causes phase separation to set in at a lower doping level and the uniform high-pairing state becomes inaccesssible. At $`t/J=3`$, our model shows a phase-separated ground state for $`x0.26`$ in the absence of Coulomb interaction. Coulomb interaction suppresses phase separation and produces two compromises – stripes and a high-pairing-scale $`d_{x^2y^2}`$ superconductor.
Fluctuations
The low-energy excitations often appear in the form of the fluctuation of the order parameters. In the present context these include: phase fluctuation ($`\theta _b`$) of the Bose condensate $`|\mathrm{\Psi }_b`$, phase fluctuation ($`\theta _p`$) in the pairing condensate of $`|\mathrm{\Psi }_f`$, orientation fluctuation of the magnetic moments ($`\widehat{\mathrm{\Omega }}`$), shape fluctuation and displacement of the stripes, gapless (neutral) fermion excitation in the case of $`d`$-wave pairing, and most importantly the “gauge fluctuation” inherent in the slave-boson theory.
Due to the occupancy constraint and the form of Eq. (2) there exists an internal gauge symmetry under a local phase change, $`b_ie^{i\theta _i}b_i,f_{i\alpha }e^{i\theta _i}f_{i\alpha }.`$ Such symmetry is broken by most of the mean-field vacua, giving rise to a fluctuating gauge field as the soft mode. Recently one of us were able to integrate out the gauge fluctuations exactly in the non-magnetic, uniform $`d`$-wave superconducting phase corresponding to Fig. 1(e). The result is the confinement of two Goldstone modes, $`\theta _b`$ and $`\theta _p`$, into one $`\varphi =2\theta _b\theta _p`$, which is the phase of the electron superconducting condensate. The final low-energy dynamics is that of a phase-fluctuating superconductor with gapless fermion excitations: $`=_\varphi +_\psi +_{int}`$, where
$`_\varphi ={\displaystyle \frac{K}{2}}(\varphi )^2+{\displaystyle \frac{1}{2u}}(_t\varphi )^2+i\overline{\rho }_0\varphi `$ (7)
$`_\psi ={\displaystyle \underset{n=1}{\overset{2}{}}}\overline{\psi }_{n\alpha }(_tiv_{xn}\tau _x_xiv_{yn}\tau _z_y)\psi _{n\alpha }`$ (8)
$`_{int}=iz_\mu j_{\psi \mu }_\mu \varphi .`$ (9)
In the above $`\overline{\rho }`$ is the average Cooper pair density, $`\psi _n`$ is the fermion field associated with the $`n`$th gap node, $`\tau _z`$ is the third Pauli matrix, $`v_{xn},v_{yn}`$ specify the linear dispersion of the gapless fermions, and $`j_{\psi \mu }=\frac{1}{2}(_n\overline{\psi }_{n\sigma }\tau _z\psi _{n\sigma },iv_{x1}\overline{\psi }_{1\sigma }\psi _{1\sigma },iv_{y2}\overline{\psi }_{2\sigma }\psi _{2\sigma })`$, is the fermion 3-current. Due to the gauge fluctuation the parameters $`K,u,z_\mu `$ are strongly renormalized. In particular $`K`$ is proportional to $`x`$, which accounts for the low superfluid density in spite of the high pairing scale in the pseudogap regime. Similar treatment of gauge fluctuations has been done for each of the mean-field phases discussed above. The results are somewhat technical and will be reported elsewhere.
For the antiferromagnet the mean-field state satisfies the occupation constraint exactly. In this phase the only low-energy degree of freedom is the fluctuation of the direction of the local spin. The spin degrees of freedom is gauge-neutral and hence unaffected by strong gauge fluctuations. The interaction between the spin waves is described by the familiar non-linear sigma model
$$_\mathrm{\Omega }=\frac{K_\sigma }{2}\left[\frac{1}{v^2}|_t\widehat{\mathrm{\Omega }}|^2+|\widehat{\mathrm{\Omega }}|^2\right].$$
(10)
In two space dimensions, it is well known that the spin-wave fluctuation does not destroy the antiferromagnetic long-range order as long as the spin stiffness $`K_\sigma `$ is sufficiently big. This certainly seems to be the case for the undoped cuprates.
In the uniform $`d`$-wave phase corresponding to Fig. 1(d) there exists residual antiferromagnetic moments. The low-energy degrees of freedom are those of the non-magnetic superconductor plus the spin fluctuation. Since the wavevector associated with the commensurate antiferromagnetic ordering, $`𝐤=(\pi ,\pi )`$, mismatches the momentum connecting the gap nodes, the magnetic degrees of freedom decouple from the low-energy fermions. The effective action is then simply the sum of Eq. (9) and Eq. (10). Due to the smallness of the magnetic moments in this phase, however, the quantum fluctuation is likely to wash out the long-range correlation of the residual magnetism. In that case, the distinction between magnetic and non-magnetic superconductors becomes obscure.
The fluctuations in the stripe phase is the richest. The low-energy degrees of freedom include the phase of the superconducting condensate and fermion quasiparticle excitations inside the stripes, the fluctuation of the magnetic moment, and the displacement and shape fluctuation of the stripes. Unlike the antiferromagnetic and the uniform superconducting phases, the fluctuation of stripes can not only modify the properties of a given stripe phase, but also change the energy ordering between the in-phase and anti-phase stripes. Despite some progress, a complete picture which involves all of the above fluctuations is still lacking. The subject is currently under investigation.
Acknowledgment We are indebted to Steve Kivelson for numerous helpful discussions. We also thank Eduardo Fradkin, Z.-X. Shen and Ned Wingreen for insightful remarks and questions. Part of the numerical calculations are carried out with the computing facility at NEC research. DHL is supported by NSF grant DMR 99-71503. QHW is supported by the National Natural Science Foundation of China, the National Centre for Research and Development of Superconductivity, and the Berkeley Scholars Program.
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# New Production Mechanism of Neutral Higgs Bosons with Right scalar tau neutrino as the LSP
## I Introduction
Supersymmetry (SUSY) has been largely studied as a possible framework for the theory beyond the Standard Model (SM). It provides a natural solution to the hierarchy problem, the generation of the electroweak symmetry breaking, as well as the grand unification of the gauge coupling constants . Among the various supersymmetric models, the minimal supersymmetric Standard Model (MSSM) is the one studied most extensively in the literature . In addition to the SM particles, it consists of the supersymmetric partners of the SM particles, called sparticles. All renormalizable interactions, including both SUSY conserving and (soft) SUSY breaking terms, are assumed to conserve the $`BL`$ global symmetry, which then results in the conservation of the $`R`$-parity, $`R=(1)^{3(BL)+2S}`$. ($`B`$ denotes the baryon number, $`L`$ the lepton number, and $`S`$ the spin of the (s)particle.) Accordingly, all the SM particles have even $`R`$-parity and the sparticles have odd $`R`$-parity. This fact has an important consequence, namely, if the initial particles of a scattering process are the SM particles, then sparticles can only be produced in pairs. $`R`$-parity conservation also implies the existence of a stable sparticle, called the lightest supersymmetric particle (LSP). The LSP is absolutely stable and cannot decay.
Recently, the Super-K neutrino experiment presented the high precision neutrino oscillation data which strongly suggests the existence of the neutrino mass . Many SUSY models have been proposed to account for a reasonable set of neutrino masses with bi-maximal mixing among the three family neutrinos . The low energy effective theory of these SUSY models can be summarized in the supersymmetric extension of the see-saw model, in which a right-handed neutrino superfield is added for each family in the framework of the MSSM . As proposed in , the observed bi-maximal mixing in neutrino data can be explained by the single right-handed neutrino dominance mechanism, which assumes that the light effective Majorana matrix come predominantly from a single right-handed neutrino which generates some particular textures of the Yukawa couplings in the superpotential of the low energy effective theory. With that example in mind, we assume the bi-maximal mixing among the three family neutrinos does not necessarily imply a large mixing among different flavor sneutrinos because the attendant R-parity conserving soft-supersymmetry breaking terms are not fixed by the generation of neutrino masses or mixings. Furthermore, the trilinear scalar coupling terms introduced by supersymmetry breaking can be large in some SUSY models. For simplicity, we shall only consider a one family model, though its phenomenology is expected to be applicable to three family models after properly including possible mixing factors.
Despited that the conventional supersymmetric see-saw models predict heavy right sneutrinos, we presume the existence of SUSY models, in which the interactions of right sneutrinos with left sneutrinos at the weak scale are described by Eqs.(7), (8) and (9) with R-parity conservation in the framework of MSSM, which may or may not include lepton number violation interactions. We study the scenario that the lightest sneutrino, mostly the right tau sneutrino $`\stackrel{~}{\nu }_{\tau _R}`$, is the LSP. To further simplify our discussion, we also assume that the mixings among the three generation sneutrinos are small enough that the dominant effect to collider phenomenology comes from the interaction of left and right tau sneutrinos. In section II, we give our assumptions and formalism for the $`\stackrel{~}{\nu }_{\tau _1}`$-LSP scenario. In section III, we discuss the possible large production rate of the lightest neutral Higgs boson predicted by this model at electron and hadron colliders. We also show that with a small left-right mixing, a $`\stackrel{~}{\nu }_{\tau _R}`$-like $`\stackrel{~}{\nu }_{\tau _1}`$ can be a good candidate for cold dark matter. In section IV, we consider in details some $`e^+e^{}`$ collider phenomenology for the $`\stackrel{~}{\nu }_{\tau _1}`$-LSP model. Section V contains our conclusion.
## II $`\stackrel{~}{\nu }_{\tau _1}`$-LSP scenario
In our one family model, the scalar tau neutrinos, like scalar quarks and scalar leptons, can mix and form mass eigenstates, and
$`\left(\begin{array}{c}\stackrel{~}{\nu }_{\tau _2}\\ \stackrel{~}{\nu }_{\tau _1}\end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cos}\theta _{\stackrel{~}{\nu }}& \mathrm{sin}\theta _{\stackrel{~}{\nu }}\\ \mathrm{sin}\theta _{\stackrel{~}{\nu }}& \mathrm{cos}\theta _{\stackrel{~}{\nu }}\end{array}\right)\left(\begin{array}{c}\stackrel{~}{\nu }_{\tau _L}\\ \stackrel{~}{\nu }_{\tau _R}^{}\end{array}\right),`$ (7)
with
$$\mathrm{tan}2\theta _{\stackrel{~}{\nu }}=\frac{2\mathrm{\Delta }m^2}{m_{\stackrel{~}{\nu }_{\tau _L}}^2m_{\stackrel{~}{\nu }_{\tau _R}}^2},$$
(8)
where $`\stackrel{~}{\nu }_{\tau _L}`$ and $`\stackrel{~}{\nu }_{\tau _R}`$ stand for the left and right tau sneutrinos, and $`m_{\stackrel{~}{\nu }_{\tau _L}}^2`$ and $`m_{\stackrel{~}{\nu }_{\tau _R}}^2`$ are the corresponding soft masses, respectively. We consider the case that the parameter $`\mathrm{\Delta }m^2`$ mainly comes from the soft SUSY breaking effect originated from the trilinear scalar couplings in
$$\mathrm{\Delta }=A_\text{v}H_2\stackrel{~}{L}_3\stackrel{~}{\nu }_{\tau _R}+C_\text{v}H_1^{}\stackrel{~}{L}_3\stackrel{~}{\nu }_{\tau _R}+\text{c.c.}$$
(9)
where $`A_\text{v}(C_\text{v})`$ is the standard (non-standard) trilinear scalar coupling, $`H_1`$ and $`H_2`$ are the two Higgs doublets and $`\stackrel{~}{L}_3`$ is the third generation scalar lepton doublet. (To simplify our discussion, we have assumed that the contribution to $`\mathrm{\Delta }m^2`$ from the $`\mu `$-term in the superpotential is much smaller than that from $`\mathrm{\Delta }`$, because of the small Yukawa coupling of tau neutrino.) The parameter $`\mathrm{\Delta }m^2`$ arises after the Higgs doublets $`H_1`$ and $`H_2`$ acquiring their vacuum expectation values $`v_1`$ and $`v_2`$, and
$`\mathrm{\Delta }m^2`$ $``$ $`v\left(A_\text{v}\mathrm{sin}\beta C_\text{v}\mathrm{cos}\beta \right),`$ (10)
with $`v=\sqrt{v_1^2+v_2^2}176\text{ GeV}`$ and the angle $`\beta =\mathrm{tan}^1(v_2/v_1)`$. As shown in Eq.(9), we have included the non-standard trilinear scalar coupling $`C_\text{v}`$ in the soft-supersymmetry-breaking scalar potential to extend the applicability of our effective model. As to be shown later, a special value of $`C_\text{v}`$ can lead to distinct collider signatures. The masses of the mass eigenstates $`\stackrel{~}{\nu }_{\tau _1}`$ and $`\stackrel{~}{\nu }_{\tau _2}`$ are given by $`m_{\stackrel{~}{\nu }_{\tau _1},\stackrel{~}{\nu }_{\tau _2}}^2=(m_{\stackrel{~}{\nu }_{\tau _L}}^2+m_{\stackrel{~}{\nu }_{\tau _R}}^2\sqrt{(m_{\stackrel{~}{\nu }_{\tau _L}}^2m_{\stackrel{~}{\nu }_{\tau _R}}^2)^2+4(\mathrm{\Delta }m^2)^2})/2`$, with $`m_{\stackrel{~}{\nu }_{\tau _1}}<m_{\stackrel{~}{\nu }_{\tau _2}}`$. In terms of the sneutrino masses, the sneutrino mixing factor $`\mathrm{sin}2\theta _{\stackrel{~}{\nu }}`$ can be written as:
$$\mathrm{sin}2\theta _{\stackrel{~}{\nu }}=\frac{2\mathrm{\Delta }m^2}{m_{\stackrel{~}{\nu }_{\tau _2}}^2m_{\stackrel{~}{\nu }_{\tau _1}}^2}.$$
(11)
Since we require a positive mass for the lightest tau sneutrino $`\stackrel{~}{\nu }_{\tau _1}`$, we must have the mass parameters satisfy the following constraint:
$$m_{\stackrel{~}{\nu }_{\tau _L}}^2m_{\stackrel{~}{\nu }_{\tau _R}}^2>(\mathrm{\Delta }m^2)^2.$$
(12)
Under the $`\stackrel{~}{\nu }_{\tau _1}`$-LSP scenario, the LSP can be generated either through the direct productions in high energy collision or through the decay of heavy sparticles, such as the next-to-lightest supersymmetric particle (NLSP) and the heavy tau sneutrinos $`\stackrel{~}{\nu }_{\tau _2}`$. It can also be pair-produced through the $`Z`$ boson decay if the $`\stackrel{~}{\nu }_{\tau _1}`$ mass is smaller than one half of the $`Z`$ mass. The partial decay width for the $`Z`$ boson decaying into a $`\stackrel{~}{\nu }_{\tau _1}\stackrel{~}{\nu }_{\tau _1}^{}`$ pair is given at the tree level as:
$$\mathrm{\Gamma }(Z\stackrel{~}{\nu }_{\tau _1}\stackrel{~}{\nu }_{\tau _1}^{})=\frac{\alpha _{em}(\mathrm{cot}\theta _W+\mathrm{tan}\theta _W)^2\mathrm{sin}^4\theta _{\stackrel{~}{\nu }}}{48}m_Z\left[1\frac{4m_{\stackrel{~}{\nu }_{\tau _1}}^2}{m_Z^2}\right]^{3/2},$$
(13)
where $`\theta _W`$ is the weak mixing angle and $`\alpha _{em}`$ is the fine structure constant evaluated at the $`Z`$-mass scale $`m_Z`$. Clearly, the partial decay width $`\mathrm{\Gamma }(Z\stackrel{~}{\nu }_{\tau _1}\stackrel{~}{\nu }_{\tau _1}^{})`$ depends on the $`\stackrel{~}{\nu }_{\tau _1}`$ mass and the mixing angle $`\mathrm{sin}\theta _{\stackrel{~}{\nu }}`$. From the experimental data at LEP and SLC, the invisible decay channel of $`Z`$ boson is bounded from above, and $`\mathrm{\Delta }\mathrm{\Gamma }(Z\text{invisible})<2\text{ MeV}`$ . Though this experimental constraint is automatically satisfied for the lightest tau sneutrino mass $`m_{\stackrel{~}{\nu }_{\tau _1}}`$ to be larger than one half of the $`Z`$-boson mass, it does not excludes the possibility that $`m_{\stackrel{~}{\nu }_{\tau _1}}`$ can be very small as compared to $`m_Z/2`$. Indeed, we find that $`m_{\stackrel{~}{\nu }_{\tau _1}}`$ can take any small value as long as the mixing angle $`\mathrm{sin}\theta _{\stackrel{~}{\nu }}`$ is smaller than about $`0.39`$. Because $`\stackrel{~}{\nu }_{\tau _L}`$ carries the electroweak quantum number and $`\stackrel{~}{\nu }_{\tau _R}`$ is a singlet field, a SUSY model usually predicts the soft mass parameter $`m_{\stackrel{~}{\nu }_{\tau _L}}`$ to be larger than $`m_{\stackrel{~}{\nu }_{\tau _R}}`$ at the weak scale, after including the running effect of the renormalization group equations. This case is assumed hereafter. From Eqs.(8) and (12), we conclude that the soft mass parameter $`\mathrm{\Delta }m`$ should be less than $`m_{\stackrel{~}{\nu }_{\tau _L}}\sqrt{\mathrm{tan}\theta _{\stackrel{~}{\nu }}}`$, and $`m_{\stackrel{~}{\nu }_{\tau _L}}<m_{\stackrel{~}{\nu }_{\tau _R}}\mathrm{cot}\theta _{\stackrel{~}{\nu }}`$. Consequently, as an example, for $`m_{\stackrel{~}{\nu }_{\tau _L}}`$ to be 200 GeV, the upper bound on $`\mathrm{\Delta }m`$ is about $`130\text{ GeV}`$, which imposes a constraint on the values of $`A_\text{v}`$ and $`C_\text{v}`$ for a fixed $`\mathrm{tan}\beta `$. (Following the method in Ref. , we have checked that in this model, there is no useful bound on the value of $`A_\text{v}`$, assuming $`C_\text{v}=0`$, from requiring the absence of dangerous charge and color breaking minima or unbounded from below directions.) For $`\mathrm{sin}\theta _{\stackrel{~}{\nu }}>0.39`$, a larger $`\mathrm{sin}\theta _{\stackrel{~}{\nu }}`$ requires a larger $`m_{\stackrel{~}{\nu }_{\tau _1}}`$, e.g., when $`\mathrm{sin}\theta _{\stackrel{~}{\nu }}`$ is 0.5, the minimal allowed value for $`m_{\stackrel{~}{\nu }_{\tau _1}}`$ is $`32\text{ GeV}`$. Our results are summarized in Fig.1, which shows the allowed parameter space on the $`(\mathrm{sin}\theta _{\stackrel{~}{\nu }},m_{\stackrel{~}{\nu }_{\tau _1}})`$ parameter plane. (Only the allowed range for $`m_{\stackrel{~}{\nu }_{\tau _1}}<m_Z/2`$ is shown. For $`m_{\stackrel{~}{\nu }_{\tau _1}}>m_Z/2`$, $`\mathrm{sin}\theta _{\stackrel{~}{\nu }}`$ can take any value within 1.)
## III New Production Mechanism of Higgs Bosons at Colliders
The trilinear scalar couplings for the neutral Higgs bosons and tau sneutrinos can be derived from Eqs.(7) and (9) as:
$`\mathrm{\Delta }`$ $`=`$ $`\{H^0(A_\text{v}\mathrm{sin}\alpha C_\text{v}\mathrm{cos}\alpha )+h^0(A_\text{v}\mathrm{cos}\alpha +C_\text{v}\mathrm{sin}\alpha )\}\{\mathrm{sin}2\theta _{\stackrel{~}{\nu }}(\stackrel{~}{\nu }_{\tau _2}\stackrel{~}{\nu }_{\tau _2}^{}\stackrel{~}{\nu }_{\tau _1}\stackrel{~}{\nu }_{\tau _1}^{})`$ (14)
$`+`$ $`\mathrm{cos}2\theta _{\stackrel{~}{\nu }}(\stackrel{~}{\nu }_{\tau _1}\stackrel{~}{\nu }_{\tau _2}^{}+\stackrel{~}{\nu }_{\tau _2}\stackrel{~}{\nu }_{\tau _1}^{})\}/\sqrt{2}+iA^0(A_\text{v}\mathrm{cos}\beta +C_\text{v}\mathrm{sin}\beta )(\stackrel{~}{\nu }_{\tau _2}\stackrel{~}{\nu }_{\tau _1}^{}\stackrel{~}{\nu }_{\tau _1}\stackrel{~}{\nu }_{\tau _2}^{})/\sqrt{2},`$ (15)
where $`h^0`$ is the lightest CP-even neutral Higgs boson, $`H^0`$ denotes the other CP-even neutral Higgs boson and $`A^0`$ is the CP-odd neutral Higgs particle. The phase angle $`\alpha `$ defines the mass eigenstates of $`h^0`$ and $`H^0`$.
When the mixing angle $`\theta _{\stackrel{~}{\nu }}`$ is small, the lightest sneutrino $`\stackrel{~}{\nu }_{\tau _1}`$ (which is almost $`\stackrel{~}{\nu }_{\tau _R}`$-like) predominantly interacts with $`\stackrel{~}{\nu }_{\tau _2}`$ (which is almost $`\stackrel{~}{\nu }_{\tau _L}`$-like) via the scalar interactions, cf. $`\mathrm{\Delta }`$. Hence, to study the production of the LSP $`\stackrel{~}{\nu }_{\tau _1}`$, it is desirable to first examine the decay and the production of $`\stackrel{~}{\nu }_{\tau _2}`$.
When $`\stackrel{~}{\nu }_{\tau _2}`$ is produced at colliders, it may decay into the tau neutrino $`\nu _\tau `$ and the neutralino $`\stackrel{~}{\chi }_1^0`$ (or $`\stackrel{~}{\chi }_2^0`$), or into the Higgs particle $`h^0`$ and the lightest sneutrino $`\stackrel{~}{\nu }_{\tau _1}`$, assuming the other modes either are forbidden by mass relation or have negligible partial decay widths.) The branching ratios (BR) of these two decay modes depend on the SUSY parameters. For illustration, we give their tree level partial decay widths as follows:
$`\mathrm{\Gamma }(\stackrel{~}{\nu }_{\tau _2}\stackrel{~}{\nu }_{\tau _1}h^0)`$ $`=`$ $`{\displaystyle \frac{\mathrm{cos}^22\theta _{\stackrel{~}{\nu }}|A_\text{v}\mathrm{cos}\alpha +C_\text{v}\mathrm{sin}\alpha |^2}{32\pi m_{\stackrel{~}{\nu }_{\tau _2}}}}`$ (16)
$`\times `$ $`\sqrt{\left[\left(1+{\displaystyle \frac{m_{h^0}}{m_{\stackrel{~}{\nu }_{\tau _2}}}}\right)^2{\displaystyle \frac{m_{\stackrel{~}{\nu }_{\tau _1}}^2}{m_{\stackrel{~}{\nu }_{\tau _2}}^2}}\right]\left[\left(1{\displaystyle \frac{m_{h^0}}{m_{\stackrel{~}{\nu }_{\tau _2}}}}\right)^2{\displaystyle \frac{m_{\stackrel{~}{\nu }_{\tau _1}}^2}{m_{\stackrel{~}{\nu }_{\tau _2}}^2}}\right]},`$ (17)
$`\mathrm{\Gamma }(\stackrel{~}{\nu }_{\tau _2}\nu _\tau \stackrel{~}{\chi }_j^0)`$ $`=`$ $`{\displaystyle \frac{\alpha _{em}m_{\stackrel{~}{\nu }_{\tau _2}}\mathrm{cos}^2\theta _{\stackrel{~}{\nu }}}{8}}\left|{\displaystyle \frac{V_{1j}}{\mathrm{cos}\theta _W}}{\displaystyle \frac{V_{2j}}{\mathrm{sin}\theta _W}}\right|^2\left(1{\displaystyle \frac{m_{\stackrel{~}{\chi }_j^0}^2}{m_{\stackrel{~}{\nu }_{\tau _2}}^2}}\right)^2,`$ (18)
where $`V_{1j}`$ and $`V_{2j}`$, for $`j=1,2,3,4`$, denote the matrix elements of the diagonalizing matrix for the neutralino mass matrix . Based on Eqs.(17) and (18), we plot in Fig.2(a) the branching ratio BR($`\stackrel{~}{\nu }_{\tau _2}\stackrel{~}{\nu }_{\tau _1}+h^0`$) for $`m_{\stackrel{~}{\nu }_{\tau _2}}=200`$ GeV, $`m_{\stackrel{~}{\nu }_{\tau _1}}=20`$ GeV and $`m_{h^0}=130`$ GeV, as a function of $`A_\text{v}`$, with $`C_\text{v}=0`$ and $`\mathrm{tan}\beta =2`$. (The constraint from the invisible decay width of $`Z`$ boson requires $`A_\text{v}<96\text{ GeV}`$, with $`\mathrm{sin}\theta _{\stackrel{~}{\nu }}<0.42`$, cf. Fig.1 and Eq.(11).) Four curves for different neutralino mixing scenarios are plotted. The first two curves (B1 and B2) are for $`\stackrel{~}{\chi }_1^0`$-NLSP to be Bino-like, in which we have assumed the common soft SUSY breaking masses $`m_1=\text{100 GeV}`$, $`m_2=\text{200 GeV}`$ and $`\mathrm{tan}\beta =2`$, but with different $`\mu `$ values: $`\mu =500,500`$ GeV. The third curve (M) is for the mixed-type $`\stackrel{~}{\chi }_1^0`$-NLSP scenario, in which $`\mu =100`$ GeV. The last one (H) is for the Higgsino-like $`\stackrel{~}{\chi }_1^0`$-NLSP scenario, in which $`m_1=200\text{ GeV}`$, $`m_2=400\text{ GeV}`$ and $`\mu =100\text{ GeV}`$. As shown in the figure, the branching ratio increases as $`A_\text{v}`$ increases. Furthermore, due to the small Yukawa coupling of the $`\stackrel{~}{\nu }_{\tau _2}`$-higgsino-$`\nu _\tau `$ interaction, the branching ratio BR($`\stackrel{~}{\nu }_{\tau _2}\stackrel{~}{\nu }_{\tau _1}+h^0`$) increases rapidly under the Higgsino-like $`\stackrel{~}{\chi }_1^0`$-NLSP scenario.
The amusing feature of the model is that the lightest neutral Higgs boson $`h^0`$ can be largely produced at either electron or hadron colliders from the decay of $`\stackrel{~}{\nu }_{\tau _2}`$, which is mainly $`\stackrel{~}{\nu }_{\tau _L}`$ when the left-right mixing angle is small. Since a left sneutrino $`\stackrel{~}{\nu }_{\tau _L}`$ carries electroweak quantum number, it can be produced directly in collisions, or indirectly from the decay of charginos, neutralinos, sleptons, and the cascade decay of squarks and gluinos, etc. There are plenty of studies in the literature to show that the production rates of the above mentioned sparticles at the current and future colliders can be very large, depending on the SUSY parameters . In that case, our model would predict a large production rate of events including either single $`h^0`$ or multiple $`h^0`$’s, provided that BR($`\stackrel{~}{\nu }_{\tau _2}\stackrel{~}{\nu }_{\tau _1}+h^0`$) is large enough. Many of the events including $`h^0`$’s can also contain single or multiple isolated leptons and/or photons with large transverse momentum, so it will not be difficult to trigger on such events experimentally. Due to the limited space in this short Letter, we cannot explore all the interesting possibilities in details for various colliders. Instead, we shall illustrate the above observation for the future Linear Collider (LC) in the next section.
Before closing this section, we remark that, in contrast to the $`\stackrel{~}{\nu }_{\tau _L}`$-LSP scenario, a $`\stackrel{~}{\nu }_{\tau _R}`$-like $`\stackrel{~}{\nu }_{\tau _1}`$ can be a good candidate for the cold dark matter (CDM). The left sneutrino had been suggested in the literature to be the LSP in the MSSM . They can annihilate rapidly in the early universe via s-channel $`Z`$-boson and t-channel neutralino and chargino exchanges. To reduce the LSP annihilation and obtain an acceptable relic abundance, it was proposed that the left sneutrinos should be either as light as $`m_{\stackrel{~}{\nu }}2\text{ GeV}`$ or as heavy as $`550\text{ GeV}<m_{\stackrel{~}{\nu }}<2300\text{ GeV}`$ . However, both of these proposals have been excluded by experiments. The light $`\stackrel{~}{\nu }_L`$ scenario was excluded by the measurements of $`Z`$ decay width, and the heavy $`\stackrel{~}{\nu }_L`$ scenario was excluded by the Heidelberg-Moscow direct detection experiment . Since a $`\stackrel{~}{\nu }_{\tau _R}`$-like $`\stackrel{~}{\nu }_{\tau _1}`$ interacts with other particles mainly through the left-right sneutrino mixing or the trilinear scalar coupling $`\stackrel{~}{\nu }_{\tau _R}`$-$`\stackrel{~}{\nu }_{\tau _L}`$-$`h^0`$, the LSP annihilation cross sections are generically small due to the presence of the small $`\mathrm{sin}^4\theta _{\stackrel{~}{\nu }}`$ or $`\mathrm{sin}^2\theta _{\stackrel{~}{\nu }}`$ factors coming from the mixing effect or the couplings. Comparing to the $`\stackrel{~}{\nu }_L`$-LSP annihilation in the ordinary MSSM , the $`\stackrel{~}{\nu }_{\tau _1}`$-LSP annihilation rate via the exchange of $`Z`$-boson, neutralinos, or charginos, is suppressed by a factor of $`\mathrm{sin}^4\theta _{\stackrel{~}{\nu }}`$ because the mixing of $`\stackrel{~}{\nu }_{\tau _R}`$ and $`\stackrel{~}{\nu }_{\tau _L}`$ yields a factor of $`\mathrm{sin}\theta _{\stackrel{~}{\nu }}`$ in the scattering amplitude. In addition to the usual MSSM processes, $`\stackrel{~}{\nu }_{\tau _1}`$ can also annihilate via an s-channel Higgs boson to produce light fermion pairs, whose scattering amplitude is suppressed by a factor of $`\mathrm{sin}\theta _{\stackrel{~}{\nu }}`$. Notice that the above rate can strongly depend on $`\mathrm{tan}\beta `$ because of the coupling of Higgs boson and fermions (such as bottom quarks). In lack of a complete SUSY model which gives the mass spectrum of the sparticles, and the $`\stackrel{~}{\nu }_{\tau _1}`$ annihilation rate depends on the details of the MSSM parameters, we only remark that with the additional suppression factor discussed above, a $`\stackrel{~}{\nu }_{\tau _R}`$-like LSP $`\stackrel{~}{\nu }_{\tau _1}`$ can be a good candidate for CDM.
## IV $`\stackrel{~}{\nu }_{\tau _1}`$-LSP phenomenology at the LC
In this section, we consider a simple example to illustrate the interesting phenomenology of our model expected at high energy colliders.
The tree-level cross section for the production of $`\stackrel{~}{\nu }_{\tau _2}`$ pair in $`e^+e^{}`$ collision is
$`\sigma (e^+e^{}\stackrel{~}{\nu }_{\tau _2}\stackrel{~}{\nu }_{\tau _2}^{})={\displaystyle \frac{\pi \alpha _{em}^2\mathrm{cos}^2\theta _{\stackrel{~}{\nu }}}{6S}}{\displaystyle \frac{14\mathrm{sin}^2\theta _W+8\mathrm{sin}^4\theta _W}{\mathrm{sin}^42\theta _W}}\left[1{\displaystyle \frac{4m_{\stackrel{~}{\nu }_{\tau _2}}^2}{S}}\right]^{3/2}\left[1{\displaystyle \frac{m_Z^2}{S}}\right]^2,`$ (19)
through the s-channel $`Z`$-exchange diagram. ($`\sqrt{S}`$ is the center-of-mass energy of the collider.) As a comparison, the tree-level cross sections for the direct productions of $`\stackrel{~}{\nu }_{\tau _1}`$ in $`e^+e^{}`$ collision associated with $`\stackrel{~}{\nu }_{\tau _2}^{}`$ or $`\stackrel{~}{\nu }_{\tau _1}^{}`$ are
$`\sigma (e^+e^{}\stackrel{~}{\nu }_{\tau _1}\stackrel{~}{\nu }_{\tau _2}^{})`$ $`=`$ $`\sigma (e^+e^{}\stackrel{~}{\nu }_{\tau _1}^{}\stackrel{~}{\nu }_{\tau _2})`$ (20)
$`=`$ $`{\displaystyle \frac{\pi \alpha _{em}^2\mathrm{sin}^22\theta _{\stackrel{~}{\nu }}}{24S}}{\displaystyle \frac{14\mathrm{sin}^2\theta _W+8\mathrm{sin}^4\theta _W}{\mathrm{sin}^42\theta _W}}`$ (21)
$`\times `$ $`\left[\left(1{\displaystyle \frac{m_{\stackrel{~}{\nu }_{\tau _2}}^2m_{\stackrel{~}{\nu }_{\tau _1}}^2}{S}}\right)^2{\displaystyle \frac{4m_{\stackrel{~}{\nu }_{\tau _1}}^2}{S}}\right]^{3/2}\left[1{\displaystyle \frac{m_Z^2}{S}}\right]^2,`$ (22)
$`\sigma (e^+e^{}\stackrel{~}{\nu }_{\tau _1}\stackrel{~}{\nu }_{\tau _1}^{})`$ $`=`$ $`{\displaystyle \frac{\pi \alpha _{em}^2\mathrm{sin}^4\theta _{\stackrel{~}{\nu }}}{6S}}{\displaystyle \frac{14\mathrm{sin}^2\theta _W+8\mathrm{sin}^4\theta _W}{\mathrm{sin}^42\theta _W}}`$ (23)
$`\times `$ $`\left[1{\displaystyle \frac{4m_{\stackrel{~}{\nu }_{\tau _1}}^2}{S}}\right]^{3/2}\left[1{\displaystyle \frac{m_Z^2}{S}}\right]^2.`$ (24)
Hence, the direct productions of $`\stackrel{~}{\nu }_{\tau _1}`$ could be highly suppressed when the factor $`\mathrm{sin}^4\theta _{\stackrel{~}{\nu }}`$ is much smaller than 1. Apart from a different phase space factor, the direct production of $`\stackrel{~}{\nu }_{\tau _1}\stackrel{~}{\nu }_{\tau _1}^{}`$ is smaller than the production of $`\stackrel{~}{\nu }_{\tau _2}\stackrel{~}{\nu }_{\tau _2}^{}`$ by a factor of $`\mathrm{sin}^4\theta _{\stackrel{~}{\nu }}/\mathrm{cos}^2\theta _{\stackrel{~}{\nu }}`$.
As mentioned in the previous section, $`\stackrel{~}{\nu }_{\tau _2}`$ could decay into $`\stackrel{~}{\nu }_{\tau _1}h^0`$ or $`\nu _\tau \stackrel{~}{\chi }_j^0`$. With a large BR($`\stackrel{~}{\nu }_{\tau _2}\stackrel{~}{\nu }_{\tau _1}h^0`$), it may be possible to observe the $`2h^0+E_T/`$ signal originated from the production of the $`\stackrel{~}{\nu }_{\tau _2}`$ pair. To test this scenario, we need to calculate the SM rate for the process $`e^+e^{}h^0h^0\nu _i\overline{\nu }_i`$, where $`\nu _i`$ ($`\overline{\nu }_i`$) is the left-handed neutrino (anti-neutrino) for the $`i`$th family. Since the final state neutrinos, like the LSPs, carry away energy, the above SM process is the intrinsic background to the detection of the signal event $`e^+e^{}\stackrel{~}{\nu }_{\tau _2}\stackrel{~}{\nu }_{\tau _2}^{}h^0h^0\stackrel{~}{\nu }_{\tau _1}\stackrel{~}{\nu }_{\tau _1}^{}`$, because both processes produce the event signature of $`e^+e^{}2h^0+E_T/`$. Due to the large suppression factor from the 4-body phase space, the cross section for the SM process is typically small. For example, when the Higgs mass is 130 GeV, the SM rate is about $`0.03\text{fb}`$, for $`\sqrt{S}=500\text{GeV}`$. On the other hand, for a 200 GeV $`\stackrel{~}{\nu }_{\tau _2}`$, with small left-right tau sneutrino mixing effect (i.e. $`\mathrm{sin}\theta _{\stackrel{~}{\nu }}0`$), the tree-level cross section for the $`\stackrel{~}{\nu }_{\tau _2}`$-pair production is about 12 fb. This relatively large production rate can lead to an enhancement in the Higgs boson pair signal, provided BR($`\stackrel{~}{\nu }_{\tau _2}\stackrel{~}{\nu }_{\tau _1}h^0`$) is large enough, cf. Fig.2.
It is also interesting to consider a special case with $`C_\text{v}=A_\text{v}\mathrm{tan}\beta `$, so that $`\mathrm{\Delta }m^2`$ vanishes and $`\stackrel{~}{\nu }_{\tau _L}`$ does not mix with $`\stackrel{~}{\nu }_{\tau _R}`$. Hence, the direct production of $`\stackrel{~}{\nu }_{\tau _1}`$ (now a pure $`\stackrel{~}{\nu }_{\tau _R}`$) vanishes and $`\stackrel{~}{\nu }_{\tau _1}`$ is predominantly produced through the decay of $`\stackrel{~}{\nu }_{\tau _2}\stackrel{~}{\nu }_{\tau _1}h^0`$. As shown in Fig.2(b), the branching ratio BR($`\stackrel{~}{\nu }_{\tau _2}\stackrel{~}{\nu }_{\tau _1}+h^0`$) is generally higher than that in Fig.2(a) because of the inclusion of the non-standard soft SUSY breaking parameter $`C_\text{v}`$. For instance, assuming the same mass parameters as the Higgsino-like $`\stackrel{~}{\chi }_1^0`$-NLSP scenario given above, the BR($`\stackrel{~}{\nu }_{\tau _2}\stackrel{~}{\nu }_{\tau _1}+h^0`$) is about 0.1 for $`A_\text{v}=10\text{ GeV}`$, and 0.9 for $`A_\text{v}=90\text{ GeV}`$ with $`\mathrm{tan}\beta =2`$, which results in the production cross section of $`e^+e^{}\stackrel{~}{\nu }_{\tau _2}\stackrel{~}{\nu }_{\tau _2}^{}h^0h^0\stackrel{~}{\nu }_{\tau _1}\stackrel{~}{\nu }_{\tau _1}^{}`$ to be about 0.12 fb and 11 fb, respectively. When comparing to the SM rate, the signal rate of our model can be larger by about a factor of 300. With a larger $`\mathrm{tan}\beta `$, the branching ratio of $`\stackrel{~}{\nu }_{\tau _2}\stackrel{~}{\nu }_{\tau _1}+h^0`$ increases and more $`2h^0+E_T/`$ signal events are expected.
Thus far, we have only shown that the $`\stackrel{~}{\nu }_{\tau _1}`$-LSP signal rate of $`2h^0+E_T/`$ production can be much larger than the SM rate. However, to distinguish our $`\stackrel{~}{\nu }_{\tau _1}`$-LSP scenario from the ordinary MSSM scenario with $`\stackrel{~}{\chi }^0`$-LSP, we also need to compare our signal rate with that predicted by the ordinary MSSM. When the decay mode $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_1^0+h^0`$ is available, the MSSM $`2h^0+E_T/`$ signature mainly comes from the pair production of $`\chi _2^0`$ in $`e^+e^{}`$ collisions. The event rate of $`e^+e^{}\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_2^0h^0h^0\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0`$ depends on the detail of the MSSM parameters. For example, assuming the usual scenario that $`\stackrel{~}{\chi }_1^0`$ is almost Bino-like, the partial decay width of $`\stackrel{~}{\chi }_2^0h^0\stackrel{~}{\chi }_1^0`$ will come from one-loop corrections. Since the difference in the masses of $`\chi _2^0`$ and $`\chi _1^0`$ is large enough (greater than $`m_{h^0}`$) for the other tree-level decay modes of $`\chi _2^0`$ to be opened, the branching ratio BR($`\stackrel{~}{\chi }_2^0h^0\stackrel{~}{\chi }_1^0`$) is generally small. Furthermore, if $`\stackrel{~}{\chi }_2^0`$ is Higgsino-like, the cross section of $`e^+e^{}\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_2^0`$ at a 500 GeV collider is expected to be small, at the order of fb or smaller, for MSSM mass parameters to be at the order of a few hundreds of GeV. For a gaugino-like $`\stackrel{~}{\chi }_2^0`$, the above cross section can increase by a factor of 10. A detailed comparison between our model and the ordinary MSSM for the event rate of $`e^+e^{}2h^0E_T/`$ is beyond the scope of this short Letter.
## V conclusion
Motivated by the neutrino oscillation data, we study a low energy effective theory in which the interactions of right sneutrinos with left sneutrinos at the weak scale are described by Eqs.(7), (8) and (9) with R-parity conservation in the framework of MSSM. For simplicity, we assume the mixings among the three generation sneutrinos are small enough that the dominant effect to collider phenomenology comes from the interaction of the left and right tau sneutrinos, which mix via the soft SUSY breaking effect arising from the trilinear scalar couplings of scalar Higgs doublets and sleptons. We find that the mass of the lightest tau sneutrino $`\stackrel{~}{\nu }_{\tau _1}`$ can take any value (even smaller than $`m_Z/2`$) to agree with the $`Z`$ decay width measurement, provided that the sneutrino mixing parameter $`\mathrm{sin}\theta _{\stackrel{~}{\nu }}`$ is smaller than $`0.39`$. In that case, $`\stackrel{~}{\nu }_{\tau _1}`$ is almost the right tau sneutrino $`\stackrel{~}{\nu }_{\tau _R}`$, and becomes a good candidate to be the LSP of the model as well as the cold dark matter. Because $`\stackrel{~}{\nu }_{\tau _R}`$ mainly interacts with $`\stackrel{~}{\nu }_{\tau _L}`$ through Higgs boson, the branching ratio BR($`\stackrel{~}{\nu }_{\tau _2}\stackrel{~}{\nu }_{\tau _1}+h^0`$) can be large. It results in a large production rate of single $`h^0`$ or multiple $`h^0`$’s in electron or hadron collisions via the decay of $`\stackrel{~}{\nu }_{\tau _2}`$, which is $`\stackrel{~}{\nu }_{\tau _L}`$-like and can be copiously produced from the decay of charginos, neutralinos, sleptons, and the cascade decay of squarks and gluinos, etc. Hence, the events including Higgs boson(s) can also contain single or multiple isolated leptons and/or photons with large transverse momentum, which can make it easy to trigger on such events experimentally. Given the possible large production rate of the above mentioned sparticles, the production rate of the lightest neutral Higgs boson is expected to be largely enhanced from its SM rate. For example, with the trilinear couplings $`C_\text{v}=A_\text{v}\mathrm{tan}\beta `$, the left and right sneutrino do not mix, and the BR($`\stackrel{~}{\nu }_{\tau _2}\stackrel{~}{\nu }_{\tau _1}+h^0`$) is approaching to 1 as $`A_\text{v}`$ increases. In that case, the signal rate ($`12`$ fb) of $`e^+e^{}2h^0+E_T/`$ at a 500 GeV LC is enhanced by a factor of 400, as compared to its SM rate.
## Acknowledgments
CPY thanks L. Diaz-Cruz and Y. Okada for discussions, and G.L. Kane for useful suggestions and a critical reading of the manuscript. We are also grateful to the warm hospitality of the Center for Theoretical Science in Taiwan where part of this work was completed. This work is in part supported by the National Science Council in Taiwan and the National Science Foundation in the USA under the grant PHY-9802564.
Note Added:
After posting this manuscript to the xxx-archives, stamped as hep-ph/0006313, we noticed that in a new paper, hep-ph/0006312, several supersymmetry breaking mechanisms were proposed for generating light sneutrinos. CPY thanks N. Weiner for explaining the SUSY models proposed in Ref. , and pointing out an error in the previous version of the manuscript.
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# USACH-FM-00/07 Deformed Heisenberg Algebra with Reflection, Anyons and Supersymmetry of Parabosons1footnote 11footnote 1Talk given at the International Conference on “Spin-Statistics Connection and Commutation Relations”, Centro Internazionale per la Cultura Scientifica dell’Università di Napoli, Anacapri, Capri Island, Italy – May 31-June 3, 2000
## Introduction
Generalized statistics was introduced in physics in the form of parastatistics as an exotic possibility extending the Bose and Fermi statistics . It was closely related with the discovery of color in the context of the theory of strong interactions. Nowadays generalized statistics finds applications in the physics of the quantum Hall effect and (probably) it is relevant to high temperature superconductivity . Supersymmetry, instead, unifies Bose and Fermi statistics and its development lead to the construction of field and string theories with exceptional properties , that transformed the same idea of supersymmetry in one of the cornerstones of modern theoretical physics. Supersymmetry was observed in nuclear physics in the form of dynamical symmetry , whereas its manifestation as a fundamental symmetry of elementary particle physics still waits for experimental confirmation.
Though supersymmetry and generalized statistics may be unified in the form of parasupersymmetry , nevertheless, by the construction, the two concepts seem to be independent. Recently, the existence of intimate relation between generalized statistics and supersymmetry was established by observation of hidden supersymmetric structure in purely parabosonic and purely pafermionic systems. The key tool with which the observation of close relationship between generalized statistics and supersymmetry was realized is the so called deformed Heisenberg algebra with reflection, or the R-deformed Heisenberg algebra (RDHA) . The algebraic construction of RDHA appeared in the work of Wigner , where he investigated the problem of correlation of equations of motion with quantum mechanical commutation relations and proposed the generalized quantization schemes which subsequently lead to the theoretical discovery of parastatistics (in this context, see also refs. ). RDHA represents, probably, one of the first examples of deformation of bosonic harmonic oscillator which, as it was shown recently, possesses some universality being also related to parafermions , to (2+1)-dimensional anyons , and to the bosonized form of supersymmetric quantum mechanics . Besides, the RDHA structure underlies the construction of fractional supersymmetry . In this talk we review the application of RDHA for the universal description of ordinary spin-$`j`$ and anyon fields in 2+1 dimensions by means of first order linear differential equations , and discuss the exotic supersymmetry of purely parabosonic systems .
## Supersymmetry of parabosons
The deformed Heisenberg algebra with reflection is generated by the creation-annihilation operators $`a^+`$, $`a^{}`$ and by the reflection operator $`R`$ satisfying the relations
$$[a^{},a^+]=1+\nu R,\{R,a^\pm \}=0,R^2=1,$$
(1)
where $`\nu `$ is a real deformation parameter. Due to these basic relations, the creation-annihilation operators satisfy the trilinear parabosonic commutation relations $`[\{a^{},a^+\},a^\pm ]=\pm 2a^\pm `$. Introducing the Fock vacuum state, $`a^{}|0=0`$, and fixing the action of operator $`R`$ on it as $`R|0=|0`$, we arrive at the relation $`a^{}a^+|0=(1+\nu )|0`$, which together with trilinear commutation relation means that at $`\nu =p1`$, $`p=1,2,\mathrm{}`$, the operators $`a^\pm `$ have the sense of single-mode creation-annihilation operators of paraboson of order $`p`$. Vice versa, one can show that parabosonic trilinear commutation relations themselves give rise to RDHA . In general case, the number operator is realized in the form $`N=\frac{1}{2}\{a^{},a^+\}\frac{1}{2}(1+\nu )`$, and the reflection operator can be represented as $`R=(1)^N=\mathrm{cos}\pi N`$.
The reflection operator introduces a natural $`Z_2`$ grading structure in the Fock space and its presence in the definition of RDHA can be considered as an indication on possible relationship between parabosons and supersymmetry. To reveal a supersymmetry of parabosons, one notes that RDHA can also be given by the relations
$$a^+a^{}=F(N),a^{}a^+=F(N+1),[N,a^\pm ,]=\pm a^\pm ,$$
(2)
where $`F(N)=N+\nu \mathrm{sin}^2\frac{\pi N}{2}`$ is the characteristic function satisfying for $`\nu >1`$ the relations $`F(0)=0`$, $`F(n)>0`$, $`n=1,2,\mathrm{}`$. These relations mean, in particular, that for $`\nu >1`$ the corresponding representations of RDHA are unitary and infinite-dimensional. On the other hand, in the case $`\nu =(2p+1)`$, $`p=1,2,\mathrm{}`$, the characteristic function possesses the property $`F(2p+1)=0`$ underlying the existence of $`(2p+1)`$-dimensional (non-unitary) irreducible representations of RDHA, which are associated with the deformed parafermions of order $`2p=2,4,\mathrm{}`$ . Let us restrict ourselves here by the case of unitary infinite-dimensional representations ($`\nu >1`$). When $`\nu =2k+1`$, $`k=0,1,\mathrm{}`$, the structure function satisfies the relation $`F(2n+1)=F(2n+\nu +1)`$, $`n=0,1,\mathrm{}`$. Therefore, in the case of parabosonic systems of even order $`p=2(k+1)`$, the spectrum of the quadratic Hamiltonian $`H=a^+a^{}`$ (or of $`H=a^{}a^+`$), reveals doubling of all higher-lying levels. This indicates on existence of supersymmetry in such purely parabosonic systems. As it follows from the explicit form of the characteristic function, in the case of paraboson of order $`p=2(k+1)`$, $`k=0,1,2\mathrm{}`$, and Hamiltonian $`H=a^+a^{}`$, the supersymmetry is characterized by the presence of $`k+1`$ singlet states with energies $`E=0,2,\mathrm{},2k`$. Therefore, only in the case $`k=0`$, the spectrum has one singlet state of zero energy, whereas all other cases are characterized by the presence of $`k`$ higher-lying singlet states of nonzero energy in addition to the zero energy ground state. The corresponding supercharges are the infinite series operators in the corresponding paraboson operators:
$$Q_+=(a^+)^{2k+1}\mathrm{sin}^2\pi J_0,Q_{}=(a^{})^{2k+1}\mathrm{cos}^2\pi J_0,J_0=\frac{1}{4}\{a^+,a^{}\}.$$
(3)
They together with the Hamiltonian satisfy the polynomial superalgebra:
$`Q_\pm ^2=0,[H,Q_\pm ]=0,`$
$`\{Q_+,Q_{}\}=(H2k)(H2k+2)\mathrm{}(H+2k2)(H+2k),`$ (4)
which in the case $`k=0`$ is reduced to the conventional $`N=1`$ linear superalgebra. The role of the grading operator in such purely parabosonic supersymmetric systems belongs to the reflection operator $`R=\mathrm{cos}\pi N`$.
The case of paraboson system of order $`p=2`$ ($`\nu =1`$) given by the Hamiltonian $`H=a^{}a^+`$ is characterized by the $`N=1`$ spontaneously broken linear supersymmetry: all the states are paired in supersymmetric doublets with the lowest energy level $`E=2`$. The corresponding supercharges have the form (3) with $`k=0`$ and with operators $`a^+`$ and $`a^{}`$ changed in their places . The systems of parabosons of order $`p=4,6,\mathrm{}`$, given by the Hamiltonian $`H=a^{}a^+2`$ possess nonlinear (polynomial) supersymmetry of the form similar to (4) .
It was shown in ref. that the supersymmetry of purely parabosonic systems can be understood as the supersymmetry of Calogero-like systems with exchange interaction and that in principle it can be realized in one-dimensional systems of identical fermions. Besides, it was demonstrated that nonlinear parabosonic supersymmetry can be obtained via appropriate modification of the classical analog of usual supersymmetric quantum mechanics.
## RDHA and anyons
The parabosonic supersymmetry structures corresponding to $`H=a^+a^{}`$ and $`H=a^{}a^+`$ can be unified and extended to the $`osp(2|2)`$ superalgebraic structure . The operators $`T_3=\frac{1}{2}\{a^+,a^{}\}`$, $`T_\pm =\frac{1}{2}(a^\pm )^2`$ and $`I=\frac{1}{2}(\nu +R)`$ have a sense of even generators of $`osp(2|2)`$ forming $`sl(2)\times u(1)`$ subalgebra, whereas the operators $`Q^\pm =a^\pm \mathrm{\Pi }_\pm `$ and $`S^\pm =a^\pm \mathrm{\Pi }_{}`$ are its odd generators, where $`\mathrm{\Pi }_\pm =\frac{1}{2}(1\pm R)`$ are the projectors on even and odd subspaces of the Fock space. On the other hand, the operators $`J_\mu `$, $`\mu =0,1,2`$, $`J_0=\frac{1}{2}T_3`$, $`J_1\pm iJ_2=T_\pm `$, and $`_\alpha `$, $`\alpha =1,2`$, $`_1=(a^++a^{})/\sqrt{2}`$, $`_2=i(a^+a^{})/\sqrt{2}`$, can be considered as even and odd generators of $`osp(1|2)`$ superalgebra: $`[J_\mu ,J_\nu ]=iϵ_{\mu \nu \lambda }J^\lambda `$, $`[J_\mu ,_\alpha ]=\frac{1}{2}(\gamma _\mu )_\alpha {}_{}{}^{\beta }_{\beta }^{}`$, $`\{_\alpha ,_\beta \}=4i(J\gamma )_{\alpha \beta }`$; here (2+1)-dimensional $`\gamma `$-matrices appear in the Majorana representation, see ref. . Since the Casimir operator $`𝒞=J_\mu J^\mu \frac{i}{8}^\alpha _\alpha `$ takes the fixed value $`𝒞=\frac{1}{16}(1\nu ^2)`$, this means that any irreducible representation of RDHA carries the corresponding irreducible representation of $`osp(1|2)`$, which, in turn, is a direct sum of two irreducible representations of $`so(2,1)`$ with the Casimir operator $`C=J^\mu J_\mu `$ taking the value $`C=\alpha _+(\alpha _+1)`$, $`\alpha _+=\frac{1}{4}(1+\nu )`$, on even ($`R=1`$) subspace of the Fock space, and $`C=\alpha _{}(\alpha _{}1)`$, $`\alpha _{}=\alpha _++\frac{1}{2}`$, on odd ($`R=1`$) subspace. The $`osp(1|2)`$ structure associated with RDHA can be exploited to describe anyons by means of covariant linear differential equations. For the purpose, let us consider the field $`\mathrm{\Psi }^n(x)`$ depending on space-time point $`x_\mu `$ in 2+1 dimensions and carrying infinite- ($`\nu >1`$, $`n=0,1,\mathrm{}`$) or finite- dimensional ($`\nu =(2p+1)`$, $`p=1,2,\mathrm{}`$, $`n=0,\mathrm{},2p`$) representation of RDHA. Then the following spinor set of linear differential equations describes universally ordinary spin-$`j`$ fields and anyons :
$$D_\alpha \mathrm{\Psi }(x)=0,D_\alpha =R𝒫_\alpha +m_\alpha ,$$
(5)
where $`R`$ is the reflection operator of RDHA, $`m`$ is a mass parameter, and $`𝒫_\alpha =(i\gamma _\mu ^\mu )_\alpha {}_{}{}^{\beta }_{\beta }^{}`$. The condition of integrability of two equations (5) is equivalent to the equations $`(^2+m^2)\mathrm{\Psi }_+(x)=0`$, $`(i^\mu J_\mu sm)\mathrm{\Psi }_+(x)=0`$, $`s=\frac{1}{4}(1+\nu )`$, for the even part (in the RDHA sense) of the field, $`R\mathrm{\Psi }_+(x)=\mathrm{\Psi }_+(x)`$, whereas the solution to equations (5) in odd subspace, $`R\mathrm{\Psi }_{}(x)=\mathrm{\Psi }_{}(x)`$, is trivial, $`\mathrm{\Psi }_{}(x)=0`$. The parameter $`\nu `$ fixes the value of spin $`s`$, and one concludes that in the case of finite-dimensional representations ($`\nu =(2p+1)`$), the corresponding field $`\mathrm{\Psi }_+(x)`$ carries integer or half-integer spin $`s=j`$, whereas the case of infinite-dimensional representations of RDHA ($`\nu >1`$) corresponds to the field of arbitrary spin $`s>0`$ (anyon). The case of anyon with $`s<0`$ can be obtained by a simple change $`mm`$ in (5). In the case of infinite-dimensional unitary representations of RDHA, the linear differential equation $`(i^\mu J_\mu sm)\mathrm{\Psi }_+(x)=0`$ is the (2+1)-dimensional analog of the $`(3+1)D`$ infinite-component Majorana equation, whose fundamental role for the description of anyons was established in ref. under investigation of the $`(2+1)D`$ model of relativistic particle with torsion (see also refs. ).
Varying the deformation parameter in the region $`\nu >1`$, one can obtain the fields of integer spin ($`\nu =4n1`$, $`s=n`$, $`n=1,2,\mathrm{}`$) as well as of half-integer spin ($`\nu =4n+1`$, $`s=n+\frac{1}{2}`$, $`n=0,1,\mathrm{}`$). However, such fields of integer and half-integer spin have a nature to be essentially different from the nature of usual spin-$`j`$ fields appearing in the case of finite-dimensional representations of RDHA ($`\nu =(2p+1)`$) since they have hidden nonlocality. In the rest frame, the solution to the Klein-Gordon and Majorana equations has only one nontrivial component in correspondence with the pseudoscalar (helicity) nature of spin in 2+1 dimensions . But the Lorentz boost enlivens all the infinite number of components of the field $`\mathrm{\Psi }_+(x)`$ in the case $`\nu >1`$ . There is the analog of coordinate representation for RDHA, in which
$$a^\pm =\frac{1}{\sqrt{2}}(qi𝒟_\nu ),𝒟_\nu =i\left(\frac{d}{dq}\frac{\nu }{2q}R\right),$$
(6)
and $`R\psi (q)=\psi (q)`$. In such representation the fields $`\mathrm{\Psi }_+(x)`$ have a structure of the functions to be even in continuous variable $`q`$: $`\mathrm{\Psi }_+(x,q)=\mathrm{\Psi }_+(x,q)`$. This means that the corresponding solutions to the equations (5) in the case $`\nu >1`$ have the hidden half-infinite nonlocality ($`q0`$) which is analogous to the string-like nonlocality of anyon fields in other approaches .
To conclude, RDHA finds various theoretical applications including the described two. It would be interesting to look for the experimental manifestation of the exotic supersymmetry of parabosons. In general, the existence of the polynomial supersymmetry is characterized by the presence of several singlet states which could be considered as an indication on parabosonic-like excitations (quasiparticles) in the system. The search for possible experimental manifestation of the nonlocal (in internal variable $`q`$) fields of integer and half-integer spin associated with parabosons seems to be another interesting problem.
### Acknowledgments.
I am grateful to I. Bandos, M. Rausch de Traubenberg and D. Sorokin for useful discussions, and to G. Marmo and E. C. G. Sudarshan for bringing refs. to my attention. The work was supported in part by FONDECYT (Chile) under grant 1980619 and by DICYT (USACH).
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# Radiating dipoles in photonic crystals
## I Introduction
Photonic crystals (PCs) have been the subject of intensive research over the past decade. These are fabricated periodic dielectric arrays that employ a combination of (i) Mie scattering from individual elements of the array, and (ii) Bragg scattering from the periodic lattice, to induce a band structure for photon propagation. This band structure is, in many ways, analogous to electronic band structure in a semiconductor. Through a judicious selection of materials and of the periodicity of the lattice, the photonic dispersion relation and the associated electromagnetic (EM) mode structure of a PC can be adapted to a variety of device applications. The most dramatic modification of the photon dispersion occurs when the linear propagation of a photon in a PC is prohibited in all directions for a range of frequencies, giving rise to a complete photonic band gap (PBG).
The radiative dynamics of an optically active material placed within or near a PC can be dramatically modified from free space. This is a result of the “colored” electromagnetic reservoir provided by the solutions to the electromagnetic field equations within a PC. In the optical domain, theoretical studies of atomic transitions coupled to the EM modes of a PC with an optical PBG predict a number of novel quantum optical phenomena. These phenomena include: the suppression or enhancement of spontaneous emission and the associated fractional localization of light near radiating atoms ; rapid all–optical switching ; and anomalous superradiant emission, as well as low–threshold lasing near the edge of a PBG. Microfabrication of PCs with complete PBGs at optical wavelengths has proven to be a difficult task, both because the lattice periodicity should be comparable to the wavelength of the light under consideration, and because a high dielectric contrast between the elements of the lattice is required. Fortunately, recent advances in microlithography and in semiconductor infiltration in colloidal crystals have produced materials with significant pseudo–gaps in their photonic band structure . The development of materials with complete PBGs in the optical regime appears imminent.
High–quality PBG materials at microwave frequencies have been available for some time . Sizeable band–gaps with center frequencies ranging from a few GHz up to 2 THz have been reported; these crystals have thus proven the soundness of the concept of the PBG. Microwave PBG materials may be relatively easily manufactured using micro-machining techniques, and are currently of interest for applications such as the shielding of human tissue from microwave radiation, and for improving the radiation characteristics of microwave antennæ. Although PBG materials at microwave frequencies have been extensively studied, the behavior of radiating dipolar antennæ embedded in microwave PCs has not received the same degree of attention. This is despite the fact that such antennæ would share many properties in common with atomic emission in a PC. In the microwave domain, a dipole antenna could take the form of an electrically excited metallic pin with a high $`Q`$ (quality) factor.
The radiative dynamics of the above system can be described by a charged, one–dimensional simple harmonic oscillator (SHO). Such an electric dipole oscillator can also provide an excellent description of the radiation of single or multiple two–level atoms in the optical domain. This description is valid provided that the total excitation energy of the atoms is well below an energy where saturation (nonlinear) effects become important. Moreover, the radiation reservoir can itself be modeled as a bath of many independent SHOs: Radiative damping arises from a linear coupling between the system SHO and the large number of reservoir oscillator modes. The similarities between the microwave and optical systems, coupled with the mature state of microwave technology, suggests that many of the predicted effects for atomic dipoles in the optical domain could be realized and studied first in the microwave domain.
Analytical techniques exist for treating certain forms of coupling between the dipole and reservoir for certain modal distributions of the reservoir. However, PCs present coupling distributions and spectral properties which defy analytical methods. This is due to the presence of a restricted and rapidly–varying reservoir mode distribution, which renders invalid the usual Born–Markov type approximation schemes for the system–reservoir interaction. To obtain accurate results, we solve the system numerically for a large, but finite, number of oscillators in the reservoir by discretizing the modes of the reservoir following the approach of Ullersma . In dealing with our system, there are crucial issues concerning obtaining the correct coupling strength between the oscillator and the reservoir modes, as well as in employing the proper renormalization and mode sampling in numerical simulations. When these criteria are satisfied, the SHO method comprises a powerful approach to treating radiative dynamics.
Here, we develop a rigorous quantitative treatment of the radiative dynamics of an electric dipole oscillator coupled to the electromagnetic reservoir within a model PC. In the process, we provide a sound theoretical basis for this and other approaches to non-Markovian radiative dynamics which involve the discretization of a model electromagnetic reservoir. Additionally, we show how our method can be applied to realistic PC’s with complicated dispersion relations and EM mode structures. The paper is organized as follows. In Section II, we develop a classical field theory for electromagnetic field modes in PCs, and we derive the coupling constants for the coupling between a radiating dipole and these Bloch modes. This leads to the Hamiltonian of the coupled system and the associated equations of motion. Renormalization issues arising from the non–relativistic nature of our theory are discussed in Section III, whereas Section IV describes the discretization of the reservoir and the numerical solution of the equations of motion. In Section V, these techniques are applied to a highly computationally challenging model, that of a three–dimensional, isotropic dispersion relation with a complete PBG. The demonstration of fractional localization and related phenomena validates the SHO approach to modeling radiative dynamics in PCs. In Section VI we summarize the results and emphasize the possibilities for testing these predictions experimentally in the microwave domain. The two appendices are concerned with the details of the field theory for the PC and with the details of the model of the one–sided, isotropic PBG, respectively.
## II Classical field theory
In this section, we derive the equations governing the dynamics of a radiating dipole oscillator located inside a PC. Typically the equation of motion for a damped oscillator, with time-dependent coordinate $`q(t)`$, is written as the second–order differential equation
$$\ddot{q}(t)+\gamma \dot{q}(t)+\omega _0^2q(t)=F(t).$$
(1)
Here, we have introduced a damping constant $`\gamma `$, the natural frequency $`\omega _0`$ and the driving field $`F(t)`$ for the amplitude $`q`$ of the linear oscillator. For instance, for a freely oscillating $`RLC`$ circuit with ohmic resistance $`R`$, capacitance $`C`$ and inductance $`L`$, we have $`\gamma =R/L`$, $`\omega _0^2=1/LC`$, $`F(t)=0`$, and $`q(t)`$ is the electric charge. Eq. (1) is, however, not the most general way of incorporating damping into the equations of motion for a harmonic oscillator. This description can break down if, for example, there is a suppression of modes in the reservoir to which the dipole oscillator can couple. Such a suppression of modes is a feature of the EM reservoir present in a PC. A more general description of damping forces acting on the harmonic oscillator therefore requires a precise knowledge of the mode structure of its environment, and the corresponding coupling of the system oscillator to these modes. In the case of a radiating dipole located in a PC, it is then appropriate to model its emission dynamics with a SHO coupled to a reservoir of SHOs. The essential difference between the vacuum and a PC is then contained in the spectral distribution, or density of states (DOS), of the reservoir oscillators, and in the coupling constants between the reservoir modes and the system oscillator.
The characterization of the reservoir is carried out in detail in Appendix A; here we only summarize the salient results. Given a radiating dipole with a natural frequency $`\omega _0`$, we obtain the classical Hamiltonian
$$H=H_{\mathrm{dip}}+H_{\mathrm{res}}+H_{\mathrm{ct}}+H_{\mathrm{int}}.$$
(2)
The first term on the right–hand side of the Hamiltonian is the energy of the dipole oscillator itself,
$$H_{\mathrm{dip}}=\xi \omega _0|\alpha |^2.$$
(3)
The natural frequency of the isolated oscillator is $`\omega _0`$, and $`\xi `$ is a constant with the dimension of energy $`\times `$ time. This permits us to write the energy of a SHO in units of its natural frequency $`\omega `$, i.e., $`E(\omega )=\xi \omega `$. The system oscillator’s complex amplitude is given by the dimensionless, time–dependent quantity $`\alpha `$, defined with respect to the coordinate $`q(t)`$ of Eq. (1) as
$$\alpha (t)\sqrt{\frac{L\omega _0}{2\xi }}q(t)+ı\sqrt{\frac{1}{2\xi L\omega _0}}(L\dot{q}(t)).$$
(4)
The next term in the Hamiltonian (2) corresponds to the free evolution of the radiation reservoir, which is modeled as a bath of independent SHOs,
$$H_{\mathrm{res}}=\underset{\mu }{}\xi \omega _\mu |\beta _\mu |^2.$$
(5)
The natural electromagnetic modes of the PC are Bloch modes (see Appendix A), labeled with the index $`\mu (n\stackrel{}{k})`$, where $`n`$ stands for the band index and $`\stackrel{}{k}`$ is a reciprocal lattice vector that lies in the first Brillouin zone (BZ). Their dispersion relation, $`\omega _\mu `$, is different from the vacuum case, and may have complete gaps and/or the corresponding density of states may exhibit appreciable pseudogap structure, the manifestation of multiple (Bragg) scattering effects in periodic media.
As we are working within the framework of a non-relativistic field theory, we have introduced a mass renormalization counter term $`H_{\mathrm{ct}}=\xi \mathrm{\Delta }|\alpha |^2`$ that cancels unphysical UV-divergent terms . The quantity $`\mathrm{\Delta }`$ is specified in Section III.
The interaction between the oscillator and the reservoir is given by a linear coupling term. As the oscillator frequency is quite large, and the effective linewidth of the oscillation is relatively small, it is possible to simplify the interaction by applying the rotating–wave approximation. In this approximation, couplings in the Hamiltonian of the form $`\alpha ^{}\beta _\mu ^{}`$ and its complex conjugate are neglected, as these terms oscillate very rapidly compared to the terms of the type $`\alpha ^{}\beta _\mu `$ and its conjugate. Hence, the interaction Hamiltonian can be expressed as
$$H_{\mathrm{int}}=ı\xi \underset{\mu }{}\left(\alpha ^{}g_\mu ^{}\beta _\mu \alpha g_\mu \beta _\mu ^{}\right).$$
(6)
In the case of a point dipole, i.e., when its spatial extent $`a`$ is much smaller than the wavelength corresponding to its natural frequency, $`\lambda _0=2\pi \omega _0/c`$, the coupling constants $`g_\mu `$ can be derived from (i) the magnitude of the dipole moment, $`d(t)=aq(t)`$, located at $`\stackrel{}{r}_0`$, and (ii) the dipole orientation, $`\widehat{d}`$, relative to that of the Bloch modes, $`\stackrel{}{E}_\mu (\stackrel{}{r}_0)`$:
$$g_\mu g_\mu (\stackrel{}{r}_0)=ac\sqrt{\frac{\pi }{L\omega _0\omega _\mu }}\left(\widehat{d}\stackrel{}{E}_\mu ^{}(\stackrel{}{r}_0)\right).$$
(7)
This dependence of the coupling constant on the dipole’s precise location within the PC is the second essential difference from the free–space case. As shown in Refs. , this position dependence may be quite strong, thus making its incorporation a sine qua non for any quantitative theory of of radiating antennæ or fluorescence phenomena in realistic PCs.
The emission dynamics can be evaluated from the Poisson brackets of the oscillator amplitudes and their initial values, $`\alpha (0)=1`$ and $`\beta _\mu (0)=0`$ $`(\mu )`$. Our choice of $`\alpha (0)`$ and $`\beta _\mu (0)`$ corresponds to the initial condition of an excited dipole antenna and a completely de–excited bath. The only non–zero Poisson brackets are
$$\{\alpha ,\alpha ^{}\}=\{\beta _\mu ,\beta _\mu ^{}\}=\frac{ı}{\xi }.$$
(8)
Eqs. (2), (7) and (8), together with the initial values for the oscillator amplitudes, completely determine the emission dynamics of a radiating dipole embedded in a PC. In the following sections, we solve the corresponding equations of motion. This task is complicated by the nature of the reservoir’s excitation spectrum: as discussed, the non-smooth density of states prohibits the use of a Markovian approximation and its appealing simplifying features . Instead, we have to revert to a solution of the full non–Markovian problem. This is accomplished by firstly rearranging the reservoir modes in a manner more suitable to both analytical as well as numerical solutions, and subsequently solving the equations of motion. In what follows, we bridge the gap between previous studies of simplified model dispersion relations and band structure computations .
Although we will formally develop our theory for an LC circuit in a microwave PC, we emphasize that the formalism applies equally well to a semiclassical Lorentz oscillator model of an excited two-level atom, i.e., an electron with charge $`e`$ and mass $`m`$ which is bound to a stationary nucleus, for which the energy of excitation is well below that required for saturation effects to become relevant. The oscillator coordinate $`q(t)`$ may then be identified with the deviation of the electron’s position from its equilibrium value, $`\gamma `$ is the inverse life time of the excited state, and $`\omega _0`$ denotes the frequency for transitions between excited and ground state of the two-level atom. This corresponds to making the substitutions:
$$Lm,(L\dot{q})p,\xi \mathrm{},$$
(9)
where $`h=2\pi \mathrm{}`$ is Planck’s constant.
## III Projected Local Density of States, Mass renormalization and Lamb shift
From the Hamiltonian (2) we derive the equations of motion for the amplitudes
$`\dot{\alpha }(t)`$ $`=`$ $`ı\left(\omega _0\mathrm{\Delta }\right)\alpha (t)ı\xi {\displaystyle \underset{\mu }{}}g_\mu ^{}\beta _\mu (t)`$ (10)
$`\dot{\beta _\mu }(t)`$ $`=`$ $`ı\omega _\mu \beta _\mu (t)+g_\mu \alpha (t),`$ (11)
for which we seek a solution with initial conditions $`\alpha (0)=1`$ and $`\beta _\mu (0)=0`$ ($`\mu )`$. Our formalism however requires that we first determine the mass renormalization counter term $`\mathrm{\Delta }`$. This is most conveniently done in a rotating frame with slowly varying amplitudes $`a(t)`$ and $`b(t)`$, defined as $`\alpha (t)=a(t)e^{ı\omega _0t}`$ and $`\beta (t)=b(t)e^{ı\omega _\mu t}`$ respectively:
$`\dot{a}(t)`$ $`=`$ $`ı\xi {\displaystyle \underset{\mu }{}}g_\mu ^{}e^{ı(\omega _0\omega _\mu )t}b_\mu (t)+ı\mathrm{\Delta }a(t)`$ (12)
$`\dot{b}(t)`$ $`=`$ $`g_\mu e^{ı(\omega _o\omega _\mu )t}a(t).`$ (13)
Conversely, Eqs. (12) and (13) comprise a stiff set of differential equations making their solution a difficult task. Numerical solution of the problem is more easily performed in the non-rotating frame, to which we return in Sect. IV.
Eq. (13) may be formally integrated,
$$b_\mu (t)=g_\mu _0^t𝑑t^{}e^{ı(\omega _0\omega _\mu )t^{}}a(t^{}),$$
(14)
and inserted into Eq. (12) to yield
$$\dot{a}(t)=_0^{\mathrm{}}𝑑t^{}G(tt^{})a(t^{})+ı\mathrm{\Delta }a(t),$$
(15)
where the Green function $`G(\tau )`$ contains all the information about the reservoir and is the subject of our studies for the remainder of this section. It is defined as
$$G(\tau )\mathrm{\Theta }(\tau )\underset{\mu }{}|g_\mu |^2e^{ı(\omega _0\omega _\mu )\tau }.$$
(16)
Here, $`\mathrm{\Theta }(\tau )`$ denotes the Heaviside step function, which ensures the causality of $`G(\tau )`$. We now proceed to evaluate $`G(\tau )`$ for the form of the coupling constants $`g_\mu `$ given in Eq. (7). To this end, we introduce the projected local DOS (PLDOS) $`N(\stackrel{}{r}_0,\widehat{d},\omega )`$ through
$$N(\stackrel{}{r}_0,\widehat{d},\omega )=\underset{n}{}_{\mathrm{BZ}}\frac{d^3k}{(2\pi )^3}\delta (\omega \omega _{n\stackrel{}{k}})|\widehat{d}\stackrel{}{E}_{n\stackrel{}{k}}(\stackrel{}{r}_0)|^2,$$
(17)
where we have replaced the symbolic sum over $`\mu `$ by its proper representation as a sum over bands plus a wave vector integral over the BZ. With these changes, we may rewrite $`G(\tau )`$ compactly as
$$G(\tau )=\beta \mathrm{\Theta }(\tau )_0^{\mathrm{}}𝑑\omega \frac{N(\stackrel{}{r}_0,\widehat{d},\omega )}{\omega }e^{ı(\omega _0\omega )\tau }.$$
(18)
Here, we have abbreviated $`\beta =(\pi a^2c^2)/(L\omega _0)`$. Eq. (18) makes more explicit what we have argued before: The spontaneous emission dynamics of active media in Photonic Crystals are completely determined by the PLDOS, $`N(\stackrel{}{r}_0,\widehat{d},\omega )`$. As the PLDOS may be drastically different from location to location within the Wigner–Seitz cell of the PC , it is imperative to have detailed knowledge about where in the PC the dipole is situated in order to understand and predict the outcome of corresponding experiments.
One additional point deserves special attention: the total DOS, $`N(\omega )`$, is related to the local DOS via
$`N(\omega )`$ $`=`$ $`{\displaystyle \frac{1}{V}}{\displaystyle _V}d^3r{\displaystyle 𝑑\mathrm{\Omega }_{\widehat{d}}ϵ_p(\stackrel{}{r})N(\stackrel{}{r},\widehat{d},\omega )}`$ (19)
$``$ $`{\displaystyle \frac{1}{V}}{\displaystyle _V}d^3r{\displaystyle 𝑑\mathrm{\Omega }_{\widehat{d}}N(\stackrel{}{r},\widehat{d},\omega )},`$ (20)
where $`V`$ is the volume of the Wigner–Seitz cell, and $`𝑑\mathrm{\Omega }_{\widehat{d}}`$ is the average over all possible orientations of the dipole. Strictly speaking, it is not possible to base conclusions about the radiation dynamics on the total DOS. This is a direct consequence of the fact that the natural modes of PCs are Bloch waves rather than plane waves as in free space. Depending on the band index, these Bloch modes prefer to “reside” predominantly in either low or high dielectric index regions (so-called air and dielectric bands respectively). Only in the case of very low index contrast (“nearly free photons”) may the total DOS be viewed as a reliable guide to interpreting radiative dynamics within a PC. The total DOS is, nevertheless, an adequate rule-of-thumb estimator.
From Eq. (17) we can now obtain the Fourier transform of the Green function, $`G(\mathrm{\Omega }\omega _0)`$, centered around the atom’s bare transition frequency $`\omega _0`$:
$`G(\mathrm{\Omega }\omega _0)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑te^{ı(\mathrm{\Omega }\omega _0)t}G(t)`$ (21)
$`=`$ $`\pi \beta {\displaystyle \frac{N(\stackrel{}{r}_0,\widehat{d},\mathrm{\Omega })}{\mathrm{\Omega }}}\mathrm{\Theta }(\mathrm{\Omega })`$ (23)
$`+ı\beta {\displaystyle _0^{\mathrm{}}}𝑑\omega {\displaystyle \frac{N(\stackrel{}{r}_0,\widehat{d},\omega )}{\omega }}\mathrm{}\left({\displaystyle \frac{1}{\mathrm{\Omega }\omega }}\right),`$
where $`\mathrm{}`$ stands for the principal value.
For large $`\omega `$, we have $`N(\stackrel{}{r}_0,\widehat{d},\omega )\omega ^2`$. The imaginary part of $`G(\mathrm{\Omega }\omega _0)`$ apparently contains a linear divergence in the UV. This divergence is to be expected for a non-relativistic theory, analogous to the problem of spontaneous emission in vacuum , and is removed from the theory by using the mass counter renormalization term $`\mathrm{\Delta }`$, as first pointed out by Bethe . Consequently, we decompose the imaginary part of $`G(\mathrm{\Omega }\omega _0)`$ into
$$\mathrm{}\left(G(\mathrm{\Omega }\omega _0)\right)(\mathrm{\Delta }+\delta _{\mathrm{vac}}+\delta _a),$$
(24)
where we have used the notation:
$`\mathrm{\Delta }`$ $`=`$ $`\beta {\displaystyle _0^{\mathrm{}}}𝑑\omega {\displaystyle \frac{N(\stackrel{}{r}_0,\widehat{d},\omega )}{\omega ^2}}`$ (25)
$`\delta _{\mathrm{vac}}`$ $`=`$ $`{\displaystyle \frac{\beta \omega _0}{\pi ^2c^3}}{\displaystyle _0^{\mathrm{\Omega }_c}}𝑑\omega \mathrm{}\left({\displaystyle \frac{1}{\omega _0\omega }}\right)`$ (26)
$`\delta _\mathrm{a}`$ $`=`$ $`{\displaystyle \frac{\beta \omega _0}{\pi ^2c^3}}{\displaystyle _0^{\mathrm{\Omega }_c}}d\omega \mathrm{}\left({\displaystyle \frac{1}{\omega _0\omega }}\right)\times `$ (28)
$`\times {\displaystyle \frac{N(\stackrel{}{r}_0,\widehat{d},\omega )N^{(\mathrm{vac})}(\omega )}{\omega ^2}}.`$
Here, we have performed a Wigner–Weisskopf-type approximation on the vacuum and anomalous Lamb shifts , $`\delta _{\mathrm{vac}}`$ and $`\delta _\mathrm{a}`$, respectively. This approximation is justified by the fact that, despite its highly non–Markovian nature, a radiating dipole in a PC is still a weak coupling problem, as can be seen, for instance, by estimating the coupling constant
$$gd_0\omega _0\sqrt{\frac{2\pi }{V\xi \omega _{n\stackrel{}{k}}}}$$
(29)
in the Lorentz oscillator model. Here, $`V\overline{a}^3`$ is the volume of the Wigner–Seitz cell of the PC ($`\overline{a}`$ is the corresponding lattice constant) and $`d_0=ea_0`$ is the oscillator’s dipole moment for the elementary charge $`e`$ and Bohr atomic radius, $`a_0`$. At optical frequencies ($`\omega 10^{15}`$ s<sup>-1</sup>), a silicon inverted opal has a PBG at the frequency $`\overline{a}\omega /2\pi c0.8`$, so that we obtain $`10^7g/\omega _010^61`$, thus justifying our Wigner–Weisskopf approximation. As a consequence, we must treat the real part of $`G(\mathrm{\Omega }\omega _0)`$ exactly, but are still allowed to tackle the imaginary part of $`G(\mathrm{\Omega }\omega _0)`$ using standard perturbation methods of QED. In addition, we have introduced the vacuum or free–space DOS $`N^{(\mathrm{vac})}(\omega )=\omega ^2/(\pi ^2c^3)`$, and a cutoff frequency $`\mathrm{\Omega }_c\omega _0`$, which is chosen large enough that the results of the following analysis remain independent of the precise value of $`\mathrm{\Omega }_c`$. In a Lorentz oscillator model, for instance, $`\mathrm{\Omega }_c`$ can be identified with the Compton frequency $`\mathrm{\Omega }_cmc^2/\mathrm{}`$, as $`\omega >\omega _c`$ probes the relativistic aspects of the oscillating charge, which are beyond the scope of the model.
With the foregoing analysis, we have determined the mass renormalization counter term $`\mathrm{\Delta }`$. In addition, we have derived an explicit expression for the anomalous Lamb shift $`\delta _a`$ which originates in the “reshuffling” of the reservoir’s spectral weight by the PC.
## IV Discretization of the reservoir
To solve the equation of motion for the amplitude of the system oscillator, let us rewrite Eq. (15) in a more explicit form:
$`\dot{a}(t)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑\omega N(\stackrel{}{r}_0,\widehat{d},\omega )g^2(\omega ){\displaystyle _0^t}𝑑t^{}e^{ı(\omega _0\omega )(tt^{})}a(t^{})`$ (31)
$`+ı\mathrm{\Delta }a(t),`$
where $`g^2(\omega )=\beta /\omega `$, and the mass renormalization counter term $`\mathrm{\Delta }`$ is given by
$$\mathrm{\Delta }=\beta _0^{\mathrm{}}𝑑\omega \frac{N(\stackrel{}{r}_0,\widehat{d},\omega )}{\omega ^2}.$$
(32)
We remind the reader that $`a(0)=1`$.
We are now in a position to comment on the origin of the linear damping term $`\gamma \dot{q}(t)`$ that appears in Eq. (1): If we consider the long time limit, i. e., $`t1/\omega _0`$, and assume that the PLDOS $`N(\stackrel{}{r}_0,\widehat{d},\omega )`$ is a smooth function for frequencies around $`\omega _0`$, we can approximate the frequency integral in Eq. (31) by $`\left[2\pi \beta N(\stackrel{}{r}_0,\widehat{d},\omega _0)/\omega _0\right]\delta (tt^{})`$, which leads to
$$\dot{a}(t)=\gamma a(t),$$
(33)
where the decay constant is defined as
$$\gamma =2\pi \beta N(\stackrel{}{r}_0,\widehat{d},\omega _0)/\omega _0.$$
(34)
This approximation is is valid only for long times relative to $`1/\omega _0`$, and for a sufficiently smooth density of states. However, in the case of a PC, the PLDOS may have sharp discontinuities and gaps, thus requiring that the full equations of motion instead.
To solve the integro-differential equation (31) in a PC, we appeal to the literal meaning of the PLDOS as a density of states: $`N(\stackrel{}{r}_0,\widehat{d},\omega )`$ may be interpreted as an unnormalized probability density of finding a reservoir oscillator with frequency $`\omega `$ at position $`\stackrel{}{r}_0`$ and orientation $`\widehat{d}`$. Consequently, we transform Eq. (31) back to a system of coupled differential equations by employing a Monte Carlo integration scheme for an arbitrary function $`f(\omega )`$ according to
$`{\displaystyle _0^{\mathrm{}}}𝑑\omega N(\stackrel{}{r}_0,\widehat{d},\omega )f(\omega )`$ $``$ $`{\displaystyle _0^{\mathrm{\Omega }_c}}𝑑\omega N(\stackrel{}{r}_0,\widehat{d},\omega )f(\omega )`$ (35)
$``$ $`{\displaystyle \frac{N_0}{M}}{\displaystyle \underset{i=1}{\overset{M}{}}}f(\omega _i),`$ (36)
where the normalization constant
$$N_0=_0^{\mathrm{\Omega }_c}𝑑\omega N(\stackrel{}{r}_0,\widehat{d},\omega )$$
(37)
depends on the cutoff frequency, $`\mathrm{\Omega }_c`$. There are $`M1`$ bath oscillators, contained within a set of frequencies $`\{\omega _i,1iM\}`$, the frequencies of which are obtained by randomly sampling the interval $`[0,\mathrm{\Omega }_c]`$ according to the probability density $`p(\stackrel{}{r}_0,\widehat{d},\omega )=N(\stackrel{}{r}_0,\widehat{d},\omega )/N_0`$. Note that the $`\omega _i`$ may be degenerate, as prescribed by $`p(\stackrel{}{r}_0,\widehat{d},\omega )`$.
Applying this Monte Carlo scheme to Eq. (31) and transforming back to a non-rotating frame in order to avoid having to solve a numerically stiff problem, we obtain
$`\dot{\alpha }(t)`$ $`=`$ $`ı\left(\omega _0\mathrm{\Delta }\right)\alpha (t)ı\xi {\displaystyle \underset{i=1}{\overset{N}{}}}g_i\beta _i(t)`$ (38)
$`\dot{\beta _i}(t)`$ $`=`$ $`ı\omega _i\beta _i(t)+g_i\alpha (t),`$ (39)
where $`g_i=g(\omega _i),1iM`$, and the mass renormalization counter term is evaluated up to the cutoff frequency $`\mathrm{\Omega }_c`$, i.e., $`\mathrm{\Delta }=_0^{\mathrm{\Omega }_c}𝑑\omega N(\stackrel{}{r}_0,\widehat{d},\omega )/\omega ^2`$.
When comparing Eqs. (38) and (39) to our initial equations of motion, Eqs. (10) and (11), we observe that the considerations in the previous section have allowed us to rearrange the three-dimensional wave vector sum over the modes $`\mu (n\stackrel{}{k})`$ into a simple one-dimensional sum over a set of frequencies $`\{\omega _i\}`$ with a probability distribution $`p(\stackrel{}{r}_0,\widehat{d},\omega )`$ that is easily determined through standard photonic band structure computation . In the following section, we give the solutions of (38) and (39) for a model system which has previously been treated by other methods. In particular, we will demonstrate that known results for the radiative dynamics can be recaptured and do not depend on the the value of the cutoff frequency $`\mathrm{\Omega }_c`$ and the number $`M`$ of reservoir oscillators once these quantities are large enough such that all the relevant features of $`N(\stackrel{}{r}_0,\widehat{d},\omega )`$, are adequately represented.
## V Numerical Results for a model system
In order to establish the validity of our approach, we now solve Eqs. (38) and (39) for a generic model of a PBG, the three-dimensional isotropic, one-sided PBG . In Appendix B, we outline the construction of the model’s dispersion relation and how to obtain the corresponding model DOS, $`N_m(\omega )`$. We note that we do not appeal to an effective mass approximation in the dispersion relation , as is done in most treatments of band–edge dynamics. This allows us to recover the correct form of the large frequency behavior of the photon density of states.
In Fig. 1, we show the behavior of $`N_m(\omega )`$ as a function of frequency for values of the relevant parameters, the gap size parameter $`\eta =0.8`$ and the normalized center frequency $`\omega _ca/2\pi c=0.5`$ (see Appendix B). The DOS exhibits a square-root singularity at the band edge $`\omega _ua/2\pi c=0.6`$, as well as a UV divergence, $`N_m(\omega )\omega ^2`$, as $`\omega \mathrm{}`$; these are the characteristic features of this model. Due to the simultaneous presence of both divergences, this model clearly represents a severe numerical test of our approach. In order to test the method, we thus replace the PLDOS entering Eqs. (38) and (39) by $`N_m(\omega )`$.
In Fig. 2, we present the results of our numerical solution for the radiation dynamics of a dipole oscillator with frequency $`\omega _0`$ that is coupled to the modes of a PC, as described by Eqs. (38) and (39), for various values of the bare oscillator frequency, $`\omega _0a/2\pi c`$, relative to the bandedge at $`\omega _ua/2\pi c=0.6`$. The coupling strength has been chosen such that $`g(\omega _0)=10^4`$, corresponding to $`\beta =10^8\times \omega _0^3`$.
Clearly visible are normal mode oscillations, also referred to as vacuum Rabi oscillations, and the fractional localization of the oscillator’s energy at long times near the photonic band–edge . As expected, for frequencies deep in the photonic band–gap ($`\omega _0a/2\pi c=0.58`$), where the system oscillator is effectively decoupled from the bath oscillators, we find no noticeable decay of the oscillator amplitude. Deep in the photonic conduction band ($`\omega _0a/2\pi c=0.62`$), the system oscillator is coupled to a bath with a smooth and slowly–varying mode density, as in free space. We therefore observe exponential decay of the oscillator amplitude, though with a time scale that differs significantly from that in free space. Due to the large value of the DOS close to the photonic band edge, the initial decay is faster for bare oscillator frequencies close to this edge than for frequencies deep inside the allowed photonic band. These results were obtained for a smooth exponential cutoff for the DOS around $`\mathrm{\Omega }_ca/2\pi c=3.0`$ and $`M=2.5\times 10^5`$ oscillators representing the modes of the PC. We also performed numerical simulations between all combinations of $`\mathrm{\Omega }_c`$ and $`M`$ with values $`\mathrm{\Omega }_ca/2\pi c=3.0,6.0,9.0`$ and $`M=2.5\times 10^5,5\times 10^5,10^6`$ and found that the numerical values differ by at most 0.2% of the values shown in Fig 1. This demonstrates that, despite the presence of the singularities in the DOS, our approach still provides accurate and convergent results.
## VI Discussion
In summary, we have developed a realistic field theory for an oscillating electric dipole located in a PC. The theory is based on the natural modes of the PC, the Bloch waves, and allows the direct incorporation of realistic band structure calculations in order to obtain quantitative results for the radiation dynamics of the dipole antenna. We have shown how the theory must be renormalized in order to account for unphysical divergences and have identified the classical analogue of the Lamb shift of the dipole’s natural radiation frequency. Finally, we have developed a reliable numerical scheme based on a probability interpretation of the PLDOS that solves the equations of motion for the dipole oscillator coupled to the electromagnetic mode reservoir of the PC.
The viability of this approach was demonstrated for an isotropic model DOS for which we have derived well-known results for radiating atomic systems in the context of a radiating classical dipole. The model considered contains two divergences, one square-root-divergence at the photonic band edge and a quadratic UV-divergence, and therefore clearly comprises the most serious test of our approach. More realistic models of a three dimensional photonic band-edge take into account the anisotropy of the BZ, and therefore do not suffer from a band–edge singularity . As a result, our formalism is clearly more than capable of treating more realistic descriptions of the electromagnetic reservoir within a PC.
Though we have developed our theory for an LC circuit in a microwave PC, we have pointed out in Section II that the formalism applies equally well to a semiclassical Lorentz oscillator model of an excited two-level atom. Therefore, our approach is applicable to both microwave antennæ and to optical atomic transitions. However, technological constraints suggest that microwave experiments will likely be easier to perform than optical experiments involving single atoms. As discussed, an appropriate microwave antenna could, for example, take the form of a high-Q metallic pin placed in or near a PC. The pin can then be excited by a focused ultrashort laser pulse that generates free carriers at one end; these carriers then undergo several oscillations across the pin before re–establishing charge equilibrium. The resulting signal could be easily detected and compared with the emission from such an antenna positioned in free space, or within a homogeneous sample of the dielectric material that makes up the backbone of the PC under consideration.
In its own right, such a microwave system could have considerable applications in radio science and microwave technology. For example, the PBG can be used as a frequency filter, and can be used to fine tune the bandwidth of a dipole emitter with a resonant frequency near the edge of the gap. It may also be possible to actively modify the photonic band structure, effectively changing the radiation pattern of a dipole emitter. A feasible scheme for active band structure modification has recently been proposed in the context of optical PCs , in which the PC is infiltrated with a liquid crystalline material whose nematic director is aligned using applied electric fields. By rotating the director, it was found that the band structure could be significantly modified, and that PBGs may be opened and closed altogether. Similar methods may be applied to the case of microwave PCs.
Although we have concentrated specifically on the linear model, the method of coupled oscillators can be extended to treat a nonlinear antennæ , or a collection of two–level atoms in a regime where saturation effects arise. As we have shown here, this method of coupled classical oscillators reproduces effects normally associated with quantum optical calculations. We expect that a nonlinear oscillator model will reproduce some of the effects studied for a single two–level atom coupled to the modes of a PC without the need for quantizing the field. However, a classical treatment would need to be abandoned if multiphoton excitations are non–negligible . Given that multiphoton effects are difficult to observe in the microwave domain and even more challenging in the optical domain, it is reasonable to expect that a classical model of radiative dynamics in a PC should be sufficient for foreseeable experiments.
## ACKNOWLEDGMENTS
We are grateful to K.-J. Boller and R. Beigang for stimulating discussions concerning the experimental realization of radiating dipoles in microwave PCs. The work of KB the was supported by the Deutsche Forschungsgemeinschaft (DFG) under Grant Bu 1107/1-1. NV acknowledges support from the Ontario Graduate Scholarship Program. BCS acknowledges support from the Department of Physics, University of Toronto, and the support of an Australian Research Council Large Grant. This work was supported in part by the New Energy and Industrial Technology Development Organization of Japan and by Photonics Research Ontario.
## A Classical field theory for Photonic Crystals
In this Appendix, we present a self-contained formulation of a classical field theory for the Bloch modes of a PC, and we develop the Hamiltonian describing the coupling of a radiating dipole couples to these modes. As a first step, we review the computation of dispersion relations, and of electric and magnetic field modes from band structure calculations . We then demonstrate how the results of such band structure calculations can be used to construct the corresponding vector potentials and free field Hamiltonian. Finally, we derive the full minimal coupling Hamiltonian for a classical radiating dipole embedded in a PC. This may be compared to the formulation of a general, quantized field theory for EM modes in nonuniform dielectric media in terms of a normal mode expansion in the context of quantum optics .
### 1 Review of band structure calculations
We develop our theory in terms of the magnetic field $`\stackrel{}{H}`$ rather than in terms of the electric or displacement fields because (i) $`\stackrel{}{H}=0`$ and, (ii) the transverse and longitudinal components of the magnetic field are continuous across the dielectric boundaries. This leads to more rapid convergence of the relevant Fourier series expansions.
In a three-dimensional PC, we can write the eigenvalue equation for the magnetic field $`\stackrel{}{H}`$ as
$$\times \left(\eta _p(\stackrel{}{r})\times \stackrel{}{H}\right)+\frac{\omega ^2}{c^2}\stackrel{}{H}=\stackrel{}{0}$$
(A1)
with $`\eta _p(\stackrel{}{r})`$ the inverse of the periodic dielectric permittivity,
$$ϵ_p(\stackrel{}{r})=ϵ_b+(ϵ_aϵ_b)\underset{\stackrel{}{n}𝒵^3}{}S(\stackrel{}{r}\stackrel{}{n}𝑨).$$
(A2)
The medium is assumed to consist of a background material with bulk permittivity $`ϵ_b`$ and a set of scatterers, with bulk permittivity $`ϵ_a`$. The shape of the scatterers is described by the function $`S`$, i. e.,, $`S(\stackrel{}{r})=1`$ if $`\stackrel{}{r}`$ lies inside the scatterer and zero elsewhere, distributed periodically at positions
$$\left\{\stackrel{}{R}\right\}=\left\{\underset{i=1}{\overset{3}{}}n_i\stackrel{}{a}_i\right|n_i𝒵\}.$$
(A3)
The notation of Eq. (A2) is obtained by defining the matrix $`𝑨=(\stackrel{}{a}_1\stackrel{}{a}_2\stackrel{}{a}_3)`$ and $`𝒵^3=𝒵𝒵𝒵`$. The dielectric permittivity is spatially periodic modula $`\stackrel{}{n}𝑨`$. The assumption of a scalar permittivity is reasonable for bulk materials which are not birefringent but in no way restricts the considerations below. Chromatic dispersion effects are considered to be negligible, thus allowing the time-dependence of the permittivity to be ignored. Let us define the dual matrix $`𝑩=2\pi (𝑨^1)^T`$. For $`𝑩=(\stackrel{}{b}_1\stackrel{}{b}_2\stackrel{}{b}_3)`$, this definition leads to the orthogonality relation
$$\stackrel{}{a}_i\stackrel{}{b}_j=2\pi \delta _{ij}.$$
(A4)
Whereas the points $`\stackrel{}{n}𝑨`$ are the real space lattice vectors, the points $`\stackrel{}{m}𝑩`$, for $`\stackrel{}{m}𝒵^3`$ are the reciprocal lattice vectors. The inverse permittivity can be expanded in the dual basis as
$$\eta _p(\stackrel{}{r})=\underset{\stackrel{}{m}𝒵^3}{}\eta _\stackrel{}{m}e^{ı\stackrel{}{m}𝑩\stackrel{}{r}}.$$
(A5)
The differential equation (A1) has periodic coefficients. By the Bloch-Floquet theorem we can expand the magnetic field as
$$\stackrel{}{H}_\stackrel{}{k}=e^{ı\stackrel{}{k}\stackrel{}{r}}\stackrel{}{u}_\stackrel{}{k}(\stackrel{}{r})$$
(A6)
where $`\stackrel{}{u}_\stackrel{}{k}`$ is spatially periodic modulo $`𝑨`$; that is,
$$\stackrel{}{u}_\stackrel{}{k}(\stackrel{}{r})=\stackrel{}{u}_\stackrel{}{k}(\stackrel{}{r}+\stackrel{}{n}𝑨).$$
(A7)
The set $`\{\stackrel{}{k}\}`$ labeling the solutions can be restricted to lie within in the irreducible part of the first Brillouin zone (BZ), since any value of $`\stackrel{}{k}`$ can then be obtained through a combination of group transformations with respect to an operation from the point group of the crystal and translations with respect to a reciprocal lattice vector. We can therefore express each wavevector $`\stackrel{}{k}`$ as
$$\stackrel{}{k}\stackrel{}{k}_{𝑻,\stackrel{}{m}}\stackrel{}{k}_{}𝑻+\stackrel{}{m}𝑩,$$
(A8)
where $`\stackrel{}{k}_{}`$ is an element of the irreducible part of the 1. BZ and $`𝑻`$ an element of the crystal’s point group.
Applying the Bloch-Floquet theorem, Eq. (A6), the magnetic field can be expanded as
$$\stackrel{}{H}_\stackrel{}{k}=e^{ı\stackrel{}{k}\stackrel{}{r}}\underset{\stackrel{}{m}}{}\underset{\lambda =1}{\overset{2}{}}h_\stackrel{}{m}^{\stackrel{}{k},\lambda }\widehat{e}_\stackrel{}{m}^{\stackrel{}{k},\lambda }e^{ı\stackrel{}{m}𝑩\stackrel{}{r}}.$$
(A9)
Here $`\lambda `$ is the index of polarization and the vectors
$$\{\widehat{e}_\stackrel{}{m}^{\stackrel{}{k},1},\widehat{e}_\stackrel{}{m}^{\stackrel{}{k},2},\frac{\stackrel{}{k}+\stackrel{}{m}𝑩}{|\stackrel{}{k}+\stackrel{}{m}𝑩|}\}$$
(A10)
form an orthonormal right-handed triad. This expansion inserted into Eq (A1) yields an infinite eigenvalue problem which is then solved numerically by a suitable truncation. Typically the cardinality of the set $`\{\stackrel{}{m}\}`$ is on the order of $`10^3`$ . For any given $`\stackrel{}{k}_{}`$ we obtain a discrete set of eigenfrequencies $`\omega _{n\stackrel{}{k}}`$ and corresponding eigenfunctions $`H_{n\stackrel{}{k}}`$ which we label by the band index $`n𝒩`$. It is important to note that the expression for the electric field can be recovered from the magnetic field via
$$\stackrel{}{E}_{n\stackrel{}{k}}(\stackrel{}{r})=i\frac{c}{\omega _{n\stackrel{}{k}}ϵ_p(\stackrel{}{r})}\times \stackrel{}{H}_{n\stackrel{}{k}}(\stackrel{}{r})$$
(A11)
In addition, the Bloch waves obey the following orthogonality relations:
$`{\displaystyle d^3r\stackrel{}{H}_{n\stackrel{}{k}}^{}(\stackrel{}{r})\stackrel{}{H}_{m\stackrel{}{k}^{}}(\stackrel{}{r})}`$ $``$ $`\delta _{nm}\delta (\stackrel{}{k}\stackrel{}{k}^{}),`$ (A12)
$`{\displaystyle d^3rϵ_p(\stackrel{}{r})\stackrel{}{E}_{n\stackrel{}{k}}^{}(\stackrel{}{r})\stackrel{}{E}_{m\stackrel{}{k}^{}}(\stackrel{}{r})}`$ $``$ $`\delta _{nm}\delta (\stackrel{}{k}\stackrel{}{k}^{}),`$ (A13)
where the integration is over all space in both cases. We are free to choose the constants of proportionality in the above relations, and do so in the next subsection.
### 2 Free–field Hamiltonian
Based on the above considerations, we are now in a position to derive the general expressions for the scalar and vector potential, $`\varphi (\stackrel{}{r},t)`$ and $`\stackrel{}{A}(\stackrel{}{r},t)`$ respectively, for the classical Hamiltonian of the free radiation field. We find that the expressions become particularly transparent in the Dzyaloshinsky gauge, i.e., when $`\varphi (\stackrel{}{r},t)0`$. Then,
$`\stackrel{}{E}(\stackrel{}{r},t)`$ $`=`$ $`{\displaystyle \frac{1}{c}}{\displaystyle \frac{\stackrel{}{A}(\stackrel{}{r},t)}{t}},`$ (A14)
$`\stackrel{}{H}(\stackrel{}{r},t)`$ $`=`$ $`\times \stackrel{}{A}(\stackrel{}{r},t),`$ (A15)
and the gauge condition $`\left(ϵ_p(\stackrel{}{r})\stackrel{}{A}(\stackrel{}{r},t)\right)=0`$, reveals that in a PC the natural modes of the radiation field are no longer transverse. This is of importance when quantizing the field theory . Given Eqs. (A1), (A11), (A14) and (A15), it is now straightforward to derive the following expansion of the vector potential $`\stackrel{}{A}(\stackrel{}{r},t)`$
$`\stackrel{}{A}(\stackrel{}{r},t)`$ $`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle _{\mathrm{BZ}}}{\displaystyle \frac{d^3k}{(2\pi )^3}}\sqrt{{\displaystyle \frac{2\pi \xi c^2}{\omega _{n\stackrel{}{k}}}}}`$ (A16)
$`\left(\beta _{n\stackrel{}{k}}(t)\stackrel{}{A}_{n\stackrel{}{k}}(\stackrel{}{r})+\beta _{n\stackrel{}{k}}^{}(t)\stackrel{}{A}_{n\stackrel{}{k}}^{}(\stackrel{}{r})\right),`$ (A17)
where the time evolution of the free field is described by $`\beta _{n\stackrel{}{k}}(t)=\beta _{n\stackrel{}{k}}(0)e^{ı\omega _{n\stackrel{}{k}}t}`$. The field modes $`\stackrel{}{A}_{n\stackrel{}{k}}(\stackrel{}{r})`$ obey
$$\times \times \stackrel{}{A}_{n\stackrel{}{k}}(\stackrel{}{r})=\frac{\omega _{n\stackrel{}{k}}^2}{c^2}ϵ_p(\stackrel{}{r})\stackrel{}{A}_{n\stackrel{}{k}}(\stackrel{}{r}),$$
(A18)
which is the same equation as that for the electric field modes $`\stackrel{}{E}_{n\stackrel{}{k}}(\stackrel{}{r})`$ of Eq. (A11). We now choose the normalization of $`\stackrel{}{A}_{n\stackrel{}{k}}`$ such that
$$d^3rϵ_p(\stackrel{}{r})\stackrel{}{A}_{n\stackrel{}{k}}(\stackrel{}{r})\stackrel{}{A}_{m\stackrel{}{k}^{}}(\stackrel{}{r})=\delta _{nm}\delta (\stackrel{}{k}\stackrel{}{k}^{}),$$
(A19)
$`{\displaystyle d^3r\left(\times \stackrel{}{A}_{n\stackrel{}{k}}(\stackrel{}{r})\right)\left(\times \stackrel{}{A}_{m\stackrel{}{k}^{}}(\stackrel{}{r})\right)}`$ $`=`$ (A20)
$`{\displaystyle \frac{\omega _{n\stackrel{}{k}}^2}{c^2}}\delta _{nm}\delta (\stackrel{}{k}\stackrel{}{k}^{}).`$ (A21)
This also fixes the normalization in Eqs. (A12) and (A13). As a consequence, the total electric and magnetic field are now given by
$`\stackrel{}{E}(\stackrel{}{r},t)`$ $`=`$ $`ı{\displaystyle \underset{n}{}}{\displaystyle _{\mathrm{BZ}}}{\displaystyle \frac{d^3k}{(2\pi )^3}}\sqrt{{\displaystyle \frac{2\pi \xi c^2}{\omega _{n\stackrel{}{k}}}}}`$ (A23)
$`\left(\beta _{n\stackrel{}{k}}(t)\stackrel{}{E}_{n\stackrel{}{k}}(\stackrel{}{r})\beta _{n\stackrel{}{k}}^{}(t)\stackrel{}{E}_{n\stackrel{}{k}}^{}(\stackrel{}{r})\right),`$
$`\stackrel{}{H}(\stackrel{}{r},t)`$ $`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle _{\mathrm{BZ}}}{\displaystyle \frac{d^3k}{(2\pi )^3}}\sqrt{{\displaystyle \frac{2\pi \xi c^2}{\omega _{n\stackrel{}{k}}}}}`$ (A25)
$`\left(\beta _{n\stackrel{}{k}}(t)\stackrel{}{H}_{n\stackrel{}{k}}(\stackrel{}{r})+\beta _{n\stackrel{}{k}}^{}(t)\stackrel{}{H}_{n\stackrel{}{k}}^{}(\stackrel{}{r})\right),`$
where we have re–introduced the electric and magnetic field modes, $`\stackrel{}{E}_{n\stackrel{}{k}}(\stackrel{}{r})=(\omega _{n\stackrel{}{k}}/c)\stackrel{}{A}_{n\stackrel{}{k}}(\stackrel{}{r})`$ and $`\stackrel{}{H}_{n\stackrel{}{k}}(\stackrel{}{r})=\times \stackrel{}{A}_{n\stackrel{}{k}}(\stackrel{}{r})`$, respectively. Eqs. (A23) and (A25) finally lead us to the free field Hamiltonian
$$H_{\mathrm{res}}=\underset{n}{}_{\mathrm{BZ}}d^3k\xi \omega _{n\stackrel{}{k}}|\beta _{n\stackrel{}{k}}|^2.$$
(A26)
The only nonzero Poisson brackets are $`\{\beta _{n\stackrel{}{k}},\beta _{n\stackrel{}{k}}^{}\}=ı/w`$.
### 3 Radiating dipole embedded in a Photonic Crystal
We consider the insertion of a point dipole into a PBG structure at a prescribed location $`\stackrel{}{r}_0`$. The free dipole oscillator is described by the Hamiltonian $`H_{\mathrm{dip}}`$
$$H_{\mathrm{dip}}=\frac{L\dot{q}^2}{2L}+\frac{1}{2}L\omega _0^2q^2=\xi \omega _0|\alpha |^2,$$
(A27)
where the dipole’s natural frequency is $`\omega _0=1/LC`$ and the complex oscillator amplitude $`\alpha `$ is given in terms of the charge $`q`$ and “current” $`L\dot{q}`$ as $`\alpha (t)=q(t)\sqrt{L\omega _0/2w}+ı(L\dot{q}(t))/\sqrt{2\xi L\omega _0}`$, with Poisson brackets $`\{\alpha ,\alpha ^{}\}=ı/\xi `$. The point dipole couples to the electric field via its dipole moment $`d(t)=aq(t)`$ with orientation $`\widehat{d}`$, which yields the interaction energy
$$H_{\mathrm{int}}=aq(t)\left(\widehat{d}\stackrel{}{E}(\stackrel{}{r}_0,t)\right).$$
(A28)
In the rotating wave approximation to the interaction term, the final minimal coupling Hamiltonian for a radiating dipole in a PC is
$$H=H_{\mathrm{dip}}+H_{\mathrm{res}}+H_{\mathrm{ct}}+H_{\mathrm{int}}.$$
(A29)
Collecting all the above results we obtain
$`H`$ $`=`$ $`\xi \omega _0|\alpha |^2+{\displaystyle \underset{\mu }{}}\xi \omega _\mu |\beta _\mu |^2+H_{\mathrm{ct}}`$ (A31)
$`ı\xi {\displaystyle \underset{\mu }{}}\left(\alpha ^{}g_\mu ^{}\beta _\mu \alpha g_\mu \beta _\mu ^{}\right).`$
Here, we have introduced the symbolic index $`\mu (n\stackrel{}{k})`$ and the coupling constants $`g_\mu `$
$$g_\mu g_\mu (\stackrel{}{r}_0)=ac\sqrt{\frac{\pi }{L\omega _0\omega _\mu }}\left(\widehat{d}\stackrel{}{E}_\mu ^{}(\stackrel{}{r}_0)\right).$$
(A32)
In addition, in Eq. (A29) we have introduced a mass renormalization counter term, $`H_{\mathrm{ct}}=\xi \mathrm{\Delta }|\alpha |^2`$ in order to cancel unphysical UV-divergent terms of our non-relativistic theory, as discussed in the main text. For completeness, we list here only the nonzero Poisson brackets and initial values for an initially excited radiating dipole coupled to the Bloch waves of a PC. This, together with the Hamilton function $`H`$ in Eq. (A29) completely defines our problem:
$$\{\alpha ,\alpha ^{}\}=\{\beta _\mu ,\beta _\mu ^{}\}=\frac{ı}{\xi },$$
(A33)
where $`\alpha (0)=1`$ and $`\beta _\mu (0)=0`$ for all $`\mu `$.
## B Model dispersion relation and density of states
A particularly stringent test of our approach’s ability to describe the dynamics of a radiating dipole in a PC comes from its application to a dipole coupled to a 3D isotropic photon dispersion model for the electromagnetic reservoir. In this model, the coherent scattering condition that characterizes the photonic band edge is assumed to occur at the same frequency for all directions of propagation. Clearly this is not the case in a real crystal, whose Brillouin zone cannot have full rotational symmetry. As a result, the isotropic model overestimates the electromagnetic modes available at a band–edge, so that, for example, near the upper photonic band edge at frequency $`\omega _u`$, the corresponding DOS exhibits a divergence of the form $`N(\omega )1/\sqrt{\omega \omega _u}`$. Conversely, for large frequencies ($`\omega \omega _u`$) the DOS will exhibit a UV-divergence, i.e., $`N(\omega )\omega ^2`$, as is the case in free space. More realistic LDOS’ coming from full photonic band structure computations do not suffer from the pathological band edge divergence of the isotropic model. However, by solving the model of a 3D isotropic photonic band gap, we make contact with previous results based on the isotropic model in the effective mass approximation .
Consider a 1D photonic dispersion relation in the extended zone scheme. In order to describe a PBG at wave number $`k_0`$ with central frequency $`\omega _c=ck_0=(\omega _u+\omega _l)/2`$ and upper and lower band edge at $`\omega _u`$ and $`\omega _l`$, respectively, we use the following Ansatz
$$\omega (k)=\{\begin{array}{c}\omega _++c_+\sqrt{(kk_0)^2+\gamma _+^2}\text{for}k>0\\ \omega _{}+c_{}\sqrt{(kk_0)^2+\gamma _{}^2}\text{for}k<0\end{array}.$$
(B1)
Using the requirements $`\omega (k=0)=0`$, $`\omega (k=k_00_+)=\omega _l`$, $`\omega (k=k_0+0_+)=\omega _u`$, $`_k\omega (k=0)=_k\omega (k\mathrm{})=c`$, and $`_k\omega (k=k_00_+)=_k\omega (k=k_0+0_+)=0`$, the unknown parameters in (B1) can easily be expressed in terms of a single parameter $`\eta =\omega _l/\omega _c`$, $`1/2<\eta 1`$ that describes the size of the photonic bandgap, giving: $`\omega _+=\omega _c`$, $`c_+=c`$, $`\gamma _+=k_0(1\eta )`$, $`\omega _{}=\omega _c(\eta ^2)/(2\eta 1)`$, $`c_{}=c\eta /\sqrt{2\eta 1}`$, and $`\gamma _{}=k_0(1\eta )/\sqrt{2\eta 1}`$.
From the dispersion relation (B1), the corresponding DOS, i. e., $`N(\omega )=d^3k\delta (\omega \omega (k))`$ is given by
$$N_m(\omega )=\{\begin{array}{c}4\pi c_{}^2\frac{[k_0\sqrt{(\omega \omega _{})^2/c_{}^2\gamma _{}^2}]^2(\omega _{}\omega )}{\sqrt{(\omega \omega _{})^2/c_{}^2\gamma _{}^2}}\\ \text{for}0\omega \omega _l\\ 4\pi c_+^2\frac{[k_0+\sqrt{(\omega \omega _+)^2/c_+^2\gamma _+^2}]^2(\omega \omega _+)}{\sqrt{(\omega \omega _+)^2/c+^2\gamma _+^2}}\\ \text{for}\omega _u\omega <\mathrm{}\end{array}$$
(B2)
For sufficiently large gaps $`(\eta 0.9)`$ and bare eigenfrequencies $`\omega _0`$ of the radiating dipole close to the upper band edge, it is an excellent approximation to ignore the lower branch of the photon dispersion relation, i. e., for $`kk_0`$. The resulting DOS for this so-called three-dimensional isotropic, one-sided bandgap model is shown in Fig. 1 for a value of gap width parameter $`\eta =0.8`$ and the gap center frequency $`\omega _ca/2\pi c=0.5`$. The square-root singularity at the band edge as well as the UV divergence $`N_m(\omega )\omega ^2`$ as $`\omega \mathrm{}`$ are clearly visible.
Figure 1: The DOS for the three-dimensional, isotropic one-sided bandgap model of a PC. The parameters (see appendix B) are $`\eta =0.8`$ and $`\omega _ca/2\pi c=0.5`$.
Figure 2: The radiation dynamics resulting from the three-dimensional, isotropic one-sided bandgap model DOS as shown in Fig. 1 for various values of the bare dipole oscillator frequency $`\omega _0`$ relative to the upper photonic bandedge $`\omega _u`$. The photonic bandedge is siutated at $`\omega _ua/2\pi c=0.6`$ and the bare dipole oscillator frequencies are (a) $`\omega _0a/2\pi c=0.58`$, (b) $`\omega _0a/2\pi c=0.595`$, (c) $`\omega _0a/2\pi c=0.599`$, (d) $`\omega _0a/2\pi c=0.6`$, (e) $`\omega _0a/2\pi c=0.601`$, (f) $`\omega _0a/2\pi c=0.605`$, and (g) $`\omega _0a/2\pi c=0.62`$. Clearly visible are normal–mode oscillations, or vacuum Rabi oscillations, and the fractional localization of radiation near the photonic bandedge. The coupling strength has been chosen such that $`g(\omega _0)=10^4`$. For frequencies deep in the photonic bandgap ($`\omega _0a/2\pi c=0.58)`$ and deep in the photonic conduction band ($`\omega _0a/2\pi c=0.62)`$ we observe negligible and exponential decay, respectively.
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# The transformations of the electromagnetic potentials under translations
## Abstract
I consider infinitesimal translations $`x^\alpha =x^\alpha +\delta x^\alpha `$ and demand that Noether’s approach gives a symmetric electromagnetic energy-momentum tensor as it is required for gravitational sources. This argument determines the transformations of the electromagnetic potentials under infinitesimal translations to be $`A_\gamma ^{}(x^{})=A_\gamma (x)+_\gamma [\delta x_\beta A^\beta (x)]`$, which differs from the usually assumed invariance $`A_\gamma ^{}(x^{})=A_\gamma (x)`$, by the gauge transformation $`_\gamma [\delta x_\beta A^\beta (x)]`$.
In relativistic field theory it is well known, and often referred to as Noether’s theorem , that each independent infinitesimal symmetry implies a conserved current with an associated constant or motion. Here we are interested in the symmetry under infinitesimal translations
$$x^\alpha =x^\alpha +\delta x^\alpha $$
(1)
for which Noether’s theorem yields the conserved currents of the energy-momentum tensor with the associated constants of motion being energy and momenta. However, the energy-momentum tensor $`T^{\alpha \beta }`$ obtained in this way from the electromagnetic Lagrangian
$$=\frac{1}{16\pi }F_{\alpha \beta }F^{\alpha \beta }\mathrm{with}F^{\alpha \beta }=^\alpha A^\beta ^\beta A^\alpha $$
(2)
is not symmetric, whereas everyone believes that the correct result ought to be symmetric. The standard procedure is to add a total divergence such that the final result becomes the desired symmetric tensor $`\theta ^{\alpha \beta }`$. While this procedure is acceptable within electrodynamics, it becomes questionable as soon as one is concerned about gravity. The electromagnetic energy-momentum distributions of $`T^{\alpha \beta }`$ and $`\theta ^{\alpha \beta }`$ differ and this changes the implied gravitational force. This is in principle observable , although in practice presumably not, unless someone identifies suitable cosmological field distributions. On the theory side, in gravity the symmetry transformations of general covariance yield the symmetric energy-momentum tensor as source of the gravitational field, see for instance , and this is presumably the strongest evidence underlining that the energy-momentum distribution of $`\theta ^{\alpha \beta }`$ is the correct one.
In view of the arguments in favor of the symmetric energy-momentum tensor it is astonishing that Noether’s theorem leads to a tensor $`T^{\alpha \beta }`$ with an obviously incorrect energy-momentum distribution. In particularly, one should bear in mind that the derivation of $`T^{\alpha \beta }`$ relies only on the symmetry transformations of a four vector under translations
$$A_\gamma ^{}(x^{})=A_\gamma (x)$$
(3)
and on factoring out the local translation $`\delta x^\alpha `$. Due to the locality of the procedure, it is hard to imagine that it could lead to an incorrect energy momentum distribution when the fundamental assumptions are sound.
In this letter I give a simple solution to the problem. The non-symmetric energy momentum tensor $`T^{\alpha \beta }`$ is obtained under the assumption that $`A_\gamma `$ transforms as a four vector (3) under translations (1). Due to the gauge invariance of the electromagnetic Lagrangian (2) this does not have to be the case. We may allow for more general transformations which differ from (3) by gauge transformations, i.e.
$$A_\gamma ^{}(x^{})=A_\gamma (x)+_\gamma \mathrm{\Lambda }(x).$$
(4)
That nature may employ such a transformation behavior instead of (3) is not entirely a surprise. On the quantum level the electromagnetic fields rely on superpositions of massless creation and annihilation operators and Weinberg points out to us that such fields do not allow for representations of the (proper) Lorentz group, but only for transformations which differ from those by a gauge transformation. Therefore, it is quite natural to conjecture that nature uses gauge transformation also for translations. Repeating the arguments of Noether’s theorem with the ansatz (4) and requesting a symmetric energy-momentum tensor leads to the unique solution
$$A_\gamma ^{}(x^{})=A_\gamma (x)+_\gamma [\delta x_\beta A^\beta (x)]$$
(5)
which is conjectured to be the transformation law realized by nature for infinitesimal translations of electromagnetic potentials. The remainder of the paper is devoted to the derivation of this equation. Up to some notational changes and adaptions to the case at hand, my arguments follow closely chapter 1 of Bogoliubov and Shirkov .
First, let us quickly recall how relativistic field equations are derived from the action principle. The action is a four dimensional integral over a scalar Lagrangian density
$$𝒜=d^4x(\psi _k,_\alpha \psi _k)$$
(6)
and, therefore, by itself a scalar under the connected part of the Lorentz group. Variations of the fields are defined as functions
$$\delta \psi _k(x)=\psi _k^{}(x)\psi _k(x)$$
(7)
which are non-zero for some localized space-time region. The action is required to vanish under such variations
$$0=\delta 𝒜=$$
$$\underset{k}{}d^4x\left[(\delta \psi _k)\frac{}{\psi _k}+(\delta _\alpha \psi _k)\frac{}{(_\alpha \psi _k)}\right].$$
(8)
Integration by parts gives
$$0=\underset{k}{}d^4x(\delta \psi _k)\left[\frac{}{\psi _k}_\alpha \frac{}{(_\alpha \psi _k)}\right],$$
(9)
where we used that the surface terms vanish. As the variations $`\delta \psi _k`$ are independent, the integrand in (9) has to vanish for each $`k`$ and we arrive at the Euler-Lagrange equations
$$\frac{}{\psi _k}_\alpha \frac{}{(_\alpha \psi _k)}=0$$
(10)
for relativistic fields. For the electrodynamic Lagrangian (2) they yield $`_\alpha F^{\alpha \beta }=0`$.
Noether’s theorem applies to transformations of the coordinates for which the transformations of the field functions are also known. Such transformations constitute a symmetry of the theory when the corresponding variation of the action vanishes. The theorem states that to each such symmetry a combination of the field functions exists which defines a conserved current. For this purpose we introduce, in addition to (7), a second type of variations which combines space-time and their corresponding field variations
$$\overline{\delta }\psi _k(x)=\psi _k^{}(x^{})\psi _k(x).$$
(11)
Using (note $`\delta x^\alpha _\alpha \psi _k^{}=\delta x^\alpha _\alpha \psi _k`$ because $`\delta ^2`$ variations disappear)
$$\psi _k^{}(x^{})=\psi _k^{}(x)+\delta x^\alpha _\alpha \psi _k(x)$$
we find a relation between the variations (11) and (7)
$$\overline{\delta }\psi _k(x)=\delta \psi _k(x)+\delta x^\alpha _\alpha \psi _k(x).$$
(12)
For a scalar field $`\psi `$ (as well as for ordinary four vector fields) symmetry under translations means
$$\overline{\delta }\psi (x)=\psi ^{}(x^{})\psi (x)=0.$$
(13)
But for the electromagnetic potentials we allow (4)
$$\overline{\delta }A_\gamma (x)=A_\gamma ^{}(x^{})A_\gamma (x)=_\gamma \mathrm{\Lambda }(x).$$
(14)
With these symmetries equation (12) reduces for a scalar field to
$$\delta \psi =\delta x^\alpha _\alpha \psi (x)$$
(15)
and for the electromagnetic potentials to
$$\delta A_\gamma =_\gamma \mathrm{\Lambda }(x)\delta x^\alpha _\alpha A_\gamma (x).$$
(16)
As the Lagrange density is a scalar, we get for its combined variation (11)
$$0=\overline{\delta }=^{}(x^{})(x)=\delta +\delta x^\alpha _\alpha $$
(17)
where besides (13) we used the relation (12). Our aim is to factor an over-all variation $`\delta x^\alpha `$ out. For $`\delta `$ we proceed as in equation (8), where the fields $`\psi _k`$ are now replaced by the gauge potentials $`A_\gamma `$
$$\delta =(\delta A_\gamma )\frac{}{A_\gamma }+(\delta _\alpha A_\gamma )\frac{}{(_\alpha A_\gamma )}.$$
Using the Euler-Lagrange equation (10) to eliminate $`/A_\gamma `$, we get (the calculation remains valid in our case where $``$ does not depend on $`A_\gamma `$)
$$\delta =(\delta A_\gamma )_\alpha \frac{}{(_\alpha A_\gamma )}+(\delta _\alpha A_\gamma )\frac{}{(_\alpha A_\gamma )}$$
$$=_\alpha \left[(\delta A_\gamma )\frac{}{(_\alpha A_\gamma )}\right].$$
Let us collect all terms which contribute to $`\overline{\delta }`$ in equation (17). For this, note that $`_\beta \delta x^\alpha =0`$ holds for all combinations of indices $`\alpha `$, $`\beta `$. (Namely, for $`\alpha =\beta `$ we are led to $`\delta 1=0`$ and for $`\beta \alpha `$ the variations $`\delta x^\alpha `$ are then independent of the coordinates $`x^\beta `$). We find
$$0=\overline{\delta }=_\alpha \left[(\delta A_\gamma )\frac{}{(_\alpha A_{\gamma )}}+\delta x^\alpha \right]=$$
$$_\alpha \left[\left(_\gamma \mathrm{\Lambda }(x)(\delta x_\beta )^\beta A_\gamma \right)\frac{}{(_\alpha A_\gamma )}+g^{\alpha \beta }\delta x_\beta \right]$$
where equation (16) was used for the last step. To be able to factor $`\delta x_\beta `$ out of the bracket, one has to request
$$\mathrm{\Lambda }(x)=\delta x_\beta B^\beta (x)$$
(18)
where $`B^\beta (x)`$ is a not yet determined potential field. With this we get
$$0=\delta x_\beta _\alpha \left[(^\beta A_\gamma _\gamma B^\beta )\frac{}{(_\alpha A_\gamma )}g^{\alpha \beta }\right].$$
As the variations $`\delta x_\beta `$ are independent, the energy-momentum tensor
$$\theta ^{\alpha \beta }=\frac{}{(_\alpha A_\gamma )}(^\beta A_\gamma _\gamma B^\beta )g^{\alpha \beta }$$
(19)
gives the conserved currents
$$_\alpha \theta ^{\alpha \beta }=0.$$
(20)
Let us demand that the energy-momentum tensor (19) is symmetric. This leads to the requirement
$$B^\beta (x)=A^\beta (x)$$
(21)
for which
$$\theta ^{\alpha \beta }=\frac{1}{4\pi }\left(F^{\alpha \gamma }F_\gamma ^\beta +\frac{1}{4}g^{\alpha \beta }F_{\gamma \delta }F^{\gamma \delta }\right)$$
(22)
is symmetric because of
$$F^{\alpha \gamma }F_\gamma ^\beta =F^{\beta \gamma }F_\gamma ^\alpha .$$
Indeed, equation (22) is the symmetric tensor of the textbooks . It differs from other versions of (19) by total divergencies.
In conclusion, I have presented an argument in favor of the transformation behavior (5) and it appears that the question of the energy-momentum distribution of the electromagnetic field may finally be put at rest with the expected result. Noether’s theorem alone has no predictive power about whether the energy-momentum tensor is the symmetric or not, but evrything is consistent in the sense that a transformation law exists which gives the tensor (22). In addition, when symmetry of the tensor is assumed the result for $`\theta ^{\alpha \beta }`$ is unique. It is shown in a forthcoming paper that this approach works also for non-abelian gauge theories.
Note added
After posting this manuscript Prof. Jackiw kindly informed me that my result is a special case of his work , see for details. Prof. Hehl communicated that the use of 1-Forms leads directly to a symmetric energy-momentum tensor, see for instance .
###### Acknowledgements.
This work was in part supported by the US Department of Energy under contract DE-FG02-97ER41022.
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# Cooling sequences and color-magnitude diagrams for cool white dwarfs with hydrogen-atmospheres
## 1 Introduction
The possible interpretation of the faint blue objects in the Hubble Deep Field as halo white dwarfs (Hansen 1998; Ibata et al. 1999; Méndez & Minniti 2000), consistent with the interpretation of the observed microlensing events towards the Large Magellanic Clouds as stellar remnants (Chabrier, Segretain, & Méra 1996; Adams & Laughlin 1996; Graff, Laughlin & Freese 1998; Chabrier 1999), and the spectroscopic identification of very cool ($`T_{\mathrm{eff}}<\mathrm{\hspace{0.17em}4000}`$ K), high proper motion white dwarfs (WD) (Hodgkin et al. 2000; Ibata et al. 2000) has triggered interest in the study of old, cool WDs and stressed the need for accurate cooling sequences and predicted observational signatures for these objects. This implies a correct cooling theory for crystallized white dwarfs and reliable atmosphere models and photometric predictions.
The basic physics entering WD evolution has evolved significantly since the pioneering work of Mestel & Ruderman (1967) and the first detailed evolutionary calculations by Lamb & Van Horn (1975). Noteworthy advances have been made on the fronts of the conductive opacity (Itoh et al. 1983; Itoh, Hatashi, & Kohyama 1993; Potekhin et al. 1999), the radiative opacity (Lenzuni, Chernoff, & Salpeter 1991; Rogers & Iglesias 1992), the envelope equation of state (EOS) (Saumon, Chabrier, & Van Horn 1995, SCVH), and the detailed description of the thermodynamic properties of the dense, fully ionized interior plasma, including the main effects of ion crystallization, namely the latent heat and the chemical fractionation (Segretain et al. 1994 and references therein). In the meantime, substantial improvement in the theory of the atmosphere of cool WDs has been accomplished (Bergeron, Wesemael, & Fontaine 1991; Bergeron, Saumon, & Wesemael 1995, hereafter BSW; Bergeron, Wesemael, & Beauchamp 1995, hereafter BWB).
These studies were devoted primarily to the characterization of the disk WD population and were thus restricted to effective temperatures $`T_{\mathrm{eff}}>4000`$ K. As identified initially in BSW, the onset of molecular recombination below $`5000`$ K and the pending Collision-Induced Absorption (CIA) due to H<sub>2</sub>-H<sub>2</sub> and H<sub>2</sub>-He collisions in such dense atmospheres ($`g10^8`$ to $`10^9`$ cm s<sup>-2</sup>) results in important departures from a blackbody energy distribution, with an increasing absorption of the flux longward of 1 $`\mu `$m. The WD atmosphere calculations were extended recently to lower temperatures by Hansen (1998, 1999) and by Saumon & Jacobson (1999, hereafter SJ), reaching the effective temperature range characteristic of the halo or globular cluster old WD population. As shown by these authors, the strong CIA in the infrared redistributes the flux toward shorter wavelengths so that the emergent flux peaks in optical passbands, regardless of $`T_{\mathrm{eff}}`$. This effect has been confirmed recently with the spectroscopic observation of H<sub>2</sub> CIA in LHS 1126 (Bergeron, Ruiz & Leggett, 1997), LHS 3250 (Harris et al. 1999) and WD0346+246 (Hodgkin et al. 2000). As demonstrated initially by Hansen (1998), the colors of very cool WDs are consistent with the unidentified faint blue objects in the Hubble Deep Field (HDF) (Bahcall et al. 1994; Méndez et al. 1996; Nelson et al. 1996; Ibata et al. 1999). A preliminary set of cooling sequences incorporating the Saumon & Jacobson (1999) synthetic colors for cool WDs with pure H atmospheres provides luminosity and discovery functions in various passbands for a dark halo WD population consistent with the microlensing experiments (Chabrier 1999).
Current cooling models (Wood 1995; Segretain et al. 1994; Salaris et al. 1997; Montgomery et al. 1999) appear to be reasonably consistent with each other for WDs with $`T_{\mathrm{eff}}>\mathrm{\hspace{0.17em}4000}`$ K, i.e. younger than $`10`$ Gyr for a 0.6 $`M_{}`$ WD. For these temperatures, hydrogen molecular recombination is absent or negligible and the spectral energy distribution is moderatly affected by H<sub>2</sub>-H<sub>2</sub> or H<sub>2</sub>-He CIA (see, e.g., BSW). The aim of the present paper is to extend cooling calculations to older WDs that are characteristic of the disk, spheroid or dark halo population. To this end, we use upgraded interior physics, synthetic spectra, and interior-envelope relations. A description of the various physical inputs entering the calculations, and a comparison with previous calculations are given in §2. Cooling sequences, color-magnitude and color-color diagrams in various optical and infrared passbands are presented in §3, and the remaining uncertainties in the models are examined. Conclusions are presented in Section 4.
## 2 Model calculations
In the present study, we will concentrate on the so-called DA WDs, i.e. those having either pure hydrogen atmospheres or atmospheres with a small admixture of helium ($`[N(\mathrm{He})]/[N(\mathrm{H})]1`$). The main reason is the fact that detectable halo WDs will have very likely hydrogen-rich atmospheres. It is indeed well known that white dwarfs with pure helium outer layers must evolve more rapidly than their DA counterparts at the faint end of the cooling sequence because of the extreme transparency of these layers. Explicit calculations, assuming such pure helium envelopes, show that helium-atmosphere halo WDs, if they exist, would escape detection since they reach $`M_{bol}>\mathrm{\hspace{0.17em}19}`$ after $`8`$ Gyr (Chabrier 1997; Hansen 1998). We note, however, that even very small traces of hydrogen or metals in their envelopes –possibly due to accretion, microscopic diffusion, or convective dredge-up– could change this picture dramatically. Indeed, such traces increase significantly the opacity of the otherwise transparent neutral helium at low effective temperatures, slowing down the cooling of the star. If hydrogen itself is present, even as a trace, the emergent spectrum is brought closer to that of a DA star.
Note also that reliable calculations of pure helium model atmospheres at very low effective temperatures remain to be done. For instance accurate treatments of helium pressure-ionization, as raised initially by Böhm et al. (1977; see also BSW and Hansen, 1998) and of $`He^{}`$ free-free absorption cross-section at high densities are still lacking. While realistic model atmospheres of non-DA WDs are available for $`T_{\mathrm{eff}}>4000`$ K (BSW, BWB) and reasonable estimates of their cooling timescales have been published (Wood 1995; Segretain et al. 1994; Salaris et al. 1997), the same cannot be said at the cooler end of the sequence. Fortunately, as mentioned above, the most probable cases of detection of halo WDs are identified with DA stars.
### 2.1 Core Physics
#### 2.1.1 Equation of state
In the present paper, we consider only WDs with carbon-oxygen cores, which restricts the mass range to $`0.5M_{}M1.2M_{}`$. We define the core as the region where the C/O plasma is fully ionized and the electron gas is fully degenerate. The large electron conductivity ensures that the core is very nearly isothermal. This essentially isothermal domain encompasses typically 99.99% of the mass of cool WDs (see, e.g., Figure 14 of Tassoul, Fontaine, & Winget 1990). In this region, we adopt the equation of state (EOS) described in §2 of Segretain et al. (1994) for the liquid and the solid phases:
$$U_L(x_i,\rho ,T)=U_i^{id}+U_i^{ex}+U_e^{id}+U_e^x+U_{ie}$$
(1)
$$U_S(x_i,\rho ,T)=U_i^{th}+U_i^{anh}+U_i^{qm}+U_i^{Mad}+U_e^{id}+U_e^x+U_{ie}$$
(2)
where the different contributions are described in Segretain et al. (1994) ($`x_i=N_i/N`$ denotes the number-fraction of each element, $`C`$ and $`O`$).
In the present calculations, we introduce a major improvement with respect to Segretain et al. (1994) which is essential for cool, crystallized WDs. In Segretain et al. (1994) and in further calculations based on this EOS, the ion-electron screening contribution, $`U_{ie}`$, was taken from Yakovlev & Shalybkov (1989, hereafter YS). As stressed by these authors, their calculations are valid only in the fluid phase, where the plasma ion coupling parameter $`\mathrm{\Gamma }=2.275\times 10^5Z^{5/3}\frac{(\rho Y_e)^{1/3}}{T}<\mathrm{\hspace{0.17em}200}`$, where $`Y_e=Z/A`$ is the average electron molar fraction and $`Z^{5/3}=_ix_iZ_i^{5/3}`$. The arrow in Figure 1 indicates the luminosity at which the ion coupling parameter at the center of the star reaches the afore-mentioned value 200 for a 0.6 $`M_{}`$ WD. As shown in the figure, extrapolation of the YS fit in the solid phase ($`\mathrm{\Gamma }>200`$) yields an increasingly inaccurate cooling sequence. Note also that the YS fit is derived from calculations of the thermodynamic properties of a classical ionic plasma, whereas the C/O solid core is in a quantum state, with an ion diffraction parameter $`\eta =\mathrm{}\mathrm{\Omega }_P/kT>1`$, where $`\mathrm{\Omega }_P`$ denotes the ion plasma frequency (see Segretain et al. 1994). The Segretain et al. (1994) study was primarily devoted to disk WDs, i.e., objects brighter than $`\mathrm{log}L/L_{}4.5`$. This corresponds approximately to the end of the crystallization process in WD interiors. Halo WDs, however, are fainter than $`10^{4.5}L_{}`$ and $`\mathrm{\Gamma }>200`$ throughout a large fraction of the star.
In fully crystallized WDs, there is no further contribution to the luminosity from crystallization-induced chemical fractionation of C and O (see Segretain et al. 1994; Chabrier 1997; Isern et al. 1997). The only contributions to the luminosity come from the thermal reservoir of the star and from the residual gravitational contraction of the outermost layers (see e.g. Koester & Chanmugam, 1990; D’Antona & Mazzitelli, 1990):
$$dL/dt=_0^MC_V\frac{dT}{dt}𝑑m_0^M(T\frac{dP}{dT})_V\frac{dv}{dt}𝑑m.$$
(3)
Note that the thermal energy (first term on the r.h.s. of equation (3)) is comparable to the change in gravitational energy $`\mathrm{\Delta }\mathrm{\Omega }`$, from the virial theorem. For a substantially (entirely) crystallized WD, most of (all) the thermal energy stems from the specific heat of the quantum solid. Although for very cool WDs the thermal contribution of the central regions, where $`\mathrm{}\mathrm{\Omega }_P/kT>>1`$, becomes increasingly small, the contribution of the outer layers, where $`\mathrm{}\mathrm{\Omega }_P/kT1`$, remains substantial. Indeed, in cool WDs, the main contributions to the internal energy, namely the zero-temperature electron gas kinetic energy $`U_e^{id}`$ and the ionic electrostatic energy $`U_i^{Mad}`$ in equation (2) (see Figure 1 of Segretain et al., 1994), do not depend on temperature, so that the heat capacity, and thus the cooling, is entirely determined by the small temperature-dependent corrections to the energy. In fact the ion-electron screening contribution, which stems from the polarization of the electrons by the ionic field, becomes the dominant contribution to the specific heat at low temperature since it decreases as $`C_{V_{ie}}\eta ^1T`$, whereas the ionic crystal (Debye) contribution decreases as $`C_{V_{ii}}\eta ^3T^3`$ (Potekhin & Chabrier 2000, hereafter PC).
For fully crystallized WDs, the central density is of the order of $`\rho _c10^6\mathrm{g}\mathrm{cm}^3`$, so that the Fermi parameter $`x=p_F/m_ec=1.01\times 10^2(\rho Y_e)^{1/3}>1`$. Therefore we must consider the energetic contribution due to the ion-electron interaction of a polarizable, relativistic electron gas immersed in a quantum Coulomb crystal. The Thomas-Fermi approximation (Salpeter 1961) is valid only in the asymptotic limit of an infinite ionic charge $`Z\mathrm{}`$ and is not valid for a C<sup>6+</sup>/O<sup>8+</sup> plasma (see, e.g., YS). The ion-ion and ion-electron contributions to the EOS of a quantum electron-ion solid plasma for a finite ionic charge under the conditions of interest have been calculated recently by Potekhin & Chabrier (2000). To the best of our knowledge these are the only available calculations for such plasma conditions. Figure 1 compares the evolution of a 0.6 $`M_{}`$ crystallized WD for the same $`T_{\mathrm{eff}}`$-$`T_c`$ condition (as described below) with the solid EOS and specific heat incorporating (i) the proper ion-electron screening treatment, and (ii) an extrapolation of the YS fit. As shown, the extrapolation yields increasingly shorter cooling times for WDs older than $`8`$ Gyr.
The contribution of crystallization along evolution is treated as described in §4 of Segretain et al. (1994), from the evolution of the binding energy, $`dB(T)/dt`$. Recent complete evolutionary calculations by Montgomery et al. (1999) confirm the validity of this method for cool WDs, where there is no neutrino or nuclear luminosity. The crystallization of <sup>22</sup>Ne, which yields an azeotropic diagram (Segretain & Chabrier 1993), and was originally thought to produce a significant time-delay (Isern et al 1991; Segretain et al. 1994) has not been considered in the present calculations. Indeed, consistent calculations of the three-body Ne/C/O phase diagram (Segretain 1996) show that when Ne-crystallization sets in, a substantial fraction of oxygen has already crystallized, yielding an O-rich core, so that the remaining amount of energy due to <sup>22</sup>Ne crystallization becomes fairly small. At the crystallization luminosity found in Segretain et al. (1994), the induced time delay is at most $`0.3`$ Gyr, well within the remaining uncertainties in our calculations (phase diagram, opacities, etc). In any event, since we are presently mainly interested in old WDs originating from low-metallicity progenitors, the Ne abundance is likely to be too small to have any measurable effect.
#### 2.1.2 Initial carbon-oxygen profiles
Substantial uncertainty remains on the $`{}_{}{}^{12}C(\alpha ,\gamma )^{16}O`$ reaction which determines the final state of the He-burning AGB phase and thus the initial C/O abundance-profile of the WD. Modern experimental data and updated values of the astrophysical $`S`$-factor (see, for example, Arnould et al. 1999) suggest a rate about a factor 1.5 to 2 larger than that of Caughlan & Fowler (1988,CF88), closer to the previous Caughlan et al. (1985) value. This yields a substantially O-enriched initial profile in the WD core (see e.g. Salaris et al. 1997). We have adopted these initial profiles in our calculations. However, in order to examine the uncertainties due to the interior composition, we have also conducted calculations for a 0.6 $`M_{}`$ WD with an initial C/O distribution resulting from the CF88 lower rate (see Figure 2 of Salaris et al. 1997).
### 2.2 The Luminosity-Central Temperature Relation
The binding energy method that we use to compute cooling ages relies on the availability of $`L`$-$`T_c`$ relations which govern the cooling rates of WDs. These relations are particularly sensitive to the constitutive physics of the outer partially ionized, partially degenerate envelope which connects the nearly isothermal core of a WD to its surface (the atmospheric layers). Previous calculations based on this method (Segretain et al. 1994 and references therein; Salaris et al. 1997) used $`L`$-$`T_c`$ relations provided by independent model calculations such as, for example, those of Wood & Winget (1989) or Wood (1995). This approach necessarily introduces some inconsistencies in the calculations: the chemical composition of the core of the models used to derive the $`L`$-$`T_c`$ relation differs from the variable core composition of the nearly isothermal structure that undergoes phase separation, and the constitutive physics is generally different in the two sets of models. Moreover, for masses not directly available from independent models, $`L`$-$`T_c`$ relations were scaled on the mass, which, at best, provides a rough estimate of the correct relations.
We have improved on this front in the present paper by computing $`L`$-$`T_c`$ relations based on state-of-the-art constitutive physics and by considering several individual masses. We are still left with the inherent inconsistency of the binding energy approach due to the fixed interior composition in the calculation of the $`L`$-$`T_c`$ relation, but this shortcoming is largely offset by our ability to describe in accurate details the physics of crystallization (e.g., the phase diagram) or the radiative and conductive opacities at high density, or by the afore-mentioned uncertainties in the stratified C/O profile. Note also that the treatment of crystallization-induced fragmentation can be easily implemented in the binding energy method, whereas it is a complicated task to include it into a standard evolutionary code. At this level, the binding energy method shows an appreciable advantage. Moreover, as mentioned above, the validity of this method for cool WDs has been assessed recently by comparisons with complete evolutionary calculations (Montgomery et al. 1999).
The $`L`$-$`T_c`$ relations for various masses in the range of interest were computed with an upgraded version of the stellar model building code briefly described in Brassard & Fontaine (1994, 1997). These are full stellar models that describe the complete structure of a static WD from the center to the high atmosphere ($`\tau _R10^6`$, where $`\tau _R`$ denotes the Rosseland optical depth). The envelope calculation incorporates the SCVH EOS for H and He and the Fontaine, Graboske, & Van Horn (1977) EOS for carbon and oxygen. The radiative opacities include the OPAL 1996 data (Iglesias & Rogers 1996) complemented at low temperatures by the Rosseland opacities of H and He computed with the model atmosphere code of BSW down to 1500 K for H and 2500 K for He. These latter opacities include the complete CIA processes (see BSW). For the conductive opacities, we use a large table incorporating the Hubbard & Lampe (1969) and Itoh et al. (1983, 1993) calculations (Brassard & Fontaine 1994), which covers the entire density and temperature range relevant to the present calculations. Convection is described with the standard mixing-length theory. As shown in Tassoul et al. (1990), the $`L`$-T<sub>c</sub> relationship is insensitive to the assumed convective treatment, since convection is essentially adiabatic when it breaks through the degenerate core.
We note that Hansen (1998, 1999) has recently emphasized the importance of treating in detail the atmospheric layers in the context of the evolution of very cool WDs. As discussed initially by Fontaine & Van Horn (1976; see also Tassoul et al. 1990), the most important consequence of such a detailed description of the atmosphere on the cooling is its effect on the location of the base of the convection zone in the deeper, optically-thick envelope. Indeed, below a luminosity $`\mathrm{log}(L/L_{})3.7`$ ($`T_{\mathrm{eff}}55006500`$ K), the base of the hydrogen superficial convection zone reaches into the degenerate core, thus coupling, for the first time during the evolution, the atmospheric layers with the central thermal reservoir. Since, by then, the stratification of the envelope is fully convective and highly adiabatic, small changes at the base of the atmosphere produce corresponding changes at the base of the convection zone. A correct determination of the base of that zone is thus essential to calculate an accurate $`L`$-$`T_c`$ relation for cool WDs.
It is important to stress that, although the atmospheric structures entering the present stellar models are gray, they take into account the feedback effect of convection on the atmospheric structure and thus are not based on a Rosseland mean all the way through, unlike those used in all previous evolutionary calculations, with the exception of Hansen’s (1998; 1999) computations. These modified model atmospheres, to be described elsewhere (Brassard & Fontaine, in preparation), reproduce almost exactly the stratification of a detailed model atmosphere (from, e.g., BSW in the present context) at large optical depths. They reproduce well, in particular, the main effect due to nongrayness, namely the upward shift of the convection zone. This phenomenon is well known in stellar atmosphere theory and has been described, in a white dwarf context, by Böhm et al (1977) among others. Therefore, the boundary conditions provided by the detailed model atmosphere at $`\tau _R=100`$ are essentially the same as those provided by the gray model with convective feedback. This is what matters for the cooling; for nearly identical boundary conditions at the base of the convective atmosphere, we find the same location of the base of the full envelope convection zone, and hence the same value of $`T_c`$. In contrast, standard gray atmosphere stratification (not taking into account the feedback due to convection) leads to an incorrect determination of the base of the convective zone, overestimating its penetration and thus yielding a faster cooling (see §3.1 below).
A central parameter in the envelope calculation, and in the resulting cooling time, is the amount of hydrogen and helium present in the envelope. For the present calculations, we have used DA white dwarf standard “thick” layers, $`\mathrm{log}q(H)=4.0`$, $`\mathrm{log}q(He)=2.0`$, where $`q(X)=M(X)/M_{}`$ denotes the mass fraction of element $`X`$. In practice, the masses of the hydrogen and helium layers on top of the degenerate C/O interior depend on the WD mass, through the AGB and post-AGB evolutionary phase (see e.g. Blöcker et al., 1997). Note, however, that (i) calculations during the AGB phase strongly depend on ill-constrained parameters (overshooting, mass loss,…) and (ii) present calculations consider the evolution of solar-metallicity AGB stars, whereas the progenitors of very cool WDs are metal-depleted. Given these uncertainties in the exact amount of $`q(H)`$ and $`q(He)`$, we elected to conduct our calculations with the afore-mentioned ”standard” values for 0.6 $`M_{}`$ WD (see Fontaine & Wesemael 1997 for a general discussion of this unsettled issue). This uncertainty in the exact amount of H and He in WD envelopes certainly represents one of the major uncertainties in present WD cooling calculations. Some scaling relations, however, can be used. Indeed, the effect of the thickness of the helium layer on the cooling time has been examined in detail by Tassoul et al. (1990), Wood (1992) and Montgomery et al. (1999, §5.1). Thicker helium layers result in more transparent envelopes (since $`\kappa _{\mathrm{He}}\kappa _{\mathrm{Carbon}}`$ under WD conditions), decreasing the temperature gradient between the core and the surface, so that the central temperature decreases faster with decreasing luminosity. Therefore, models with thicker He layers are younger for a given mass and luminosity, with a $`0.75`$ Gyr decrease in the age for each order of magnitude in $`M_{He}`$ (Montgomery et al. 1999).
### 2.3 Model atmospheres and photometric colors
The observable properties of the cooling WD, such as the emergent spectrum and the photometric colors are obtained with atmosphere models. The present calculations include the colors and bolometric corrections for pure hydrogen atmospheres calculated by Bergeron et al. (1995a) above 4000 K, extended down to $`T_{\mathrm{eff}}=1500`$ K by Saumon & Jacobson (1999). As mentioned previously, for very cool WDs ($`T_{\mathrm{eff}}<5000`$ K), molecular hydrogen becomes stable and the main source of opacity in the infrared is the CIA by H<sub>2</sub>. This opacity forces the stellar flux to emerge at shorter wavelengths, with a peak near 1$`\mu `$m. This increased flux in the $`R`$ and $`I`$ bands and decreased flux in infrared results in increasingly blue color indices for cooler WDs (Hansen 1998, 1999; SJ).
Although similar to those of Hansen (1998, 1999), the present calculations include a more detailed treatment of the microphysics entering the atmosphere (BSW and SJ). Indeed, an important feature in these cool and dense atmospheres is the effect of the surrounding particles on the partition function of an atom, which eventually leads to the pressure-ionization of hydrogen. These modified internal partition functions imply a different (non-ideal) ionization equilibrium, in particular for the abundances of H<sub>2</sub>, H$`{}_{}{}^{+}{}_{2}{}^{}`$ and H$`{}_{}{}^{+}{}_{3}{}^{}`$. This in turn modifies the abundance of free electrons, and thus of H<sup>-</sup> ions, a dominant source of opacity in the optical. The atmosphere models of SJ were used in the range $`1500T_{\mathrm{eff}}4000`$ K and $`7.5\mathrm{log}g9.0`$, complemented by BSW for $`4000T_{\mathrm{eff}}10000`$ K, with a pure hydrogen composition. These synthetic spectra and atmosphere models successfully reproduce spectroscopic and photometric observations of cool H-rich (DAs) WDs above 4000 K (Bergeron et al. 1997).
Optical colors for cool WDs with a small but non-zero $`N(\mathrm{He})/N(\mathrm{H})`$ composition ratio will be only slightly different ($`<\mathrm{\hspace{0.17em}0.2}`$ mag) from the colors of pure hydrogen atmospheres presented here (see BSW; Bergeron et al. 1997). The color indices of Table 2 of BSW can then be used above 4000 K for mixed composition atmospheres.
## 3 Results
### 3.1 Cooling curves
For the purposes of comparisons, we define a set of reference models which include the physics described in the previous section, including crystallization-induced fragmentation (see below), with the Salaris et al. (1997) high-rate stratified initial C/O profiles. The $`L`$-$`T_c`$ relations that we used for these reference models have been computed as described above, from full static stellar models with chemical layering defined by $`\mathrm{log}q(\mathrm{H})=4.0`$ and $`\mathrm{log}q(\mathrm{He})=2.0`$. The photometric colors are taken from pure-H atmosphere calculations.
#### 3.1.1 Comparison with existing calculations
We first compare, in Figure 2, our $`L`$-$`T_c`$ relation for a 0.6 $`M_{}`$ WD with the one obtained by Wood (1995), Montgomery et al. (1999, Table 1) and another one kindly provided by Hansen (1999) for the same values of $`q(\mathrm{H})`$ and $`q(\mathrm{He})`$. The four curves agree reasonably well, but the relatively small differences shown here translate into significant differences on the cooling time (see below). Unfortunately, it is currently impossible to account in detail for the differences in the $`L`$-$`T_c`$ relations illustrated in Figure 2; differences in the constitutive physics used by the various groups are probably at the origin of most of the deviations. The $`L`$-$`T_c`$ relation is particularly sensitive to the opacity profile throughout the star, so slightly different implementations of the conductive and radiative opacities, as well as the use of different generations of these data, could very well account for most of the discrepancies. As expected, the Wood (1995) and the Montgomery et al. (1999) relations are very similar, since they rely on the same input physics, yielding similar cooling sequences.
As mentioned in §2.2., the change of slope observed in all three curves is a well known phenomenon and is due to convection breaking into the degenerate thermal reservoir at sufficiently low luminosities. When that occurs, there is a flattening of the temperature gradient between the core and the surface due to the larger efficiency of convection as compared to radiation. Initially, this leads to a slowdown of the cooling process as an excess energy is liberated through the more transparent convective envelope, but this is rapidly followed by a phase whereby convection speeds up the cooling process as compared to the case of purely radiative models (see Figure 3 of Tassoul et al. 1990 and related discussion). For luminosities lower than that of the breaks in slope shown in Figure 2, the $`L`$-$`T_c`$ relation becomes sensitive, among other things, to the details of the atmosphere. From this point of view, the Hansen (1999) calculations and our own improve upon those of Wood (1995) or Montgomery et al. (1999) because the latter ones rely on a standard gray atmosphere strategy which overestimates the penetration of the superficial convection zone and leads to a value of $`T_c`$ slightly lower than it should be, as discussed previously.
On the other hand, for luminosities larger than those of the changes of slope, it is well established that the relation is completely insensitive to the stratification of the upper envelope and, in particular, to the details of the atmosphere (see Tassoul et al. 1990 for a complete discussion). The discrepancies shown in Figure 2 between the three curves for these higher luminosities are then the explicit proof that there are indeed significant differences in the implementation/calculation of the constitutive physics between the three groups, notably at the level of the conductive opacities. So, to a certain extent, the differences found in Figure 2 result from the fact that we are comparing models with different physical inputs.
Figure 3 illustrates the consequences of these differences on the cooling time of the star. For each $`L`$-$`T_c`$ relation, we show two cases: (i) taking into account the chemical fractionation of C and O at crystallization (rightmost curves), (ii) ignoring this process (leftmost curves). Interestingly enough, we find that in the phases of interest for the present study $`[\mathrm{log}(L/L_{})<4.5]`$), our results (solid curves) agree reasonably well with those obtained when using the $`L`$-$`T_c`$ relation provided by Hansen (long-dashed curves). The results derived on the basis of the Wood (1995) (or Montgomery et al. 1999) $`L`$-$`T_c`$ relation (dotted curves) provide, in comparison, shorter cooling timescales, as anticipated from the discussion above. The differences illustrated in Figure 3 simply amplify the differences already observed in the $`L`$-$`T_c`$ relations of Figure 2.
Another interesting comparison is provided by the short-dashed curve in Figure 3 which presents the cooling curve obtained by Hansen himself (1999), which is supposed to include the complete treatment of crystallization and thus must be compared with the right long-dash curve. The difference between the two curves is noticeable ($`>\mathrm{\hspace{0.17em}0.5}`$ Gyr) and can not be totally ignored. Since such a comparison eliminates the effect of the energy transfer problem (the $`L`$-$`T_c`$ relations are the same), the remaining discrepancies must be blamed mostly on the different treatments of the thermodynamics of the ionized interior, outlined in §2.1.1. Since Hansen has not detailed his interior EOS, it is not possible to pursue this point further. Note, however, that the treatment of crystallization in Hansen’s calculations has been demonstrated to be incorrect (Isern et al. 2000). We thus disagree with him about the importance of chemical fractionation at crystallization, as discussed in the next subsection. Note also that apparently Hansen’s calculations do not extend beyond $`\mathrm{log}L/L_{}5.0`$, which corresponds to an age of $`t12`$ Gyr for a 0.6$`M_{}`$ WD.
#### 3.1.2 Effects of internal composition and crystallization
Figure 3 illustrates also the effect of initial C/O stratification on the cooling time for a 0.6 $`M_{}`$ H-atmosphere WD, as examined in detail by Salaris et al. (1997). The rightmost dot-dash curve displays the cooling sequence of our reference model with an initial C/O profile resulting from a low (CF88) $`{}_{}{}^{12}C(\alpha ,\gamma )^{16}O`$ reaction rate, to be compared with the sequence obtained with a high rate (Caughlan et al. 1985) induced profile (rightmost solid line). A larger rate yields a larger initial oxygen-enriched core and affects the cooling in several ways (Segretain et al. 1994; Salaris et al. 1997): (i) the larger oxygen content corresponds to a smaller heat capacity $`[C_V_i(X_i/A_i)]`$ and thus to a faster cooling prior to crystallization, (ii) the gravitational energy release at crystallization, $`\mathrm{\Delta }E\frac{\mathrm{\Delta }\rho }{\rho }Mg`$, is smaller because of the spindle form of the phase diagram (Segretain & Chabrier 1993; Segretain et al. 1994; Chabrier 1997) but (iii) the total amount of oxygen to be differentiated is larger \[cf. (i)\], and (iv) crystallization occurs earlier in the evolution since oxygen crystallizes at a higher temperature than carbon. Whether this corresponds to a larger effective temperature depends on the $`L`$-$`T_c`$ (and thus $`T_{\mathrm{eff}}`$-$`T_c`$) relation for each WD mass (Figure 2). As shown by Salaris et al. (1997), effects (ii) through (iv) more or less compensate for both reaction rate induced profiles, and the time delay induced by crystallization $`\mathrm{\Delta }\tau =\mathrm{\Delta }E/L`$ is about the same for these two stratified profiles. The difference in cooling times thus stems primarily from the available heat content \[point (i)\]. For example, at a luminosity $`\mathrm{log}(L/L_{})=4.5`$ (resp. -5.0), our reference 0.6 $`M_{}`$ model has an age of $`t=10.3`$ (resp. 12.7) Gyr whereas the same luminosity corresponds to an age $`t=10.7`$ (resp. 13.3) Gyr for a low-rate initial profile (see Figure 3), confirming the Salaris et al. (1997) analysis. This illustrates the present uncertainty in cooling times due to uncertainties in the <sup>12</sup>C$`(\alpha ,\gamma )^{16}`$O reaction rate and induced WD initial internal composition.
Another uncertainty affecting the internal composition is the exact shape of the C/O crystallization diagram. In the present calculations, we use the Segretain & Chabrier (1993) spindle diagram, calculated within the framework of the density-functional theory of freezing. The Barrat, Hansen, & Mochkovitch (1988) calculations were based on previous (obsolete) values of the plasma parameter $`\mathrm{\Gamma }`$ at crystallization but would yield similar results if updated (see Segretain & Chabrier 1993), whereas the Ichimaru, Iyetomi, & Ogata (1988) azeotropic C/O diagram is demonstrably erroneous (DeWitt, Slattery, & Chabrier 1996). Figures 3 and and 4 illustrate the effect of chemical fractionation at crystallization, due to the difference of abundance of carbon and oxygen in the fluid and in the solid phase, on a 0.6 $`M_{}`$ WD. This induces a variation of chemical potential at constant volume and temperature (Chabrier 1997) which provides an additional source of energy (Mochkovitch 1983; García-Berro et al. 1988; Segretain et al. 1994). Although this energy amounts to only $`1\%`$ of the binding energy of the star, it is released at a low luminosity and thus lengthens appreciably the lifetime of the star (Chabrier 1997; Isern et al. 1997). For example, a luminosity $`\mathrm{log}(L/L_{})=4.5`$ (resp. -5.0) for our stratified profile corresponds to an age $`t=8.8`$ (resp. 11.4) Gyr if fractionation is ignored and $`t=10.3`$ (12.7) Gyr if it is taken into account. As seen in figure 3, once crystallization has proceeded throughout the entire star, the time delay remains constant, and amounts from $`1`$ Gyr for the least (0.5 $`M_{}`$) and most (1.2 $`M_{}`$) massive WDs, to a maximum value of about 1.5 Gyr for our reference 0.6 to 0.8 $`M_{}`$ WDs. The effect thus remains substantial and must be properly taken into account in WD cooling theory. This is at odds with Hansen’s (1999) results, as explained in §3.1.1. As seen in Figure 3, crystallization in our calculations occurs at a later age and a fainter luminosity than in the Wood (1995) or Montgomery et al. (1999) calculations. This reflects the faster cooling for Wood and Montgomery’s calculations, as discussed in §3.1.1. This results in a smaller crystallization-induced delay in this latter case since $`\mathrm{\Delta }\tau 1/L`$.
Although the crystallization model for binary ionic mixtures has not yet been verified observationally in WDs, it has been studied for a long time by geophysicists (see, for example, Loper 1984; Buffett et al. 1992). Although the nature of the plasma (or alloy) is different, the physics of the process (thermodynamics and energy transport) is exactly the same. Note, however, that these crystallization-induced delays represent upper limits. Indeed, the present calculations assume complete mixing of the C-enriched fluid layers, i.e. a maximum efficiency for this process. This is supported by the fact that the mixing instability timescale is much shorter than the evolutionary timescale (Mochkovitch, 1983).
Figure 4 displays the cooling sequences $`M_{bol}(t)`$ for different masses for our reference model calculations.
### 3.2 Cooling times
Figure 5 displays the mass-$`T_{\mathrm{eff}}`$ relation, or equivalently the radius- or surface gravity-$`T_{\mathrm{eff}}`$ relation (see Table 1-5), for several constant WD cooling times. Before crystallization sets in, massive WDs evolve more slowly, because of their greater energy content ($`C_V`$) and their smaller radiative surface areas<sup>*</sup><sup>*</sup>*Remember the objects under consideration are cool enough so that the neutrino luminosity is completely negligible.. Since they are hotter and brighter at a given $`T_c`$ than less-massive WDs, and since crystallization occurs always at the same internal temperature $`T_c`$, massive WDs crystallize earlier and at a higher $`T_{\mathrm{eff}}`$ and $`L`$. At this stage, the crystallized core enters the Debye cooling regime ($`C_VT^3`$) and cools more quickly. This crystallization process causes the bending on the cooling times downward.
Figure 6 displays the same constant cooling times as a function of absolute $`M_V`$ magnitude. The bulk of all hydrogen-atmosphere WDs, those with $`m<\mathrm{\hspace{0.17em}0.8}M_{}`$ (assuming a WD mass distribution similar to the one observed in the disk), remains brighter than $`M_V18`$ after 14 Gyr.
### 3.3 Color-magnitude and color-color diagrams
Figures 7 and 8 display cooling sequences in $`M_V`$ vs. $`(VI`$) and $`M_J`$ vs. $`(JK)`$ diagrams, respectively. WDs with a small admixture of helium in the atmosphere will have similar colors (see Tables 1 and 2 of BSW). The triangles illustrate the coolest disk WDs identified spectroscopically as H-rich atmosphere WDs by Leggett, Ruiz, & Bergeron (1998). Collision-induced absorption by H<sub>2</sub> at $`T_{\mathrm{eff}}4500`$-5000 K causes the turnover to bluer $`JK`$ at $`M_J13`$-$`15`$, depending on the mass. At these temperatures, the effect is still modest and a large fraction of the flux emerges longward of 1 $`\mu `$m (see Figure 5 of BSW). As $`T_{\mathrm{eff}}`$ decreases, however, the CIA becomes stronger, causing a turnover in the optical colors near $`M_V16`$-$`17`$ at later stages of the evolution. The 12 Gyr curve forms a loop in the $`M_V`$ vs. $`(VI)`$ diagram. This stems from the very different cooling rates of WDs for different masses. As mentioned above, less massive WDs cool faster initially because of their smaller heat content, whereas massive WDs crystallize earlier and then enter the rapid Debye cooling regime. The combination of these two effects yields the slowest cooling rate for $`0.8M_{}`$ WDs after $`8`$ Gyr, whereas both the least and most massive objects cool faster, reaching fainter magnitudes and bluer colors at younger ages.
Figure 9 shows a color-color diagram comparing optical and near-IR colors for a 0.6 and a 1.2 $`M_{}`$ H-atmosphere WD. Note that gravity does not affect the qualitative behavior of the diagram. Indeed increasing gravity has an effect similar to decreasing the effective temperature on the spectrum (see, e.g., Saumon et al. 1994). Naturally the age-dependence is different, as shown on the figure by the solid circles and squares, respectively. As already mentioned in Chabrier (1999), we note from the previous figures the strong dependence of colors upon age for $`t>\mathrm{\hspace{0.17em}12}`$ Gyr. If the gravity, i.e., the mass of such a WD, can be determined independently, this provides a powerful tool to determine the age of the Galactic halo.
Tables 1-5 give the characteristic properties of the present WD H-rich atmosphere cooling sequences in various broadband filters for the mass-range characteristic of WDs with C/O cores.
## 4 Conclusion
We have computed evolutionary sequences for cool ($`T_{\mathrm{eff}}<\mathrm{10\hspace{0.17em}000}`$ K) pure hydrogen atmosphere WDs, although the calculations can be applied to WDs with a small admixture of He in the atmosphere as well, using appropriate color indices (BSW). These models are primarily aimed at identifying old, cool WDs in the Galactic old disk, spheroid or dark halo, or in globular clusters once fainter detection limits can be achieved. We have first improved upon the equation of state of the fully ionized interior to calculate accurate cooling sequences for crystallized WDs, namely objects fainter than $`\mathrm{log}(L/L_{})4.5`$ ($`t>\mathrm{\hspace{0.17em}9}`$ Gyr for a 0.6 $`M_{}`$). These WDs are in a quantum regime. The temperature-dependent contributions, in particular the ion-electron energy of a quantum, relativistic ion-electron plasma, are the only sources of specific heat and thus determine entirely the thermal reservoir to be radiated into space, even though they provide a negligible contribution to the internal energy. We have next obtained improved $`L`$-$`T_c`$ relations which govern the rate of cooling of white dwarfs. The envelope and atmosphere calculations include state-of-the-art radiative opacities down to $`T_{\mathrm{eff}}=1500`$ K for pure hydrogen. We have examined and quantified the uncertainties in these cooling sequences arising from (i) the initial C/O stratification of the WD, which follows from the <sup>12</sup>C$`(\alpha ,\gamma )^{16}`$O reaction rate, (ii) the $`LT_c`$ relation (the core-surface boundary condition), and (iii) the crystallization-induced gravitational energy release. The uncertainty in the <sup>12</sup>C$`(\alpha ,\gamma )^{16}`$O reaction rate translates into a $`<\mathrm{\hspace{0.17em}0.5}`$ Gyr difference in cooling time. The uncertainty due to different $`L`$-$`T_c`$ relations when using accurate boundary conditions yields substantial differences in the range $`(L/L_{})10^{4.0}`$$`10^{4.5}`$, when, for the first time, the superficial convection zone reaches the degenerate core. This is probably due to the detailed physics entering the core (EOS) and envelope (EOS, $`_{\mathrm{rad}}`$, $`_{\mathrm{cond}}`$) calculations. These differences, however, vanish almost entirely below this limit, i.e., for $`t>\mathrm{\hspace{0.17em}10}`$ Gyr. The maximum time delay induced by chemical fractionation at crystallization amounts to $`1`$ to 1.5 Gyr, depending on the WD mass, and represents one of the major uncertainties in dating very cool WDs from their observed luminosity. Future laser-driven experiments on the crystallization of dense plasmas should shed light on this complicated physics problem, which bears major consequences in the present astrophysical context. The other major uncertainty in WD cooling stems from the ill-determined mass fraction of hydrogen and helium in the external envelope, in particular for old WDs originating from metal-depleted progenitors.
We find noticeable differences ($`>\mathrm{\hspace{0.17em}1}`$ Gyr) between our and Hansen’s (1999) cooling sequences. Part of it very likely stems from the different treatments of the interior EOS and envelope EOS and opacities, but also from the demonstrably underestimated crystallization-induced time delay in Hansen’s calculations. A $`0.5`$-1 mag difference in colors occurs also for $`VI<\mathrm{\hspace{0.17em}1}`$ between Hansen’s and our calculations, which most likely arises from details of the calculation of the atmosphere models, including the non-ideal effects (see Figure 3 of SJ). This shows that, although a consistent general theory for the cooling of cool WDs is emerging, work remains to be done to reach more robust results. The theory is still very uncertain for very cool helium-rich atmosphere WDs. A correct calculation of helium pressure ionization and of the He<sup>-</sup> free-free absorption cross-section at high densities in the atmosphere of these objects remains to be done. Attention must also be devoted to dense atmospheres including a trace of metals since the subsequent increase in opacity will affect dramatically the cooling of the star. Work in this direction is in progress.
Meanwhile, the present calculations should provide what we believe to be presently the most accurate calculations of very cool white dwarf sequences. They provide a useful basis to search for and to identify faint, old WDs either in the field or in globular clusters. By allowing the determination of the mass and age of possible halo WDs (Hodgkin et al. 2000, Ibata et al. 2000), they will also provide important constraints on the age of the Galactic disk and halo, and on the Galaxy initial mass function.
Note: The present models in various filters are available upon request to Gilles Chabrier (chabrier@ens-lyon.fr).
We thank M. A. Wood and B. Hansen for kindly providing models and results from their calculations, M. Hernanz for providing tables of the stratified profiles and A. Potekhin for stimulating discussions. This work was supported in part by NSF grant AST-9731438 to D.S.
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# Local cohomology of generalized Okamoto–Painlevé pairs and Painlevé equations
## 1. Introduction
To study the spaces of initial conditions of Painlevé equations constructed by Okamoto \[O1\] \[O2\] \[O3\], we introduced the notion of generalized Okamoto-Painlevé pair $`(S,Y)`$ in \[STT\]. This is a pair of a complex projective surface $`S`$ and an anti-canonical divisor $`Y|K_S|`$ of $`S`$ satisfying the following conditions: For the irreducible decomposition $`Y=_{i=1}^rm_iY_i`$, one has $`YY_i=\mathrm{deg}Y|_{Y_i}=0`$ for $`1ir`$. In addition, if $`S`$ is a rational surface, $`(S,Y)`$ is called a generalized rational Okamoto-Painlevé pair.
The generalized rational Okamoto-Painlevé pairs of non-fibered type are classified into three types; elliptic, multiplicative and additive. Moreover, the generalized rational Okamoto-Painlevé pairs of additive type such that $`D:=Y_{red}=_{i=1}^rY_i`$ is a divisor with only normal crossings correspond to the spaces of initial conditions of the Painlevé equations (cf. \[STT\] \[Sakai\]) as follows:
| type of $`(S,Y)`$ | $`\stackrel{~}{E_8}`$ | $`\stackrel{~}{E_7}`$ | $`\stackrel{~}{D_8}`$ | $`\stackrel{~}{D_7}`$ | $`\stackrel{~}{D_6}`$ | $`\stackrel{~}{E_6}`$ | $`\stackrel{~}{D_5}`$ | $`\stackrel{~}{D_4}`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| Painlevé equations | $`P_I`$ | $`P_{II}`$ | $`P_{III}^{\stackrel{~}{D_8}}`$ | $`P_{III}^{\stackrel{~}{D_7}}`$ | $`P_{III}=P_{III}^{\stackrel{~}{D_6}}`$ | $`P_{IV}`$ | $`P_V`$ | $`P_{VI}`$ |
In what follows, $`(S,Y)`$ is a generalized rational Okamoto-Painlevé pair of non-fibered type satisfying the condition: $`D=Y_{red}`$ is a normal crossing divisor with at least two irreducible components so that all irreducible components of $`Y_{red}`$ are smooth rational curves. Let $`\mathrm{\Theta }_S(\mathrm{log}D)`$ be the sheaf of regular vector fields which have logarithmic zero along $`D`$. Here and after, all sheaves of $`𝒪_S`$-modules are considered in algebraic category.
Then we have the following key exact sequence:
$$0H_D^1(\mathrm{\Theta }_S(\mathrm{log}D))H^1(S,\mathrm{\Theta }_S(\mathrm{log}D))\stackrel{res}{}H^1(SD,\mathrm{\Theta }_S(\mathrm{log}D)),$$
where $`H^1(S,\mathrm{\Theta }_S(\mathrm{log}D))`$ and $`H^1(SD,\mathrm{\Theta }_S(\mathrm{log}D))`$ are the space of infinitesimal deformations of the pair $`(S,Y)`$ (cf. \[Kaw\]) and the space of infinitesimal deformations of $`SY`$, respectively.
In \[STT\], we show that the directions corresponding to local cohomology $`H_D^1(\mathrm{\Theta }_S(\mathrm{log}D))`$ in the deformation of a pair $`(S,Y)`$ induce differential equations on $`SY`$, by generalizing the Kodaira-Spencer theory to the open surface $`SY`$.
In this paper, we will show
$$dimH_D^1(\mathrm{\Theta }_S(\mathrm{log}D))1,$$
when $`(S,Y)`$ is of additive type with the normal crossing divisor $`D=Y_{red}`$ (Theorem 2.1). This result is natural since $`SY`$ corresponds to the space of initial conditions of the Painlevé equations in this case. On the other hand, this is not always the case. In fact, we will prove that
$$H_D^1(\mathrm{\Theta }_S(\mathrm{log}D))=\{0\},$$
for pairs $`(S,Y)`$ of type $`\stackrel{~}{A_8}`$ which is classified as a multiplicative type (Proposition 4.1). This means that there does not exist differential equation on $`SY`$ by the way above. All these computations are carried out through calculations of Čech cohomologies by taking a coordinate system explicitly.
## 2. Local cohomology sequences and Time variables
Let $`(S,Y)`$ be a generalized rational Okamoto–Painlevé pair and set $`D=Y_{red}`$. Moreover, in this section, we assume that
1. $`(S,Y)`$ is of non-fibered type and
2. $`Y_{red}`$ is a simple normal crossing divisor with at least two irreducible components, i.e. $`(r2)`$ so that all irreducible components of $`Y_{red}`$ are smooth rational curves.
Here $`(S,Y)`$ is called of fibered type if $`S`$ has a structure of an elliptic surfacefibration $`f:S^1`$ with $`f^{}(\mathrm{})=nY`$ for some $`n1`$. If $`(S,Y)`$ is not of fibered type, we call $`(S,Y)`$ of non-fibered type. (cf. Definition 1.3, \[STT\]).
In what follows, $`𝒪_S`$ and $`𝒪_{SD}`$ denote the sheaves of germs of algebraic regular functions on $`S`$ and $`SD`$ respectively. Moreover all sheaves of $`𝒪_S`$-modules are considered in algebraic category unless otherwise stated. Let us consider the following exact sequence of local cohomology groups (\[Corollary 1.9, \[Gr\]\])
$$\begin{array}{cccccc}H^0(S,\mathrm{\Theta }_S(\mathrm{log}D))& & H^0(SD,\mathrm{\Theta }_S(\mathrm{log}D))& & H_D^1(\mathrm{\Theta }_S(\mathrm{log}D))& \\ H^1(S,\mathrm{\Theta }_S(\mathrm{log}D))& \stackrel{\mu }{}& H^1(SD,\mathrm{\Theta }_S(\mathrm{log}D)).& & & \end{array}$$
Since $`(S,Y)`$ is of non-fibered type, from \[(2), Proposition 2.1 \[STT\]\], we see that
$$H^0(SD,\mathrm{\Theta }_S(\mathrm{log}D))=H^0(SD,\mathrm{\Theta }_S)=\{0\}.$$
Hence we have the following exact sequence:
(1)
$$\begin{array}{cccccc}0& H_D^1(\mathrm{\Theta }_S(\mathrm{log}D))& & H^1(S,\mathrm{\Theta }_S(\mathrm{log}D))& \stackrel{\mu }{}& H^1(SD,\mathrm{\Theta }_S(\mathrm{log}D)).\end{array}$$
The following theorem is the main statement in this paper.
###### Theorem 2.1.
Let $`(S,Y)`$ be a generalized rational Okamoto-Painlevé pair $`(S,Y)`$ with the condition above. Moreover assume that $`D=Y_{red}`$ is of additive type. Then we have
$$dimH^0(D,\mathrm{\Theta }_S(\mathrm{log}D)N_D)=1.$$
Here we put $`N_D=𝒪_S(D)/𝒪_S`$.
Since we have a natural inclusion
$$H^0(D,\mathrm{\Theta }_S(\mathrm{log}D)N_D)H_D^1(\mathrm{\Theta }_S(\mathrm{log}D)),$$
we obtain
$$dimH_D^1(\mathrm{\Theta }_S(\mathrm{log}D))1.$$
This theorem plays an important role to understand the Painlevé equation related to $`(S,Y)`$.(cf. \[STT\]). Though we will not investigate the further structure of local cohomology here. Instead, we propose the following
###### Conjecture 2.1.
Under the same notation and assumption as in Theorem 2.1,
$$H_D^1(\mathrm{\Theta }_S(\mathrm{log}D))H^0(D,\mathrm{\Theta }_S(\mathrm{log}D)N_D).$$
From the exact sequence (1), we see that the subspace $`H_D^1(S,\mathrm{\Theta }_S(\mathrm{log}D))`$ of $`H^1(S,\mathrm{\Theta }_S(\mathrm{log}D))`$ coincides with the kernel of $`\mu `$. This implies that:
$$H_D^1(S,\mathrm{\Theta }_S(\mathrm{log}D))\left\{\begin{array}{c}\text{Infinitesimal deformations of }(S,D)\text{ whose restriction}\\ \text{to }SD\text{ induces the trivial deformation}\end{array}\right\}.$$
In \[§6 \[STT\]\], we show that any direction corresponding to a non-zero element of the local cohomology group $`H_D^1(S,\mathrm{\Theta }_S(\mathrm{log}D))`$ induces a differential equation (at least locally) by using Čech coboundaries.
At this moment, we can not prove Conjecture 2.1 with full generality. However, we see that the one dimensional vector subspace $`H^1(D,\mathrm{\Theta }_S(\mathrm{log}D)N_D)`$ of $`H_D^1(\mathrm{\Theta }_S(\mathrm{log}D))H^1(\mathrm{\Theta }_S(\mathrm{log}D))`$ really corresponds to the time variable $`t`$ in the known Painlevé equation. It is unlikely that we will have more time variables, so this gives an evidence of Conjecture 2.1.
###### Remark 2.1.
We will consider $`(S,Y)`$ of multiplicative type later, where the situation is different. (cf. Proposition 4.1).
Let us make preparations for the proof of Theorem 2.1.
Recall that
$$H_D^1(S,\mathrm{\Theta }_S(\mathrm{log}D))=\underset{}{\mathrm{lim}}\mathrm{Ext}^1(𝒪_{nD},\mathrm{\Theta }_S(\mathrm{log}D))$$
where $`𝒪_{nD}=𝒪_S/𝒪_S(nD)`$ (cf. \[Theorem 2.8, \[Gr\]\]).
On the other hand, since $`\mathrm{\Theta }_S(\mathrm{log}D)`$ is a locally free $`𝒪_S`$-module, we see that
$$om(𝒪_{nD},\mathrm{\Theta }_S(\mathrm{log}D))=0,$$
and
$$xt^1(𝒪_{nD},\mathrm{\Theta }_S(\mathrm{log}D))=\mathrm{\Theta }_S(\mathrm{log}D)N_{nD},$$
where $`N_{nD}=𝒪_S(nD)/𝒪_S`$. By an argument using a spectral sequence, we see that
$$H_D^1(S,\mathrm{\Theta }_S(\mathrm{log}D))=\underset{}{\mathrm{lim}}H^0(\mathrm{\Theta }_S(\mathrm{log}D)N_{nD}).$$
Hence we have a natural inclusion
$$H^0(\mathrm{\Theta }_S(\mathrm{log}D)N_D)H_D^1(S,\mathrm{\Theta }_S(\mathrm{log}D)).$$
###### Lemma 2.1.
Let $`(S,Y)`$ be a generalized rational Okamoto–Painlevé pair as above and set $`D=Y_{red}`$. Then we have the following exact sequences :
(2)
$$0\mathrm{\Theta }_DN_D\mathrm{\Theta }_SN_D\nu _{}(_{i=1}^rN_{Y_i/S})N_D0,$$
(3)
$$0\nu _{}(_{i=1}^rN_{Y_i/S})\mathrm{\Theta }_S(\mathrm{log}D)N_D\mathrm{\Theta }_DN_D0.$$
Here $`\mathrm{\Theta }_D`$ denotes the tangent sheaf of $`D`$ and $`\nu :\stackrel{~}{D}=_{i=1}^9Y_iD`$ the normalization map.
Proof. The first exact sequence (2) follows from \[(1.9), \[B-W\]\].
Let us consider the following diagram:
$$\begin{array}{ccccccc}& 0& & 0& & \mathrm{ker}\lambda & \\ & & & & & & \\ & & & & & & \\ 0& \mathrm{\Theta }_S(\mathrm{log}D)& & \mathrm{\Theta }_S(\mathrm{log}D)𝒪_S(D)& & \mathrm{\Theta }_S(\mathrm{log}D)N_D& 0\\ & & & & & \lambda & \\ 0& \mathrm{\Theta }_S& & \mathrm{\Theta }_S𝒪_S(D)& & \mathrm{\Theta }_SN_D& 0\\ & & & & & & \\ & \nu _{}(_{i=1}^rN_{Y_i/S})& \stackrel{\mu }{}& \nu _{}(_{i=1}^rN_{Y_i/S})N_D& & coker\lambda & 0\\ & & & & & & \\ & 0& & 0& & 0& .\end{array}$$
By the snake lemma, we obtain the exact sequence
$$0\mathrm{ker}\lambda \nu _{}(_{i=1}^rN_{Y_i/S})\stackrel{\mu }{}\nu _{}(_{i=1}^rN_{Y_i/S})N_Dcoker\lambda 0.$$
From a local consideration, we see that the map $`\mu `$ is a zero map, hence
$$\mathrm{ker}\lambda \nu _{}(_{i=1}^rN_{Y_i/S}),\text{and}\nu _{}(_{i=1}^rN_{Y_i/S})N_Dcoker\lambda .$$
Moreover since $`Im\lambda \mathrm{ker}[\mathrm{\Theta }_SN_D\nu _{}(_{i=1}^rN_{Y_i/S})N_D]`$, from the exact sequence (2), we obtain the exact sequence (3).∎
Note that since $`N_{Y_i/S}=𝒪_{Y_i}(2)`$, we have
$$H^0(_{i=1}^rN_{Y_i/S})=\{0\},H^1(_{i=1}^rN_{Y_i/S})^r.$$
Moreover one can easily see that
$$\mathrm{\Theta }_D\nu _{}(_{i=1}^r\mathrm{\Theta }_{Y_i}(t_i))\nu _{}((_{i=1}^r𝒪_{Y_i}(2t_i))$$
where $`t_i`$ is the number of intersections of $`Y_i`$ with the other $`Y_j`$. On the other hand, since $`DY_i=t_i2`$ and $`\nu `$ is a finite morphism, we see that
$$\begin{array}{cc}H^0(D,\mathrm{\Theta }_DN_D)& H^0(D,\nu _{}(_{i=1}^r\mathrm{\Theta }_{Y_i}(t_i))N_D)\hfill \\ & H^0(\stackrel{~}{D},(_{i=1}^r\mathrm{\Theta }_{Y_i}(t_i))\nu ^{}(N_D))\hfill \\ & _{i=1}^rH^0(Y_i,\mathrm{\Theta }_{Y_i}(t_i)N_D)\hfill \\ & _{i=1}^rH^0(Y_i,𝒪_{Y_i})\hfill \\ & ^r.\hfill \end{array}$$
Proof of Theorem 2.1.
From the exact sequence (3), one can obtain
$$H^0(_{i=1}^rN_{Y_i/S})H^0(\mathrm{\Theta }_S(\mathrm{log}D)N_D)H^0(\mathrm{\Theta }_DN_D)\stackrel{\delta }{}H^1(_{i=1}^rN_{Y_i/S})$$
where $`\delta `$ denotes the connected homomorphism.
From this sequence, the connecting homomorphism $`\delta `$
$$\delta :H^0(\mathrm{\Theta }_DN_D)_{i=1}^rH^1(N_{Y_i/S})$$
can be identified with a linear map $`\delta :^r^r`$ and we have an isomorphism
$$H^0(D,\mathrm{\Theta }_S(\mathrm{log}D)N_D)\mathrm{ker}\delta .$$
Now we state the following proposition, the proof of which is given in §3.
###### Proposition 2.1.
Let $`(S,Y)`$ be a generalized rational Okamoto–Painlevé pair of additive type and set $`D=Y_{red}`$. Assume that $`D`$ is a normal crossing divisor, then we can choose a basis of $`H^0(\mathrm{\Theta }_DN_D)`$ and $`_{i=1}^rH^1(N_{Y_i})`$ so that the linear map $`\delta :H^0(\mathrm{\Theta }_DN_D)_{i=1}^rH^1(N_{Y_i})`$ is represented by the intersection matrix $`((Y_iY_j))_{1i,jr}`$ with respect to these basis.
In this case, it is well-known that the rank of the intersection matrix $`((Y_iY_j))_{1i,jr}`$ is $`r1`$. Hence from the Proposition 2.1, we have
$$dim_{}H^0(\mathrm{\Theta }_S(\mathrm{log}D)N_D)=dim_{}\mathrm{ker}\delta =1.$$
## 3. Proof of Proposition 2.1
In this section, we shall prove Proposition 2.1.
Here we give a detailed proof only for the case of $`\stackrel{~}{E_7}`$. The proof of other cases are similar.
Let $`(S,Y)`$ be a generalized rational Okamoto–Painlevé pair of type $`\stackrel{~}{E_7}`$. Then according to the results in Appendix B of \[Sakai\], $`(S,Y)`$ can be obtained by blowings up 9-points of $`^2`$ as follows.
Let $`[x:y:z]`$ be the homogeneous coordinates of $`^2`$.
$$\begin{array}{c}p_1:(0:1:0)p_2:(\frac{x}{y},\frac{z}{x})=(0,0)p_3:(\frac{x}{y},\frac{yz}{x^2})=(0,0)p_6:(\frac{x}{y},\frac{y^2z}{x^3})=(0,1)\hfill \\ p_7:(\frac{x}{y},\frac{y(y^2zx^3)}{x^4})=(0,0)p_8:(\frac{x}{y},\frac{y^2(y^2zx^3)}{x^5})=(0,s)\hfill \\ p_9:(\frac{x}{y},\frac{y(y^2(y^2zx^3)+sx^5)}{x^6})=(0,\alpha _0),\hfill \\ \\ p_4:(0:0:1),\hfill \\ p_5:(0:\alpha _1:1).\hfill \end{array}$$
Note that there exist three deformation parameters $`(\alpha _0,\alpha _1,s)`$ of the blowings-up. Moreover there exists a $`^\times `$-action on the family of surfaces by
$$(\mu ,(\alpha _0,\alpha _1;s,[x:y:z]))(\mu ^3\alpha _0,\mu ^3\alpha _1,\mu ^2s;[x:\mu y:\mu ^2z]).$$
If we set $`\lambda =\alpha _0+\alpha _1`$ and $`\mu =s/\lambda `$, then
$$t=\frac{s^3}{\lambda ^2},a_0=\frac{s^3\alpha _0}{\lambda ^3},a_1=\frac{s^3\alpha _1}{\lambda ^3},X=x,Y=\frac{sy}{\lambda },Z=\frac{\lambda ^2z}{s^2}$$
are invariant under the $`^\times `$-action, and we have the relation
$$a_0+a_1=t.$$
Now we can introduce the affine open covering of $`S`$ and affine coordinates by the explicit blowings-up of $`^2`$. The following diagram shows how one can perform blowing-ups and introduce the new coordinates $`(x_i,y_i)(1i13),(u_j,v_j)(j=14,15,16)`$.
$$\begin{array}{c}E_i^1,E_i^2=1,\\ Y_i^1,Y_i^2=2.\end{array}$$
$`U_i`$ $`=`$ $`Spec[x_i,y_i]^2(i=1,2,\mathrm{},7).`$
$`U_8`$ $`=`$ $`Spec[x_8,y_8,{\displaystyle \frac{1}{1+x_8}}]^2\{1+x_8=0\}.`$
$`U_9`$ $`=`$ $`Spec[x_9,y_9,{\displaystyle \frac{1}{1+x_{9}^{}{}_{}{}^{2}y_9}}]^2\{1+x_{9}^{}{}_{}{}^{2}y_9=0\}.`$
$`U_{10}`$ $`=`$ $`Spec[x_{10},y_{10},{\displaystyle \frac{1}{1+x_{10}y_{10}^{}{}_{}{}^{2}}}]^2\{1+x_{10}y_{10}^{}{}_{}{}^{2}=0\}.`$
$`U_{11}`$ $`=`$ $`Spec[x_{11},y_{11},{\displaystyle \frac{1}{1tx_{11}^{}{}_{}{}^{2}y_{11}^{}{}_{}{}^{2}+x_{11}^{}{}_{}{}^{3}y_{11}^{}{}_{}{}^{2}}}]^2\{1tx_{11}^{}{}_{}{}^{2}y_{11}^{}{}_{}{}^{2}+x_{11}^{}{}_{}{}^{3}y_{11}^{}{}_{}{}^{2}=0\}.`$
$`U_{12}`$ $`=`$ $`Spec[x_{12},y_{12},{\displaystyle \frac{1}{1ty_{12}^{}{}_{}{}^{2}+x_{12}y_{12}^{}{}_{}{}^{3}}}]^2\{1ty_{12}^{}{}_{}{}^{2}+x_{12}y_{12}^{}{}_{}{}^{3}=0\}.`$
$`U_{13}`$ $`=`$ $`Spec[x_{13},y_{13},{\displaystyle \frac{1}{1+tx_{13}^{}{}_{}{}^{2}y_{13}^{}{}_{}{}^{2}+a_0x_{13}^{}{}_{}{}^{3}y_{13}^{}{}_{}{}^{3}x_{13}^{}{}_{}{}^{4}y_{13}^{}{}_{}{}^{3}}}]`$
$``$ $`^2\{1+tx_{13}^{}{}_{}{}^{2}y_{13}^{}{}_{}{}^{2}+a_0x_{13}^{}{}_{}{}^{3}y_{13}^{}{}_{}{}^{3}x_{13}^{}{}_{}{}^{4}y_{13}^{}{}_{}{}^{3}=0\}.`$
$`U_j`$ $`=`$ $`Spec[u_j,v_j]^2(j=14,15).`$
$`U_{16}`$ $`=`$ $`Spec[u_{16},v_{16},{\displaystyle \frac{1}{1+tu_{16}^2+a_0u_{16}^3u_{16}^4v_{16}}}]^2\{1+tu_{16}^2+a_0u_{16}^3u_{16}^4v_3=0\}.`$
$$\begin{array}{cc}Y_1=\{x_1=0,x_2=0,y_3=0\},\hfill & Y_2=\{x_3=0,y_4=0\},\hfill \\ Y_3=\{x_4=0,y_5=0\},\hfill & Y_4=\{x_5=0,y_6=0,y_8=0\},\hfill \\ Y_5=\{x_6=0,y_7=0\},\hfill & Y_6=\{x_8=0,y_9=0\},\hfill \\ Y_7=\{x_9=0,y_{10}=0,y_{11}=0\},\hfill & Y_8=\{x_{11}=0,y_{12}=0,y_{13}=0\}.\hfill \end{array}$$
$$\begin{array}{c}S=\underset{i=1}{\overset{16}{}}U_i.\\ Y=Y_1+2Y_2+3Y_3+4Y_4+2Y_5+3Y_6+2Y_7+Y_8,D=\underset{i=1}{\overset{8}{}}Y_i.\\ SY=U_{14}U_{15}U_{16}.\end{array}$$
More explicitly, new coordinates can be given by the following formula
$$\begin{array}{c}\{\begin{array}{c}x_1=\frac{X}{Y}\hfill \\ y_1=\frac{Y}{Z}\hfill \end{array}\{\begin{array}{c}x_2=\frac{X}{Ya_1Z}\hfill \\ y_2=\frac{Ya_1Z}{Z}\hfill \end{array}\{\begin{array}{c}x_3=\frac{Z}{Y}\hfill \\ y_3=\frac{X}{Z}\hfill \end{array}\\ \{\begin{array}{c}x_4=\frac{Z}{X}\hfill \\ y_4=\frac{X^2}{YZ}\hfill \end{array}\{\begin{array}{c}x_5=\frac{YZ}{X^2}\hfill \\ y_5=\frac{X^3}{Y^2Z}\hfill \end{array}\{\begin{array}{c}x_6=\frac{Y^2Z}{X^3}\hfill \\ y_6=\frac{X}{Y}\hfill \end{array}\\ \{\begin{array}{c}x_7=\frac{Y}{X}\hfill \\ y_7=\frac{Z}{X}\hfill \end{array}\{\begin{array}{c}x_8=\frac{Y^2ZX^3}{X^3}\hfill \\ y_8=\frac{X^4}{Y(Y^2ZX^3)}\hfill \end{array}\{\begin{array}{c}x_9=\frac{Y(Y^2ZX^3)}{X^4}\hfill \\ y_9=\frac{X^5}{Y^2(Y^2ZX^3)}\hfill \end{array}\\ \{\begin{array}{c}x_{10}=\frac{Y^2(Y^2ZX^3)}{X^5}\hfill \\ y_{10}=\frac{X}{Y}\hfill \end{array}\{\begin{array}{c}x_{11}=\frac{Y^2(Y^2ZX^3)+tX^5}{X^5}\hfill \\ y_{11}=\frac{X^6}{Y(Y^2(Y^2ZX^3)+tX^5)}\hfill \end{array}\\ \{\begin{array}{c}x_{12}=\frac{Y(Y^2(Y^2ZX^3)+tX^5)}{X^6}\hfill \\ y_{12}=\frac{X}{Y}\hfill \end{array}\{\begin{array}{c}x_{13}=\frac{Y(Y^2(Y^2ZX^3)+tX^5)+a_0X^6}{X^6}\hfill \\ y_{13}=\frac{X^7}{Y(Y(Y^2(Y^2ZX^3)+tX^5)+a_0X^6)}\hfill \end{array}\\ \{\begin{array}{c}u_{14}=\frac{X}{Z}\hfill \\ v_{14}=\frac{Y}{X}\hfill \end{array}\{\begin{array}{c}u_{15}=\frac{X}{Z}\hfill \\ v_{15}=\frac{Ya_1Z}{X}\hfill \end{array}\{\begin{array}{c}u_{16}=\frac{X}{Y}\hfill \\ v_{16}=\frac{Y(Y(Y^2(Y^2ZX^3)+tX^5)+a_0X^6)}{X^7}.\hfill \end{array}\end{array}$$
From these formula, we can determine the coordinate transformation between $`(x_i,y_i)`$’s and $`(u_j,v_j)`$’s.
For later use, we need only the coordinate transformations near each component $`Y_i`$. Here we will list up the coordinate transformations only for a neighborhood of each $`Y_i`$.
$$\begin{array}{cc}Y_1:\{\begin{array}{c}x_1=\frac{x_2y_2}{a_1+y_2}\hfill \\ y_1=a_1+y_2\hfill \end{array}\{\begin{array}{c}x_1=x_3y_3\hfill \\ y_1=\frac{1}{x_3}\hfill \end{array}\hfill & Y_5:\{\begin{array}{c}x_6=x_{7}^{}{}_{}{}^{2}y_7\hfill \\ y_6=\frac{1}{x_7}\hfill \end{array}\hfill \\ Y_2:\{\begin{array}{c}x_3=x_{4}^{}{}_{}{}^{2}y_4\hfill \\ y_3=\frac{1}{x_4}\hfill \end{array}\hfill & Y_6:\{\begin{array}{c}x_8=x_{9}^{}{}_{}{}^{2}y_9\hfill \\ y_8=\frac{1}{x_9}\hfill \end{array}\hfill \\ Y_3:\{\begin{array}{c}x_4=x_{5}^{}{}_{}{}^{2}y_5\hfill \\ y_4=\frac{1}{x_5}\hfill \end{array}\hfill & Y_7:\{\begin{array}{c}x_9=x_{10}y_{10}\hfill \\ y_9=\frac{1}{x_{10}}\hfill \end{array}\{\begin{array}{c}x_9=x_{11}(t+x_{11})y_{11}\hfill \\ y_9=\frac{1}{t+x_{11}}\hfill \end{array}\hfill \\ Y_4:\{\begin{array}{c}x_5=x_6y_6\hfill \\ y_5=\frac{1}{x_6}\hfill \end{array}\{\begin{array}{c}x_5=x_8(1+x_8)y_8\hfill \\ y_5=\frac{1}{1+x_8}\hfill \end{array}\hfill & Y_8:\{\begin{array}{c}x_{11}=x_{12}y_{12}\hfill \\ y_{11}=\frac{1}{x_{12}}\hfill \end{array}\{\begin{array}{c}x_{11}=x_{13}(a_0+x_{13})y_{13}\hfill \\ y_{11}=\frac{1}{a_0+x_{13}}.\hfill \end{array}\hfill \end{array}$$
Let us consider the sheaf $`\mathrm{\Theta }_S(\mathrm{log}D)`$ and the sheaf exact sequence
$$0\nu _{}(_{i=1}^8N_{Y_i/S})\mathrm{\Theta }_S(\mathrm{log}D)N_D\mathrm{\Theta }_DN_D0.$$
We will analyse the edge homomorphism
(4)
$$\delta :H^0(\mathrm{\Theta }_DN_D)H^1(\nu _{}(_{i=1}^8N_{Y_i/S}))$$
by using the Čech cocycles.
Noting that $`\nu `$ is a finite morphism, and $`\mathrm{\Theta }_DN_D\nu _{}(_{i=1}^8\mathrm{\Theta }_{Y_i}(t_i))N_D`$ where $`t_i`$ is the number of intersections of $`Y_i`$ with other components, we see that
(5)
$$H^0(D,\mathrm{\Theta }_DN_D)H^0(D,\nu _{}(_{i=1}^8\mathrm{\Theta }_{Y_i}(t_i))N_D)^8.$$
For each $`i(1i8)`$, we introduce a generator $`\theta _i`$ of the cohomology group in (5) as follows.
| $`\theta _1`$ | $`\left\{\begin{array}{c}\theta _1^1=\frac{a_1+y_1}{x_1}\frac{}{y_1}\text{on}U_1Y_1,\theta _1^2=\frac{a_1+y_2}{x_2}\frac{}{y_2}\text{on}U_2Y_1,\hfill \\ \theta _1^3=\frac{1+a_1x_3}{y_3}\frac{}{x_3}\text{on}U_3Y_1\hfill \end{array}\right\}`$ |
| --- | --- |
| $`\theta _2`$ | $`\left\{\theta _2^3=\frac{1}{x_3}\frac{}{y_3}\text{on}U_3Y_2,\theta _2^4=\frac{1}{y_4}\frac{}{x_4}\text{on}U_4Y_2\right\}`$ |
| $`\theta _3`$ | $`\left\{\theta _3^4=\frac{1}{x_4}\frac{}{y_4}\text{on}U_4Y_3,\theta _3^5=\frac{1}{y_5}\frac{}{x_5}\text{on}U_5Y_3\right\}`$ |
| $`\theta _4`$ | $`\left\{\begin{array}{c}\theta _4^5=\frac{1y_5}{x_5}\frac{}{y_5}\text{on}U_5Y_4,\theta _4^6=\frac{1x_6}{y_6}\frac{}{x_6}\text{on}U_6Y_4,\hfill \\ \theta _4^8=\frac{1}{y_8}\frac{}{x_8}\text{on}U_8Y_4\hfill \end{array}\right\}`$ |
| $`\theta _5`$ | $`\left\{\theta _5^6=\frac{1}{x_6}\frac{}{y_6}\text{on}U_6Y_5,\theta _5^7=\frac{1}{y_7}\frac{}{x_7}\text{on}U_7Y_5\right\}`$ |
| $`\theta _6`$ | $`\left\{\theta _6^8=\frac{1}{x_8}\frac{}{y_8}\text{on}U_8Y_6,\theta _6^9=\frac{1}{y_9}\frac{}{x_9}\text{on}U_8Y_6\right\}`$ |
| $`\theta _7`$ | $`\left\{\begin{array}{c}\theta _7^9=\frac{1+ty_9}{x_9}\frac{}{y_9}\text{on}U_9Y_7,\theta _7^{10}=\frac{t+x_{10}}{y_{10}}\frac{}{x_{10}}\text{on}U_{10}Y_7,\hfill \\ \theta _7^{11}=\frac{1}{y_{11}}\frac{}{x_{11}}\text{on}U_{11}Y_7\hfill \end{array}\right\}`$ |
| $`\theta _8`$ | $`\left\{\begin{array}{c}\theta _8^{11}=\frac{1+a_0y_{11}}{x_{11}}\frac{}{y_{11}}\text{on}U_{11}Y_8,\theta _8^{12}=\frac{a_0+x_{12}}{y_{12}}\frac{}{x_{12}}\text{on}U_{12}Y_8,\hfill \\ \theta _8^{13}=\frac{1}{y_{13}}\frac{}{x_{13}}\text{on}U_{13}Y_8\hfill \end{array}\right\}`$ |
On the other hand, for each $`i(1i8)`$, we have a generator $`\eta _iH^1(Y_i,N_{Y_i/S})`$ as follows.
| $`\eta _1`$ | $`\left\{\eta _1^{12}=0\text{on}U_1U_2Y_1,\eta _1^{13}=\frac{1}{x_3}\frac{}{y_3}\text{on}U_1U_3Y_1\right\}`$ |
| --- | --- |
| $`\eta _2`$ | $`\left\{\eta _2^{34}=\frac{1}{x_4}\frac{}{y_4}\text{on}U_3U_4Y_2\right\}`$ |
| $`\eta _3`$ | $`\left\{\eta _3^{45}=\frac{1}{x_5}\frac{}{y_5}\text{on}U_4U_5Y_3\right\}`$ |
| $`\eta _4`$ | $`\left\{\eta _4^{56}=0\text{on}U_5U_6Y_4,\eta _4^{58}=\frac{1}{x_8}\frac{}{y_8}\text{on}U_5U_8Y_4\right\}`$ |
| $`\eta _5`$ | $`\left\{\eta _5^{67}=\frac{1}{x_7}\frac{}{y_7}\text{on}U_6U_7Y_5\right\}`$ |
| $`\eta _6`$ | $`\left\{\eta _6^{89}=\frac{1}{x_9}\frac{}{y_9}\text{on}U_8U_9Y_6\right\}`$ |
| $`\eta _7`$ | $`\left\{\eta _7^{\mathrm{9\hspace{0.17em}10}}=0\text{on}U_9U_{10}Y_7,\eta _7^{\mathrm{9\hspace{0.17em}11}}=\frac{1}{x_{11}}\frac{}{y_{11}}\text{on}U_9U_{11}Y_7\right\}`$ |
| $`\eta _8`$ | $`\left\{\eta _8^{\mathrm{11\hspace{0.17em}13}}=\frac{1}{x_{13}}\frac{}{y_{13}}\text{on}U_{11}U_{13}Y_8,\eta _8^{\mathrm{12\hspace{0.17em}13}}=0\text{on}U_{12}U_{13}Y_8\right\}`$ |
We take $`\{\theta _i\}`$ and $`\{\eta _i\}`$ as basis of $`H^0(\mathrm{\Theta }_DN_D)`$ and $`_{i=1}^8H^1(N_{Y_i/M})`$ respectively.
By using these bases, we compute the matrix representing the connecting homomorphism $`\delta `$.
For that purpose, let us lift 0-cocycle $`\theta _1`$ to 0-cochains of $`\mathrm{\Theta }_S(\mathrm{log}D)N_D`$ as
$$\stackrel{~}{\theta _1^1}=\frac{a_1+y_1}{x_1}\frac{}{y_1}\mathrm{on}U_1,\stackrel{~}{\theta _1^2}=\frac{a_1+y_2}{x_2}\frac{}{y_2}\mathrm{on}U_2,\stackrel{~}{\theta _1^3}=\frac{1+a_1x_3}{y_3}\frac{}{x_3}\mathrm{on}U_3,$$
$$\stackrel{~}{\theta _1^i}=0\mathrm{on}U_i(i=4,5,\mathrm{},16).$$
Other 0-cocycles can be lifted in a similar way.
We first compute $`\delta (\theta _1)`$.
From the definition of $`\delta `$, we have $`\delta (\theta _1)=\{\delta (\theta _1)_{ij}\mathrm{on}U_iU_jD\}`$ with
$$\begin{array}{c}\{\begin{array}{c}\delta (\theta _1)_{12}=(\stackrel{~}{\theta _1^2}\stackrel{~}{\theta _1^1})|_{Y_1}=\left(\frac{a_1+y_2}{x_2}\frac{}{y_2}\frac{a_1+y_1}{x_1}\frac{}{y_1}\right)|_{Y_1}=\frac{a_1}{y_2}\frac{}{x_2}\hfill \\ \delta (\theta _1)_{13}=(\stackrel{~}{\theta _1^3}\stackrel{~}{\theta _1^1})|_{Y_1}=\left(\frac{1+a_1x_3}{y_3}\frac{}{x_3}\frac{a_1+y_1}{x_1}\frac{}{y_1}\right)|_{Y_1}=\frac{1+a_1x_3}{x_3}\frac{}{y_3}\hfill \end{array}\hfill \\ \delta (\theta _1)_{34}=(\stackrel{~}{\theta _1^4}\stackrel{~}{\theta _1^3})|_{Y_2}=\left(0\frac{1+a_1x_3}{y_3}\frac{}{x_3}\right)|_{Y_2}=\frac{1}{x_4}\frac{}{y_4}=\eta _2^{34}\hfill \end{array}$$
Other $`\delta (\theta _1)_{ij}`$’s are zero.
Obviously $`\{\delta (\theta _1)_{34}\}=\{\eta _2^{34}\}=\eta _2`$.
In order to see that $`\{\delta (\theta _1)_{12},\delta (\theta _1)_{13}\}=2\eta _1`$, we set $`\tau =\{\tau _1=\frac{}{x_1},\tau _2=\frac{}{x_2},\tau _3=a_1\frac{}{y_3}\}C^0(N_{Y_1/S})`$.
$$\{\begin{array}{c}\delta (\theta _1)_{12}+2\eta _1^{12}=\frac{a_1}{y_2}\frac{}{x_2}=\tau _2\tau _1\hfill \\ \delta (\theta _1)_{13}+2\eta _1^{13}=\frac{1+a_1x_3}{x_3}\frac{}{y_3}=\tau _3\tau _1\hfill \end{array}$$
This implies $`\{\delta (\theta _1)_{12},\delta (\theta _1)_{13}\}+2\eta _1=\delta \tau `$. Then we have $`\delta (\theta _1)=2\eta _1+\eta _2`$.
Other $`\delta (\theta _i)`$’s can be treated in a similar way. In what follows, we just list up the results of computations.
$`\delta (\theta _2)=\eta _12\eta _2+\eta _3`$
$$\begin{array}{c}\{\begin{array}{c}\delta (\theta _2)_{12}=\stackrel{~}{\theta _2^2}\stackrel{~}{\theta _2^1}|_{Y_1}=00=\eta _1^{12}\hfill \\ \delta (\theta _2)_{13}=\stackrel{~}{\theta _2^3}\stackrel{~}{\theta _2^1}|_{Y_1}=\frac{1}{x_3}\frac{}{y_3}0=\eta _1^{13}\hfill \end{array}\hfill \\ \delta (\theta _2)_{34}=\stackrel{~}{\theta _2^4}\stackrel{~}{\theta _2^3}|_{Y_2}=\frac{1}{y_4}\frac{}{x_4}\frac{1}{x_3}\frac{}{y_3}=2\frac{1}{x_4}\frac{}{y_4}=2\eta _2^{34}\hfill \\ \delta (\theta _2)_{45}=\stackrel{~}{\theta _2^5}\stackrel{~}{\theta _2^4}|_{Y_3}=0\left(\frac{1}{y_4}\frac{}{x_4}\right)=\frac{1}{x_5}\frac{}{y_5}=\eta _3^{45}\hfill \end{array}$$
$`\delta (\theta _3)=\eta _22\eta _3+\eta _4`$
$$\begin{array}{c}\delta (\theta _3)_{34}=\stackrel{~}{\theta _3^4}\stackrel{~}{\theta _3^3}|_{Y_2}=\frac{1}{x_4}\frac{}{y_4}0=\frac{1}{x_4}\frac{}{y_4}=\eta _2^{34}\hfill \\ \delta (\theta _3)_{45}=\stackrel{~}{\theta _3^5}\stackrel{~}{\theta _3^4}|_{Y_3}=\frac{1}{y_5}\frac{}{x_5}\frac{1}{x_4}\frac{}{y_4}=2\frac{1}{x_5}\frac{}{y_5}=2\eta _3^{45}\hfill \\ \{\begin{array}{c}\delta (\theta _3)_{56}=\stackrel{~}{\theta _3^6}\stackrel{~}{\theta _3^5}|_{Y_4}=0\left(\frac{1}{y_5}\frac{}{x_5}\right)=\frac{}{y_6}\hfill \\ \delta (\theta _3)_{58}=\stackrel{~}{\theta _3^8}\stackrel{~}{\theta _3^5}|_{Y_4}=0\left(\frac{1}{y_5}\frac{}{x_5}\right)=\frac{1}{x_8}\frac{}{y_8}\hfill \end{array}\hfill \end{array}$$
Set $`\{\tau _5=0,\tau _6=\frac{}{y_6},\tau _8=0\}C^0(N_{Y_4/S})`$.
$$\{\begin{array}{c}\delta (\theta _3)_{56}\eta _4^{56}=\frac{}{y_6}0=\tau _6\tau _5\hfill \\ \delta (\theta _3)_{58}\eta _4^{58}=\frac{1}{x_8}\frac{}{y_8}\frac{1}{x_8}\frac{}{y_8}=0=\tau _8\tau _5\hfill \end{array}$$
Thus we have $`\{\delta (\theta _3)_{56},\delta (\theta _3)_{58}\}=\eta _4`$.
$`\delta (\theta _4)=\eta _32\eta _4+\eta _5+\eta _6`$
$$\begin{array}{c}\delta (\theta _4)_{45}=\stackrel{~}{\theta _4^5}\stackrel{~}{\theta _4^4}|_{Y_3}=\frac{1y_5}{x_5}\frac{}{y_5}0=\frac{1}{x_5}\frac{}{y_5}=\eta _3^{45}\hfill \\ \{\begin{array}{c}\delta (\theta _4)_{56}=\stackrel{~}{\theta _4^6}\stackrel{~}{\theta _4^5}|_{Y_4}=\frac{1x_6}{y_6}\frac{}{x_6}\frac{1y_5}{x_5}\frac{}{y_5}=\frac{1x_6}{x_6}\frac{}{y_6}\hfill \\ \delta (\theta _4)_{58}=\stackrel{~}{\theta _4^8}\stackrel{~}{\theta _4^5}|_{Y_4}=\frac{1}{y_8}\frac{}{x_8}\frac{1y_5}{x_5}\frac{}{y_5}=\frac{1+2x_8}{x_8(1+x_8)}\frac{}{y_8}\hfill \end{array}\hfill \\ \delta (\theta _4)_{67}=\stackrel{~}{\theta _4^7}\stackrel{~}{\theta _4^6}|_{Y_5}=0\frac{1x_6}{y_6}\frac{}{x_6}=\frac{1}{x_7}\frac{}{y_7}=\eta _5^{67}\hfill \\ \delta (\theta _4)_{89}=\stackrel{~}{\theta _4^9}\stackrel{~}{\theta _4^8}|_{Y_6}=0\left(\frac{1}{y_8}\frac{}{x_8}\right)=\frac{1}{x_9}\frac{}{y_9}=\eta _6^{89}\hfill \end{array}$$
Set $`\{\tau _5=\frac{}{x_5},\tau _6=\frac{}{y_6},\tau _8=0\}C^0(N_{Y_4/S})`$.
$$\{\begin{array}{c}\delta (\theta _4)_{56}+2\eta _4^{56}=\frac{1x_6}{x_6}\frac{}{y_6}=\tau _6\tau _5\hfill \\ \delta (\theta _4)_{58}+2\eta _4^{56}=\frac{1}{x_8(1+x_8)}\frac{}{y_8}=\tau _8\tau _5\hfill \end{array}$$
Thus we have $`\{\delta (\theta _4)_{56},\delta (\theta _4)_{58}\}=2\eta _4`$.
$`\delta (\theta _5)=\eta _42\eta _5`$
$$\begin{array}{c}\{\begin{array}{c}\delta (\theta _5)_{56}=\stackrel{~}{\theta _5^6}\stackrel{~}{\theta _5^5}|_{Y_4}=\frac{1}{x_6}\frac{}{y_6}0=\frac{1}{x_6}\frac{}{y_6}\hfill \\ \delta (\theta _5)_{58}=\stackrel{~}{\theta _5^8}\stackrel{~}{\theta _5^5}|_{Y_4}=00=0\hfill \end{array}\hfill \\ \delta (\theta _5)_{67}=\stackrel{~}{\theta _5^7}\stackrel{~}{\theta _5^6}|_{Y_5}=\frac{1}{y_7}\frac{}{x_7}\left(\frac{1}{x_6}\frac{}{y_6}\right)=2\frac{1}{x_7}\frac{}{y_7}=2\eta _5^{67}\hfill \end{array}$$
Set $`\{\tau _5=\frac{}{x_5},\tau _6=0,\tau _8=0\}C^0(N_{Y_4/S})`$.
$$\{\begin{array}{c}\delta (\theta _5)_{56}\eta _4^{56}=\frac{1}{x_6}\frac{}{y_6}=\tau _6\tau _5\hfill \\ \delta (\theta _5)_{58}\eta _4^{58}=\frac{1}{x_8(1+x_8)}\frac{}{y_8}=\tau _8\tau _5\hfill \end{array}$$
Thus we have $`\{\delta (\theta _5)_{56},\delta (\theta _5)_{58}\}=\eta _4`$.
$`\delta (\theta _6)=\eta _42\eta _6+\eta _7`$
$$\begin{array}{c}\{\begin{array}{c}\delta (\theta _6)_{56}=\stackrel{~}{\theta _6^6}\stackrel{~}{\theta _6^5}|_{Y_4}=00=\eta _4^{56}\hfill \\ \delta (\theta _6)_{58}=\stackrel{~}{\theta _6^8}\stackrel{~}{\theta _6^5}|_{Y_4}=\frac{1}{x_8}\frac{}{y_8}0=\eta _4^{58}\hfill \end{array}\hfill \\ \delta (\theta _6)_{89}=\stackrel{~}{\theta _6^9}\stackrel{~}{\theta _6^8|_{Y_6}}=\frac{1}{y_9}\frac{}{x_9}\frac{1}{x_8}\frac{}{y_8}=2\frac{1}{x_9}\frac{}{y_9}=2\eta _6^{89}\hfill \\ \{\begin{array}{c}\delta (\theta _6)_{\mathrm{9\hspace{0.17em}10}}=\stackrel{~}{\theta _6^{10}}\stackrel{~}{\theta _6^9}|_{Y_7}=0\left(\frac{1}{y_9}\frac{}{x_9}\right)=\frac{}{y_{10}}\hfill \\ \delta (\theta _6)_{\mathrm{9\hspace{0.17em}11}}=\stackrel{~}{\theta _6^{11}}\stackrel{~}{\theta _6^9}|_{Y_7}=0\left(\frac{1}{y_9}\frac{}{x_9}\right)=\frac{1}{x_{11}}\frac{}{y_{11}}\hfill \end{array}\hfill \end{array}$$
Set $`\{\tau _9=0,\tau _{10}=\frac{}{y_{10}},\tau _{11}=0\}C^0(N_{Y_7/S})`$.
$$\{\begin{array}{c}\delta (\theta _6)_{\mathrm{9\hspace{0.17em}10}}\eta _7^{\mathrm{9\hspace{0.17em}10}}=\frac{}{y_{10}}=\tau _{10}\tau _9\hfill \\ \delta (\theta _6)_{\mathrm{9\hspace{0.17em}11}}\eta _7^{\mathrm{9\hspace{0.17em}11}}=0=\tau _{11}\tau _9\hfill \end{array}$$
Hence we have $`\{\delta (\theta _6)_{\mathrm{9\hspace{0.17em}10}},\delta (\theta _6)_{\mathrm{9\hspace{0.17em}11}}\}=\eta _7`$.
$`\delta (\theta _7)=\eta _62\eta _7+\eta _8`$
$$\begin{array}{c}\delta (\theta _7)_{89}=\stackrel{~}{\theta _7^9}\stackrel{~}{\theta _7^8}|_{Y_6}=\frac{1+ty_9}{x_9}\frac{}{y_9}0=\frac{1}{x_9}\frac{}{y_9}=\eta _6^{89}\hfill \\ \{\begin{array}{c}\delta (\theta _7)_{\mathrm{9\hspace{0.17em}10}}=\stackrel{~}{\theta _7^{10}}\stackrel{~}{\theta _7^9}|_{Y_7}=\frac{t+x_{10}}{y_{10}}\frac{}{x_{10}}\frac{1+ty_9}{x_9}\frac{}{y_9}=\frac{t+x_{10}}{x_{10}}\frac{}{y_{10}}\hfill \\ \delta (\theta _7)_{\mathrm{9\hspace{0.17em}11}}=\stackrel{~}{\theta _7^{11}}\stackrel{~}{\theta _7^9}|_{Y_7}=\frac{1}{y_{11}}\frac{}{x_{11}}\frac{1+ty_9}{x_9}\frac{}{y_9}=\frac{t2x_{11}}{x_{11}(t+x_{11})}\frac{}{y_{11}}\hfill \end{array}\hfill \\ \{\begin{array}{c}\delta (\theta _7)_{\mathrm{11\hspace{0.17em}12}}=\stackrel{~}{\theta _7^{12}}\stackrel{~}{\theta _7^{11}}|_{Y_8}=0\left(\frac{1}{y_{11}}\frac{}{x_{11}}\right)=\frac{}{y_{12}}=\eta _{\mathrm{11\hspace{0.17em}12}}\hfill \\ \delta (\theta _7)_{\mathrm{11\hspace{0.17em}13}}=\stackrel{~}{\theta _7^{13}}\stackrel{~}{\theta _7^{11}}|_{Y_8}=0\left(\frac{1}{y_{11}}\frac{}{x_{11}}\right)=\frac{1}{x_{13}}\frac{}{y_{13}}=\eta _8^{\mathrm{11\hspace{0.17em}13}}\hfill \end{array}\hfill \end{array}$$
Set $`\{\tau _9=t\frac{}{x_9},\tau _{10}=\frac{}{y_{10}},\tau _{11}=0\}C^0(N_{Y_7/S})`$.
$$\{\begin{array}{c}\delta (\theta _7)_{\mathrm{9\hspace{0.17em}10}}+2\eta _7^{\mathrm{9\hspace{0.17em}10}}=\frac{t+x_{10}}{x_{10}}\frac{}{y_{10}}=\tau _{10}\tau _9\hfill \\ \delta (\theta _7)_{\mathrm{9\hspace{0.17em}11}}+2\eta _7^{\mathrm{9\hspace{0.17em}11}}=\frac{t}{x_{11}(tx_{11})}\frac{}{y_{11}}=\tau _{11}\tau _9\hfill \end{array}$$
Hence we have $`\{\delta (\theta _7)_{\mathrm{9\hspace{0.17em}10}},\delta (\theta _7)_{\mathrm{9\hspace{0.17em}11}}\}=2\eta _7`$.
$`\delta (\theta _8)=\eta _72\eta _8`$
$$\begin{array}{c}\{\begin{array}{c}\delta (\theta _8)_{\mathrm{9\hspace{0.17em}10}}=\stackrel{~}{\theta _8^{10}}\stackrel{~}{\theta _8^9}|_{Y_7}=00=\eta _7^{\mathrm{9\hspace{0.17em}10}}\hfill \\ \delta (\theta _8)_{\mathrm{9\hspace{0.17em}11}}=\stackrel{~}{\theta _8^{11}}\stackrel{~}{\theta _8^9}|_{Y_7}=\frac{1+a_0y_{11}}{x_{11}}\frac{}{y_{11}}0=\eta _7^{\mathrm{9\hspace{0.17em}11}}\hfill \end{array}\hfill \\ \{\begin{array}{c}\delta (\theta _8)_{\mathrm{11\hspace{0.17em}12}}=\stackrel{~}{\theta _8^{12}}\stackrel{~}{\theta _8^{11}}|_{Y_8}=\frac{a_0+x_{12}}{y_{12}}\frac{}{x_{12}}\frac{1+a_0y_{11}}{x_{11}}\frac{}{y_{11}}=\frac{a_0+x_{12}}{x_{12}}\frac{}{y_{12}}\hfill \\ \delta (\theta _8)_{\mathrm{11\hspace{0.17em}13}}=\stackrel{~}{\theta _8^{13}}\stackrel{~}{\theta _8^{11}}|_{Y_8}=\frac{1}{y_{13}}\frac{}{x_{13}}\frac{1+a_0y_{11}}{x_{11}}\frac{}{y_{11}}=\frac{a_0+2x_{13}}{x_{13}(a_0x_{13})}\frac{}{y_{13}}\hfill \end{array}\hfill \end{array}$$
Set $`\{\tau _{11}=a_0\frac{}{x_{11}},\tau _{12}=\frac{}{y_{12}},\tau _{13}=0\}C^0(N_{Y_8/S})`$.
$$\{\begin{array}{c}\delta (\theta _8)_{\mathrm{11\hspace{0.17em}12}}+2\eta _{\mathrm{11\hspace{0.17em}12}}=\frac{a_0+x_{12}}{x_{12}}\frac{}{y_{12}}=\tau _{12}\tau _{11}\hfill \\ \delta (\theta _8)_{\mathrm{11\hspace{0.17em}13}}+2\eta _8^{\mathrm{11\hspace{0.17em}13}}=\frac{a_0}{x_{13}(a_0x_{13})}\frac{}{y_{13}}=\tau _{13}\tau _{11}\hfill \end{array}$$
Hence we have $`\{\delta (\theta _8)_{\mathrm{11\hspace{0.17em}12}},\delta (\theta _8)_{\mathrm{11\hspace{0.17em}13}}\}=2\eta _8`$.
Summing up all the computations, we see that the matrix of the linear map $`\delta `$ is given by
$$\left(\begin{array}{cccccccc}2& 1& 0& 0& 0& 0& 0& 0\\ 1& 2& 1& 0& 0& 0& 0& 0\\ 0& 1& 2& 1& 0& 0& 0& 0\\ 0& 0& 1& 2& 1& 1& 0& 0\\ 0& 0& 0& 1& 2& 0& 0& 0\\ 0& 0& 0& 1& 0& 2& 1& 0\\ 0& 0& 0& 0& 0& 1& 2& 1\\ 0& 0& 0& 0& 0& 0& 1& 2\end{array}\right).$$
Since it is well-known that this matrix coincides with the intersection matrix $`((Y_i,Y_j))_{1i,jr}`$ of type $`\stackrel{~}{E_8}`$, this completes the proof of Proposition 2.1. ∎
## 4. Local cohomology of generalized Okamoto–Painlevé pair of multiplicative type
Let $`(S,Y)`$ be a generalized Okamoto-Painlevé pair as in §2.
For a pair $`(S,Y)`$ of additive type, the result of §2 shows the existence of differential equations on $`SY`$ (cf. \[STT\]). Even for $`(S,Y)`$ of multiplicative type, if $`dimH_D^1(\mathrm{\Theta }_S(\mathrm{log}D))1`$, we can derive a differential equation in the same way as in \[STT\]. Unfortunately, we can prove $`H_D^1(\mathrm{\Theta }_S(\mathrm{log}D))=\{0\}`$ for a pair $`(S,Y)`$ of $`\stackrel{~}{A_8}`$. (For other multiplicative types, we expect that $`H_D^1(\mathrm{\Theta }_S(\mathrm{log}D))=\{0\}`$.) This means that, for a pair $`(S,Y)`$ of multiplicative type, there is no differential equation on $`SY`$ in the sense as in \[STT\].
In this section, we will calculate the local cohomology group $`H_D^1(\mathrm{\Theta }_S(\mathrm{log}D))`$ of pair $`(S,Y)`$ of $`\stackrel{~}{A_8}`$ type.
$``$ Construction of $`(S,Y)`$ of $`\stackrel{~}{A_8}`$ type
Now, we consider $`(S,Y)`$ of $`\stackrel{~}{A_8}`$ type as an example of multiplicative type.
According to \[Sakai\], any $`\stackrel{~}{A_8}`$-surface is obtained by blowing up $`^2`$ at the following 9 points given by
$$\begin{array}{c}p_1:(1:0:0)p_2:(\frac{z}{x},\frac{y}{z})=(0,0)p_6:(\frac{z}{x},\frac{xy}{z^2})=(0,1),\hfill \\ p_3:(0:0:1)p_4:(\frac{y}{z},\frac{x}{y})=(0,0)p_5:(\frac{y}{z},\frac{zx}{y^2})=(0,b),\hfill \\ p_7:(0:1:0)p_8:(\frac{x}{y},\frac{z}{x})=(0,0)p_9:(\frac{x}{y},\frac{yz}{x^2})=(0,c).\hfill \end{array}$$
Moreover there exists a $`^\times `$-action on the family of surfaces by
$$(\mu ,(b,c,[x:y:z]))(\mu ^3b,\mu ^3c,[\mu x:\mu ^1y:z]).$$
By putting $`t=bc`$, we can normalize this description.
We can choose the following coordinate system of $`\stackrel{~}{A_8}`$-surfaces $`S_t`$ parameterized by t.
$`U_i`$ $`=`$ $`Spec[x_i,y_i]^2(i=0,1,\mathrm{},8).`$
$`U_9`$ $`=`$ $`Spec[x_9,y_9,{\displaystyle \frac{1}{t+x_9}}]^2\{t+x_9=0\}.`$
$`U_i`$ $`=`$ $`Spec[x_i,y_i,{\displaystyle \frac{1}{1+x_i}}]^2\{1+x_i=0\}(i=10,11).`$
$`U_{12}`$ $`=`$ $`Spec[u_{12},v_{12},{\displaystyle \frac{1}{t+u_{12}v_{12}}}]^2\{t+u_{12}v_{12}=0\}.`$
$`U_i`$ $`=`$ $`Spec[u_i,v_i,{\displaystyle \frac{1}{1+u_iv_i}}]^2\{1+u_iv_i=0\}(i=13,14).`$
where $`t^\times `$.
$$\begin{array}{cc}Y_1=\{x_0=0,y_1=0\},\hfill & Y_2=\{x_1=0,y_2=0,y_9=0\},\hfill \\ Y_3=\{x_2=0,y_3=0\},\hfill & Y_4=\{x_3=0,y_4=0\},\hfill \\ Y_5=\{x_4=0,y_5=0,y_{10}=0\},\hfill & Y_6=\{x_5=0,y_6=0\},\hfill \\ Y_7=\{x_6=0,y_7=0\},\hfill & Y_8=\{x_7=0,y_8=0,y_{11}=0\},\hfill \\ Y_9=\{x_8=0,y_0=0\},\hfill & \end{array}$$
$$S_t=\underset{i=0}{\overset{14}{}}U_i,$$
$$Y_t=\underset{i=1}{\overset{9}{}}Y_i,D=Y,$$
$$S_tY_t=U_{12}U_{13}U_{14}.$$
For later use, we need only the coordinate transformations near each component $`Y_i`$. Here we will list up the coordinate transformations only for a neighborhood of each $`Y_i`$.
$$\begin{array}{ccc}Y_1:\{\begin{array}{c}x_0=x_{1}^{}{}_{}{}^{2}y_1\hfill \\ y_0=\frac{1}{x_1}\hfill \end{array}\hfill & Y_2:\{\begin{array}{c}x_1=x_2y_2\hfill \\ y_1=\frac{1}{x_2}\hfill \end{array}\{\begin{array}{c}x_1=x_9(t+x_9)y_9\hfill \\ y_1=\frac{1}{(t+x_9)}\hfill \end{array}\hfill & Y_3:\{\begin{array}{c}x_2=x_{3}^{}{}_{}{}^{2}y_3\hfill \\ y_2=\frac{1}{x_3}\hfill \end{array}\hfill \\ Y_4:\{\begin{array}{c}x_3=x_{4}^{}{}_{}{}^{2}y_4\hfill \\ y_3=\frac{1}{x_4}\hfill \end{array}\hfill & Y_5:\{\begin{array}{c}x_4=x_5y_5\hfill \\ y_4=\frac{1}{x_5}\hfill \end{array}\{\begin{array}{c}x_4=x_{10}(1+x_{10})y_{10}\hfill \\ y_4=\frac{1}{(1+x_{10})}\hfill \end{array}\hfill & Y_6:\{\begin{array}{c}x_5=x_{6}^{}{}_{}{}^{2}y_6\hfill \\ y_5=\frac{1}{x_6}\hfill \end{array}\hfill \\ Y_7:\{\begin{array}{c}x_6=x_{7}^{}{}_{}{}^{2}y_7\hfill \\ y_6=\frac{1}{x_7}\hfill \end{array}\hfill & Y_8:\{\begin{array}{c}x_7=x_8y_8\hfill \\ y_7=\frac{1}{x_8}\hfill \end{array}\{\begin{array}{c}x_7=x_{11}(1+x_{11})y_{11}\hfill \\ y_7=\frac{1}{(1+x_{11})}\hfill \end{array}\hfill & Y_9:\{\begin{array}{c}x_8=x_{0}^{}{}_{}{}^{2}y_0\hfill \\ y_8=\frac{1}{x_0}.\hfill \end{array}\hfill \end{array}$$
This gives a generalized Okamoto–Painlevé pair $`(S_t,Y_t)`$ of type $`\stackrel{~}{A}_8`$.
###### Proposition 4.1.
Let $`(S_t,Y_t)`$ be as above and set $`D_t=(Y_t)_{red}`$.
If $`t`$ is not a root of unity, then we have
$$H_{D_t}^1(\mathrm{\Theta }_{S_t}(\mathrm{log}D_t))=0.$$
###### Remark 4.1.
We expect that, if $`t`$ is a root of unity then $`(S_t,Y_t)`$ is of fibered type.
Proof. From the diagram
$$\begin{array}{ccccccc}& & & 0& & 0& \\ & & & & & & \\ 0& \mathrm{\Theta }_S(\mathrm{log}D)& & \mathrm{\Theta }_S(\mathrm{log}D)𝒪_S((n1)D)& & \mathrm{\Theta }_S(\mathrm{log}D)N_{(n1)D}& 0\\ & & & & & \mu & \\ 0& \mathrm{\Theta }_S(\mathrm{log}D)& & \mathrm{\Theta }_S(\mathrm{log}D)𝒪_S(nD)& & \mathrm{\Theta }_S(\mathrm{log}D)N_{nD}& 0\\ & & & & & & \\ & & & \mathrm{\Theta }_S(\mathrm{log}D)N_D^n& & coker\mu & \\ & & & & & & \\ & & & 0& & 0& ,\end{array}$$
we obtain the exact sequence
$$0\mathrm{\Theta }_S(\mathrm{log}D)N_{(n1)D}\mathrm{\Theta }_S(\mathrm{log}D)N_{nD}\mathrm{\Theta }_S(\mathrm{log}D)N_D^n0.$$
Therefore, we have the following sequence for each $`n2`$.
$$0H^0(\mathrm{\Theta }_S(\mathrm{log}D)N_{(n1)D})H^0(\mathrm{\Theta }_S(\mathrm{log}D)N_{nD})H^0(\mathrm{\Theta }_S(\mathrm{log}D)N_D^n).$$
If $`H^0(\mathrm{\Theta }_S(\mathrm{log}D)N_D^n)=0`$ for any $`n1`$, then we have
$$H^0(\mathrm{\Theta }_S(\mathrm{log}D)N_{(n1)D})H^0(\mathrm{\Theta }_S(\mathrm{log}D)N_{nD}).$$
Since $`H^0(\mathrm{\Theta }_S(\mathrm{log}D)N_D)=0`$, we have
$$H^0(\mathrm{\Theta }_S(\mathrm{log}D)N_{nD})=0(n1).$$
Therefore, noting that (cf. §2)
$$H_D^1(S,\mathrm{\Theta }_S(\mathrm{log}D))=\underset{}{\mathrm{lim}}H^0(\mathrm{\Theta }_S(\mathrm{log}D)N_{nD}),$$
we obtain
$$H_D^1(S,\mathrm{\Theta }_S(\mathrm{log}D))=0.$$
Now let us calculate $`H^0(\mathrm{\Theta }_S(\mathrm{log}D)N_D^n)`$.
Recall that we have the exact sequence
$$0\nu _{}(_{i=1}^9N_{Y_i})\mathrm{\Theta }_S(\mathrm{log}D)N_D\mathrm{\Theta }_DN_D0,$$
where $`\nu :\stackrel{~}{D}=_{i=1}^9Y_iD`$ is the normalization map (Lemma 2.1 (3)). Tensoring this sequence with $`N_D^{n1}`$, we obtain
$$0\nu _{}(_{i=1}^9N_{Y_i})N_D^{n1}\mathrm{\Theta }_S(\mathrm{log}D)N_D^n\mathrm{\Theta }_DN_D^n0.$$
Therefore we have
$$\begin{array}{c}H^0(D,\nu _{}(_{i=1}^9N_{Y_i})N_D^{n1})H^0(D,\mathrm{\Theta }_S(\mathrm{log}D)N_D^n)H^0(D,\mathrm{\Theta }_DN_D^n)\hfill \\ H^1(D,\nu _{}(_{i=1}^9N_{Y_i})N_D^{n1}).\hfill \end{array}$$
In $`\stackrel{~}{A_8}`$ case, we see that $`t_i=2`$, where $`t_i`$ is the number of intersections of $`Y_i`$ with other components. From this, we obtain
$$\mathrm{\Theta }_D\nu _{}(_{i=1}^9\mathrm{\Theta }_{Y_i}(t_i))\nu _{}(_{i=1}^9𝒪_{Y_i}(2t_i))\nu _{}(_{i=1}^9𝒪_{Y_i}).$$
Since $`\nu `$ is a finite morphism, we have
$$\begin{array}{c}H^i(D,\nu _{}(_{i=1}^9N_{Y_i})N_D^{n1})H^i(\stackrel{~}{D},(_{i=1}^9N_{Y_i})\nu ^{}(N_D^{n1}))(i=0,1),\hfill \\ H^0(D,\mathrm{\Theta }_DN_D^n)H^0(D,\nu _{}(_{i=1}^9𝒪_{Y_i})N_D^n)H^0(\stackrel{~}{D},(_{i=1}^9𝒪_{Y_i})\nu ^{}(N_D^n)).\hfill \end{array}$$
Note that $`DY_i=2t_i=0`$, we see that $`N_D|_{Y_i}𝒪_{Y_i}`$. Therefore we have
$$(_{i=1}^9N_{Y_i})\nu ^{}(N_D^{n1})_{i=1}^9N_{Y_i},(_{i=1}^9𝒪_{Y_i})\nu ^{}(N_D^n)_{i=1}^9𝒪_{Y_i}.$$
Summing up the arguments above, we obtain
$$\begin{array}{ccccccc}0& & H^0(\mathrm{\Theta }_S(\mathrm{log}D)N_D^n)& & _{i=1}^9H^0(Y_i,𝒪_{Y_i})& \stackrel{\delta }{}& _{i=1}^9H^1(N_{Y_i})\\ & & & & & & \\ & & & & ^9& \stackrel{\delta }{}& ^9.\end{array}$$
We will analyse the edge homomorphism $`\delta `$ by using the Čech cocycles.
Noting that
(6)
$$H^0(\mathrm{\Theta }_DN_D^n)H^0(\mathrm{\Theta }_{Y_i}(2)N_D^n)^9.$$
For each $`i(1i9)`$, we introduce a generator $`\theta _i`$ of the cohomology group in (6) as follows.
| $`\theta _1`$ | $`\left\{\theta _1^0=\frac{y_0}{x_0^ny_0^n}\frac{}{y_0}\text{on}U_0Y_1,\theta _1^1=\frac{x_1}{x_1^ny_1^n}\frac{}{x_1}\text{on}U_1Y_1\right\}`$ |
| --- | --- |
| $`\theta _2`$ | $`\left\{\begin{array}{c}\theta _2^1=\frac{(1ty_1)^ny_1}{x_1^ny_1^n}\frac{}{y_1}\text{on}U_1Y_2,\theta _2^2=\frac{(x_2t)^nx_2}{x_2^ny_2^n}\frac{}{x_2}\text{on}U_2Y_2,\\ \theta _2^9=\frac{1}{(t+x_9)^{n1}y_9^n}\frac{}{x_9}\text{on}U_9Y_2\end{array}\right\}`$ |
| $`\theta _3`$ | $`\left\{\theta _3^2=\frac{y_2}{x_2^ny_2^n}\frac{}{y_2}\text{on}U_2Y_3,\theta _3^3=\frac{x_3}{x_3^ny_3^n}\frac{}{x_3}\text{on}U_3Y_3\right\}`$ |
| $`\theta _4`$ | $`\left\{\theta _4^3=\frac{y_3}{x_3^ny_3^n}\frac{}{y_3}\text{on}U_3Y_4,\theta _4^4=\frac{x_4}{x_4^ny_4^n}\frac{}{x_4}\text{on}U_4Y_4\right\}`$ |
| $`\theta _5`$ | $`\left\{\begin{array}{c}\theta _5^4=\frac{(1x_4)^ny_4}{x_4^ny_4^n}\frac{}{y_4}\text{on}U_4Y_5,\theta _5^5=\frac{(x_51)^nx_5}{x_5^ny_5^n}\frac{}{x_5}\text{on}U_5Y_5,\\ \theta _5^{10}=\frac{1}{(1+x_{10})^{n1}y_{10}^n}\frac{}{x_{10}}\text{on}U_{10}Y_5\end{array}\right\}`$ |
| $`\theta _6`$ | $`\left\{\theta _6^5=(1)^n\frac{y_5}{x_5^ny_5^n}\frac{}{y_5}\text{on}U_5Y_6,\theta _6^6=(1)^n\frac{x_6}{x_6^ny_6^n}\frac{}{x_6}\text{on}U_6Y_6\right\}`$ |
| $`\theta _7`$ | $`\left\{\theta _7^6=(1)^n\frac{y_6}{x_6^ny_6^n}\frac{}{y_6}\text{on}U_6Y_7,\theta _7^7=(1)^n\frac{x_7}{x_7^ny_7^n}\frac{}{x_7}\text{on}U_7Y_7\right\}`$ |
| $`\theta _8`$ | $`\left\{\begin{array}{c}\theta _8^7=(1)^n\frac{(1x_7)^ny_7}{x_7^ny_7^n}\frac{}{y_7}\text{on}U_7Y_8,\theta _8^8=(1)^n\frac{(x_81)^nx_8}{x_8^ny_8^n}\frac{}{x_8}\text{on}U_8Y_8,\\ \theta _8^{11}=(1)^n\frac{1}{(1+x_{11})^{n1}y_{11}^n}\frac{}{x_{11}}\text{on}U_{11}Y_8\end{array}\right\}`$ |
| $`\theta _9`$ | $`\left\{\theta _9^8=\frac{y_8}{x_8^ny_8^n}\frac{}{y_8}\text{on}U_8Y_9,\theta _9^0=\frac{x_0}{x_0^ny_0^n}\frac{}{x_0}\text{on}U_0Y_9\right\}`$ |
On the other hand, for each $`i(1i9)`$, we have a generator $`\eta _iH^1(Y_i,N_{Y_i/S}N_D^{n1})`$ as follows.
| $`\eta _1`$ | $`\left\{\eta _1^{01}=\frac{1}{x_1x_1^{n1}y_1^{n1}}\frac{}{y_1}\text{on}U_0U_1Y_1\right\}`$ |
| --- | --- |
| $`\eta _2`$ | $`\left\{\eta _2^{12}=\frac{1}{y_2^{n1}}\frac{}{y_2}\text{on}U_1U_2Y_2,\eta _2^{29}=0\text{on}U_2U_9Y_2\right\}`$ |
| $`\eta _3`$ | $`\left\{\eta _3^{23}=\frac{1}{x_3x_3^{n1}y_3^{n1}}\frac{}{y_3}\text{on}U_2U_3Y_3\right\}`$ |
| $`\eta _4`$ | $`\left\{\eta _4^{34}=\frac{1}{x_4x_4^{n1}y_4^{n1}}\frac{}{y_4}\text{on}U_3U_4Y_4\right\}`$ |
| $`\eta _5`$ | $`\left\{\eta _5^{45}=\frac{1}{y_5^{n1}}\frac{}{y_5}\text{on}U_4U_5Y_5,\eta _5^{\mathrm{5\hspace{0.17em}10}}=0\text{on}U_5U_{10}Y_5\right\}`$ |
| $`\eta _6`$ | $`\left\{\eta _6^{56}=(1)^n\frac{1}{x_6x_6^{n1}y_6^{n1}}\frac{}{y_6}\text{on}U_5U_6Y_6\right\}`$ |
| $`\eta _7`$ | $`\left\{\eta _7^{67}=(1)^n\frac{1}{x_7x_7^{n1}y_7^{n1}}\frac{}{y_7}\text{on}U_6U_7Y_7\right\}`$ |
| $`\eta _8`$ | $`\left\{\eta _8^{78}=(1)^n\frac{1}{y_8^{n1}}\frac{}{y_8}\text{on}U_7U_8Y_8,\eta _8^{\mathrm{8\hspace{0.17em}11}}=0\text{on}U_8U_{11}Y_8\right\}`$ |
| $`\eta _9`$ | $`\left\{\eta _9^{80}=\frac{1}{x_0x_0^{n1}y_0^{n1}}\frac{}{y_0}\text{on}U_8U_0Y_9\right\}`$ |
We take $`\{\theta _i\}`$ and $`\{\eta _i\}`$ as basis of $`H^0(\mathrm{\Theta }_DN_D^n)`$ and $`_{i=1}^9H^1(N_{Y_i/S}N_D^{n1})`$ respectively. By using these bases, we compute the matrix representing the connecting homomorphisn $`\delta `$.
For that purpose, let us lift 0-cocycle $`\theta _1`$ to 0-cochains of $`\mathrm{\Theta }_S(\mathrm{log}D)N_D`$ as
$$\stackrel{~}{\theta _1^0}=\frac{y_0}{x_0^ny_0^n}\frac{}{y_0}\mathrm{on}U_0,\stackrel{~}{\theta _1^1}=\frac{x_1}{x_1^ny_1^n}\frac{}{x_1}\mathrm{on}U_1,$$
$$\stackrel{~}{\theta _1^i}=0\mathrm{on}U_i(i=2,3,\mathrm{},14).$$
Other 0-cocycles can be lifted in a similar way.
We first compute $`\delta (\theta _1)`$.
From the definition of $`\delta `$, we have $`\delta (\theta _1)=\{\delta (\theta _1)_{ij}\mathrm{on}U_iU_jD\}`$ with
$$\begin{array}{c}\delta (\theta _1)_{80}=\stackrel{~}{\theta _1^8}+\stackrel{~}{\theta _1^0}|_{Y_9}=0+\frac{y_0}{x_0^ny_0^n}\frac{}{y_0}=\eta _9^{80}\hfill \\ \delta (\theta _1)_{01}=\stackrel{~}{\theta _1^0}+\stackrel{~}{\theta _1^1}|_{Y_1}=\frac{y_0}{x_0^ny_0^n}\frac{}{y_0}+\left(\frac{x_1}{x_1^ny_1^n}\frac{}{x_1}\right)=2\frac{1}{x_1x_1^{n1}y_1^{n1}}\frac{}{y_1}=2\eta _1^{01}\hfill \\ \{\begin{array}{c}\delta (\theta _1)_{12}=\stackrel{~}{\theta _1^1}+\stackrel{~}{\theta _1^2}|_{Y_2}=\left(\frac{x_1}{x_1^ny_1^n}\frac{}{x_1}\right)+0=\frac{1}{y_2^{n1}}\frac{}{y_2}=\eta _2^{12}\hfill \\ \delta (\theta _1)_{29}=\stackrel{~}{\theta _1^2}+\stackrel{~}{\theta _1^9}|_{Y_2}=0+0=\eta _2^{29}\hfill \end{array}\hfill \end{array}$$
Other $`\delta (\theta _1)_{ij}`$’s are zero.
Obviously $`\delta (\theta _1)=\eta _92\eta _1+\eta _2`$.
Other $`\delta (\theta _i)`$’s can be treated in a similar way. In what follows, we just list up a few results of computations.
$`\delta (\theta _2)=\eta _12\eta _2+(t)^n\eta _3`$
$$\begin{array}{c}\delta (\theta _2)_{01}=\stackrel{~}{\theta _2^0}+\stackrel{~}{\theta _2^1}|_{Y_1}=0+\frac{(1ty_1)^ny_1}{x_1^ny_1^n}\frac{}{y_1}=\frac{1}{x_1x_1^{n1}y_1^{n1}}\frac{}{y_1}=\eta _1^{01}(y_1=0\mathrm{on}U_1Y_1)\hfill \\ \{\begin{array}{c}\delta (\theta _2)_{12}=\stackrel{~}{\theta _2^1}+\stackrel{~}{\theta _2^2}|_{Y_2}=\frac{(1ty_1)^ny_1}{x_1^ny_1^n}\frac{}{y_1}+\left(\frac{(x_2t)^nx_2}{x_2^ny_2^n}\frac{}{x_2}\right)=\frac{(x_2t)^n}{x_2x_2^{n1}y_2^{n1}}\frac{}{y_2}\hfill \\ \delta (\theta _2)_{29}=\stackrel{~}{\theta _2^2}+\stackrel{~}{\theta _2^9}|_{Y_2}=\left(\frac{(x_2t)^nx_2}{x_2^ny_2^n}\frac{}{x_2}\right)+\left(\frac{1}{(t+x_9)^{n1}y_9^n}\frac{}{x_9}\right)=\frac{1}{x_9(t+x_9)^{n1}y_9^{n1}}\frac{}{y_9}\hfill \end{array}\hfill \\ \delta (\theta _2)_{23}=\stackrel{~}{\theta _2^2}+\stackrel{~}{\theta _2^3}|_{Y_3}=\left(\frac{(x_2t)^nx_2}{x_2^ny_2^n}\frac{}{x_2}\right)\frac{}{x_2}+0=(t)^n\frac{1}{x_3x_3^{n1}y_3^{n1}}\frac{}{y_3}=(t)^n\eta _3^{23}\hfill \end{array}$$
Set $`\{\tau _1=\frac{(1ty_1)^n+(1ty_1)^{n1}2}{y_1x_1^{n1}y_1^{n1}}\frac{}{x_1},\tau _2=\frac{(x_2t)^{n1}}{x_2^{n1}y_2^{n1}}\frac{}{y_2},\tau _9=0\}C^0(N_{Y_2/M}N_D^{n1})`$.
$$\{\begin{array}{c}\delta (\theta _2)_{12}+2\eta _2^{12}=\frac{(x_2t)^n}{x_2x_2^{n1}y_2^{n1}}\frac{}{y_2}+2\frac{1}{y_2^{n1}}\frac{}{y_2}=\tau _1+\tau _2\hfill \\ \delta (\theta _2)_{29}+2\eta _2^{29}=\frac{1}{x_9(t+x_9)^{n1}y_9^{n1}}\frac{}{y_9}0=\tau _2+\tau _9\hfill \end{array}$$
Thus we have $`\{\delta (\theta _2)_{12},\delta (\theta _2)_{29}\}=2\eta _2`$.
$`\delta (\theta _3)=\frac{1}{(t)^n}\eta _22\eta _3+\eta _4`$
$$\begin{array}{c}\{\begin{array}{c}\delta (\theta _3)_{12}=\stackrel{~}{\theta _3^1}+\stackrel{~}{\theta _3^2}|_{Y_2}=0+\frac{y_2}{x_2^ny_2^n}\frac{}{y_2}=\frac{y_2}{x_2^ny_2^n}\frac{}{y_2}\hfill \\ \delta (\theta _3)_{29}=\stackrel{~}{\theta _3^2}+\stackrel{~}{\theta _3^9}|_{Y_2}=\frac{y_2}{x_2^ny_2^n}\frac{}{y_2}+0=\frac{1}{x_9^n(t+x_9)^ny_9^{n1}}\frac{}{y_9}\hfill \end{array}\hfill \\ \delta (\theta _3)_{23}=\stackrel{~}{\theta _3^2}+\stackrel{~}{\theta _3^3}|_{Y_3}=\frac{y_2}{x_2^ny_2^n}\frac{}{y_2}+\left(\frac{x_3}{x_3^ny_3^n}\frac{}{x_3}\right)=2\frac{1}{x_3x_3^{n1}y_3^{n1}}\frac{}{y_3}=2\eta _3^{23}\hfill \\ \delta (\theta _3)_{34}=\stackrel{~}{\theta _3^3}+\stackrel{~}{\theta _3^4}|_{Y_4}=\left(\frac{x_3}{x_3^ny_3^n}\frac{}{x_3}\right)+0=\frac{1}{x_4x_4^{n1}y_4^{n1}}\frac{}{y_4}=\eta _4^{34}\hfill \end{array}$$
Set $`\{\tau _1=\frac{1(1ty_1)^n}{(t)^ny_1x_1^{n1}y_1^{n1}}\frac{}{x_1},\tau _2=\frac{(t)^n(x_2t)^n}{(t)^nx_2x_2^{n1}y_2^{n1}}\frac{}{y_2},\tau _9=\frac{1}{(t)^n(x_9+t)^ny_9^{n1}}\frac{}{y_9}\}C^0(N_{Y_2/M}N_D^{n1})`$.
$$\{\begin{array}{c}\delta (\theta _3)_{12}\frac{1}{(t)^n}\eta _2^{12}=\frac{y_2}{x_2^ny_2^n}\frac{}{y_2}\frac{1}{(t)^n}\frac{1}{y_2^{n1}}\frac{}{y_2}=\tau _1+\tau _2\hfill \\ \delta (\theta _3)_{29}\frac{1}{(t)^n}\eta _2^{29}=\frac{1}{x_9^n(t+x_9)^ny_9^{n1}}\frac{}{y_9}0=\tau _2+\tau _9\hfill \end{array}$$
Thus we have $`\{\delta (\theta _3)_{12},\delta (\theta _3)_{29}\}=\frac{1}{(t)^n}\eta _2`$.
Summing up all the computations, we see that the matrix of the linear map $`\delta `$ is given by
$$\delta =\left(\begin{array}{ccccccccc}2& 1& 0& 0& 0& 0& 0& 0& 1\\ 1& 2& (t)^n& 0& 0& 0& 0& 0& 0\\ 0& (t)^n& 2& 1& 0& 0& 0& 0& 0\\ 0& 0& 1& 2& 1& 0& 0& 0& 0\\ 0& 0& 0& 1& 2& 1& 0& 0& 0\\ 0& 0& 0& 0& 1& 2& 1& 0& 0\\ 0& 0& 0& 0& 0& 1& 2& 1& 0\\ 0& 0& 0& 0& 0& 0& 1& 2& 1\\ 1& 0& 0& 0& 0& 0& 0& 1& 2\end{array}\right)$$
Since
$$det\delta =\frac{((t)^n1)^2}{(t)^n},$$
we see
$$rank\delta =\{\begin{array}{c}9(t\mathrm{is}\mathrm{not}\mathrm{a}\mathrm{root}\mathrm{of}\mathrm{unity})\hfill \\ 8(t\mathrm{is}\mathrm{a}\mathrm{root}\mathrm{of}\mathrm{unity}).\hfill \end{array}$$
Therefore, if $`t`$ is not a root of unity, then
$$H^0(\mathrm{\Theta }_S(\mathrm{log}D)N_D^n)=0.$$
Acknowledgements
The author is deeply grateful to Professor Masa-Hiko Saito, who gave support to the whole of this work. We would also like to thank Kenji Iohara for many useful advices.
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# Contents
## 1 Introduction
At the end of LEP2 operation the total cross section for the process $`e^{}e^+f\overline{f}+n\gamma `$ will have to be calculated with the precision $`0.2\%1\%`$, depending on event selection. The arbitrary differential distributions have to be also calculated with the corresponding precision. In future linear colliders (LC’s) the precision requirement can be even more demanding. This is especially true for high luminosity linear colliders, like in the case of TESLA. The above new requirements necessitate development of the new calculational framework for the QED corrections and the construction of new dedicated MC programs. The present work is a part of an effort in this direction.
The main limiting factor preventing us from getting more precise theoretical predictions for the $`e^{}e^+f\overline{f}+n\gamma `$ process is higher-order QED radiative corrections (QED part of electroweak Standard Model). In order to achieve the 0.2% precision tag, the virtual corrections have to be calculated up to two-three loops and the multiple bremsstrahlung up to two-three hard photons, integrating exactly the multiphoton phase-space for the arbitrary event selection (phase-space limits).
For any realistic kinematical cuts, one cannot get the precise theoretical predictions for $`e^{}e^+f\overline{f}+n\gamma `$ at the above ambitious precision level without Monte Carlo (MC) event generators. It is therefore mandatory to formulate perturbative Standard Model (SM) calculations in a formulation friendly to their use within the Monte Carlo event generator.
Let us stress that the Monte Carlo method is for us nothing more and nothing less than the numerical integration over the Lorentz invariant phase-space. It is therefore an exercise in the applied mathematics. In the present work we shall not, however, elaborate on the methods of the Monte Carlo phase-space integration and construction of the Monte Carlo event generator. This is delegated to ref. which describes the new Monte Carlo event generator $`𝒦𝒦`$ in which the matrix element of the present paper is implemented and all numerical results presented here are calculated using the latest version 4.13 of $`𝒦𝒦`$.
In the present work we concentrate on the definition and construction of the matrix element for the process $`e^{}e^+f\overline{f}`$ within Standard Model. We shall especially address the problem of the higher-order QED corrections. This work is a continuation of two recent papers and .
### 1.1 Two types of QED matrix element and exponentiations
In the $`𝒦𝒦`$ Monte Carlo and in this paper we use two types of matrix element with two types of exponentiation: exclusive exponentiation nicknamed EEX and coherent exclusive exponentiation referred to as CEEX. Both are termed as “exclusive” as opposed to “inclusive”, see also the discussion in . The exclusivity means that the procedure of exponentiation, that is summing up the infrared (IR) real and virtual contribution, within the standard perturbative scheme of quantum field theory, is done at the level of the fully differential (multiphoton) cross section, or even better, at the level of the scattering matrix element (spin amplitudes), before any phase-space integration over photon momenta is done.
The other “inclusive” exponentiation is an ad hoc procedure of summing up IR corrections after phase-space integration over photon momenta, that is for inclusive distributions. In spite of its weak theoretical basis the inclusive exponentiation is very commonly done routinely in all semi-analytical approaches like that in ref. . In Section 5.1 we shall come back to inclusive exponentiation and show how to justify it theoretically.
The two exclusive exponentiations EEX and CEEX are well suited for the fully exclusive Monte Carlo event generators in which four-momenta of all final-state particles are available. Historically EEX was formulated for the first time in ref. for the initial-state radiation (ISR) and an improved version was presented in ref. . It follows very closely the Yennie-Frautschi-Suura (YFS) exponentiation of the classical ref. . The extension of EEX to the final-state radiation (FSR) was done shortly thereafter , but it was actually never fully published. The computer program YFS3, in which EEX for FSR was implemented, was incorporated in KORALZ and some numerical results were published in , without actually giving details of the QED matrix element. The present work gives in fact the first full account of the EEX matrix element for ISR and FSR for the process $`e^{}e^+f\overline{f}+n\gamma ,fe`$. This is to be contrasted with the situation for small angle Bhabha scattering (the well-known LEP-SLC luminosity process) where the EEX type matrix element was fully documented in refs. .
CEEX is a new version of the exclusive exponentiation, generally more efficient for calculations beyond first-order, facilitating inclusion of full spin polarization, narrow resonances and any kind of interferences. Its first version, limited to first-order, was presented in ref. . In the present work we extend it to (still incomplete) second-order.
Let us characterize briefly the main features of EEX and CEEX. EEX is formulated in terms of spin summed/averaged differential distributions, this is the source of some advantages and disadvantages which may be summarized as follows:
* Differential distributions in practice are given analytically in terms of Mandelstam variables and scattering angles, they are therefore easy to inspect by human eye and to check correctness of certain important limits like leading-logarithmic and soft limits.
* Analytical representation of the differential distributions allows for analytical phase-space integration and development of the semi-analytical formulas, which are useful for cross-check with the MC results.
* Spin effects are difficult to add already at $`𝒪(\alpha ^1)`$, because one is forced to calculate radiative corrections to spin density matrices, not an easy task.
* Squaring sums of spin amplitudes from groups of Feynman diagrams leads to many interference terms which in the exponentiation procedure are handled analytically and individually. Because of that interference terms can be dealt with efficiently in EEX only for simple processes dominated by a small number of Feynman diagrams and only up to first-order.
CEEX is formulated in terms of spin amplitudes and this is also the source of some advantages and disadvantages:
* Differential distributions are calculated out of spin amplitudes numerically – spin amplitudes are generally simpler/smaller objects, especially beyond $`𝒪(\alpha ^1)`$.
* Since an analytical representation for differential distributions is not available semi-analytical integration over the phase-space is practically impossible.
* Spin effects are added relatively easily, during numerical evaluation of the differential distributions out of spin amplitudes. Adding higher-order corrections does no make the treatment of spin polarization more difficult.
* Inclusion of all kind of interference effects (among real photon emissions, many Feynman diagrams etc.) comes almost for free – it is done numerically in the process of summing and squaring various contributions to spin amplitudes.
As we see CEEX has many advantages over EEX, so why do we keep EEX? There are important reasons:
* Generally, CEEX is a relatively new invention, the older and more primitive but well established EEX is a useful reference for numerical tests of CEEX.
* EEX is better suited for semi-analytical integration over the phase-space, and can be tested with these semi-analytical results.
* In the present $`𝒦𝒦`$ MC the $`𝒪(\alpha ^3)`$ leading logarithmic corrections are available for EEX and are not yet available for CEEX.
Summarizing, we see that it make sense to keep EEX as a backup solution even if we already rely on CEEX as a default and leading solution.
### 1.2 Notation, terminology
It is useful to introduce certain notation and terminology already at this stage. In particular, the most common perturbative calculation (no exponentiation) is “order-by-order”. That is all terms beyond a certain order are set to zero. In Fig. 1 that means we end at certain row – at $`𝒪(\alpha ^2)`$ we include the first three rows. Exponentiation is blurring this picture because a certain class of terms is summed up to infinite order and the meaning of the $`rth`$ order exponentiation is that we truncate to $`𝒪(\alpha ^r)`$ the infrared (IR) finite components, the so-called $`\beta `$’s. On the other hand, in the leading-logarithmic approximation the focus is on summing up first the contributions like $`\alpha ^nL^n`$ and later $`\alpha ^nL^{n1}`$, that is in Fig. 1 we sum up in column-wise order, neglecting terms far away from the first column which represents the so-called LL-approximation. Taking the actual value of $`\alpha /<pi1/400`$ and of the big logarithm $`L=\mathrm{ln}(s/m_f^2)10`$, we discover quickly that in Fig. 1 the limiting line following the numerical importance of the terms is neither row-wise nor column-wise but diagonal-wise. This is why we shall often use $`𝒪(\alpha ^r)_{prag}`$ $`r=1,2,3`$ approximation, depicted also in Fig. 1, in which we use (exponentiated or not) $`𝒪(\alpha ^r)`$ calculation in which we use incomplete sub-leading terms, in the sense of the LL approximation. Note that for the LL approximation we shall never use the strict collinear (zero $`p_T`$) approximation. The LL approximation will be done at the level of the differential distributions (or spin amplitudes) without forcing $`p_T=0`$ on photons. Just to give a rough idea, the precision level of order $`0.51\%`$ corresponds to $`𝒪(\alpha ^1)_{prag}`$, $`0.10.5\%`$ to $`𝒪(\alpha ^2)_{prag}`$ and going below $`0.05\%`$ will require $`𝒪(\alpha ^3)_{prag}`$. The above is true for the exponentiated calculation. Lack of exponentiation makes the calculation less precise by a factor $`25`$. The pure non-logarithmic terms of order $`𝒪(\alpha ^2)`$ are negligible ($`<10^5`$) for any foreseeable practical application.
### 1.3 Outline
The outline of the paper is the following. In Section 2 we describe in detail the SM/QED matrix element for the exclusive exponentiation (EEX) based on the Yennie-Frautschi-Suura (YFS) work of ref. , that is of the type of matrix element defined for the first time in ref. . In Section 2 we describe the new second-order matrix element with coherent exclusive exponentiation (CEEX), which is the default matrix element in $`𝒦𝒦`$ MC. Its first-order variant was given in , and is also defined here for the sake of completeness. In Section 3 we elaborate on how do we combine the electroweak corrections of refs. with the QED corrections within EEX and CEEX. In Section 4 we discuss the differences between EEX and CEEX. In Section 5 we integrate analytically over the phase-space for the EEX matrix element in the case of very simple kinematical cuts. The resulting analytical results are used in Section 6 where numerical results from $`𝒦𝒦`$ MC are presented. The most important task in Section 6 is, however, the determination of the physical and technical precision for the total cross section and charge asymmetry at the $`Z`$-peak, LEP2 and 500 GeV. In particular we discuss the contribution from the initial-final state interference (IFI) which is included in our new CEEX matrix element (IFI is neglected in EEX). In the last Section 7 we summarize our work. In Appendix A we define the Weyl-spinor techniques used in construction of CEEX multi-photon spin amplitudes.
## 2 Amplitudes for Exclusive Exponentiation
As it was already indicated, the role of the EEX matrix element described in this section is to provide a testing environment for the new more sophisticated matrix element of the CEEX class, which will be defined in the next section.
The kinematics of the process $`e^{}e^+f\overline{f}+n\gamma `$ is depicted in fig. 2. In the case of the EEX matrix element presented here we neglect the initial-final state interference (IFI). Consequently, we are allowed in the following to distinguish among photons emitted from the initial- and final-state fermions. The four-momentum
$`X=p_1+p_2{\displaystyle \underset{j=1}{\overset{n}{}}}k_j=q_1+q_2+{\displaystyle \underset{l=1}{\overset{n^{}}{}}}k_l^{}`$ (1)
of the $`s`$-channel virtual boson $`Z+\gamma `$ is then well defined. Let us denote the rest frame of $`X`$ as XMS.
### 2.1 Master formula
Denoting Lorentz invariant phase-space by
$$d^n\mathrm{Lips}(P;p_1,p_2,\mathrm{},p_n)=\underset{j=1}{\overset{n}{}}\frac{d^3p_j}{p_j^0}\delta ^{(4)}\left(P\underset{j=1}{\overset{n}{}}p_j\right)$$
(2)
we define for the process $`e^{}(p_1)+e^+(p_2)f(q_1)+\overline{f}(q_2)+n\gamma (k_j)+n^{}\gamma (k_l^{})`$ the $`𝒪(\alpha ^r)`$ total cross-section
$$\sigma _{EEX}^{(r)}=\underset{n=0}{\overset{\mathrm{}}{}}\underset{n^{}=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\frac{1}{n^{}!}d^{n+n^{}}\mathrm{Lips}(p_1+p_2;q_1,q_2,k_1\mathrm{},k_n,k_1^{}\mathrm{},k_n^{}^{})\rho _{EEX}^{(r)},r=0,1,2,3$$
(3)
in terms of the fully differential multiphoton distribution
$$\begin{array}{cc}\hfill \rho & {}_{EEX}{}^{(r)}(p_1,p_2,q_1,q_2,k_1\mathrm{},k_n,k_1^{}\mathrm{},k_n^{})=e^{Y_e(\mathrm{\Omega }_I;p_1,p_2)+Y_f(\mathrm{\Omega }_F;q_1,q_2)}\hfill \\ & \underset{j=1}{\overset{n}{}}\stackrel{~}{S}_I(k_j)\overline{\mathrm{\Theta }}(\mathrm{\Omega }_I;k_j)\underset{l=1}{\overset{n^{}}{}}\stackrel{~}{S}_F(k_l^{})\overline{\mathrm{\Theta }}(\mathrm{\Omega }_F;k_l^{})\{\overline{\beta }_0^{(r)}(X,p_1,p_2,q_1,q_2)\hfill \\ & +\underset{j=1}{\overset{n}{}}\frac{\overline{\beta }_{1I}^{(2)}(X,p_1,p_2,q_1,q_2,k_j)}{\stackrel{~}{S}_I(k_j)}+\underset{l=1}{\overset{n^{}}{}}\frac{\overline{\beta }_{1F}^{(2)}(X,p_1,p_2,q_1,q_2,k_l)}{\stackrel{~}{S}_F(k_l)}\hfill \\ & +\underset{nj>k1}{}\frac{\overline{\beta }_{2II}^{(2)}(X,p_1,p_2,q_1,q_2,k_j,k_k)}{\stackrel{~}{S}_I(k_j)\stackrel{~}{S}_I(k_k)}+\underset{n^{}l>m1}{}\frac{\overline{\beta }_{2FF}^{(2)}(X,p_1,p_2,q_1,q_2,k_l,k_m)}{\stackrel{~}{S}_F(k_l)\stackrel{~}{S}_F(k_m)}\hfill \\ & +\underset{j=1}{\overset{n}{}}\underset{l=1}{\overset{n^{}}{}}\frac{\overline{\beta }_{2IF}^{(2)}(X,p_1,p_2,q_1,q_2,k_j,k_l)}{\stackrel{~}{S}_I(k_j)\stackrel{~}{S}_F(k_l)}+\underset{nj>k>l1}{}\frac{\overline{\beta }_{3III}^{(3)}(X,p_1,p_2,q_1,q_2,k_j,k_k,k_l)}{\stackrel{~}{S}_I(k_j)\stackrel{~}{S}_I(k_k)\stackrel{~}{S}_I(k_l)}\}.\hfill \end{array}$$
(4)
Let us explain the notation and physics content in the above expression. The YFS soft factors for real photons emitted from the initial- and final-state fermions read
$$\stackrel{~}{S}_I(k_j)=Q_e^2\frac{\alpha }{4\pi ^2}\left(\frac{p_1}{k_jp_1}\frac{p_2}{k_jp_2}\right)^2,\stackrel{~}{S}_F(k_l)=Q_f^2\frac{\alpha }{4\pi ^2}\left(\frac{q_1}{k_lq_1}\frac{q_2}{k_lq_2}\right)^2,$$
(5)
where electric charges of the electron and fermion $`f`$ are $`Q_e`$ and $`Q_f`$. The $`Y`$-function in the exponential YFS form factor is defined as in ref. :
$$\begin{array}{cc}\hfill Y_f(\mathrm{\Omega },p,\overline{p})& 2Q_f^2\alpha \stackrel{~}{B}(\mathrm{\Omega },p,\overline{p})+2Q_f^2\alpha \mathrm{}B(p,\overline{p})\hfill \\ \hfill & 2Q_f^2\alpha \frac{1}{8\pi ^2}\frac{d^3k}{2k^0}\mathrm{\Theta }(\mathrm{\Omega };k)\left(\frac{p}{kp}\frac{\overline{p}}{k\overline{p}}\right)^2\hfill \\ & +2Q_f^2\alpha \mathrm{}\frac{d^4k}{k^2}\frac{i}{(2\pi )^3}\left(\frac{2pk}{2kpk^2}\frac{2\overline{p}k}{2k\overline{p}k^2}\right)^2.\hfill \end{array}$$
(6)
The above form factor is infrared-finite and depends explicitly on the soft photon domains $`\mathrm{\Omega }=\mathrm{\Omega }_I,\mathrm{\Omega }_F`$ which includes (surrounds) the IR divergence point $`k=0`$. We define $`\mathrm{\Theta }(\mathrm{\Omega };k)=1`$ for $`k\mathrm{\Omega }`$ and $`\mathrm{\Theta }(\mathrm{\Omega };k)=0`$ for $`k\mathrm{\Omega }`$. Contributions from the real photons inside $`\mathrm{\Omega }`$ are summed to infinite-order and combined with the analogous virtual contributions forming the exponential YFS form factor. In the Monte Carlo we generate photons $`k\mathrm{\Omega }`$ characterized by the function $`\overline{\mathrm{\Theta }}(\mathrm{\Omega },k)=1\mathrm{\Theta }(\mathrm{\Omega },k)`$. We require, as usual, that $`\mathrm{\Omega }_I`$ and $`\mathrm{\Omega }_F`$ are small enough (they can be chosen arbitrarily small) such that the total cross section as defined in eq. (4) and any other physically meaningful observable do not depend on the actual choice of them, i.e. $`\mathrm{\Omega }_{I,F}`$ are dummy parameters in the calculation! If we neglect the initial-final state interference then we may choose $`\mathrm{\Omega }_I\mathrm{\Omega }_F`$. Let us define $`\mathrm{\Omega }_I`$ with the $`k^0<E_{min}`$ condition in the centre of the mass system of incoming $`e^\pm `$ beams and $`\mathrm{\Omega }_F`$ with $`k^0<E_{min}^{}`$ in the centre of the mass of the outgoing fermions $`f\overline{f}`$. The two domains differ because the Lorentz frames in which they are defined are different. The above choice is the easiest for the Monte Carlo generation but in the later discussion we shall describe in detail how do we implement the $`\mathrm{\Omega }_i=\mathrm{\Omega }_F`$ option in our Monte Carlo. The actual YFS form factors for the above choices are well known :
$$\begin{array}{cc}\hfill Y_e(\mathrm{\Omega }_I;p_1,p_2)& =\gamma _e\mathrm{ln}\frac{2E_{min}}{\sqrt{2p_1p_2}}+\frac{1}{4}\gamma _e+Q_e^2\frac{\alpha }{\pi }\left(\frac{1}{2}+\frac{\pi ^2}{3}\right),\hfill \\ \hfill Y_f(\mathrm{\Omega }_F;q_1,q_2)& =\gamma _f\mathrm{ln}\frac{2E_{min}}{\sqrt{2q_1q_2}}+\frac{1}{4}\gamma _f+Q_f^2\frac{\alpha }{\pi }\left(\frac{1}{2}+\frac{\pi ^2}{3}\right),\hfill \end{array}$$
(7)
where
$$\gamma =\gamma _e=2Q_e^2\frac{\alpha }{\pi }\left(\mathrm{ln}\frac{2p_1p_2}{m_e^2}1\right),\gamma _f=2Q_f^2\frac{\alpha }{\pi }\left(\mathrm{ln}\frac{2q_1q_2}{m_f^2}1\right).$$
(8)
### 2.2 Pure virtual corrections
The perturbative QED matrix element is located in the $`\overline{\beta }`$-functions. The $`\overline{\beta }_0`$ function is “proportional” to the Born $`e^{}e^+f\overline{f}`$ differential cross section $`d\sigma ^{\mathrm{Born}}(s,\vartheta )/d\mathrm{\Omega }`$ and it contains (infrared-finite) corrections calculable order by order. According to our general strategy we shall calculate $`\overline{\beta }_0`$ and other $`\overline{\beta }`$’s in the $`𝒪(\alpha ^i)_{prag}`$, $`i=0,1,2`$.
The $`𝒪(\alpha ^i)_{prag}`$ expressions for $`\overline{\beta }_0^{(i)},i=0,1,2`$ read<sup>1</sup><sup>1</sup>1 It may look that we miss pure $`(\alpha /\pi )`$ term in $`\delta _{I,F}^{(1)}`$. The calculation shows that such a non-logarithmic contribution is accidentally equal zero.
$`\overline{\beta }_0^{(r)}(X,p_1,p_2,q_1,q_2)=(1+\delta _I^{(r)})(1+\delta _F^{(r)}){\displaystyle \frac{1}{4}}{\displaystyle \underset{k,l=1,2}{}}{\displaystyle \frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}}(X^2,\vartheta _{kl})`$ (9)
$`\delta _I^{(0)}=0,\delta _I^{(1)}={\displaystyle \frac{1}{2}}\gamma ,\delta _I^{(2)}=\delta _I^{(1)}+{\displaystyle \frac{1}{8}}\gamma ^2,\delta _I^{(3)}=\delta _I^{(2)}+{\displaystyle \frac{1}{48}}\gamma ^3,`$
$`\delta _F^{(0)}=0,\delta _F^{(1)}={\displaystyle \frac{1}{2}}\gamma _f,\delta _F^{(2)}=\delta _I^{(1)}+{\displaystyle \frac{1}{8}}\gamma _f^2,\delta _F^{(3)}=\delta _I^{(2)}+{\displaystyle \frac{1}{48}}\gamma _f^3,`$
where
$$\vartheta _{11}=\mathrm{}(\stackrel{}{p}_1,\stackrel{}{q}_1),\vartheta _{12}=\mathrm{}(\stackrel{}{p}_1,\stackrel{}{q}_2),\vartheta _{21}=\mathrm{}(\stackrel{}{p}_2,\stackrel{}{q}_1),\vartheta _{22}=\mathrm{}(\stackrel{}{p}_2,\stackrel{}{q}_2),$$
(10)
with all 3-vectors taken in the rest frame of the four-momentum $`X`$, that is in the frame XMS.
Let us first explain the fact that instead of having a single $`d\sigma ^{\mathrm{Born}}/d\mathrm{\Omega }(\vartheta )`$ with a single $`\vartheta `$ we take an average over four $`\vartheta _{kl}`$. In fact we could adopt one $`\vartheta `$, for example $`\vartheta _0=\mathrm{}(\stackrel{}{p}_1\stackrel{}{p}_2,\stackrel{}{q}_1\stackrel{}{q}_2)`$ where all three-momenta are taken in XMS. The main reason for our apparently more complicated choice is related to the implementation of the first and higher-order real photon contributions in the next subsections. More precisely, it is well known that the exact single photon ISR matrix element can be cast as a linear combination of the two $`d\sigma ^{\mathrm{Born}}/d\mathrm{\Omega }(\vartheta _k),k=1,2`$ distributions. The same is true for FSR . (Our implementation of the leading-logarithmic (LL) matrix element for 2 and 3 real photons will also involve the linear combination of this type.) It is therefore logical and practical to use a similar solution already for $`\overline{\beta }_0`$. One should also keep in mind that in the soft limit, when all photons are soft, then all four angles $`\vartheta _{kl}`$ are identical and the averaging over them is a spurious operation anyway.
The reader not familiar with exponentiation may have an even more elementary question: Why do we have a freedom of defining $`\vartheta `$ in $`d\sigma ^{\mathrm{Born}}/d\mathrm{\Omega }(\vartheta )`$ in first place? Is this ambiguity dangerous? These questions are already discussed in refs. . The answer is the following: Strictly speaking the differential cross section $`d\sigma ^{\mathrm{Born}}(s,\vartheta )/d\mathrm{\Omega }`$ and $`\overline{\beta }_0^{(i)}`$ are defined within the two body phase-space. Later on they are used, however, in eq. (4) and in the definitions of $`\overline{\beta }^{(i)},i=1,2,\mathrm{}`$ all over the phase-space with additional soft and/or hard photons. This requires some kind of extrapolation of $`\overline{\beta }_0`$ and $`d\sigma ^{\mathrm{Born}}(s,\vartheta )/d\mathrm{\Omega }`$ beyond the two body phase-space. In ref. this extrapolation was done using manipulations on the four-momenta and in ref. it was done as an extrapolation in the Mandelstam variables $`s,t,u`$. Here we present another solution which is somewhere in between the previous two ones. What is really important, however, is that the effect due to change from one particular choice of extrapolation to another is always, for the entire calculation, a kind of “higher-order effect”. For instance at $`𝒪(\alpha ^1)`$ changing the type of extrapolation is an $`𝒪(\alpha ^2)`$ effect! Of course, it is always wise to use some kind of “smooth” extrapolation which is able to minimize the higher-order effects.
Another possible question is: Why we did not write down the second-order virtual correction factor in an additive way, like for instance $`(1+\delta _I^{(2)}+\delta _F^{(2)}+\delta _I^{(1)}\delta _F^{(1)})`$? We have opted for factorized form because it is generally known that the factorized form is closer to reality at higher perturbative orders. Another important reason is that the factorized form is easier for semi-analytical integrations over the phase-space in the next section.
### 2.3 One real photon with virtual corrections
The contributions $`\overline{\beta }_1^{(2)}`$ are needed directly in eq. (4) and the $`𝒪(\alpha ^1)_{prag}`$ version of $`\overline{\beta }_1^{(1)}`$ enters indirectly as a construction element in $`\overline{\beta }_2`$. They are constructed from QED distributions with a single real photon emission and up to one virtual photon contribution. They are defined separately for initial- and final-state photons
$$\begin{array}{cc}\hfill \overline{\beta }_{1I}^{(i)}(X,p_1,p_2,q_1,q_2,k_j)=& D_{1I}^{(i)}(X,p_1,p_2,q_1,q_2,k_j)\stackrel{~}{S}_I(k_j)\overline{\beta }_0^{(i1)}(X,p_1,p_2,q_1,q_2),\hfill \\ \hfill \overline{\beta }_{1F}^{(i)}(X,p_1,p_2,q_1,q_2,k_l^{})=& D_{1F}^{(i)}(X,p_1,p_2,q_1,q_2,k_l^{})\stackrel{~}{S}_F(k_l^{})\overline{\beta }_0^{(i1)}(X,p_1,p_2,q_1,q_2),\hfill \end{array}$$
(11)
where $`i=1,2`$. Let us define first all ingredients for the initial-state contribution. The single initial-state photon emission differential distribution at $`𝒪(\alpha ^r)`$, $`r=1,2,3`$, with the eventual additional up to two-loop virtual correction from the initial- and/or final-state photon reads
$$\begin{array}{cc}\hfill D_{1I}^{(r)}(X,& p_1,p_2,q_1,q_2,k_j)=Q_e^2\frac{\alpha }{4\pi ^2}\frac{2p_1p_2}{(k_jp_1)(k_jp_2)}W_e(\widehat{\alpha }_j,\widehat{\beta }_j)\hfill \\ & \left\{\frac{(1\widehat{\alpha }_j)^2}{2}\underset{r=1,2}{}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(X^2,\vartheta _{1r})+\frac{(1\widehat{\beta }_j)^2}{2}\underset{r=1,2}{}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(X^2,\vartheta _{2r})\right\}\hfill \\ & \left(1+\mathrm{\Delta }_I^{(r1)}(z_j)\right)(1+\delta _F^{(r1)}),\hfill \end{array}$$
(12)
where
$$\begin{array}{cc}& \widehat{\alpha }_j=\frac{k_jp_2}{p_1p_2},\widehat{\beta }_j=\frac{k_jp_1}{p_1p_2},z_j=(1\widehat{\alpha }_j)(1\widehat{\beta }_j),\hfill \\ & \mathrm{\Delta }_I^{(0)}(z)0,\mathrm{\Delta }_I^{(1)}(z)\frac{1}{2}\gamma \frac{1}{4}\gamma \mathrm{ln}(z),\hfill \\ & \mathrm{\Delta }_I^{(2)}(z)\mathrm{\Delta }_I^{(1)}(z)+\frac{1}{8}\gamma ^2\frac{1}{8}\gamma ^2\mathrm{ln}(z)+\frac{1}{24}\gamma ^2\mathrm{ln}^2(z),\hfill \\ & W_e(a,b)1\frac{m_e^2}{2p_1p_2}\frac{(1a)(1b)}{(1a)^2+(1b)^2}\left(\frac{a}{b}+\frac{b}{a}\right).\hfill \end{array}$$
(13)
Again the question arises why the averaging over $`r`$ in $`\vartheta _{kr}`$ is introduced? In the case of just one ISR hard photon the averaging trivially disappears because $`\vartheta _{k1}=\vartheta _{k2}`$ and in this case our formula coincides with the exact $`𝒪(\alpha ^1)`$ result, see , as it should. In the less trivial case of the presence of the additional hard photons there is an ambiguity in defining $`D_{1I}^{(r)}`$ which is reflected in our “averaging” procedure; however, it is harmless i.e. the effect is of $`𝒪(\alpha ^2)`$.,
It is necessary and interesting to check the soft limit. If in the presence of many additional photons ($`n>1`$) we take the soft limit $`k_j0`$, keeping momenta of other photons constant, then $`\vartheta _{kr}`$ are in general all different. However, in eq. (12) the sums over $`d\sigma ^{\mathrm{Born}}/d\mathrm{\Omega }`$ combine into a simple average over all four angles, as in eq. (LABEL:beta0) – in fact the single photon distribution reduces to
$$D_{1I}^{(2,1)}(X,p_1,p_2,q_1,q_2,k_j)\stackrel{~}{S}_I(k_j)\overline{\beta }_0^{(1,0)}(X,p_1,p_2,q_1,q_2)$$
and therefore $`\overline{\beta }_{1I}^{(2,1)}(X,p_1,p_2,q_1,q_2,k_j)`$ is infrared-finite as required. The above argument shows that the extrapolations for $`\overline{\beta }_0`$ and $`\overline{\beta }_1`$ have to be of the same type. If we have opted for another extrapolation in eq. (12), for example without averaging, with a single angle $`\vartheta _{kr}\vartheta _k`$, then the extrapolation in eqs. (LABEL:beta0) would need to be changed appropriately.
Another interesting limit is the collinear limit. If all (possibly hard) photons are collinear to initial- or final-fermions then all angles $`\vartheta _{sr},s,r=1,2`$ are identical and equal to the familiar leading-logarithmic effective scattering angle for the hard process in the “reduced frame” XMS. This will facilitate introduction of the higher-order LL corrections in the following.
Another remark on eq. (12) is in order: There are many equivalent ways, modulo term of $`𝒪(m^2/s)`$, of writing the single bremsstrahlung spin summed differential distribution . Our choice follows the representation implemented in the Monte Carlo programs YFS2 , KORALZ and MUSTRAAL , because it minimizes the machine rounding errors (quite important due to the smallness of electron mass), and it is explicitly expressed in terms of Born differential cross sections – this feature facilitates introduction of electroweak corrections.
The virtual correction $`(1+\mathrm{\Delta }_I^{(1)}(\widehat{\alpha }_j,\widehat{\beta }_j))`$ is taken in the leading logarithmic approximation (sufficient for our $`𝒪(\alpha ^2)_{prag}`$ approach) and it agrees with the corresponding distribution in ref. . In the $`k_j0`$ limit we have $`\mathrm{\Delta }_I^{(1)}(\widehat{\alpha }_j,\widehat{\beta }_j)\delta _I^{(1)}`$ as expected, and as required for infrared finiteness of $`\overline{\beta }_{1F}^{(2)}`$. The other factor $`(1+\delta _F^{(1)})`$ represents the contribution from the simultaneous emission of the real initial and the virtual final photon. We again prefer the factorized form over an additive one $`(1+\mathrm{\Delta }_I^{(1)}+\delta _F^{(1)})`$.
The essential ingredients for the $`𝒪(\alpha ^r)`$ final-state $`\overline{\beta }_{1F}^{(r)},r=1,2`$, is the single final-state photon emission matrix element with up to one-loop virtual initial- or final-state photon corrections
$$\begin{array}{cc}\hfill D_{1F}^{(r)}(X,& p_1,p_2,q_1,q_2,k_l^{})=Q_f^2\frac{\alpha }{4\pi ^2}\frac{2q_1q_2}{(k_l^{}q_1)(k_l^{}q_2)}W_f(\widehat{\eta }_l,\widehat{\zeta }_l)\hfill \\ & \left\{\frac{(1\widehat{\eta }_l)^2}{2}\underset{r=1,2}{}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(X^2,\vartheta _{r1})+\frac{(1\widehat{\zeta }_l)^2}{2}\underset{r=1,2}{}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(X^2,\vartheta _{r2})\right\}\hfill \\ & \left(1+\mathrm{\Delta }_F^{(r1)}(z_l)\right)(1+\delta _I^{(r1)})\hfill \end{array}$$
(14)
where
$$\begin{array}{cc}& \eta _l=\frac{k_l^{}q_2}{q_1q_2},\zeta _l=\frac{k_l^{}q_1}{q_1q_2},\widehat{\eta }_l=\frac{\eta _l}{1+\eta _l+\zeta _l},\widehat{\zeta }_l=\frac{\zeta _l}{1+\eta _l+\zeta _l},\hfill \\ & z_l=(1\widehat{\eta }_l)(1\widehat{\zeta }_l)\hfill \\ & \mathrm{\Delta }_F^{(0)}(z)0,\mathrm{\Delta }_F^{(1)}(z)\frac{1}{2}\gamma _f+\frac{1}{4}\gamma _f\mathrm{ln}(z),\hfill \\ & W_f(a,b)1\frac{m_f^2}{2q_1q_2}\frac{(1a)(1b)}{(1a)^2+(1b)^2}\left(\frac{a}{b}+\frac{b}{a}\right),\hfill \end{array}$$
(15)
All discussion on the ISR distribution of eq. (12) applies also to the above FSR distribution.
### 2.4 Two real photons with virtual corrections
The contributions $`\overline{\beta }_{II}^{(2)},\overline{\beta }_{FF}^{(2)}`$ and $`\overline{\beta }_{IF}^{(2)}`$ are related to emission of the real two initial, two final and one initial and one final photons correspondingly. They are genuine $`𝒪(\alpha ^2)`$ objects because they appear in this order for the first time. For the same reason they do not include any virtual contributions. They are defined formally in the usual way
$`\begin{array}{cc}\hfill \overline{\beta }_{II}^{(r)}(X& ,p_1,p_2,q_1,q_2,k_j,k_k)=D_{II}^{(r)}(X,p_1,p_2,q_1,q_2,k_j,k_k)\hfill \\ & \stackrel{~}{S}_I(k_j)\overline{\beta }_{1I}^{(r1)}(X,p_1,p_2,q_1,q_2,k_k)\stackrel{~}{S}_I(k_k)\overline{\beta }_{1I}^{(r1)}(X,p_1,p_2,q_1,q_2,k_j)\hfill \\ & \stackrel{~}{S}_I(k_j)\stackrel{~}{S}_I(k_k)\overline{\beta }_0^{(r2)}(X,p_1,p_2,q_1,q_2),r=2,3,\hfill \end{array}`$ (16)
$`\begin{array}{cc}\hfill \overline{\beta }_{FF}^{(r)}(X& ,p_1,p_2,q_1,q_2,k_l^{},k_m^{})=D_{FF}^{(r)}(X,p_1,p_2,q_1,q_2,k_l^{},k_m^{})\hfill \\ & \stackrel{~}{S}_F(k_l^{})\overline{\beta }_{1F}^{(r1)}(X,p_1,p_2,q_1,q_2,k_m^{})\stackrel{~}{S}_F(k_m^{})\overline{\beta }_{1F}^{(r1)}(X,p_1,p_2,q_1,q_2,k_l^{})\hfill \\ & \stackrel{~}{S}_F(k_l^{})\stackrel{~}{S}_F(k_m^{})\overline{\beta }_{r2}^{(r2)}(X,p_1,p_2,q_1,q_2),r=2,3,\hfill \end{array}`$
$`\begin{array}{cc}\hfill \overline{\beta }_{IF}^{(r)}(X& ,p_1,p_2,q_1,q_2,k_j,k_l^{})=D_{IF}^{(r)}(X,p_1,p_2,q_1,q_2,k_j,k_l^{})\hfill \\ & \stackrel{~}{S}_I(k_j)\overline{\beta }_{1F}^{(r1)}(X,p_1,p_2,q_1,q_2,k_l^{})\stackrel{~}{S}_F(k_l^{})\overline{\beta }_{1I}^{(r1)}(X,p_1,p_2,q_1,q_2,k_j)\hfill \\ & \stackrel{~}{S}_I(k_j)\stackrel{~}{S}_F(k_l^{})\overline{\beta }_{r2}^{(0)}(X,p_1,p_2,q_1,q_2),r=2,3.\hfill \end{array}`$
The new objects in the above expressions are the differential distributions $`D_{II}^{(2)},D_{FF}^{(2)}`$ and $`D_{IF}^{(2)}`$ for double bremsstrahlung. They are not taken directly from Feynman diagrams but they are constructed in such a way that:
* If one photon is hard and one is soft then the single bremsstrahlung expression of eqs. (12,14) are recovered
* If both photons are hard and collinear then the proper LL limit, which we know from the double or triple convolution of the Altarelli-Parisi kernels, is also recovered.
The resulting expressions are rather compact and the LL limit is manifest, this is not necessarily true for the exact double bremsstrahlung spin amplitudes (see next section). The method is similar to that of refs. . In the case of ISR we shall also include one-loop virtual corrections read from the triple convolution of the Altarelli-Parisi kernels, see ref. .
Our construction in the case of the double real ISR reads as follows
$$\begin{array}{cc}\hfill D& {}_{II}{}^{(2)}(X,p_1,p_2,q_1,q_2,k_1,k_2)\hfill \\ & Q_e^4\frac{\alpha }{4\pi ^2}\frac{2p_1p_2}{(k_1p_1)(k_1p_2)}\frac{\alpha }{4\pi ^2}\frac{2p_1p_2}{(k_2p_1)(k_2p_2)}W_e(\widehat{\alpha }_1,\widehat{\beta }_1)W_e(\widehat{\alpha }_2,\widehat{\beta }_2)\hfill \\ & \{\mathrm{\Theta }(v_1v_2)(1+\mathrm{\Delta }_{II}^{(r1)}(z_1,z_{12}))(1+\delta _F^{(r1)})\hfill \\ & \left[\chi _2(\widehat{\alpha }_1;\widehat{\alpha }_2^{},\widehat{\beta }_2^{})\underset{r=1,2}{}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(X^2,\vartheta _{1r})+\chi _2(\widehat{\beta }_1;\widehat{\alpha }_2^{},\widehat{\beta }_2^{})\underset{r=1,2}{}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(X^2,\vartheta _{2r})\right]\hfill \\ & +\mathrm{\Theta }(v_2v_1)\left(1+\mathrm{\Delta }_{II}^{(r1)}(z_2,z_{21})\right)(1+\delta _F^{(r1)})\hfill \\ & [\chi _2(\widehat{\alpha }_2;\widehat{\alpha }_1^{},\widehat{\beta }_1^{})\underset{r=1,2}{}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(X^2,\vartheta _{1r})+\chi _2(\widehat{\beta }_2;\widehat{\alpha }_1^{},\widehat{\beta }_1^{})\underset{r=1,2}{}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(X^2,\vartheta _{2r})]\},\hfill \end{array}$$
(17)
where
$$\begin{array}{cc}& \widehat{\alpha }_1^{}=\frac{\widehat{\alpha }_1}{1\widehat{\alpha }_2},\widehat{\alpha }_2^{}=\frac{\widehat{\alpha }_2}{1\widehat{\alpha }_1},\widehat{\beta }_1^{}=\frac{\widehat{\beta }_1}{1\widehat{\beta }_2},\widehat{\beta }_2^{}=\frac{\widehat{\beta }_2}{1\widehat{\beta }_1},\hfill \\ & v_i=\widehat{\alpha }_i+\widehat{\beta }_i,z_i=(1\widehat{\alpha }_i)(1\widehat{\beta }_i),z_{ij}=(1\widehat{\alpha }_i\widehat{\alpha }_j)(1\widehat{\beta }_i\widehat{\beta }_j),\hfill \\ & \chi _2(u;a,b)\frac{1}{4}(1u)^2\left[(1a)^2+(1b)^2\right],\hfill \\ & \mathrm{\Delta }_{II}^{(0)}=0,\mathrm{\Delta }_{II}^{(1)}(z_i,z_{ij})=\frac{1}{2}\gamma \frac{1}{6}\gamma \mathrm{ln}(z_i)\frac{1}{6}\gamma \mathrm{ln}(z_{ij}).\hfill \end{array}$$
(18)
The variables $`\widehat{\alpha }_i,\widehat{\beta }_i`$ for $`i`$-th photon are defined as in eq. (12).
In order to understand our construction let us examine how the LL collinear limit is realized in the exact single bremsstrahlung matrix element of eq. (12). If the photon carrying the fraction $`x_1`$ of the beam energy is collinear, let us say, with $`p_1`$ then $`\widehat{\alpha }_1x`$, $`\widehat{\beta }_10`$, all four angles are the same $`\vartheta _{sr}\vartheta ^{}`$ and we recover immediately the correct LL formula
$$\frac{1}{2}(1\widehat{\alpha }_1)^2\underset{r=1,2}{}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(\vartheta _{1r})+\frac{1}{2}(1\widehat{\beta }_1)^2\underset{r=1,2}{}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(\vartheta _{2r})\frac{1}{2}(1+(1x)^2)\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(\vartheta ^{}).$$
It is therefore natural to employ for the double emission the angular dependent Altarelli-Parisi (AP) factors of the kind
$$\frac{1}{2}[(1\widehat{\alpha }_2)^2+(1\widehat{\beta }_2)^2]\frac{1}{2}[(1\widehat{\alpha }_1)^2+(1\widehat{\beta }_1)^2].$$
The above formula is too simple, however, to reproduce correctly the result of the double convolution of the AP kernels in the case when both photons are collinear with the same fermion
$$\frac{1}{2}(1+(1x_1)^2)\frac{1}{2}(1+(1[x_2/(1x_1)]^2))\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(\vartheta ^{}).$$
where $`x_2^{}=x_1/(1x_1)`$ reflects the loss of energy in the emission cascade due to emission of $`k_1`$. In order to match the above cascade limit we construct a better angular dependent AP factor as
$$\frac{1}{2}[(1\widehat{\alpha }_1)^2+(1\widehat{\beta }_1)^2]\frac{1}{2}[(1\widehat{\alpha }_2^{})^2+(1\widehat{\beta }_2^{})^2].$$
The above fulfils both types of LL collinear limit, when two photons are collinear with a single beam or each of them follows different beam. Finally, let us reproduce the limit in which one photon, let us say the 1-st, is hard and the other, the 2-nd, is soft, $`v_2=\widehat{\alpha }_2+\widehat{\beta }_20`$. In this case it is logical to split the above double bremsstrahlung angular dependent AP factor into two pieces
$$\begin{array}{cc}\hfill \chi _2(\widehat{\alpha }_1;\widehat{\alpha }_2^{},\widehat{\beta }_2^{})& =\frac{1}{2}(1\widehat{\alpha }_1)^2\frac{1}{2}[(1\widehat{\alpha }_2^{})^2+(1\widehat{\beta }_2^{})^2],\hfill \\ \hfill \chi _2(\widehat{\beta }_1;\widehat{\alpha }_2^{},\widehat{\beta }_2^{})& =\frac{1}{2}(1\widehat{\beta }_1)^2\frac{1}{2}[(1\widehat{\alpha }_2^{})^2+(1\widehat{\beta }_2^{})^2]\hfill \end{array}$$
and associate each one with the corresponding $`d\sigma ^{\mathrm{Born}}/d\mathrm{\Omega }`$, following eq. (12). The order in the cascade does not matter. We simply symmetrize over the two orderings in the cascade – it is essentially Bose-Einstein symmetrization.
The above construction clearly provides the correct limit $`D_{II}^{(2)}(k_1,k_2)\stackrel{~}{S}(k_2)D_{1I}^{(1)}(k_2)`$ for $`v_1=const`$ and $`v_20`$. As a consequence $`\overline{\beta }_{II}^{(2)}(X,p_1,p_2,q_1,q_2,k_1,k_2)`$ is finite in the limit of one or both photon momenta tending to zero.
The construction of eq. (17) will be inadequate if both photons are hard and at least one has high transverse momentum. It reflects the fact that we do not control fully in EEX the second-order NLL, $`𝒪(\alpha ^2L)`$, contributions. However, we have known since a long time that the construction of the type of eq. (17) agrees rather well with the exact double bremsstrahlung matrix element calculated using spinor techniques, see . For both photons having high transverse momenta there is only about 20% disagreement for the approximate and exact results (integrated over the double photon phase-space). This result is confirmed in the present work by the numerical comparisons of EEX and CEEX, where the double bremsstrahlung matrix element is exact.
The double final-state bremsstrahlung distribution is defined/constructed in an analogous way
$$\begin{array}{cc}& D_{FF}^{(r)}(X,p_1,p_2,q_1,q_2,k_1^{},k_2^{})=\hfill \\ & Q_f^4\frac{\alpha }{4\pi ^2}\frac{2q_1p_2}{(k_1^{}q_1)(k_1^{}p_2)}\frac{\alpha }{4\pi ^2}\frac{2q_1p_2}{(k_2^{}q_1)(k_2^{}p_2)}W_f(\widehat{\eta }_1,\widehat{\zeta }_1)W_f(\widehat{\eta }_2,\widehat{\zeta }_2)\hfill \\ & \{\mathrm{\Theta }(v_1^{}v_2^{})[\chi _2(\eta _1;\eta _2^{},\zeta _2^{})\underset{r=1,2}{}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(X^2,\vartheta _{1r})+\chi _2(\zeta _1;\eta _2^{},\zeta _2^{})\underset{r=1,2}{}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(X^2,\vartheta _{2r})]\hfill \\ & +\mathrm{\Theta }(v_2^{}v_1^{})[\chi _2(\eta _2;\eta _1^{},\zeta _1^{})\underset{r=1,2}{}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(X^2,\vartheta _{1r})+\chi _2(\zeta _2;\eta _1^{},\zeta _1^{})\underset{r=1,2}{}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(X^2,\vartheta _{2r})]\}\hfill \\ & \left(1+\mathrm{\Delta }_I^{(r1)}(z_j)\right),\hfill \end{array}$$
(19)
where
$$\eta _1^{}=\frac{\eta _1}{1+\eta _2},\eta _2^{}=\frac{\eta _2}{1+\eta _1},\zeta _1^{}=\frac{\zeta _1}{1+\zeta _2},\zeta _2^{}=\frac{\zeta _2}{1+\zeta _1}.$$
(20)
The “primed” Sudakov variables are here defined differently than in the ISR case because the fermion momenta $`q_{1,2}`$ get affected by photon emission. Virtual corrections are absent because we restrict FSR to $`𝒪(\alpha ^2)`$<sub>LL</sub>. The above expression is tagged with $`r=2,3`$ for $`𝒪(\alpha ^r)`$, however, FSR we implement essentially only in $`𝒪(\alpha ^2)`$ and the only correction in $`𝒪(\alpha ^3)`$ is the ISR one-loop correction.
The distribution for one photon from the initial-state and one photon from the final-state at $`𝒪(\alpha ^r)`$ $`r=1,2`$ we construct as follows
$$\begin{array}{cc}\hfill D_{IF}^{(r)}& (X,p_1,p_2,q_1,q_2,k_j,k_l^{})=\hfill \\ & Q_e^2\frac{\alpha }{4\pi ^2}\frac{2p_1p_2}{(k_jp_1)(k_jp_2)}W_e(\widehat{\alpha }_j,\widehat{\beta }_j)Q_f^2\frac{\alpha }{4\pi ^2}\frac{2p_1p_2}{(k_l^{}p_1)(k_l^{}p_2)}W_f(\widehat{\eta }_l,\widehat{\zeta }_l)\hfill \\ & \{\frac{(1\widehat{\alpha }_j)^2}{2}\frac{(1\widehat{\eta }_l)^2}{2}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(X^2,\vartheta _{11})+\frac{(1\widehat{\alpha }_j)^2}{2}\frac{(1\widehat{\zeta }_l)^2}{2}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(X^2,\vartheta _{12})\hfill \\ & +\frac{(1\widehat{\beta }_j)^2}{2}\frac{(1\widehat{\eta }_l)^2}{2}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(X^2,\vartheta _{21})+\frac{(1\widehat{\beta }_j)^2}{2}\frac{(1\widehat{\zeta }_l)^2}{2}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(X^2,\vartheta _{22})\}\hfill \\ & \left(1+\mathrm{\Delta }_I^{(r1)}(z_1)\right)\left(1+\mathrm{\Delta }_F^{(r1)}(z_2^{})\right)\hfill \end{array}$$
(21)
where the variables $`\widehat{\alpha }_j,\widehat{\beta }_j,\widehat{\eta }_l,\widehat{\zeta }_l`$ and other components are defined as in eqs. (12,14). The above construction is in fact the easiest because two photons cannot be emitted in a cascade from one line and we fully exploit the four scattering angles in the Born differential cross sections. It is trivial to check that all soft and collinear limits are correct.
### 2.5 Three real photons
The differential distribution for of 3 real ISR photons is essentially obtained by the triple convolution of the AP kernel, for each beam separately and the the two results are combined with help of additional convolution. This exercise was done for the collinear sub-generator of BHLUMI and we exploit here these results. Even though the collinear limit is of primary importance, we have to be very careful in construction of the fully differential triple photon distribution to preserve all soft limits: when all three photons are soft, when two of them are soft , and only one of them is soft. In these limits the three-photon differential distribution has to reproduce smoothly the previously defined Born, single and double bremsstrahlung distributions times the appropriate soft factor(s). Otherwise we may have a problem with IR finiteness of
$$\begin{array}{cc}\hfill \overline{\beta }& {}_{III}{}^{(3)}(X,p_i,q_j,k_1,k_2,k_3)=D_{III}^{(r)}(X,p_i,q_j,k_1,k_2,k_3)\hfill \\ & \stackrel{~}{S}_I(k_1)\overline{\beta }_{1I}^{(2)}(X,p_i,q_j,k_2,k_3)\stackrel{~}{S}_I(k_2)\overline{\beta }_{1I}^{(2)}(X,p_i,q_j,k_1,k_3)\stackrel{~}{S}_I(k_3)\overline{\beta }_{1I}^{(2)}(X,p_i,q_j,k_1,k_2)\hfill \\ & \stackrel{~}{S}_I(k_1)\stackrel{~}{S}_I(k_2)\overline{\beta }_{1I}^{(1)}(X,p_i,q_j,k_3)\stackrel{~}{S}_I(k_3)\stackrel{~}{S}_I(k_1)\overline{\beta }_{1I}^{(1)}(X,p_i,q_j,k_2)\hfill \\ & \stackrel{~}{S}_I(k_2)\stackrel{~}{S}_I(k_3)\overline{\beta }_{1I}^{(1)}(X,p_i,q_j,k_1)\stackrel{~}{S}_I(k_1)\stackrel{~}{S}_I(k_2)\stackrel{~}{S}_I(k_3)\overline{\beta }_0^{(0)}(X,p_i,q_j).\hfill \end{array}$$
(22)
It is therefore not completely straightforward to turn the strictly collinear expression for three real photon distributions of ref. into the fully differential (finite $`p_T`$) triple photon distribution which we need. As in the case of double real ISR the guiding principle is that (i) the hardest photon decides which of the angles is used in $`d\sigma ^{\mathrm{Born}}/d\mathrm{\Omega }(X^2,\vartheta _{lr})`$ and (ii) we have to perform Bose symmetrization, that is sum over all orderings in a cascade emission of several photons from one beam. For three real photons there are no virtual corrections.
Our construction in the case of the triple real ISR reads as follows
$$\begin{array}{cc}\hfill D_{II}^{(3)}(X,p_1,p_2,q_1,q_2,k_1,k_2,k_3)& \underset{l=1,3}{}Q_e^2\frac{\alpha }{4\pi ^2}\frac{2p_1p_2}{(k_lp_1)(k_lp_2)}W_e(\widehat{\alpha }_l,\widehat{\beta }_l)\hfill \\ \hfill \left\{\mathrm{\Theta }(v_1v_2)\mathrm{\Theta }(v_2v_3)\right[& \chi _3(\widehat{\alpha }_1;\widehat{\alpha }_2^{},\widehat{\beta }_2^{},\widehat{\alpha }_3^{\prime \prime },\widehat{\beta }_3^{\prime \prime })\underset{r=1,2}{}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(X^2,\vartheta _{1r})\hfill \\ & +\chi _3(\widehat{\beta }_1;\widehat{\alpha }_2^{},\widehat{\beta }_2^{},\widehat{\alpha }_3^{\prime \prime },\widehat{\beta }_3^{\prime \prime })\underset{r=1,2}{}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}(X^2,\vartheta _{2r})]\hfill \\ \hfill +\mathrm{remaining}& \mathrm{five}\mathrm{permutations}\mathrm{of}(1,2,3)\},\hfill \end{array}$$
(23)
where
$$\begin{array}{cc}& \chi _3(u_1;a_2,b_2,a_3,b_3)\frac{1}{8}(1u_1)^2\left[(1a_2)^2+(1b_2)^2\right]\left[(1a_3)^2+(1b_3)^2\right],\hfill \\ & \widehat{\alpha }_3^{\prime \prime }=\frac{\widehat{\alpha }_3}{1\widehat{\alpha }_1\widehat{\alpha }_2},\widehat{\beta }_3^{\prime \prime }=\frac{\widehat{\beta }_3}{1\widehat{\beta }_1\widehat{\beta }_2},\hfill \end{array}$$
(24)
In most cases such an approach should be enough; however, in some special cases with two hard photons explicitly tagged it may not be sufficient. We have programmed and run special tests (unpublished) relying on the up to 3 hard-photon ISR amplitudes constructed with the methods similar to these in ref., in order to get additional confidence in the approximate real emission distrubutions presented in this Section.
## 3 Amplitudes for Coherent Exclusive Exponentiation
The Coherent Exclusive Exponentiation (CEEX) was introduced for the first time in ref. . It is deeply rooted in the Yennie-Frautschi-Suura (YFS) exponentiation . It applies in particular to processes with narrow resonances where it is related also to works of Greco et.al. . The exponentiation procedure, that is a reorganisation of the QED perturbative series such that infrared (IR) divergences are summed up to infinite-order is done at the spin-amplitude level for both real and virtual IR singularities. This is to be contrasted with traditional YFS exponentiation, on which our EEX is based, where isolating the real IR divergences is done for squared spin-summed spin amplitudes, that is for differential distributions and spin density matrices<sup>2</sup><sup>2</sup>2 The realization of EEX for spin density matrices is an obvious generalisation of the EEX/YFS exponentiation which, however, was never fully implemented in practice..
Our calculations of the spin amplitudes for fermion pair production in electron positron scattering is done with the help of the powerful Weyl spinor (WS) techniques. There are several variants of WS techniques. We have opted for the method of Kleiss and Stirling (KS) , which we found the best suited for our CEEX. In ref. the KS spinor technique for massless and massive fermions was reviewed and appended with the rules for controlling their complex phases, or equivalently, the fermion rest frame (all three axes) in which the fermion spin is quantised – this is a critical point if we want to control fully the spin density matrix of the fermions. This fermion rest frame we call the GPS frame and the rule for finding it we call the GPS rule. For the sake of completeness we include definitions of the KS spinors, photon polarization vectors, and our GPS rules in Appendix A.
The very interesting feature of CEEX is that, although it is formulated entirely in terms of the spin-amplitudes, the IR cancellations in CEEX occur for the integrated cross sections (probabilities), as usual; in practice they are realised numerically. There is no contradiction in the above statement. In order to avoid any confusion on this point, we shall provide the new detailed proof of IR cancellations in CEEX scheme in one of the following subsections.
### 3.1 Master formula
Defining the Lorentz invariant phase-space as
$$𝑑\mathrm{Lips}_n(P;p_1,p_2,\mathrm{},p_n)=(2\pi )^4\delta (P\underset{i=1}{\overset{n}{}}p_i)\underset{i=1}{\overset{n}{}}\frac{d^3p}{(2\pi )^32p_i^0}$$
(25)
we write the CEEX total cross section for the process
$$e^{}(p_a)+e^+(p_b)f(p_c)+\overline{f}(p_d)+\gamma (k_1)+\gamma (k_2)+\mathrm{}+\gamma (k_n),n=0,1,2,\mathrm{},\mathrm{}$$
(26)
with polarized beams and decays of unstable final fermions being sensitive to fermion spin polarizations, following refs. , as follows:
$$\sigma ^{(r)}=\frac{1}{\mathrm{flux}(s)}\underset{n=0}{\overset{\mathrm{}}{}}𝑑\mathrm{Lips}_{n+2}(p_a+p_b;p_c,p_d,k_1,\mathrm{},k_n)\rho _{\mathrm{CEEX}}^{(r)}(p_a,p_b,p_c,p_d,k_1,\mathrm{},k_n)$$
(27)
where, in the CMS $`\mathrm{flux}(s)2s`$,
$$\begin{array}{cc}\hfill \rho _{\mathrm{CEEX}}^{(r)}& (p_a,p_b,p_c,p_d,k_1,k_2,\mathrm{},k_n)=\frac{1}{n!}e^{Y(\mathrm{\Omega };p_a,\mathrm{},p_d)}\overline{\mathrm{\Theta }}(\mathrm{\Omega })\underset{\sigma _i=\pm 1}{}\underset{\lambda _i,\overline{\lambda }_i=\pm 1}{}\hfill \\ & \underset{i,j,l,m=0}{\overset{3}{}}\widehat{\epsilon }_a^i\widehat{\epsilon }_b^j\sigma _{\lambda _a\overline{\lambda }_a}^i\sigma _{\lambda _b\overline{\lambda }_b}^j𝔐_n^{(r)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2}\mathrm{}{}_{\sigma _n}{}^{k_n})\left[𝔐_n^{(r)}({}_{\overline{\lambda }}{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2}\mathrm{}{}_{\sigma _n}{}^{k_n})\right]^{}\sigma _{\overline{\lambda }_c\lambda _c}^l\sigma _{\overline{\lambda }_d\lambda _d}^m\widehat{h}_c^l\widehat{h}_d^m,\hfill \end{array}$$
(28)
and assuming domination of the $`s`$-channel exchanges, including resonances, the complete set of spin amplitudes for emission of $`n`$ photons we define in $`𝒪(\alpha ^r)_{\mathrm{CEEX}}`$ $`r=0,1,2`$ as follows:
$`𝔐_n^{(0)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}\mathrm{}{}_{\sigma _n}{}^{k_n})={\displaystyle \underset{\mathrm{}\{I,F\}^n}{}}{\displaystyle \underset{i=1}{\overset{n}{}}}𝔰_{[i]}^{\{\mathrm{}_i\}}\widehat{\beta }_0^{(0)}\left({}_{\lambda }{}^{p};X_{\mathrm{}}\right),`$
$`𝔐_n^{(1)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}\mathrm{}{}_{\sigma _n}{}^{k_n})={\displaystyle \underset{\mathrm{}\{I,F\}^n}{}}{\displaystyle \underset{i=1}{\overset{n}{}}}𝔰_{[i]}^{\{\mathrm{}_i\}}\{\widehat{\beta }_0^{(1)}\left({}_{\lambda }{}^{p};X_{\mathrm{}}\right)+{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \frac{\widehat{\beta }_{1\{\mathrm{}_j\}}^{(1)}\left({}_{\lambda }{}^{p}{}_{\sigma _j}{}^{k_j};X_{\mathrm{}}\right)}{𝔰_{[j]}^{\{\mathrm{}_j\}}}}\},`$
$`𝔐_n^{(2)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}\mathrm{}{}_{\sigma _n}{}^{k_n})=`$
$`={\displaystyle \underset{\mathrm{}\{I,F\}^n}{}}{\displaystyle \underset{i=1}{\overset{n}{}}}𝔰_{[i]}^{\{\mathrm{}_i\}}\{\widehat{\beta }_0^{(2)}\left({}_{\lambda }{}^{p};X_{\mathrm{}}\right)+{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \frac{\widehat{\beta }_{1\{\mathrm{}_j\}}^{(2)}\left({}_{\lambda }{}^{p}{}_{\sigma _j}{}^{k_j};X_{\mathrm{}}\right)}{𝔰_{[j]}^{\{\mathrm{}_j\}}}}+{\displaystyle \underset{1j<ln}{}}{\displaystyle \frac{\widehat{\beta }_{2\{\mathrm{}_j\mathrm{}_l\}}^{(2)}\left({}_{\lambda }{}^{p}{}_{\sigma _j}{}^{k_j}{}_{\sigma _l}{}^{k_l};X_{\mathrm{}}\right)}{𝔰_{[j]}^{\{\mathrm{}_j\}}𝔰_{[l]}^{(\mathrm{}_l)}}}\},`$
In the following subsections we shall explain all basic notation, then in the next section we shall discuss in detail the IR structure in CEEX, effectively deriving all the above formulas. At $`𝒪(\alpha ^r)`$ we have to provide for functions $`\widehat{\beta }_k^{(r)},k=0,1,\mathrm{},r`$ from Feynman diagrams, which are infrared-finite by construction . Their actual precise definitions will be given in the following. We shall define/calculate them explicitly up to $`𝒪(\alpha ^2)`$.
#### 3.1.1 Spin notation
In order to shorten our many formulas, we use a compact collective notations
$$({}_{\lambda }{}^{p})=({}_{\lambda _a}{}^{p_a}{}_{\lambda _b}{}^{p_b}{}_{\lambda _c}{}^{p_c}{}_{\lambda _d}{}^{p_d})$$
for fermion four-momenta $`p_A,A=a,b,c,d`$ (i.e., $`p_1=p_a,p_2=p_b,q_1=p_c,q_2=p_d`$) and helicities $`\lambda _A,A=a,b,c,d`$. For $`k=1,2,3`$, $`\sigma ^k`$ are Pauli matrices and $`\sigma _{\lambda ,\mu }^0=\delta _{\lambda ,\mu }`$ is the unit matrix. The components $`\widehat{\epsilon }_1^j,\widehat{\epsilon }_2^k,j,k=1,2,3`$ are the components of the conventional spin polarization vectors of $`e^{}`$ and $`e^+`$ respectively, defined in the so-called GPS fermion rest frames (see Appendix A and ref. for the exact definition of these frames). We define $`\widehat{\epsilon }_A^0=1`$ in a non-standard way (i.e. $`p_A\widehat{\epsilon }_A=m_e,A=a,b`$). The polarimeter vectors $`\widehat{h}_C`$ are similarly defined in the appropriate GPS rest frames of the final unstable fermions ($`p_C\widehat{h}_C=m_f,C=c,d`$). Note that, in general, $`\widehat{h}_C`$ may depend in a non-trivial way on the momenta of all decay products, see refs. for details. We did not introduce polarimeter vectors for bremsstrahlung photons, i.e. we take advantage of the fact that luckily all high-energy experiments are completely blind to photon spin polarizations.
#### 3.1.2 IR regulators and YFS form-factor
Here we introduce/explain our notation for IR integration limits for real photons in eqs. (27) and (28) and in the following sections. In general, the factor $`\overline{\mathrm{\Theta }}(\mathrm{\Omega })`$ in eq. (27) defines the infrared (IR) integration limits for all real photons. More precisely for a single photon, $`\mathrm{\Omega }`$ is the domain surrounding the IR divergence point $`k=0`$, which is in fact excluded from the MC phase-space. In CEEX there is no real distinction among ISR and FSR photons, $`\mathrm{\Omega }`$ is therefore necessarily the same for all photons. We define a characteristic function $`\mathrm{\Theta }(\mathrm{\Omega },k)`$ of the IR domain $`\mathrm{\Omega }`$ as $`\mathrm{\Theta }(\mathrm{\Omega },k)=1`$ for $`k\mathrm{\Omega }`$ and $`\mathrm{\Theta }(\mathrm{\Omega },k)=0`$ for $`k\mathrm{\Omega }`$. The characteristic function for the part of the phase-space included in the MC integration for a single real photon is $`\overline{\mathrm{\Theta }}(\mathrm{\Omega },k)=1\mathrm{\Theta }(\mathrm{\Omega },k)`$. The analogous characteristic function for all real photons is, of course, the following product
$$\overline{\mathrm{\Theta }}(\mathrm{\Omega })=\underset{i=1}{\overset{n}{}}\overline{\mathrm{\Theta }}(\mathrm{\Omega },k).$$
(29)
In the present calculation corresponding to the $`𝒦𝒦`$ Monte Carlo program we opt for $`\mathrm{\Omega }`$ defined traditionally with the photon energy cut condition $`k^0<E_{\mathrm{min}}`$.
The YFS form factor for $`\mathrm{\Omega }`$ defined with the condition $`k^0<E_{\mathrm{min}}`$ reads
$$\begin{array}{cc}& Y(\mathrm{\Omega };p_a,\mathrm{},p_d)=Q_e^2Y_\mathrm{\Omega }(p_a,p_b)+Q_f^2Y_\mathrm{\Omega }(p_c,p_d)\hfill \\ & +Q_eQ_fY_\mathrm{\Omega }(p_a,p_c)+Q_eQ_fY_\mathrm{\Omega }(p_b,p_d)Q_eQ_fY_\mathrm{\Omega }(p_a,p_d)Q_eQ_fY_\mathrm{\Omega }(p_b,p_c).\hfill \end{array}$$
(30)
where
$$\begin{array}{cc}\hfill Y_\mathrm{\Omega }(p,q)& 2\alpha \stackrel{~}{B}(\mathrm{\Omega },p,q)+2Q_f^2\alpha \mathrm{}B(p,q)\hfill \\ \hfill & 2\alpha \frac{1}{8\pi ^2}\frac{d^3k}{k^0}\mathrm{\Theta }(\mathrm{\Omega };k)\left(\frac{p}{kp}\frac{q}{kq}\right)^2\hfill \\ & +2\alpha \mathrm{}\frac{d^4k}{k^2}\frac{i}{(2\pi )^3}\left(\frac{2pk}{2kpk^2}\frac{2qk}{2kqk^2}\right)^2\hfill \end{array}$$
(31)
is given analytically in terms of logarithms and Spence functions. As we see, the above YFS form factor includes terms due to the initial-final state interference (IFI). The above form-factor will be derived in the following. The additional contribution to the YFS form-factor due to the narrow $`Z`$-resonance will be discussed in detail separately.
#### 3.1.3 Partitions and $`𝔰`$-factors
The coherent sum is taken over the set $`\{\mathrm{}\}=\{I,F\}^n`$ of all $`2^n`$ partitions – the single partition $`\mathrm{}`$ is defined as a vector $`(\mathrm{}_1,\mathrm{}_2,\mathrm{},\mathrm{}_n)`$ where $`\mathrm{}_i=I`$ for an ISR photon and $`\mathrm{}_F=F`$ for a FSR photon, see the analogous construction in refs. . The set of all partitions is explicitly the following
$$\{\mathrm{}\}=\{(I,I,I,\mathrm{},I),(F,I,I,\mathrm{},I),(I,F,I,\mathrm{},I),(F,F,I,\mathrm{},I),\mathrm{}(F,F,F,\mathrm{},F)\}.$$
The $`s`$-channel four-momentum in the (possibly) resonant $`s`$-channel propagator is $`X_{\mathrm{}}=p_a+p_b\underset{\mathrm{}_i=I}{}k_i.`$
The soft (eikonal) amplitude factors $`𝔰_{[i]}^{(\omega )},\omega =I,F`$, are complex numbers and they are defined as follows
$`𝔰_{[i]}^{\{I\}}𝔰_{\sigma _i}^{\{I\}}(k_i)=eQ_e{\displaystyle \frac{b_\sigma (k,p_a)}{2k_ip_a}}+eQ_e{\displaystyle \frac{b_\sigma (k_i,p_b)}{2k_ip_b}},\left|𝔰_{[i]}^{\{I\}}\right|^2={\displaystyle \frac{e^2Q_e^2}{2}}\left({\displaystyle \frac{p_a}{k_ip_a}}{\displaystyle \frac{p_b}{k_ip_b}}\right)^2,`$ (32)
$`𝔰_{[i]}^{\{F\}}𝔰_{\sigma _i}^{\{F\}}(k_i)=+eQ_f{\displaystyle \frac{b_\sigma (k_i,p_c)}{2kp_c}}eQ_f{\displaystyle \frac{b_\sigma (k_i,p_d)}{2k_ip_d}},\left|𝔰_{[i]}^{\{F\}}\right|^2={\displaystyle \frac{e^2Q_f^2}{2}}\left({\displaystyle \frac{p_c}{k_ip_c}}{\displaystyle \frac{p_d}{k_ip_d}}\right)^2,`$
$`b_\sigma (k,p)=\sqrt{2}{\displaystyle \frac{\overline{u}_\sigma (k)\overline{)}p𝔲_\sigma (\zeta )}{\overline{u}_\sigma (k)𝔲_\sigma (\zeta )}}=\sqrt{2}\sqrt{{\displaystyle \frac{2\zeta p}{2\zeta k}}}s_\sigma (k,\widehat{p}),`$
see also Appendix A for more details. As indicated above, the moduli squared of the CEEX soft factors coincide up to a normalization constant with the corresponding EEX real photon soft factors $`\stackrel{~}{S}(k_i)`$.
#### 3.1.4 Born
The simplest IR-finite $`\widehat{\beta }`$-function $`\widehat{\beta }_0^{(0)}`$ is just the Born spin amplitude times a certain kinematical factor (see the next subsection)
$$\widehat{\beta }_0^{(0)}\left({}_{\lambda }{}^{p};X\right)=𝔅\left({}_{\lambda }{}^{p};X\right)\frac{X^2}{(p_c+p_d)^2}.$$
(33)
The Born spin amplitude $`𝔅\left({}_{\lambda }{}^{p};X\right)`$ is a basic building block in the construction of all of our spin amplitudes – let us define it already at this point. The many equivalent notations for $`𝔅`$ will be introduced for flexibility – in view of its role as a basic building block in the calculation of the multi-bremsstrahlung amplitudes. Using Feynman rules and our basic massive spinors with definite GPS helicities of Appendix A, Born spin amplitudes for<sup>3</sup><sup>3</sup>3 For the moment we require $`fe`$. the $`e^{}(p_a)e^+(p_b)f(p_c)\overline{f}(p_d)`$ process are given by
$$\begin{array}{cc}& 𝔅\left({}_{\lambda }{}^{p};X\right)=𝔅\left({}_{\lambda _a}{}^{p_a}{}_{\lambda _b}{}^{p_b}{}_{\lambda _c}{}^{p_c}{}_{\lambda _d}{}^{p_d};X\right)=𝔅[{}_{\lambda _b}{}^{p_b}{}_{\lambda _a}{}^{p_a}][{}_{\lambda _c}{}^{p_c}{}_{\lambda _d}{}^{p_d}](X)=𝔅_{[bc][cd]}(X)=\hfill \\ & =ie^2\underset{B=\gamma ,Z}{}\mathrm{\Pi }_B^{\mu \nu }(X)(G_{e,\mu }^B)_{[ba]}(G_{f,\nu }^B)_{[cd]}H_B=\underset{B=\gamma ,Z}{}𝔅_{[bc][cd]}^B(X),\hfill \\ & (G_{e,\mu }^B)_{[ba]}\overline{v}(p_b,\lambda _b)G_{e,\mu }^Bu(p_a,\lambda _a),(G_{f,\mu }^B)_{[cd]}\overline{u}(p_c,\lambda _c)G_{f,\mu }^Bv(p_d,\lambda _d),\hfill \\ & G_{e,\mu }^B=\gamma _\mu \underset{\lambda =\pm }{}\omega _\lambda g_\lambda ^{B,e},G_{f,\mu }^B=\gamma _\mu \underset{\lambda =\pm }{}\omega _\lambda g_\lambda ^{B,f},\omega _\lambda =\frac{1}{2}(1+\lambda \gamma _5),\hfill \\ & \mathrm{\Pi }_B^{\mu \nu }(X)=\frac{g^{\mu \nu }}{X^2M_{B}^{}{}_{}{}^{2}+i\mathrm{\Gamma }_BX^2/M_B},\hfill \end{array}$$
(34)
where $`g_\lambda ^{B,f}`$ are the usual chiral ($`\lambda =+1,1=R,L`$) coupling constants of the vector boson $`B=\gamma ,Z`$ to fermion $`f`$ in units of the elementary charge $`e`$. If not specified otherwise, the “hook function” $`H_B`$ is trivial $`H_\gamma =H_Z=1`$. It will be used to introduce special effects into Born spin amplitudes, like running coupling constants or an additional form-factor due to a narrow resonance.
Spinor products are reorganized with the help of the Chisholm identity, see eq (207) in the Appendix A, which applies assuming that electron spinors are massless, and the inner product of eq. (204), also in the Appendix A:
$$𝔅_{[ba][cd]}^B(X)=2ie^2\frac{\delta _{\lambda _a,\lambda _b}\left[g_{\lambda _a}^{B,e}g_{\lambda _a}^{B,f}T_{\lambda _c\lambda _a}T_{\lambda _b\lambda _d}^{}+g_{\lambda _a}^{B,e}g_{\lambda _a}^{B,f}U_{\lambda _c\lambda _b}^{}U_{\lambda _a\lambda _d}\right]}{X^2M_{B}^{}{}_{}{}^{2}+i\mathrm{\Gamma }_BX^2/M_B},$$
(35)
where
$$\begin{array}{cc}\hfill T_{\lambda _c\lambda _a}=& \overline{u}(p_c,\lambda _c)u(p_a,\lambda _a)=S(p_c,m_c,\lambda _c,p_a,0,\lambda _a),\hfill \\ \hfill T_{\lambda _b\lambda _d}^{}=& \overline{v}(p_b,\lambda _b)v(p_d,\lambda _d)=S(p_b,0,\lambda _b,p_d,m_d,\lambda _d),\hfill \\ \hfill U_{\lambda _c\lambda _b}^{}=& \overline{u}(p_c,\lambda _c)v(p_b,\lambda _b)=S(p_c,m_c,\lambda _c,p_b,0,\lambda _b),\hfill \\ \hfill U_{\lambda _a\lambda _d}=& \overline{u}(p_a,\lambda _a)v(p_d,\lambda _d)=S(p_a,0,\lambda _a,p_d,m_d,\lambda _d).\hfill \end{array}$$
(36)
Note that the use of the Chisholm identity is a technical detail which should not obscure the generality of our approach. What we need in practice is any numerical method of evaluation of the Born spin amplitudes defined in eq. (34), and Chisholm identity is just one possibility.
#### 3.1.5 Off-space extrapolation
In eq. (3.1) Born spin amplitudes are obviously used for $`p_i`$ which do not necessarily obey the four-momentum conservation $`p_a+p_b=p_c+p_d`$. In the exclusive exponentiation this is natural and necessary because, in the presence of the bremsstrahlung photons, the relation $`X=p_a+p_b=p_c+p_d`$ may not hold. In eq. (3.1) only fermion momenta enter as an argument of the Born spin amplitudes. Photon momenta play only an indirect role, they disturb fermion momenta through energy and momentum conservation (sometimes referred to as a “recoil effect”). The self-suggesting questions are: Is this acceptable? Is this dangerous? Can this be avoided? The clear answer is: It is unavoidable and natural feature of the exclusive exponentiation that certain scattering matrix elements originally defined within $`n`$-body phase-space are in fact used in the phase-space with more particles. Let us call it off-space extrapolation, analogously to off-shell extrapolation<sup>4</sup><sup>4</sup>4In the off-shell case particles do not obey $`p^2=m^2`$, here we also modify the dimension of the phase-space.. It surely makes sense, and in principle is not dangerous, provided it is done with a little bit of care.
A technical remark: In the actual calculations of the multiphoton spin amplitudes fermion momenta $`p_i`$ in eq. (35) may be replaced, and occasionally will be replaced, by the momentum $`k`$ of one of the photons. This will be due to purely technical reasons (specific to the method of calculating multiphoton spin amplitudes). In such a case, the spinor into which $`k`$ enters as an argument is always understood to be massless.
#### 3.1.6 Pseudo-flux factor
One demonstration of the “off-space extrapolation” is the presence of the auxiliary factor $`F=X_{\mathrm{}}/(p_c+p_d)^2`$. In the framework of CEEX, its presence is not really mandatory and it disappears in the “in-space” situation $`p_a+p_b=p_c+p_d`$. In other words, the $`F`$-factor does not affect the soft limit; it really matters if at least one hard FSR photon is present. It is not related to narrow resonances, but rather to the leading-logarithmic (LL) structure of the higher-orders. Nevertheless, the $`F`$-factor is useful, because it is already implicitly present in the photon emission matrix element at $`𝒪(\alpha ^1)`$ and in all higher-orders, as can be seen in the LL approximation. It is therefore natural to include it at the early stage, already in the $`𝒪(\alpha ^0)`$ exponentiation. If we do not include it at the $`𝒪(\alpha ^0)`$ then it will be included order by order anyway. However, in such a case, the convergence of perturbative expansion will be deteriorated. As we shall see below, the introduction of the $`F`$-factor will slightly complicate the higher-order exponentiation and construction of the $`\widehat{\beta }`$ functions, but the gain is worth the effort. Furthermore, the $`F`$-factor has also been always present in the “crude distribution” in the YFS-type Monte Carlo generators, see for instance ref. , so it also improves the variance of the MC weight, especially for $`𝒪(\alpha ^0)_{\mathrm{CEEX}}`$.
### 3.2 IR structure in CEEX
Let us discuss in detail the origin of the $`𝒪(\alpha ^r)_{\mathrm{CEEX}}`$ expressions eqs. (27-28) and the mechanism of the IR cancellations. Our real starting point is the infinite order perturbative expression for the total cross section given by the standard quantum-mechanical expression of the type “matrix element squared modulus times phase-space”
$$\sigma ^{(\mathrm{})}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}d\tau _n(p_a+p_b;p_c,p_d,k_1,\mathrm{},k_n)\frac{1}{4}\underset{\lambda ,\sigma _i,\mathrm{},\sigma _n=\pm }{}|_n({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2}\mathrm{}{}_{\sigma _n}{}^{k_n})|^2,$$
(37)
where $`d\tau _n`$ is the respective $`n\gamma +2f`$ Lorentz invariant phase space, and $`_n`$ are the corresponding spin amplitudes. To simplify the discussion we take the unpolarized case, without narrow resonances.
#### 3.2.1 IR virtual factorization to infinite-order
According to the Yennie-Frautschi-Suura fundamental factorization theorem , all virtual IR corrections can be re-located into an exponential form-factor<sup>5</sup><sup>5</sup>5 In the LL approximation it is, of course, the doubly-logarithmic Sudakov form-factor. order by order and in infinite order
$$_n^{(\mathrm{})}=e^{\alpha B_4(p_a,p_b,p_c,p_d)}𝔐_n^{(\mathrm{})}.$$
(38)
As the convergence of the perturbative series is questionable, the above equation is in practice treated as a symbolic representation of the order-by-order relation which at $`𝒪(\alpha ^r)`$ reads
$$_n^{(r)}=\underset{l=0}{\overset{rn}{}}\frac{(\alpha B_4)^{rl}}{(rl)!}𝔐_n^{[l+n]},(nr),$$
(39)
where the index $`l`$ is the number of loops in $`𝔐_n^{[l+n]}`$. The above identity is quite powerful because $`𝔐_n^{[l+n]}`$ are not only free of the virtual IR-divergences, but are also universal: they are the same in every perturbative order $`r`$ – for example for one photon, the one-loop (IR-subtracted) component, $`𝔐_1^{(1)}`$, is the same in the fifth-order and, let us say, in the second order, where it appears for the first time. The above identity can also be reformulated as follows
$$𝔐_n^{(r)}=\underset{l=0}{\overset{rn}{}}𝔐_n^{[l+n]}=\left[e^{\alpha B_4(p_a,p_b,p_c,p_d)}_n^{(r)}\right]|_{𝒪(\alpha ^r)},$$
(40)
where, $`_n^{(r)}`$ has to be calculated from Feynman diagrams in at least<sup>6</sup><sup>6</sup>6 The use of $`𝔐_n^{(r+m)}`$ at $`𝒪(\alpha ^{(r+m)})`$, $`m>0`$ will yield the same result – this is another way of stating the universality property. $`𝒪(\alpha ^r)`$. The above steps are exactly the same as in .
The YFS form-factor $`B_4`$ for $`e^{}(p_a)+e^+(p_b)f(p_c)+\overline{f}(p_d)+n\gamma `$ reads
$$\begin{array}{cc}& \alpha B_4(p_a,p_b,p_c,p_d)=\frac{d^4k}{k^2m_\gamma ^2+iϵ}\frac{i}{(2\pi )^3}\left|J_I(k)J_F(k)\right|^2,\hfill \\ & J_I=eQ_e(\widehat{J}_a(k)\widehat{J}_b(k)),J_F=eQ_f(\widehat{J}_c(k)\widehat{J}_d(k)),\widehat{J}_f^\mu (k)=\frac{2p_f^\mu +k^\mu }{k^2+2kp_f+iϵ}.\hfill \end{array}$$
(41)
Using the identity $`(_kZ_kJ_k)^2=_{i>k}Z_iZ_k(J_iJ_k)^2`$, valid for $`Z_k=0`$, where $`Z_k`$ is the charge or minus charge of the particle in the initial- or final-state respectively, we may cast (see ref. ) $`B_4`$ into a sum of the simpler dipole components
$`\begin{array}{cc}& B_4(p_a,p_b,p_c,p_d)=Q_e^2B_2(p_a,p_b)+Q_f^2B_2(p_c,p_b)\hfill \\ & +Q_eQ_fB_2(p_a,p_c)+Q_eQ_fB_2(p_b,p_d)Q_eQ_fB_2(p_a,p_d)Q_eQ_fB_2(p_b,p_c),\hfill \end{array}`$ (42)
$`B_2(p_i,p_j){\displaystyle \frac{d^4k}{k^2m_\gamma ^2+iϵ}\frac{i}{(2\pi )^3}\left(\widehat{J}(p_i,k)\widehat{J}(p_j,k)\right)^2}.`$
In the above we assume that IR singularities are regularized with a finite photon mass $`m_\gamma `$ which enters into all $`B_2`$’s and implicitly into $`𝔰`$-factors (and in the real photon phase-space integrals, see the following discussion).
#### 3.2.2 IR real factorization to infinite-order
The next step is isolation of the real IR singularities and it is worth to elaborate on this point because here the CEEX method differs in essential details from the original YFS method . We use again results of the basic analysis of real IR singularities of ref. , the essential difference is that we do not square the amplitudes immediately – it is done numerically at the later stage. The validity of the whole basic analysis of the IR cancellations in ref. remains, however, useful because it is done in terms of the currents
$$j_f^\mu (k)=\frac{2p_f^\mu }{2p_fk},f=a,b,c,d.$$
(43)
The above currents are simply related to our $`𝔰`$-factors:
$$\begin{array}{cc}& 𝔰_\sigma ^{\{I\}}(k)=const\times Q_e(j_aj_b)ϵ_\sigma (\beta ),\hfill \\ & 𝔰_\sigma ^{\{F\}}(k)=const\times Q_f(j_cj_d)ϵ_\sigma (\beta ).\hfill \end{array}$$
(44)
It is important to remember that the whole structure of the real IR divergences is entirely controlled by the squares of the currents $`|j(k)|^2`$, for $`j=j_aj_b`$ or $`j=j_cj_d`$, independently whether we prefer to work with the amplitudes or their squares, because only the squares $`|j(k)|^2`$ are IR divergent and the other contractions do not matter (as was already stressed in ref. ). Similarly, if we express spin amplitudes in terms of $`𝔰`$-factors, only the squares $`|𝔰(k)|^2`$ are IR divergent and not the interference terms like $`\mathrm{}\{𝔰(k)(\mathrm{})^{}\}`$.
Having the above in mind we may proceed using results of ref. and we see that for instance the most IR divergent part of $`_n`$ is proportional to the products of $`n`$ $`𝔰`$-factors
$$\begin{array}{cc}\hfill 𝔐_n& ({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2}\mathrm{}{}_{\sigma _n}{}^{k_n})\widehat{\beta }_0\left({}_{\lambda }{}^{p};X\right)𝔰_{\sigma _1}(k_1)𝔰_{\sigma _2}(k_2)\mathrm{}𝔰_{\sigma _n}(k_n)\hfill \end{array}$$
(45)
where the function $`\widehat{\beta }_0`$ is not IR divergent any more, and we assumed for the moment the absence of the narrow resonances, using the sum of ISR and FSR $`𝔰`$-factors<sup>7</sup><sup>7</sup>7In the non-resonant case we may set $`X=p_a+p_b`$, for example.
$$𝔰_\sigma (k)𝔰_\sigma ^{\{F\}}(k)+𝔰_\sigma ^{\{I\}}(k).$$
(46)
However, there are also non-leading IR singularities. Suppressing inessential spin indices the whole real IR structure is revealed in the following decomposition :
$$\begin{array}{cc}\hfill 𝔐& {}_{n}{}^{(\mathrm{})}(k_1,k_2,k_3,\mathrm{},k_n)=\widehat{\beta }_0\underset{s=1}{\overset{n}{}}𝔰(k_s)+\underset{j=1}{\overset{n}{}}\widehat{\beta }_1(k_j)\underset{sj}{}𝔰(k_s)\hfill \\ & +\underset{j_1>j_2}{}\widehat{\beta }_2(k_{j_1},k_{j_2})\underset{sj_1,j_2}{}𝔰(k_s)+\underset{j_1>j_2>j_3}{}\widehat{\beta }_2(k_{j_1},k_{j_2},k_{j_3})\underset{sj_1,j_2,j_3}{}𝔰(k_s)+\mathrm{}\hfill \\ & +\underset{j=1}{\overset{n}{}}\widehat{\beta }_{n1}(k_1,\mathrm{}k_{j1},k_{j+1},\mathrm{},k_n)𝔰(k_j)+\widehat{\beta }_n(k_1,k_2,k_3,\mathrm{},k_n)\hfill \end{array}$$
(47)
where functions $`\widehat{\beta }_i`$ are IR free and include finite loop corrections to infinite-order. Let us stress that these functions $`\widehat{\beta }_i`$ are genuinely new objects. They were not used and even not considered in ref. .
#### 3.2.3 Finite-order $`\widehat{\beta }`$’s
The decomposition of eq. (47) has also its order-by-order representation, which at $`𝒪(\alpha ^r)`$, $`r=n+l`$, reads as follows:
$$\begin{array}{cc}\hfill 𝔐& {}_{n}{}^{(n+l)}(k_1,k_2,k_3,\mathrm{},k_n)=\widehat{\beta }_0^{(l)}\underset{s=1}{\overset{n}{}}𝔰(k_s)+\underset{j=1}{\overset{n}{}}\widehat{\beta }_1^{(1+l)}(k_j)\underset{sj}{}𝔰(k_s)\hfill \\ & +\underset{j_1<j_2}{}\widehat{\beta }_2^{(2+l)}(k_{j_1},k_{j_2})\underset{sj_1,j_2}{}𝔰(k_s)+\underset{j_1<j_2<j_3}{}\widehat{\beta }_2^{(3+l)}(k_{j_1},k_{j_2},k_{j_3})\underset{sj_1,j_2,j_3}{}𝔰(k_s)+\mathrm{}\hfill \\ & +\underset{j=1}{\overset{n}{}}\widehat{\beta }_{n1}^{(n1+l)}(k_1,\mathrm{}k_{j1},k_{j+1},\mathrm{},k_n)𝔰(k_j)+\widehat{\beta }_n^{(n+l)}(k_1,k_2,k_3,\mathrm{},k_n)\hfill \\ \hfill =& \underset{s=1}{\overset{n}{}}𝔰(k_s)\{\widehat{\beta }_0^{(l)}+\underset{j=1}{\overset{n}{}}\frac{\widehat{\beta }_1^{(1+l)}(k_j)}{𝔰(k_j)}+\underset{j_1<j_2}{}\frac{\widehat{\beta }_2^{(2+l)}(k_{j_1},k_{j_2})}{𝔰(k_{j_1})𝔰(k_{j_2})}+\underset{j_1<j_2<j_3}{}\frac{\widehat{\beta }_2^{(3+l)}(k_{j_1},k_{j_2},k_{j_3})}{𝔰(k_{j_1})𝔰(k_{j_2}𝔰(k_{j_3})}\hfill \\ & +\underset{j=1}{\overset{n}{}}\frac{\widehat{\beta }_{n1}^{(n1+l)}(k_1,\mathrm{}k_{j1},k_{j+1},\mathrm{},k_n)}{\underset{sj}{}𝔰(k_s)}+\frac{\widehat{\beta }_n^{(n+l)}(k_1,k_2,k_3,\mathrm{},k_n)}{\underset{s}{}𝔰(k_s)}\}.\hfill \end{array}$$
(48)
The new functions $`\widehat{\beta }_n^{(n+l)}(k_1,k_2,k_3,\mathrm{},k_n)`$ contain up to $`l`$-loop corrections, and are not only completely IR-finite, but are also universal: for instance the $`\widehat{\beta }_1^{(2)}(k)`$, which appears for the first time in decomposition of $`𝔐_1^{(2)}(k)`$, is functionally the same when decomposing $`𝔐_2^{(3)}(k_1,k_2)`$ or any higher-order $`𝔐_n^{(n+l)}`$. This feature is essential for reversing the relations of eq. (48), that is for practical order-by-order calculations of $`\widehat{\beta }_n^{(n+l)}`$ from $`𝔐_n^{(r)}`$, obtained directly from the Feynman rules:
$$\begin{array}{cc}& \widehat{\beta }_0^{(l)}=𝔐_0^{(l)},\hfill \\ & \widehat{\beta }_1^{(1+l)}(k_1)=𝔐_1^{(1+l)}(k_1)\widehat{\beta }_0^{(l)}𝔰(k_1),\hfill \\ & \widehat{\beta }_2^{(2+l)}(k_1,k_2)=𝔐_2^{(2+l)}(k_1,k_2)\widehat{\beta }_1^{(1+l)}(k_1)𝔰(k_2)\widehat{\beta }_1^{(1+l)}(k_2)𝔰(k_1)\widehat{\beta }_0^{(l)}𝔰(k_1)𝔰(k_2),\hfill \\ & \widehat{\beta }_3^{(3+l)}(k_1,k_2,k_3)=𝔐_3^{(3+l)}(k_1,k_2,k_3)\hfill \\ & \widehat{\beta }_2^{(2+l)}(k_1,k_2)𝔰(k_3)\widehat{\beta }_2^{(2+l)}(k_1,k_3)𝔰(k_2)\widehat{\beta }_2^{(2+l)}(k_2,k_3)𝔰(k_1)\hfill \\ & \widehat{\beta }_1^{(1+l)}(k_1)𝔰(k_2)𝔰(k_3)\widehat{\beta }_1^{(1+l)}(k_2)𝔰(k_1)𝔰(k_3)\widehat{\beta }_1^{(1+l)}(k_3)𝔰(k_1)𝔰(k_2)\hfill \\ & \widehat{\beta }_0^{(l)}𝔰(k_1)𝔰(k_2)𝔰(k_3),\mathrm{},\hfill \\ & \widehat{\beta }_n^{(n+l)}(k_1,\mathrm{},k_n)=𝔐_n^{(n+l)}(k_1,\mathrm{},k_n)\underset{j=1}{\overset{n}{}}\widehat{\beta }_{n1}^{(n1+l)}(k_1,\mathrm{}k_{j1},k_{j+1},\mathrm{},k_n)𝔰(k_j)\hfill \\ & \underset{j_1<j_2}{}\widehat{\beta }_{n2}^{(n2+l)}(k_1,\mathrm{}k_{j_11},k_{j_1+1},\mathrm{}k_{j_21},k_{j_2+1},\mathrm{},k_n)𝔰(k_{j_1})𝔰(k_{j_2})\mathrm{}\hfill \\ & \underset{j_1<j_2}{}\widehat{\beta }_2^{(1+l)}(k_{j_1},k_{j_2})\underset{sj_1,j_2}{}𝔰(k_s)\underset{j=1}{\overset{n}{}}\widehat{\beta }_1^{(1+l)}(k_j)\underset{sj}{}𝔰(k_s)\widehat{\beta }_0^{(l)}\underset{s=1}{\overset{n}{}}𝔰(k_s).\hfill \end{array}$$
(49)
The above set of equations is a recursive rule, i.e., higher-order $`\widehat{\beta }`$’s are constructed in terms of lower-order ones. In practical calculations we do not go to infinite-order but we stop at some $`𝒪(\alpha ^r)`$ and the above set of equations is truncated for $`\widehat{\beta }_n^{(n+l)}`$ by the requirement $`n+lr`$. The above truncation is harmless from the point of view of IR cancellations because we omit higher-order $`\widehat{\beta }`$’s which are IR-finite. As a consequence of the above fixed-order truncation eq. (47) takes the following form:
$$\begin{array}{cc}\hfill 𝔐& {}_{n}{}^{(r)}(k_1,k_2,k_3,\mathrm{},k_n)=\hfill \\ \hfill =& \underset{s=1}{\overset{n}{}}𝔰(k_s)\{\widehat{\beta }_0^{(r)}+\underset{j=1}{\overset{n}{}}\frac{\widehat{\beta }_1^{(r)}(k_j)}{𝔰(k_j)}+\underset{j_1<j_2}{}\frac{\widehat{\beta }_2^{(r)}(k_{j_1},k_{j_2})}{𝔰(k_{j_1})𝔰(k_{j_2})}+\underset{j_1<j_2<j_3}{}\frac{\widehat{\beta }_2^{(r)}(k_{j_1},k_{j_2},k_{j_3})}{𝔰(k_{j_1})𝔰(k_{j_2}𝔰(k_{j_3})}\hfill \\ & +\underset{j_1<j_2<\mathrm{}<j_r}{}\frac{\widehat{\beta }_r^{(r)}(k_{j_1},k_{j_2},\mathrm{},k_{j_r})}{𝔰(k_{j_1})𝔰(k_{j_2})\mathrm{}𝔰(k_{j_r})}\},\hfill \end{array}$$
(50)
where, contrary to eq. (48), we now allow only for $`r<n`$; in such a case the sum has $`r+1`$ terms instead of $`n`$.
The above formula represents the general finite-order $`𝒪(\alpha ^r)_{exp}`$ case while for $`r=0`$ only the first term survives, and in our $`𝒪(\alpha ^2)`$ case there are three terms. The CEEX spin amplitudes in our master formula eq. (3.1) represent the cases of $`r=0,1,2`$.
Just to give an explicit example, in the recursive calculation of $`\widehat{\beta }`$’s in $`𝒪(\alpha ^3)`$ we would need to calculate $`\widehat{\beta }_0^{(l)},l=0,1,2,3`$, $`\widehat{\beta }_1^{(1+l)},l=0,1,2`$, $`\widehat{\beta }_2^{(2+l)},l=0,1`$ and $`\widehat{\beta }_3^{(3)}`$. In the present work, at $`𝒪(\alpha ^r)`$, $`r=0,1,2`$, we shall employ the following set of recursive definitions based on eqs. (49)
$$\begin{array}{cc}& \widehat{\beta }_0^{(l)}({}_{\lambda }{}^{p})=𝔐_0^{(l)}({}_{\lambda }{}^{p}),l=0,1,2,\hfill \\ & \widehat{\beta }_1^{(1+l)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1})=𝔐_1^{(1+l)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1})\widehat{\beta }_0^{(l)}({}_{\lambda }{}^{p})𝔰_{\sigma _1}(k_1),l=0,1,\hfill \\ & \widehat{\beta }_2^{(2)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2})=𝔐_2^{(2)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2})\widehat{\beta }_1^{(1)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1})𝔰_{\sigma _2}(k_2)\widehat{\beta }_1^{(1)}({}_{\lambda }{}^{p}{}_{\sigma _2}{}^{k_2})𝔰_{\sigma _1}(k_1)\widehat{\beta }_0^{(0)}({}_{\lambda }{}^{p})𝔰_{\sigma _1}(k_1)𝔰_{\sigma _2}(k_2),\hfill \end{array}$$
(51)
where the $`𝔐`$-amplitude is given by eq. (40). Here we restored spin indices but we still specialize to the non-resonant case, and our $`\widehat{\beta }`$’s do not have the partition dependent $`X_{\mathrm{}}`$ argument as in $`\widehat{\beta }`$’s of eqs. (3.1-LABEL:eq:ceex-master2). We shall provide a definition for $`\widehat{\beta }`$’s in the resonant case in the following section 3.3.4.
#### 3.2.4 IR cancellations in CEEX
At fixed-order $`𝒪(\alpha ^r)_{\mathrm{CEEX}}`$, and remembering that $`|\mathrm{exp}(B_4)|^2=\mathrm{exp}(2\mathrm{}B_4)`$, we have obtained
$$\sigma ^{(r)}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}𝑑\tau _n(p_1+p_2;p_3,p_4,k_1,\mathrm{},k_n)e^{2\alpha \mathrm{}B_4(p_a,\mathrm{},p_d)}\frac{1}{4}\underset{\mathrm{spin}}{}\left|𝔐_n^{(r)}(k_1,k_2,\mathrm{}k_n)\right|^2,$$
(52)
where $`𝔐_n^{(r)}`$ is defined in eq. (50) and we factorize out the $`𝔰`$-factors
$$\begin{array}{cc}\hfill \frac{1}{4}\underset{\mathrm{spin}}{}\left|𝔐_n^{(r)}(k_1,k_2,k_3,\mathrm{},k_n)\right|^2& =d_n(k_1,k_2,k_3,\mathrm{},k_n)\underset{s=1}{\overset{n}{}}|𝔰(k_s)|^2,\hfill \\ \hfill d_n(k_1,k_2,k_3,\mathrm{},k_n)=|& \widehat{\beta }_0^{(r)}+\underset{j=1}{\overset{n}{}}\frac{\widehat{\beta }_1^{(r)}(k_j)}{𝔰(k_j)}+\underset{j_1<j_2}{}\frac{\widehat{\beta }_2^{(r)}(k_{j_1},k_{j_2})}{𝔰(k_{j_1})𝔰(k_{j_2})}+\underset{j_1<j_2<j_3}{}\frac{\widehat{\beta }_2^{(r)}(k_{j_1},k_{j_2},k_{j_3})}{𝔰(k_{j_1})𝔰(k_{j_2}𝔰(k_{j_3})}\hfill \\ & +\underset{j_1<j_2<\mathrm{}<j_r}{}\frac{\widehat{\beta }_r^{(r)}(k_{j_1},k_{j_2},\mathrm{},k_{j_r})}{𝔰(k_{j_1})𝔰(k_{j_2})\mathrm{}𝔰(k_{j_r})}|^2,\hfill \end{array}$$
(53)
In the above the function $`d_n(k_1,k_2,k_3,\mathrm{},k_n)`$ is IR-finite and we are allowed set $`m_\gamma 0`$ in it. Apart from $`2\alpha \mathrm{}B_4`$ the IR regulator $`m_\gamma `$ still remains in all $`𝔰(k_i)`$-factors and in the lower phase-space boundary of all real photons in $`d^3k/2k^0`$.
The above total cross section is perfectly IR-finite, as can be checked with a little bit of effort by analytical partial differentiation<sup>8</sup><sup>8</sup>8 This method of validating IR-finiteness was noticed by G. Burgers . The classical method of ref. relies on the techniques of the Melin transform, which could be also used here. with respect the photon mass
$$\begin{array}{cc}& \frac{}{m_\gamma }\sigma ^{(r)}=\hfill \\ & =\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}𝑑\tau _n(P;p_3,p_4,k_1,\mathrm{},k_n)e^{2\alpha \mathrm{}B_4}\frac{}{m_\gamma }\{2\alpha \mathrm{}B_4\}\frac{1}{4}\underset{\mathrm{spin}}{}\left|𝔐_n^{(r)}(k_1,k_2,\mathrm{}k_n)\right|^2\hfill \\ & +\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{s=1}{\overset{n}{}}𝑑\tau _{n1}(P;p_3,p_4,k_1,\mathrm{},k_{s1},k_{s+1},\mathrm{},k_n)e^{2\alpha \mathrm{}B_4}\hfill \\ & \times \frac{}{m_\gamma }\left\{\frac{d^3k_s}{2k_s^0}|𝔰(k_s)|^2\right\}\underset{js}{}|𝔰(k_j)|^2d_n(k_1,k_2,\mathrm{},k_s,\mathrm{},k_n)\hfill \end{array}$$
(54)
It is now necessary to notice that
$$\frac{}{m_\gamma }\left\{\frac{d^3k_s}{2k_s^0}|𝔰(k_s)|^2\right\}$$
is a $`\delta `$-like function concentrated at $`k_s=0`$ and we may therefore use the limit
$$d_n(k_1,\mathrm{},k_s,\mathrm{},k_n)d_n(k_1,k_2,\mathrm{},k_{s1},0,k_{s+1},\mathrm{},k_n)d_{n1}(k_1,k_2,\mathrm{},k_{s1},k_{s+1},\mathrm{},k_n)$$
The above helps us to notice that all terms in $`_{s=1}^n`$ are identical and we may sum them up, (after formally renaming the photon integration variables in the second integral) and rewrite eq. (54) as follows
$$\begin{array}{cc}\hfill \frac{}{m_\gamma }\sigma ^{(r)}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}𝑑\tau _n& (P;p_3,p_4,k_1,\mathrm{},k_n)e^{2\alpha \mathrm{}B_4}\frac{1}{4}\underset{\mathrm{spin}}{}\left|𝔐_n^{(r)}(k_1,k_2,\mathrm{}k_n)\right|^2\hfill \\ & \times \frac{}{m_\gamma }\{2\alpha \mathrm{}B_4+\frac{d^3k_s}{2k_s^0}|𝔰(k_s)|^2\}=0,\hfill \end{array}$$
(55)
where the independence on $`m_\gamma `$ of the sum of the 1-photon real and virtual integrals is due to the usual cancellation of the IR-divergences in the YFS scheme, shown explicitly many times.
The integral of eqs. (37) and (52) are perfectly implementable in the Monte Carlo form, with small $`m_\gamma `$ being the IR regulator, using a method very similar to that in ref. . Traditionally, however, the lower boundary on the real soft photons is defined using the energy cut condition $`k^0>\epsilon \sqrt{s}/2`$ in the laboratory frame. The practical advantage of such a cut is the lower photon multiplicity in the MC simulation, and consequently a faster computer program<sup>9</sup><sup>9</sup>9 The disadvantage of the cut $`k^0>\epsilon \sqrt{s}/2`$ is that in the MC it has to be implemented in different reference frames for ISR and for FSR – this costs the additional delicate procedure of bringing these two boundaries together, see ref. and/or discussion in the analogous $`t`$-channel case in ref. .. If the above energy cut on the photon energy is adopted, then the real soft-photon integral between the lower LIPS boundary defined by $`m_\gamma `$ and that defined by $`\epsilon `$ can be evaluated by hand and summed up rigorously (the only approximation is $`m_\gamma /m_e0`$) in the following.
#### 3.2.5 Explicit IR boundary for real photons
A general notation for the IR domain $`\mathrm{\Omega }`$ was already introduced, see eq. (29). Let us now exclude the $`\mathrm{\Omega }`$ domain from the real photon phase space (integrate out analytically). Splitting the real photon integration phase space we rewrite the eq. (52) as follows
$$\begin{array}{cc}\hfill \sigma ^{(r)}=& \underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{j=1}{\overset{n}{}}\left\{\frac{d^3k_j}{2k_j^0}|𝔰(k_j)|^2\mathrm{\Theta }(\mathrm{\Omega },k_j)+\frac{d^3k_j}{2k_j^0}|𝔰(k_j)|^2\overline{\mathrm{\Theta }}(\mathrm{\Omega },k_j)\right\}\hfill \\ & 𝑑\tau _2(P\underset{j=1}{\overset{n}{}}k_j;p_3,p_4)e^{2\alpha \mathrm{}B_4}d_n(k_1,k_2,\mathrm{},k_n).\hfill \end{array}$$
(56)
After expanding the binomial product into $`2^n`$ terms let us consider for instance the sum of all $`(\genfrac{}{}{0pt}{}{n}{1})=n`$ terms in which one photon is in $`\mathrm{\Omega }`$ and the other ones are not:
$$\begin{array}{cc}\hfill \frac{1}{n!}& \underset{s=1}{\overset{n}{}}\frac{d^3k_s}{2k_s^0}|𝔰(k_s)|^2\mathrm{\Theta }(\mathrm{\Omega },k_s)\underset{js}{\overset{n}{}}\frac{d^3k_j}{2k_j^0}|𝔰(k_j)|^2\overline{\mathrm{\Theta }}(\mathrm{\Omega },k_j)\hfill \\ & 𝑑\tau _2(P\underset{j=1}{\overset{n}{}}k_j;p_3,p_4)e^{2\alpha \mathrm{}B_4}d_n(k_1,k_2,\mathrm{},k_{s1},0,k_{s+1},\mathrm{},k_n)\hfill \\ \hfill =& \frac{1}{n!}\left(\genfrac{}{}{0pt}{}{n}{1}\right)\frac{d^3k}{2k^0}|𝔰(k)|^2\mathrm{\Theta }(\mathrm{\Omega },k)\hfill \\ & 𝑑\tau _{n+1}(P;p_3,p_4,k_1,k_2,\mathrm{},k_{n1})\underset{j=1}{\overset{n1}{}}\overline{\mathrm{\Theta }}(\mathrm{\Omega },k_j)|𝔰(k_j)|^2d_{n1}(k_1,k_2,\mathrm{},k_{n1})\hfill \end{array}$$
(57)
A similar summation is performed for the $`(\genfrac{}{}{0pt}{}{n}{s})`$ terms where $`s`$ photons are in $`\mathrm{\Omega }`$ giving rise to
$$\begin{array}{cc}& \sigma ^{(r)}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{s=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{s}\right)\left(\frac{d^3k}{2k^0}|𝔰(k)|^2\mathrm{\Theta }(\mathrm{\Omega },k)\right)^s\hfill \\ & 𝑑\tau _{2+ns}(P;p_3,p_4,k_1,k_2,\mathrm{},k_s)\underset{j=1}{\overset{ns}{}}\left\{|𝔰(k_j)|^2\overline{\mathrm{\Theta }}(\mathrm{\Omega },k_j)\right\}e^{2\alpha \mathrm{}B_4}d_{ns}(k_1,k_2,\mathrm{},k_{ns})\hfill \\ & =\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}𝑑\tau _{2+n}(P;p_3,p_4,k_1,k_2,\mathrm{},k_n)\mathrm{exp}\left(\frac{d^3k_j}{2k_j^0}|𝔰(k_j)|^2\mathrm{\Theta }(\mathrm{\Omega },k_j)\right)e^{2\alpha \mathrm{}B_4(p_1,\mathrm{},p_4)}\hfill \\ & \times \underset{j=1}{\overset{n}{}}\{|𝔰(k_j)|^2\overline{\mathrm{\Theta }}(\mathrm{\Omega },k_j)\}d_n(k_1,k_2,\mathrm{},k_n).\hfill \end{array}$$
(58)
The additional overall exponential factor contains the well known function
$$\begin{array}{cc}& \stackrel{~}{B}_4(p_1,\mathrm{},p_4)=\frac{d^3k_j}{2k_j^0}|𝔰(k_j)|^2\mathrm{\Theta }(\mathrm{\Omega },k_j)=Q_e^2\stackrel{~}{B}_2(p_1,p_2)+Q_f^2\stackrel{~}{B}_2(p_3,p_4)\hfill \\ & +Q_eQ_f\stackrel{~}{B}_2(p_1,p_3)+Q_eQ_f\stackrel{~}{B}_2(p_2,p_4)Q_eQ_f\stackrel{~}{B}_2(p_1,p_4)Q_eQ_f\stackrel{~}{B}_2(p_2,p_3),\hfill \\ & \stackrel{~}{B}_2(p,q)\frac{d^3k}{2k^0}\mathrm{\Theta }(\mathrm{\Omega },k_j)\left(j_p(k)j_q(k)\right)^2\frac{d^3k}{2k^0}\mathrm{\Theta }(\mathrm{\Omega },k_j)\frac{(1)}{8\pi ^2}\left(\frac{p}{kp}\frac{q}{kq}\right)^2,\hfill \end{array}$$
(59)
which forms together with $`2\alpha \mathrm{}B_4(p_1,\mathrm{},p_4)`$ the conventional YFS form-factor
$$Y(\mathrm{\Omega };p_1,\mathrm{},p_4)=2\alpha \stackrel{~}{B}_4(p_1,\mathrm{},p_4)+2\alpha \mathrm{}B_4(p_1,\mathrm{},p_4)$$
(60)
in our master eqs. (27,28). The dependence on $`m_\gamma `$ in $`Y`$ cancels out. Photon mass gets effectively replaced by the size of $`\mathrm{\Omega }`$ in its role of the IR regulator. The YFS form-factor $`Y`$ can be decomposed into six dipole components, see eq. (30) and can be calculated analytically in terms of logs and Spence functions, see refs. keeping all fermion masses exactly.
As already indicated, in the MC with the YFS exponentiation it would be possible to do without $`\mathrm{\Omega }`$ (declare it as empty) and rely uniquely on the IR regularization with a small photon mass $`m_\gamma `$ only . In such a case the formulas (31) for YFS form factor would include only the second virtual photon integral part.
For the sake of the completeness of the discussion it is necessary to examine once again the IR cancellations in the total cross section with $`\mathrm{\Omega }`$ as the new IR-regulator
$$\begin{array}{cc}\hfill \sigma ^{(r)}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}𝑑\tau _{2+n}& (P;p_3,p_4,k_1,k_2,\mathrm{},k_n)\underset{j=1}{\overset{n}{}}\left\{|𝔰(k_j)|^2\overline{\mathrm{\Theta }}(\mathrm{\Omega },k_j)\right\}\hfill \\ & \times e^{\stackrel{~}{B}_4(\mathrm{\Omega };p_1,\mathrm{},p_4)+2\alpha \mathrm{}B_4(p_1,\mathrm{},p_4)}d_n(k_1,k_2,\mathrm{},k_n).\hfill \end{array}$$
(61)
IR-finiteness of the total cross section now simply translates into independence on the $`\mathrm{\Omega }`$ domain, (assuming, as usual, that the size of $`\mathrm{\Omega }`$ is very small)
$$\frac{\delta }{\delta \mathrm{\Omega }}\sigma ^{(r)}=0.$$
(62)
The proof can be done along the same lines as the previous one for the photon mass. Let us assume that we want to vary $`\mathrm{\Omega }\mathrm{\Omega }^{}=\mathrm{\Omega }+\delta \mathrm{\Omega }`$, that is $`\overline{\mathrm{\Omega }}^{}=\overline{\mathrm{\Omega }}\delta \mathrm{\Omega }`$. Note that $`\mathrm{\Omega }^{}`$ can be much bigger or smaller than of $`\mathrm{\Omega }`$, the only requirement is that both are very small<sup>10</sup><sup>10</sup>10$`\delta \mathrm{\Omega }`$ does not need to be infinitesimal with respect to $`\mathrm{\Omega }`$. We proceed as follows
$$\begin{array}{cc}\hfill \sigma ^{(r)}& =\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{j=1}{\overset{n}{}}\left\{\frac{d^3k_j}{k_j^0}|𝔰(k_j)|^2\overline{\mathrm{\Theta }}(\mathrm{\Omega }^{},k_j)+\frac{d^3k_j}{k_j^0}|𝔰(k_j)|^2\mathrm{\Theta }(\delta \mathrm{\Omega },k_j)\right\}\hfill \\ & \times d\tau _2(Pk_j;p_3,p_4)e^{\stackrel{~}{B}_4(\mathrm{\Omega };p_1,\mathrm{},p_4)+2\alpha \mathrm{}B_4(p_1,\mathrm{},p_4)}d_n(k_1,k_2,\mathrm{},k_n)\hfill \\ & =\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{s=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{s}\right)\left\{\frac{d^3k}{k^0}|𝔰(k)|^2\mathrm{\Theta }(\delta \mathrm{\Omega },k)\right\}^s𝑑\tau _{2+ns}(P;p_3,p_4,k_1,\mathrm{},k_{ns})\hfill \\ & \times \underset{j=1}{\overset{ns}{}}\{|𝔰(k_j)|^2\overline{\mathrm{\Theta }}(\mathrm{\Omega }^{},k_j)\}e^{\stackrel{~}{B}_4(\mathrm{\Omega };p_1,\mathrm{},p_4)+2\alpha \mathrm{}B_4(p_1,\mathrm{},p_4)}d_{ns}(k_1,k_2,\mathrm{},k_{ns})\hfill \\ & =\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}𝑑\tau _{2+n}(P;p_3,p_4,k_1,\mathrm{},k_n)e^{{\scriptscriptstyle {\scriptscriptstyle \frac{d^3k}{k^0}}|𝔰(k)|^2\mathrm{\Theta }(\delta \mathrm{\Omega },k)}+\stackrel{~}{B}_4(\mathrm{\Omega };p_1,\mathrm{},p_4)+2\alpha \mathrm{}B_4(p_1,\mathrm{},p_4)}\hfill \\ & \times \underset{j=1}{\overset{n}{}}\{|𝔰(k_j)|^2\overline{\mathrm{\Theta }}(\mathrm{\Omega }^{},k_j)\}d_n(k_1,k_2,\mathrm{},k_n).\hfill \end{array}$$
(63)
recovering the same expression as (61), but with $`\mathrm{\Omega }^{}`$ instead of $`\mathrm{\Omega }`$.
### 3.3 Narrow neutral resonance in CEEX
The main new feature of CEEX in comparison with EEX is that the separation of the IR real singularities is done at the spin amplitude level and after squaring and spin summing them (numerically) the higher order terms are retained while in CEEX they are truncated. For more detailed discussion of see section 4.3, where we show explicitly the relations among $`\widehat{\beta }^{}s`$ of EEX and $`\overline{\beta }`$’s of EEX. Keeping the above in mind, we still have at least three possible versions of CEEX. In the following we shall describe them, concentrating mostly on the third one which is designed for the neutral resonances<sup>11</sup><sup>11</sup>11 The case of exponentiation for charged resonances like $`W^\pm `$ resonances is not yet covered in the literature. and which is the principal version implemented in the $`𝒦𝒦`$ Monte Carlo. Let us stress immediately that the resonance may be arbitrarily narrow. However, our approach works without any modification for any value of the resonance width.
#### 3.3.1 General discussion
We believe that CEEX is the only workable technique for treatment of narrow resonances in the exclusive MC. To understand the essential difference among three possible formulations of CEEX it is enough to limit the discussion to the simplest case of the $`𝒪(\alpha ^0)`$. The three possible options are:
* Version for the non-resonant Born without partitions:
$$𝔐_n^{(0)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2}\mathrm{}{}_{\sigma _n}{}^{k_n})=\underset{i=1}{\overset{n}{}}(𝔰_{\sigma _i}^I(k_i)+𝔰_{\sigma _i}^F(k_i))𝔅_{[ba][cd]}$$
(64)
* Version for the non-resonant Born with partitions:
$$𝔐_n^{(0)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2}\mathrm{}{}_{\sigma _n}{}^{k_n})=\underset{\mathrm{}\{I,F\}^n}{}\underset{i=1}{\overset{n}{}}𝔰_{\sigma _i}^{\{\mathrm{}_i\}}(k_i)𝔅_{[ba][cd]}(X_{\mathrm{}})$$
(65)
* Version for the resonant Born:
$$𝔐_n^{(0)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2}\mathrm{}{}_{\sigma _n}{}^{k_n})=\underset{\mathrm{}\{I,F\}^n}{}\underset{i=1}{\overset{n}{}}𝔰_{\sigma _i}^{\{\mathrm{}_i\}}(k_i)\frac{X_{\mathrm{}}^2}{(p_3+p_4)^2}\underset{R=\gamma ,Z}{}𝔅_{[ba][cd]}^B(X_{\mathrm{}})e^{\alpha B_4^R(X_{\mathrm{}})}$$
(66)
Let us immediately define the additional form-factor for the $`Z`$ resonance (case (C))
$$\alpha B_4^Z(X)=\frac{d^4k}{k^2m_\gamma ^2+iϵ}\frac{i}{(2\pi )^3}J_{I}^{}{}_{\mu }{}^{}(k)(J_F^\mu (k))^{}\left(\frac{(X)^2\overline{M}^2}{(Xk)^2\overline{M}^2}1\right),$$
(67)
where $`\overline{M}^2=M_Z^2iM_Z\mathrm{\Gamma }_Z`$, the currents $`J^\mu `$ are defined in (41), while for the non-resonant part we have $`B_4^\gamma (X)=0`$. The $`B_4^Z(X)`$ form-factor sums up to infinite order the virtual $`\alpha \mathrm{ln}(\mathrm{\Gamma }_Z/M_Z)`$ contributions – we postpone discussion of its origin and importance to the latter part of this section.
Coming back to the more elementary level we see that the case (B) becomes (A) if we can neglect the partition dependence of the four momentum in the Born: $`𝔅_{[ba][cd]}(X_{\mathrm{}})𝔅_{[ba][cd]}(P)`$, where $`P=p_a+p_b`$ or $`P=p_c+p_d`$ or any other choice which does not depend on momenta of the individual photons. This is thanks to the identity:
$$\underset{i=1}{\overset{n}{}}(𝔰_\sigma ^{\{F\}}(k_i)+𝔰_\sigma ^{\{I\}}(k_i))\underset{\mathrm{}\{I,F\}}{}\underset{i=1}{\overset{n}{}}𝔰_{\sigma _i}^{\{\mathrm{}_i\}}(k_i)$$
(68)
Only case (C) is efficient for the resonant process, so obviously (A) and (B) are limited to non-resonant processes. The immediate question is: which of them is better? If (A) is not summing higher order much better than (B), then it has the clear advantage of being simpler – summation over partitions makes the computer code more complicated and adds heavily to the consumption of CPU time. The answer is that, although we did not investigate quantitatively the differences between (A) and (B), we think that (B) sums up the LL higher orders more efficiently than (A) and is therefore better, even if there is no resonance. In our case, since we want to cover the resonant process anyway, it is a natural choice to use (B) for the non-resonant background component of the spin amplitudes (off-shell $`\gamma `$ exchange) even if it is not vital. Once summation over partitions is in place, it is the easiest to use it for the non-resonant background as well. The additional bonus of better higher order convergence provides an extra justification. Summarizing, if (C) is implemented then (B) comes for free.
Having discussed the differences among the three options let us now concentrate on option (C) for the resonant process, remembering that for the non-resonant background component it becomes automatically (B). First of all, for the narrow neutral resonance (the $`Z`$ boson in our case) the emission of the photons in the production and the decay processes are well separated by a long time interval, and are therefore completely independent and uncorrelated. In the perturbative QED this simple physical fact is reflected in a certain specific class of cancellations among the ISR and FSR photons on one hand and among virtual and real corrections on the other hand. For the inclusive observables like the total cross section or charge asymmetry the effects of the ISR-FSR interference (IFI) in the non-resonant case are of order $`\alpha /\pi `$, typically up to $`1\%`$, as can be seen from many example of the explicit $`𝒪(\alpha ^1)`$ calculations. The IFI effect will be of order $`(\alpha /\pi )(E_{\mathrm{max}}/E_{beam})`$, when the experimental cut on photon energy is $`E_{\mathrm{max}}`$. Note that the IFI effect is not directly enhanced by the big mass-logarithms like $`\mathrm{ln}s/m_e^220`$. For the resonant process the IFI effects in the inclusive observables are of order $`(\alpha /\pi )(\mathrm{\Gamma }/M)`$ and are therefore often negligible on the scale of the experimental error. One has to remember, however, that the additional suppression factor $`\mathrm{\Gamma }/M`$ disappears if the experimental cut on photon energy is of order of the resonance width, $`E_{\mathrm{max}}/E_{beam}\mathrm{\Gamma }/M`$, and for an even stronger cut $`E_{\mathrm{max}}<\mathrm{\Gamma }`$ the IFI effect becomes of order $`(\alpha /\pi )(E_{\mathrm{max}}/\mathrm{\Gamma })`$.
If $`\mathrm{\Gamma }/M`$ is extremely small, like for the $`\tau `$ lepton, the IFI cancellation can be taken for granted and the photon emission interference between production and decay can be neglected whatsoever. In the case of the $`Z`$ boson close to the $`Z`$ resonance (LEP1) the IFI effect is detectable experimentally but it is small enough that it can be omitted in the Monte Carlo programs used for correcting for the detector acceptance only. In this case KORALZ/YFS3 with the EEX matrix element was the acceptable solution.
The most convenient solution is the universal Monte Carlo in which IFI is included, which can evaluate IFI effects near the resonance, far from the resonance, for inclusive quantities and for strong energy cuts $`E_{\mathrm{max}}\mathrm{\Gamma }`$. This is exactly what our CEEX offers.
#### 3.3.2 Derivation of the resonance formfactor
As we have already pointed out (following refs. ), in the presence of narrow resonances it is not enough to sum up coherently the real emissions, taking properly into account energy shift in the resonance propagator (only due to ISR photons). It is also necessary to do the same for the virtual emission, and also sum them up to infinite-order – this is why the resonance form factor $`\mathrm{exp}(B_4^Z)`$ has to be included, see eqs. (66) and (67). In the following we shall derive eq. (67) for $`B_4^Z`$ and show analytically that the IFI cancellations do really work, as expected, to infinite order.
Let us write again the formula for standard YFS function in eq. (41) in a slightly modified notation
$$\begin{array}{cc}& B_4(p_a,\mathrm{},p_d)=\frac{d^4k}{k^2m_\gamma ^2+iϵ}\frac{i}{(2\pi )^3}S(k),\hfill \\ & S(k)=S_I(k)+S_F(k)+S_{Int}(k),\hfill \\ & S_I(k)=|J_I(k)|^2,S_F(k)=|J_F(k)|^2,S_{Int}(k)=2\mathrm{}(J_I(k)J_F^{}(k))\hfill \end{array}$$
(69)
In the presence of the narrow resonance, the YFS factorization of the virtual IR contributions has to take into account the dependence of the scalar part of the resonance propagator on photon energies of order $`\mathrm{\Gamma }`$ (the numerator treated in soft photon approximation as usual). The relevant integrals with $`n`$ virtual photons look as follows:
$$I=(P^2\overline{M}^2)\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{\mathrm{}𝒫_n}{}\underset{i=0}{\overset{n}{}}\frac{i}{(2\pi )^3}\frac{d^4k_i}{k_i^2m_\gamma ^2}S_\mathrm{}_i(k_i)\frac{1}{P_{\mathrm{}}^2\overline{M}^2},$$
(70)
where $`\overline{M}^2=\overline{M}^2iM\mathrm{\Gamma }`$, and $`𝒫_n`$ is set of all $`3^n`$ partitions $`(\mathrm{}_1,\mathrm{}_2,\mathrm{},\mathrm{}_n)`$ with $`\mathrm{}_i=I,F,Int`$, and $`P_{\mathrm{}}P\underset{\mathrm{}_i=Int}{}k_i`$ includes only momenta of photons in $`S_{Int}`$ and not of photons in $`S_I`$ or $`S_F`$. The $`(P^2\overline{M}^2)`$ factor is conventional, to make the integral dimensionless. We shall show that the above integral factorizes into the conventional YFS formfactor (dependent on the photon mass $`m_\gamma `$) and the additional non-IR factor due to the resonance $`R=Z`$
$$I=\mathrm{exp}(B_4^R(m_\gamma ,s,\overline{M}))=\mathrm{exp}(B_4(m_\gamma ,s)+\mathrm{\Delta }B_4^R(s,\overline{M})).$$
(71)
Our aim is to find the analytical form of the the additional function $`\mathrm{\Delta }B_4^R`$. In the current calculation we use the following approximate formula, also used by Greco et.al. ,
$$\alpha \mathrm{\Delta }B_4^R(s^{})=2Q_eQ_f\frac{\alpha }{\pi }\mathrm{ln}\left(\frac{t}{u}\right)\mathrm{ln}\left(\frac{\overline{M}^2s}{\overline{M}^2}\right)=\frac{1}{2}\gamma _{Int}\mathrm{ln}\left(\frac{\overline{M}^2s}{\overline{M}^2}\right).$$
(72)
In the following:
* We shall derive the above approximate result and
* show explicitly that the above approximate virtual interference part of the formfactor cancels exactly with the corresponding real interference contributions.
Since soft virtual photons entering into $`S_I`$ and $`S_F`$ in eq. (70) do not enter the resonance propagator, we may therefore factorize and sum up the contributions with $`S_I`$ and $`S_F`$:
$$\begin{array}{cc}\hfill I& =\underset{n_1=0}{\overset{\mathrm{}}{}}\frac{1}{n_1!}\underset{i_1=0}{\overset{n_1}{}}\frac{i}{(2\pi )^3}\frac{d^4k_{i_1}}{k_{i_1}^2m_\gamma ^2}S_I(k_{i_1})\underset{n_2=0}{\overset{\mathrm{}}{}}\frac{1}{n_2!}\underset{i_2=0}{\overset{n_2}{}}\frac{i}{(2\pi )^3}\frac{d^4k_{i_2}}{k_{i_2}^2m_\gamma ^2}S_F(k_{i_2})\hfill \\ & \times \underset{n_3=0}{\overset{\mathrm{}}{}}\frac{1}{n_3!}\underset{i_3=0}{\overset{n_3}{}}\frac{i}{(2\pi )^3}\frac{d^4k_{i_3}}{k_{i_3}^2m_\gamma ^2}S_{Int}(k_{i_3})\frac{1}{(P_{j=1}^{n_3}k_j)^2\overline{M}^2}\hfill \\ & e^{\alpha B_I+\alpha B_F}\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{i=0}{\overset{n}{}}\frac{i}{(2\pi )^3}\frac{d^4k_i}{k_i^2m_\gamma ^2}S_{Iin}(k_i)\frac{1}{(P_{j=1}^nk_j)^2\overline{M}^2}\hfill \end{array}$$
(73)
Now we neglect the quadratic terms in photon energies $`𝒪(k_ik_j)`$
$$\begin{array}{cc}& \frac{1}{(P_{j=1}^nk_j)^2\overline{M}^2}\frac{1}{P^22P_{j=1}^nk_j\overline{M}^2}=\frac{1}{P^2\overline{M}^2}\frac{1}{1_{j=1}^n\frac{2Pk_j}{P^2\overline{M}^2}}\hfill \\ & \frac{1}{P^2\overline{M}^2}\underset{j=1}{\overset{n}{}}\frac{1}{1\frac{2Pk_j}{P^2\overline{M}^2}}\frac{1}{P^2\overline{M}^2}\underset{j=1}{\overset{n}{}}\frac{P^2\overline{M}^2}{(Pk_j)^2\overline{M}^2}\hfill \end{array}$$
(74)
and this leads to
$$\begin{array}{cc}& I=e^{\alpha B_I+\alpha B_F}\mathrm{exp}\left(\frac{i}{(2\pi )^3}\frac{d^4k}{k_i^2m_\gamma ^2+iϵ}S_{Iin}(k)\frac{P^2\overline{M}^2}{(Pk)^2\overline{M}^2}\right)=e^{B_4(m_\gamma )+\mathrm{\Delta }B_4^R(\mathrm{\Gamma })}\hfill \\ & \mathrm{\Delta }B_4^R(\mathrm{\Gamma })=\frac{i}{(2\pi )^3}\frac{d^4k}{k^2}S_{Iin}(k)\left(\frac{P^2\overline{M}^2}{(Pk)^2\overline{M}^2}1\right)\hfill \end{array}$$
(75)
How solid is the above “derivation”? Strictly speaking it is justified in the limit where we follow Yennie, Frautschi and Suura in ref. and express the $`k0`$ emission amplitude as
$$\frac{1}{k}\left(\epsilon _1+𝒪(k/\overline{M})+\frac{k}{\mathrm{\Gamma }_Z}(\epsilon _2+𝒪(k/\overline{M}))\right),$$
where $`\epsilon _{1,2}`$ are constants independent of $`k`$, so that
$$\left|2Pk_j/(P^2\overline{M}^2)\right|1,$$
that is if photon energy is below the resonance width. This restriction is thus entirely analogous to the usual YFS expansion into IR-singular part and the rest. We note that Greco et.al. in refs. have also pointed out that the result for $`\mathrm{\Delta }B_4^R(\mathrm{\Gamma })`$ in eq. (75) follows from the YFS expansion; here shall show how this happens in detail.
The best situation would be to have a more precise evaluation of the integral of eq. (70) (the integral is probably calculable analytically). For the moment, however, following refs. : we choose an easier “pragmatic” approach based on the fact that the virtual and real contributions from IFI for photons $`E_\gamma >\mathrm{\Gamma }`$ do cancel, as a consequence of the time separation between production and decay, and we shall check that the above cancellation really works. In this way we trade analytical evaluation of the more difficult multiphoton virtual integral for an easier evaluation of the multiphoton real integral.
#### 3.3.3 Cancellation of the virtual formfactor with the real emissions
Let us therefore examine analytically the real multi-photon emission contribution from the IFI<sup>12</sup><sup>12</sup>12 In the practical CEEX calculation the contribution from IFI is evaluated numerically, inside the MC program.. The starting point is the integral in which the total photon energy $`K=_{j=1}^nk_j`$ is kept below $`E_{\mathrm{max}}=v_{\mathrm{max}}\sqrt{s}`$, where $`\mathrm{\Gamma }<E_{\mathrm{max}}<<\sqrt{s}`$:
$$\begin{array}{cc}& \sigma =\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{i=1}{\overset{n}{}}\frac{d^3k_i}{2k_i^0}\underset{\sigma _1\mathrm{}\sigma _n}{}\left|\underset{\mathrm{}\{I,F\}^n}{}\underset{j=1}{\overset{n}{}}𝔰_{[j]\{\mathrm{}_j\}}\frac{1}{X_{\mathrm{}}^2\overline{M}^2}e^{\alpha B_4^R(X_{\mathrm{}})}\right|^2\mathrm{\Theta }(E_{\mathrm{max}}\underset{j=1}{\overset{n}{}}k_j)\hfill \\ & =\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{K^0<v\sqrt{s}}{}\underset{i=1}{\overset{n}{}}\frac{d^3k_i}{2k_i^0}\underset{\sigma _1\mathrm{}\sigma _n}{}\underset{\mathrm{},\mathrm{}^{}\{I,F\}^n}{}\underset{j=1}{\overset{n}{}}𝔰_{[j]\{\mathrm{}_j\}}𝔰_{[j]\{\mathrm{}_j^{}\}}^{}\frac{e^{\alpha B_4^R(X_{\mathrm{}})}}{X_{\mathrm{}}^2\overline{M}^2}\left(\frac{e^{\alpha B_4^R(X_{\mathrm{}^{}})}}{X_{\mathrm{}^{}}^2\overline{M}^2}\right)^{}\hfill \\ & =\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{K^0<v\sqrt{s}}{}\underset{i=1}{\overset{n}{}}\frac{d^3k_i}{2k_i^0}\underset{\mathrm{},\{I^2,F^2,IF,FI\}^n}{}\underset{\mathrm{}_j=I^2}{}\stackrel{~}{S}_I(k_j)\underset{\mathrm{}_j=F^2}{}\stackrel{~}{S}_F(k_j)\hfill \\ & \underset{\mathrm{}_j=IF}{}\stackrel{~}{S}_{Int}(k_j)\underset{\mathrm{}_j=FI}{}\stackrel{~}{S}_{Int}(k_j)\frac{e^{\alpha B_4^R(PK_IK_{IF})}}{(PK_IK_{IF})^2\overline{M}^2}\left(\frac{e^{\alpha B_4^R(PK_IK_{FI})}}{(PK_IK_{FI})^2\overline{M}^2}\right)^{}\hfill \end{array}$$
(76)
where the Born amplitude we have simplified to the level of the scalar part of the resonance propagator and we denote
$$\begin{array}{cc}& \stackrel{~}{S}_I(k_j)=\underset{\sigma _j}{}|𝔰_{[j]}^{\{0\}}|^2,\stackrel{~}{S}_F(k_j)=\underset{\sigma _j}{}|𝔰_{[j]}^{\{1\}}|^2,\hfill \\ & \stackrel{~}{S}_{Int}(k_j)=\underset{\sigma _j}{}𝔰_{[j]}^{\{0\}}(𝔰_{[j]\{1\}})^{}=\underset{\sigma _j}{}𝔰_{[j]}^{\{1\}}(𝔰_{[j]\{0\}})^{},\hfill \\ & K_{I^2}=\underset{\mathrm{}_j=I^2}{}k_j,K_{F^2}=\underset{\mathrm{}_j=F^2}{}k_j,\hfill \\ & K_{IF}=\underset{\mathrm{}_j=IF}{}k_j,K_{FI}=\underset{\mathrm{}_j=FI}{}k_j,K=K_{I^2}+K_{F^2}+K_{IF}+K_{FI}\hfill \end{array}$$
(77)
As we see, the product of two sums, each over $`2^n`$ partitions $`\mathrm{},\mathrm{}^{}\{I,F\}^n`$, is now replaced by the single sum over $`4^n`$ partitions $`\mathrm{},\{I^2,F^2,IF,FI\}^n`$, where $`IF,FI`$ represent the interference terms.
Keeping track of the dependence of the propagators on $`K_{I^2}`$, $`K_{IF}`$ and $`K_{FI}`$, the summation over the number of photons can be reorganised, leading us back to the following factorized formula
$$\begin{array}{cc}& \sigma (v_{\mathrm{max}})=\underset{n_1=0}{\overset{\mathrm{}}{}}\frac{1}{n_1!}\underset{i_1=1}{\overset{n_1}{}}\frac{d^3k_{i_1}}{2k_{i_1}^0}\mathrm{\hspace{0.33em}2}\stackrel{~}{S}_I(k_{i_1})\underset{n_2=0}{\overset{\mathrm{}}{}}\frac{1}{n_2!}\underset{i_2=1}{\overset{n_2}{}}\frac{d^3k_{i_2}}{2k_{i_2}^0}\mathrm{\hspace{0.33em}2}\stackrel{~}{S}_F(k_{i_2})\hfill \\ & \underset{n_3=0}{\overset{\mathrm{}}{}}\frac{1}{n_3!}\underset{i_3=1}{\overset{n_3}{}}\frac{d^3k_{i_2}}{2k_{i_3}^0}\mathrm{\hspace{0.33em}2}\stackrel{~}{S}_{Int}(k_{i_3})\frac{e^{\alpha B_4^R(PK_{I^2}K_{IF})}}{(PK_{I^2}K_{IF})^2\overline{M}^2}\hfill \\ & \underset{n_4=0}{\overset{\mathrm{}}{}}\frac{1}{n_4!}\underset{i_4=1}{\overset{n_3}{}}\frac{d^3k_{i_4}}{2k_{i_4}^0}\mathrm{\hspace{0.33em}2}\stackrel{~}{S}_{Int}(k_{i_4})\left(\frac{e^{\alpha B_4^R(PK_{I^2}K_{FI})}}{(PK_{I^2}K_{FI})^2\overline{M}^2}\right)^{}\hfill \\ & \mathrm{\Theta }(E_{\mathrm{max}}K_{I^2}K_{F^2}K_{IF}K_{FI}),\hfill \end{array}$$
(78)
where $`K_{I^2}=_{i_1}k_{i_1}`$, $`K_{F^2}=_{i_2}k_{i_2}`$, $`K_{IF}=_{i_3}k_{i_3}`$ and $`K_{FI}=_{i_4}k_{i_4}`$. The sums over the pure initial and final state contributions, and the interference contributions are now well factorized and can be performed analytically. As a first step we integrate and sum up contributions from the very soft photons below $`\epsilon \sqrt{s}`$, similarly to what was shown in ref. :
$$\begin{array}{cc}& \sigma (v_{\mathrm{max}})=_0^{E_{\mathrm{max}}}𝑑E\delta (EE_IE_FE_{Int})_0^{E_{\mathrm{max}}}𝑑E_I𝑑E_F𝑑E_{IF}𝑑E_{FI}\hfill \\ & \underset{n_1=0}{\overset{\mathrm{}}{}}\frac{1}{n_1!}\underset{i_1=1}{\overset{n_1}{}}_{k_{i_1}^0>\epsilon E}\frac{d^3k_{i_1}}{2k_{i_1}^0}2\stackrel{~}{S}_I(k_{i_1})e^{2\alpha \stackrel{~}{B}_I(\epsilon E)+2\alpha \mathrm{}B_I}\delta (E_I\underset{i_1}{}k_{i_1})\hfill \\ & \underset{n_2=0}{\overset{\mathrm{}}{}}\frac{1}{n_2!}\underset{i_2=1}{\overset{n_2}{}}_{k_{i_2}^0>\epsilon E}\frac{d^3k_{i_2}}{2k_{i_2}^0}2\stackrel{~}{S}_F(k_{i_2})e^{2\alpha \stackrel{~}{B}_F(\epsilon E)+2\alpha \mathrm{}B_F}\delta (E_F\underset{i_2}{}k_{i_2})\hfill \\ & \underset{n_3=0}{\overset{\mathrm{}}{}}\frac{1}{n_3!}\underset{i_3=1}{\overset{n_3}{}}_{k_{i_3}^0>\epsilon E}\frac{d^3k_{i_2}}{2k_{i_3}^0}2\stackrel{~}{S}_{Int}(k_{i_3})\frac{e^{\alpha \mathrm{\Delta }B_4^R(PK_{I^2}K_{IF})}}{(PK_{I^2}K_{IF})^2\overline{M}^2}e^{\alpha \stackrel{~}{B}_{Int}(\epsilon E)+\alpha \mathrm{}B_{Int}}\hfill \\ & \underset{n_4=0}{\overset{\mathrm{}}{}}\frac{1}{n_4!}\underset{i_3=1}{\overset{n_4}{}}_{k_{i_3}^0>\epsilon E}\frac{d^3k_{i_2}}{2k_{i_3}^0}2\stackrel{~}{S}_{Int}(k_{i_3})\left(\frac{e^{\alpha \mathrm{\Delta }B_4^R(PK_{I^2}K_{FI})}}{(PK_{I^2}K_{FI})^2\overline{M}^2}\right)^{}e^{\alpha \stackrel{~}{B}_{Int}(\epsilon E)+\alpha \mathrm{}B_{Int}}\hfill \\ & e^{2\alpha \mathrm{}\mathrm{\Delta }B_4^R}\delta (E_{Int}\underset{i_3}{}k_{i_3}),\hfill \end{array}$$
(79)
where $`E=\frac{\sqrt{s}}{2}`$. The integration over photon momenta can be performed without any approximation leading to the following result
$$\begin{array}{cc}& \sigma (v_{\mathrm{max}})=_0^{v_{\mathrm{max}}}𝑑v\delta (vv_Iv_Fv_{IF}v_{FI})\hfill \\ & 𝑑v_IF(\gamma _I)\gamma _Iv_I^{\gamma _I1}e^{2\alpha \stackrel{~}{B}_I(E)+2\alpha \mathrm{}B_I}𝑑v_FF(\gamma _F)\gamma _Fv_F^{\gamma _F1}e^{2\alpha \stackrel{~}{B}_F(E)+2\alpha \mathrm{}B_F}\hfill \\ & 𝑑v_{IF}F\left(\frac{\gamma _{Int}}{2}\right)\frac{1}{2}\gamma _{Int}v_{IF}^{\frac{1}{2}\gamma _{IF}1}\left(\frac{e^{\alpha \mathrm{\Delta }B_4^R(s(1v_I)(1v_{IF}))}}{s(1v_I)(1v_{IF})\overline{M}^2}\right)e^{\alpha \stackrel{~}{B}_{Int}(E)+\alpha \mathrm{}B_{Int}}\hfill \\ & 𝑑v_{FI}F\left(\frac{\gamma _{Int}}{2}\right)\frac{1}{2}\gamma _{Int}v_{FI}^{\frac{1}{2}\gamma _{FI}1}\left(\frac{e^{\alpha \mathrm{\Delta }B_4^R(s(1v_I)(1v_{FI}))}}{s(1v_I)(1v_{FI})\overline{M}^2}\right)^{}e^{\alpha \stackrel{~}{B}_{Int}(E)+\alpha \mathrm{}B_{Int}}\hfill \end{array}$$
(80)
in which is explicitly free of any IR divergences.
The essential question is whether we have perfect cancellations of the $`\mathrm{ln}(\mathrm{\Gamma }/M_Z)`$ terms in the interference subintegral
$$\begin{array}{c}\hfill I_{Int}=\mathrm{}_0^{v_{max}v_Iv_Fv_{FI}}𝑑v_{IF}F\left(\frac{\gamma _{IF}}{2}\right)\frac{1}{2}\gamma _{IF}v_{IF}^{\frac{1}{2}\gamma _{Int}1}\frac{e^{\alpha \mathrm{\Delta }B_4^R(s^{}(1v_{IF}))}}{s^{}(1v_{IF})\overline{M}^2}\end{array}$$
(81)
We omit from consideration the constant IR-finite factor $`e^{\alpha \stackrel{~}{B}_{Int}(E)+\alpha \mathrm{}B_{Int}}`$ because it does not depend on resonance parameters. The bulk of the integral comes from the neighbourhood of $`v_{IF}=0`$ and the integrand is $`1/v^2`$ at large $`v`$ due to the resonance; we can therefore extend the integration limit to $`_0^{\mathrm{}}𝑑v_{Int}`$ at the expense of an error of $`𝒪\left(\frac{\mathrm{\Gamma }}{M_Z}\right)`$. One possible evaluation method is to use the standard techniques of the complex functions. First, we reformulate the integral as an integral over the discontinuity $`C_1`$ along the real axis<sup>13</sup><sup>13</sup>13 We have also pulled out of the integral the $`e^{\alpha \mathrm{\Delta }B_4^R}`$ factor, because the most of integral comes from the neighbourhood of the singularity at $`v_{IF}=0`$.
$$\begin{array}{c}\hfill I_{Int}=F\left(\frac{\gamma _{IF}}{2}\right)e^{\alpha \mathrm{\Delta }B_4^R(s^{})}\frac{1}{i\mathrm{sin}(\pi \frac{1}{2}\gamma _{Int})}_{C_1}𝑑z\frac{1}{2}\gamma _{Int}(z)^{\frac{1}{2}\gamma _{Int}1}\frac{1}{s^{}\overline{M}^2s^{}z}\end{array}$$
(82)
Since the contour can be closed in a standard way with the big circle, the integral is given by the value of the residue at $`z=1\overline{M}^2/s^{}`$.
$$\begin{array}{cc}& I_{Int}=F\left(\frac{\gamma _{IF}}{2}\right)e^{\alpha \mathrm{\Delta }B_4^R(s^{})}\frac{\pi \frac{1}{2}\gamma _{Int}}{\mathrm{sin}(\pi \frac{1}{2}\gamma _{Int})}\left(\frac{\overline{M}^2s^{}}{s^{}}\right)^{\gamma _{Int}1}\frac{1}{s^{}}\hfill \\ & =\frac{1}{\overline{M}^2s^{}}F\left(\frac{\gamma _{Int}}{2}\right)\frac{\pi \frac{1}{2}\gamma _{Int}}{\mathrm{sin}(\pi \frac{1}{2}\gamma _{Int})}e^{\alpha \mathrm{\Delta }B_4^R(s^{})}\left(\frac{\overline{M}^2s^{}}{s^{}}\right)^{\frac{1}{2}\gamma _{Int}}\hfill \\ & =\frac{1}{\overline{M}^2s^{}}(1+𝒪(\gamma _{Int}))\hfill \end{array}$$
(83)
The above is true because
$$\begin{array}{c}\hfill \alpha \mathrm{\Delta }B_4^R(s^{})=2Q_eQ_f\frac{\alpha }{\pi }\mathrm{ln}\left(\frac{t}{u}\right)\mathrm{ln}\left(\frac{\overline{M}^2s^{}}{\overline{M}^2}\right)=\frac{1}{2}\gamma _{Int}\mathrm{ln}\left(\frac{\overline{M}^2s^{}}{\overline{M}^2}\right)\end{array}$$
(84)
We have therefore proven the full cancellation of the dependence on the resonance parameters for the integrated cross section.
#### 3.3.4 Definitions of $`\widehat{\beta }`$’s with partitions
The $`𝒪(\alpha ^r)`$, $`r=0,1,2`$, $`\widehat{\beta }`$-functions for the variant of the CEEX with summation over the partitions, as in eqs. (3.1-LABEL:eq:ceex-master2), are derived with the recursive relations of eqs. (49) (similar to those of eqs. (51)). The only additional complication is that we must keep track of the indices which say whether an external real photon is of ISR or FSR type and of the total photon momentum after emission of the ISR photons (the one which enters resonance propagator, if such a resonance is present)
$$\begin{array}{cc}& \widehat{\beta }_0^{(l)}\left({}_{\lambda }{}^{p};P\right)=𝔐_0^{(l)}\left({}_{\lambda }{}^{p};P\right),l=0,1,2,\hfill \\ & \widehat{\beta }_{1\{I\}}^{(1+l)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};Pk_1)=𝔐_{1\{I\}}^{(1+l)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};Pk_1)\widehat{\beta }_0^{(l)}({}_{\lambda }{}^{p};Pk_1)𝔰_{\sigma _1}^{\{I\}}(k_1),l=0,1,\hfill \\ & \widehat{\beta }_{1\{F\}}^{(1+l)}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};P\right)=𝔐_{1\{F\}}^{(1+l)}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};P\right)\widehat{\beta }_0^{(l)}\left({}_{\lambda }{}^{p};P\right)𝔰_{\sigma _1}^{\{F\}}(k_1),l=0,1,\hfill \\ & \widehat{\beta }_{2\{\omega _1,\omega _2\}}^{(2)}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2};X_\omega \right)=𝔐_{2\{\omega _1,\omega _2\}}^{(2)}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2};X_\omega \right)\hfill \\ & \widehat{\beta }_{1\{\omega _1\}}^{(1)}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};X_\omega \right)𝔰_{\sigma _2}^{\{\omega _2\}}(k_2)\widehat{\beta }_{1\{\omega _2\}}^{(1)}\left({}_{\lambda }{}^{p}{}_{\sigma _2}{}^{k_2};X_\omega \right)𝔰_{\sigma _1}^{\{\omega _1\}}(k_1)\widehat{\beta }_0^{(0)}\left({}_{\lambda }{}^{p};X_\omega \right)𝔰_{\sigma _1}^{\{\omega _1\}}(k_1)𝔰_{\sigma _2}^{\{\omega _2\}}(k_2),\hfill \end{array}$$
(85)
where $`X_\omega =P\underset{\omega _i=I}{}k_i`$, $`P=p_a+p_b`$. Introduction of the partition index $`\omega _i`$ defining whether a photon belongs to ISR or FSR is in a sense not such a deep and great complication – it is now just another (third) attribute of the photon like its helicity.
Let us look closer into the structure of the term like $`\widehat{\beta }_{1\{\omega _1\}}^{(1)}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};X_\omega \right)𝔰_{\sigma _2}^{\{\omega _2\}}(k_2)`$. For example $`\omega _1=F`$ and $`\omega _2=I`$ it reads $`\widehat{\beta }_{1\{F\}}^{(1)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};Pk_2)𝔰_{\sigma _2}^{\{I\}}(k_2)`$, that is the total shift in $`X`$ in $`\widehat{\beta }^{(1)}`$ depends not only on the type $`\omega _1`$ of “its own photon” but also on the type $`\omega _2`$ of the photon in $`𝔰^{\{\omega _2\}}`$ factor which multiplies it!
The $`𝔐`$-amplitude in eq. (85) is given essentially by eq. (40) with the formfactor including the resonance part (if present)
$$𝔐_{n\{\omega \}}^{(r)R}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}\mathrm{}{}_{\sigma _n}{}^{k_n};X_\omega \right)=\left[e^{\alpha B_4\alpha B_4^R(X_\omega )}_{n\{\omega \}}^{(r)R}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}\mathrm{}{}_{\sigma _n}{}^{k_n};X_\omega \right)\right]|_{𝒪(\alpha ^r)},$$
(86)
As we see the type $`R=\gamma ,Z`$ of the “resonance” formfactor $`B_4^R`$ has to be adjusted to the type of the component in $`^{(r)R}`$ (we have temporarily introduced an explicit index $`R`$ into $``$ and $`𝔐`$ and $`\gamma `$ is essentially a “resonance” with the zero width).
### 3.4 Virtual corrections, no real photons
We now start to accumulate the actual formulas for the $`\widehat{\beta }`$-functions entering the CEEX amplitudes of in eqs. (3.1-LABEL:eq:ceex-master2) with the case of no real photons and up to two virtual photons. The “raw material” are the $``$-amplitudes from Feynman diagrams which are turned into $`\widehat{\beta }`$-functions using the recursive relations of eqs. (85).
#### 3.4.1 Photonic corrections
Let us start with the simple case of $`𝒪(\alpha ^1)`$ spin amplitudes with one virtual and zero real photon coming directly from Feynman diagrams, which will be used to obtain the first order $`\widehat{\beta }_0^{(1)}`$. The relevant spin amplitudes are
$$_0^{(1)}\left({}_{\lambda }{}^{p};X\right)=𝔅\left({}_{\lambda }{}^{p};X\right)[1+Q_e^2F_1(s,m_e,m_\gamma )][1+Q_f^2F_1(s,m_f,m_\gamma )]+_{\mathrm{box}}^{(1)}\left({}_{\lambda }{}^{p};X\right),$$
(87)
where $`F_1`$ is the standard electric form-factor regularized with a photon mass, see fig. 3. We omit, for the moment, the magnetic form-factor $`F_2`$; this is justified for light final fermions. It will be restored in the future. In $`F_1`$ we keep the exact final fermion mass. If not stated otherwise, the four-momentum conservation $`p_a+p_b=p_c+p_d`$ holds.
In the present work we use spin amplitudes for $`\gamma `$-$`\gamma `$ and $`\gamma `$-$`Z`$-boxes in the small mass approximation $`m_e^2/s0,m_f^2/s0`$, see fig. 3, following refs. ,
$$\begin{array}{cc}\hfill _{\mathrm{Box}}^{(1)}\left({}_{\lambda }{}^{p};X\right)& =2ie^2\underset{B=\gamma ,Z}{}\frac{g_{\lambda _a}^{B,e}g_{\lambda _a}^{B,f}T_{\lambda _c\lambda _a}T_{\lambda _b\lambda _d}^{}+g_{\lambda _a}^{B,e}g_{\lambda _a}^{B,f}U_{\lambda _c\lambda _b}^{}U_{\lambda _a\lambda _d}}{X^2M_{B}^{}{}_{}{}^{2}+i\mathrm{\Gamma }_BX^2/M_B}\delta _{\lambda _a,\lambda _b}\delta _{\lambda _c,\lambda _d}\hfill \\ & \frac{\alpha }{\pi }Q_eQ_f\left[\delta _{\lambda _a,\lambda _c}f_{\mathrm{BDP}}(\overline{M}_B^2,m_\gamma ,s,t,u)\delta _{\lambda _a,\lambda _c}f_{\mathrm{BDP}}(\overline{M}_B^2,m_\gamma ,s,u,t)\right],\hfill \end{array}$$
(88)
where
$$\begin{array}{cc}& f_{\mathrm{BDP}}(\overline{M}_B^2,m_\gamma ,s,u,t)=\mathrm{ln}\left(\frac{t}{u}\right)\mathrm{ln}\left(\frac{m_\gamma ^2}{(tu)^{1/2}}\right)2\mathrm{ln}\left(\frac{t}{u}\right)\mathrm{ln}\left(\frac{\overline{M}_Z^2s}{\overline{M}_Z^2}\right)\hfill \\ & +\mathrm{Li}_2\left(\frac{\overline{M}_Z^2+u}{\overline{M}_Z^2}\right)\mathrm{Li}_2\left(\frac{\overline{M}_Z^2+t}{\overline{M}_Z^2}\right)\hfill \\ & +\frac{(\overline{M}_Z^2s)(ut\overline{M}_Z^2)}{u^2}\left\{\mathrm{ln}\left(\frac{t}{s}\right)\mathrm{ln}\left(\frac{\overline{M}_Z^2s}{\overline{M}_Z^2}\right)+\mathrm{Li}_2\left(\frac{\overline{M}_Z^2+t}{\overline{M}_Z^2}\right)\mathrm{Li}_2\left(\frac{\overline{M}_Z^2s}{\overline{M}_Z^2}\right)\right\}\hfill \\ & +\frac{(\overline{M}_Z^2s)(\overline{M}_Z^2s)}{us}\mathrm{ln}\left(\frac{\overline{M}_Z^2s}{\overline{M}_Z^2}\right)+\frac{\overline{M}_Z^2s}{u}\mathrm{ln}\left(\frac{t}{\overline{M}_Z^2}\right),\hfill \end{array}$$
(89)
$`\overline{M}_Z^2=M_Z^2iM_Z\mathrm{\Gamma }_Z`$, $`\overline{M}_\gamma ^2=m_\gamma ^2`$, and the function $`f_{\mathrm{BDP}}`$ is that of eq. (11) of ref. . The standard Mandelstam variables $`s,t`$ and $`u`$ are defined as usual: $`s=(p_a+p_b)^2,t=(p_ap_c)^2,t=(p_ap_d)^2`$. Since in the rest of our calculation we do not use $`m_f^2/s0`$, we therefore intend to replace the above box spin amplitudes with the finite-mass results<sup>14</sup><sup>14</sup>14 For the $`\gamma `$-$`\gamma `$ box we use the spin amplitudes with the exact final fermion mass. It seems, however, that the $`\gamma `$-$`Z`$ box for the heavy fermion is missing in the literature. that were given in ref. .)
Now using eq. (86) we determine
$$\widehat{\beta }_0^{(1)}\left({}_{\lambda }{}^{p};X\right)=𝔅\left({}_{\lambda }{}^{p};X\right)(1+\delta _{Virt}^{(1)e}(s))(1+\delta _{Virt}^{(1)f}(s))+_{\mathrm{Box}}^{(1)}\left({}_{\lambda }{}^{p};X\right)$$
(90)
where
$$\begin{array}{cc}& \delta _{Virt}^{(1)e}(s)=Q_e^2F_1(s,m_e,m_\gamma )Q_e^2\alpha B_2(p_a,p_b,m_\gamma )=Q_e^2\frac{\alpha }{\pi }\frac{1}{2}\overline{L}_e,\hfill \\ & \delta _{Virt}^{(1)f}(s)=Q_f^2F_1(s,m_f,m_\gamma )Q_f^2\alpha B_2(p_c,p_d,m_\gamma )=Q_f^2\frac{\alpha }{\pi }\frac{1}{2}\overline{L}_f,\hfill \\ & \overline{L}_e=\mathrm{ln}\left(\frac{s}{m_e^2}\right)+i\pi 1,\overline{L}_f=\mathrm{ln}\left(\frac{s}{m_f^2}\right)+i\pi 1.\hfill \end{array}$$
(91)
Note that we departed in eq. (90) from the strict $`𝒪(\alpha ^1)`$ by retaining the $`\delta _{Virt}^{(1)e}(s)\delta _{Virt}^{(1)f}(s)`$ term, i.e., replacing the “additive” form $`1+\delta _{Virt}^{(1)e}(s)+\delta _{Virt}^{(1)f}(s)`$ with the “factorized” form $`(1+\delta _{Virt}^{(1)e}(s))(1+\delta _{Virt}^{(1)f}(s))`$. The above does not need really much justification – it is obviously closer to the reality of the higher-orders, so the “factorized” form is preferable. The only question is whether the above method does not disturb IR-cancellations. It does not, as it is seen from the definitions of $`\delta _{Virt}^{(1)e}(s)`$ and $`\delta _{Virt}^{(1)f}(s)`$.
The IR-subtraction in $`_{\mathrm{Box}}^{(1)}`$ using eq. (86) at $`𝒪(\alpha ^1)`$ leads to the IR-finite $`_{\mathrm{Box}}`$. The above subtraction is equivalent to the following substitution
$$f_{\mathrm{BDP}}(\overline{M}_B^2,m_\gamma ,s,t,u)f_{\mathrm{BDP}}(\overline{M}_B^2,m_\gamma ,s,t,u)f_{\mathrm{IR}}(m_\gamma ,t,u),$$
(92)
where
$$f_{\mathrm{IR}}(m_\gamma ,t,u)=\frac{2}{\pi }B_2(p_a,p_c,m_\gamma )\frac{2}{\pi }B_2(p_a,p_d,m_\gamma )=\mathrm{ln}\left(\frac{t}{u}\right)\mathrm{ln}\left(\frac{m_\gamma ^2}{\sqrt{tu}}\right)+\frac{1}{2}\mathrm{ln}\left(\frac{t}{u}\right),$$
(93)
and the additional resonance factor $`\mathrm{exp}\left(\alpha B_4^Z(s)\right)`$ in eq. (86) induces the additional subtraction in the $`\gamma `$-Z box part:
$$f_{\mathrm{BDP}}(s,t,u)f_{\mathrm{BDP}}(s,t,u)\alpha B_4^Z(s),$$
(94)
see eq. (72) for the definition of $`\alpha B_4^Z`$.
Our $`𝒪(\alpha ^2)`$ expressions for $`\widehat{\beta }_0^{(2)}`$ are still incomplete. We base them on the graphs depicted in fig. 4. (In fig. 4 we omitted some trivial transpositions of the diagrams.) Following again the eq. (86), we obtain
$$\widehat{\beta }_0^{(2)}\left({}_{\lambda }{}^{p};X\right)=𝔅\left({}_{\lambda }{}^{p};X\right)(1+\delta _{Virt}^{(2)e}(s,m_e))(1+\delta _{Virt}^{(2)f}(s,m_f))+_{\mathrm{Box}}^{(2)}\left({}_{\lambda }{}^{p};X\right)$$
(95)
In the present calculation we neglect the two-loop double-box contributions in $`_{\mathrm{Box}}^{(2)}`$, depicted in the first row in fig. 5 and vertex-box type of diagrams, see examples of diagrams in the second row in fig. 5 <sup>15</sup><sup>15</sup>15 In fact the two-loop double-box contributions became known recently , so there is a chance to include it in the future.. In fact we keep only the first order box contribution $`_{\mathrm{Box}}^{(1)}`$ in our incomplete $`𝒪(\alpha ^2)`$ type matrix element.
Two remarks: in spite of the temporary lack of the above contribution we are not stuck because what we neglect is IR-finite! This statement is not so trivial as it may look because in the calculation without exponentiation neglecting such contributions would violate IR cancellations, and correcting for such a violation would be rather complicated and physically dangerous. Secondly, what we neglect is expected to be numerically small, of $`𝒪(\alpha ^2L^1)`$ and therefore it does not make much harm to our overall physical precision.
Coming back to the $`𝒪(\alpha ^2)`$ corrections to the electric form factor from the diagrams in fig. 4, they are well known since they were calculated in refs. and they contribute as follows
$$\begin{array}{c}\hfill \delta _{Virt}^{(2)e}(s,m_e)=\delta _{Virt}^{(1)e}(s)+\left(\frac{\alpha }{\pi }\right)^2\left(\frac{\overline{L_e}^2}{8}+\overline{L_e}\left(\frac{3}{32}\frac{3}{4}\zeta _2+\frac{3}{2}\zeta _3\right)\right),\\ \hfill \delta _{Virt}^{(2)f}(s,m_f)=\delta _{Virt}^{(1)f}(s)+\left(\frac{\alpha }{\pi }\right)^2\left(\frac{\overline{L_f}^2}{8}+\overline{L_f}\left(\frac{3}{32}\frac{3}{4}\zeta _2+\frac{3}{2}\zeta _3\right)\right),\end{array}$$
(96)
In the above we kept terms of $`𝒪(\alpha ^2L^2)`$ and $`𝒪(\alpha ^2L^1)`$, and neglected the known negligible terms of $`𝒪(\alpha ^2L^0)`$.
#### 3.4.2 Electroweak corrections
In the not so interesting case of the absence of the electroweak (EW) corrections the couplings of two neutral bosons $`\gamma `$ and $`Z`$ are defined in a conventional way:
$$\begin{array}{cc}& G_\lambda ^{Z,f}=g_V^{Z,f}\lambda g_A^{Z,f}(\mathrm{?}),G_\lambda ^{\gamma ,f}=g_V^{Z,f},\lambda =+,=R,L,\hfill \\ & g_V^{\gamma ,e}=Q_e=1,g_{V,f}^\gamma =Q_f,g_A^{\gamma ,e}=0,g_{A,f}^\gamma =0,\hfill \\ & g_V^{Z,e}=\frac{2T_e^34Q_e\mathrm{sin}^2\theta _W}{16\mathrm{sin}^2\theta _W\mathrm{cos}^2\theta _W},g_V^{Z,f}=\frac{2T_f^34Q_f\mathrm{sin}^2\theta _W}{16\mathrm{sin}^2\theta _W\mathrm{cos}^2\theta _W},\hfill \\ & g_A^{Z,e}=\frac{2T_e^3}{16\mathrm{sin}^2\theta _W\mathrm{cos}^2\theta _W},g_A^{Z,f}=\frac{2T_f^3}{16\mathrm{sin}^2\theta _W\mathrm{cos}^2\theta _W},\hfill \end{array}$$
(97)
where $`T_f^3`$ is the isospin of the left-handed component of the fermion ($`T_d^3=1/2,T_e^3=1/2`$).
The actual implementation of EW corrections is practically the same as in KORALZ . It goes as follows: The $`\gamma `$ and $`Z`$-propagators are multiplied by the corresponding hook-functions (scalar form-factors) due to vacuum polarizations
$$H_\gamma H_\gamma \times \frac{1}{2\mathrm{\Pi }_\gamma },H_ZH_Z\times 16\mathrm{sin}^2\theta _W\mathrm{cos}^2\theta _W\frac{G_\mu M_Z^2}{\alpha _{_{\mathrm{QED}}}8\pi \sqrt{2}}\rho _{\mathrm{EW}}.$$
(98)
In addition the vector couplings of the $`Z`$ get multiplied by extra form factors. First of all we replace
$$\begin{array}{cc}& g_V^{Z,e}=\frac{2T_e^34Q_e\mathrm{sin}^2\theta _W}{16\mathrm{sin}^2\theta _W\mathrm{cos}^2\theta _W}=>\frac{2T_e^34Q_e\mathrm{sin}^2\theta _WF_{EW}^e(s)}{16\mathrm{sin}^2\theta _W\mathrm{cos}^2\theta _W}\hfill \\ & g_V^{Z,f}=\frac{2T_f^34Q_f\mathrm{sin}^2\theta _W}{16\mathrm{sin}^2\theta _W\mathrm{cos}^2\theta _W}=>\frac{2T_f^34Q_f\mathrm{sin}^2\theta _WF_{EW}^f(s)}{16\mathrm{sin}^2\theta _W\mathrm{cos}^2\theta _W}\hfill \end{array}$$
(99)
where $`F_{EW}^e(s)`$ and $`F_{EW}^f(s)`$ are electroweak form factors provided by the DIZET package , which is part the ZFITTER semianalytical code and correspond to electroweak vertex corrections.
The electroweak box diagrams require more complicated treatment. In the Born spin amplitudes we have essentially two products of the coupling constants
$$\begin{array}{cc}& g_\lambda ^{Z,e}g_\lambda ^{Z,f}=(g_V^{Z,e}\lambda g_A^{Z,e})(g_V^{Z,f}+\lambda g_A^{Z,f})=g_V^{Z,e}g_V^{Z,f}\lambda g_A^{Z,e}g_V^{Z,f}+\lambda g_V^{Z,e}g_A^{Z,f}g_A^{Z,e}g_A^{Z,f},\hfill \\ & g_\lambda ^{Z,e}g_\lambda ^{Z,f}=(g_V^{Z,e}\lambda g_A^{Z,e})(g_V^{Z,f}\lambda g_A^{Z,f})=g_V^{Z,e}g_V^{Z,f}\lambda g_A^{Z,e}g_V^{Z,f}\lambda g_V^{Z,e}g_A^{Z,f}+g_A^{Z,e}g_A^{Z,f}.\hfill \end{array}$$
(100)
In the above the following modification is done for the doubly-vector component
$$g_V^{Z,e}g_V^{Z,f}=>\frac{4T_e^3T_f^38T_e^3Q_fF_{EW}^f(s)8T_f^3Q_eF_{EW}^e(s)+16Q_fQ_fF_{EW}^{ef}(s,t)}{(16\mathrm{sin}^2\theta _W\mathrm{cos}^2\theta _W)^2},$$
(101)
where the new form factor $`F_{EW}^{ef}(s,t)`$ corresponds to electroweak boxes and is angle dependent. The Born spin amplitudes modified in the above way are used also in the case of the presence of the single and multiple real photons, see next sections.
### 3.5 One real photon
The discussion of the $`\widehat{\beta }_1`$ tensors corresponding to emission of a single real photon we start with the tree level case (zero virtual photons). The starting point is the well known $`𝒪(\alpha ^1)`$ split amplitude for the single bremsstrahlung which we shall reconsider separately first in the case of the emission from the initial state beams (ISR) and later for emission from the final state fermions (FSR). This will be the “raw material” for obtaining $`\widehat{\beta }_1^{(0)}`$ using eqs. (85).
The first-order, 1-photon, ISR matrix element from the Feynman diagrams depicted in fig. 6 reads
$$\begin{array}{cc}\hfill _{1\{I\}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1})=& eQ_e\overline{v}(p_b,\lambda _b)𝐌_1\frac{\overline{)}p_a+m\overline{)}k_1}{2k_1p_a}\overline{)}ϵ_{\sigma _1}^{}(k_1)u(p_a,\lambda _a)\hfill \\ \hfill +& eQ_e\overline{v}(p_b,\lambda _b)\overline{)}ϵ_{\sigma _1}^{}(k_1)\frac{\overline{)}p_b+m+\overline{)}k_1}{2k_1p_b}𝐌_{\{I\}}u(p_a,\lambda _a),\hfill \end{array}$$
(102)
where
$$𝐌_{\{I\}}=ie^2\underset{B=\gamma ,Z}{}\mathrm{\Pi }_B^{\mu \nu }(X)G_{e,\mu }^B(G_{f,\nu }^B)_{[cd]},$$
(103)
is the annihilation scattering spinor matrix, including final-state spinors. The above expression we split into soft IR parts<sup>16</sup><sup>16</sup>16 This kind of separation was already exploited in ref. . We thank E. Richter-Wa̧s for attracting our attention to this method. proportional to $`(\overline{)}p\pm m)`$ and non-IR parts proportional to $`\overline{)}k_1`$. Employing the completeness relations of eq. (LABEL:transition-defs) in the Appendix A to those parts we obtain:
$$\begin{array}{cc}\hfill _{1\{I\}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1})=& \frac{eQ_e}{2k_1p_a}\underset{\rho }{}𝔅[{}_{\lambda _b}{}^{p_b}{}_{\rho _a}{}^{p_a}]{}_{[cd]}{}^{}U[{}_{\rho _a}{}^{p_a}{}_{\sigma _1}{}^{k_1}{}_{\lambda _a}{}^{p_a}]+\frac{eQ_e}{2k_1p_b}\underset{\rho }{}V[{}_{\lambda _b}{}^{p_b}{}_{\sigma _1}{}^{k_1}{}_{\rho _b}{}^{p_b}]𝔅[{}_{\rho _b}{}^{p_b}{}_{\lambda _a}{}^{p_a}]_{[cd]}\hfill \\ \hfill +& \frac{eQ_e}{2k_1p_a}\underset{\rho }{}𝔅[{}_{\lambda _b}{}^{p_b}{}_{\rho }{}^{k_1}]{}_{[cd]}{}^{}U[{}_{\rho }{}^{k_1}{}_{\sigma _1}{}^{k_1}{}_{\lambda _a}{}^{p_a}]\frac{eQ_e}{2k_1p_b}\underset{\rho }{}V[{}_{\lambda _b}{}^{p_b}{}_{\sigma _1}{}^{k_1}{}_{\rho }{}^{k_1}]𝔅[{}_{\rho }{}^{k_1}{}_{\lambda _a}{}^{p_a}]{}_{[cd]}{}^{}.\hfill \end{array}$$
(104)
The summation in the first two terms gets eliminated due to the diagonality property of $`U`$ and $`V`$, see eq. (LABEL:diagonality) in the Appendix A, and leads to
$$\begin{array}{cc}& ^{1\{I\}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1})=𝔰_{\sigma _1}^{\{I\}}(k_1)𝔅[{}_{\lambda }{}^{p}]+r_{\{I\}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}),\hfill \\ & r_{\{I\}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1})=+\frac{eQ_e}{2k_1p_a}\underset{\rho }{}𝔅[{}_{\lambda _b}{}^{p_b}{}_{\rho }{}^{k_1}]{}_{[cd]}{}^{}U[{}_{\rho }{}^{k_1}{}_{\sigma _1}{}^{k_1}{}_{\lambda _a}{}^{p_a}]\frac{eQ_e}{2k_1p_b}\underset{\rho }{}V[{}_{\lambda _b}{}^{p_b}{}_{\sigma _1}{}^{k_1}{}_{\rho }{}^{k_1}]𝔅[{}_{\rho }{}^{k_1}{}_{\lambda _a}{}^{p_a}]{}_{[cd]}{}^{},\hfill \\ & 𝔰_{\sigma _1}^{\{I\}}(k_1)=eQ_e\frac{b_{\sigma _1}(k_1,p_a)}{2k_1p_a}+eQ_e\frac{b_{\sigma _1}(k_1,p_b)}{2k_1p_b},\hfill \end{array}$$
(105)
The soft part is now clearly separated and the remaining non-IR part, necessary for the CEEX, is obtained.
The case of final-state one real photon emission (FSR), see fig. 7, can be analysed in a similar way. The first-order FSR, 1-photon, matrix element is
$$\begin{array}{cc}\hfill _{1\{F\}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1})& =eQ_f\overline{u}(p_c,\lambda _c)\overline{)}ϵ_{\sigma _1}^{}(k_1)\frac{\overline{)}p_c+m+\overline{)}k_1}{2k_1p_c}𝐌_0v(p_d,\lambda _d)\hfill \\ & +eQ_f\overline{u}(p_c,\lambda _c)𝐌_{\{F\}}\frac{\overline{)}p_d+m\overline{)}k_1}{2k_1p_d}\overline{)}ϵ_{\sigma _1}^{}(k_1)v(p_d,\lambda _d),\hfill \end{array}$$
(106)
where
$$𝐌_{\{F\}}=ie^2\underset{B=\gamma ,Z}{}\mathrm{\Pi }_B^{\mu \nu }(X)(G_{e,\mu }^B)_{[ba]}G_{f,\nu }^B,$$
(107)
is spinor matrix for annihilation scattering,including initial spinors. Similarly, the expansion into soft and non-IR parts for the FSR spin amplitudes is done in the way completely analogous to the ISR case
$$\begin{array}{cc}& _{1\{F\}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1})=𝔰_{\sigma _1}^{\{F\}}(k_1)𝔅({}_{\lambda }{}^{p})+r_{\{F\}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}),\hfill \\ & r_{\{F\}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1})=\frac{eQ_f}{2k_1p_c}\underset{\rho }{}U[{}_{\lambda _c}{}^{p_c}{}_{\sigma _1}{}^{k_1}{}_{\rho }{}^{k_1}]𝔅{}_{[ba]}{}^{}[{}_{\rho }{}^{k_1}{}_{\lambda _d}{}^{p_d}]\frac{eQ_f}{2k_1p_d}\underset{\rho }{}𝔅{}_{[ba]}{}^{}[{}_{\lambda _c}{}^{p_c}{}_{\rho }{}^{k_1}]V[{}_{\rho }{}^{k_1}{}_{\sigma _1}{}^{k_1}{}_{\lambda _d}{}^{p_d}],\hfill \\ & 𝔰_{\sigma _1}^{\{F\}}(k_1)=eQ_f\frac{b_{\sigma _1}(k_1,p_c)}{2k_1p_c}eQ_f\frac{b_{\sigma _1}(k_1,p_d)}{2k_1p_d}.\hfill \end{array}$$
(108)
For the purpose of the following discussion of the remaining non-IR terms it is useful to introduce an even more compact tensor notation:
$$U[{}_{\lambda _f}{}^{p_f}{}_{\sigma _i}{}^{k_i}{}_{\sigma _j}{}^{k_j}]U_{[f,i,j]},𝔅[{}_{\lambda _b}{}^{p_b}{}_{\lambda _a}{}^{p_a}][{}_{\lambda _c}{}^{p_c}{}_{\lambda _d}{}^{p_d}]𝔅_{[ba][cd]},$$
(109)
etc. For the “primed” indices we understand contractions, for instance
$$U_{[a,i,j^{}]}V_{[j^{},j,b]}\underset{\sigma _j^{}=\pm }{}U[{}_{\lambda _a}{}^{p_a}{}_{\sigma _i}{}^{k_i}{}_{\sigma _j^{}}{}^{k_j}]V[{}_{\sigma _j^{}}{}^{k_j}{}_{\sigma _j}{}^{k_j}{}_{\lambda _b}{}^{p_b}].$$
(110)
Using the above notation, the complete $`𝒪(\alpha ^1)`$ spin amplitude for 1-photon ISR+FSR, coming directly from Feynman diagrams, with the explicit split into IR and non-IR parts, ISR and FSR parts, reads
$$\begin{array}{cc}& 𝔐_1^{(1)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1})=𝔐_{1\{I\}}^{(1)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1})(Pk_1)+𝔐_{1\{F\}}^{(1)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1})(P)\hfill \\ & =𝔰_{[1]}^{\{I\}}𝔅({}_{\lambda }{}^{p};Pk_1)+r_{\{I\}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};Pk_1)+𝔰_{[1]}^{\{F\}}𝔅\left({}_{\lambda }{}^{p};P\right)+r_{\{F\}}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};P\right),\hfill \\ & r_{\{I\}}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};X\right)=\frac{eQ_e}{2kp_a}𝔅_{[b1^{}cd]}(X)U_{[1^{}1a]}\frac{eQ_e}{2kp_b}V_{[b11^{}]}𝔅_{[1^{}acd]}(X)\hfill \\ & r_{\{F\}}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};X\right)=\frac{eQ_f}{2kp_c}U_{[c11^{}]}𝔅_{[ba1^{}d]}(X)\frac{eQ_f}{2kp_d}𝔅_{[bac1^{}]}(X)V_{[1^{}1d]}\hfill \end{array}$$
(111)
In the lowest-order the Born spin amplitudes $`𝔅`$ are defined in eq. (35), and we show explicitly as an argument the four-momentum $`X`$ which enters the propagator of the $`s`$-channel exchange. Note that the formulas here differ by an overall sign from those of ref.
#### 3.5.1 First- and second-order $`\widehat{\beta }_1`$
Now we employ the tree level, $`𝒪(\alpha ^1)`$ variant of eqs. (85) getting the following results:
$$\begin{array}{cc}& \widehat{\beta }_{1\{I\}}^{(1)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};Pk_1)r_{\{I\}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};Pk_1)\hfill \\ & \widehat{\beta }_{1\{F\}}^{(1)}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};P\right)r_{\{F\}}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};P\right)+(\frac{(p_c+p_d+k_1)^2}{(p_c+p_d)^2}1)𝔅\left({}_{\lambda }{}^{p};X\right),\hfill \end{array}$$
(112)
The “context dependent” reduced total momentum $`X`$ (the total four-momentum in the resonance propagator, if present) is in the above definition uniquely defined as $`X=Pk_1`$ in the case of ISR, and $`X=P`$ in the case of FSR. In the general context of the CEEX amplitude of eqs. (3.1-LABEL:eq:ceex-master2), that is in presence of the additional “spectator” ISR photons in a given term, $`X`$ is also defined quite unambiguously: $`X_{\mathrm{}}`$ includes not only $`k_1`$ but also all additional ISR momenta in the process. For the pseudo-flux factor there is some ambiguity, however. In the presence of the additional “spectator” ISR photons it can be defined either as $`(p_a+p_bk_1)^2/(p_a+p_b)^2`$ or $`(p_c+p_d+k_1)^2/(p_c+p_d)^2`$. We are free to choose any of them and we opted for the second choice (it seems to lead to more stable MC weights).
The one-loop level, $`𝒪(\alpha ^2)`$ case of $`\widehat{\beta }_1^{(2)}`$ is quite interesting because this is the first time that we deal with the nontrivial case of the simultaneous emission of virtual and real photons. It is therefore instructive to write the formal definitions of $`\widehat{\beta }_1^{(2)}`$ following eqs. (86) and (85) in this particular case:
$`𝔐_{1\{\omega \}}^{(2)}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};X_\omega \right)`$ $`=\left\{e^{\alpha B_4\alpha B_4^R(X_\omega )}_{1\{\omega \}}^{(2)}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};X_\omega \right)\right\}|_{𝒪(\alpha ^2)},\omega =I,F,R=\gamma ,Z,`$
$`\widehat{\beta }_{1\{I\}}^{(2)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};Pk_1)`$ $`=𝔐_{1\{I\}}^{(2)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};Pk_1)𝔰_{\sigma _1}^{\{I\}}(k_1)\widehat{\beta }_0^{(1)}({}_{\lambda }{}^{p};Pk_1)`$
$`\widehat{\beta }_{1\{F\}}^{(2)}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};P\right)`$ $`=𝔐_{1\{F\}}^{(2)}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};P\right)𝔰_\sigma ^{\{F\}}(k_1)\widehat{\beta }_0^{(1)}\left({}_{\lambda }{}^{p};P\right).`$
What is presently available from the Feynman diagrams? For the moment we have at our disposal the amplitudes corresponding to vertex-like diagrams in fig. (8), and we miss diagrams of the “5-box” type shown in the third (bottom) row in fig. 5. More precisely, after applying the IR virtual subtraction of eq. (3.5.1) we expand in the number of loops, keeping track of the initial/final state attachment of the virtual photon:
$$\begin{array}{cc}& 𝔐_{1\{\omega \}}^{(2)}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};X\right)=𝔐_{1\{\omega \}}^{(1)}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};X\right)+\alpha Q_e^2𝔐_{1\{\omega \},I^2}^{[1]}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};X\right)+\alpha Q_f^2𝔐_{1\{\omega \},F^2}^{[1]}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};X\right)\hfill \\ & +\alpha Q_eQ_f𝔐_{1\{\omega \},Box5}^{[1]}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};X\right).\hfill \end{array}$$
(113)
In the above expression the first term describes the already discussed tree level single bremsstrahlung, the next two correspond to vertex-like diagrams in fig. 8, and the last one represents the “5-box” type diagrams in the third row of fig. 5. In the present version we temporarily omit from the calculation the contribution to $`\widehat{\beta }_1^{(2)}`$ from the last, “5-box” term which looks as follows:
$$\widehat{\beta }_{1\{\omega \},Box5}^{(2)}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};X\right)=\alpha Q_eQ_f𝔐_{1\{\omega \},Box5}^{[1]}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};X\right)𝔰_{[1]}^{\{I\}}_{\mathrm{Box}}^{(1)}\left({}_{\lambda }{}^{p};X\right)𝔰_{[1]}^{\{F\}}_{\mathrm{Box}}^{(1)}\left({}_{\lambda }{}^{p};X\right).$$
(114)
As we see the trivial IR-part, which we remove, is proportional to the ordinary box contributions already discussed before. We expect the above to contribute in the integrated cross section to be at most of $`𝒪(\alpha ^2L^1)`$, and in the resonance scattering it will be suppressed by the additional $`\mathrm{\Gamma }/M`$ factor.
Limiting ourselves to the pure “vertex-like” diagrams of fig. 8, for one real ISR $`(\omega =I)`$ photon we obtain from the Feynman rules the following $`𝒪(Q_e^2\alpha ^2)`$ result
$$\begin{array}{cc}\hfill \widehat{\beta }_{1\{I\}}^{(2)}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};X\right)& r_{\{I\}}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};X\right)(1+\delta _{Virt}^{(1)e}(s)+\rho _{Virt}^{(2)e}(s,\stackrel{~}{\alpha }_1,\stackrel{~}{\beta }_1))(1+\delta _{Virt}^{(1)f}(s))\hfill \\ & +𝔅\left({}_{\lambda }{}^{p};X\right)𝔰_{\sigma _1}^{\{I\}}(k_1)\rho _{Virt}^{(2)e}(s,\stackrel{~}{\alpha },\stackrel{~}{\beta })\hfill \end{array}$$
(115)
where
$$\begin{array}{cc}& \rho _{Virt}^{(2)e}(s,\stackrel{~}{\alpha },\stackrel{~}{\beta })=\frac{\stackrel{~}{\alpha }}{\pi }Q_e^2\frac{1}{2}(V(s,\stackrel{~}{\alpha },\stackrel{~}{\beta })+V(s,\stackrel{~}{\beta },\stackrel{~}{\alpha })),\hfill \\ & V(s,\stackrel{~}{\alpha },\stackrel{~}{\beta })=\mathrm{ln}(\stackrel{~}{\alpha })\mathrm{ln}(1\stackrel{~}{\beta })\hfill \\ & +\mathrm{Li}_2(\stackrel{~}{\alpha })\frac{1}{2}\mathrm{ln}^2(1\stackrel{~}{\alpha })+\frac{3}{2}\mathrm{ln}(1\stackrel{~}{\alpha })+\frac{1}{2}\frac{\stackrel{~}{\alpha }(1\stackrel{~}{\alpha })}{(1+(1\stackrel{~}{\alpha })^2)}\hfill \end{array}$$
(116)
and we use Sudakov variables
$$\stackrel{~}{\alpha }_i=\frac{2k_ip_b}{2p_ap_b},\stackrel{~}{\beta }_i=\frac{2k_ip_a}{2p_ap_b}.$$
(117)
Let us make a number of observations concerning eq. (115):
* The terms of $`𝒪(\alpha ^4)`$ like $`|𝔰_\sigma ^{\{I\}}\rho _{Virt}^{(2)e}|^2`$ in the cross section, although beyond $`𝒪(\alpha ^2)`$, are not rejected, as it would be the case in the ordinary $`𝒪(\alpha ^2)`$ calculation without exponentiation. They are included in the process of numerical evaluation of the differential cross sections out of spin amplitudes. (It is essential that they are IR-finite.)
* The term $`r_{\{I\}}\delta _{Virt}^{(1)e}`$ contributes $`𝒪(\alpha ^2L^2)`$ to the integrated cross section – one $`L^1`$ is explicit (from the virtual photon) and another $`L^1`$ is from the integration over the angle of the real photon.
* The term $`\mathrm{ln}(\stackrel{~}{\alpha })\mathrm{ln}(1\stackrel{~}{\beta })`$ contributes a correction of $`𝒪(\alpha ^2L^2)`$ to the integrated cross section with the double logarithm $`L^2`$ resulting directly from the integration over the angle of the real photon:
$$\frac{dk^3}{k^0}\mathrm{}[\rho _{Virt}^{(2)e}(k)\{\widehat{\beta }_0𝔰_\sigma ^{\{I\}}(k)\}^{})]Q_e^2\alpha ^2_{m_e^2/s}\frac{d\stackrel{~}{\alpha }}{\stackrel{~}{\alpha }}\mathrm{ln}(\stackrel{~}{\alpha })Q_e^2\alpha ^2\mathrm{ln}^2\frac{s}{m_e^2}$$
* The other terms in $`\widehat{\beta }_{1\{I\}}^{(2)}`$ contribute at most $`𝒪(\alpha ^2L^1)`$.
* The FSR virtual corrections are included multiplicatively through the factor $`(1+\delta _{Virt}^{(1)f}(s))`$ and not additively like $`(1+\delta _{Virt}^{(1)e}(s)+\delta _{Virt}^{(1)f}(s))`$. This is our deliberate choice.
* The subleading term $`\stackrel{~}{\alpha }(1\stackrel{~}{\alpha })/(1+(1\stackrel{~}{\alpha })^2)`$ has in fact a more complicated spin structure than that of the Born amplitude (it should be restored in future). The unpolarized integrated cross section is however correct in $`𝒪(\alpha ^2L^1)`$.
The analogous $`𝒪(Q_f^2\alpha ^2)`$ contribution for one real FSR $`(\omega =0)`$ photon is
$$\begin{array}{cc}\hfill \widehat{\beta }_{1\{F\}}^{(2)}\left({}_{\lambda }{}^{p}{}_{\sigma }{}^{k};X\right)& r_{\{F\}}\left({}_{\lambda }{}^{p}{}_{\sigma }{}^{k};X\right)(1+\delta _{Virt}^{(1)e}(s))(1+\delta _{Virt}^{(1)f}(s)+\rho _{Virt}^{(2)f}(s,\stackrel{~}{\alpha }^{},\stackrel{~}{\beta }^{}))\hfill \\ & +𝔅\left({}_{\lambda }{}^{p};X\right)𝔰_\sigma ^{\{F\}}(k)\rho _{Virt}^{(2)f}(s,\stackrel{~}{\alpha }^{},\stackrel{~}{\beta }^{})\hfill \\ & +𝔅\left({}_{\lambda }{}^{p};X\right)𝔰_\sigma ^{\{F\}}(k)(1+\delta _{Virt}^{(1)e}(s))(1+\delta _{Virt}^{(1)f}(s))(1\frac{(p_c+p_d+k)^2}{(p_c+p_d)^2})\hfill \end{array}$$
(118)
where
$$\begin{array}{cc}& \rho _{Virt}^{(2)f}(s,\stackrel{~}{\alpha }^{},\stackrel{~}{\beta }^{})=\frac{\stackrel{~}{\alpha }}{\pi }Q_f^2\frac{1}{4}\overline{L}_f(\mathrm{ln}(1\stackrel{~}{\alpha }^{\prime \prime })+\mathrm{ln}(1\stackrel{~}{\beta }^{\prime \prime })),\hfill \\ & \stackrel{~}{\alpha }^{}=\frac{2kp_d}{2p_cp_d},\stackrel{~}{\beta }^{}=\frac{2kp_c}{2p_cp_d},\stackrel{~}{\alpha }^{\prime \prime }=\frac{\stackrel{~}{\alpha }^{}}{1+\stackrel{~}{\alpha }^{}+\stackrel{~}{\beta }^{}},\stackrel{~}{\beta }^{\prime \prime }=\frac{\stackrel{~}{\beta }^{}}{1+\stackrel{~}{\alpha }^{}+\stackrel{~}{\beta }^{}}.\hfill \end{array}$$
(119)
In the above FSR amplitudes we keep only the LL part averaged over the photon angles, similarly as in EEX. This corresponds to present status of our CEEX amplitudes implemented in $`𝒦𝒦`$ MC version 4.13, and we expect this to be improved in the future.
### 3.6 2-real photons
In the $`𝒪(\alpha ^2)`$ contributions from two real photons are completely at the tree level, without virtual corrections (in the future $`𝒪(\alpha ^3)`$ version we shall include virtual corrections to double bremsstrahlung in the LL approximation). The double bremsstrahlung is considered in three separate cases, two ISR photons, two FSR photons and one ISR plus one FSR photon. The corresponding spin amplitudes will be given without any approximation, in particular we do not use the small mass approximation $`m_f/\sqrt{s}<<1`$. The main problems to be solved will be
* to write all spin amplitudes in a form easy for numerical evaluation, that is in terms of $`U`$ and $`V`$ matrices,
* to extract $`\widehat{\beta }_2`$ functions by means of removing IR-singular parts.
#### 3.6.1 2-real ISR photons
The second-order, two-photon, ISR matrix element from the Feynman rules, see Fig. 9, reads as follows
$$\begin{array}{cc}\hfill _{2\{II\}}^{(2)}& ({}_{\lambda _a}{}^{p_a}{}_{\lambda _b}{}^{p_b}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2};Pk_1k_2)=ie^2\underset{B=\gamma ,Z}{}\mathrm{\Pi }_B^{\mu \nu }(Pk_1k_2)(G_{f,\nu }^B)_{[cd]}(eQ_e)^2\overline{v}(p_b,\lambda _b)\{\hfill \\ & G_{e,\mu }^B\frac{(\overline{)}p_a+m)\overline{)}k_1\overline{)}k_2}{2k_1p_a2k_2p_a+2k_1k_2}\overline{)}ϵ_{\sigma _1}^{}(k_1)\frac{(\overline{)}p_a+m)\overline{)}k_2}{2k_2p_a}\overline{)}ϵ_{\sigma _2}^{}(k_2)\hfill \\ \hfill +& \overline{)}ϵ_{\sigma _1}^{}(k_1)\frac{(\overline{)}p_b+m)+\overline{)}k_1}{2k_1p_b}\overline{)}ϵ_{\sigma _2}^{}(k_2)\frac{(\overline{)}p_b+m)+\overline{)}k_1+\overline{)}k_2}{2k_1p_b2k_2p_b+2k_1k_2}G_{e,\mu }^B\hfill \\ \hfill +& \overline{)}ϵ_{\sigma _1}^{}(k_1)\frac{(\overline{)}p_b+m)+\overline{)}k_1}{2k_1p_b}G_{e,\mu }^B\frac{(\overline{)}p_a+m)\overline{)}k_2}{2k_2p_a}\overline{)}ϵ_{\sigma _2}^{}(k_2)+(12)\}u(p_a,\lambda _a).\hfill \end{array}$$
(120)
We shall use eq (85) which in this case reads
$$\begin{array}{cc}& \widehat{\beta }_{2\{II\}}^{(2)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2};Pk_1k_2)=𝔐_{2\{II\}}^{(2)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2};Pk_1k_2)\widehat{\beta }_{1\{I\}}^{(1)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};Pk_1k_2)𝔰_{\sigma _2}^{\{I\}}(k_2)\hfill \\ & \widehat{\beta }_{1\{I\}}^{(1)}({}_{\lambda }{}^{p}{}_{\sigma _2}{}^{k_2};Pk_1k_2)𝔰_{\sigma _1}^{\{I\}}(k_1)\widehat{\beta }_0^{(0)}({}_{\lambda }{}^{p};Pk_1k_2)𝔰_{\sigma _1}^{\{I\}}(k_1)𝔰_{\sigma _2}^{\{I\}}(k_2).\hfill \end{array}$$
(121)
We shall proceed similarly as in 1-photon case, we shall isolate from the above expression the group of terms containing two factors $`(\overline{)}p+m)`$, then the group containing single factor $`(\overline{)}p+m)`$ and finally the rest. Such a split represents almost exactly the split in eq. (50) into contribution with two $`𝔰`$-factors (double IR singularity), with single $`𝔰`$-factor (single IR singularity) and the IR-finite remnant $`\widehat{\beta }_2^{(2)}`$ which is our primary goal. In other words, we decompose $`𝔐_{2\{II\}}^{(2)}`$ into several terms/parts, as described above, and we apply the IR-subtraction of eq. (121) term-by-term.
Let us discuss first the doubly IR-singular part proportional to two factors $`(\overline{)}p+m)`$. To simplify maximally the discussion let us neglect for the moment $`2k_1k_2`$ in the propagator. Using the completeness relations of eq. (LABEL:transition-defs) and the diagonality property of eq. (LABEL:diagonality) in Appendix A, we can factorize soft factors exactly and completely
$$\begin{array}{cc}\hfill (eQ_e)^2\overline{v}(p_b,\lambda _b)& \{G_{e,\mu }^B\frac{(\overline{)}p_a+m)}{2k_1p_a+2k_2p_a}\overline{)}ϵ_{\sigma _1}^{}(k_1)\frac{(\overline{)}p_a+m)}{2k_2p_a}\overline{)}ϵ_{\sigma _2}^{}(k_2)\hfill \\ & +\overline{)}ϵ_{\sigma _1}^{}(k_1)\frac{(\overline{)}p_b+m)}{2k_1p_b}\overline{)}ϵ_{\sigma _2}^{}(k_2)\frac{(\overline{)}p_b+m)}{2k_1p_b+2k_2p_b}G_{e,\mu }^B\hfill \\ & +\overline{)}ϵ_{\sigma _1}^{}(k_1)\frac{(\overline{)}p_b+m)}{2k_1p_b}G_{e,\mu }^B\frac{(\overline{)}p_a+m)}{2k_2p_a}\overline{)}ϵ_{\sigma _2}^{}(k_2)+(12)\}u(p_a,\lambda _a)\hfill \\ \hfill =(G_{e,\mu }^B)_{[ba]}(eQ_e)^2& \{\frac{b_{\sigma _1}(k_1,p_a)}{2k_1p_a+2k_2p_a}\frac{b_{\sigma _2}(k_2,p_a)}{2k_2p_a}+\frac{b_{\sigma _1}(k_1,p_b)}{2k_1p_b}\frac{b_{\sigma _2}(k_2,p_b)}{2k_1p_b+2k_2p_b}\hfill \\ & \frac{b_{\sigma _1}(k_1,p_b)}{2k_1p_b}\frac{b_{\sigma _2}(k_2,p_a)}{2k_2p_a}+(12)\}\hfill \\ \hfill =(G_{e,\mu }^B)_{[ba]}& 𝔰_{\sigma _1}^{\{I\}}(k_1)𝔰_{\sigma _2}^{\{I\}}(k_2),\hfill \end{array}$$
(122)
where the identity
$$\frac{1}{2k_1p_a+2k_2p_a}\frac{1}{2k_1p_a}+\frac{1}{2k_1p_a+2k_2p_a}\frac{1}{2k_2p_a}=\frac{1}{2k_1p_a}\frac{1}{2k_2p_a}$$
(123)
was instrumental.
If we restore the terms $`2k_1k_2`$ in the propagator the corresponding analog of (122) $`_{2\{II\}}^{\mathrm{DoubleIR}}`$ leads to
$$\begin{array}{cc}& \widehat{\beta }_{2\{II\}}^{(2)\mathrm{Double}}[{}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2}]=_{2\{II\}}^{\mathrm{DoubleIR}}[{}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2}]𝔰_{\sigma _1}^{\{I\}}(k_1)𝔰_{\sigma _2}^{\{I\}}(k_2)𝔅[{}_{\lambda }{}^{p}]\hfill \\ & =(𝔰_{[1]}^{(a)}𝔰_{[2]}^{(a)}\mathrm{\Delta }_a+𝔰_{[1]}^{(b)}𝔰_{[2]}^{(b)}\mathrm{\Delta }_b)𝔅[{}_{\lambda }{}^{p}],\hfill \\ & 𝔰_{\sigma _i}^{(a)}(k_i)𝔰_{[i]}^{(a)}=eQ_e\frac{b_{\sigma _1}(k_i,p_a)}{2k_ip_a},𝔰_{\sigma _i}^{(b)}(k_i)𝔰_{[i]}^{(b)}=+eQ_e\frac{b_{\sigma _1}(k_i,p_b)}{2k_ip_b},\hfill \\ & 𝔰_{\sigma _i}^{\{I\}}(k_i)𝔰_{[i]}^{(a)}+𝔰_{[i]}^{(b)}𝔰_{\sigma _i}^{(a)}(k_i)+𝔰_{\sigma _i}^{(b)}(k_i),\hfill \\ & \mathrm{\Delta }_f=\frac{2k_1p_f+2k_2p_f}{2k_1p_f+2k_2p_f2k_1k_2}1=\frac{\pm 2k_1k_2}{2k_1p_f+2k_2p_f2k_1k_2},f=a,b,c,d\hfill \end{array}$$
(124)
and the upper sign should be taken for $`f=a,b`$. Obviously $`\widehat{\beta }^{(2)\mathrm{Double}}`$ is IR-finite because of the $`\mathrm{\Delta }_f`$ factor. In the above we have introduced a more compact notation for $`𝔰`$-factors. In addition from now on we shall use the following shorthand notation
$$r_{if}=2k_ip_f,r_{ij}=2k_ik_j,f=a,b,c,d,i,j,=1,2,\mathrm{}n.$$
(125)
The next class of terms which we are going to consider carefully is the one in which we sum terms with a single $`(\overline{)}p+m)`$, more precisely, let us include terms, which may lead to a single IR singularity (if $`k_1<<k_2`$ or $`k_2<<k_1`$), that is with $`(\overline{)}p+m)`$ next to a spinor, at the end of the fermion line:
$$\begin{array}{cc}\hfill _{2\{II\}}^{\mathrm{SingleIR}}[{}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2}]=ie^2\underset{B=\gamma ,Z}{}& \mathrm{\Pi }_B^{\mu \nu }(X)(G_{f,\nu }^B)_{[cd]}\hfill \\ \hfill \times (eQ_e)^2\overline{v}(p_b,\lambda _b)& \{G_{e,\mu }^B\frac{\overline{)}k_1\overline{)}k_2}{r_{1a}r_{2a}+r_{12}}\overline{)}ϵ_{\sigma _1}^{}(k_1)\frac{(\overline{)}p_a+m)}{r_{2a}}\overline{)}ϵ_{\sigma _2}^{}(k_2)\hfill \\ \hfill +& \overline{)}ϵ_{\sigma _1}^{}(k_1)\frac{(\overline{)}p_b+m)}{r_{1b}}\overline{)}ϵ_{\sigma _2}^{}(k_2)\frac{\overline{)}k_1+\overline{)}k_2}{r_{1b}r_{2b}+r_{12}}G_{e,\mu }^B\hfill \\ \hfill +& \overline{)}ϵ_{\sigma _1}^{}(k_1)\frac{(\overline{)}p_b+m)}{r_{1b}}G_{e,\mu }^B\frac{\overline{)}k_2}{r_{2a}}\overline{)}ϵ_{\sigma _2}^{}(k_2)\hfill \\ \hfill +& \overline{)}ϵ_{\sigma _1}^{}(k_1)\frac{\overline{)}k_1}{r_{1b}}G_{e,\mu }^B\frac{(\overline{)}p_a+m)}{r_{2a}}\overline{)}ϵ_{\sigma _2}^{}(k_2)+(12)\}u(p_a,\lambda _a).\hfill \end{array}$$
(126)
Using the compact notation, already introduced when (re)calculating single bremsstrahlung, we express $`_{2\{II\}}^{\mathrm{SingleIR}}`$ in a form friendly for numerical evaluation, that is in terms of $`U`$ and $`V`$ matrices,
$$\begin{array}{cc}& _{2\{II\}}^{\mathrm{SingleIR}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2})=\hfill \\ & =eQ_e\frac{𝔅_{[b1^{}][cd]}U_{[1^{}1a]}𝔅_{[b2^{}][cd]}U_{[2^{}1a]}}{r_{1a}r_{2a}+r_{12}}𝔰_{[2]}^{(a)}+eQ_e𝔰_{[1]}^{(b)}\frac{V_{[b22^{}]}𝔅_{[2^{}a][cd]}+V_{[b21^{}]}𝔅_{[1^{}a][cd]}}{r_{1a}r_{2a}+r_{12}}\hfill \\ & eQ_e𝔰_{[1]}^{(b)}𝔅_{[b2^{}][cd]}\frac{U_{[2^{}2a]}}{r_{2a}}+eQ_e\frac{V_{[b11^{}]}}{r_{1b}}𝔅_{[1^{}a][cd]}𝔰_{[2]}^{(a)}+(12)\hfill \end{array}$$
(127)
On the other hand the single-IR part to be eliminated is
$$\begin{array}{cc}\hfill \widehat{\beta }_{1(1)[1]}^{(1)}𝔰_{[2]}^{\{I\}}& +\widehat{\beta }_{1(1)[2]}^{(1)}𝔰_{[1]}^{\{I\}}=r_{[1]}^{\{I\}}𝔰_{[2]}^{\{I\}}+r_{[2]}^{\{I\}}𝔰_{[1]}^{\{I\}}\hfill \\ & =(eQ_e𝔅_{[b1^{}][cd]}\frac{U_{[1^{}1a]}}{r_{1a}}eQ_e\frac{V_{[b11^{}]}}{r_{1a}}𝔅_{[1^{}a][cd]})𝔰_{[2]}^{\{I\}}+(12).\hfill \end{array}$$
(128)
Altogether we get
$$\begin{array}{cc}& \widehat{\beta }_{2\{II\}}^{\mathrm{Single}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2})=_{2\{II\}}^{\mathrm{SingleIR}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2})\widehat{\beta }_{1(1)[1]}^{(1)}𝔰_{[2]}^{\{I\}}\widehat{\beta }_{1(1)[2]}^{(1)}𝔰_{[1]}^{\{I\}}\hfill \\ & =eQ_e𝔅_{[b2^{}][cd]}\frac{U_{[2^{}1a]}}{r_{1a}r_{2a}+r_{12}}𝔰_{[2]}^{(a)}+eQ_e𝔰_{[1]}^{(b)}\frac{V_{[b21^{}]}}{r_{1a}r_{2a}+r_{12}}𝔅_{[1^{}a][cd]}\hfill \\ & eQ_e𝔅_{[b1^{}][cd]}\left(\frac{U_{[1^{}1a]}}{r_{1a}r_{2a}+r_{12}}\frac{U_{[1^{}1a]}}{r_{1a}}\right)𝔰_{[2]}^{(a)}\hfill \\ & +eQ_e𝔰_{[1]}^{(b)}(\frac{V_{[b22^{}]}}{r_{1a}r_{2a}+r_{12}}\frac{V_{[b22^{}]}}{r_{2b}})𝔅_{[2^{}a][cd]}+(12).\hfill \end{array}$$
(129)
It is rather straightforward to see that the above is IR-finite.
Finally, we have to include all remaining terms from eq. (122) which have not yet included in $`\widehat{\beta }_{2\{II\}}`$. They are IR-finite (in the case of only soft photon energy) and they read
$$\begin{array}{cc}\hfill \widehat{\beta }_{2\{II\}}^{\mathrm{Rest}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2})& =ie^2\underset{B=\gamma ,Z}{}\mathrm{\Pi }_B^{\mu \nu }(X)(G_{f,\nu }^B)_{[cd]}(eQ_e)^2\overline{v}(p_b,\lambda _b)\{\hfill \\ & G_{e,\mu }^B\frac{(\overline{)}p_a+m)\overline{)}k_1\overline{)}k_2}{r_{1a}r_{2a}+r_{12}}\overline{)}ϵ_{\sigma _1}^{}(k_1)\frac{\overline{)}k_2}{r_{2a}}\overline{)}ϵ_{\sigma _2}^{}(k_2)\hfill \\ \hfill +& \overline{)}ϵ_{\sigma _1}^{}(k_1)\frac{\overline{)}k_1}{r_{1b}}\overline{)}ϵ_{\sigma _2}^{}(k_2)\frac{(\overline{)}p_b+m)+\overline{)}k_1+\overline{)}k_2}{r_{1b}r_{2b}+r_{12}}G_{e,\mu }^B\hfill \\ \hfill +& \overline{)}ϵ_{\sigma _1}^{}(k_1)\frac{\overline{)}k_1}{r_{1b}}G_{e,\mu }^B\frac{\overline{)}k_2}{r_{2a}}\overline{)}ϵ_{\sigma _2}^{}(k_2)+(12)\}u(p_a,\lambda _a),\hfill \end{array}$$
(130)
Using tensor notation in the fermion helicity indices the above can be expressed in terms of $`U`$ and $`V`$ matrices as follows
$$\begin{array}{cc}\hfill \widehat{\beta }_{2\{II\}}^{\mathrm{Rest}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2})=& (eQ_e)^2\frac{𝔅_{[ba^{}][cd]}U_{[a^{}12^{\prime \prime }]}𝔅_{[b1^{}][cd]}U_{[1^{}12^{\prime \prime }]}𝔅_{[b2^{}][cd]}U_{[2^{}12^{\prime \prime }]}}{r_{1a}r_{2a}+r_{12}}\frac{U_{[2^{\prime \prime }2a]}}{r_{2a}}\hfill \\ \hfill +& (eQ_e)^2\frac{V_{[b11^{\prime \prime }]}}{r_{1b}}\frac{V_{[1^{\prime \prime }2b^{}]}𝔅_{[b^{}a][cd]}+V_{[1^{\prime \prime }21^{}]}𝔅_{[1^{}a][cd]}+V_{[1^{\prime \prime }22^{}]}𝔅_{[2^{}a][cd]}}{r_{1b}r_{2b}+r_{12}}\hfill \\ \hfill +& (eQ_e)^2\frac{V_{[b11^{}]}}{r_{1b}}𝔅_{[1^{}2^{}][cd]}\frac{U_{[2^{}2a]}}{r_{2a}}+(12).\hfill \end{array}$$
(131)
The total ISR $`\widehat{\beta }_{2\{II\}}`$ is the sum of the three
$$\widehat{\beta }_{2\{II\}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2})=\widehat{\beta }_{2\{II\}}^{\mathrm{Double}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2})+\widehat{\beta }_{2\{II\}}^{\mathrm{Single}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2})+\widehat{\beta }_{2\{II\}}^{\mathrm{Rest}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2}).$$
(132)
#### 3.6.2 2-real FSR photons
The case of final-state double real photon emission can be analysed in a similar way. The second-order FSR, two-photon, matrix element is
$$\begin{array}{cc}\hfill _{2\{FF\}}^{(2)}& \left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2};P\right)=ie^2\underset{B=\gamma ,Z}{}\mathrm{\Pi }_B^{\mu \nu }(P)(G_{e,\mu }^B)_{[ba]}(eQ_f)^2\overline{u}(p_c,\lambda _c)\{\hfill \\ & \overline{)}ϵ_{[1]}^{}\frac{(\overline{)}p_c+m)+\overline{)}k_1}{2k_1p_c}\overline{)}ϵ_{[2]}^{}\frac{(\overline{)}p_c+m)+\overline{)}k_1+\overline{)}k_2}{2k_1p_c+2k_2p_c+2k_1k_2}G_{f,\nu }^B\hfill \\ & +G_{f,\nu }^B\frac{(\overline{)}p_d+m)\overline{)}k_1\overline{)}k_2}{2k_1p_d+2k_2p_d+2k_1k_2}\overline{)}ϵ_{[1]}^{}\frac{(\overline{)}p_d+m)\overline{)}k_2}{2k_2p_d}\overline{)}ϵ_{[2]}^{}\hfill \\ & +\overline{)}ϵ_{[1]}^{}\frac{(\overline{)}p_c+m)+\overline{)}k_1}{2k_1p_c}G_{f,\nu }^B\frac{(\overline{)}p_d+m)\overline{)}k_2}{2k_2p_d}\overline{)}ϵ_{[2]}^{}+(12)\}v(p_d,\lambda _d),\hfill \end{array}$$
(133)
Similarly, the expansion into soft and non-IR parts for FSR spin amplitudes is done in the way completely analogous to the ISR case. The subtraction formula is now
$$\begin{array}{cc}& \widehat{\beta }_{2\{FF\}}^{(2)}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2};P\right)=𝔐_{2\{FF\}}^{(2)}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2};P\right)\widehat{\beta }_{1\{F\}}^{(1)}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};P\right)𝔰_{\sigma _2}^{\{F\}}(k_2)\hfill \\ & \widehat{\beta }_{1\{F\}}^{(1)}\left({}_{\lambda }{}^{p}{}_{\sigma _2}{}^{k_2};P\right)𝔰_{\sigma _1}^{\{F\}}(k_1)\widehat{\beta }_0^{(0)}\left({}_{\lambda }{}^{p};P\right)𝔰_{\sigma _1}^{\{F\}}(k_1)𝔰_{\sigma _2}^{\{F\}}(k_2).\hfill \end{array}$$
(134)
First we obtain the contribution from terms with two $`(\overline{)}pm)`$ factors
$$\begin{array}{cc}& \widehat{\beta }_{2\{FF\}}^{(2)\mathrm{Double}}[{}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2}]=_{2\{FF\}}^{\mathrm{DoubleIR}}[{}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2}]𝔰_{[1]}^{\{F\}}𝔰_{[2]}^{\{F\}}𝔅_{[ba][cd]}\frac{(p_c+p_d+k_1+k_2)^2}{(p_c+p_d)^2}\hfill \\ & =\left(\mathrm{\Delta }_c𝔰_{[1]}^{(c)}𝔰_{[2]}^{(c)}+\mathrm{\Delta }_d𝔰_{[1]}^{(d)}𝔰_{[2]}^{(d)}\right)𝔅_{[ba][cd]}\hfill \\ & 𝔰_{[1]}^{\{F\}}𝔰_{[2]}^{\{F\}}𝔅_\mathrm{F}[{}_{\lambda _b}{}^{p_b}{}_{\lambda _a}{}^{p_a}](\frac{(p_c+p_d+k_1+k_2)^2}{(p_c+p_d)^2}1)\hfill \\ & 𝔰_{\sigma _i}^{(c)}(k_i)𝔰_{[i]}^{(c)}=+eQ_f\frac{b_{\sigma _i}(k_i,p_c)}{r_{ic}},𝔰_{\sigma _i}^{(d)}(k_i)𝔰_{[i]}^{(d)}=eQ_f\frac{b_{\sigma _i}(k_i,p_d)}{r_{id}},\hfill \\ & 𝔰_{\sigma _i}^{\{F\}}(k_i)𝔰_{\sigma _i}^{(c)}(k_i)+𝔰_{\sigma _i}^{(d)}(k_i)𝔰_{[i]}^{(c)}+𝔰_{[i]}^{(d)},\hfill \end{array}$$
(135)
and is explicitly IR-finite. The second group of terms with only one $`(\overline{)}pm)`$ factor at the end of the fermion line is
$$\begin{array}{cc}\hfill & {}_{2\{FF\}}{}^{\mathrm{SingleIR}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2})=ie^2\underset{B=\gamma ,Z}{}\mathrm{\Pi }_B^{\mu \nu }(X)(G_{e,\mu }^B)_{[ba]}(eQ_f)^2\overline{u}(p_c,\lambda _c)\{\hfill \\ & \overline{)}ϵ_{[1]}^{}\frac{(\overline{)}p_c+m)}{r_{1c}}\overline{)}ϵ_{[2]}^{}\frac{\overline{)}k_1+\overline{)}k_2}{r_{1c}+r_{2c}+r_{12}}G_{f,\nu }^B+G_{f,\nu }^B\frac{\overline{)}k_1\overline{)}k_2}{r_{1d}+r_{2d}+r_{12}}\overline{)}ϵ_{[1]}^{}\frac{(\overline{)}p_d+m)}{r_{2d}}\overline{)}ϵ_{[2]}^{}\hfill \\ \hfill +& \overline{)}ϵ_{[1]}^{}\frac{(\overline{)}p_c+m)}{r_{1c}}G_{f,\nu }^B\frac{\overline{)}k_2}{r_{2d}}\overline{)}ϵ_{[2]}^{}+\overline{)}ϵ_{[1]}^{}\frac{\overline{)}k_1}{r_{1c}}G_{f,\nu }^B\frac{(\overline{)}p_d+m)}{r_{2d}}\overline{)}ϵ_{[2]}^{}+(12)\}v(p_d,\lambda _d),\hfill \end{array}$$
(136)
and it translates in the matrix notation (in fermion spin indices) into
$$\begin{array}{cc}\hfill & {}_{2\{FF\}}{}^{\mathrm{SingleIR}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2})=\hfill \\ \hfill =& eQ_f𝔰_{[1]}^{(c)}\frac{U_{[c21^{}]}}{r_{1c}+r_{2c}+r_{12}}𝔅_{[ba][1^{}d]}+eQ_f𝔰_{[1]}^{(c)}\frac{U_{[c22^{}]}}{r_{1c}+r_{2c}+r_{12}}𝔅_{[ba][2^{}d]}\hfill \\ \hfill +& eQ_f𝔅_{[ba][c1^{}]}\frac{V_{[1^{}1d]}}{r_{1d}+r_{2d}+r_{12}}𝔰_{[2]}^{(d)}+eQ_f𝔅_{[ba][c2^{}]}\frac{V_{[2^{}1d]}}{r_{1d}+r_{2d}+r_{12}}𝔰_{[2]}^{(d)}\hfill \\ \hfill +& eQ_f𝔰_{[1]}^{(c)}𝔅_{[ba][c2^{}]}\frac{V_{[2^{}2d]}}{r_{2d}}+eQ_f\frac{U_{[c11^{}]}}{r_{1c}}𝔅_{[ba][1^{}d]}𝔰_{[2]}^{(d)}+(12),\hfill \end{array}$$
(137)
On the other hand the single-IR part to be eliminated is
$$\begin{array}{cc}\hfill \widehat{\beta }_{1(0)[1]}^{(1)}& 𝔰_{[2]}^{\{F\}}+\widehat{\beta }_{1(0)[2]}^{(1)}𝔰_{[1]}^{\{F\}}=r_{[1]}^{\{F\}}𝔰_{[2]}^{\{F\}}+r_{[2]}^{\{F\}}𝔰_{[1]}^{\{F\}}\hfill \\ \hfill =& (+eQ_e𝔅_{[ba][1^{}d]}\frac{U_{[c11^{}]}}{r_{1c}}eQ_e\frac{V_{[1^{}1d]}}{r_{1d}}𝔅_{[ba][c1^{}]})𝔰_{[2]}^{\{F\}}+(12)\hfill \\ & 𝔅_{[ba][cd]}(\frac{(p_c+p_d+k_1)^2}{(p_c+p_d)^2}1)𝔰_{[1]}^{\{F\}}𝔰_{[2]}^{\{F\}}+(12)\hfill \end{array}$$
(138)
Altogether we get
$$\begin{array}{cc}\hfill \widehat{\beta }_{2\{FF\}}^{\mathrm{Single}}& ({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2})=_{2\{FF\}}^{\mathrm{SingleIR}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2})\widehat{\beta }_{1(0)}^{(1)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1})𝔰^{\{F\}}[{}_{\sigma _2}{}^{k_2}]\widehat{\beta }_{1(0)}^{(1)}({}_{\lambda }{}^{p}{}_{\sigma _2}{}^{k_2})𝔰^{\{F\}}[{}_{\sigma _1}{}^{k_1}]\hfill \\ \hfill =& eQ_f𝔰_{[1]}^{(c)}\left\{\left(\frac{U_{[c22^{}]}}{r_{2c}+r_{1c}+r_{12}}\frac{U_{[c22^{}]}}{r_{2c}}\right)𝔅_{[ba][2^{}d]}+\frac{U_{[c21^{}]}}{r_{2c}+r_{1c}+r_{12}}𝔅_{[ba][1^{}d]}\right\}\hfill \\ \hfill +& eQ_f\left\{𝔅_{[ba][c1^{}]}\left(\frac{V_{[1^{}1d]}}{r_{1d}+r_{2d}+r_{12}}\frac{V_{[1^{}1d]}}{r_{1d}}\right)+\frac{V_{[2^{}1d]}}{r_{1d}+r_{2d}+r_{12}}𝔅_{[ba][c2^{}]}\right\}𝔰_{[2]}^{(d)}\hfill \\ & +𝔅_{[ba][cd]}(\frac{(p_c+p_d+k_1)^2}{(p_c+p_d)^2}1)𝔰_{[1]}^{\{F\}}𝔰_{[2]}^{\{F\}}+(12),\hfill \end{array}$$
(139)
Finally we include the remaining terms in eq. (133)
$$\begin{array}{cc}\hfill _{2\{FF\}}^{\mathrm{Rest}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2})=ie^2\underset{B=\gamma ,Z}{}& \mathrm{\Pi }_B^{\mu \nu }(X)(G_{e,\mu }^B)_{[ba]}(eQ_f)^2\overline{u}(p_c,\lambda _c)\{\hfill \\ & \overline{)}ϵ_{[1]}^{}\frac{\overline{)}k_1}{r_{1c}}\overline{)}ϵ_{[2]}^{}\frac{(\overline{)}p_c+m)+\overline{)}k_1+\overline{)}k_2}{r_{1c}+r_{2c}+r_{12}}G_{f,\nu }^B\hfill \\ & +G_{f,\nu }^B\frac{(\overline{)}p_d+m)\overline{)}k_1\overline{)}k_2}{r_{1d}+r_{2d}+r_{12}}\overline{)}ϵ_{[1]}^{}\frac{\overline{)}k_2}{r_{2d}}\overline{)}ϵ_{[2]}^{}\hfill \\ & +\overline{)}ϵ_{[1]}^{}\frac{\overline{)}k_1}{r_{1c}}G_{f,\nu }^B\frac{\overline{)}k_2}{r_{2d}}\overline{)}ϵ_{[2]}^{}+(12)\}v(p_d,\lambda _d),\hfill \end{array}$$
(140)
which in the programmable matrix notation looks as follows
$$\begin{array}{cc}\hfill \widehat{\beta }_{2\{FF\}}^{\mathrm{Rest}}& ({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2})=_{2\{FF\}}^{\mathrm{Rest}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2})=\hfill \\ \hfill =& (eQ_f)^2\frac{U_{[c11^{\prime \prime }]}}{r_{1c}}\frac{U_{[1^{\prime \prime }2c^{}]}𝔅_{[ba][c^{}d]}+U_{[1^{\prime \prime }21^{}]}𝔅_{[ba][1^{}d]}+U_{[1^{\prime \prime }22^{}]}𝔅_{[ba][2^{}d]}}{r_{1c}+r_{2c}+r_{12}}\hfill \\ \hfill +& (eQ_f)^2\frac{𝔅_{[ba][cd^{}]}V_{[d^{}12^{\prime \prime }]}𝔅_{[ba][c1^{}]}V_{[1^{}12^{\prime \prime }]}𝔅_{[ba][c2^{}]}V_{[2^{}12^{\prime \prime }]}}{r_{1d}+r_{2d}+r_{12}}\frac{V_{[2^{\prime \prime }2d]}}{r_{2d}}\hfill \\ \hfill +& (eQ_f)^2\frac{U_{[c11^{}]}}{r_{1c}}𝔅_{[ba][1^{}2^{}]}\frac{V_{[2^{}2d]}}{r_{2d}}+(12)\hfill \end{array}$$
(141)
The total contribution from double FSR real photon emission is
$$\widehat{\beta }_{2\{FF\}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2})=\widehat{\beta }_{2\{FF\}}^{\mathrm{Double}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2})+\widehat{\beta }_{2\{FF\}}^{\mathrm{Single}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2})+\widehat{\beta }_{2\{FF\}}^{\mathrm{Rest}}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2}).$$
(142)
#### 3.6.3 1-real ISR and 1-real FSR photon
As we have seen in the previous cases of double real emission most complications are due to simultaneous emission from one fermion “leg”. The case of one real ISR and one real FSR photon is easier because there is at most one photon on one leg:
$$\begin{array}{cc}\hfill _{2\{IF\}}^{(2)}& ({}_{\lambda _a}{}^{p_a}{}_{\lambda _b}{}^{p_b}{}_{\lambda _c}{}^{p_c}{}_{\lambda _d}{}^{p_d}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2};Pk_1)=ie^2\underset{B=\gamma ,Z}{}\mathrm{\Pi }_B^{\mu \nu }(Pk_1)\hfill \\ & eQ_e\overline{v}(p_b,\lambda _b)\left(G_{e,\mu }^B\frac{\overline{)}p_a+m\overline{)}k_1}{2k_1p_a}\overline{)}ϵ_{[1]}^{}+\overline{)}ϵ_{[1]}^{}\frac{\overline{)}p_b+m+\overline{)}k_1}{2k_1p_b}G_{e,\mu }^B\right)u(p_a,\lambda _a)\hfill \\ & eQ_f\overline{u}(p_c,\lambda _c)\left(G_{f,\nu }^B\frac{\overline{)}p_d+m\overline{)}k_2}{2k_2p_d}\overline{)}ϵ_{[2]}^{}+\overline{)}ϵ_{[2]}^{}\frac{\overline{)}p_c+m+\overline{)}k_2}{2k_2p_c}G_{f,\nu }^B\right)v(p_d,\lambda _d)\hfill \end{array}$$
(143)
and the subtraction formula is now
$$\begin{array}{cc}& \widehat{\beta }_{2\{IF\}}^{(2)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2};Pk_1)=𝔐_{2\{IF\}}^{(2)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2};Pk_1)\widehat{\beta }_{1\{I\}}^{(1)}({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1};Pk_1)𝔰_{\sigma _2}^{\{F\}}(k_2)\hfill \\ & \widehat{\beta }_{1\{F\}}^{(1)}({}_{\lambda }{}^{p}{}_{\sigma _2}{}^{k_2};Pk_1)𝔰_{\sigma _1}^{\{I\}}(k_1)\widehat{\beta }_0^{(0)}({}_{\lambda }{}^{p};Pk_1)𝔰_{\sigma _1}^{\{I\}}(k_1)𝔰_{\sigma _2}^{\{F\}}(k_2).\hfill \end{array}$$
(144)
The simplicity of this contribution is manifest in the fact that $`\widehat{\beta }_{2\{IF\}}`$ is obtained by simple subtraction (omission) of all terms proportional to one or two $`(\overline{)}pm)`$ factors
$$\begin{array}{cc}\hfill \widehat{\beta }_{2\{IF\}}& \left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2};X\right)=ie^2\underset{B=\gamma ,Z}{}\mathrm{\Pi }_B^{\mu \nu }(X)\hfill \\ & eQ_e\overline{v}(p_b,\lambda _b)\left(G_{e,\mu }^B\frac{\overline{)}k_1}{r_{1a}}\overline{)}ϵ_{[1]}^{}+\overline{)}ϵ_{[1]}^{}\frac{\overline{)}k_1}{r_{1b}}G_{e,\mu }^B\right)u(p_a,\lambda _a)\hfill \\ & eQ_f\overline{u}(p_c,\lambda _c)\left(G_{f,\nu }^B\frac{\overline{)}k_2}{r_{2d}}\overline{)}ϵ_{[2]}^{}+\overline{)}ϵ_{[2]}^{}\frac{\overline{)}k_2}{r_{2c}}G_{f,\nu }^B\right)v(p_d,\lambda _d).\hfill \end{array}$$
(145)
In the computation-friendly matrix notation it reads
$$\begin{array}{cc}& \widehat{\beta }_{2\{IF\}}\left({}_{\lambda }{}^{p}{}_{\sigma _1}{}^{k_1}{}_{\sigma _2}{}^{k_2};X\right)=ie^2\underset{B=\gamma ,Z}{}\mathrm{\Pi }_B^{\mu \nu }(X)eQ_eeQ_f\hfill \\ & \times \left((G_{e,\mu }^B)_{[b1^{}]}\frac{U_{[1^{}1a]}}{r_{1a}}+\frac{V_{[b11^{}]}}{r_{1b}}(G_{e,\mu }^B)_{[1^{}a]}\right)\left((G_{f,\nu }^B)_{[c2^{}]}\frac{V_{[2^{}2d]}}{r_{2d}}+\frac{U_{[c22^{}]}}{r_{2c}}(G_{f,\nu }^B)_{[2^{}d]}\right)\hfill \\ & =eQ_eeQ_f(𝔅_{[b1^{}][c2^{}]}(X)\frac{U_{[1^{}1a]}}{r_{1a}}\frac{V_{[2^{}2d]}}{r_{2d}}+\frac{U_{[c22^{}]}}{r_{2c}}𝔅_{[b1^{}][2^{}d]}(X)\frac{U_{[1^{}1a]}}{r_{1a}}\hfill \\ & +\frac{V_{[b11^{}]}}{r_{1b}}𝔅_{[1^{}a][c2^{}]}(X)\frac{V_{[2^{}2d]}}{r_{2d}}+\frac{V_{[b11^{}]}}{r_{1b}}\frac{U_{[c22^{}]}}{r_{2c}}𝔅_{[1^{}a][2^{}d]}(X))\hfill \end{array}$$
(146)
## 4 Relations between CEEX and EEX
Having shown the CEEX and EEX schemes in a detail, we would like to compare certain important/interesting fearures of both schemes in a more detail. In particular we would like to show how the two examples of the EEX scheme can be obtained as a limiting case of CEEX, and to show the exact relation among $`\overline{\beta }`$’s of EEX and $`\widehat{\beta }`$’s of CEEX. From these considerations it will be clear that the CEEX scheme is more general than the EEX scheme.
### 4.1 Neglecting partition dependence
Let us first examine the interesting limit of CEEX in which which we drop the dependence on the partition index $`X_{\mathrm{}}P`$, where $`P=p_a+p_b`$, for example. Note that it is not in the EEX class. In this limit, in the simplest case of the $`𝒪(\alpha ^0)`$ exponentiation we have:
$$\underset{\mathrm{}𝒫}{}e^{\alpha B_4^{}(X_{\mathrm{}})}\frac{X_{\mathrm{}}^2}{s_{cd}}𝔅\left({}_{\lambda }{}^{p};X_{\mathrm{}}\right)\underset{i=1}{\overset{n}{}}𝔰_{[i]}^{\{\mathrm{}_i\}}e^{\alpha B_4}𝔅\left({}_{\lambda }{}^{p};P\right)\underset{i=1}{\overset{n}{}}(𝔰_{[i]}^{\{0\}}+𝔰_{[i]}^{\{1\}}),$$
(147)
because of the identity
$$\underset{\mathrm{}𝒫}{}\underset{i=1}{\overset{n}{}}𝔰_{[i]}^{\{\mathrm{}_i\}}\underset{i=1}{\overset{n}{}}\left(𝔰_{[i]}^{\{0\}}+𝔰_{[i]}^{\{1\}}\right).$$
The relevance, advantages and disadvanteges of this scenario were already discussed in Section ??? Is this realized in $`𝒦𝒦`$MC????? Note that in the above transition we keep the ISR$``$FSR interference contribution.
### 4.2 Neglecting IFI
The second important case we would like to discuss is the case of the very narrow resonances, when the ISR$``$FSR interference contribution to any physicaal observable is so small that it can be neglected whatsoever. This corresponds to a well defined limit in the CEEX scheme. In this limit, in the simplest case of the $`𝒪(\alpha ^0)`$ exponentiation we have:
$$\begin{array}{cc}\hfill |_n^{(0)}|^2& =\underset{\mathrm{}𝒫}{}\underset{\mathrm{}^{}𝒫}{}e^{\alpha B_4^{}(X_{\mathrm{}})}e^{\alpha (B_4^{}(X_{\mathrm{}^{}}))^{}}𝔅\left({}_{\lambda }{}^{p};X_{\mathrm{}}\right)𝔅\left({}_{\lambda }{}^{p};X_{\mathrm{}^{}}\right)^{}\underset{i=1}{\overset{n}{}}𝔰_{[i]}^{\{\mathrm{}_i\}}\underset{j=1}{\overset{n}{}}𝔰_{[j]}^{\{\mathrm{}_i^{}\}}^{}\hfill \\ & e^{2\alpha \mathrm{}B_2(p_a,p_b)}e^{2\alpha \mathrm{}B_2(p_c,p_d)}\underset{\mathrm{}𝒫}{}\left|𝔅\left({}_{\lambda }{}^{p};X_{\mathrm{}}\right)|^2\underset{i=1}{\overset{n}{}}\right|𝔰_{[i]}^{\{\mathrm{}_i\}}|^2.\hfill \end{array}$$
(148)
What we did in the above transition, we have neglected the ISR$``$FSR interferences entirely by dropping non-diagonal terms $`\mathrm{}\mathrm{}^{}`$ in the double sum over partitions, and we have replaced the resonance-type form-factor by the sum of the traditional YFS formfactors for ISR and FSR (no interferenence). In this way we have got the $`𝒪(\alpha ^0)_{\mathrm{EEX}}`$ which at this order is identical to $`𝒪(\alpha ^0)_{\mathrm{CEEX}}`$. At the $`𝒪(\alpha ^r)_{\mathrm{CEEX}}`$ $`r=1,2`$, in order to get from $`𝒪(\alpha ^r)_{\mathrm{CEEX}}`$ to $`𝒪(\alpha ^r)_{\mathrm{EEX}}`$ we have in addition to truncate $`\widehat{\beta }`$’s down to $`\overline{\beta }`$’, as will be shown in the next subsection.
The $`𝒪(\alpha ^r)_{\mathrm{EEX}}`$ $`r=1,2`$ neglecting the ISR$``$FSR interferences was used in YFS2/3 of KORALZ and it is well justified close to $`Z`$ resonance position at LEP1, see also relevant numerical results in the next Section. At LEP2 the above approximation cannot be justified any more.
### 4.3 Relation among $`\overline{\beta }`$’s for EEX and $`\widehat{\beta }`$’s of CEEX
For the sake of completeness of the discussion, it is necessary to find out the relation of the $`\beta `$’s defined at the amplitude level to the older EEX/YFS $`\overline{\overline{\beta }}`$’s defined at the level of the differential distributions. Let us suppress all spin indices, understanding that for every term like $`|\mathrm{}|^2`$ or $`\mathrm{}[AB^{}]`$ the appropriate spin sum/average is done. The traditional $`\overline{\beta }`$’s of the YFS scheme at the $`𝒪(\alpha ^2)`$ level are
$$\begin{array}{cc}& \overline{\beta }_0^{(l)}=\left|𝔐_0^{(l)}\right|_{(\alpha ^l)}^2,l=0,1,2,\hfill \\ & \overline{\beta }_1^{(l)}(k)=\left|𝔐_1^{(l)}(k)\right|_{(\alpha ^{l+1})}^2\overline{\beta }_0^{(l)}|𝔰(k)|^2,l=0,1,\hfill \\ & \overline{\beta }_2^{(2)}(k_1,k_2)=\left|𝔐_1^{(2)}(k_1,k_2)\right|^2\overline{\beta }_1^{(1)}(k_1)|𝔰(k_2)|^2\overline{\beta }_1^{(1)}(k_2)|𝔰(k_1)|^2\overline{\beta }_0^{(0)}|𝔰(k_1)||𝔰(k_2)|^2,\hfill \end{array}$$
(149)
where subscript $`|_{(\alpha ^r)}`$ means truncation to $`𝒪(\alpha ^r)`$. Now for each $`𝔐_n^{(n+l)}`$we substitute its expansion in terms of $`\widehat{\beta }`$’s according to eq. (48) getting the following relation
$$\begin{array}{cc}& \overline{\beta }_0^{(l)}=|\widehat{\beta }_0^{(l)}|_{(\alpha ^l)}^2,l=0,1,2,\hfill \\ & \overline{\beta }_1^{(l)}(k)=|\widehat{\beta }_1^{(l)}(k)|^2+2\mathrm{}[\widehat{\beta }_0^{(l)}(\widehat{\beta }_1^{(l)}(k))^{}]_{(\alpha ^{l+1})},l=0,1,\hfill \\ & \overline{\beta }_2^{(2)}(k_1,k_2)=|\widehat{\beta }_2^{(2)}(k_1,k_2)|^2+2\mathrm{}[\widehat{\beta }_1^{(1)}(k_1)𝔰(k_2)\{\widehat{\beta }_1^{(1)}(k_2)𝔰(k_1)\}^{}]\hfill \\ & +2\mathrm{}[\widehat{\beta }_2^{(2)}(k_1,k_2)\{\widehat{\beta }_1^{(1)}(k_1)𝔰(k_2)\widehat{\beta }_1^{(1)}(k_2)𝔰(k_1)+\widehat{\beta }_1^{(0)}𝔰(k_1)𝔰(k_2)\}^{}],\hfill \end{array}$$
(150)
As we see, the relation is not completely trivial; there are some extra terms on the RHS, which are all IR-finite. From the above exercise it is obvious that $`\overline{\beta }`$’s are generally more complicated objects than the $`\widehat{\beta }`$’s and that for example the inclusion of the spin density matrix formalism into the $`\overline{\beta }`$’s would be a quite nontrivial exercise – the great advantage of the CEEX scheme is that this is done numerically. It is also seen that in the $`\overline{\beta }_0`$ and $`\overline{\beta }_1`$ some higher-order virtual terms are unnecessarily truncated, which probably is worsening perturbative convergence of the EEX/YFS scheme in comparison with CEEX. The above formula shows in a most clear and clean way the difference between the EEX and CEEX exponentiation schemes.
## 5 Semi-analytical approach
In this section we shall present results of semi-analytical calculations which reproduce in the $`𝒪(\alpha ^2)_{prag}`$, or even in the $`𝒪(\alpha ^3)_{prag}`$, certain selected results, mainly integrated cross-sections, of the Monte Carlo calculation.
The semi-analytical approach in which one integrates over the phase space analytically, often leaving the last one- or two-dimensional integrations for numerical treatment (usually non Monte Carlo), is the oldest one. Four decades ago there were no computers powerful enough even to dream about the numerical integration over the complete multiphoton phase space. Even now, in spite of proliferation of the MC programs, the non Monte Carlo semianalytical programs are still very popular and useful, especially programs used to calculate the total cross section and charge asymmetry for the fermion-pair final state, near the $`Z`$ resonance, like ZFITTER and TOPAZ0 . Semi-analytical programs have certain advantages over the MC programs – they are generally faster in terms of computer CPU time and are therefore better suited for fitting input parameters of the Standard Model, like the Higgs mass<sup>17</sup><sup>17</sup>17 It is definitely possible to fit input SM parameters with help of the Monte Carlo event generators, as it is currently done in the $`W`$-mass measurement in LEP2.. Nevertheless, semi-analytical calculations have also two long-standing important disadvantages:
* They are able to provide predictions for rather very primitive or absent experimental cut-offs. In practice they always have to be used in combination with the MC event generators anyway. The MC is used to remove or “straighten” the real experimental cuts to be closer to these which are practically implementable in the semianalytical calculation. Obviously this introduces additional systematic errors.
* They are prohibitively complicated beyond the three-body final state, that is they are relatively easy up to $`𝒪(\alpha ^1)`$, (single photon emission in the fermion pair production). In the collinear approximation (structure function method) one is able to add the effects due to emission of the second photon and further photons; however, this can improve the precision only within the leading logarithmic scope and makes the introduction of the realistic cuts even harder.
The other important role of the semianalytical calculations is to provide the numerical tests of the Monte Carlo programs. They are typically used to check technical correctness of the phase space integration, the so-called technical precision, and also correctness of the implementation of the SM matrix element. In the following we shall see examples of both kinds of tests. In particular, we shall see the test of the technical precision of $`𝒦𝒦`$MC at the $`210^4`$ level based of the semianalytical formula obtained in this section, in the case of a single kinematical cut on the total energy of all photons.
The role of semianalytical calculations as a test of the Monte Carlo programs is important but limited. One can easily imagine the situation in which the numerical problem shows up not for simple kinematical cuts, for which the “calibration” with help of the semianalytical program has been done, but for more realistic and complicated cuts. It is therefore always true that the ultimate test of the MC calculation is always the comparison of two MC programs, because it can be done for arbitrary cut-offs<sup>18</sup><sup>18</sup>18 This kind of test was for instance done for the first modern $`𝒪(\alpha ^1)`$ Monte Carlo event generator MUSTRAAL of ref. , with the very high precision at that time of 1%.. This recipe may look often prohibitively expensive in the effort required to realize it; however, at the sub-permille precision level the amount of work required for the realization of the semianalytical formulas and testing the semianalytical program is probably comparable. So altogether, for the sub-permille precision prediction the approach with two independent MC programs seems to be the most economical one. In any case the semi-analytical calculations will be always very useful especially when the precision requirements are not excessive and we do not deal with the observables involving complicated experimental cuts.
### 5.1 Inclusive exponentiation, IEX
An important ingredient in many semi-analytical calculations is the “exponentiation”. The meaning and the technique of exponentiation in the context of the semi-analytical calculations is however different from the exclusive exponentiation discussed in most of this paper. Let us elaborate more on what the exponentiation really means in the semianalytical approach.
As already discussed in ref. , in the typical semianalytical approach one is practicing what we call an “ad-hoc exponentiation” or “naive exponentiation” procedure in which one takes the QED finite-order, let us say $`𝒪(\alpha )`$ or $`𝒪(\alpha ^2)`$, analytical result for a certain one- or two-dimensional inclusive distribution and this result is “improved” by hand, in such a way that the soft limit (the limit in which hard photons are eliminated) in the resulting distribution conforms to the Yennie-Frautschi-Suura work . The well known examples of the ad-hoc exponentiation are presented in refs. and later in ref. for the initial-state bremsstrahlung in $`e^+e^{}`$ annihilation. There are also many other examples of the ad-hoc exponentiation, including calculations for the deep inelastic and Bhabha scattering processes. The ad-hoc exponentiation procedure may get improved gradually by taking into account higher-order effects. For example, the $`𝒪(\alpha )`$ procedure of ref. was extended to $`𝒪(\alpha ^2)`$ in ref. and later to $`𝒪(\alpha ^3)`$ in ref. . The problem is that the procedure of ad-hoc exponentiation is essentially rather art than science – one may regard it at best as a “by hand interpolation” between two kinds of analytical formulas – one valid in the soft photon limit and another in the hard photon limit, see ref. for more detailed discussion on this interpretation. Since this approach is not systematic, it is therefore difficult to estimate the uncertainty of the results and it has to be “reinvented” for each perturbative order and for each inclusive observable again and again.
The self-suggesting question is therefore: Is there any better and more systematic way of re-formulating the ad-hoc exponentiation for any inclusive distributions in the analytical form? I would be also desirable to find a direct connection to the exclusive exponentiation YFS exponentiation (of the EEX or CEEX type) which is discussed and implemented in this work. There is an obvious hint in which direction to go. If all soft photons are soft, then we know exactly the analytical formula for the multi-photon phase space integral,
$$\begin{array}{cc}\hfill f(\gamma ,V)& =e^{\gamma \mathrm{ln}\epsilon }\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{i=1}{\overset{n}{}}\underset{k_0>2\epsilon \sqrt{s}/2}{}\frac{d^3k_i}{k_i^0}\stackrel{~}{S}(p_1;p_2;k_i)\delta \left(V\frac{1}{s}(p_1+p_2)\left(\underset{i=1}{\overset{n}{}}k_i\right)\right)\hfill \\ & =\frac{e^{C\gamma }}{\mathrm{\Gamma }(1+\gamma )}\gamma V^{\gamma 1}=F(\gamma )\gamma V^{\gamma 1},C=0.57721566\mathrm{},\hfill \end{array}$$
already obtained in the original YFS paper . So why not include hard photons in the game? Let us therefore define the “YFS inclusive exponentiation” as a result of the analytical phase space integration of the distributions of the YFS exclusive exponentiation:
YFS inclusive exponentiation $``$ Analytical integration of YFS multiphoton integrals.
In this way we have got a clear and clean connection between the YFS exponentiation (as implemented in the Monte Carlo) and the “YFS inclusive exponentiation” – the connection is simply the analytical phase-space integration! As a result, we do not need any obscure “recipes” any more and what we only need to know is how to integrate (analytically) the phase-space.
The above looks promising but the YFS exclusive exponentiation involves non-trivial integrals over the multiphoton phase-space, this is why it is implemented in the form of the Monte Carlo integration/simulation numerical program. The relevant integrals over $`n`$ real photon phase space do not seem at first sight to be treatable analytically at all.
The situation is not so hopeless, however, as it may look at first sight, and in the following we shall present the solution. We start, as promised, from the full YFS expression, the same as in the Monte Carlo and are able to do “the impossible” — that is to perform the phase-space integral analytically. We calculate the phase-space integral approximately. We shall do it, however, in such a way that in the approximate method becomes exact in the soft photon limit, The soft photon contributions are there integrated exactly and only the remaining “non-infrared” contribution will be calculate using approximate methods, typically the leading-logarithmic (LL) collinear approximation. The LL approximations in non-IR parts may concern both the phase-space parametrization and the matrix element<sup>19</sup><sup>19</sup>19 Let us note that the LL evaluation of the phase-space integral was already employed to some extent in the original YFS work . At that time, because of lack of fast computers, it was the only accessible method..
The profit from the above approach is two-fold: Contrary to the traditional ad-hoc exponentiation we gain, for a given exponentiated inclusive distribution, a clear and clean connection between the YFS exponentiation (as implemented in the Monte Carlo) and the “YFS inclusive exponentiation”, we do not need any obscure “recipes” any more. It means that the resulting inclusive exponentiation (IEX) is now a systematic order-by-order procedure, this feature is simply inherited from the exclusive YFS exponentiation.
As already stressed, the inclusive YFS exponentiation will never fully replace the YFS Monte Carlo because it is possible to deal analytically with only a very limited number of the distributions, without cuts or with very simple cuts, while in the Monte Carlo one may calculate an arbitrary distribution in the presence of the most complicated cuts.
In the following subsections we shall show explicitly the analytical integrations leading to $`𝒪(\alpha ^2)_{prag}`$ IEX results. We shall typically compare the Monte Carlo with EEX matrix element and IEX formulas, both in the $`𝒪(\alpha ^2)_{prag}`$ class. Their difference will be then necessarily of $`𝒪(\alpha ^3)_{prag}`$, i.e. up to factor 10 smaller – quite a strong test of both calculations. On one occasion, we shall go to a more difficult level of the $`𝒪(\alpha ^3)_{prag}`$, in which case the difference between MC and IEX is of order $`𝒪(\alpha ^4)_{prag}`$.
Finally, let us note that the set of IEX formulas presented in this section was used over many years as a basic test of the precision of the YFS2 and KORALZ/YFS3 programs. Some of the IEX results were already shown in ref. and . Most of them are, however, shown here for the first time. In the mean time the analogous set of IEX results was obtained and published for the $`t`$-channel dominated process . In fact ref. shows an even more sophisticated case in the $`𝒪(\alpha ^3)_{prag}`$ class than the $`𝒪(\alpha ^2)_{prag}`$ results presented here. Using the experience of ref. it would be definitely possible, for the $`s`$-channel process of this paper, to upgrade systematically the calculation of this section to $`𝒪(\alpha ^3)_{prag}`$, both for ISR and FSR.
### 5.2 Semi-analytical formulas for ISR
We shall start the construction in IEX expressions with the ISR case, first showing the basic techniques working out the example with the $`𝒪(\alpha ^0)`$ EEX matrix element. In this case the multiphoton differential distribution is just the Born cross section times real soft-factors. While for the other IEX formulas the phase space will be integrated basically in the $`𝒪(\alpha ^2)_{prag}`$, the case of the of $`𝒪(\alpha ^0)_{\mathrm{EEX}}`$ we shall make more effort and do it in the $`𝒪(\alpha ^3)_{prag}`$, like in ref. . Let us attract attention of the reader to the fact that we have the matrix element in the $`𝒪(\alpha ^0)_{\mathrm{EEX}}`$ and the phase space integration is in $`𝒪(\alpha ^2)_{prag}`$ or $`𝒪(\alpha ^3)_{prag}`$. There is no contradiction in this, as we shall see in the following.
#### 5.2.1 Baseline high precision results for $`𝒪(\alpha ^{(0)})_{\mathrm{EEX}}`$
The complete $`𝒪(\alpha ^2)_{prag}`$ calculation/exponentiation according to the rules layed down in the beginning of this Section will be rather involved, let us therefore illustrate our calculational methods with the simplest possible example. Even this simple example features some nontrivial technical features and we shall therefore present two versions of the calculation.
The basic example discussed in the following is the $`𝒪(\alpha ^0)`$ initial-state YFS inclusive exponentiation. In the master equation (4) we set the charge of the final fermion to zero, $`Q_f=0`$ and we replace the sum of $`\overline{\beta }`$’s with the $`𝒪(\alpha ^0)`$ version of $`\overline{\beta }_0`$ proportional to the Born differential cross section
$$\overline{\beta }_0^{(0)}(q_1,q_2)=\frac{2}{\beta _f}\frac{d\sigma ^{\mathrm{Born}}}{d\mathrm{\Omega }}((q_1+q_2)^2,\vartheta ),\beta _f=\left(14m_f^2/(q_1+q_2)^2\right)^{1/2},$$
(151)
where the normalization is such that
$$\frac{d^3q_1}{q_1^0}\frac{d^3q_2}{q_2^0}\delta ^{(4)}(Xq_1q_2)\overline{\beta }_0^{(0)}(q_1,q_2)=\sigma ^{\mathrm{Born}}\left((q_1+q_2)^2\right).$$
(152)
The initial-state $`𝒪(\alpha ^0)`$ YFS formula reads
$$\begin{array}{cc}\hfill \sigma _0& =\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\frac{d^3q_1}{q_1^0}\frac{d_2^q}{q_2^0}\underset{i=1}{\overset{n}{}}\frac{d^3k_i}{k_i^0}\stackrel{~}{S}_I(k_i)\mathrm{\Theta }\left(k_i^0\frac{1}{2}\epsilon \sqrt{s}\right)\hfill \\ & \delta ^{(4)}(p_1+p_2q_1q_2\underset{j}{}k_j)e^{Y_I(\epsilon )}\overline{\beta }_0^{(0)}((q_1+q_2)^2,\vartheta _0)\hfill \end{array}$$
(153)
Integration over the final-state fermion 2-body phase space is done trivially leading to
$$\sigma _0=_0^1𝑑v\sigma ^{\mathrm{Born}}(s(1v))e^{\delta _{YFS}}\rho _0(v)$$
(154)
where the essential multiphoton integral
$$\rho _0(v)=e^{\gamma \mathrm{ln}\epsilon }\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\frac{d^3q_1}{q_1^0}\frac{d^3q_2}{q_2^0}\underset{i=1}{\overset{n}{}}\underset{k_i^0>\epsilon \sqrt{s}/2}{}\frac{d^3k_i}{k_i^0}\stackrel{~}{S}_I(k_i)\delta \left(1v\frac{1}{s}(p_1+p_2\underset{j}{}k_j)^2\right)$$
(155)
is the main object of our interest. Note that we have split $`Y_I(\epsilon )=\gamma \mathrm{ln}\epsilon +\delta _{YFS}`$.
The QED matrix element, beyond the soft photon integral, is in this simplified case totally absent. The inclusive YFS exponentiation, as defined above, amounts to calculating analytically the multiphoton phase-space integral for $`\rho _0(v)`$. As explained above, we shall do it in $`𝒪(\alpha ^2)_{prag}`$, but we shall keep the proper soft limit undestroyed. Let us note first that in the soft limit $`v0`$ the function $`\rho _0(v)`$ coincides with the soft integral of eq. (5.1), i.e. $`\rho (v)f(\gamma ,v)`$. Since the most singular part in this limit is known we isolated it and we expect the $`𝒪(\alpha ^2)_{prag}`$ result to be in form
$$\rho _0(v)=f(\gamma ,v)(1+v\gamma f_1(v))$$
(156)
where $`f_1(v)`$ is nonsingular. How does one find the function $`f_1(v)`$? Let us inspect the difference
$$\begin{array}{cc}\hfill d_0(v)& =\rho (v)f(\gamma ,v)\hfill \\ & =\frac{1}{2!}\frac{d^3k_1}{k_1^0}\stackrel{~}{S}_I(k_1)\frac{d^3k_2}{k_2^0}\stackrel{~}{S}_I(k_2)\hfill \\ & \left[\delta \left(1v\frac{1}{s}(p_1+p_2\underset{j}{}k_j)^2\right)\delta \left(v\frac{1}{s}(p_1+p_2)\left(\underset{i=1}{\overset{n}{}}k_i\right)\right)\right]\hfill \end{array}$$
(157)
This new object has rather interesting properties. First of all, the $`𝒪(\alpha ^1)`$ integrals cancel exactly and the first nontrivial integral is of $`𝒪(\alpha ^2)`$. This second-order integral is not, however, infrared divergent! According to our rules we are therefore allowed, without any danger of spoiling the soft limit, to calculate it in the leading-logarithmic approximation.
Let us present now our first of two methods of calculating $`\rho _0(v)`$. In the LL approximation we replace collinear singularities in the photon angle $`\vartheta _\gamma =0,\pi `$ by $`\delta `$-like peaks
$$\begin{array}{cc}\hfill \frac{d^3k_i}{k_i^0}\stackrel{~}{S}_I(k_i)& =\frac{\alpha }{2\pi ^2}\underset{0}{\overset{1}{}}\frac{dx_i}{x_i}\underset{1}{\overset{1}{}}𝑑c_i\frac{s_i^2}{(1\beta _e^2c_i^2)^2}\underset{0}{\overset{2\pi }{}}𝑑\varphi _i\hfill \\ & \underset{0}{\overset{1}{}}\frac{dx_i}{x_i}𝑑c_i\left[\frac{1}{2}\gamma \delta (c_i1)+\frac{1}{2}\gamma \delta (c_i+1)\right]\hfill \end{array}$$
(158)
where
$$\beta _e=(14m_e^2/s)^{1/2},c_i=\mathrm{cos}\theta _i,s_i=\mathrm{sin}\theta _i,i=1,2,$$
and using the above LL substitution we get
$$d_0(v)=\underset{\epsilon 0}{lim}\frac{\gamma ^2}{4}\underset{\epsilon }{\overset{1}{}}\frac{dx_1}{x_1}\underset{\epsilon }{\overset{1}{}}\frac{dx_2}{x_2}\left[\delta \left(v(1x_1)(1x_2)\right)\delta (vx_1x_2)\right]$$
(159)
Two immediate remarks are in order: Out of four terms in the product $`[\delta (c_11)+\delta (c_1+1)][\delta (c_21)\delta (c_2+1)]`$ only two contribute, these with two anticollinear photons, $`c_1=1,c_2=1`$ and $`c_1=1,c_2=1`$. The result of integration depends critically on the careful regularization and for this reason we keep explicitly the $`\epsilon `$ infrared regulator. A quick calculation gives a zero value for the integral. The very similar phenomenon is present in the calculation of $`f(\gamma ,v)`$ where a naive calculation up to second-order gives the $`v^{\gamma 1}`$ factor only. The remaining $`F(\gamma )=1\frac{\pi ^2}{12}\gamma ^2+\mathrm{}`$ factor comes from careful consideration of the $`k^0>\epsilon \sqrt{s}/2`$ condition for two photons. With our proper regularization we obtain
$$d_0(v)=\frac{\gamma ^2}{4}\frac{\mathrm{ln}(1v)}{v},$$
(160)
which is finite in the $`v0`$ limit.
Now comes the second calculation method which will be often employed in the following. In this variant we take into account the influence of additional soft photons (in addition to the two hard ones). They do not change the second-order result but provide the proper infrared regulation replacing the former $`\epsilon `$ regulator. The LL treatment of the phase-space will be a little bit different. Starting from eq. (5.1) we split (in CMS frame) the photon integration into its forward and backward hemisphere parts
$$\frac{d^3k}{k^0}=\underset{\theta >\pi /2}{}\frac{d^3k}{k^0}+\underset{\theta <\pi /2}{}\frac{d^3k}{k^0}$$
and after changing the summation order we get
$$\begin{array}{cc}\hfill f(\gamma ,v)& =e^{\gamma \mathrm{ln}\epsilon }\underset{n}{}\underset{n^{}}{}\frac{1}{n!}\frac{1}{n^{}!}\underset{i=1}{\overset{n}{}}\underset{\theta _i>\pi /2}{}\frac{d^3k_i^+}{k_i^{+0}}\stackrel{~}{S}_I(k_i^+)\mathrm{\Theta }\left(k_i^{+0}\epsilon \frac{\sqrt{s}}{2}\right)\hfill \\ & \underset{j=1}{\overset{n^{}}{}}\underset{\theta _j<\pi /2}{}\frac{d^3k_j^{}}{k_j^0}\stackrel{~}{S}_I(k_j^{})\mathrm{\Theta }\left(k_j^0\epsilon \frac{\sqrt{s}}{2}\right)\delta \left(v\frac{2}{s}P(K^++K^{})\right)\hfill \end{array}$$
(161)
where $`P=p_1+p_2,K^+=\underset{i=1}{\overset{n}{}}k_i^+,K^{}=\underset{j=1}{\overset{n^{}}{}}k_j^{}`$. The above sum of integrals factorizes into two sums. Each of the sums can be evaluated exactly leading to the following simple convolution
$$f(\gamma ,v)=𝑑v_+𝑑v_{}\delta (vv_{}v_+)f(\frac{\gamma }{2},v_+)f(\frac{\gamma }{2},v_{}),$$
(162)
This identity holds for the integration result anyway, but we have also obtained it through the direct phase-space integration. So far all calculations were exact and we only reorganized the phase-space integration which will be useful in the next step. Let us consider the $`d_0(v)`$ difference again
$$\begin{array}{cc}\hfill d_0(v)& =e^{\gamma \mathrm{ln}\epsilon }\underset{n}{}\underset{n^{}}{}\frac{1}{n!}\frac{1}{n^{}!}\underset{i=1}{\overset{n}{}}\underset{\theta _i>\pi /2}{}\frac{d^3k_i^+}{k_i^{+0}}\stackrel{~}{S}_I(k_i^+)\mathrm{\Theta }\left(k_i^{+0}\epsilon \frac{\sqrt{s}}{2}\right)\hfill \\ & \underset{j=1}{\overset{n^{}}{}}\underset{\theta _j<\pi /2}{}\frac{d^3k_j^{}}{k_j^0}\stackrel{~}{S}_I(k_j^{})\mathrm{\Theta }(k_j^0\epsilon \frac{\sqrt{s}}{2})[\delta (v1+\frac{1}{s}(PK^+K^{})^2)\hfill \\ & \delta (v\frac{2}{s}P(K^++K^{}))].\hfill \end{array}$$
(163)
As before, the whole integral is finite in $`v0`$ limit and it gets the first non-zero contribution in the second-order. From the previous exercises we know that the essential second-order leading-logarithmic contribution comes from two anticollinear photons – this is why we divided photon phase-space into two hemispheres. Now, the LL approximation is realized by substituting in the first $`\delta `$
$$K^{\pm \mu }=(K^{\pm 0},0,0,\pm |K^{\pm 0}|).$$
Note that, contrary to the previous calculation, we did not modify the $`\stackrel{~}{S}`$ factors, we did not introduce collinear $`\delta `$’s in the photon angular distribution and we keep infinite numbers of photons. In spite of the apparent increase of the complication level, the integral reduces to a nice form
$$f(\gamma ,v)=𝑑v_+𝑑v_{}\left[\delta (vv_{}v_++v_+v_{})\delta (vv_{}v_+)\right]f(\frac{\gamma }{2},v_+)f(\frac{\gamma }{2},v_{}),$$
(164)
which is calculable analytically! Neglecting terms $`𝒪(\gamma ^3)`$ we obtain
$$d_0(v)=\frac{e^{C\gamma }}{\mathrm{\Gamma }(1+\gamma )}\gamma v^{\gamma 1}\frac{1}{4}\gamma \mathrm{ln}(1v).$$
(165)
Note that since in the present variant of the calculation we have treated soft photons more friendly we recovered the proper soft factor $`f(\gamma ,v)`$ as a factor in the solution.
Summarizing, the $`𝒪(\alpha ^2)_{prag}`$ phase-space integration result is
$$\rho _0(v)=\frac{e^{C\gamma }}{\mathrm{\Gamma }(1+\gamma )}\gamma v^{\gamma 1}\left(1\frac{1}{4}\gamma \mathrm{ln}(1v)\right)$$
(166)
and the corresponding cross section reads
$$\sigma _0(v_{max})=e^{\delta _{YFS}}\frac{e^{C\gamma }}{\mathrm{\Gamma }(1+\gamma )}\underset{0}{\overset{v_{max}}{}}𝑑v\sigma ^{\mathrm{Born}}(s(1v))\gamma v^{\gamma 1}\left(1\frac{1}{4}\gamma \mathrm{ln}(1v)\right)$$
(167)
The above integration methods provide us with the $`𝒪(\alpha ^2)_{prag}`$ phase-space integration result for any of the $`\overline{\beta }_0`$ contributions as listed in eqs. (LABEL:beta0). For example the contribution from $`\overline{\beta }_0^{(2)}`$ reads
$$\begin{array}{cc}\hfill \sigma _0^{(2)}& =\underset{0}{\overset{1}{}}𝑑v\sigma ^{\mathrm{Born}}\left(s(1v)\right)\rho _0^{(2)},\hfill \\ \hfill \rho _0^{(2)}& =F(\gamma )e^{\delta _{YFS}}\gamma v^{\gamma 1}(1+\delta _I^{(1)}+\delta _I^{(2)})\left(1\frac{1}{4}\gamma \mathrm{ln}(1v)\right)\hfill \end{array}$$
(168)
From now one we shall not restrict ourselves to $`𝒪(\alpha ^0)_{prag}`$ EEX matrix elements, but rather consider the complete EEX-class $`𝒪(\alpha ^2)_{prag}`$ EEX matrix elements as defined in Section 2.
The practical significance of IEX formula of eq. 168 is rather important. The biggest terms neglected in it are of $`𝒪(\gamma ^3)`$ and $`𝒪(\alpha \gamma )`$ and we expect them to stay below 0.1%. (This will be true provided they are no extra enhancement factors, see discussion below.) In other words we expect for the $`\overline{\beta }_0^{(2)}`$ contribution in the EEX matrix element in Section 2 the result of the Monte Carlo phase-space integration will agree with the formula (168) to within about 0.1% for an arbitrary cut $`v_{max}`$.
Let us check the above conjecture with the numerical exercise. In the numerical test we shall already include at this moment not only ISR $`\overline{\beta }_0^{(2)}`$ contribution of eq. (168), see also Table 2, but also the analogous FSR $`\overline{\beta }_0^{(2)}`$ contribution which will be calculated<sup>20</sup><sup>20</sup>20 We could present results of the numerical tests (which we have done) for ISR alone. They look however very much the same like as simultaneous ISR and FSR so we decided not to present their figures. in the next sub-section, see eq.(LABEL:eq:IEX-FSR) and Table 3. We consider the total cross section with the cut on the total photon energy defined by $`v_{max}`$ as follows
$$\sigma _{\overline{\beta }_0^{(2)}\overline{\beta }_0^{(2)}}^{(2)}=_0^{v_{\mathrm{max}}}𝑑v_0^{v/(1u_{\mathrm{max}})}𝑑u\sigma _{\mathrm{Born}}^f\left(s(1u)(1v)\right)\rho _{I\overline{\beta }_0^{(2)}}^{(2)}(v)\rho _{F\overline{\beta }_0^{(2)}}^{(2)}(u)$$
(169)
In order to get a clearer picture about the magnitude of the discrepancy between EEX MC and IEX formula we use the artificially flat Born cross-section $`\sigma _{\mathrm{Born}}^f(s(1u)(1v))\sigma _{\mathrm{Born}}^f(s)`$ in both. Results of the comparison are presented in Figure 1. Following our expectation the difference is well below 0.1% for the entire range of the photon energy cutoff $`v_{\mathrm{max}}`$.
The situation does not look as good when we switch-on the $`s`$-dependence in Born cross-section. In Figure 10(a) we see the relevant comparison. At the CMS energy of 189GeV the position of the Z radiative return is at $`v=0.75`$ and we clearly see a worsening there with respect to the previous case in Figure 1 where the discrepancy is now almost 0.2% (0.4% in terms of $`\sigma _{\mathrm{Born}}`$). The situation is even more dramatic in the last bin which corresponds to $`v_{\mathrm{max}}=14m_\mu ^2/s`$ and here the discrepancy among $`𝒪(\alpha ^2)_{prag}`$ IEX and $`𝒪(\alpha ^2)_{prag}`$ MC EEX is -2% of the total cross section, that is -7% in terms of the Born cross-section! This is, of course due to the $`Z`$ resonance and $`1/s`$ behaviour of the Born cross-section at very low $`s`$ (especially for the case of the $`\mu `$ channel shown in Figure 10). In order to be sure that the above effect is not due to some technical problem in the MC integration we had to improve our IEX formula and upgrade the analytical phase space integration for ISR to the level of $`𝒪(\alpha ^3)_{prag}`$. The comparison with the $`𝒪(\alpha ^3)_{prag}`$ IEX for the same EEX $`𝒪(\alpha ^2)_{prag}`$ MC we show in Figure 10(b). Now the difference is reduced to less than 0.1% everywhere, an in the last bin it is reduced from -2% to +0.2%, as expected. The additional terms of the $`𝒪(L^1\alpha ^2)`$ and $`𝒪(L^3\alpha ^3)`$ are shown in Table 2 at the end of this Section. We do not show the details on the phase space integration which provide these two additional terms. The method is generally rather similar to the one used in this Section and in ref. .
#### 5.2.2 Beta-bar-one, $`\overline{\beta }_1`$
In the following step our aim is to calculate analytically the ISR contribution to the total cross section from $`\overline{\beta }_{1I}^{(2)}`$ as given by
$$\begin{array}{cc}& \sigma =\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\frac{d^3q_1}{q_1^0}\frac{d^3q_2}{q_2^0}\underset{j=1}{\overset{n}{}}\frac{d^3k_j}{k_j^0}\stackrel{~}{S}_I(p_1,p_2;k_j)(1\mathrm{\Theta }(\mathrm{\Omega }_I;k_j))\hfill \\ & e^{Y(\mathrm{\Omega }_I)}\underset{j=1}{\overset{n}{}}\overline{\beta }_{1I}^{(2)}(X,p_1,p_2,q_1,q_2,k_j)/\stackrel{~}{S}_I(k_j)\delta ^{(4)}\left(p_1+p_2q_1q_2\underset{j=1}{\overset{n}{}}k_j\right)\hfill \\ & =\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\frac{d^3q_1}{q_1^0}\frac{d^3q_2}{q_2^0}\underset{j=1}{\overset{n}{}}\frac{d^3k_j}{k_j^0}\stackrel{~}{S}_I(p_1,p_2;k_j)(1\mathrm{\Theta }(\mathrm{\Omega }_I;k_j))e^{Y(\mathrm{\Omega }_I)}\hfill \\ & \frac{d^3k}{k^0}(1\mathrm{\Theta }(\mathrm{\Omega }_I;k))\delta ^{(4)}\left(p_1+p_2q_1q_2k\underset{j=1}{\overset{n}{}}k_j\right)\overline{\beta }_{1I}^{(2)}(q_1+q_2,p_1,p_2,q_1,q_2,k).\hfill \end{array}$$
(170)
We start again from the EEX $`𝒪(\alpha ^2)_{prag}`$ matrix element for the initial-state bremsstrahlung and we shall perform the phase-space integration also in $`𝒪(\alpha ^2)_{prag}`$. We integrate first over final-state fermion four-momenta
$$\begin{array}{cc}& \frac{d^3q_1}{q_1^0}\frac{d^3q_2}{q_2^0}\delta ^{(4)}\left(p_1+p_2q_1q_2k\right)\overline{\beta }_{1I}^{(2)}(q_1+q_2,p_1,p_2,q_1,q_2,k)\hfill \\ & =B_1^{(2)}(p_1,p_2,k)\sigma ^{\mathrm{Born}}((q_1+q_2)^2)\hfill \end{array}$$
(171)
where
$$\begin{array}{cc}\hfill B_1^{(2)}(p_1,p_2,k)=& \frac{\alpha }{4\pi ^2}\frac{2p_1p_2}{(kp_1)(kp_2)}W_e(\widehat{\alpha },\widehat{\beta })(1+\mathrm{\Delta }_I^{(1)}(\widehat{\alpha },\widehat{\beta }))\frac{1}{2}\left\{(1\widehat{\alpha })^2+(1\widehat{\beta })^2\right\}\hfill \\ & \stackrel{~}{S}_I(p_1,p_2,k)(1+\delta _I^{(1)})\hfill \end{array}$$
(172)
and obtain
$$\begin{array}{cc}\hfill \sigma & =\underset{0}{\overset{1}{}}𝑑ve^{\delta _{YFS}+\gamma \mathrm{ln}\epsilon }\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{k_j^0>\epsilon \frac{\sqrt{s}}{2}}{}\underset{j=1}{\overset{n}{}}\frac{d^3k_j}{k_j^0}\stackrel{~}{S}_I(p_1,p_2;k_j)\underset{k_j^0>\epsilon \frac{\sqrt{s}}{2}}{}\frac{d^3k}{k^0}\hfill \\ & B_1^{(2)}(p_1,p_2,k)\sigma ^{\mathrm{Born}}(s(1v))\delta \left(v\frac{1}{s}(P\underset{j}{}k_jk)^2\right)\hfill \\ & \underset{0}{\overset{1}{}}𝑑v\rho _1^{(2)}(v)\sigma ^{\mathrm{Born}}(s(1v))\hfill \end{array}$$
(173)
In the calculation of $`\rho _1^{(2)}`$ we could follow the first of the methods employed for $`\overline{\beta }_0`$. Let us describe it briefly without going into details of the calculation. We calculate the first two non-zero integrals $`𝒪(\alpha )`$ and $`𝒪(\alpha ^2)`$. The first one $`𝒪(\alpha )`$ has to be calculated keeping both the leading $`𝒪(L\alpha )`$ and the subleading term $`𝒪(L^0\alpha )`$. This can be done following the well known $`𝒪(\alpha )`$ analytical calculations . The $`𝒪(\alpha ^2)`$ integral with two real photons can be treated in LL approximation, i.e. keeping only $`𝒪(L^2\alpha ^2)`$ terms. This can be done introducing collinear peaks in the photon angles as demonstrated in the case of $`\overline{\beta }_0`$. Both integrals are connected due to infrared regulation with $`\epsilon `$. The first one is proportional to $`e^{\gamma \mathrm{ln}\epsilon }1+\gamma \mathrm{ln}\epsilon `$ and the term $`\gamma \mathrm{ln}\epsilon `$ from the first one cancels the infrared divergence in the second (independently of the LL approximation). As in the case of $`\overline{\beta }_0`$ one has to pay attention to the subtle “edge effects” in the $`\epsilon `$ regularization<sup>21</sup><sup>21</sup>21 Generally, the calculation for $`\overline{\beta }_1`$ is more difficult than for $`\overline{\beta }_0`$ and $`\overline{\beta }_2`$ because this is the only case in $`𝒪(\alpha ^2)`$ where we deal with simultaneous real and virtual photon emission. .
Let us describe in detail the second method in which soft photons provide the convenient infrared regulation. The main $`𝒪(\alpha )`$ contribution comes from the configuration in which we have $`k^0v\sqrt{s}/2`$ and one or more soft photons. This part has to be calculated exactly in $`𝒪(\alpha )`$. We split, as before,
$$\rho _1^{(2)}(v)=f_1^{(2)}(v)+d_1^{(2)}(v)$$
(174)
in such a way that $`d_1^{(2)}(v)`$ is vanishing in $`𝒪(\alpha )`$ – it can be therefore calculated in second-order LL while $`f_1^{(2)}(v)`$ is simple enough to be calculated exactly in the $`𝒪(\alpha )`$. We define
$$\begin{array}{cc}& f_1^{(2)}(v)=e^{\delta _{YFS}}\frac{d^3k}{k^0}e^{\gamma \mathrm{ln}\epsilon }\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{k_j^0>\epsilon \frac{\sqrt{s}}{2}}{}\underset{j=1}{\overset{n}{}}\frac{d^3k_j}{k_j^0}\stackrel{~}{S}_I(k_j)\hfill \\ & \delta \left(v\frac{2}{s}P(\underset{j}{}k_j+k)\right)B_1^{(1)}(p_1,p_2,k)=e^{\delta _{YFS}}\frac{d^3k}{k^0}f(\gamma ,v\frac{2}{s}Pk)B_1^{(1)}(p_1,p_2,k)\hfill \end{array}$$
(175)
where
$$B_1^{(1)}(p_1,p_2,k)=\frac{\alpha }{4\pi ^2}\frac{2p_1p_2}{(kp_1)(kp_2)}W_e(\widehat{\alpha },\widehat{\beta })\frac{1}{2}\left\{(1\widehat{\alpha })^2+(1\widehat{\beta })^2\right\}\stackrel{~}{S}_I(p_1,p_2,k).$$
(176)
The remarkable feature of $`f_1^{(2)}`$ is that we could integrate over spectator photons exactly. Note that the $`\epsilon `$ regulator has disappeared from the $`k`$integral. In the next step we integrate exactly over photon angles following the old $`𝒪(\alpha )`$ calculations and we are left with the single integral over the photon energy $`x=2k^0/\sqrt{s}`$, with the strongest singularity $`(vx)^{\gamma 1}`$ being regularized nicely by soft photons
$$\begin{array}{cc}\hfill f_1^{(2)}(v)& =e^{\delta _{YFS}}F(\gamma )\underset{0}{\overset{v}{}}𝑑x\gamma (vx)^{\gamma 1}\gamma \left[1+\frac{1}{2}x\right]\hfill \\ & =e^{\delta _{YFS}}F(\gamma )\gamma v^\gamma \left[1+\frac{1}{2}v\frac{1}{2}\gamma v\right]+𝒪(\gamma ^3).\hfill \end{array}$$
(177)
Now we shall calculate the remaining part $`d_1^{(2)}`$ of $`\rho _1^{(2)}`$. Since it vanishes at $`𝒪(\alpha )`$ we may calculate it in LL approximation. Although strictly speaking it is not necessary, we treat photons gently (as in $`\overline{\beta }_0`$ example) such that we do not use the crude collinear approximation. As before, we split the photon angular integration into forward and backward hemispheres and we integrate immediately over the final fermion momenta
$$\begin{array}{cc}& d_1^{(2)}(v)=e^{\gamma \mathrm{ln}\epsilon }\underset{n}{}\underset{n^{}}{}\frac{1}{n!}\frac{1}{n^{}!}\mathrm{\hspace{0.33em}2}\underset{\theta <\pi /2}{}\frac{d^3k}{k^0}e^{\delta _{YFS}}\hfill \\ & \underset{i=1}{\overset{n}{}}\underset{\theta _i>\pi /2}{}\frac{d^3k_i^+}{k_i^{+0}}\stackrel{~}{S}_I(k_i^+)\mathrm{\Theta }\left(k_i^{+0}\epsilon \frac{\sqrt{s}}{2}\right)\underset{j=1}{\overset{n^{}}{}}\underset{\theta _j<\pi /2}{}\frac{d^3k_j^{}}{k_j^0}\stackrel{~}{S}_I(k_j^{})\mathrm{\Theta }\left(k_j^0\epsilon \frac{\sqrt{s}}{2}\right)\hfill \\ & \{[\delta (v1+\frac{(PkK^+K^{})^2}{s})\hfill \\ & \delta (v\frac{2P(k+K^++K^{})}{s})]B_1^{(1)}(p_1,p_2,k)\hfill \\ & +\delta (v1+\frac{1}{s}(PkK^+K^{})^2)[B_1^{(2)}(p_1,p_2,k)B_1^{(1)}(p_1,p_2,k)]\}.\hfill \end{array}$$
(178)
Using the collinear replacement $`K^\pm =(K^{\pm 0},0,0,\pm |K^{\pm 0}|)`$ in $`\delta `$’s allows us to integrate over spectator multiple photons
$$\begin{array}{cc}& d_1^{(2)}(v)=\underset{0}{\overset{1}{}}𝑑v_+\underset{0}{\overset{1}{}}𝑑v_{}\underset{\theta <\pi /2}{}\frac{d^3k}{k^0}e^{\delta _{YFS}}f(\frac{\gamma }{2},v_+)f(\frac{\gamma }{2},v_{})\hfill \\ & \{[\delta (v1+(1xv_+)(1v_{}))\delta (vxv_+v_{})]B_1^{(1)}(p_1,p_2,k)\hfill \\ & +\delta (v1+(1xv_+)(1v_{}))[B_1^{(2)}(p_1,p_2,k)B_1^{(1)}(p_1,p_2,k)]\}.\hfill \end{array}$$
(179)
where $`x=2k^0/\sqrt{s}`$ and the other notation is the same as in $`\overline{\beta }_0`$ case. Integration over photon angles leads to
$$\begin{array}{cc}\hfill d_1^{(2)}(v)& =\underset{0}{\overset{1}{}}𝑑v_+\underset{0}{\overset{1}{}}𝑑v_{}\underset{0}{\overset{1}{}}𝑑xe^{\delta _{YFS}}f(\frac{\gamma }{2},v_+)f(\frac{\gamma }{2},v_{})\hfill \\ & \{[\delta (v1+(1xv_+)(1v_{}))\delta (vxv_+v_{})]\gamma b_1(x)\hfill \\ & +\delta (v1+(1xv_+)(1v_{}))\gamma ^2b_2(x)\},\hfill \\ \hfill b_1(x)& =1+\frac{1}{2}x,b_2(x)=1+\frac{1}{2}x\frac{1}{8}[1+(1x)^2]\frac{\mathrm{ln}(1x)}{x}\hfill \end{array}$$
(180)
Let us show very briefly the calculation of the part proportional to the difference of $`\delta `$’s which is somewhat more tricky. We convolute first $`b_1`$ with photons in the same hemisphere and next with photons from opposite hemisphere
$$\begin{array}{cc}& d_{1A}^{(2)}(v)=𝑑V𝑑v_{}\left[\delta \left(v1+(1V)(1v_{})\right)\delta \left(vVv_{}\right)\right]\hfill \\ & e^{\delta _{YFS}}f(\frac{\gamma }{2},v_{})𝑑x𝑑v_+\delta \left(Vxv_+\right)f(\frac{\gamma }{2},v_+)\gamma b_1(x)\hfill \\ & =e^{\delta _{YFS}}F^2\left(\frac{\gamma }{2}\right)\gamma v^{\gamma 1}\underset{0}{\overset{1}{}}𝑑yy^{\frac{1}{2}\gamma }(1y)^{\frac{1}{2}\gamma 1}\left\{(1vy)^{\frac{1}{2}\gamma }1\right\}\left[1+\frac{vy}{2}\left(1\frac{\gamma }{2}\right)\right]\hfill \\ & =e^{\delta _{YFS}}F(\gamma )\gamma v^\gamma \left(1+\frac{1}{2}v\right)\left(\frac{1}{2}\gamma \mathrm{ln}(1v)\right)+𝒪(\gamma ^3)\hfill \end{array}$$
(181)
The remaining part of $`d_1^{(2)}`$ is easier to calculate because it is explicitly of $`𝒪(\gamma ^2)`$
$$d_{1B}^{(2)}(v)=e^{\delta _{YFS}}F(\gamma )\gamma v^\gamma \gamma \left\{\frac{1}{2}\left(1+\frac{1}{2}v\right)\frac{1}{8}\left(1+(1v)^2\right)\frac{\mathrm{ln}(1v)}{v}\right\}+𝒪(\gamma ^3).$$
(182)
The contribution from the initial-state $`\overline{\beta }_1`$ with $`𝒪(\alpha ^2)_{prag}`$ QED matrix element and with $`𝒪(\alpha ^2)_{prag}`$ analytical integration over the multiphoton phase-space reads
$$\begin{array}{cc}\hfill \rho _1^{(2)}(v)=& e^{\delta _{YFS}}F(\gamma )\gamma v^{\gamma 1}\{\frac{1}{2}(1+\frac{1}{2}v)\hfill \\ & +\gamma [\frac{1}{2}v\frac{1}{4}v^2+\frac{1}{8}(1+3(1v)^2)\mathrm{ln}(1v)]\}+𝒪(\gamma ^3).\hfill \end{array}$$
(183)
The contribution with $`𝒪(\alpha ^1)_{prag}`$ QED matrix element and with analytical $`𝒪(\alpha ^2)_{prag}`$ multiphoton phase-space integration is obtained by retaining only $`d_{1A}^{(2)}`$ and it reads
$$\begin{array}{cc}\hfill \rho _1^{(1)}(v)=& e^{\delta _{YFS}}F(\gamma )\gamma v^{\gamma 1}\{\frac{1}{2}(1+\frac{1}{2}v)\hfill \\ & +\gamma [\frac{1}{2}v^2\frac{1}{2}(1+\frac{1}{2}v)\mathrm{ln}(1v)]\}+𝒪(\gamma ^3).\hfill \end{array}$$
(184)
#### 5.2.3 Beta-bar-two, $`\overline{\beta }_2`$
In the following step our aim is to calculate analytically the contribution to the total cross section from $`\overline{\beta }_{1I}^{(2)}`$ as given by
$$\begin{array}{cc}& \sigma _2=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\frac{d^3q_1}{q_1^0}\frac{d^3q_2}{q_2^0}\underset{j=1}{\overset{n}{}}\frac{d^3k_j}{k_j^0}\stackrel{~}{S}_I(p_1,p_2;k_j)(1\mathrm{\Theta }(\mathrm{\Omega }_I;k_j))e^{Y(\mathrm{\Omega }_I)}\hfill \\ & \underset{nj>k1}{}\frac{\overline{\beta }_{2II}^{(2)}(X,p_1,p_2,q_1,q_2,k_j,k_k)}{\stackrel{~}{S}_I(k_j)\stackrel{~}{S}_I(k_k)}\delta ^{(4)}\left(p_1+p_2q_1q_2\underset{j=1}{\overset{n}{}}k_j\right)\hfill \end{array}$$
(185)
This contribution is in a sense more trivial than the previous two because it is pure $`𝒪(\alpha ^2)`$, it does not have any infrared singularity in the two-photon phase-space integral.
We can calculate the contribution from $`\overline{\beta }_2`$ with the same methods as in the case of $`\overline{\beta }_0`$ or $`\overline{\beta }_1`$. The integral is reorganized easily such that the integration over photon momenta in the $`\overline{\beta }_{2II}^{(2)}`$ is isolated and we are able to integrate over final-state fermion momenta bringing the integral to the standard form
$$\sigma _2=\underset{0}{\overset{1}{}}𝑑v\rho _2^{(2)}(v)\sigma ^{\mathrm{Born}}\left(s(1v)\right).$$
(186)
The function $`\rho _2^{(2)}(v)`$ can be calculated in the LL approximation with either of the two methods (keeping an additional spectator photon or not) and after integration over photon angles the integral boils down to the following integral over longitudinal photon momenta, separately for the case with two collinear and two anticollinear photons
$$\begin{array}{cc}\hfill \rho _2^{(2)}(v)=& \underset{0}{\overset{1}{}}dv_{}dv_+\frac{\gamma ^2}{4}\delta (vv_+v_{})[\frac{1}{2v_+v_{}}\chi (v_+)\chi \left(\frac{v_{}}{1v_+}\right)\hfill \\ & +\frac{1}{2v_+v_{}}\chi (v_{})\chi \left(\frac{v_+}{1v_{}}\right)\frac{1}{v_+}\omega (v_{})\frac{1}{v_{}}\omega (v_+)\frac{1}{v_+v_{}}]\hfill \\ \hfill +& \underset{0}{\overset{1}{}}𝑑v_{}𝑑v_+\frac{\gamma ^2}{4}\delta (v1+(1v_+)(1v_{}))\hfill \\ & \left[\frac{1}{v_+v_{}}\chi (v_+)\chi (v_{})\frac{1}{v_+}\omega (v_{})\frac{1}{v_{}}\omega (v_+)\frac{1}{v_+v_{}}\right]\hfill \\ \hfill =& \gamma ^2\frac{1}{4}v,\hfill \end{array}$$
(187)
where $`\chi (x)=(1+(1x)^2)/2`$ and $`\omega (x)=1+x/2`$.
By eventually keeping additional soft photons in the calculation we get our final result for the initial-state $`𝒪(\alpha ^2)_{prag}`$ contribution from $`\overline{\beta }_2`$ in a more elegant form
$$\rho _2^{(2)}(v)=e^{\delta _{YFS}}F(\gamma )\gamma v^{\gamma 1}\left\{\frac{1}{4}\gamma v^2\right\}+𝒪(\gamma ^3)$$
(188)
We have compared numerically the above formula with the $`𝒦𝒦`$MC in the case of FSR switched off and found an agreement better than 0.1%. In Figure 12(a) we present the comparison in which, as in the case of the previous $`\overline{\beta }`$’s, FSR is switched on. In Figure 12(a) we compare the convolution of the ISR $`\overline{\beta }_{2I}^{(2)}`$ and the FSR $`\overline{\beta }_{0F}^{(2)}`$:
$$\sigma _{\overline{\beta }_2^{(2)}\overline{\beta }_0^{(2)}}^{(2)}=_0^{v_{\mathrm{max}}}𝑑v_0^{v/(1u_{\mathrm{max}})}𝑑u\sigma _{\mathrm{Born}}^f\left(s(1u)(1v)\right)\rho _{I\overline{\beta }_2^{(2)}}^{(2)}(v)\rho _{F\overline{\beta }_0^{(2)}}^{(2)}(u)$$
(189)
The above IEX result is compared with the $`𝒦𝒦`$MC results, and they agree within 0.2%. In Figure 12(b) we show the analogous comparison for the convolution of the FSR $`\overline{\beta }_{2F}^{(2)}`$ and the ISR $`\overline{\beta }_{0I}^{(2)}`$ (anticipating IEX results for FSR $`\overline{\beta }_{2F}^{(2)}`$ in the next section) and we find the similar agreement. Finally, there is another more trivial contribution in the $`\overline{\beta }^{(2)}`$ family which correspond to the case with one real photon emitted in the initial state and one in the final state. This case does not require a separate analytical phase space integration effort because the relevant IEX formula involves the convoluton of the already known expression for the ISR $`\overline{\beta }_{1I}^{(2)}`$ and the FSR $`\overline{\beta }_{1F}^{(2)}`$. The corresponding numerical comparison of the IEX and EEX MC is shown in Figure Figure 12(c). In fact the IEX matrix element was deliberatery constructed in such a way (factorizing virtual corrections) that the above convolution-type IEX formula results.
#### 5.2.4 Summary on IEX for ISR
The entire initial-state $`𝒪(\alpha ^2)_{prag}`$ integrated cross section is obtained by combining contributions from all three $`\overline{\beta }`$’s and it reads
$$\begin{array}{cc}& \sigma _I^{(2)}=\underset{0}{\overset{1}{}}𝑑v\rho _I^{(2)}(v)\sigma ^{\mathrm{Born}}(s(1v)),\hfill \\ & \rho _I^{(2)}(v)=e^{\delta _{YFS}}F(\gamma )\gamma v^{\gamma 1}\{1+\frac{\gamma }{2}+\frac{\gamma }{8}\hfill \\ & +v(1+\frac{1}{2})+\gamma [\frac{v}{2}\frac{1+3(1v)^2}{8}\mathrm{ln}(1v)]\}+𝒪(\gamma ^3)+𝒪(\gamma \alpha )\hfill \end{array}$$
(190)
This above ISR formula has been obtained as a result of ad-hoc exponentiation (interpolation) in ref. and was used there as a numerical parametrization/testing of the cross section from the Monte Carlo program YFS2. It is is now derived starting from YFS exclusive exponentiation by means of direct phase-space integration<sup>22</sup><sup>22</sup>22 Ad-hoc exponentiation is of course easier to do and in ref. even the $`𝒪(\alpha ^3)_{prag}`$ formula for the initial-state bremsstrahlung was given but the derivation method presented here is much better founded and the result does not depend on any kind of interpolation or guesswork.!
Summarizing our IEX calculations for ISR, we have obtained through the analytical integration over the ISR multiphoton phase space the inclusive exponentiated cross section for the IEX matrix elements in the $`𝒪(\alpha ^0)_{prag}`$, $`𝒪(\alpha ^1)_{prag}`$ and $`𝒪(\alpha ^2)_{prag}`$ for each $`\overline{\beta }_i`$ $`i=0,1,2`$ separately. The phase-space integration was always done analytically within the $`𝒪(\alpha ^2)_{prag}`$. All results from the above extensive study are summarized in Table 2 where we list the two functions $`d_S`$ and $`\mathrm{\Delta }_H(v)`$ in the formula
$$\begin{array}{cc}& \sigma _I=\underset{0}{\overset{1}{}}𝑑v\rho _I(v)\sigma ^{\mathrm{Born}}(s(1v)),\rho _I(v)=e^{\delta _{YFS}}F(\gamma )\gamma v^{\gamma 1}\left(d_S+\mathrm{\Delta }_H(v)\right)\hfill \\ & \delta _{YFS}=\frac{\gamma }{4}+\frac{\alpha }{\pi }\left(\frac{1}{2}+\frac{\pi ^2}{3}\right),\gamma =2\frac{\alpha }{\pi }\left(\mathrm{ln}\frac{s}{m_e^2}1\right),F(\gamma )=\frac{e^{C\gamma }}{\mathrm{\Gamma }(1+\gamma )}.\hfill \end{array}$$
(191)
All notation is recalled for the convenience of the reader.
### 5.3 Semi-analytical formulas for FSR
The calculation of the $`𝒪(\alpha ^2)_{prag}`$ IEX formula for the FSR, with the $`u_{\mathrm{max}}`$ cutoff, that is $`u=1s^{}/s<u_{\mathrm{max}}`$ is quite similar to one in the ISR case and we do not enter into details. We only discuss the basic differences between the ISR and FSR cases and present the final result.
If we switch off the ISR completely then the FSR integrated cross section for the $`𝒪(\alpha ^r)_{prag}`$ $`r=0,1,2`$ EEX matrix element reads
$$\begin{array}{cc}& \sigma _F(u_{\mathrm{max}})=\sigma _{\mathrm{Born}}\underset{0}{\overset{u_{\mathrm{max}}}{}}𝑑u\rho _F(u),\rho _F(u)=e^{\delta _{YFS}}F(\gamma _f)\gamma _fu^{\gamma _f1}\left(d_S^{}+\mathrm{\Delta }_H^{}(u)\right)\hfill \\ & \delta _{YFS}^{}=\frac{\gamma _f}{4}\frac{1}{2}\gamma _f\mathrm{ln}(1u)+\frac{\alpha }{\pi }\left(\frac{1}{2}+\frac{\pi ^2}{3}\right),\gamma _f=2\frac{\alpha }{\pi }\left(\mathrm{ln}\frac{s}{m_f^2}1\right),\hfill \end{array}$$
(192)
where the functions $`d_S^{}`$ and $`\mathrm{\Delta }_H^{}(u)`$ obtained with analytical integration of the phase space using the $`𝒪(\alpha ^2)_{prag}`$ approximation, are listed in Table 3.
The main difference and a complication in the phase space analytical integration with respect to the case of ISR is that the YFS formfactor $`\delta _{YFS}^{}`$ depends in the case of FSR on the integration variable $`u`$. This is why terms of $`𝒪(L^2\alpha ^2)`$ are different in the two cases. In Table 3 we show separately the contributions from each $`\overline{\beta }`$. Note that in the case of FSR we did not integrate analytically the phase space for $`\overline{\beta }_0`$ at the $`𝒪(\alpha ^3)_{prag}`$, like in the case of ISR. (It was not necessary in order to reach the precision level of 0.2%.) We have checked numerically the agreement of the $`𝒦𝒦`$MC with the eq. (LABEL:eq:IEX-FSR) separately for each type of the $`\overline{\beta }`$, with the ISR switched off (plots are not shown). We have already presented the complete set of numerical results in the case of the ISR switched off in this section, for each combination of the ISR and FSR $`\overline{\beta }`$’s.
### 5.4 Semi-analytical IEX for ISR and FSR
The last numerical test which we show in Figure 4 is the case in which we switch on all ISR and FSR $`\overline{\beta }`$’s listed in both tables 2 and 3.
$$\sigma _{\mathrm{tot}.}=\underset{0}{\overset{v_{\mathrm{max}}}{}}𝑑v\underset{0}{\overset{v/(1u_{\mathrm{max}})}{}}𝑑u\sigma _{\mathrm{Born}}^f\left(s(1u)(1v)\right)\rho _F(u)\rho _F(v).$$
(193)
It is done for the constant Born cross section, the case with the variable cross section will be shown in the next section. We use the IEX formula of the pure $`𝒪(\alpha ^2)_{prag}`$ type (without $`𝒪(\alpha ^3)_{prag}`$ improvements for ISR). The everall agreements between IEX and $`𝒦𝒦`$MC is within the advertised 0.2%. By looking into all previous figures in this and the previous subsection it is interesting to note that this difference does not come from one particular combination of the ISR and FSR $`\overline{\beta }`$’s, but from several ones.
The reader may wonder why we elaborate so much in this section for the IEX semianalytical formula which are related to the EEX type of the matrix element in $`𝒦𝒦`$MC if in fact the main matrix element in $`𝒦𝒦`$MC is now CEEX. The main reason is that historically the EEX was the first available example of the exclusive exponentiation, and the IEX semianalytical formula were developed in parallel, providing the valuable cross-check of the MC. At this stage, as we see in the next section, both IEX and EEX provide the reference calculation and valuable test for the CEEX. The precision of the present $`𝒪(\alpha ^2)_{prag}`$ IEX is limited, but it could be improved to the full $`𝒪(\alpha ^3)_{prag}`$ if necessary. More important limitation in the present $`𝒪(\alpha ^2)_{prag}`$ IEX as a test of the CEEX model is the lack of the ISR$``$FSR interference. We believe that this effect can be included in the semianalytical IEX if necessary. The ad-hoc variant of the $`𝒪(\alpha ^1)`$ exponentiation including the ISR$``$FSR interference is already available in refs. .
## 6 Numerical results and tests
In this section we shall mainly present the numerical results from the $`𝒦𝒦`$MC in which the Standard Model amplitudes for the process $`e^{}e^+f\overline{f}+n\gamma `$ of the previous Section 2 (EEX) and Section 3 (CEEX) are implemented. The analytical results of the Section 5 will be also exploited to obtain numerical results from the semianalytical program $`𝒦𝒦`$sem. These results are mainly for the $`\mu ^{}\mu ^+`$ final state. For more results on the quark final states and other interesting numerical results from $`𝒦𝒦`$MC we refer the reader to forthcoming proceedings of the LEP2 Monte Carlo Workshop.
The general structure of the $`𝒦𝒦`$MC code is depicted in the Figure 13. The program is divided in the two distinct parts (levels):
* Phase-space Monte-Carlo integration engine with the common importance sampling for the entire family of QED distributions (EEX and CEEX)
* Collection (library) of programs for the SM/QED spin amplitudes and differential distributions, at various orders, with various styles of exponentiation.
In this work we do not enter into the description of the MC integration algorithm in the universal MC integration engine. The Monte Carlo method of the phase space integration is fully documented (for the first time) in ref. and some aspects of the phase-space parametrization are documented in the forthcoming ref. . Here we regard this low-level MC program as a black box capable to integrate the phase space exactly (up to a statistical error).
The life is however not that simple and the numerical program which “in principle” is doing something “exactly/rigorously” may still give imprecise results due to programming bugs and numerical instabilities, especially in a program as complicated as $`𝒦𝒦`$MC is. This is why we always introduce a concept of the technical precision of the given program/calculation, see below. The basic aim of our numerical exercises presented in this Section is the determination of the total theoretical precision associated with our calculation of Standard Model predictions for experimental observables (although we limit the discussion to QED part of SM for most of our discussion). As for observables we shall concentrate mainly on the total cross section and charge asymmetry at LEP1, LEP2 and Linear Collider energies.
What are the technical and physical precisions? The technical precision we define as all uncertainties related to pure numerical problems like programming bugs, numerical instabilities, numerical approximation, etc. In our case the question of the technical precision will mainly concern the MC integration engine. It is important to determine it at the early stage of the work and it should be generally much better than the physical precision. On the other hand, the physical precision is the total uncertainty related to neglected higher orders in coupling constant $`\alpha `$ or in other expansion parameters like the inverse of the big-log $`1/L`$, or the ratio of the width to mass $`\mathrm{\Gamma }/M`$ for a narrow resonance. For physical precision we understand that the above truncations are done in the spin amplitudes and/or differential cross section. If some of them are done in the phase space integration then we tend to associate them with the technical precision (as phase space integration is a technical problem).
We shall start this section with the basic discussion of the technical precision, then we proceed to a subsection elaborating on the physical precision for the EEX matrix element, based on comparisons of the $`𝒦𝒦`$MC and semi-analytical results, and later we discuss the physical precision for the case of the full CEEX matrix element. In this Section we also present numerical results and a rather complete discussion of the effects due to the ISR-FSR interference in the fermion pair production process.
We note that it would be good to include also more numerical tests and lower energies $`10GeV`$, and for very high energies $`1TeV`$, and some more tests specific to spin effects. However, the basic pattern of the spin correlations in double $`\tau `$ decay was already cross-checked in ref. .
### 6.1 Basic test of technical precision
The best way to determine technical precision is to compare results of the two or even more independent calculations which implement the same physics model but differ in technical details of the actual implementation like the method of phase space integration, independent coding, etc. The best two possible methods are: (a) to compare two independent Monte Carlo calculations or (b) to compare Monte Carlo results with the results of a semi-analytical calculation. The method (a) is generally better because it can be done for arbitrary kinematical selections (cuts) and for the simplified QED matrix element, while method (b) is limited to a simple or absent kinematical selections. In the following we shall use method (b).
For our basic test of the technical precision we use the simplest possible variant of the QED model, that is of the type $`𝒪(\alpha ^0)_{\mathrm{EEX}}`$ defined in Section 2. For this type of QED matrix element we were able to integrate analytically the total cross section in the Section 5. The relevant formula can be read from the first row in Tables 2 and 3. For the sake of completeness we write down the complete expression explicitly:
$$\begin{array}{cc}& \sigma _{\mathrm{SAN}}^f=_0^{v_{\mathrm{max}}}𝑑v\sigma _{\mathrm{Born}}^f(s(1u)(1v))\rho _I^{(0)}(v)\rho _F^{(0)}(u),\hfill \\ & \rho _I^{(0)}(v)=F(\gamma _e)e^{\frac{1}{4}\gamma _e+\frac{\alpha }{\pi }\left(\frac{1}{2}+\frac{\pi ^2}{3}\right)}\gamma _ev^{\gamma _e1}\left(1\frac{1}{4}\gamma _e\mathrm{ln}(1v)\frac{1}{2}\frac{\alpha }{\pi }\mathrm{ln}^2(1v)+0\gamma _e^2\right),\hfill \\ & \rho _F^{(0)}(u)=F(\gamma _f)e^{\frac{1}{4}\gamma _f\frac{1}{2}\gamma _f\mathrm{ln}(1u)+\frac{\alpha }{\pi }\left(\frac{1}{2}+\frac{\pi ^2}{3}\right)}\gamma _fu^{\gamma _f1}\left(1\frac{1}{4}\gamma _f\mathrm{ln}(1u)\right),\hfill \end{array}$$
(194)
As we remember the coefficient in front of the $`𝒪(L^3\alpha ^3)`$ term is zero, as marked explicitly. It was essential to calculate analytically and introduce the ISR term of $`𝒪(L^1\alpha ^2)`$ because it amounts numerically to several per cent for the cross section located close to $`v=1`$.
In fig. 14 we present the comparison of the $`𝒦𝒦`$ MC with the semianalytical formula of eq. (194). The difference between the MC result and the semianalytical result is divided by the semianalytical result and as we see it is remarkably small! The comparison is done for the $`\mu ^+\mu ^{}`$ final state at $`\sqrt{s}=189`$ GeV, as a function of $`v_{\mathrm{max}}`$. In the last point (bin) the entire phase space is covered, $`v_{\mathrm{max}}=14m_\mu ^2/s`$.
The conclusion from the above exercise is that we control the phase-space integration at the level of $`2\times 10^4`$ for $`v_{\mathrm{max}}<0.999`$, including the $`Z`$ radiative return, and at the level of $`3\times 10^3`$ for no cuts at all.
The possible loophole in this estimate of precision is that it may break down when we cut the transverse momenta of the real photons, or switch to a more sophisticated QED model. The second is very unlikely as the phase space and the actual SM model matrix element are separated into completely separate modules in the program. The question of the cut transverse momenta of the real photons requires further discussion. Here, it has to be stressed that in our MC the so-called big-logarithm
$$L=\mathrm{ln}\left(\frac{s}{m_f^2}\right)1$$
(195)
is the result of the phase space integration and if this integration were not correct then we would witness the breakdown of the infrared (IR) cancellation and the fermion mass cancellation for FSR. We do not see anything like that at the 0.02% precision level. In addition there is a wealth of comparisons with many independent codes of the phase space integration for $`n_\gamma =1,2,3`$ real photons, with and without cuts on photon $`p_T`$. It should be remembered that the multiphoton phase space integration module/code in $`𝒦𝒦`$MC is unchanged since last 10 years. For ISR it is based on YFS2 algorithm of ref. and for FSR on YFS3 algorithm of ref. , these modules/codes were part of the KORALZ multiphoton MC from the very beginning, already at the time of the LEP1 1989 workshop , and they were continuously tested since then. The phase space integration for $`n_\gamma =1`$ was tested very early by the authors of YFS2/YFS3 against the older MC programs MUSTRAAL and KORALB and with analytical calculations, at the precision level $`<0.1\%`$, with and without cuts on photon $`p_T`$. The phase space integration for $`n_\gamma =2,3`$ with cuts on photon $`p_T`$ was tested very many times over the years by the authors of the YFS2/YFS3/KORALZ and independently by all four LEP collaborations, using other integration programs like COMPHEP, GRACE and other ones, in the context of the search of the anomalous $`2\gamma `$ and $`3\gamma `$ events. Another important series of tests was done in ref. for ISR $`n_\gamma =1,2`$ photons (with cuts sensitive to $`p_T`$ of photons), comparing KORALZ/YFS2 with the other independent MC’s for the $`\nu \overline{\nu }\gamma (\gamma )`$ final states. Typically, these tests, in which QED matrix element was programmed in several independent ways, showed agreement at the level of 10% for the cross section for $`n_\gamma =2`$ which was of order 0.1% of the Born, or 0.2-0.5% for $`n_\gamma =1`$ which was of order 1% of the Born, so they never invalidated our present technical precision of 0.02% in terms of Born cross section (or total cross section in terms of Z-inclusive cut).
We conclude therefore that the technical precision of $`𝒦𝒦`$MC due to phase space integration is 0.02% of the integrated cross section, for any cuts on photon energies Z-inclusive and Z-exclusive, stronger than<sup>23</sup><sup>23</sup>23It downgrades to 0.5% for $`M_{inv}(\mu \overline{\mu })2m_\mu `$, i.e. full phase space. $`M_{inv}(f\overline{f})>0.1\sqrt{s}`$ and any mild cut on the transverse photon energies due to any typical realistic experimental cuts. For the cross sections with a single photon tagged it is about 0.2-0.5% and with two photons tagged it is $`10\%`$ of the corresponding integrated cross section. These conclusions are based on the comparisons with at least six other independent codes.
### 6.2 Physical precision, the case of EEX
We now start the presentation of the numerical results from $`𝒦𝒦`$MC run in the EEX mode with semianalytical calculations based of the results in Section 5. Note that the EEX matrix element of Section 2 is very similar (basically the same) to that implemented since many years in KORALZ program . We do for two reasons: (a) these tests were historically the first, (they existed unpublished for many years giving us confidence that the KORALZ/YFS3 program provides correct results) and (b) they are now still useful as a reference calculation for the newer CEEX scheme. They will also allow us to introduce some notations and gradually introduce the reader to the subject of the discussion of the theoretical precision of our results. Of course, we shall remember that in the case of EEX we do not include the ISR-FSR interferences (IFI).
In Figure 15 we show the dependence of the total cross section on the cut on the total photon energy (ISR+FSR). The comparison is done for the $`\mu ^+\mu ^{}`$ final state at $`\sqrt{s}=189`$ GeV, as a function of $`v_{\mathrm{max}}`$. In the last point (bin) the entire phase space is covered, i.e. $`v_{\mathrm{max}}=14m_\mu ^2/s`$. The very striking (and well known) phenomenon is that the total cross section due to huge ISR correction is almost 3 times the Born cross section, in the absence of any kinematical cuts. Part of this ISR contribution is located close to $`v=1`$, $`s^{}4m_\mu ^2/s`$, let us call it the $`\gamma \gamma ^{}`$ process, it amounts to as much as the Born cross section itself, $`\sigma _{\gamma \gamma ^{}}\sigma _{\mathrm{Born}}`$, while dominant part of the cross section $`\sigma _{\mathrm{ZRR}}2\sigma _{\mathrm{Born}}`$ is concentrated close to $`v=1M_Z^2/s0.75`$, and is associated with the so called “$`Z`$ radiative return” (ZRR) process, that is the resonant production of Z, after emission of rather hard ISR photon, usually lost in the beam pipe. In the experiment the $`\gamma \gamma ^{}`$ process is almost always eliminated from the data, and the ZRR process is also not very often included in the data sample. The typical experimental cut is situated somewhere in the range $`0.1<v_{\mathrm{max}}<0.3`$. As we see in Figure 15 (a), the total QED corrections $`(\sigma (v_{\mathrm{max}})\sigma _{\mathrm{Born}})/\sigma _{\mathrm{Born}}`$ is in this case quite close to zero, in fact slightly negative.
In Figure 15(b) we compare the $`𝒦𝒦`$MC with the semianalytical expression based on the phase-space integration in Section 5. In the MC calculation we use the second order EEX type of the QED model EEX2$``$$`𝒪(1,\alpha ,\alpha L,\alpha ^2L^2)_{\mathrm{EEX}}`$, defined in Section 2. The semianalytical formula used in Figure 15(b) is also in the class EEX2. It is defined as follows
$$\sigma _{\mathrm{SAN}}^f=_0^{v_{\mathrm{max}}}𝑑v\sigma _{\mathrm{Born}}^f(s(1u)(1v))\rho _I^{(2)}(v)\rho _F^{(2)}(u),$$
(196)
where the distributions $`\rho _I^{(2)}`$ and $`\rho _F^{(2)}`$ are from the Tables 2 and 3 in Section 5.
What kind of lesson can we draw from Figure 15(b)? We treat the result in Figure 15(b) as an indication that, the contribution from QED (non-IFI) photonic corrections to combined physical and technical precision in the EEX2-class integrated cross section for the standard cut $`v_{\mathrm{max}}0.2`$ is about 0.2%, for the ZRR process it is 0.7% and for the $`\gamma \gamma ^{}`$ process it is 3%. We are here talking about the technical precision of the coding of the EEX2 matrix element, not associated with the phase space integration (covered in the previous section).
In the next Figure 16 we make an attempt to estimate the physical precision of the QED model in the EEX class. Specifically, we look into difference between EEX2 (as defined above) and EEX1, with the EEX1 being the $`𝒪(\alpha ^1)_{\mathrm{EEX}}`$ of Section 2, EEX1$``$$`𝒪(1,\alpha ,\alpha L)_{\mathrm{EEX}}`$. It is plotted in Figure 16(a) both from $`𝒦𝒦`$MC and semianalytical formula. Taking conservatively (see the discussion below) half of the difference between EEX2 and EEX1 as an estimate of the physical precision of EEX2 we arrive to a similar estimate of about 0.2% for the standard cut $`v_{\mathrm{max}}0.2`$, 0.7% for the ZRR process and up to 3% for the $`\gamma \gamma ^{}`$ process.
The other useful piece of information comes from Figure 16(b), where we plot the difference EEX3$``$EEX2, with EEX3$``$$`𝒪(1,\alpha ,\alpha L,\alpha ^2L^2,\alpha ^3L^3)_{\mathrm{EEX}}`$, provides direct insight into the neglected third order LL contributions. As we see it is always below $`310^4`$, (This estimate will be also useful for the case of CEEX.) If the $`𝒪(L^3\alpha ^3)`$ corrections is of this size, then necessarily the main contribution to the above estimate of the theoretical error is coming from the $`𝒪(L^1\alpha ^2)`$ corrections!
In fact the lack of the $`𝒪(\alpha ^2L^1)`$ corrections in both EEX2 and EEX1 is the main deficiency of the above tests, so they cannot pin down directly the size of this contribution. Keeping this limitation in mind, from the above test we nevertheless estimate tentatively the combined physical and technical precision in the integrated EEX3-class cross section of the $`𝒦𝒦`$MC to be 0.2% for the standard cut $`v_{\mathrm{max}}0.2`$, 0.7% for the ZRR process and about to 1.5% for the $`\gamma \gamma ^{}`$ process. The caveat of this exercise is that we know retrospectively the QED non-IFI component of the precision on the KORALZ/YFS3 Monte Carlo at LEP2 energies because the EEX of KORALZ and the EEX of $`𝒦𝒦`$MC are practically the same<sup>24</sup><sup>24</sup>24KORALZ version 4.02 and earlier have EEX implemented differently from $`𝒦𝒦`$MC.. The above does help indeed, in spite of the fact that the neglected IFI contribution to integrated cross section is of order 1%, because KORALZ in the non-exponentiated $`𝒪(\alpha )`$ mode can calculate IFI separately, see discussion in the following subsections.
Let us finally make an ultimate effort to estimate the total precision, staying all the time within the EEX model. As we have already noted the most important missing contribution seems to be the $`𝒪(L^1\alpha ^2)`$, most probably the ISR part of it. In the the semianalytical formula for the total cross section we are able to add it, since it is known from ref. . The $`𝒪(L^3\alpha ^3)`$ corrections we may add as well and in this way we replace $`\rho _I^{(2)}`$ by the $`\rho _I^{(3)}`$ of ref. which is the true $`𝒪(\alpha ^3)_{prag}`$ for ISR (according to the terminology in the Introduction) and $`𝒪(\alpha ^2)_{prag}`$ for FSR (no IFI). Let us call it EEX3best$``$$`𝒪(1,\alpha ,\alpha L,\alpha ^2L^2,\alpha ^2L^1,\alpha ^3L^3)_{\mathrm{EEX}}`$. The difference between semianalytical EEX3best and EEX3 from $`𝒦𝒦`$MC is plotted in Figure 17. As we see this final test confirms the previous estimate of the physical precision of the EEX type of the matrix element.
### 6.3 Physical precision, the case of CEEX
Quantitative determination of the physical precision should be based on the comparison of the calculations in two consecutive orders in the expansion parameters, for instance by comparing results from $`𝒪(\alpha ^r)`$ and $`𝒪(\alpha ^{r1})`$ calculation, or $`𝒪(L^r\alpha ^n)`$ versus $`𝒪(L^{r1}\alpha ^n)`$, etc. For example when only the Born and $`𝒪(\alpha ^1)`$ results are available one should take the difference between the two (or some fraction of it) as an estimate of the physical precision. The above conservative recipe gives solid estimate of the physical precision and we shall employ it as our basic method in the following. In most cases in the literature, however, authors try to estimate the uncalculated higher order effects with some “rule of thumb”. For instance in the case when Born and $`𝒪(\alpha ^1)`$ results are known they take $`\frac{1}{2}L\frac{\alpha }{\pi }`$ as an estimate of the missing/uncalculated $`𝒪(\alpha ^2)`$ corrections. This has to be done with care because one may easily overlook some “enhancement factor”. For example the cross section close to a resonance can be modified by additional powers of the big logarithm $`\mathrm{ln}\frac{\mathrm{\Gamma }}{M}`$. In most cases these “enhancement factors” are already seen in the $`𝒪(\alpha ^1)`$ calculation so it is not difficult to trace them.
We are in rather comfortable situation because for QED “photonic” corrections we have at our disposal $`𝒪(\alpha ^0)`$, $`𝒪(\alpha ^1)`$ and $`𝒪(\alpha ^2)`$ calculations (at least for ISR, where they are the biggest). We can therefore afford to take half of the difference between the $`𝒪(\alpha ^1)`$ and $`𝒪(\alpha ^2)`$ calculations as a conservative estimate of the physical precision due to QED “photonic” corrections. We also profit from the fact that the exponentiation speeds up considerably the convergence of the perturbative series by “advanced summation” of certain class of corrections to infinite order, and by not introducing additional spurious cut-off parameters dividing real emissions into soft and hard ones which are typical for the calculations without exponentiation (see discussion on the famous $`k_0`$ parameter in the 1989 LEP workshop ).
Let us mention that in our estimates of the physical precision we omit from the discussion the $`𝒪(\alpha ^2)`$ effects due to an additional fermion pair, either real or virtual. We do it because: (a) there are many MC programs which implement production of the four fermion final states (often with additional ISR) and (b) in the experiment this contribution can be eliminated at the early stage from the data in the experimental data analysis aimed at single fermion pair production, see for example ref. . In fact this point is still under debate, see forthcoming proceedings of the LEP2 Monte Carlo workshop . It was proposed that in the final combined LEP2 data certain the so called non-singlet initial and final state secondary pair contribution will be kept in the data, as done by OPAL, see refs. . We have recently included the virtual corrections of the “vacuum polarization” type with the fermionic bubble in the $`𝒪(\alpha ^2)`$ photonic contributions like vertex corrections in yet unpublished version of 4.14 of $`𝒦𝒦`$MC. This is done having in mind combining results of $`𝒦𝒦`$MC with the other MC program for four-fermion production process like KORALW . The tandem of $`𝒦𝒦`$MC and KORALW programs will be able to realize any possible scenario of the treatment of the soft/light pair corrections in the LEP2 data.
In Fig. 18 we present the numerical results on which we base our quantitative estimate of the physical precision due to photonic QED corrections. In this figure we plot the difference between $`𝒪(\alpha ^2)_{\mathrm{CEEX}}`$ and $`𝒪(\alpha ^1)_{\mathrm{CEEX}}`$ for the total cross section and charge asymmetry at 189GeV as a function of the cut on the total energy emitted by all ISR and FSR photons for $`\mu ^+\mu ^{}`$ final state. The cut is formulated with the $`s^{}>s_{\mathrm{min}}^{}`$ or equivalently $`v<v_{\mathrm{max}}`$ condition, where $`s^{}`$ is the effective mass squared of the $`\mu ^+\mu ^{}`$ pair and $`v=1s^{}/s`$, as usual. One should remember that the actual experimental cut is around $`v_{\mathrm{max}}0.2`$ (eliminating $`Z`$ radiative return) in the case of the standard data analysis, and sometimes around $`v_{\mathrm{max}}0.9`$ in the case when $`Z`$ radiative return is admitted in the data. The “kink” around $`v_{\mathrm{max}}0.75`$ is at the position of the $`Z`$ radiative return. In either case, whether we admit or eliminate the $`Z`$ radiative return, that is for $`v_{\mathrm{max}}0.9`$, the difference between $`𝒪(\alpha ^2)_{\mathrm{CEEX}}`$ and $`𝒪(\alpha ^1)_{\mathrm{CEEX}}`$ in total cross section is below 0.4% and for the charge asymmetry it is below 0.002.
Taking conservatively half of this difference among $`𝒪(\alpha ^2)_{\mathrm{CEEX}}`$ and $`𝒪(\alpha ^1)_{\mathrm{CEEX}}`$ as an estimate of the neglected $`𝒪(\alpha ^3)_{\mathrm{CEEX}}`$ and higher orders we conclude that the physical precision due to photonic QED corrections of our $`𝒪(\alpha ^2)_{\mathrm{CEEX}}`$ calculation for all possible cutoffs within $`0<v_{\mathrm{max}}<0.9`$ range is 0.2% in the total cross reaction and 0.001 in the charge asymmetry. This estimate would be even a factor of two better, if we restricted ourselves to the most typical cut-off range $`0.1<v_{\mathrm{max}}<0.3`$. The above estimate will be confirmed by more auxiliary tests in the following.
As we see we have improved on the physical precision estimate with respect to the previous estimates for the EEX model – in addition we do include IFI all the time. For the ZRR process we now quote for the integrated cross-section 0.2% instead of the previous 0.7% and for $`\gamma \gamma ^{}`$ we have something like 0.3% instead of the previous 1.5%. This we interpret as a result of inclusion of the $`𝒪(L\alpha ^2)`$ ISR correction in our CEEX spin amplitudes.
We have to stress very strongly that the estimate of the physical precision depends on the type of the observable (we took $`\sigma `$ and $`A_{\mathrm{FB}}`$), the type of the final state (we took $`\mu `$ pair final state; for the quark final state the QED FSR effects are smaller due to the smaller electric charges of quarks) and on many other input parameters, for example the total CMS energy. The great thing about the Monte Carlo is that the type of the evaluation we proposed and implemented in this Section (half of difference $`𝒪(\alpha ^2)`$$``$$`𝒪(\alpha ^1)`$) can be repeated for any observable, any final state and any energy. For example in Figure 19 we repeat our evaluation of the physical precision for $`\sigma `$ and $`A_{\mathrm{FB}}`$ at the Linear Collider energy 500GeV. As we see the resulting precision is worse, the worsening is negligible for a mild cut of order $`v_{\mathrm{max}}<0.5`$ and almost factor two for the $`Z`$ radiative return, which is now placed close the $`v=0.95`$.
### 6.4 Absolute predictions, more on physical/technical precision
In this Section we shall present the SM absolute predictions for the total cross section and charge asymmetry at LEP2 (189GeV) and at the Linear Collider (500GeV). We compare them with our own semianalytical program $`𝒦𝒦`$sem, with KORALZ and in some cases with ZFITTER . They may not improve our basic estimates of the technical and physical precision from the previous sections, but they can confirm them (or disprove them!).
In Table 5 we show numerical results for the total cross-section $`\sigma (v_{\mathrm{max}})`$ and charge asymmetry $`A_{\mathrm{FB}}(v_{\mathrm{max}})`$ as a function of the cut $`v_{\mathrm{max}}`$ on the total photon energy (the cut-off parameter $`v_{\mathrm{max}}`$ is defined as in the previous subsection). Generally, in Table 5 we show results with the ISR-FSR interference (IFI) switched on and off. The $`𝒦𝒦`$sem semianalytical program (part of the $`𝒦𝒦`$MC package) provides “reference results” for $`\sigma `$ and $`A_{\mathrm{FB}}`$, see the first column in Table 5, which are without IFI, obtained from using EEX3best formula defined in the previous Section 6.2. For the charge asymmetry we use the convolution-type semianalytical formula like that of eq. (196). (In fact we use this formula separately for the cross section in the forward and backward hemispheres and then we calculate $`A_{\mathrm{FB}}`$ from these partial integrals.) Results from the $`𝒦𝒦`$MC in Table 5 are shown for two types of QED matrix element: $`𝒪(\alpha ^2)_{\mathrm{CEEX}}`$ with and without IFI. In addition we include results from KORALZ are for the $`𝒪(\alpha ^1)`$ matrix element with and without IFI which will be discussed in the next Section.
As tables with list of numbers are difficult to comprehend, we present the essential results of the Table 5 in Figure 20, where they are all plotted as a difference with the reference results of our semianalytical program $`𝒦𝒦`$sem. (In other words the results from $`𝒦𝒦`$sem are exactly on the $`x`$-axis.)
In the case of IFI switched on $`𝒦𝒦`$sem cannot be used as a cross-check of the $`𝒦𝒦`$MC. Remembering that IFI in KORALZ in the $`𝒪(\alpha ^1)`$ mode (without exponentiation) is very well tested we combine the $`𝒪(\alpha ^1)`$ IFI contribution with the CEEX result without IFI. Such a hybrid solution denoted in Fig. 20 as “CEEX2+IFI at $`𝒪(\alpha ^1)`$” us used as our primary test of the full CEEX matrix element with IFI switched on. The above procedure is done separately for cross sections in the forward and backward hemispheres such that the prediction for charge asymetry is also available.
It is worth to mention that the above hybrid solution was already successfully used in refs. for the study of the IFI contribution at Z peak, imposing strong acollinearity cut. It is also implemented in a semianalytical form in ZFITTER 6.x. On general ground we expect this recipe to be rather good, because IFI correction itself at $`𝒪(\alpha ^1)`$ does not contain any large mass logarithm and is relatively small and can be handled additively.
In Figure 20 we also show the numerical results from $`𝒦𝒦`$MC in the EEX3 mode (no IFI), from KORALZ in the EEX2 mode (no IFI), which are not included in Table 5.
Let us now comment on the results in Figure 20. The EEX3 from $`𝒦𝒦`$MC differs from EEX3best of $`𝒦𝒦`$sem (no IFI in both) by about 0.7% for the ZRR process, as we have already seen, and we interpret this difference as the result of the missing $`𝒪(L^1\alpha ^2)`$. The EEX2 of KORALZ 4.03 is closer to the EEX3best of $`𝒦𝒦`$sem for ZRR process – we do not see any contradiction in this since the implementation of EEX in KORALZ and $`𝒦𝒦`$MC differs in the details (causing a difference of $`𝒪(L^1\alpha ^2)`$ in the integrated cross section.)
In the case of IFI switched off, the CEEX2 result, corresponding exactly to $`𝒪(\alpha ^2)_{\mathrm{CEEX}}`$, defined in Section 3, as implemented in $`𝒦𝒦`$MC 4.13, agrees very well with the EEX3best of $`𝒦𝒦`$sem. This result is compatible with the total theoretical precision of 0.2% for the integrated cross-section, even including the ZRR process.
In the case of IFI switched on, the hybrid solution “CEEX2+IFI at $`𝒪(\alpha ^1)`$” also agrees with the full CEEX2 result confirming the total theoretical precision of 0.2% for the integrated cross-section, including the ZRR process.
For the charge asymmetry in Figure 20 situation is quite similar. The IFI effect is up to 4% for strong cuts. In the case of IFI switched off, the CEEX2 result agrees with the EEX3best of $`𝒦𝒦`$sem to within 0.2%. For IFI included, the CEEX2 agrees with the hybrid solution rather well, to within 0.4%. Note that in the above Monte Carlo exercise we have used the symmetric definition of the scattering angle $`\theta ^{}`$ of ref. (which is close to what is used in the LEP experiments).
Summarising, the numerical results in Figure 20 establish our basic estimate of the theoretical precision of the $`𝒦𝒦`$MC, due to QED effects, at LEP2 energies of about 0.2% for total cross section and 0.2-0.4% (depending on cut-offs) for charge asymmetry.
Finally, we examine the analogous results from the $`𝒦𝒦`$MC at 500GeV in Figure 21. In this case we include only results from the $`𝒦𝒦`$MC and $`𝒦𝒦`$sem. The pattern of agreement is up to a factor two the same as at 189GeV.
### 6.5 Initial-final state interference
The control of the initial-final state interference correction down to the precision of 0.2% in the integrated cross-section and in the charge asymmetry is rather important – this is why we dedicate this section to a more detailed study of this QED correction. In particular we would like to answer the following questions:
* How big is the ISR$``$FSR interference in $`\sigma _{tot}`$, $`A_{FB}`$?
* Do we know ISR$``$FSR at $`𝒪(\alpha ^1)`$?
* Do we know ISR$``$FSR beyond $`𝒪(\alpha ^1)`$?
* How sensitive is ISR$``$FSR to cut-off changes?
KORALZ is the best starting point and reference for the problem of calculating the ISR$``$FSR. In Figure 22 we show results from the $`𝒪(\alpha ^1)`$ KORALZ (no exponentiation) for the $`\mu ^+\mu ^{}`$ final state at $`\sqrt{s}`$=189GeV. Angular distributions from KORALZ, pure $`𝒪(\alpha ^1)`$ (without exponentiation), were verified very precisely at the level of $`0.01\%`$ using a special analytical calculation, see ref. , so we know ISR$``$FSR at $`𝒪(\alpha ^1)`$ very precisely. As we see the ISR$``$FSR contribution to the integrated cross-section is about 3% and about 0.03 to $`A_{\mathrm{FB}}`$. This is definitely above the ultimate experimental error tag for the combined LEP2 data at the end of LEP2 operation. The energy cut on the total photon energy is fixed in the results of Figure 22 to just one value $`v<v_{\mathrm{max}}=0.1`$ (where $`v_{\mathrm{max}}=1s^{}/s`$ is defined as usual). This is close to the usual value in the experimental LEP2 data analysis. We introduce also the angular cut $`|\mathrm{cos}\theta |<\mathrm{cos}\theta _{\mathrm{max}}`$ and vary the value of $`\mathrm{cos}\theta _{\mathrm{max}}`$, see Figure 22(b), the value used in the experimental LEP2 data analysis is around $`\mathrm{cos}\theta _{\mathrm{max}}=0.9`$; this corresponds to two bins before the last one in Figure 22(b) (the last point in the plot is for $`\mathrm{cos}\theta _{\mathrm{max}}=1`$). In this way already we have answered the first two questions from the list above.
In Figure 23 we present similar results from the $`𝒦𝒦`$MC which will help us to answer whether we know ISR$``$FSR beyond $`𝒪(\alpha ^1)`$ and inspect in a more detail the dependence on cut-offs. In Figure 23(a-b) we essentially repeat the exercise of Figure 22 finding out the ISR$``$FSR contribution to the angular distribution and $`A_{\mathrm{FB}}`$ for the same energy-cut using $`𝒦𝒦`$MC instead of KORALZ. As we see the results change slightly, the ISR$``$FSR effect is about 20%-30% smaller. We attribute it mainly to (a) different (better) treatment of the ISR in $`𝒦𝒦`$MC (b) exponentiation of the ISR$``$FSR effect is in $`𝒦𝒦`$MC. As it is well known, in $`𝒪(\alpha ^1)`$, the ISR$``$FSR contributes like $`4Q_eQ_f\frac{\alpha }{\pi }\mathrm{ln}\frac{1\mathrm{cos}\theta }{1+\mathrm{cos}\theta }`$ to the angular distribution – this even causes the angular distribution to be negative close to $`\mathrm{cos}\theta =1`$. In the CEEX exponentiation the above singularity is summed up to infinite order and the angular distribution near $`|\mathrm{cos}\theta |=1`$ is not singular any more. (This kind of exponentiation will be implemented in the next version of ZFITTER, see for first numerical results.) The typical experimental cut $`|\mathrm{cos}\theta |<0.9`$ eliminates most of the above trouble anyway – what is probably more important is the correct “convolution” of the IR-finite $`𝒪(\alpha ^1)`$ ISR$``$FSR with the $`𝒪(\alpha ^2)`$ ISR. In the $`𝒦𝒦`$MC this is done in a maximally clean way from the theoretical/physical point of view (at the amplitude level) while in the semianalytical programs like ZFITTER this is done in a more “ad hoc” manner. Let us remind the reader that we still lack the genuine IR-finite $`𝒪(\alpha ^2)`$ corrections in the ISR$``$FSR class from diagrams like 2-boxes and 5-boxes, see Section 3. These contributions are most likely negligible, of order $`𝒪(L^1\alpha ^2)`$ at most.
In Figure 23(c-d) we make the energy cut more loose, $`v_{\mathrm{max}}=0.9`$, thus admitting the ZRR into the available phase-space. As a result the relative ISR$``$FSR decreases by a factor 3, simply because it gets “diluted” in the factor 3 bigger integrated cross section, while ZRR does not contribute to ISR$``$FSR because of its narrow-resonance character, already discussed in Section 3 at length. The fact that the ZRR does not contribute to the ISR$``$FSR can be seen explicitly in Figure 24 where we plot the ISR$``$FSR contribution to $`A_{\mathrm{FB}}`$ “bin-per-bin”, that is calculated in each bin separately. As we see the contribution from the ZRR which at this energy (189GeV) is located at $`v=0.75`$ is very small, smaller than from all other $`v`$’s where there is no Z resonance.
In the above exercises and also in the following we use always the energy cut on the $`v=1s^{}/s`$ variable defined in terms of the effective mass of the “bare” final fermions, that is without any attempt of combining them with the collinear FSR photons. This is experimentally well justified for the $`\mu `$-pair final states but not for $`\tau `$-pairs or quarks. It is possible and in fact rather easy to define a “propagator” or “reduced” $`s_p^{}`$ which takes into account the loss of energy due to ISR but not FSR. In other words the $`s_p^{}`$ effective mass squared sums up FSR photons. One can ask a legitimate question: If we would cut not on the “bare” final fermion variable $`s^{}`$, but instead on the “propagator” $`s_p^{}`$ then perhaps the estimate of the ISR$``$FSR contribution would be dramatically different, for instance it would be much smaller? In Figure 26 we show a numerical exercise in which we employ the energy cut in terms of $`v_p=1s_p^{}/s`$. One can construct such a $`s_p^{}`$ looking into angles of the outgoing fermions. This type of variable was used in ref . In Figure 26 we use the definition of $`s_p^{}`$ of ALEPH . As we see in this Figure, the result is not dramatically different from what we have seen in Figure 23. The magnitude of ISR$``$FSR contribution is close to what we could see if we applied the same value of the energy cut for the “bare” $`s^{}`$ (as we have checked independently).
We shall now examine the dependence of the ISR$``$FSR contribution on the energy-cut $`v_{\mathrm{max}}`$ in a more detail. In Figure 26 we show the ISR$``$FSR contribution to $`A_{\mathrm{FB}}`$ as a function of energy-cut $`v_{\mathrm{max}}`$ at two energies (a) 189GeV and (b) at the $`Z`$ peak $`\sqrt{s}=M_Z`$. No cut on $`\mathrm{cos}\theta `$ is applied. In addition to $`𝒦𝒦`$MC results we show the results from $`𝒪(\alpha ^1)`$ mode of KORALZ and from ZFITTER. At 189GeV and for the typical energy cut $`0.2<v_{\mathrm{max}}<0.3`$ all three programs agree very well. This cut is relatively “inclusive” such that exponentiation effects are not so important and ISR is eliminated in a “gentle” way (the total cross section is close to the Born value). For stronger cuts $`v_{\mathrm{max}}<0.2`$ we see a large (factor 2) discrepancy among the results from $`𝒦𝒦`$MC and both KORALZ and ZFITTER, because of the lack of exponentiation in KORALZ and ZFITTER. (in ZFITTER ISR$``$FSR is taken without exponentiation and combined with ISR “additively”). We also observe the discrepancy of about 0.2% among $`𝒦𝒦`$MC on one hand and both KORALZ and ZFITTER on the other hand for the ZRR. Our guess is that it is due to the difference in the method of combining ISR$``$FSR with the second order ISR (of course, we believe that the CEEX method of doing it at the amplitude level is the best one can do). In Figure 26(b) we see, first of all, the well known phenomenon of the strong suppression of the ISR$``$FSR contribution at the resonance, especially for the loose cut-off. Even for a strong cut, $`v_{\mathrm{max}}=0.1`$, the ISR$``$FSR contribution is about 0.01, about factor 30 smaller than in the off-resonance case. Here, $`𝒦𝒦`$MC agrees rather well with KORALZ and ZFITTER. The differences are generally<sup>25</sup><sup>25</sup>25 The difference between KORALZ and ZFITTER should be perhaps smaller because both are $`𝒪(\alpha ^1)`$? May be the difference is due to the angle definition? up to 0.0015.
In Figure 27(a) we examine the ISR$``$FSR contribution to integrated cross section as a function of the energy cut $`v_{\mathrm{max}}`$. At 189GeV and for the typical energy cut $`0.1<v_{\mathrm{max}}<0.6`$ all three programs agree reasonably well, KORALZ and ZFITTER are generally closer to each other than to $`𝒦𝒦`$MC. After admitting the ZRR, $`v_{\mathrm{max}}>0.8`$ all three programs agree even better. For a very strong cut, $`v_{\mathrm{max}}<0.1`$ KORALZ and ZFITTER disagree dramatically with $`𝒦𝒦`$MC because of the lack of exponentiation in KORALZ and ZFITTER. In Figure 27(b) we see, again the strong suppression of the ISR$``$FSR contribution at the resonance, especially for the loose cut-off. Suppression is cut-off dependent and generally stronger for KORALZ and ZFITTER than for $`𝒦𝒦`$MC. Most of the comments which we made on ISR$``$FSR contribution to $`A_{\mathrm{FB}}`$ apply also here.
Finally in Figure 28 we go back to a vicinity of the $`Z`$ peak (LEP1) and we show the magnitude of the ISR$``$FSR contribution to the integrated cross section as a function of the CMS energy, for the $`\mu ^{}\mu ^+`$ final state and for all five quark final states taken together (the so-called hadronic cross section) from the $`𝒦𝒦`$Monte Carlo. No angular cut or energy cut is applied (full phase space). For the $`\mu ^{}\mu ^+`$ final state we also include results from the KORALZ $`𝒪(\alpha ^1)`$ and ZFITTER/TOPAZ0 . Results from quarks are multiplied by a factor 10 to be visible, because ISR$``$FSR contribution in this case is small. It is not only suppressed by the smallness of the quark charge, but we also have partial cancellation among up- and down-type quarks, see ref . However, the ISR$``$FSR contribution to hadronic cross section has to be known much more precisely (factor $`3`$) because it is measured much more precisely, due to higher statistics. In Fig. 28 we see that the suppression of ISR$``$FSR is much weaker as we go away from the centre of the resonance, and it changes the resonance curve in such a way that it affects the fitted mass of the $`Z`$. The actual size of the shift of $`M_Z`$ was studied in ref. and it was found to be 0.15MeV. Results of the $`𝒦𝒦`$MC are smaller about 10-20% than the $`𝒪(\alpha ^1)`$ estimates of KORALZ $`𝒪(\alpha ^1)`$ and ZFITTER, away from the $`Z`$ peak. This is compatible with the 10-20% size of the $`𝒪(L^2\alpha ^2)`$ ISR corrections with respect to $`𝒪(L^1\alpha ^1)`$ corrections, which are included in $`𝒦𝒦`$MC and are not included in KORALZ and apparently also not included in ZFITTER/TOPAZ0 (which agree very well with KORALZ). Our last comment concerns the reliability of our estimate for the ISR$``$FSR contribution in the absence of the correct implementation of the simultaneous emission of the FSR photon and FSR gluon. We think that through the usual arguments, see ref. , we can neglect from consideration emission of the FSR single gluon, as long as we stick to very a inclusive cross section, like the total cross section in Figure 28. For stronger angular cuts, or events with definite jet multiplicity, we would need to improve our calculation.
We summarize the results of this section on ISR$``$FSR as follows:
* For a typical expt. energy cut of 0.3 ISR$``$FSR int. is about 1.5% in $`\sigma _{tot}`$ and $`A_{FB}`$.
* For the energy cut 0.1 it is a factor 2 bigger.
* The cut $`|\mathrm{cos}\theta |<0.9`$ makes it 25% smaller.
* The $`𝒪(\alpha ^1)`$ ISR$``$FSR int. is under total control using KORALZ and $`𝒦𝒦`$ Monte Carlo for arbitrary cuts.
* Effects beyond $`𝒪(\alpha ^1)`$ are negligible, ($`<`$20% of $`𝒪(\alpha ^1)`$), except when the energy cut is stronger than 0.1.
* ISR$``$FSR int. at the Z radiative return is very small, as expected.
* Changing from $`s^{}`$ to $`Q^2`$-propagator in the energy cut has no effect.
### 6.6 Total theoretical precision
Let us summarize the total theoretical precision:
* For the most typical cut-off range $`0.1<v_{\mathrm{max}}<0.3`$ excluding the Z radiative return we quote for CEEX the total precision 0.2% for LEP2 and for the LC at 0.5TeV.
* For a cut-off including ZRR we quote 0.2% total precision for LEP2 and 0.4% total precision for the LC at 0.5TeV.
* For $`\gamma \gamma ^{}`$ we quote 0.3% at LEP2 (no firm result for LC).
In the above estimates the technical component was significantly below the physical one. Restrictions apply: No light fermion pairs (pure photonic QED), no EW component.
## 7 Outlook and summary
The most important new features in the present CEEX are the ISR-FSR interference, the second-order subleading corrections, and the exact matrix element for two hard photons. This makes CEEX already a unique source of SM predictions for the LEP2 physics program and for the LC physics program. Note that for these the electroweak correction library has to be reexamined at LC energies. The most important omission in the present version is the lack of neutrino and electron channels. Let us stress that the present program is an excellent starting platform for the construction of the second-order Bhabha MC generator based on CEEX exponentiation. We hope to be able to include the Bhabha and neutrino channels soon, possibly in the next version. The other important directions for the development are the inclusion of the exact matrix element for three hard photons, together with virtual corrections up to $`𝒪(\alpha ^3L^3)`$ and the emission of the light fermion pairs. The inclusion of the $`W^+W^{}`$ and $`t\overline{t}`$ final states is still in a farther perspective.
## Acknowledgements
Two of us (SJ and BFLW) would like to thank the CERN EP and TH Divisions. We are grateful to all four LEP Collaborations and their members for support. In particular we would like to thank Dr. D. Schlatter of ALEPH for continuous support and help. One of us (S.J.) would like to thank the DESY Directorate for its generous support in the critical stage of the beginning of this project. We would like to express our gratitude to W. Płaczek, E. Richter-Wa̧s, M. Skrzypek and S. Yost for valuable comments.
## 8 Appendix A: Basic KS/GPS spinors and photon polarizations
The arbitrary massless spinor $`u_\lambda (p)`$ of momentum $`p`$ and chirality $`\lambda `$ is defined according to KS methods . In the following we follow closely the notation of ref. (in particular we also use $`\zeta =\zeta _{}`$). In the above framework every spinor is transformed out of the two constant basic spinors $`𝔲_\lambda (\zeta )`$, of opposite chirality $`\lambda =\pm `$, as follows
$$u_\lambda (p)=\frac{1}{\sqrt{2p\zeta }}\overline{)}p𝔲_\lambda (\zeta ),𝔲_+(\zeta )=\overline{)}\eta 𝔲_{}(\zeta ),\eta ^2=1,(\eta \zeta )=0.$$
(197)
The usual relations hold: $`\overline{)}\zeta 𝔲_\lambda (\zeta )=0`$, $`\omega _\lambda 𝔲_\lambda (\zeta )=𝔲_\lambda (\zeta )`$, $`𝔲_\lambda (\zeta )\overline{𝔲}_\lambda (\zeta )=\overline{)}\zeta \omega _\lambda `$, $`\overline{)}pu_\lambda (p)=0`$, $`\omega _\lambda u_\lambda (p)=u_\lambda (p)`$, $`u_\lambda (p)\overline{u}_\lambda (p)=\overline{)}p\omega _\lambda `$, where $`\omega _\lambda =\frac{1}{2}(1+\lambda \gamma _5)`$. Spinors for the massive particle with four-momentum $`p`$ (with $`p^2=m^2`$) and spin projection $`\lambda /2`$ are defined similarly
$$u(p,\lambda )=\frac{1}{\sqrt{2p\zeta }}(\overline{)}p+m)u_\lambda (\zeta ),v(p,\lambda )=\frac{1}{\sqrt{2p\zeta }}(\overline{)}pm)u_\lambda (\zeta ).$$
(198)
or, equivalently, in terms of massless spinors
$$u(p,\lambda )=u_\lambda (p_\zeta )+\frac{m}{\sqrt{2p\zeta }}𝔲_\lambda (\zeta ),v(p,\lambda )=u_\lambda (p_\zeta )\frac{m}{\sqrt{2p\zeta }}𝔲_\lambda (\zeta ),$$
(199)
where $`p_\zeta \widehat{p}p\zeta m^2/(2\zeta p)`$ is the light-cone projection ($`p_\zeta ^2=0`$) of the $`p`$ obtained with the help of the constant auxiliary vector $`\zeta `$.
The above definition is supplemented in ref. with the precise prescription on spin quantization axes, translation from spin amplitudes to density matrices (also in vector notation) and the methodology of connecting production and decay for unstable fermions. We collectively call these rules global positioning of spin (GPS). Thanks to these we are able to easily introduce polarizations for beams and implement polarization effects for final fermion decays (of $`\tau `$-leptons, $`t`$-quarks), for the first time also in the presence of emission of many ISR and FSR photons!
The GPS rules determining the spin quantization frame for $`u(p,\pm )`$ and $`v(p,\pm )`$ of eq. (199) are summarized as follows:
* In the rest frame of the fermion, take the $`z`$-axis along $`\stackrel{}{\zeta }`$.
* Place the $`x`$-axis in the plane defined by the $`z`$-axis from the previous point and the vector $`\stackrel{}{\eta }`$, in the same half-plane as $`\stackrel{}{\eta }`$.
* With the $`y`$-axis, complete the right-handed system of coordinates. The rest frame defined in this way we call the GPS frame of the particular fermion.
See ref. for more details. In the following we shall assume that polarization vectors of beams and of outgoing fermions are defined in their corresponding GPS frames.
The inner product of the two massless spinors is defined as follows
$$s_+(p_1,p_2)\overline{u}_+(p_1)u_{}(p_2),s_{}(p_1,p_2)\overline{u}_{}(p_1)u_+(p_2)=(s_+(p_1,p_2))^{}.$$
(200)
The above inner product can be evaluated using the Kleiss-Stirling expression
$$s_+(p,q)=2(2p\zeta )^{1/2}(2q\zeta )^{1/2}\left[(p\zeta )(q\eta )(p\eta )(q\zeta )iϵ_{\mu \nu \rho \sigma }\zeta ^\mu \eta ^\nu p^\rho q^\sigma \right]$$
(201)
in any reference frame. In particular, in the laboratory frame we typically use $`\zeta =(1,1,0,0)`$ and $`\eta =(0,0,1,0)`$, which leads to the following “massless” inner product
$$s_+(p,q)=(q^2+iq^3)\sqrt{(p^0p^1)/(q^0q^1)}+(p^2+ip^3)\sqrt{(q^0q^1)/(p^0p^1)}.$$
(202)
Equation (199) immediately provides us also with the inner product for massive spinors
$$\begin{array}{cc}& \overline{u}(p_1,\lambda _1)u(p_2,\lambda _2)=S(p_1,m_1,\lambda _1,p_2,m_2,\lambda _2),\hfill \\ & \overline{u}(p_1,\lambda _1)v(p_2,\lambda _2)=S(p_1,m_1,\lambda _1,p_2,m_2,\lambda _2),\hfill \\ & \overline{v}(p_1,\lambda _1)u(p_2,\lambda _2)=S(p_1,m_1,\lambda _1,p_2,m_2,\lambda _2),\hfill \\ & \overline{v}(p_1,\lambda _1)v(p_2,\lambda _2)=S(p_1,m_1,\lambda _1,p_2,m_2,\lambda _2),\hfill \end{array}$$
(203)
where
$$S(p_1,m_1,\lambda _1,p_2,m_2,\lambda _2)=\delta _{\lambda _1,\lambda _2}s_{\lambda _1}(p_{1}^{}{}_{\zeta }{}^{},p_{2}^{}{}_{\zeta }{}^{})+\delta _{\lambda _1,\lambda _2}\left(m_1\sqrt{\frac{2\zeta p_2}{2\zeta p_1}}+m_2\sqrt{\frac{2\zeta p_1}{2\zeta p_2}}\right).$$
(204)
In our spinor algebra we shall exploit the completeness relations
$$\begin{array}{cc}& \overline{)}p+m=\underset{\lambda }{}u(p,\lambda )\overline{u}(p,\lambda ),\overline{)}pm=\underset{\lambda }{}v(p,\lambda )\overline{v}(p,\lambda ),\hfill \\ & \overline{)}k=\underset{\lambda }{}u(k,\lambda )\overline{u}(k,\lambda ),k^2=0.\hfill \end{array}$$
(205)
For a circularly polarized photon with four-momentum $`k`$ and helicity $`\sigma =\pm 1`$ we adopt the KS choice (see also ref. ) of polarization vector<sup>26</sup><sup>26</sup>26 Contrary to other papers on Weyl spinor techniques we keep here the explicitly complex conjugation in $`ϵ`$. This conjugation is cancelled by another conjugation following from Feynman rules, but only for outgoing photons, not for a beam photon, as in the Compton process, see ref. .
$$(ϵ_\sigma ^\mu (\beta ))^{}=\frac{\overline{u}_\sigma (k)\gamma ^\mu u_\sigma (\beta )}{\sqrt{2}\overline{u}_\sigma (k)u_\sigma (\beta )},(ϵ_\sigma ^\mu (\zeta ))^{}=\frac{\overline{u}_\sigma (k)\gamma ^\mu 𝔲_\sigma (\zeta )}{\sqrt{2}\overline{u}_\sigma (k)𝔲_\sigma (\zeta )},$$
(206)
where $`\beta `$ is an arbitrary light-like four-vector $`\beta ^2=0`$. The second choice with $`𝔲_\sigma (\zeta )`$ (not exploited in ) often leads to simplifications in the resulting photon emission amplitudes. Using the Chisholm identity<sup>27</sup><sup>27</sup>27 For $`\beta =\zeta `$ the identity is slightly different because of the additional minus sign in the “line-reversal” rule, i.e. $`\overline{u}_\sigma (k)\gamma ^\mu 𝔲_\sigma (\zeta )=\overline{𝔲}_\sigma (\zeta )\gamma ^\mu u_\sigma (k)`$, in contrast to the usual $`\overline{u}_\sigma (k)\gamma ^\mu u_\sigma (\beta )=+\overline{u}_\sigma (\beta )\gamma ^\mu u_\sigma (k).`$
$`\overline{u}_\sigma (k)\gamma _\mu u_\sigma (\beta )\gamma ^\mu `$ $`=2u_\sigma (\beta )\overline{u}_\sigma (k)+2u_\sigma (k)\overline{u}_\sigma (\beta ),`$ (207)
$`\overline{u}_\sigma (k)\gamma _\mu 𝔲_\sigma (\zeta )\gamma ^\mu `$ $`=2𝔲_\sigma (\zeta )\overline{u}_\sigma (k)2u_\sigma (k)\overline{𝔲}_\sigma (\zeta ),`$ (208)
we get two useful expressions, equivalent to eq. (206):
$$\begin{array}{cc}& (\overline{)}ϵ_\sigma (k,\beta ))^{}=\frac{\sqrt{2}}{\overline{u}_\sigma (k)u_\sigma (\beta )}\left[u_\sigma (\beta )\overline{u}_\sigma (k)+u_\sigma (k)\overline{u}_\sigma (\beta )\right]\hfill \\ & (\overline{)}ϵ_\sigma (k,\zeta ))^{}=\frac{\sqrt{2}}{\sqrt{2\zeta k}}\left[𝔲_\sigma (\zeta )\overline{u}_\sigma (k)u_\sigma (k)\overline{𝔲}_\sigma (\zeta )\right].\hfill \end{array}$$
(209)
In the evaluation of photon emission spin amplitudes we shall use the following important building block – the elements of the “transition matrices” $`U`$ and $`V`$ defined as follows
$$\begin{array}{cc}& \overline{u}(p_1,\lambda _1)\overline{)}ϵ_\sigma ^{}(k,\beta )u(p_2,\lambda _2)=U({}_{\sigma }{}^{k})[{}_{\lambda _1}{}^{p_1}{}_{\lambda _2}{}^{p_2}]=U_{\lambda _1,\lambda _2}^\sigma (k,p_1,m_1,p_2,m_2),\hfill \\ & \overline{v}(p_1,\lambda _1)\overline{)}ϵ_\sigma ^{}(k,\zeta )v(p_2,\lambda _2)=V({}_{\sigma }{}^{k})[{}_{\lambda _1}{}^{p_1}{}_{\lambda _2}{}^{p_2}]=V_{\lambda _1,\lambda _2}^\sigma (k,p_1,m_1,p_2,m_2).\hfill \end{array}$$
(210)
In the case of $`𝔲_\sigma (\zeta )`$ the above transition matrices are rather simple<sup>28</sup><sup>28</sup>28 Our $`U`$ and $`V`$ matrices are not the same as the $`M`$-matrices of ref. , but rather products of several of those. :
$`U^+(k,p_1,m_1,p_2,m_2)=\sqrt{2}\left[\begin{array}{cc}\sqrt{\frac{2\zeta p_2}{2\zeta k}}s_+(k,\widehat{p_1}),& 0\\ m_2\sqrt{\frac{2\zeta p_1}{2\zeta p_2}}m_1\sqrt{\frac{2\zeta p_2}{2\zeta p_1}},& \sqrt{\frac{2\zeta p_1}{2\zeta k}}s_+(k,\widehat{p_2})\end{array}\right],`$ (211)
$`U_{\lambda _1,\lambda _2}^{}(k,p_1,m_1,p_2,m_2)=\left[U_{\lambda _2,\lambda _1}^+(k,p_2,m_2,p_1,m_1)\right]^{},`$
$`V_{\lambda _1,\lambda _2}^\sigma (k,p_1,m_1,p_2,m_2)=U_{\lambda _1,\lambda _2}^\sigma (k,p_1,m_1,p_2,m_2).`$
The more general case with $`u_\sigma (\beta )`$ looks a little bit more complicated:
$$\begin{array}{cc}& U^+(k,p_1,m_1,p_2,m_2)=\frac{\sqrt{2}}{s_{}(k,\beta )}\times \hfill \\ & \left[\begin{array}{cc}s_+(\widehat{p}_1,k)s_{}(\beta ,\widehat{p}_2)+m_1m_2\sqrt{\frac{2\zeta \beta }{2\zeta p_1}\frac{2\zeta k}{2\zeta p_2}},& m_1\sqrt{\frac{2\zeta \beta }{2\zeta p_1}}s_+(k,\widehat{p}_2)+m_2\sqrt{\frac{2\zeta \beta }{2\zeta p_2}}s_+(\widehat{p}_1,k)\\ m_1\sqrt{\frac{2\zeta k}{2\zeta p_1}}s_{}(\beta ,\widehat{p}_2)+m_2\sqrt{\frac{2\zeta k}{2\zeta p_2}}s_{}(\widehat{p}_1,\beta ),& s_{}(\widehat{p}_1,\beta )s_+(k,\widehat{p}_2)+m_1m_2\sqrt{\frac{2\zeta \beta }{2\zeta p_1}\frac{2\zeta k}{2\zeta p_2}}\end{array}\right],\hfill \end{array}$$
(212)
with the same relations (LABEL:UVsimple2) and (LABEL:UVsimple3).
In the above the following numbering of elements in matrices $`U`$ and $`V`$ was adopted
$$\{(\lambda _1,\lambda _2)\}=\left[\begin{array}{cc}(++)& (+)\\ (+)& ()\end{array}\right].$$
(213)
When analysing (multi-) bremsstrahlung amplitudes we shall also often employ the following compact notation
$$U[{}_{\lambda _1}{}^{p}{}_{\sigma }{}^{k}{}_{\lambda _2}{}^{p}]=U_{\lambda _1,\lambda _2}^\sigma (k,p_1,m_1,p_2,m_2),V[{}_{\lambda _1}{}^{p}{}_{\sigma }{}^{k}{}_{\lambda _2}{}^{p}]=V_{\lambda _1,\lambda _2}^\sigma (k,p_1,m_1,p_2,m_2),$$
(214)
When analysing the soft real photon limit we shall exploit the following important diagonality property<sup>29</sup><sup>29</sup>29 Let us also keep in mind the relation $`b_\sigma (k,p)=(b_\sigma (k,p))^{}`$, which can save time in the numerical calculations.
$`U[{}_{\lambda _1}{}^{p}{}_{\sigma }{}^{k}{}_{\lambda _2}{}^{p}]=V[{}_{\lambda _1}{}^{p}{}_{\sigma }{}^{k}{}_{\lambda _2}{}^{p}]=b_\sigma (k,p)\delta _{\lambda _1\lambda _2},`$ (215)
$`b_\sigma (k,p)=\sqrt{2}{\displaystyle \frac{\overline{u}_\sigma (k)\overline{)}p𝔲_\sigma (\zeta )}{\overline{u}_\sigma (k)𝔲_\sigma (\zeta )}}=\sqrt{2}\sqrt{{\displaystyle \frac{2\zeta p}{2\zeta k}}}s_\sigma (k,\widehat{p}),`$
which also holds in the general case of $`u_\sigma (\beta )`$, where
$$b_\sigma (k,p)=\frac{\sqrt{2}}{s_\sigma (k,\beta )}\left(s_\sigma (\beta ,\widehat{p})s_\sigma (\widehat{p},k)+\frac{m^2}{2\zeta \widehat{p}}\sqrt{(2\beta \zeta )(2\zeta k)}\right).$$
(216)
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# Théorème 1
## 1 La trace de Dennis à coefficients.
Après avoir rappelé en 1.1 les propriétés du groupe $`K_1(A;𝐙/n)`$, on définit en 1.3 le groupe $`HH_1(A;𝐙/n)`$. La trace de Dennis à coefficients nécessite une notion adaptée de trace pour une $`1`$-forme différentielle non commutative. Cette construction est détaillée en 1.4 et la trace de Dennis à coefficients est construite en 1.5. Si l’anneau considéré est commutatif, une construction plus élémentaire est proposée en 1.6 grâce aux formes différentielles à la de Rham. Des traces d’ordres supérieurs, de source $`K_i(A;𝐙/n)`$, de but $`HC_i^{}(A;𝐙/n)`$ ou $`HH_i(A;𝐙/n)`$ ($`i1`$), sont définies en 1.7. Ces définitions 1.7 ne sont utilisées nulle part ailleurs dans le texte.
### 1.1 Le groupe $`K_1(A;𝐙/n)`$.
Soit $`A`$ un anneau. On note $`Proj(A)`$ la catégorie des $`A`$-modules à droite, projectifs et de type fini. Soit $`PProj(A)`$ et $`\alpha Aut_A(P)`$. Rappelons que le groupe de Bass $`K_1(A)`$ peut-être vu comme le quotient du groupe abélien libre engendré par les classes d’isomorphie de paires $`(P,\alpha )`$ modulo le sous-groupe engendré par les deux types d’éléments suivants:
(a) $`(P,\alpha )+(P,\beta )(P,\alpha \beta )`$
(b) $`(P_1P_2,\alpha _1\alpha _2)(P_1,\alpha _1)(P_2,\alpha _2).`$
On note $`[P,\alpha ]K_1(A)`$ la classe de $`(P,\alpha )`$
Soient $`A`$ et $`B`$ deux anneaux unitaires et soit $`\phi :Proj(A)Proj(B)`$ un foncteur additif. Le foncteur $`\phi `$ est dit cofinal (, VII.1, p. 345) si tout objet $`R`$ de $`Proj(B)`$ est facteur direct d’un objet de la forme $`\phi (P)`$ avec $`P`$ objet de $`Proj(A)`$, c’est-à-dire $`\phi (P)RS`$ avec $`SOb(Proj(B))`$.
À un tel foncteur cofinal $`\phi `$, on associe la catégorie $`𝒞(\phi )`$ dont les objets sont les triplets $`(P,\alpha ,Q)`$ avec $`POb(Proj(A))`$, $`QOb(Proj(A))`$, et où $`\alpha `$ est un isomorphisme de $`B`$-modules $`\phi (P)\phi (Q)`$.
Un morphisme dans $`𝒞(\phi )`$ de source $`(P,\alpha ,Q)`$, de but $`(P_1,\alpha _1,Q_1)`$ est un couple $`(f,g)`$ de morphismes de $`A`$-modules $`f:PP_1`$ et $`g:QQ_1`$ tels que $`\phi (g)\alpha =\alpha _1\phi (f)`$. L’ensemble des classes d’isomorphie d’objets de la catégorie $`𝒞(\phi )`$ est un monoïde abélien; on note $`K(𝒞(\phi ))`$ le groupe de Grothendieck associé.
###### Définition 7
Le quotient de $`K(𝒞(\phi ))`$ par le sous-groupe $`N`$ engendré par les éléments
$$(P,\alpha ,Q)+(Q,\beta ,R)(P,\beta \alpha ,R)$$
est noté $`K(\phi )`$. On désigne par $`[P,\alpha ,Q]`$ la classe de $`(P,\alpha ,Q)Ob(𝒞).`$
Remarque 8. Soit $`\alpha ^{}:PQ`$ un isomorphisme de $`A`$-modules. Alors l’élément $`x=[P,\phi (\alpha ^{}),Q]K(\phi )`$ est nul. En effet $`(\alpha ^{},\mathrm{id}_Q)`$ est un isomorphisme de la catégorie $`𝒞`$ de source
$`[P,\phi (\alpha ^{}),Q]`$, de but $`[Q,\mathrm{id}_{\phi (Q)},Q]`$, donc $`x=[Q,\mathrm{id}_{\phi (Q)},Q]=0`$. On en déduit qu’on peut toujours supposer qu’un élément de $`K(\phi )`$ est de la forme $`[P,\alpha ,L]`$ avec $`L`$ libre car si $`[P,\alpha ,Q]`$ appartient à $`K(\phi )`$ et si $`QR`$ est isomorphe à un module libre $`L`$, on a $`[P,\alpha ,Q]=[PR,\alpha \mathrm{id}_{nR},L]`$.
D’après , VII.5, p. 375, on a
###### Théorème 9
Soient $`A`$ et $`B`$ deux anneaux unitaires et soit $`\phi :Proj(A)Proj(B)`$ un foncteur cofinal. Alors, on a une suite exacte de groupes abéliens
où les applications $`\phi _1`$, $``$ et $`\phi _0`$ sont définies ainsi
$$\phi _1([P,\alpha ])=[\phi (P),\phi (\alpha )],$$
$$([P,\alpha ,Q])=[P][Q],$$
$$\phi _0([P])=[\phi (P)].$$
Pour $`[R,\lambda ]K_1(B)`$, l’application $`\rho `$ est définie par cofinalité de $`\phi `$ en écrivant $`RS\phi (P)`$ et en posant $`\rho ([R,\lambda ])=[P,\lambda id_S,P].`$
Remarque 10. De cette suite exacte, on déduit que le groupe $`K(\phi )`$ est l’extension
Appliquons cette construction au contexte suivant.
###### Définition 11
Soit $`n`$ un entier, $`n2`$. Le foncteur $`.n:Proj(A)Proj(A)`$ est défini par $`.n(P)=nP=P\mathrm{}P`$ et $`.n(f)=nf=f\mathrm{}f`$ ($`n`$ facteurs). Le groupe $`K(.n)`$ est noté $`K_1(A;𝐙/n)`$. Un élément de $`K_1(A;𝐙/n)`$ est de la forme $`[P,\alpha ,Q]`$, où $`P`$ et $`Q`$ sont dans $`Proj(A)`$ et où $`\alpha `$ est un isomorphisme de $`A`$-modules $`nPnQ`$.
On a par conséquent la suite exacte de $`K`$-théorie à coefficients
L’extension de la remarque ci-dessus s’écrit
$`()`$
avec $`\rho ([P,\alpha ])=[P,\alpha id_{(n1)P},P]`$ et $`([P,\alpha ,Q])=[P][Q].`$ Cette extension fournit le premier exemple de calcul de $`K`$-théorie à coefficients.
Exemple: Soit $`A`$ un anneau commutatif local. Pour $`n2`$, le groupe $`K_1(A;𝐙/n)`$ est égal à $`A^\times /(n)`$.
Revenons au cas général. On a les relations intéressantes suivantes.
###### Lemme 12
Soit $`A`$ un anneau et soit $`n`$ un entier, $`n>1`$.
Si $`n2\mathrm{mod}4`$ , on a $`nK_1(A;𝐙/n)=0`$ ;
et si $`n2\mathrm{mod}4`$ , on a $`2nK_1(A;𝐙/n)=0`$ .
En particulier, si $`p`$ est un nombre premier impair, $`K_1(A,𝐙/p)`$ est un $`𝐙/p`$-espace vectoriel.
Preuve : montrons d’abord que pour tout $`n>1`$, on a $`2nK_1(A;𝐙/n)=0`$. Si $`p`$ et $`q`$ sont deux entiers et si $`P\mathrm{Proj}(A)`$, on a un isomorphisme $`p(qP)q(pP)`$. Pour $`q=p=n`$, notons $`v_P:n^2Pn^2P`$ cet isomorphisme. La relation évidente $`v_P^2=\mathrm{id}`$ montre que l’élément $`x_P=[P,v_P,P]`$ de $`K_1(A;𝐙/n)`$ est tel que $`2x_P=0`$. Soit $`x=[P,\alpha ,Q]K_1(A;𝐙/n)`$ ; on a $`nx=[nP,\beta ,nQ]`$ avec $`\beta =v_Qn\alpha v_P`$, ce qui donne $`nx=x_Q+x^{}+x_P`$ avec $`x^{}=[nP,n\alpha ,nQ]=0`$. On en déduit $`2nx=0`$, ce qui montre $`2nK_1(A;𝐙/n)=0`$.
Par ailleurs, la suite exacte de $`K`$-théorie à coefficients donne $`n^2K_1(A;𝐙/n)`$ $`=0`$. Pour $`n`$ impair, les relations $`2nK_1(A;𝐙/n)=0`$ et $`n^2K_1(A;𝐙/n)=0`$ conduisent à $`nK_1(A;𝐙/n)=0`$.
Il reste à montrer que pour $`n0\mathrm{mod}4`$, on a également $`nK_1(A;𝐙/n)=0`$. Fixons une notation : pour $`\sigma 𝔖_n`$ et $`P\mathrm{Proj}(A)`$, on note $`\sigma `$ l’automorphisme de $`nP`$ défini pour $`(z_1,\mathrm{},z_n)nP`$ par $`\sigma (z_1,\mathrm{},z_n)=z_{\sigma (1)},\mathrm{},z_{\sigma (n)}nP`$. On vérifie facilement que pour tout $`n>0`$ et tout $`P\mathrm{Proj}(A)`$, l’isomorphisme $`v_P:n(nP)n(nP)`$ est le produit de $`{\displaystyle \frac{n(n1)}{2}}`$ transpositions à supports disjoints. Ce nombre est pair si $`n0`$ ou $`1\mathrm{mod}4`$. La proposition suivante montre que dans ce cas, l’élément $`[P,v_P,P]`$ de$`K_1(A;𝐙/n)`$ est nul et par suite que $`nK_1(A;𝐙/n)=0`$. $`\mathrm{}`$
###### Proposition 13
Soit $`n4`$ et $`\tau 𝔖_n`$ le produit de deux transpositions à supports disjoints. Alors pour tout $`P\mathrm{Proj}(A)`$ l’élément $`[P,\tau ,P]`$ de $`K_1(A;𝐙/n)`$ est nul.
Preuve : dans $`𝔖_n`$, $`\tau `$ est conjugué à $`\tau ^{}=(12)(34)`$. Dans $`K_1(A;𝐙/n)`$, on a donc $`[P,\tau ,P]=[P,\tau ^{},P]`$. La matrice $`T^{}\mathrm{GL}(𝐙)`$ de $`\tau ^{}`$ appartient à $`[\mathrm{GL}(𝐙),\mathrm{GL}(𝐙)]`$ donc $`[P,\tau ^{},P]=0`$. $`\mathrm{}`$
### 1.2 L’algèbre différentielle graduée $`\mathrm{\Omega }_{nc}^{}(A)`$.
Soit $`k`$ un anneau commutatif unitaire et soit $`A`$ une $`k`$-algèbre unitaire. Désignons par $`\mu :A_kAA`$ la multiplication de $`A`$. Posons $`\mathrm{\Omega }_{nc}^1(A):=\mathrm{ker}\mu `$. On sait que $`\mathrm{\Omega }_{nc}^1(A)`$ est le sous-bimodule de $`A_kA`$ engendré par
$$\{1aa1,aA\}$$
On introduit la 1-forme différentielle non commutative
$$d_{nc}a:=1aa1.$$
La structure de $`A`$-bimodule évident de $`\mathrm{\Omega }_{nc}^1(A)`$ se lit en terme de formes différentielles par la relation
$`()`$
$$d_{nc}(a_1)a_2=d_{nc}(a_1a_2)a_1d_{nc}a_2$$
de sorte que tout élément de $`\mathrm{\Omega }_{nc}^1(A)`$ s’écrit comme somme de 1-formes différentielles $`a_0d_{nc}a_1`$. En tant que $`k`$-module, on a $`\mathrm{\Omega }_{nc}^1(A)A_k(A/k)`$.
A partir de $`\mathrm{\Omega }_{nc}^1(A)`$, on construit une $`k`$-algèbre différentielle graduée $`\mathrm{\Omega }_{nc}^{}(A)`$ en posant $`\mathrm{\Omega }_{nc}^r(A):=\mathrm{\Omega }_{nc}^1(A)_A\mathrm{}_A\mathrm{\Omega }_{nc}^1(A)`$ ($`r`$ facteurs). La structure de $`A`$-module à droite sur $`\mathrm{\Omega }_{nc}^r(A)`$ se définit en jouant à saute-moutons à l’aide de la formule $`()`$ :
$$a_0d_{nc}a_1\mathrm{}d_{nc}a_r\alpha =a_0d_{nc}a_1\mathrm{}d_{nc}(a_r\alpha )a_0d_{nc}a_1\mathrm{}d_{nc}a_{r1}a_rd_{nc}\alpha =\mathrm{etc}.$$
Tout élément de $`\mathrm{\Omega }_{nc}^r(A)`$ est somme de $`r`$-formes différentielles du type
$$a_0d_{nc}a_1\mathrm{}d_{nc}a_r.$$
C’est encore la relation $`()`$ qui définit le produit de $`\omega _r\mathrm{\Omega }_{nc}^r(A)`$ par $`\omega _s\mathrm{\Omega }_{nc}^s(A)`$.
En tant que $`k`$-module, on a $`\mathrm{\Omega }_{nc}^r(A)=A(A/k)^r`$ et le bord de Hochschild de $`A`$ admet une expression très simple dans l’algèbre $`\mathrm{\Omega }_{nc}^{}(A)`$ : pour $`\omega _r\mathrm{\Omega }_{nc}^r(A)`$ avec $`\omega _r=\omega ^{}d_{nc}\alpha ,\alpha A`$, on a $`b(\omega _r)=(1)^{r1}(\omega ^{}\alpha \alpha \omega ^{})`$. Le complexe de Hochschild normalisé, traditionnellement noté $`(\overline{C}_{}(A),b)`$ n’est autre que $`(\mathrm{\Omega }_{nc}^{}(A),b)`$. Pour $`\omega \mathrm{\Omega }_{nc}^{}(A)`$ tel que $`b(\omega )=0`$, on note $`[\omega ]`$ sa classe d’homologie de Hochschild.
### 1.3 Le $`k`$-module gradué $`HH_{}(A;𝐙/n)`$.
Soit $`k`$ un anneau commutatif unitaire, soient $`(C_{},d)`$ et $`(C_{}^{},d^{})`$ des complexes (de chaînes) de $`k`$-modules et soit $`f:(C_{},d)(C_{}^{},d^{})`$ un morphisme de complexes. Le cône de $`f`$ est le complexe de chaînes $`(\mathrm{co}(f),)`$ défini par $`\mathrm{co}(f)_r=C_r^{}C_{r1}`$ et $`(a^{},a)=(d^{}(a^{})+f(a),d(a))`$. De la suite exacte courte de complexes
$$0(C_{}^{},d^{})(\mathrm{co}(f),)(C_{},d)_{[1]}0$$
on déduit une suite exacte longue d’homologie dont le connectant$`H_r\left((C_{},d)_{[1]}\right)H_{r1}(C_{}^{},d^{})`$ n’est autre que $`H_{r1}(f)`$.
###### Définition 14
Soient $`k`$ un anneau commutatif unitaire et $`A`$ une $`k`$-algèbre unitaire. Soit $`(\mathrm{\Omega }_{}(A),b)`$ le complexe de Hochschild de $`A`$. L’homologie du cône de l’application $`.n:(\mathrm{\Omega }_{}(A),b)(\mathrm{\Omega }_{}(A),b)`$ définie par $`.n(a_0a_1\mathrm{}a_r)=na_0a_1\mathrm{}a_r`$ s’appelle l’homologie de Hochschild à coefficients $`𝐙/n`$ de l’algèbre $`A`$. On note cette homologie $`HH_{}(A;𝐙/n)`$.
Pour $`(\omega _r,\omega _{r1})\mathrm{\Omega }^r(A)\mathrm{\Omega }^{r1}(A)`$ telle que $`(\omega _r,\omega _{r1})=0`$, on note $`[\omega _r,\omega _{r1}]`$ sa classe d’homologie de Hochschild à coefficients.
De cette définition de $`HH_{}(A;𝐙/n)`$, on extrait la suite exacte longue
avec $`\rho ([\omega _r])=[\omega _r\mathrm{,0}]`$ et $`([\omega _r,\omega _{r1}])=[\omega _{r1}]`$.
### 1.4 Calcul différentiel non commutatif d’ordre 1.
Soit $`P`$ un $`A`$-module projectif, à droite, de type fini. Un système de coordonnées de $`P`$ (cf. , II.46) est une suite $`𝒮=(x_j,\phi _j)_{1jr}`$ avec $`x_jP`$, $`\phi _jP^{}=\mathrm{Hom}_A(P,A)`$ telle que pour tout $`xP`$, on ait $`x=_{j=1}^rx_j\phi _j(x)`$.
Soit $`u\mathrm{End}P`$. Par définition, la matrice de $`u`$ dans le système de coordonnées $`𝒮`$ est la matrice $`U=(U_{ij})\mathrm{Mat}_{r,r}(A)`$ définie par $`U_{ij}=\phi _iu(x_j)`$. La quantité $`_{i=1}^r\phi _iu(x_i)A`$ s’appelle la trace de $`u`$ relativement au système de coordonnées $`𝒮`$. On la note $`\mathrm{tr}(u,𝒮)`$, ou $`\mathrm{tr}(u)`$. On pose $`\mathrm{rg}(P)=\mathrm{rg}(P,𝒮)=\mathrm{tr}(\mathrm{id}_P,𝒮)=_{i=1}^r\phi _i(x_i)A`$. Ces quantités dépendent du système de coordonnées $`𝒮`$ choisi.
La connexion de Levi-Civita de $`P`$ (cf. et ) est l’application $`d_P:PP_A\mathrm{\Omega }^1(A)`$ dont l’expression dans le système de coordonnées $`𝒮`$ est
$$d_P\left(\underset{j=1}{\overset{r}{}}x_j\phi _j(x)\right)=\underset{j=1}{\overset{r}{}}x_jd_{nc}(\phi _j(x)).$$
Comme toute connexion, $`d_P`$ n’est pas une application $`A`$-linéaire mais satisfait à la relation :
$$d_P(xa)=d_P(x)a+xd_{nc}a,xP,aA.$$
###### Définition 15
Soit $`\alpha :PQ`$ une application $`A`$-linéaire à droite entre les modules projectifs de type fini $`P`$ et $`Q`$. Soient $`d_P`$ et $`d_Q`$ les connexions de Levi-Civita de $`P`$ et $`Q`$. L’application $`A`$-linéaire
$`d_{nc}\alpha :PQ_A\mathrm{\Omega }^1(A)\text{ est définie par}`$
$`d_{nc}\alpha =d_Q\alpha (\alpha \mathrm{id})d_P.`$
L’abus d’écriture fréquent qui consiste à poser $`d_P=d_Q=d_{nc}`$ conduit à la formule $`d_{nc}(\alpha (x))=d_{nc}\alpha (x)+\alpha (d_{nc}x)`$, $`xP`$.
Soient $`𝒮=(x_j,\phi _j)_{1jr}`$ et $`𝒮^{}=(y_i,\psi _i)_{1is}`$ des systèmes de coordonnées respectifs de $`P`$ et $`Q`$. Il est facile de donner une interprétation matricielle de l’application $`d\alpha `$. Si $`M=(M_{ij})\mathrm{Mat}_{s,r}(A)`$ est la matrice de $`\alpha `$, exprimée dans les systèmes de coordonnées $`𝒮`$ et $`𝒮^{}`$, avec $`M_{ij}=\psi _i\alpha (x_j)`$, un calcul simple montre que pour $`x=_{j=1}^rx_j\phi _j(x)`$, on a $`d_{nc}\alpha (x)=_{i,j}y_id(M_{ij})\phi _j(x)`$, ce qui permet de dire que $`d_{nc}\alpha `$ a pour matrice $`d_{nc}M=(d_{nc}M_{ij})\mathrm{Mat}_{s,r}(\mathrm{\Omega }^1(A))`$ dans les systèmes $`𝒮`$ et $`𝒮^{}`$.
Supposons de plus que $`\alpha :PQ`$ soit un isomorphisme de $`A`$-modules. Soit $`N\mathrm{Mat}_{s,r}(A)`$ la matrice de $`\alpha ^1`$ exprimée dans les systèmes de coordonnées $`𝒮^{}`$ et $`𝒮`$. On a $`MN=\mathrm{mat}_𝒮^{}(\mathrm{id}_Q)`$ et $`NM=\mathrm{mat}_𝒮(\mathrm{id}_P)`$. L’application composée , évidemment notée$`\alpha ^1d_{nc}\alpha `$ est $`A`$-linéaire, de matrice $`NdM\mathrm{Mat}_{r,r}(\mathrm{\Omega }^1(A))`$. La trace de cette matrice, $`\mathrm{tr}(NdM)\mathrm{\Omega }^1(A)`$, s’appelle la trace de $`\alpha ^1d_{nc}\alpha `$ dans les systèmes de coordonnées $`𝒮`$ et $`𝒮^{}`$.
###### Proposition 16
Soient $`P`$ et $`Q`$ deux $`A`$-modules à droite, projectifs et de type fini, de systèmes de coordonnées respectifs $`𝒮`$ et $`𝒮^{}`$. Posons $`p=\mathrm{rg}(P,𝒮)`$ et $`q=\mathrm{rg}(Q,𝒮^{})`$. Soit $`\alpha :PQ`$ un isomorphisme de $`A`$-modules et soit $`\mathrm{tr}(\alpha ^1d_{nc}\alpha )\mathrm{\Omega }^1(A)`$ la trace de $`\alpha ^1d_{nc}\alpha `$, exprimée dans les systèmes de coordonnés $`𝒮`$ et $`𝒮^{}`$.
Alors le bord de Hochschild de $`\mathrm{tr}(\alpha ^1d_{nc}\alpha )`$ est donné par $`b(\mathrm{tr}(\alpha ^1d_{nc}\alpha ))=qpA`$.
Preuve : on a $`\mathrm{tr}(\alpha ^1d_{nc}\alpha )=\mathrm{tr}(Nd_{nc}M)`$ d’où
$`b(\mathrm{tr}(\alpha ^1d_{nc}\alpha ))=b(\mathrm{tr}Nd_{nc}M)=\mathrm{tr}NM+\mathrm{tr}MN=\mathrm{tr}(\mathrm{id}_P,𝒮)+\mathrm{tr}(\mathrm{id}_Q,𝒮^{}).`$ $`\mathrm{}`$
Remarque 17. Dans la proposition ci-dessus, si $`P=Q`$ et si $`𝒮=𝒮^{}`$, alors pour tout automorphismes $`\alpha `$ de $`P`$, $`\mathrm{tr}(\alpha ^1d_{nc}\alpha ,𝒮)`$ est un cycle de Hochschild. Dans ce cas on note $`[\mathrm{tr}(\alpha ^1d_{nc}\alpha ]HH_1(A)`$, sa classe d’homologie de Hochschild.
Nous utiliserons également le résultat suivant, de preuve facile laissée au lecteur.
###### Proposition 18
Soient $`P`$, $`Q`$ et $`R`$ des $`A`$-modules à droite, projectifs et de type fini ; soient $`\alpha \mathrm{Hom}_A(P,Q)`$, $`\beta \mathrm{Hom}_A(Q,R)`$. Alors on a la relation
$$d_{nc}(\beta \alpha )=d_{nc}\beta \alpha +\beta d_{nc}\alpha .$$
En particulier $`\alpha ^1d_{nc}\alpha =d_{nc}\alpha ^1\alpha `$.
### 1.5 La trace de Dennis à coefficients $`\overline{D}_1`$.
Pour tout entier $`r0`$, K. Dennis () a construit un morphisme $`D_r:K_r(A)HH_r(A)`$. Rappelons que si $`r=0`$ et $`[P]K_0(A)`$, on a $`D_0([P])=\mathrm{rg}(P,𝒮)\mathrm{mod}[A,A]`$ et que si $`r=1`$ et $`[P,\alpha ]K_1(A)`$, on a $`D_1([P,\alpha ])=[\mathrm{tr}(\alpha ^1d_{nc}\alpha )]HH_1(A)`$.
###### Théorème 19
Soient $`A`$ un anneau unitaire, $`n`$ un entier, $`n2`$ et soit $`x=[P_1,\alpha ,P_2]K_1(A;𝐙/n)`$. Après avoir fixé des systèmes de coordonnées $`𝒮_1`$ et $`𝒮_2`$ sur $`P_1`$ et $`P_2`$, on pose $`p_1=\mathrm{rg}(P_1,𝒮_1),p_2=\mathrm{rg}(P_2,𝒮_2)`$ et $`\overline{D}_1(x)=[\mathrm{tr}(\alpha ^1d_{nc}\alpha ),p_1p_2]HH_1(A;𝐙/n)`$.
Alors l’application $`\overline{D}_1`$ est un morphisme de groupes s’insérant dans le diagramme commutatif à lignes exactes ci-dessous
Preuve : la quantité $`c=(\mathrm{tr}(\alpha ^1d_{nc}\alpha ),p_1p2)`$ est un cycle du cône de la multiplication $`.n:(\mathrm{\Omega }^{}(A),b)(\mathrm{\Omega }^{}(A),b)`$. En effet,
$$(c)=b(\mathrm{tr}(\alpha ^1d_{nc}\alpha ))+n(p_1p_2);$$
or
$$b(\mathrm{tr}(\alpha ^1d_{nc}\alpha ))=\mathrm{rg}(nP_2,n𝒮_2)\mathrm{rg}(nP_1,n𝒮_1),$$
donc $`(c)=0.`$
Soit $`[c]=[\mathrm{tr}\alpha ^1d_{nc}\alpha ,p_1p_2]HH_1(A,𝐙/n)`$ la classe d’homologie du cycle $`c`$. Cette classe est indépendante du choix du représentant $`(P_1,\alpha ,P_2)`$ de $`x`$. Supposons d’abord $`(P_1,\alpha ,P_2)(P_1^{},\alpha ^{},P_2^{})`$ dans la catégorie $`𝒞`$. Il existe dans ce cas un couple d’isomorphismes $`f_i:P_iP_i^{}`$ tel que $`\alpha ^{}(nf_1)=(nf_2)\alpha `$. Choisissons des systèmes de coordonnées $`𝒮_{}^{}{}_{i}{}^{}`$ sur $`P_i^{}`$. Posons
$`c=(\mathrm{tr}\alpha ^1d_{nc}\alpha ,\mathrm{rg}(P_1,𝒮_1)\mathrm{rg}(P_2,𝒮_2))\text{ et}`$
$`c^{}=(\mathrm{tr}\alpha ^{}{}_{}{}^{1}d_{nc}^{}\alpha ^{},\mathrm{rg}(P_1^{},𝒮_{}^{}{}_{1}{}^{})\mathrm{rg}(P_2^{},𝒮_{}^{}{}_{2}{}^{})).`$
On a $`c=c^{}=0`$. Introduisons
$`\omega _1`$ $`=`$ $`\mathrm{tr}(f_2^1d_{nc}f_2f_1^1d_{nc}f_1)\text{ et}`$
$`\omega _2`$ $`=`$ $`\mathrm{tr}((\alpha ^{})^1d_{nc}(nf_2)d_{nc}(\alpha nf_1^1)+(nf_1)\alpha ^1d_{nc}\alpha d_{nc}(nf_1^1));`$
on a $`\omega _i\mathrm{\Omega }^i(A)`$ et $`(\omega _2,\omega _1)=c^{}c`$, c’est-à-dire $`[c^{}]=[c]`$ dans $`HH_1(A,𝐙/n)`$.
Le même argument montre que la classe de $`c`$ est indépendante du choix des systèmes de coordonnés $`𝒮_1`$ et $`𝒮_2`$ retenus sur $`P_1`$ et $`P_2`$.
Montrons à présent qu’on a un morphisme de groupes
$$D:K_0(𝒞)HH_1(A;𝐙/n)$$
en posant $`D((P_1,\alpha ,P_2)^\widehat{})=[\mathrm{tr}(\alpha ^1d_{nc}\alpha ),\mathrm{rg}(P_1,𝒮_1)\mathrm{rg}(P_2,𝒮_2)]`$. Pour $`i=\mathrm{1,2,3}`$, soient $`t_i=(P_i,\alpha _i,Q_i)^\widehat{}`$ trois éléments de $`K_0(𝒞)`$ tels que $`t_1+t_2=t_3`$. Après avoir fixé des systèmes de coordonnées $`𝒮_i`$ et $`𝒮_{}^{}{}_{i}{}^{}`$ sur $`P_i`$ et $`Q_i`$ $`(i=\mathrm{1,2})`$, on choisit $`𝒮_1𝒮_2`$ (resp. $`𝒮_{}^{}{}_{1}{}^{}𝒮_{}^{}{}_{2}{}^{}`$) comme système de coordonnées de $`P_3`$ (resp. $`Q_3`$).
Posons $`c_i=(\mathrm{tr}(\alpha _i^1d_{nc}\alpha _i),\mathrm{rg}(P_i,𝒮_i)\mathrm{rg}(Q_i,𝒮_i))`$ pour $`i=\mathrm{1,2}`$.
De $`\alpha _3=\alpha _1\alpha _2`$ et du choix proposé pour les systèmes de coordonnées, on tire immédiatement $`c_3=c_1+c_2`$. Cette égalité entre cocycles entraîne$`[c_1]+[c_2]=[c_3]`$.
Il reste à montrer $`D(N)=0`$, où $`N`$ est le sous-groupe de $`K_0(𝒞)`$ tel que $`K_1(A;𝐙/n)=K_0(𝒞)/N`$. Soient $`z_1=(P_1,\alpha ,P_2)^\widehat{}`$ et $`z_2=(P_2,\beta ,P_3)^\widehat{}`$ deux éléments de $`K_0(𝒞)`$; on pose $`z=(P_1,\beta \alpha ,P_3)^\widehat{}`$. Après avoir choisi des systèmes de coordonnées $`𝒮_i`$ sur $`P_i`$, on pose
$`c_1`$ $`=`$ $`(tr\alpha ^1d_{nc}\alpha ,\mathrm{rg}(P_1,𝒮_1)\mathrm{rg}(P_2,𝒮_2)),`$
$`c_2`$ $`=`$ $`(\mathrm{tr}\beta ^1d_{nc}\beta ,\mathrm{rg}(P_2,𝒮_2)\mathrm{rg}(P_3,𝒮_3),`$
$`c`$ $`=`$ $`(\mathrm{tr}(\beta \alpha )^1d_{nc}(\beta \alpha ),\mathrm{rg}(P_1,𝒮_1)\mathrm{rg}(P_3,𝒮_3)),`$
$`\theta _2`$ $`=`$ $`\mathrm{tr}(\alpha ^1\beta ^1d_{nc}\beta d_{nc}\alpha ),`$
$`\theta _1`$ $`=`$ $`0.`$
On a $`c_1+c_2c_3=(\theta _2,\theta _1)`$, ce qui montre que l’application $`D`$ factorise en un morphisme $`D_1^{(n)}:K_1(A,𝐙/n)HH_1(A;𝐙/n)`$. En remarquant que$`\mathrm{tr}(\alpha \mathrm{id})^1d(\alpha \mathrm{id})=\mathrm{tr}(\alpha ^1d\alpha )`$, il est immédiat de vérifier que les diagrammes commutent. $`\mathrm{}`$
###### Corollaire 20
Posons $`\stackrel{~}{HH_1}(A;𝐙/n)=HH_1(A;𝐙/n)/\mathrm{Im}(\rho D_1)`$. La “classe caractéristique secondaire”
$$\overline{d}_1:K_0(A)_{(n)}\stackrel{~}{HH_1}(A;𝐙/n).$$
définie pour $`x=(y)K_0(A)_{(n)}`$ par $`\overline{d}_1(x)=\overline{D}_1(y)\mathrm{mod}\text{Im}\rho D_1`$ est un morphisme de groupes abéliens.
Nous montrerons plus loin que cette classe caractéristique secondaire n’est pas triviale en général.
### 1.6 Le cas des algèbres commutatives.
Si l’algèbre $`A`$ est commutative, on peut simplifier notablement la construction de la trace de Dennis à coefficients en transitant par le $`A`$-module $`\mathrm{\Omega }_{dR}^1(A)`$ des différentielles de Kähler de $`A`$. Fixons pour cela un peu de vocabulaire.
Le module des différentielles de Kähler de $`A`$ est le $`A`$-module $`\mathrm{\Omega }_{dR}^1(A):=\mathrm{ker}\mu /(\mathrm{ker}\mu )^2`$, où $`\mu :A_kAA`$ est la multiplication de $`A`$. Ce $`A`$-module est engendré par les symboles $`da`$, $`aA`$, soumis aux relations $`d\lambda =0`$, $`\lambda k`$, $`d(a_0a_1)=a_0da_1+a_1da_0`$.
On sait que lorsque $`A`$ est commutative, la trace de Dennis $`D_1`$ est essentiellement la dérivée logarithmique. En effet, de la commutativité de $`A`$, on déduit
$$HH_1(A)=\mathrm{\Omega }_{nc}^1(A)/[A,\mathrm{\Omega }_{nc}^1(A)].$$
L’application $`\gamma :HH_1(A)\mathrm{\Omega }_{dR}^1(A)`$ définie par $`\gamma ([a_0d_{nc}a_1])=a_0da_1`$ est un isomorphisme. On en déduit le morphisme surjectif de $`A`$-bimodules $`\rho :\mathrm{\Omega }_{nc}^1(A)\mathrm{\Omega }_{dR}^1(A)`$ défini par $`\rho (a_0d_{nc}a_1)=a_0da_1`$.
Pour $`x=[P,\alpha ]K_1(A)`$, on obtient facilement $`\gamma D_1(x)=det(\alpha )^1d(det(\alpha )).`$ En écrivant $`K_1(A)=A^\times SK_1(A)`$, on en déduit que la restriction de $`D_1`$ à $`SK_1(A)`$ est nulle et que la restriction de la trace de Dennis à $`A^\times `$ est la dérivée logarithmique, c’est-à-dire que pour $`uA^\times `$, on a $`\gamma D_1(u)=u^1du`$.
Soit $`P`$ un $`A`$-module projectif de type fini sur $`A`$ et soit $`𝒮=(x_j,\phi _j)_{1jr}`$ un système de coordonnées sur $`P`$. La connexion de Levi-Civita commutative de $`P`$ est l’application $`d_P^{}:PP_A\mathrm{\Omega }_{dR}^1(A)`$ définie pour $`x=_{j=1}^rx_j\phi _j(x)`$ par $`d_p^{}(x)=_{j=1}^rx_jd\phi _j(x)`$. Remarquons que si $`d_P:PP_A\mathrm{\Omega }_{nc}^1(A)`$ désigne la connexion de Levi-Civita non commutative, le diagramme suivant est commutatif.
Soient $`P`$ et $`Q`$ deux $`A`$-modules projectifs de type fini et soit $`\alpha :PQ`$ une application $`A`$-linéaire. Soient
$$:PP_A\mathrm{\Omega }_{dR}^1(A)$$
et
$$^{}:QQ_A\mathrm{\Omega }_{dR}^1(A)$$
des connexions sur $`P`$ et $`Q`$ respectivement. On définit l’application $`A`$-linéaire $`d\alpha =d(\alpha ,,^{})`$ de source $`P`$, de but $`Q_A\mathrm{\Omega }_{dR}^1(A)`$ par
$$d(\alpha ,,^{})=^{}\alpha (\alpha id).$$
Soit $`\alpha :PQ`$ un isomorphisme de $`A`$-module. En choisissant des systèmes de coordonnées sur $`P`$ et $`Q`$, l’application $`A`$-linéaire
$$\alpha ^1d\alpha :=(\alpha ^1\mathrm{id})d(\alpha ,,^{})$$
admet une matrice carrée à coefficients dans $`\mathrm{\Omega }_{dR}^1(A)`$ dont la trace est notée $`tr(\alpha ^1d\alpha )`$.
###### Théorème 21
Soient $`k`$ un anneau commutatif unitaire, $`A`$ une $`k`$-algèbre commutative et soit $`n2`$ un entier.
Soit
$$D_1^{(n)}:K_1(A;𝐙/n)\mathrm{\Omega }_{dR}^1(A)/(n)$$
l’application définie pour $`x=[P,\alpha ,Q]`$ dans $`K_1(A,𝐙/n)`$ par
$$D_1^{(n)}(x)=tr(\alpha ^1d(\alpha ,n,n^{}))\mathrm{mod}n\mathrm{\Omega }_{dR}^1(A),$$
$``$ et $`^{}`$ sont des connexions sur $`P`$ et $`Q`$ respectivement.
Alors l’application $`D_1^{(n)}`$ est un morphisme de groupes abéliens.
Remarque 22. Supposons de plus que l’algèbre $`A`$ soit intègre et que $`n`$ soit un entier premier à la caractéristique de $`A`$. La suite exacte longue 1.3 conduit à l’isomorphisme $`HH_1(A;𝐙/n)HH_1(A)/(n)`$. Notons $`\overline{\gamma }:HH_1(A)/(n)\mathrm{\Omega }_{dR}^1(A)/(n)`$ l’isomorphime induit de $`\gamma :HH_1(A)\mathrm{\Omega }_{dR}^1(A)`$. A l’aide des connections de Levi-Civita, on vérifie que l’application $`D_1^{(n)}`$ s’insère dans le diagramme commutatif
Preuve du théorème 23: avec les notations de 1.1, on a $`K_1(A;𝐙/n)=K_0(𝒞)/N`$. Soit $`(P,\alpha ,Q)`$ un objet de $`𝒞`$.
Si $``$ et $`_1`$ sont deux connexions sur $`P`$ et si $`^{}`$ est une connexion sur $`Q`$, on a
$$\alpha ^1d(\alpha ,n,n^{})\alpha ^1d(\alpha ,n_1,n^{})=n(_1),$$
application $`A`$-linéaire dont la trace est congrue à $`0`$ modulo $`n\mathrm{\Omega }_{dR}^1(A).`$
Si $``$ est une connexion sur $`P`$ et si $`^{}`$ et $`_1^{}`$ sont deux connexions sur $`Q`$, on a
$$\alpha ^1d(\alpha ,n,n^{})\alpha ^1d(\alpha ,n,n_1^{})=\alpha ^1n(^{}_1^{})\alpha ,$$
application $`A`$-linéaire dont la trace, égale à celle de $`n(^{}_1^{})`$, est bien congrue à $`0`$ modulo $`n\mathrm{\Omega }_{dR}^1(A)`$.
Ces deux remarques montrent que pour $`(P,\alpha ,Q)Ob(𝒞)`$, le choix des connexions sur $`P`$ ou $`Q`$ n’intervient pas pour la définition de $`tr(\alpha ^1d\alpha )`$ modulo $`n\mathrm{\Omega }_{dR}^1(A)`$. C’est pourquoi, pour alléger, nous omettons à présent de préciser les connexions choisies.
Si les objets $`(P,\alpha ,Q)`$ et $`(P_1,\alpha _1,Q_1)`$ sont isomorphes dans la catégorie $`𝒞`$, il existe des applications $`A`$-linéaires $`f`$ et $`g`$ telles que $`\alpha _1=ng\alpha nf^1`$. Un rapide calcul donne
$$\mathrm{tr}(\alpha _1^1d\alpha _1)=n\mathrm{tr}(g^1dg)+\mathrm{tr}(\alpha ^1d\alpha )+n\mathrm{tr}(fdf^1)$$
soit
$$\mathrm{tr}(\alpha _1^1d\alpha _1)\mathrm{tr}(\alpha ^1d\alpha )\mathrm{mod}n\mathrm{\Omega }_{dR}^1(A),$$
ce qui montre que seule la classe d’isomorphie de l’objet $`(P,\alpha ,Q)`$ intervient pour la définition de la trace à coefficients. Enfin, les relations banales
$$\mathrm{tr}\left((\alpha _1\alpha _2)^1d(\alpha _1\alpha _2)\right)=\mathrm{tr}(\alpha _1^1d\alpha _1)+\mathrm{tr}(\alpha _2^1d\alpha _2)$$
et
$$\mathrm{tr}\left((\alpha \beta )^1d(\alpha \beta )\right)=\mathrm{tr}(\alpha ^1d\alpha )+\mathrm{tr}(\beta ^1d\beta )$$
montrent qu’on a un morphisme de groupes $`D_1^{(n)}`$ de source $`K_1(A;𝐙/n)`$ de but $`\mathrm{\Omega }_{dR}^1(A)/(n)`$ en posant
$$D_1^{(n)}\left([P,\alpha ,Q]\right)=\mathrm{tr}(\alpha ^1d\alpha )\mathrm{mod}n\mathrm{\Omega }_{dR}^1(A).$$
Exemple. Soit $`A`$ un anneau commutatif. On suppose $`K_0(A)=𝐙`$. Alors, $`K_1(A;𝐙/n)=K_1(A)/(n)=A^\times /(n)SK_1(A)/(n)`$. La restriction de la trace de Dennis $`D_1^{(n)}`$ au facteur $`SK_1(A)/(n)`$ est nulle tandis que la restriction de la trace de Dennis $`D_1^{(n)}`$ au facteur $`A^\times /(n)`$ est donnée par
$$D_1^{(n)}([a])=a^1da\mathrm{mod}n\mathrm{\Omega }_{dR}^1(A).$$
Cette situation s’applique en particulier lorsque $`A`$ est local.
Pour tout anneau commutatif, l’image de $`K_1(A)`$ par la trace de Dennis $`D_1`$ est le sous-groupe $`dA^\times /A^\times `$ de $`\mathrm{\Omega }_{dR}^1(A)`$ engendré par $`\{u^1du,uA^\times \}`$. Du théorème 23, on déduit:
###### Corollaire 23
Soient $`k`$ un anneau commutatif unitaire, $`A`$ une $`k`$-algèbre commutative et $`n`$ un entier. On désigne par $`A^\times `$ le groupe des unités de $`A`$ et par $`dA^\times /A^\times `$ le sous-groupe de $`\mathrm{\Omega }_{dR}^1(A)`$ engendré par $`\{u^1du,uA^\times \}`$. Soit $`S`$ le sous-groupe de $`\mathrm{\Omega }_{dR}^1(A)`$ engendré par $`n\mathrm{\Omega }_{dR}^1(A)`$ et $`dA^\times /A^\times `$. La “classe caractéristique secondaire”
$$d_1^{(n)}:\stackrel{~}{K}_0(A)_{(n)}\mathrm{\Omega }_{dR}^1(A)/S$$
définie pour $`x=(y)\stackrel{~}{K}_0(A)_{(n)}`$ par $`d_1^{(n)}(x)=D_1^{(n)}(y)\mathrm{mod}S`$ est un morphisme de groupes abéliens.
Nous verrons plus bas que ce morphisme n’est pas trivial.
### 1.7 Les traces d’ordre supérieur.
Soit $`A`$ une $`k`$-algèbre. Désignons par $`HH_{}(A)`$, $`HC_{}^{}(A)`$ et $`HC_{}^{per}(A)`$ respectivement les homologies de Hochschild, “cyclique négative” et “cyclique périodique” de $`A`$ (voir , 1.1 et 5.1 pour les définitions). Soit $`D_{}:K_{}(A)HH_{}(A)`$ la trace de Dennis et $`\gamma _{}:K_{}(A)HC_{}^{}(A)`$ le caractère de Chern universel (, II). Pour tout $`r0`$, Weibel (, p. 541) a montré qu’on a un diagramme commutatif
$`(0)_r`$
Nous montrons ici que pour tout $`r1`$ et tout $`n2`$, il existe un diagramme commutatif
$`(0)_r^n`$
reliant la $`K`$-théorie à coefficients de $`A`$ aux diverses homologies à coefficients au moyen d’applications $`\gamma _r^{(n)}`$ et $`D_r^{(n)}`$ qui sont précisées plus bas.
Notons $`C_{}(A)`$, $`CC_{}(A)^{}`$ et $`CC_{}^{per}(A)`$ les complexes de chaînes des homologies de Hochschild, cyclique négative et périodique de $`A`$. La correspondance de Dold-Kan (, 8.4) $`DK:Ch_+AbS`$ réalise une équivalence entre la catégorie $`Ch_+`$ des complexes de chaînes gradués positivement et la catégorie des groupes abéliens simpliciaux. Notons $`:AbSTop`$ le foncteur réalisation géométrique et $`X:Ch_{}Top`$ la composée de ces deux foncteurs. Les espaces $`(A)=X(C_{}(A))`$, $`^{}(A)=X(CC_{}^{}(A))`$ et $`^{per}(A)=X(CC_{}^{per}(A))`$ sont des espaces classifiants pour les homologies de Hochschild ou cycliques. Pour $`r0`$, on a $`\pi _r((A))=HH_r(A)`$, $`\pi _r(^{}(A))=HC_{}^{}(A)`$ et $`\pi _r(^{per}(A))=HC_{}^{per}(A)`$. D’après , prop. 4.3, on a
###### Proposition 24
Pour toute algèbre $`A`$, il existe dans la catégorie $`Top`$ un diagramme commutatif
$`(1)`$
tel que pour $`r0`$, en appliquant le foncteur $`\pi _r()`$ au diagramme $`(1)`$, on obtienne le diagramme $`(0)_r`$.
Il est naturel de chercher à obtenir le diagramme $`(0)_r^n`$ par application du foncteur $`\pi _r(;𝐙/n)`$ au diagramme $`(1)`$. Ce foncteur n’est défini que pour $`r2`$. Une désuspension permet l’étude du cas $`r=1`$.
Rappelons (voir ), que pour $`r2`$, on pose $`\pi _r(;𝐙/n)=[M_n^r,]`$$`M_n^r`$ est l’espace de Moore $`S^{r1}_\alpha e^r`$ avec $`\alpha `$ de degré $`n`$.
Pour $`r2`$, on pose $`K_r(A;𝐙/n)=\pi _r(BGLA^+;𝐙/n)`$. La structure de $`H`$-espace de $`BGLA^+`$ permet de définir une application $`.n:BGLA^+BGLA^+`$ induisant la multiplication par $`n`$ sur $`\pi _r(BGLA^+)`$. La fibre homotopique $``$ de $`.n:BGLA^+BGLA^+`$ est telle que pour $`r1`$, on a $`\pi _r()=\pi _{r+1}(BGLA^+;𝐙/n)`$.
Rappelons également que l’homologie à coefficients $`H_{}(C_{};𝐙/n)`$ d’un complexe de chaînes $`C_{}`$ (définie en 1.3 comme l’homologie du cône $`co(.n)`$ de la multiplication $`.n:C_{}C_{}`$) est telle que pour $`r2`$, on ait $`H_r(C_{};𝐙/n)=\pi _r(X(C_{});𝐙/n)`$.
En conclusion, on a
###### Proposition 25
Pour $`r2`$, le diagramme $`(0)_r^n`$ est obtenu par application du foncteur $`\pi _r(;𝐙/n)`$ au diagramme $`(1)`$.
Pour $`r=1`$, contentons-nous de traiter le cas de l’homologie cyclique négative en détaillant la construction de l’application $`\gamma _1^{(n)}`$ du diagramme $`(0)_1^n`$. Les constructions pour l’homologie de Hochschild ou périodique sont analogues.
Soit $`SA`$ le cône de l’anneau $`A`$ au sens de , p. 269. D’après , prop. 3.2, pour $`r1`$, on a $`K_r(SA)K_{r1}(A)`$. En particulier,$`\pi _2(BGL(SA)^+)K_1(A)`$. D’après , thm. 10.1 et sa remarque 2, pour $`r1`$, on a $`HC_r^{}(SA)HC_{r1}^{}(A)`$. En particulier, $`\pi _2(^{}(SA))HC_1^{}(A)`$.
L’application $`\gamma ^{}:BGL(SA)^+^{}(SA)`$ obtenue en appliquant la proposition 25 à l’anneau $`SA`$ permet de définir $`\gamma _1^{(n)}=\pi _1(\gamma ^{};𝐙/n)`$,$`K_1(A;𝐙/n)=\pi _2(BGL(SA)^+;𝐙/n)`$ et $`HC_1^{}(A;𝐙/n))=\pi _2(^{}(SA);𝐙/n)`$. On a $`\gamma _1^{(n)}:K_1(A;𝐙/n)HC_1^{}(A;𝐙/n).`$ Les suites exactes longues de Barratt pour l’homotopie à coefficients (, p. 3) nous donnent
###### Proposition 26
On a le diagramme commutatif naturel à lignes exactes
Remarque 27. Il est raisonnable de conjecturer que l’application composée
est celle décrite en 1.5.
## 2 Étude de l’anneau des entiers d’un corps de nombres.
Soit $`A`$ l’anneau des entiers d’un corps de nombres $`F`$. D’après , on a $`K_1(A)=A^\times `$. L’extension $`()`$ 1.1 s’écrit donc
$`()`$
ce qui montre que $`K_1(A;𝐙/n)`$ est une extension de la $`n`$-torsion du groupe des classes de $`A`$ par un quotient du groupe des unités de $`A`$.
La trace de Dennis $`D_1^{(n)}`$ et la classe caractéristique secondaire $`d_1^{(n)}`$ construites en 1.6 s’insèrent dans le diagramme commutatif
En 2.1, on montre que $`K_1(A;𝐙/n)`$ est isomorphe à un groupe noté$`𝒰(A;𝐙/n)`$ construit à partir d’idéaux fractionnaires de $`A`$. Le lemme ($`N`$-$`N_1`$) de la section 2.2 fournit un critère de construction d’éléments de $`𝒰(A;𝐙/n)`$. En 2.3, on décrit le lien entre $`K_1(A;𝐙/n)`$ et le groupe des adèles restreints $`\widehat{A}`$ de $`A`$. Les traces $`D_1^{(n)}`$ et $`d_1^{(n)}`$ sont détaillées en 2.5.
### 2.1 Description de $`K_1(A;𝐙/n)`$ en termes d’idéaux.
Soit $`A`$ un anneau de Dedekind, de corps des fractions $`F`$ et soit $`n2`$ un entier. On désigne par $`I(A)`$ le monoïde des idéaux fractionnaires de $`A`$. Pour $`xF`$, on pose $`[x]=x\mathrm{mod}F^{\times (n)}`$. Considérons le sous-groupe $`𝒰(A;𝐙/n)`$ de $`F^\times /(n)`$ défini par
$$𝒰(A;𝐙/n):=\{xF^\times /(n)II(A),xA=I^n\}$$
###### Théorème 28
Soit $`A`$ un anneau de Dedekind. Alors, on a un isomorphisme
$$K_1(A;𝐙/n)𝒰(A;𝐙/n)SK_1(A)/(n).$$
Preuve : on sait que dans un anneau de Dedekind, tout module $`P`$ projectif de type fini et de rang $`r`$ est de la forme $`P=(r1)AI`$$`I`$ est un idéal fractionnaire de $`A`$. Soit $`[P,\alpha ,L]`$ un élément de $`K_1(A;𝐙/n)`$ avec $`P=(r1)AI`$ et $`L=rA`$. L’isomorphisme $`\alpha :nPnL`$ montre qu’il existe $`xF^\times `$ tel que $`I^n=xA`$. Définissons $`det^{(n)}:K_1(A;𝐙/n)𝒰(A;𝐙/n)`$ par $`det^{(n)}([P,\alpha ,L])=[x].`$ On vérifie que $`det^{(n)}`$ est bien définie et que c’est un morphisme de groupes. Pour obtenir le facteur direct $`𝒰(A;𝐙/n)`$ de $`K_1(A;𝐙/n)`$, on construit un morphisme $`s:𝒰(A;𝐙/n)K_1(A;𝐙/n)`$ tel que $`det^{(n)}s=\mathrm{id}_{𝒰(A;𝐙/n)}`$. Soit $`[x]𝒰(A;𝐙/n)`$ avec $`xA=I^n`$$`I`$ est un idéal fractionnaire de $`A`$. On pose $`s([x])=[A,\mathrm{id}_{(n1)A}x,I]`$. On vérifie que $`s`$ est bien définie. Pour montrer $`s([x][y])=s([x])+s([y])`$, posons $`xA=I^n`$, $`yA=J^n`$, $`s([x])=u`$, $`s([y])=v`$, $`s([x][y])=w`$, avec $`u=[A,f,I]`$, $`f=\mathrm{id}_{(n1)A}x`$, $`v=[A,g,J]`$, $`g=y\mathrm{id}_{(n1)A}`$, $`w=[A,h,IJ]`$, $`h=xy\mathrm{id}_{(n1)A}`$. On a $`u+v=[AA,h_1,IJ]`$$`h_1`$ est l’isomorphisme $`\mathrm{id}_{(n1)A}xy\mathrm{id}_{(n1)A}`$. Posons $`h_2=\mathrm{id}_{nA}h`$. Dans la catégorie $`𝒞`$, (cf. 1.1), les objets $`(AA,h_1,IJ)`$ et $`(AA,\mathrm{id}_{nA}h,AIJ)`$ sont isomorphes. Dans $`K_1(A;𝐙/n)`$, on a par conséquent $`u+v=[AA,h_1,IJ]=[A,\mathrm{id}_{nA},A]+[A,h,IJ]=0+w=w`$, ce qui montre que $`s`$ est un morphisme de groupes. On vérifie la relation $`det^{(n)}s=\mathrm{id}_{𝒰(A;𝐙/n)}`$, ce qui montre
$`K_1(A;𝐙/n)=𝒰(A;𝐙/n)\mathrm{ker}det^{(n)}.`$
Introduisons le diagramme
où les applications $`\rho ^{}`$ et $`^{}`$ sont définies par $`\rho ^{}(a)=[a]`$, $`^{}([x])=[I^1].`$ Ce diagramme est commutatif à lignes exactes. La suite exacte des noyaux des flèches verticales conduit à $`\mathrm{ker}\text{det}^{(n)}SK_1(A)/(n).`$
###### Corollaire 29
Soit $`A`$ l’anneau des entiers d’un corps de nombres $`F`$. Alors
$$K_1(A;𝐙/n)𝒰(A;𝐙/n).$$
En effet, d’après \[BMS\], on a $`SK_1(A)=0`$.
### 2.2 Le lemme ($`N`$-$`N_1`$) de construction d’éléments de $`K_1(A;𝐙/n)`$.
Soit $`L/F`$ une extension de corps de nombres de degré $`\mathrm{}`$. On pose $`A=𝒪_F`$ et on note $`B`$ la fermeture intégale de $`A`$ dans $`L`$.
Pour $`zL`$, on désigne par $`\mu _z:LL`$ la multiplication par $`z`$. Les quantités $`N_j(z)L`$ sont définies par $`det(X\mathrm{id}_L\mu _z)=_{j=0}^{\mathrm{}}(1)^\mathrm{}jN_j(z)X^j`$. En particulier,
$$N_{\mathrm{}}(z)=1,N_\mathrm{}1(z)=\mathrm{tr}_{L/F}(z),N_0(z)=N_{L/F}(z)=N(z).$$
###### Proposition 30
Soit $`zL`$ (resp. $`B`$) et $`hF`$ (resp. $`A`$). Alors on a
$$N(z+h)=N(z)+N_1(z)h+h^2\epsilon (z,h)\text{ avec }\epsilon (z,h)F(\text{ resp. }A).$$
Remarquons qu’on a la formule commode
$$N_1(z)=\left(\left(\frac{d}{dh}\right)_{hF}N(z+h)\right)_{h=0}.$$
###### Lemme 31
Soit $`L/F`$ une extension de corps de nombres et soient $`A`$ l’anneau des entiers de $`F`$ et $`B`$ la fermeture intégrale de $`A`$ dans $`L`$. Soit $`n`$ un entier $`2`$. Considérons un élément $`u`$ de $`B`$ tel que
$`N(u)=ea^n\text{ avec }eA^\times ,aA`$
et $`(N(u),N_1(u))=A.`$
En désignant par $`[u]`$ la classe de $`u`$ dans $`L^\times /(n)`$ , on a alors
$$[u]K_1(B;𝐙/n).$$
Preuve : L’hypothèse $`(N(u),N_1(u))=A`$ signifie que $`u`$ est premier à tous ses conjugués. En effet, si $`\sigma :L𝐂`$ désigne un $`F`$-plongement de $`L`$ (c’est-à-dire $`\sigma _F=\mathrm{id}`$), on a $`N(u)=_\sigma \sigma (u)`$, $`N_1(u)=\left(_\sigma \sigma (u)\right)\left(_\sigma \sigma (u)^1\right)`$ et $`P_u(X)=_\sigma (X\sigma (u))`$. Soit $`𝔭`$ un idéal premier de $`A`$. On a $`𝔭(N(u),N_1(u))`$ si et seulement si $`P_u(X)X^2Q(X)\mathrm{mod}𝔭`$. Ceci signifie qu’il existe deux plongements $`\sigma _1`$ et $`\sigma _2`$ tels que $`(\sigma _1(u),\sigma _2(u))𝔭`$. En posant $`\tau =\sigma _1^1\sigma _2`$ et $`𝔮=\sigma _1^1(𝔭)`$, on en déduit $`(u,\tau (u))𝔮`$.
Pour montrer que $`[u]`$ appartient à $`K_1(A;𝐙/n)`$, on remarque que la puissance $`n`$-ième de l’idéal fractionnaire $`I=(u,a)`$ de $`B`$ est principale. En effet $`I^n=(u^n,N(u))=(u^n,u_{\sigma \mathrm{id}}\sigma (u))`$; et puisque $`(u,\sigma (u))=B`$, on en déduit $`I^n=uB`$. Ceci montre $`[u]𝒰(B;𝐙/n)`$. $`\mathrm{}`$
### 2.3 Description de $`K_1(A;𝐙/n)`$ en termes d’adèles.
Soit $`A`$ l’anneau des entiers d’un corps de nombres $`F`$. Notons $`Spec(A)`$ le spectre premier de $`A`$. Pour $`𝔭`$ dans $`Spec(A)`$, on désigne par $`\widehat{A}_𝔭`$ le complété $`𝔭`$-adique de l’anneau de valuation discrète $`A_𝔭`$. L’anneau $`\widehat{A}=_{𝔭Spec(A)}\widehat{A}_𝔭`$, appelé anneau des adèles restreints de $`A`$, s’insère dans le diagramme commutatif
où on a posé $`\widehat{F}=F_A\widehat{A}`$.
Considérons le diagramme
Les applications $`ı`$, $`j`$ et $`\overline{j}`$ sont trivialement injectives. Dans l’anneau $`\widehat{A}`$, on s’est restreint aux places archimédiennes. D’après , Chap. X.I, dans cette situation, l’application $`\overline{ı}`$ est également injective. La somme amalgamée $`\widehat{A}^\times /(n)_{\widehat{F}^\times /(n)}F^\times /(n)`$ est donc égale à $`\widehat{A}^\times /(n)F^\times /(n)`$. Soit $`xF^\times /(n)`$. L’élément $`[x]=x\mathrm{mod}F^{\times (n)}`$ de $`F^\times /(n)`$ appartient à $`\widehat{A}^\times /(n)`$ si et seulement si $`nv_𝔭(x_𝔭)`$ pour tout $`𝔭Spec(A)`$, c’est-à-dire qu’on a $`xA=I^n`$ avec $`I`$ idéal fractionnaire. En conclusion, nous avons:
###### Proposition 32
Soit $`A`$ l’anneau des entiers d’un corps de nombres $`F`$ et soit $`n2`$ un entier. Alors le groupe $`K_1(A;𝐙/n)`$ s’identifie au sous-groupe $`\widehat{A}^\times /(n)F^\times /(n)`$ de $`\widehat{F}^\times /(n)`$.
###### Corollaire 33
Soit $`A`$ l’anneau des entiers d’un corps de nombres $`F`$. Posons $`\widehat{A}=_𝔭\widehat{A}_𝔭`$. Alors l’application $`K_1(A;𝐙n)K_1(\widehat{A};𝐙/n)`$ induite par $`A\widehat{A}`$ est l’inclusion $`\widehat{A}^\times /(n)F^\times /(n)\widehat{A}^\times /(n)`$.
Preuve : Calculons $`K_1(\widehat{A};𝐙/n)`$. Rappelons pour cela que si $`(A_i)_{iI}`$ est une famille d’anneaux commutatifs de rang stable $`d2`$ au sens de , p. 231, on a $`K_1(_{iI}A_i)_{iI}K_1(A_i)`$ et $`\stackrel{~}{K}_0(_{iI}A_i)_{iI}\stackrel{~}{K}_0(A_i).`$ Les anneaux de la famille $`(\widehat{A}_𝔭)_{𝔭Spec(A)}`$ sont tous de rang stable $`d=2`$. De $`\stackrel{~}{K}_0(\widehat{A}_𝔭)_{(n)}=0`$ et $`K_1(\widehat{A}_𝔭)=\widehat{A}_𝔭^\times `$, on tire $`K_1(\widehat{A})=\widehat{A}^\times `$. L’extension ($``$ 1.1) nous mène à
$$K_1(\widehat{A};𝐙/p)=\widehat{A}^\times /(n).$$
### 2.4 Description de $`\mathrm{\Omega }_{dR}^1(A)/(n).`$
On désigne toujours par $`A`$ l’anneau des entiers d’un corps de nombres $`F`$. L’anneau $`A𝐙_p`$ est toujours considéré comme une algèbre sur l’anneau $`𝐙_p`$ des entiers $`p`$-adiques. Désignons par $`\delta `$ le discriminant du corps $`F`$. D’après , prop. 1.5, on a les égalités suivantes
$$\mathrm{\Omega }_{dR}^1(A)=\mathrm{\Omega }_{dR}^1(\widehat{A})=_{p\delta }\mathrm{\Omega }_{dR}^1(A𝐙_p).$$
Supposons $`n`$ et $`p`$ premiers entre eux, alors $`n`$ appartient à $`(𝐙_p)^\times `$ et $`\mathrm{\Omega }_{dR}^1(A𝐙_p)/(n)=0`$. On en déduit:
###### Proposition 34
Soit $`F`$ un corps de nombres d’anneaux d’entiers $`A`$ et de discriminant $`\delta `$. Soit $`n>1`$ un entier.
a) Si $`n\delta `$, alors $`\mathrm{\Omega }_{dR}^1(\widehat{A})/(n)=\mathrm{\Omega }_{dR}^1(A)/(n)=_{p(n,\delta )}\mathrm{\Omega }_{dR}^1(A𝐙_p)/(n)`$.
En particulier, si $`p`$ est un nombre premier ramifié dans $`A`$, on a
$$\mathrm{\Omega }_{dR}^1(A)/(p)\mathrm{\Omega }_{dR}^1(A𝐙_p)/(p).$$
b) si $`(n,\delta )=1`$, alors $`\mathrm{\Omega }_{dR}^1(\widehat{A})/(n)=\mathrm{\Omega }_{dR}^1(A)/(n)=0.`$
En particulier, si $`p`$ est un nombre premier non ramifié dans $`A`$, alors on a $`\mathrm{\Omega }_{dR}^1(A)/(p)=0.`$
### 2.5 Description de la trace de Dennis à coefficients
Soit $`A`$ l’anneau des entiers d’un corps de nombres $`F`$ et soit $`n`$ un diviseur du discriminant de $`F`$. Pour les éléments de $`K_1(A;𝐙/n)`$ obtenus grâce au lemme ($`N`$-$`N_1`$), il est facile de décrire la trace de Dennis à coefficients. Supposons que $`uA`$ satisfasse aux hypothèses du lemme ($`N`$-$`N_1`$) et que de plus $`N(u)`$ soit un entier premier à $`n`$. Posons $`v=_{\sigma \mathrm{id}}\sigma (u)`$. On a $`vA`$ et $`uv=N(u)`$. Dans $`\mathrm{\Omega }_{dR}^1(A)/(n)`$, on en déduit $`D_1^{(p)}([u])=N(u)^1vdu\mathrm{mod}n`$.
L’égalité $`K_1(A;𝐙/n)F^\times /(n)\widehat{A}^\times /(n)`$ permet également de décrire localement la trace de Dennis à coefficients. Pour cela, on remarque que le diagramme suivant est commutatif.
Pour connaître la trace de Dennis à coefficients de $`A`$, il suffit donc de connaître celle de $`\widehat{A}`$. Soit $`p`$ un nombre premier ramifié dans $`A`$. Localement, la trace de Dennis à coefficients est essentiellement une dérivée logarithmique modulo $`p`$. En effet, on écrit
$$K_1(\widehat{A};𝐙/p)=\underset{𝔭𝐙(p)}{}\widehat{A}_𝔭^\times /(p)\underset{𝔭𝐙=(p)}{}\widehat{A}_𝔭^\times /(p).$$
La restriction de $`D_1^{(p)}`$ au premier facteur de cette décomposition est évidemment nulle puisque pour $`𝔭𝐙(p)`$, on a $`\mathrm{\Omega }_{dR}^1(\widehat{A}_𝔭)/(p)=0`$. Pour $`𝔭𝐙=(p)`$ et $`[u_𝔭]\widehat{A}_𝔭^\times /(p)`$, d’après l’exemple du théorème 23, on a
$$D_1^{(p)}([u_𝔭])=u_𝔭^1du_𝔭\mathrm{mod}p\mathrm{\Omega }_{dR}^1(\widehat{A}_𝔭).$$
## 3 Applications aux corps de petit degré.
### 3.1 Un théorème de Y. Yamamoto.
L’égalité $`K_1(A;𝐙/n)=𝒰(A;𝐙/n)`$ et le lemme ($`N`$,$`N_1`$) permettent de retrouver un théorème montré par Y. Yamamoto à l’aide de méthodes distinctes.
###### Théorème 35
Soit $`F`$ un corps de nombres quadratique d’anneau d’entiers $`A`$, de discriminant $`\delta `$ et soit $`n`$ un entier impair. On suppose qu’il existe deux couples $`(\alpha ,b)`$ et $`(\alpha ^{},b^{})`$ dans $`𝐙^2`$ satisfaisant aux relations
$$\alpha ^24b^n=\alpha _{}^{}{}_{}{}^{2}4b_{}^{}{}_{}{}^{n}=\delta $$
avec $`(\alpha ,b)=(\alpha ^{},b^{})=1`$. On suppose de plus que pour tout diviseur premier $`p`$ de $`n`$, les conditions ci-dessous sont satisfaites.
* $`\alpha `$ (resp. $`\alpha ^{}`$) n’est pas une puissance $`p`$-ième modulo $`b`$ (resp. $`b^{}`$);
* $`(\alpha +\alpha ^{})/2`$ est une puissance $`p`$-ième modulo $`b`$ et modulo $`b^{}`$.
Alors:
Si $`\delta <4`$ le groupe des classes de $`A`$ contient un sous groupe isomorphe à $`𝐙/n𝐙/n`$.
Si $`\delta >0`$, le groupe des classes de $`A`$ contient un sous groupe isomorphe à $`𝐙/n`$.
Preuve : L’application $`f:A𝐙/b`$ définie par $`f((x+y\sqrt{\delta })/2)=(x+y\alpha )/2`$ est un morphisme d’anneaux. On note $`f_1^{(d))}:K_1(A;𝐙/d)K_1(𝐙/b;𝐙/d)`$ l’application induite par $`f`$ en $`K`$-théorie à coefficients $`d`$. Remarquons que $`K_1(𝐙/b;𝐙/d)=(𝐙/b)^\times /(d)`$. L’élément $`u=(\alpha +\sqrt{\delta })/2`$ de $`A`$ est de norme $`N(u)=b^n`$, de trace $`\text{tr}(u)=N_1(u)=\alpha `$. D’après le lemme ($`N`$-$`N_1`$), pour tout diviseur $`d`$ de $`n`$, l’élément $`[u]F^\times /(d)`$ appartient à $`K_1(A;𝐙/d)`$. Soit $`p`$ un diviseur premier de $`n`$. De $`f(u)=\alpha `$, on déduit $`f_1^{(p)}([u])=[\alpha ]`$, quantité distincte de $`1`$ d’après l’hypothèse a). Pour tout diviseur premier $`p`$ de $`n`$, l’élément $`[u]`$ de $`K_1(A;𝐙/p)`$ n’est donc pas trivial. Montrons que $`[u]F^\times /(n)`$ définit un élément d’ordre $`n`$ de $`K_1(A;𝐙/n)`$. Supposons $`[u]`$ d’ordre $`m`$ avec $`1m<n`$. Il existe un nombre premier $`p`$ tel que $`mpn`$. De $`[u]^{n/p}=1`$, on tire $`u^{n/p}=z^n`$ avec $`zA^\times `$, soit encore $`uF^{\times (p)}`$ et donc $`[u]`$ trivial dans $`K_1(A;𝐙/p)`$. On vient de montrer que ceci est impossible. On a donc $`[u]`$ d’ordre $`n`$ dans $`K_1(A;𝐙/n)`$. Le sous-groupe $`H`$ de $`K_1(A;𝐙/n)`$ engendré par $`[u]`$ est donc isomorphe à $`𝐙/n`$.
On introduit de manière analogue $`u^{}=(\alpha ^{}+\sqrt{\delta })/2`$ et on obtient de même un sous-groupe $`H^{}`$ de $`K_1(A;𝐙/n)`$, également isomorphe à $`𝐙/n`$. Pour montrer la somme directe $`HH^{}`$dans $`K_1(A;𝐙/n)`$, on remarque que $`f(u^{})=(\alpha ^{}+\alpha )/2𝐙/b`$. Si $`u^{}H`$, c’est-à-dire $`u^{}=u^m`$ avec $`1m<n`$, l’égalité $`f_1^{(p)}([u^{}])=f_1^{(p)}([u])^m`$, satisfaite pour tout diviseur premier $`p`$ de $`n`$, s’écrit encore $`[(\alpha +\alpha ^{})/2]=[\alpha ]^m`$, ce qui donne $`[\alpha ]^m=1`$ d’après l’hypothèse b). On en déduit comme ci-dessus qu’il existe un nombre premier $`p`$ tel que $`mpn`$ pour lequel $`\alpha (𝐙/b)^{\times (p)}`$, ce qui fournit la contradiction recherchée. On montre de même $`uH^{}`$. En conclusion, sous les hypothèses proposées, le groupe $`K_1(A;𝐙/n)`$ contient un sous-groupe isomorphe à $`𝐙/n𝐙/n`$. De l’extension $`()`$ p. 20, on déduit que si $`F`$ est imaginaire, $`Cl(A)_{(n)}`$ contient $`𝐙/n𝐙/n`$ en facteur direct. Si $`\delta >0`$, $`A^\times /(n)`$ est isomorphe à $`𝐙/n`$ et donc $`Cl(A)_n`$ contient $`𝐙/n`$ en facteur direct.
Remarque 36. Dans le cas où $`F`$ est réel d’unité fondamentale $`\epsilon `$ telle qu’il existe un diviseur premier $`p`$ pour lequel $`f(\epsilon )(𝐙/b)^{\times (p)}`$, le groupe des classes contient un sous-groupe isomorphe à $`𝐙/n𝐙/n`$. En effet, soit $`t=([u])`$ toujours avec $`u=(\alpha +\sqrt{\delta })/2`$. Montrons que $`t`$ est d’ordre $`n`$ dans $`Cl(A)`$. Supposons $`t^m=0`$ avec $`1m<n`$. on en déduit $`[u]^m\text{ker }`$, soit $`u^m=\epsilon ^l`$, ce qui donne $`[\alpha ]^m=f_1^{(p)}([u]^m)=f_1^{(p)}(\epsilon )^l=1`$ puisque $`f(\epsilon )(𝐙/b)^{\times (p)}`$. on en déduit $`\alpha (𝐙/b)^{\times (p)}`$, situation exclue. Le sous-groupe $`H(t)`$ engendré par $`t`$ dans $`Cl(A)`$ est donc isomorphe à $`𝐙/n`$. La fin de la démonstration est analogue à celle du théorème. Les sous-groupes $`H(t)`$ et $`H(t^{})`$ engendrés respectivement par $`t=([(\alpha +\sqrt{\delta })/2])`$ et $`t^{}=([(\alpha ^{}+\sqrt{\delta })/2])`$ sont en somme directe dans $`Cl(A)`$.
Ces éléments du groupe des classes ont été construits pour la première fois par Yamamoto (). À partir de ces éléments, cet auteur a montré que pour tout $`n>1`$, il existe une infinité de corps quadratiques réels et imaginaires dont le groupe des classes contient un facteur $`𝐙/n`$. $`\mathrm{}`$
### 3.2 Construction d’éléments non triviaux de $`Cl(A)_{(n)}`$.
Soit $`F`$ un corps de nombres d’anneaux d’entiers $`A`$. Soit $`r_1`$ (resp. $`2r_2`$) le nombre de plongements réels (resp. complexes) de $`F`$. Si $`r=r_1+r_21`$, on a $`rg(A^\times )=r`$ et $`A^\times =K_1(A)=\mu \times _{i=1}^r𝐙\epsilon _i`$$`\mu `$ est le groupe des racines de l’unité contenues dans $`A`$ et où $`\{\epsilon _i\mathrm{,\; 1}ir\}`$ est un système fondamental d’unités de $`A`$.
Lorsque $`A`$ possède “peu” d’unités, l’extension $`()`$ p. 20 permet d’obtenir des élements non triviaux de $`Cl(A)_{(n)}`$ à partir d’éléments de $`K_1(A;𝐙/n).`$
L’égalité $`K_1(A;𝐙/n)=𝒰(A;𝐙/n)`$ et le lemme ($`N`$-$`N_1`$) conduisent au résultat suivant:
###### Proposition 37
Soit $`F`$ un corps de nombres, d’anneau d’entiers $`A`$ et soit $`n`$ un entier naturel. On suppose qu’il existe $`zA`$ tel que $`N(z)=b^n`$, $`(N(z),N_1(z))=1`$ et que pour tout diviseur $`m`$ de $`n`$, $`1m<n`$, $`b^m`$ ne soit pas la norme d’un élément de $`F`$. Alors $`[z]=z\mathrm{mod}F^{\times (n)}`$ est un élément d’ordre $`n`$ de $`K_1(A;𝐙/n)`$.
Si $`r_1+r_21=0`$, on a $`Cl(A)_{(n)}0`$. Si $`r_1+r_210`$ et si $`b`$ n’est pas la norme d’un élément de $`F`$, alors $`Cl(A)_{(n)}0`$.
Preuve : d’après le lemme ($`N`$-$`N_1`$), $`[z]`$ appartient à $`K_1(A;𝐙/n)`$. Si $`[z]`$ est d’ordre $`m`$, $`1m<n`$, il existe $`uF^\times `$ tel que $`z^m=u^n`$, c’est-à-dire $`z=u^s`$ avec $`s=n/m`$. L’équation $`N(z)=N(u)^s`$ s’écrit $`b^m=N(u)`$, ce qui n’est pas. Notons comme toujours $`:K_1(A;𝐙/p)Cl(A)_{(n)}`$ et supposons à présent que $`([z])=0`$. Dans ce cas, il existe $`uF^\times `$, $`\xi \mu `$ et des entiers $`l_i`$ tels que $`z=\xi ^{l_0}\epsilon _1^{l_1}\mathrm{}\epsilon _r^{l_r}u^n`$, d’où l’on déduit $`N(z)=\pm N(u)^n`$ (avec le signe $`+`$ si $`r=0`$) soit $`b=\pm N(u)`$, ce qui n’est pas.
$`\mathrm{}`$
La classe caractéristique secondaire
$$d_1^{(n)}:Cl(A)_{(n)}\mathrm{\Omega }_{dR}^1(A)/S$$
introduite en 1.6, corollaire 23 conduit au résultat suivant.
###### Proposition 38
Soit $`F`$ un corps de nombres d’anneaux d’entiers $`A`$. On pose $`r=r_1+r_21`$ et $`A^\times =\mu \times _{i=1}^r𝐙\epsilon _i.`$ Soit $`n`$ un diviseur du discriminant du corps $`F`$.
On suppose
1) Pour tout $`\xi \mu `$, $`\xi ^1d\xi 0\mathrm{mod}n\mathrm{\Omega }_{dR}^1(A).`$
2) Pour tout $`i`$, $`1ir`$, $`\epsilon _i^1d\epsilon _i0\mathrm{mod}n\mathrm{\Omega }_{dR}^1(A).`$
3) Il existe $`uA`$ avec $`N(u)=b^n`$, $`(N(u),N_1(u))=1`$, et
$`u^1du0\mathrm{mod}n\mathrm{\Omega }_{dR}^1(A).`$ Alors $`Cl(A)_{(n)}0`$.
Preuve : les hypothèses 1 et 2 montrent que $`dA^\times /A^\times 0\mathrm{mod}n\mathrm{\Omega }_{dR}^1(A)`$. L’application $`d_1^{(n)}`$, de source $`Cl(A)_{(n)}`$ est donc de but $`\mathrm{\Omega }_{dR}^1(A)/(n)`$. D’après le lemme ($`N`$-$`N_1`$), les hypothèses 3 fournissent l’élément $`[u]=u\mathrm{mod}F^{\times (n)}`$ de $`K_1(A;𝐙/n)`$. De cet élément, on déduit $`x=([u])`$ dans $`Cl(A)_{(n)}`$. On a $`d_1^{(n)}(x)=u^1du\mathrm{mod}n\mathrm{\Omega }_{dR}^1(A)`$, quantité non nulle par hypothèse, ce qui montre que $`x`$ est non trivial. $`\mathrm{}`$
Exemple. Posons $`x={}_{}{}^{3}\sqrt{182}`$ et soit $`F=𝐐[x]`$, d’anneau d’entiers $`A=𝐙[x]`$, d’unité fondamentale $`\epsilon =173x`$. Pour $`p=3`$, l’élément $`u=52x`$ définit un élément non nul de $`Cl(A)_{(p)}`$.
### 3.3 Exemples de $`n`$-torsion du groupe des classes : cas d’un corps quadratique imaginaire.
###### Proposition 39
Soit $`F`$ un corps quadratique d’anneau d’entiers $`A`$ et de discriminant $`\delta <0`$ et soit $`n`$ un entier impair. On suppose qu’il existe $`(\alpha ,b)𝐙^2`$ tel que $`\alpha ^24b^n=\delta `$, avec $`(\alpha ,b)=1`$. On suppose de plus que pour tout diviseur $`m`$ de $`n`$, $`1m<n`$, la quantité $`\delta +4b^n`$ n’est pas un carré parfait. Alors $`Cl(A)_{(n)}0.`$
Remarque 40. Sous les mêmes hypothèses, si $`\delta `$ est positif et si de plus $`\pm b`$ n’est pas une norme, la conclusion $`Cl(A)_{(n)}0`$ subsiste.
Preuve : on applique la proposition 37 à $`z={\displaystyle \frac{\alpha +\sqrt{\delta }}{2}}.`$ L’équation $`z^m=u^n`$ conduit à $`\delta +4b^n`$ carré parfait.
Dans les quelques exemples ci-dessous, l’anneau $`A`$ des entiers du corps $`𝐐[\sqrt{\delta }]`$ est tel que $`Cl(A)_{(n)}0`$.
$`n=3`$
$`\delta =104=2^243^3=4(26)`$
$`\delta =\mathrm{5\; 320}=2^2411^3=4(1330)`$
$`\delta =\mathrm{48\; 664}=2^2423^3=4(\mathrm{12\; 166})`$
$`n=5`$
$`\delta =127=1^242^5`$
$`\delta =\mathrm{12\; 499}=1^245^5`$
$`\delta =\mathrm{31\; 103}=1^246^5`$
$`\delta =\mathrm{131\; 071}=1^248^5`$
$`\delta =\mathrm{399\; 999}=1^2410^5`$
$`n=7`$
$`\delta =511=1^242^7`$
$`\delta =\mathrm{65\; 535}=1^244^7`$
$`\delta =\mathrm{312\; 499}=1^245^7`$
$`n=9`$
$`\delta =2047=1^22^9`$
$`\delta =\mathrm{78\; 728}=2^23^9=4(\mathrm{19\; 682})`$
$`\delta =78731=1^23^9`$
$`n=11`$
$`\delta =8191=1^22^{11}`$
$`\delta =\mathrm{708\; 584}=2^23^{11}=4(\mathrm{177\; 146})`$
$`\delta =\mathrm{708\; 587}=1^23^{11}`$
### 3.4 Exemples de $`n`$-torsion ramifiée du groupe des classes: cas d’un corps quadratique
Soit $`F`$ un corps de nombres quadratique de discriminant $`\delta `$. Si $`\delta <0`$, on exclut les deux cas $`\delta =4`$ et $`\delta =3`$ pour lesquels le groupe des classes est trivial et le groupe des unités n’est pas réduit à $`𝐙/2`$.
On pose $`\omega =\frac{\sqrt{\delta }}{2}`$ ou $`\omega =\frac{1+\sqrt{\delta }}{2}`$ suivant que $`\delta 0\mathrm{mod}4`$ ou $`\delta 1\mathrm{mod}4.`$ L’anneau $`A`$ des entiers du corps $`F`$ est $`𝐙[\omega ]`$. Posons $`P=X^2\delta `$ si $`\delta 0\mathrm{mod}4`$ et $`P=X^2X+(1\delta )/4`$ sinon. L’homologie de Hochschild de $`A`$ est donnée par la
###### Proposition 41
a) Si $`\delta 1\mathrm{mod}4`$, on a $`\mathrm{\Omega }_{dR}^1(A)=𝐙/\delta d\omega `$ et $`\omega d\omega =\frac{1}{2}d\omega `$.
b) Si $`\delta 0\mathrm{mod}4`$, on a $`\mathrm{\Omega }_{dR}^1(A)=𝐙/(\delta /2)d\omega 𝐙/2\omega d\omega .`$
On en déduit
###### Proposition 42
a) Soit $`n`$ un diviseur impair du discriminant $`\delta `$ du corps quadratique $`F`$. Alors on a $`\mathrm{\Omega }_{dR}^1(A)/(n)=𝐙/nd\omega `$ avec $`\omega d\omega =\frac{1}{2}d\omega `$ si $`\delta 1\mathrm{mod}4`$ et $`\omega d\omega =0`$ si $`\delta 0\mathrm{mod}4`$.
b) on suppose $`\delta 0\mathrm{mod}4`$ et $`n`$ diviseur pair de $`\delta `$. Alors on a
$$\mathrm{\Omega }_{dR}^1(A)/(n)=𝐙/nd\omega 𝐙/2\omega d\omega .$$
Tout élément $`z`$ de $`A`$ s’écrit $`z={\displaystyle \frac{\alpha +\beta \sqrt{\delta }}{2}}`$ avec $`\alpha `$ et $`\beta `$ dans $`𝐙`$. On note $`N(z)`$ sa norme, $`N_1(z)=tr(z)`$ sa trace et $`\sigma (z)`$ son conjugué. On a $`\sigma (z)={\displaystyle \frac{\alpha \beta \sqrt{\delta }}{2}}`$, $`tr(z)=\alpha `$ et $`N(z)={\displaystyle \frac{\alpha ^2\delta \beta ^2}{4}}`$. Dans $`\mathrm{\Omega }_{dR}^1(A)`$, on a la relation $`N(z)z^1dz=\sigma (z)dz.`$
Supposons que $`n`$ soit un diviseur impair du discriminant $`\delta `$. De $`N(z)\alpha ^2/4\mathrm{mod}n𝐙`$ et de $`\sigma (z)dz\frac{\alpha \beta }{2}d\omega \mathrm{mod}n\mathrm{\Omega }_{dR}^1(A)`$, on déduit que si $`(N(z),n)=1`$, on a $`z^1dz{\displaystyle \frac{2\beta }{\alpha }}d\omega \mathrm{mod}n\mathrm{\Omega }_{dR}^1(A).`$
En particulier, si $`F`$ est réel et si $`\epsilon ={\displaystyle \frac{\epsilon _1+\epsilon _2\sqrt{\delta }}{2}}`$ est l’unité fondamentale de $`A`$, on a toujours $`(\epsilon _1,n)=1`$ et donc $`\epsilon ^1d\epsilon 0\mathrm{mod}n\mathrm{\Omega }_{dR}^1(A)`$ si et seulement si $`n`$ divise $`\epsilon _2`$. Dans ce cas, on a $`dA^\times /A^\times 0\mathrm{mod}n\mathrm{\Omega }_{dR}^1(A)`$. On en déduit
$$\mathrm{\Omega }_{dR}^1(A)/(n,dA^\times /A^\times )=\mathrm{\Omega }_{dR}^1(A)/(n)=𝐙/ndw.$$
On obtient la classe caractéristique
$$d_1^{(n)}:Cl(A)_{(n)}𝐙/ndw$$
Enfin, si $`u={\displaystyle \frac{\alpha +\beta \sqrt{\delta }}{2}}`$ est un élément de $`A`$ tel que $`(\beta ,n)=(\alpha ,n)=1`$, alors $`u^1du\frac{2\beta }{\alpha }d\omega \mathrm{mod}n`$ est une quantité non nulle de $`𝐙/nd\omega `$. De tout ceci, on déduit que la proposition 38 prend la forme:
###### Proposition 43
Soit $`F`$ un corps quadratique de discriminant $`\delta `$ et d’anneau d’entiers $`A`$. Soit $`n`$ un diviseur impair de $`\delta `$. Si $`F`$ est réel, on suppose que l’unité fondamentale $`\epsilon ={\displaystyle \frac{\epsilon _1+\epsilon _2\sqrt{\delta }}{2}}`$ est telle que $`n|\epsilon _2`$. Soient $`(\alpha ,\beta ,b))𝐙^3`$ une solution de l’équation $`\alpha ^24b^n=\delta \beta ^2`$ avec $`(b,\alpha )=(\beta ,n)=(\alpha ,n)=1.`$ Alors $`Cl(A)`$ possède un élément d’ordre $`n`$.
En se restreignant aux éléments $`u`$ de la forme $`{\displaystyle \frac{\alpha +\sqrt{\delta }}{2}}`$, on obtient
###### Proposition 44
Soient $`\alpha `$, $`b`$ et $`n`$ trois entiers avec $`n`$ impair, $`(\alpha ,b)=(\alpha ,n)=1`$. On pose $`\delta =\alpha ^24b^n`$. On suppose que $`n`$ divise $`\delta `$ et que $`\delta `$ est le discriminant d’un corps quadratique $`F`$ d’anneau d’entiers $`A`$. Si $`\delta `$ est positif, on suppose de plus que l’unité fondamentale $`\epsilon ={\displaystyle \frac{\epsilon _1+\epsilon _2\sqrt{\delta }}{2}}`$ de $`A`$ est telle que $`n|\epsilon _2`$. Alors $`Cl(A)_{(n)}0.`$
Dans les quelques exemples ci-dessous, l’anneau $`A`$ des entiers du corps $`𝐐[\sqrt{\delta }]`$ est tel que $`Cl(A)`$ possède un élément d’ordre $`n`$.
$`n=3`$
$`\delta =231=17^24(2)^3=\mathrm{𝟑}107`$, d’unité fondamentale $`\epsilon =(430+24\sqrt{\delta })/2`$.
$`\delta =231=5^244^3=\mathrm{𝟑}711`$
$`\delta =255=1^244^3=\mathrm{𝟑}517`$
$`\delta =\mathrm{16\; 383}=1^2416^3=\mathrm{𝟑}43127`$
$`\delta =\mathrm{62\; 484}=4^2425^3=4(\mathrm{𝟑}41127)`$
$`\delta =\mathrm{3\; 999\; 999}=1^24100^3=\mathrm{𝟑}23291999`$
$`n=5`$
$`\delta =\mathrm{236\; 195}=1^249^5=\mathrm{𝟓}97487`$
$`\delta =\mathrm{644\; 195}=3^2411^5=\mathrm{𝟓}19\mathrm{6\; 781}`$
$`\delta =\mathrm{9\; 904\; 380}=4^2419^2=4(3\mathbf{5}383431)`$
$`n=7`$
$`\delta =511=1^242^7=\mathrm{𝟕}23`$
$`\delta =\mathrm{65\; 527}=3^244^7=\mathrm{𝟕}112337`$
$`n=11`$
$`\delta =\mathrm{708\; 587}=1^243^{11}=\mathrm{𝟏𝟏}371741`$
$`n=15`$
$`\delta =\mathrm{4\; 294\; 967\; 295}=1^244^{15}=\mathrm{𝟑}\mathbf{5}17257\mathrm{65\; 537}`$
## 4 Applications à la cyclotomie.
### 4.1 Notations et stratégie générale.
Soit $`p`$ un nombre premier impair et soit $`\zeta =\zeta _p`$ une racine primitive $`p`$-ième de l’unité. Le corps cyclotomique $`F=𝐐[\zeta ]`$ est une extension galoisienne de degré $`p1`$ de $`𝐐`$, de groupe de Galois $`G=\left(𝐙/p\right)^\times `$. Soit $`g`$ un générateur de $`G`$. On désigne par $`s`$, $`1<sp1`$, l’entier tel que $`g\zeta =\zeta ^s`$. La conjugaison complexe $`g^{(p1)/2}`$ est notée $`\sigma `$. L’anneau $`A`$ des entiers de $`F`$ est $`𝐙[\zeta ]`$. Soient $`Cl(A)`$ le groupe des classes de $`A`$ et $`h=h_p`$ le nombre de classes de $`A`$. Le sous-corps maximal réel $`𝐐[\zeta +\zeta ^1]`$ de $`F`$ a pour nombre de classes $`h^+`$. On sait que $`h^+h`$ et que $`h^+`$ est le nombre de classes de $`A`$ invariantes par conjugaison complexe. La $`p`$-torsion du groupe des classes $`Cl(A)`$ se décompose en $`Cl(A)_{(p)}=Cl(A)_{(p)}^{}Cl(A)_{(p)}^+`$ avec $`Cl(A)_{(p)}^\pm =\mathrm{ker}(\sigma \mathrm{id})`$. On sait que si $`p^a`$ désigne le nombre d’éléments de $`Cl(A)_{(p)}^{}`$, alors $`p^a`$ divise $`h^{}=h/h^+`$.
La conjugaison complexe sur $`A`$ définit une involution toujours notée $`\sigma `$ sur $`K_1(A;𝐙/p)`$ qui s’écrit
$$K_1(A;𝐙/p)=K_1^{}(A;𝐙/p)K_1^+(A;𝐙/p)$$
avec $`K_1^\pm (A;𝐙/p)=\mathrm{ker}(\sigma \mathrm{id})`$. Par ailleurs,
$$A^\times /(p)\mu _p\times \{\pm 1\}\times \left(𝐙/p\right)^{(p3)/2}$$
$`\mu _p=\{\mathrm{exp}(2ik\pi /p)\mathrm{,\; 0}kp1\}`$. En particulier, $`\left(A^\times /(p)\right)^{}\mu _p`$. L’extension ($``$) p. 20 se scinde donc en deux parties dont la partie antisymétrique s’écrit
$`(^{})`$
Dans toute la suite de ce texte , on pose
$$d_p^{}:=dim_{𝐙/p}Cl(A)_{(p)}^{}=dim_{𝐙/p}K_1^{}(A;𝐙/p)1.$$
Rappelons qu’un nombre premier est régulier s’il ne divise pas le nombre de classes $`h_p`$: pour un nombre premier régulier, $`Cl(A)_{(p)}=Cl(A)_{(p)}^{}=0`$.
###### Proposition 45
On a $`d_p^{}=0`$ si et seulement si $`p`$ est un nombrer premier régulier.
Preuve : Si $`p`$ est régulier, l’extension $`(^{})`$ se réduit à $`K_1^{}(A;𝐙/p)=\mu _p`$. Réciproquement, si $`K_1^{}(A;𝐙/p)=\mu _p`$, alors $`h^{}=0`$ et $`p`$ ne divise pas $`h^{}`$. D’après un théorème de Kummer (, 5.6), ceci entraîne que $`p`$ ne divise pas $`h^+`$ donc $`Cl(A)_{(p)}=0`$, c’est-à-dire $`p`$ régulier. $`\mathrm{}`$
Rappelons que $`(p,a,b,c)`$ satisfont aux hypothèses du premier cas du dernier théorème de Fermat (en abrégé DTF1) si $`p`$ est un nombre premier impair et si $`a^p=b^p+c^p`$ avec $`(a,b,c)=(p,abc)=1`$ (on parle du second cas si $`p`$ divise $`abc`$).
La démarche développée dans les paragraphes qui suivent est celle-ci. L’équation $`a^p=b^p+c^p`$ permet de construire un élément $`z`$ de $`K_1^{}(A;𝐙/p)`$. La trace de Dennis à coefficients nous permet de montrer que cet élément $`z`$ n’est pas trivial. L’action du groupe de Galois fournit$`(p1)/2`$ éléments de $`K_1^{}(A;𝐙/p)`$ construits à partir de $`z`$. Grâce à la trace de Dennis, nous minorons la dimension du sous-espace vectoriel de $`K_1^{}(A;𝐙/p)`$ engendré par ces $`(p1)/2`$ éléments en termes des polynômes de Mirimanoff. On en déduit une minoration de $`d_p^{}`$.
Au vocabulaire près, le résultat suivant est bien connu.
###### Proposition 46
Soient $`p`$ un nombre premier, $`A`$ l’anneau $`𝐙[\zeta _p]`$, $`F`$ le corps $`𝐐[\zeta _p]`$. On suppose que $`(p,a,b,c)`$ satisfont les hypothèses du premier cas du dernier théorème de Fermat. Pour $`1\mathrm{}(p1)/2`$, les éléments
$$z_{\mathrm{}}=\frac{ab\zeta ^s^{\mathrm{}}}{ab\zeta ^s^{\mathrm{}}}\mathrm{mod}F^{\times (p)}$$
appartiennent alors à $`K_1^{}(A;𝐙/p)`$.
Preuve : Sous les hypothèses DTF1, les idéaux fractionnaires principaux $`(ab\zeta ^{\mathrm{}})`$, $`1\mathrm{}p1`$ sont deux à deux premiers entre eux. On en déduit que chacun de ces idéaux s’écrit sous la forme $`(ab\zeta ^{\mathrm{}})=I_{\mathrm{}}^p`$, où les $`I_{\mathrm{}}`$ sont des idéaux fractionnaires. Par conséquent, pour $`1\mathrm{}p1`$, les éléments $`ab\zeta ^{\mathrm{}}\mathrm{mod}F^{\times (p)}`$ appartiennent à $`𝒰(A;𝐙/p)`$. $`\mathrm{}`$
### 4.2 Emploi du groupe $`K_1(A/p;𝐙/p)`$.
Soit $`\phi :AA/p`$ la projection canonique. Posons $`\lambda =1\phi (\zeta )`$. Alors $`A/(p)=𝐙/p[\lambda ]`$ avec $`\lambda ^{p1}=0`$. L’anneau $`A/p`$ est local et $`\left(A/p\right)^\times =𝐙/p^\times +\lambda 𝐙/p[\lambda ]`$. On en déduit
$$K_1(A/p;𝐙/p)=\left(A/p\right)^\times /(p)=(1+\lambda 𝐙/p[\lambda ],\times ).$$
Les modules de différentielles $`\mathrm{\Omega }_{dR}^1(A)/(p)`$, $`\mathrm{\Omega }_{dR}^1(A/p)`$ et $`\mathrm{\Omega }_{dR}^1(A/p)/(p)`$ sont tous trois isomorphes à
$$𝐙/p[X]dX/(X1)^{p2}dX,$$
donc $`\mathrm{\Omega }_{dR}^1(A/p)=𝐙/p[\lambda ]d\lambda `$ avec $`\lambda ^{p1}=0`$ et $`\lambda ^{p2}d\lambda =0`$.
Par commodité, $`o(\lambda ^j)`$ désigne un élément indéterminé de $`\lambda ^{j+1}𝐙/p[\lambda ].`$ Soit $`p`$ un nombre premier impair et soient $`x`$ et $`y`$ deux éléments de $`(𝐙/p)^\times `$ tels que $`xy=1`$. Soint les éléments $`w=xy(1\lambda )`$ et $`\sigma (w)=xy(1\lambda )^1`$ de $`(A/p)^\times `$ et soit $`z^{}=z^{}(x)`$ l’élément de $`K_1^{}(A/p;𝐙/p)`$ défini par $`z^{}=w/\sigma (w)\mathrm{mod}(A/p)^{\times (p)}.`$
###### Proposition 47
Si $`p`$ est un nombre premier impair et si $`x𝐙/p\{\mathrm{0,1,1}/2\}`$, alors l’élément $`z^{}(x)`$ ci-dessus de $`K_1^{}(A/p;𝐙/p)`$ n’est pas colinéaire à l’élément $`1\lambda `$.
Preuve : Calculons les traces $`D_1^{(p)}(z^{}(x))`$ et $`D_1^{(p)}(1\lambda )`$. On a
$$D_1^{(p)}(z^{}(x))=w^1dw\sigma (w)^1d\sigma (w).$$
Puisque $`w=1+y\lambda `$, $`w^1=_{k0}(1)^ky^k\lambda ^k`$, $`dw=yd\lambda `$ et
$$w^1dw=\underset{k0}{}(1)^ky^{k+1}\lambda ^kd\lambda .$$
De $`\sigma (w)={\displaystyle \frac{1\lambda x}{1\lambda }}`$, on déduit $`\sigma (w)^1=1+_{k1}x^{k1}y\lambda ^k`$ tandis que
$$d\sigma (w)=y\underset{k1}{}k\lambda ^{k1}d\lambda $$
et par suite
$$\sigma (w)^1d\sigma (w)=yd\lambda y(y+2)\lambda d\lambda y(3+2y+xy)\lambda ^2d\lambda +o(\lambda ^2)d\lambda .$$
Ces expressions de $`w^1dw`$ et $`\sigma (w)^1d\sigma (w)`$ conduisent à
$$D_1^{(p)}(z^{}(x))=2yd\lambda +2y\lambda d\lambda +(3y+3y^2+2y^3)\lambda ^2d\lambda +o(\lambda ^2)d\lambda .$$
Par ailleurs $`D_1^{(p)}(1\lambda )=(1\lambda )^1d\lambda =_{k0}\lambda ^kd\lambda .`$ Supposons $`z^{}(x)`$ et $`1\lambda `$ colinéaires. La comparaison des coefficients en $`d\lambda `$ et en $`\lambda ^2d\lambda `$ des traces de Dennis de $`z^{}(x)`$ et $`1\lambda `$ conduit à l’égalité $`2y^3+3y^2+y=0.`$ Puisque $`y0`$, on en déduit que $`y(𝐙/p)^\times `$ est solution de l’équation $`2X^2+3X+1=0`$ dans $`𝐙/p[X]`$. Ceci conduit à $`y=1`$ ou $`y=1/2`$. Or nécessairement $`y1`$ car sinon $`x=0`$, ce qui est exclu. Par ailleurs, $`y=1/2`$ équivaut à $`x=1/2`$, situation également exclue. $`\mathrm{}`$
###### Proposition 48
On suppose que $`(p,a,b,c)`$ satisfait DTF1 avec $`p>3`$. Alors, $`d_p^{}1`$.
Preuve : Désignons par $`\phi _1:K_1(A;𝐙/p)K_1(A/p;𝐙/p)`$ l’application induite par $`\phi `$ en $`K`$-théorie à coefficients. L’élément
$$z=\frac{ab\zeta }{ab\zeta ^1}\mathrm{mod}F^{\times (n)}$$
de $`K_1^{}(A;𝐙/p)`$ est tel que
$$\phi _1(z)=\frac{xy(1\lambda )}{xy(1\lambda )^1}\mathrm{mod}(A/p)^{\times (p)}$$
avec $`x=\overline{a}/\overline{c}`$ et $`y=x1`$. On a nécessairement $`x0`$. Si $`x=1/2`$, l’élément
$$z_1=\frac{ac\zeta }{ac\zeta ^1}\mathrm{mod}F^{\times (n)}$$
est tel que $`\phi _1(z_1)=z^{}(x_1)`$ avec $`x_1=\overline{c}/\overline{b}`$. Les hypothèses DTF1 montrent que pour $`p>3`$, il est impossible d’avoir simultanément $`x=x_1=1/2`$. La proposition précédente s’applique donc pour l’un des deux éléments $`z`$ ou $`z_1`$$`\mathrm{}`$
On a remarqué plus haut que $`d_p^{}=0`$ caractérise les nombres premiers réguliers. On a donc montré:
###### Corollaire 49
(Kummer, 1847)
Soit $`p`$ un nombre premier régulier. Alors le premier cas du dernier théorème de Fermat est satisfait pour $`p`$.
### 4.3 Emploi du groupe $`K_1(R;𝐙/p)`$ .
Dans tout ce paragraphe, $`x`$ et $`y`$ désignent deux éléments de $`(𝐙/p)^\times `$ tels que $`xy=1.`$
L’action du groupe de Galois $`G=Gal(F/𝐐)`$ sur $`\mathrm{\Omega }_{dR}^1(A)/(p)=𝐙/p[\lambda ]d\lambda `$ est peu lisible car $`g\lambda =1(1\lambda )^s`$. C’est pourquoi nous introduisons les anneaux $`R^{}=𝐙[X]/(X^p1)=𝐙[t]`$ et $`R=R^{}/p`$ avec $`t=X\mathrm{mod}(X^p1)`$. Le groupe $`G`$ opère sur $`R^{}`$ par $`gt=t^s`$ (où $`g`$ est un générateur de $`G`$ et $`(s,p)=1`$). Remarquons que l’involution $`\sigma =g^{(p1)/2}`$ est telle que $`\sigma (t)=t^1`$.
On a $`R=𝐙/p[1t]`$ avec $`(1t)^p=0`$. L’anneau $`R`$ est local, $`R^\times =(𝐙/p)^\times (1t)𝐙/p[1t]`$ et
$$K_1(R;𝐙/p)=(1+(1t)𝐙/p[1t],\times ).$$
Les modules de différentielles de Kähler $`\mathrm{\Omega }_{dR}^1(R^{})`$, $`\mathrm{\Omega }_{dR}^1(R^{})/(p)`$, $`\mathrm{\Omega }_{dR}^1(R)`$ et $`\mathrm{\Omega }_{dR}^1(R)/(p)`$ sont tous quatre isomorphes à
$$𝐙/p[X]dX/(X1)^pdX.$$
L’action de $`G`$ sur $`\mathrm{\Omega }_{dR}^1(R)`$ est donnée par $`g(t^idt)=g(t)^idg(t)=st^{s(i+1)1}dt.`$ Pour $`1kp1`$, les relations
$$g(t^{s^k}t^1dt)=st^{s^{k+1}}t^1dt,\sigma (t^{s^k}t^1dt)=t^{s^k}t^1dt$$
et
$$g(t^1dt)=st^1dt,\sigma (t^1dt)=t^1dt$$
conduisent à la décomposition commode suivante.
###### Proposition 50
Posons $`f_0^{}=t^1dt`$, et pour $`1\mathrm{}(p1)/2`$,$`f_{\mathrm{}}^\pm =\left(t^s^{\mathrm{}}t^s^{\mathrm{}}\right)t^1dt.`$
On a alors
$$\mathrm{\Omega }_{dR}^1(R)=\mathrm{\Omega }_{dR}^{}(R)\mathrm{\Omega }_{dR}^+(R)$$
$`\mathrm{\Omega }_{dR}^{}(R)`$ est de dimension $`(p+1)/2`$, de base $`(f_0^{},f_1^{},\mathrm{},f_{(p1)/2}^{})`$ et où $`\mathrm{\Omega }_{dR}^+(R)`$ est de dimension $`(p1)/2`$, de base $`(f_1^+,\mathrm{},f_{(p1)/2}^+)`$.
De plus, en désignant par $`g`$ un générateur du groupe de Galois $`G=Gal(F/𝐐)`$ et en notant $`\sigma `$ l’involution $`g^{(p1)/2}`$, on a les relations $`g(f_0^{})=sf_0^{}`$, $`g(f_{\mathrm{}}^\pm )=sf_{\mathrm{}+1}^\pm `$, $`1\mathrm{}<(p1)/2`$, $`g(f_{(p1)/2}^\pm )=f_1^\pm `$
et $`\sigma (f_0^{})=f_0^{}`$, $`\sigma (f_{\mathrm{}}^\pm )=\pm f_{\mathrm{}}^\pm `$ $`(1\mathrm{}(p1)/2`$.
###### Définition 51
Pour $`x𝐙/p\{\mathrm{0,1}\}`$, on pose $`y=x1`$ et pour $`1k(p1)/2`$, on introduit les éléments
$$\alpha _k=(x/y)^{s^{k1}}+(y/x)^{s^{k1}}$$
de $`𝐙/p`$ et les éléments suivants de $`K_1(R;𝐙/p)`$: $`v_k(x)=xyt^{s^k}\mathrm{mod}R^{\times (p)}`$, $`\sigma (v_k(x))=xyt^{s^k}\mathrm{mod}R^{\times (p)}`$ et
$$z_k(x)=v_k(x)/\sigma (v_k(x)).$$
###### Proposition 52
Dans la base $`(f_0^{},\mathrm{},f_{(p1)/2}^{})`$ de $`\mathrm{\Omega }_{dR}^{}(R)`$, la trace de Dennis de $`z_1(x)`$ s’écrit
$$D_1^{(p)}(z_1(x))=s(x1)\left(2f_0^{}+\underset{k=1}{\overset{(p1)/2}{}}\alpha _kf_k^{}\right).$$
Preuve : On a $`D_1^{(p)}(z_1(x))=v_1^1(x)dv_1(x)\sigma (v_1(x))^1d\sigma (v_1(x))`$. Pour calculer $`v_1^1(x)`$, écrivons $`v_1(x)=yt^s(1(x/y)t^s)`$. L’identité
$$\left(1(x/y)t^s\right)\left(1+(x/y)t^s+\mathrm{}+(x/y)^{p1}t^{(p1)s}\right)=1x/y=1/y$$
conduit à
$$v_1^1(x)=t^s\left(1+(x/y)t^s+\mathrm{}+(x/y)^{p1}t^{(p1)s}\right).$$
Puisque $`dv_1(x)=syt^st^1dt`$, on obtient
$$v_1^1(x)dv_1(x)=sy\left(t^1dt+\underset{i=1}{\overset{p1}{}}(x/y)^it^{is}t^1dt\right).$$
On transforme cette quantité en écrivant
$$v_1^1(x)dv_1(x)=sy\left(t^1dt+\underset{k=1}{\overset{p1}{}}(x/y)^{s^{k1}}t^{s^k}t^1dt\right)$$
soit encore
$$v_1^1(x)dv_1(x)=sy\left(t^1dt+\underset{k=1}{\overset{(p1)/2}{}}(x/y)^{s^{k1}}t^{s^k}t^1dt+\underset{k=1}{\overset{(p1)/2}{}}(x/y)^{s^{k1}}t^{s^k}t^1dt\right).$$
Pour obtenir l’expression de $`v_1^1(x)dv_1(x)`$ dans la base proposée de $`\mathrm{\Omega }_{dR}^{}(R)`$, introduisons $`\beta _k=(x/y)^{s^{k1}}(y/x)^{s^{k1}}`$. On a
$$v_1^1(x)dv_1(x)=sy\left(f_0^{}+\frac{1}{2}\underset{k=1}{\overset{(p1)/2}{}}\alpha _kf_k^{}+\frac{1}{2}\underset{k=1}{\overset{(p1)/2}{}}\beta _kf_k^+\right).$$
Le calcul de $`\sigma (v_1(x))^1d\sigma (v_1(x))`$ se déduit immédiatement de cette dernière relation car $`D_1^{(p)}`$ est équivariante, $`\sigma (f_0^{})=f_0^{}`$, $`\sigma (f_k^\pm )=\pm f_k^\pm `$. On obtient ainsi
$$\sigma (v^1(x))^1d\sigma (v_1(x))=sy\left(f_0^{}\frac{1}{2}\underset{k=1}{\overset{(p1)/2}{}}\alpha _kf_k^{}+\frac{1}{2}\underset{k=1}{\overset{(p1)/2}{}}\beta _kf_k^+\right),$$
d’où finalement
$$D_1^{(p)}(z_1(x))=sy\left(2f_0^{}+\underset{k=1}{\overset{(p1)/2}{}}\alpha _kf_k^{}\right).$$
$`\mathrm{}`$
###### Définition 53
On note $`V(x)`$ le sous-espace vectoriel de $`K_1^{}(R;𝐙)`$ engendré par l’orbite de $`z_1(x)`$ sous l’action du groupe de Galois $`G`$, c’est-à-dire
$$V(x)=\text{Vect}_{𝐙/p}(z_k(x)\mathrm{,\; 1}k(p1)/2).$$
###### Proposition 54
Soit $`C=C(x)`$ la matrice circulante d’ordre $`(p1)/2`$ à coefficients dans $`𝐙/p`$
$$C=C(x)=\left(\begin{array}{cccc}\alpha _1,& \alpha _2,& \mathrm{},& \alpha _{\frac{p1}{2}}\\ \alpha _{\frac{p1}{2}},& \alpha _1,& \mathrm{},& \alpha _{\frac{p3}{2}}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \alpha _2,& \alpha _3,& \mathrm{},& \alpha _1\end{array}\right)$$
Alors
$$dim_{𝐙/p}V(x)\text{rg}(C(x)).$$
Preuve : à une constante près, les composantes de $`D_1^{(p)}(z_1(x))`$ dans la base $`(f_0^{},f_1^{},\mathrm{},f_{(p1)/2}^{})`$ de $`\mathrm{\Omega }_{dR}^{}(R)`$ sont $`(2,\alpha _1,\mathrm{},\alpha _{(p1)/2}).`$ Puisque $`z_k(x)=g^k(z_1(x))`$, compte tenu de l’action de $`g`$ sur les vecteurs de base$`(f_0^{},f_1^{},\mathrm{},f_{(p1)/2}^{})`$, on en déduit que la matrice des composantes respectives de $`D_1^{(p)}(z_1(x))`$, $`D_1^{(p)}(z_2(x))`$, $`\mathrm{}`$, $`D_1^{(p)}(z_{(p1)/2}(x))`$ a le même rang que la matrice
$$\left(\begin{array}{ccccc}2& \alpha _1,& \alpha _2,& \mathrm{},& \alpha _{\frac{p1}{2}}\\ 2& \alpha _{\frac{p1}{2}},& \alpha _1,& \mathrm{},& \alpha _{\frac{p3}{2}}\\ 2& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 2& \alpha _2,& \alpha _3,& \mathrm{},& \alpha _1\end{array}\right).$$
Puisque $`_{k=1}^{(p1)/2}\alpha _k=1`$, le rang de cette matrice est celui de la matrice $`C(x)`$. L’image de $`V(x)`$ par la trace de Dennis $`D_1^{(p)}`$ a pour dimension le rang de $`C(x)`$, cqfd.
Le calcul du rang de la matrice $`C(x)`$ nécessite l’introduction des polynômes de Mirimanoff.
###### Définition 55
Les polynômes de Mirimanoff $`M_k(X)𝐙/p[X]`$ sont définis pour $`1kp1`$ par
$$M_k(X)=\underset{j=1}{\overset{p1}{}}j^{k1}X^j.$$
Pour $`t𝐙/p`$, on pose
$$r_p(t)=\mathrm{\#}\{k1k(p1)/2,M_{2k+1}(t)0\}$$
C’est le nombre de polynômes de Mirimanoff $`M_j(X)`$ non nuls en la valeur $`t`$ et d’indice $`j`$ impair.
###### Proposition 56
Soient $`x`$ et $`y`$ deux éléments de $`(𝐙/p)^\times `$ tels que $`xy=1`$. Les valeurs propres de la matrice $`C(x)`$ sont $`M_{2k+1}(x/y)`$, $`1k(p1)/2`$. Le rang de la matrice $`C(x)`$ est $`r_p(x/y)`$.
Preuve : Soit $`s`$ le générateur de $`\left(𝐙/p\right)^\times `$ qui détermine l’action du groupe de Galois $`G`$ sur $`A`$ et soit $`v=s^2`$ le générateur de $`𝐙/(p1)/2\left(𝐙/p\right)^\times `$. Les valeurs propres de la matrice $`C`$ sont alors
$`\lambda _k`$ $`={\displaystyle \underset{j=1}{\overset{(p1)/2}{}}}\alpha _j(v^k)^{j1}`$
$`={\displaystyle \underset{j=1}{\overset{(p1)/2}{}}}(x/y)^{s^{j1}}\left(s^{j1}\right)^{2k}+{\displaystyle \underset{j=1}{\overset{(p1)/2}{}}}(x/y)^{s^{j1+(p1)/2}}\left(s^{j1+(p1)/2}\right)^{2k}`$
$`={\displaystyle \underset{j=1}{\overset{p1}{}}}j^{2k}(x/y)^j`$
$`=M_{2k+1}(x/y)`$
Le rang de $`C(x)`$ est le nombre de valeurs propres non nulles. Ces valeurs propres étant les $`M_{2k+1}(x/y)`$, le rang de $`C(x)`$ est bien $`r_p(x/y)`$. $`\mathrm{}`$
En résumé, nous avons montré:
###### Théorème 57
Soient $`x`$ et $`y`$ deux éléments de $`(𝐙/p)^\times `$ tels que $`xy=1`$. Alors
$$dim_{𝐙/p}V(x)r_p(x/y).$$
Remarque 58. Posons
$$r_p=min\{r_p(t),t𝐙/p\{\mathrm{0,1,1}/2\}\}.$$
Alors, pour tout $`x𝐙/p\{\mathrm{0,1,1}/2\}`$, on a
$$(p1)/2dim_{𝐙/p}V(x)r_p.$$
### 4.4 Lien avec les dérivées logarithmiques de Kummer.
Soit toujours $`A`$ l’anneau des entiers du corps cyclotomique $`F=𝐐[\zeta _p]`$ avec $`p`$ premier impair. On pose $`\lambda =1\zeta `$. Identifions $`K_1(A/p;𝐙/p)`$ au groupe multiplicatif $`(1+\lambda 𝐙/p[\lambda ],\times ).`$ Dans ses recherches sur le dernier théorème de Fermat pour les nombres premiers irréguliers, Kummer a introduit certaines “dérivées logarithmiques”. Un élément $`z=_{i=0}^{p1}a_i\zeta ^i`$ de $`A`$, non divisible par $`1\zeta `$ détermine un élément de $`K_1(A/p;𝐙/p)`$ encore noté $`z`$. Pour $`1kp2`$, la dérivée logarithmique $`\mathrm{}_k(z)`$ est définie comme la classe modulo $`p`$ de l’entier
$$\frac{d^k}{dX^k}\left(\text{log}\left(\underset{i=0}{\overset{p2}{}}a_ie^{iX}\right)\right)_{X=0}.$$
Kummer a montré que
$$\mathrm{}_k:(K_1(A/p;𝐙/p),\times )(𝐙/p,+)$$
est un morphisme de groupes.
Soient $`x`$ et $`y`$ deux éléments de $`(𝐙/p)^\times `$ tels que $`xy=1`$. L’élément $`z^{}(x)={\displaystyle \frac{xy\zeta }{xy\zeta ^1}}`$ de $`K_1^{}(A/p;𝐙/p)`$ est tel que $`\mathrm{}_{2k}(z^{}(x))=0`$ , $`\mathrm{}_{2k+1}(z^{}(x))=2\mathrm{}_{2k+1}(xy\zeta ).`$
Mirimanoff a montré (cf. , VII ou ) que pour $`1k(p3)/2`$, on a l’égalité $`\mathrm{}_{2k+1}(xy\zeta )=xM_{2k+1}(x/y)`$. Ceci permet de formuler un lien entre la trace de Dennis à coefficients et les dérivées logarithmiques de Kummer. Soient
$$z_k(x)=\frac{xyt^{s^k}}{xyt^{s^k}}\mathrm{mod}(R)^{\times (p)}$$
les éléments de $`K_1(R;𝐙/p)`$ introduits à la définition 51. Soit $`C(x)\text{Mat}_{(p1)/2}(𝐙/p)`$ la matrice des coordonnées de $`D_1^{(p)}(z_1(x))`$, $`D_1^{(p)}(z_2(x))`$, …, $`D_1^{(p)}(z_{(p1)/2}(x))`$ dans la base de $`\mathrm{\Omega }_{dR}^{}(R)`$ décrite dans la proposition 50. Alors, à une constante près, la matrice $`C(x)`$ a pour valeurs propres les dérivées logarithmiques de Kummer $`\mathrm{}_{2k+1}(xy\zeta )`$.
Signalons un autre lien entre les $`\mathrm{}_{2k+1}(xy\zeta )`$ et la trace de Dennis. Le développement limité à l’ordre 2 de $`D_1^{(p)}(z^{}(x))`$ effectué à la proposition 47 peut être précisé. On obtient
$$D_1^{(p)}(z^{}(x))=\underset{k=0}{\overset{p3}{}}\gamma _k(x)\lambda ^kd\lambda $$
avec $`\gamma _0(x)=2y`$ et
$$\gamma _k(x)=(1)^ky^{k+1}+(k+1)y+\underset{j=1}{\overset{k}{}}jy^2(1+y)^{kj}.$$
Introduisons les vecteurs colonnes $`\mathrm{}(x)`$ et $`D(x)`$ de $`(𝐙/p)^{(p1)/2}`$ définis par $`\mathrm{}(x)=\left(\mathrm{}_{2k+1}(z^{}(x))\right)_{1k(p1)/2}`$ et $`D(x)=\left(\gamma _{2k}(x)\right)_{0k(p3)/2}`$. Pour $`p13`$, on constate qu’il existe une matrice triangulaire $`A\text{GL}_{(p1)/2}(𝐙/p)`$ telle que $`\mathrm{}(x)=AD(x)`$ pour tout $`x(𝐙/p)^\times `$, ce qui montre qu’il est équivalent de connaître la trace de Dennis à coefficients ou les dérivées logarithmiques de Kummer. Il serait intéressant de savoir si cette observation se généralise à tout nombre premier.
### 4.5 Application au premier cas du dernier théorème de Fermat.
Supposons que $`(p,a,b,c)`$ satisfont aux hypothèses DTF1. Notons $`\overline{a}`$, $`\overline{b}`$ et $`\overline{c}`$ les classes respectives de $`a`$, $`b`$ et $`c`$ dans $`𝐙/p`$. Introduisons le sous-espace vectoriel $`V(p,a,b,c)`$ de $`K_1^{}(A;𝐙/p)`$ engendré par l’orbite de
$$z=z_1=\frac{ab\zeta ^s}{ab\zeta ^s}\mathrm{mod}F^{\times (p)},$$
c’est-à-dire
$$V(p,a,b,c)=\text{Vect}_{𝐙/p}(z_k\mathrm{,\; 1}k(p1)/2).$$
###### Proposition 59
On pose $`x=\overline{a}/\overline{c}`$ et $`y=1x=\overline{b}/\overline{c}`$. Avec les notations de la section précédente, on a
$$1dim_{𝐙/p}V(x)dim_{𝐙/p}V(p,a,b,c)0.$$
Preuve : soient $`\phi :AA/p`$ la surjection canonique et $`\psi :RA/p`$ le morphisme d’anneaux défini par $`\psi (t)=1\lambda `$. On désigne par $`\phi _1:K_1(A;𝐙/p)`$ et $`\psi _1:K_1(R;𝐙/p)K_1(A/p;𝐙/p)`$ les applications induites en $`K`$-théorie à coefficients. Dans $`K_1(A;𝐙/p)`$, l’image de $`V(p,a,b,c)`$ par $`\phi _1`$ coïncide avec l’image de $`V(x)`$ par $`\psi _1`$, ce qui montre que la dimension de $`V(p,a,b,c)`$ est supérieure à celle de $`\psi _1(V(x))`$. On vérifie aisément que $`\psi _1`$ est surjective de noyau de dimension $`1`$. On en déduit l’inégalité proposée. $`\mathrm{}`$
###### Théorème 60
Soient $`(p,a,b,c)`$ des entiers satisfaisant aux hypothèses DTF1. Alors, on a les inégalités
$$d_p^{}r_p(\overline{a}/\overline{c})2r_p2.$$
Preuve : D’après le théorème 50 et la proposition ci-dessus, on a les inégalités
$$\begin{array}{c}\text{ }d_p^{}dim_{𝐙/p}K_1^{}(A;𝐙/p)1dim_{𝐙/p}V(p,a,b,c)1\text{ }\hfill \\ \text{ }dim_{𝐙/p}V(\overline{a}/\overline{c})2r_p(\overline{a}/\overline{c})2r_p2.\text{ }\hfill \end{array}$$
Remarque 61. A normalisation près, les calculs ci-dessus correspondent à ceux effectués par Brückner (). Soit $`\chi ^{}`$ la restriction de la trace de Dennis $`D_1^{(p)}`$ à l’espace $`V(\overline{a}/\overline{c})`$. Notre trace $`\chi ^{}`$ est à comparer avec le morphisme $`\chi `$ de ,2.1. Les quantités $`f_i(\eta )`$ introduites en , 3.5 sont telles que $`f_i(\eta )(1)^{i1}yM_{i1}(\overline{a}/\overline{c})\mathrm{mod}p`$ et la minoration $`d_pr_p2`$ correspond à l’inégalité , 5.1. À partir de cette minoration, Brückner montre que le premier cas du dernier théorème de Fermat est vrai si $`p2^{d_p+3}2d_p3`$, où $`d_p=dim_{𝐙/p}Cl(A)_{(p)}.`$ On peut aussi exploiter l’inégalité $`d_p^{}r_p2`$ en procédant comme suit.
###### Proposition 62
Soit $`p`$ un nombre premier. On a
$$d_p^{}<\frac{p+3}{4}$$
Preuve : LA quantité $`p^{d_p^{}}`$ divise $`h^{}`$. D’après et , on a
$$h^{}2p\left(\frac{p}{24}\right)^{\frac{p1}{4}}.$$
On en déduit
$$d_p\frac{p+3}{4}\frac{\mathrm{ln}(2)}{\mathrm{ln}(p)}\frac{(p1)\mathrm{ln}(24)}{4\mathrm{ln}(p)}.$$
Le second membre de cette inégalité est négatif pour $`p2`$.
De l’inégalité $`d_p^{}<(p+3)/4`$ valable pour tout $`p`$ et de l’inégalité $`d_p^{}r_p2`$, conditionnelle à une solution à DTF1, on déduit le résultat suivant.
Scholie Soit $`p3`$ un nombre premier. Si $`r_p(p+11)/4`$, alors le premier cas du dernier théorème de Fermat est satisfait pour $`p`$.
Soulignons que le calcul numérique de $`r_p`$ est assez rapide, ce qui ne semble pas être le cas pour $`d_p`$ ou $`d_p^{}`$. Pour $`p<1000`$, un calcul sur ordinateur montre qu’on a toujours l’inégalité $`r_p>(p+11)/4`$ (le nombre maximal de valeurs nulles pour $`M_{2k+1}(t)`$ est $`7`$).
### 4.6 Lien avec les nombres de Bernoulli.
L’inégalité $`d_p^{}r_p2`$ proposée au théorème 60 peut se retrouver par un autre raisonnement. Le nombre $`r_p`$ est relié à la divisibilité des nombres de Bernoulli au moyen des congruences de Kummer. Rappelons en premier lieu que les nombres de Bernoulli $`B_k𝐐`$ sont définis par
$$\frac{X}{\mathrm{exp}(X)1}=\underset{k0}{}B_k\frac{X^k}{k!}.$$
Soit $`i(p)`$ l’indice d’irrégularité de $`p`$ défini comme le nombre de nombres de Bernoulli divisibles par $`p`$ (c’est-à-dire dont le numérateur est divisible par $`p`$). On a
$$i(p)=\mathrm{\#}\{k\mathrm{,\; 1}k(p3)/2,pB_{2k}\}.$$
Rappelons en second lieu que pour $`x𝐙/p\{\mathrm{0,1}\}`$, on dit que $`x`$ satisfait les congruences de Kummer $`(𝒦)`$ si
$`(𝒦)`$
$$B_{p(2k+1)}M_{2k+1}(x)0\mathrm{mod}p(1k(p3)/2).$$
Il est clair que si $`x`$ est solution des congruences de Kummer, on a l’inégalité
$$r_p(x)i(p).$$
Par ailleurs, Kummer a montré que si $`(p,a,b,c)`$ satisfont aux hypothèses DTF1, alors $`x=\overline{a}/\overline{c}`$ satisfait les congruences $`(𝒦)`$ (cf. , VII ou ). On en déduit $`r_pi(p)`$. Cette inégalité est conditionnelle à l’existence d’une solution à DTF1. L’inégalité $`i(p)d_p^{}`$, indépendante d’une éventuelle solution à DTF1, résulte d’un théorème de Ribet .
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# New star clusters projected close to the Galactic Centre
## 1 Introduction
The extinction in regions projected close to the Galactic Centre and Plane made difficult for many years the systematic study of the extended objects therein embedded. However with recent near infrared (NIR) surveys such as the Two Micron All Sky Survey (hereafter, 2MASS; Skrutskie et al. 1997) and the Deep NIR Southern Sky Survey (DENIS; Epchtein et al. 1997) it is becoming possible to investigate these regions in a spectral domain 10 times less extinguished by dust than the optical. The NIR surveys can provide fundamental data to study the large-scale distribution of galaxies behind the Galactic Plane (Jarret et al. 2000) and the census and distribution of galactic extended objects such as bright, dark and planetary nebulae, globular and open clusters. Harris & Racine (1979) estimated that there should be around $``$ 160-200 galactic globular clusters. However so far there are 147 known globular clusters as indicated in recent compilations (e.g. Harris 1996 and updated version in Web Interface http://physun.physics. mcmaster.ca/Globular.html). Thus, new ones could be hidden behind dust clouds in bulge and disk directions. Indeed, Hurt et al. (1999) reported a candidate globular cluster lying only 10 away from the Galactic Centre and very close to the plane (b = 0.1). On the other hand, young compact clusters close to the Galactic Nucleus such as the Arches and Quintuplet clusters (Glass et al. 1990 and Nagata et al. 1995, respectively) as well as embedded clusters in HII regions and dark clouds are interesting objects to be surveyed using NIR images.
In the present study we use the 2MASS survey in the J (1.25$`\mu `$m), H (1.65$`\mu `$m) and K<sub>s</sub> (2.17$`\mu `$m) bands to search for potential IR clusters in the central parts of the Galaxy or projected on them. In Sect. 2 we discuss the process of inspection of 2MASS JHK<sub>s</sub> images and present a list of 58 new IR clusters or candidates. In Sect. 3 we discuss the angular distribution of the sample. Finally, the concluding remarks are given in Sect. 4.
## 2 IR star clusters or candidates
The search was systematically made in the region of 5$`\times `$5 centred at 17<sup>h</sup>51<sup>m</sup>10<sup>s</sup> -281610<sup>′′</sup> close to the Galactic Centre. In addition we searched for embedded clusters in directions of HII regions and dark clouds for $`|\mathrm{}|4^{}`$. In general we considered objects with size and morphology similar to those of the Arches and Quintuplet which are the closest known clusters to the Galactic Nucleus. We examined a total of 1500 images extracted from the Survey Visualization & Image Server facility (in the Web Interface http://irsa.ipac.caltech.edu/). For each available field, we obtained a K<sub>s</sub> band image and searched for objects with dimensions of about 1 arcmin ($``$ the Arches’ diameter). We extracted new images (JHK<sub>s</sub>) with $`5^{}\times 5^{}`$ centred in the coordinates of each IR cluster candidate from the preliminary list. In this phase we excluded objects affected by artifacts or contaminated by bright stars on J images. Finally, we obtained a list of 58 objects which are given in Table 1. We determined object positions from K<sub>s</sub> images (in FITS format) using SAOIMAGE 1.27.2 developed by Doug Mink. We also measured diameters for the objects and their sizes indicate that most of them are suitable only for large ground-based telescopes or Hubble Space Telescope (HST). Schlegel et al. (1998) built a reddening map from the 100 $`\mu `$m IRAS dust emission distribution considering temperature effects using 100/240 $`\mu `$m DIRBE data. Considering our object coordinates, we extracted reddening values (E(B-V)<sub>FIR</sub>) from Schlegel et al.’s reddening maps using the software dust-getval provided by them. The optical visibility of the IR star clusters or candidates was checked by means of XDSS (Second Generation Digitized Sky Survey) images with $`5^{}\times 5^{}`$ centred in object position obtained in the Web Interface http://cadcwww.dao.nrc.ca/cadcbin/getdss.
Table 1 lists the 58 IR star clusters or candidates, as follows: (1) object identification by a running number along galactic longitude, (2) and (3) galactic coordinates, (4) and (5) equatorial coordinates (J2000 epoch), (6) and (7) the major and minor diameters, (8) optical visibility (yes or no), (9) E(B-V)<sub>FIR</sub> reddening values and (10) comments. According to comments in Table 1, we found 20 objects related to or embedded in known emission nebulae (in catalogues L - Lynds 1963, RCW - Rodgers et al. 1960 and Sh - Sharpless 1959), dark nebula (LDN - Lynds 1962) or reflection nebula (Bernes - Bernes 1976). We note that these objects have high E(B-V)<sub>FIR</sub> values. Since E(B-V)<sub>FIR</sub> values represent the integrated contribution of the dust along the pathsight in a given direction, it is expected high E(B-V)<sub>FIR</sub> values in the direction of these star forming complexes close to the Galactic Centre. However, Dutra & Bica (2000) compared reddening values derived from infrared photometry of embedded clusters in dark clouds with their E(B-V)<sub>FIR</sub> values and concluded that these reddenings are compatible, except in the Galactic Nuclear region where the temperature in the Central Molecular Zone appears to be underestimated by Schlegel et al.’s temperature maps. High E(B-V)<sub>FIR</sub> values for objects with traces of optical visibility suggest background dust sources. It is interesting to note also that we detect two IR cluster candidates (objects 45 and 46) close to the optical star concentrations NGC 6432 and NGC 6465, and two open cluster candidates (objects 18 and 27).
Figure 1 shows a $`3^{}\times 3^{}`$ K<sub>s</sub> image of the Arches cluster used as reference to search for new clusters close to the Galactic Centre. Figure 2 shows a $`3^{}\times 3^{}`$ K<sub>s</sub> image of the IR star cluster candidate number 11, which is an embedded cluster candidate in Sh2-21.
## 3 Angular distribution
Figure 3 shows the angular distribution of the IR clusters or candidates compared to that of 58 catalogued open cluster (Alter et al. 1970, Lyngå 1987, Lauberts 1982) in the 10$`{}_{}{}^{}\times 10^{}`$ region centred on the Galactic Centre. The two known massive compact young clusters Arches and Quintuplet used as references for the search are not indicated, but their galactic coordinates are ($`\mathrm{}`$ = 0.12, b = 0.01) and ($`\mathrm{}`$ = 0.16, b =–0.06), respectively. In the systematically surveyed zone (rectangular area) where we detect 58 new IR clusters or candidates there are 24 previously known open clusters (including the Arches and Quintuplet clusters). We note that there is a deficiency of catalogued open clusters in quadrant Q1, probably caused by nearby dust clouds like those studied by Cambrésy (1999).
Figure 4 shows the angular distribution of the IR clusters or candidates compared to 16 known globular clusters in the same region of Figure 3. Only three known globular clusters (Palomar 6, Terzan 9 and ESO456SC38) are in the systematically surveyed zone (rectangular area) and we have not seen any additional similar object in the area. This fact could be related to globular cluster destruction due to the tidal effects of the central mass concentration in the Galaxy (Aguilar 1993). Barbuy et al. (1998) studied the spatial distribution of the globular clusters within 5 of the Galactic Centre and estimated that there could be 15 missing globular clusters on the opposite side of the Galaxy. They also found evidences of an empty zone inside a radius of about 0.7 kpc, and that only concentrated clusters would have survived to tidal disruption and disk shocking in central parts of the Bulge.
## 4 Concluding remarks
We provide a list of 58 new IR cluster or candidates detected by means of inspections of 2MASS JHK<sub>s</sub> images in the region 5$`\times `$5 centred at 17<sup>h</sup>51<sup>m</sup>10<sup>s</sup> -281610<sup>′′</sup> close to the Galactic Centre, or in directions of HII regions and dark clouds for $`|\mathrm{}|4^{}`$. Most of the objects are structurally similar to the Arches and Quintuplet clusters. Consequently, they require deep CCD images with large ground-based telescopes or HST to establish their nature. We do not detect any new evident globular cluster in the studied region, which is probably caused by globular cluster destruction due to tidal effects near the Galactic Centre. The angular distribution of the known globular and open clusters in the 10$`{}_{}{}^{}\times 10^{}`$ region centred in the Galactic Centre shows a deficiency of clusters in quadrant Q1 (0$`{}_{}{}^{}<`$ $`\mathrm{}`$ $`<5^{}`$ and 0$`{}_{}{}^{}<`$ b $`<5^{}`$) suggesting a more obscured zone. Infrared surveys such as 2MASS are ideal tools to search for distant new IR open clusters and globular clusters in highly obscured and/or star crowded regions, in particular within 5 of the Galactic Centre.
###### Acknowledgements.
This publication makes use of data products from the Two Micron All Sky Survey, which is a joint project of the University of Massachusetts and the Infrared Processing and Analysis Center/California Institute of Technology, funded by the National Aeronautics and Space Administration and the National Science Foundation. This publication also use Digitized Sky Survey images for the analysis. The Digitized Sky Survey was produced at the Space Telescope Science Institute under U.S. Government grant NAG W-2166. The images of these surveys are based on photographic data obtained using the Oschin Schmidt Telescope on Palomar Mountain and the UK Schmidt Telescope. The plates were processed into the present compressed digital form with the permission of these institutions. We acknowledge support from the Brazilian institution CNPq.
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# A General Symbolic Method with Physical Applications
### 0.1 ACKNOWLEDGMENT
While certainly taking full responsibility for the contents of this thesis, it would yet be absurd of me to take full credit, for many people have contributed, however indirectly, to its completion. I’d like to at least take the trouble to thank a few of them.
I’d first like to note my obvious debt to those scientists of various disciplines whose ideas have influenced mine. In this connection the bibliography is intended to be fairly thorough, though not exhaustive, in indicating the works that, with respect to this paper, I’ve come into valuable contact with. I’m satisfied to leave it to the reader to arrive at his own conclusions as to the degree of originality of, or influences on, the thoughts expressed herein.
My exploration of the above works has been facilitated by the suggestions and criticisms of my committee: Professors Bringsjord, Drew, McLaughlin, and Zenzen. These criticisms have decidedly increased the cogency of the presentation which follows and, not only the content, but also the form of this work has thus benefited. I’m especially indebted to my advisor, Professor Drew, for his patience and encouragement during the many revisions this work has undergone.
Finally, I’d like to thank my parents for the innumerable efforts they’ve made on my behalf over the years. I hope they too find some measure of satisfaction in the completion of this thesis.
### 0.2 ABSTRACT
This thesis derives general physical results by an entirely formal process. Beginning with a brief examination of the notion of language itself, it next explores Physics in a schematic fashion in order to arrive at conclusions on the relationship between experience and language. This investigation leads to the hypothesis that there is no separate reality to which language refers, and therefor to the test of constructing physical theory without reference to experiment: If experience is not to direct the interpretation of language then language must yield its own interpretation.
To make such an idea acceptable it is next shown how references to such a presumably fictional entity such as an exterior reality may arise within language itself, and how such references may, and must, be retained. From this starting-point an entirely formal language is developed, along with an associated algebra and a Calculus, neither of which are restricted to finite quantities.
With the completion of the general symbolic system the derivation of both Relativity Theory and Quantum Theory, as well as the formal structures to which they apply, including space-time and sums-over-histories, follows from a purely non-empirical and finitary basis. The dynamical and thermodynamical laws yielding the phenomenological aspects of experience, such as are described by variables for pressure, volume, and temperature, as well as the divisions comprising phases of matter, are also argued to naturally follow on this basis. It is therefor plausibly claimed that the formal approach has succeeded in yielding its own interpretation and in thus reproducing what has previously been asserted to be of necessarily empirical origin.
It is, however, then found that this system is comprised of formally incompatible parts. It is thus apparently necessary to either reject the restriction to finite quantities or else accept the necessity of augmenting the formal system with a properly exterior reality, by which it is meant that “experience” must “inform” the system. Development of formal non-finitary theory is then argued to provide a plausible means of unifying the formalisms of Relativity and Quantum Theory.
## Chapter 1 INTRODUCTION AND HISTORICAL REVIEW
### 1.1 Introduction: On the Notion of Language
Every sentence ever written carries with it an implicit assertion of a power of that sentence to convey an idea and a more general affirmation of written language itself. How indeed could anyone speak of something to which language cannot refer? And language is pervasive; many kinds of languages have been created by many different peoples to serve varied ends.
Besides obvious examples, one may also speak of music as being a kind of language, albeit one with a very flexible grammar. Music, in fact, considered as a language, highlights the active role played by interpretation in determining the ideas conveyed; any given musical score may be performed in a variety of ways by different artists and received variously by distinct audiences. Along with interpretation often comes an assertion of what a language ought to mean and what its author meant. It’s often difficult for a consensus to be achieved on this last issue, especially when the author is dead !
Mathematics is said to be a language; the language of the sciences. While music often allows for a great deal of flexibility in its interpretation, Mathematics strives to be a precise language so that its utilization may have only one result - a uniquely correct interpretation. Mathematics is also often pure; It need not refer to anything outside of itself, in contrast to a musical score where each written note corresponds to some sound to be made.
Consideration of these examples motivates a preliminary stipulation. Language will be said to be in use whenever a system of written symbols has some relevance. This is a deliberately unrestrictive convention, but it has its uses.
First of all, in connection with mathematics, attention is shifted away from any so-called “mathematical objects” which might be supposed to lead an “existence” outside of math itself. Thus sets need not have any Platonic “reality”, and there need not even be mention of an outside real world to which an applied mathematics might refer. Such considerations are obviated and one may ask merely whether or not a given string of symbols is to be used.
Secondly, as illustrated by musical scores, language is then a very general notion. There need not be a uniquely determined relationship between the symbols and what the language refers to. It is merely the possibility of some connection being adopted that justifies asserting that a language is in use.
This last observation might bring to mind the search for alien intelligent life in the universe. Clearly one cannot presume, beforehand, to know the language used by an undiscovered race of beings. Nevertheless, if a radio source, for example, were to exhibit an unexpected and physically unexplained regularity it would not be surprising for there to be conjectures of an intelligent source for the “signal” and attempts to decipher the language utilized in its construction.
It seems, then, that the above definition of language is the only one which adequately addresses the way in which the term is commonly used, although it doesn’t give a constructive idea of what a language may be like or what it may be about: Language is defined by those who use it.
It is the intention of this thesis to explore language in as general a way as possible by specifying the least restrictive and most natural system of writing the author has conceived of. It is hoped that a convincing treatment of this symbolic system itself will shed some light on the ways in which one understands those “things” to which language refers.
### 1.2 Motivation and Historical Review
Before entering upon the developments particular to this thesis, an overview of some relevant prior work in the areas of Linguistics and Cognitive Science, of Metamathematics, and of Physics will be given. The results in these fields will indicate some plausible expectations for the work to be presented here, but this will also be an opportunity to indicate the points of divergence between this thesis and the perspectives advocated in other work.
#### 1.2.1 Linguistics and Cognitive Science
The form of modern Linguistics and Cognitive Science has been strongly influenced by the work of Chomsky<sup>1</sup><sup>1</sup>1 An introduction to this work, as well as references to more technical expositions, can be found in .. This work is foundational in the sense that, even for those that disagree with it, it is ultimately necessary to justify any serious theory in this area by how it relates to Chomsky’s approach and results. What, then, are some of Chomsky’s conclusions about human language?
Language, according to Chomsky, is a discrete combinatorial system; “A finite number of discrete elements (in this case, words) are sampled, combined, and permuted to create larger structures (in this case, sentences) with properties that are quite distinct from those of their elements. For example, the meaning of Man Bites Dog is different from the meaning of any of the three words inside it, and different from the meaning of the same words combined in reverse order.”<sup>2</sup><sup>2</sup>2 See pg. 84 of . This may be contrasted with so-called “blending systems” wherein “the properties of the combination lie between the properties of its elements, and the properties of the elements are lost in the average or mixture. For example, combining red paint and white paint results in pink paint.”<sup>3</sup><sup>3</sup>3 See pg. 85 of . Given that the properties of the elements of a blending system are lost in a combination, it can hardly be sensible to attempt to formulate general statements about a mixture in terms of its elements. Conversely, in a discrete combinatorial system the elements may be recovered from their combination and general rules may be given for the formation of such combinations.
The rules for the combination of the elements of a discrete combinatorial system are called its “grammar”. “A grammar specifies how words may combine to express meanings; that specification is independent of the particular meanings we typically convey or expect others to convey to us.”<sup>4</sup><sup>4</sup>4 See pg. 87 of . The utilization of a grammar is dependent upon the categorization of words into parts of speech: “A part of speech, then, is not a kind of meaning; it is a kind of token that obeys certain formal rules, like a chess piece or a poker chip. … When we construe an aspect of the world as something that can be identified and counted or measured and that can play a role in events, language often allows us to express that aspect as a noun, whether or not it is a physical object. For example, when we say I have three reasons for leaving, we are counting reasons as if they were objects. Similarly, when we construe some aspect of the world as an event or state involving several participants that affect one another, language often allows us to express that aspect as a verb. For example, when we say The situation justified drastic measures, we are talking about justification as if it were something the situation did, though again we know that justification is not something we can watch happening at some particular time and place.”<sup>5</sup><sup>5</sup>5 See pg. 106 of .
Because grammar operates on categories of words it gives sentences a modular structure which may be described, as is familiar from Computer Science<sup>6</sup><sup>6</sup>6 See ., by “production rules”. For example, a noun phrase consists of an optional determiner, followed by any number of adjectives, followed by a noun. A sentence consists of a noun phrase followed by a verb phrase. A verb phrase consists of a verb followed by a noun phrase.<sup>7</sup><sup>7</sup>7 See pg. 98 of . It will be noticed that in these rules the order of the parts of speech within the (written) sentence is crucial to the correct application of the grammar. Languages in which the grammar may really be given in this form are called “isolating languages”. Such a scheme is not always carried out, however. In so-called “inflecting” languages, such as Latin, the spelling of the words themselves is modified in order to reflect the role they play within the sentence. as is done in the conjugation of verbs. When such cases are provided it is possible to rearrange the words within a sentence and still determine the role played by each word. English is, in fact, partially isolating and inflecting. It is clear, however, that any inflected language may be rewritten in a particular order so that it may be interpreted as an isolating language, and the generality of the notion of grammar given above is thus retained in all cases.
The notion of grammar indicated so far allows a different form of production rule for each kind of phrase or part of speech within a given language, and also allows for variation in the grammars of distinct languages. It has, however, been found<sup>8</sup><sup>8</sup>8 See pg. 111 of . that all grammatical rules within a given language, as well as all grammatical rules within all languages, may be put in a standard form. This constitutes Chomsky’s theory of a Universal Grammar. This empirically discovered independence of grammar from the particular culture in which it is used has been asserted to justify the claim that language ability is an inherited and intrinsic property of individuals. The inheritance of grammatical language ability has further been argued to explain the rapidity and particular way in which children learn language. Thus one aspect of linguistic analysis indicates an inherent component of language and calls into question the degree to which one can discern the properties of a truly external world.
More generally, the inherent aspects of language may be asserted to hinge, in part, on the notion of similarity. The identification of words as being of particular parts of speech constitutes recognizing a kind of similarity, and, moreover, the generalizing of specific examples to general rules which is necessary for learning also requires the identification of a similarity. The problem is that similarity is apparently in the mind of the beholder: “Suppose we have three glasses, the first two filled with colorless liquid, the third with a bright red liquid. I might be likely to say the first two are more like each other than either is like the third. But it happens that the first glass is filled with water and the third with water colored with by a drop of vegetable dye, while the second is filled with hydrochloric acid- and I am thirsty.”<sup>9</sup><sup>9</sup>9 See pg. 416 of . This points to a boot-strap problem if one is to be able to learn and not have some inherent sense of similarity. It may be observed, then, that pre-existing notions of similarity are posited by the cognitive scientists in order to solve the old philosophical problem of induction.
Another observation in Linguistics and Cognitive Science points to a difficulty in speaking of a real external world. Here the observation has to do with the metaphors that people use in language. It has been found that there are two basic metaphors in language<sup>10</sup><sup>10</sup>10 See pg. 354 of .: that of location in space and that of force, or cause. The first may be illustrated by the sentences: “Minnie told Mary a story. Alex asked Annie a question. Carol wrote Connie a letter.” Thus “Ideas are gifts, communication is giving, the speaker is the sender, the audience is the recipient, knowing is having.”<sup>11</sup><sup>11</sup>11 See pp. 353-354 of . Force as a metaphor may be illustrated with the following pairs of sentences: “The ball was rolling along the grass.” and “The ball kept on rolling along the grass”. “John doesn’t go out of the house.” and “John can’t go out of the house.” “Larry didn’t close the door.” and “Larry refrained from closing the door.” “The difference is that the second sentence makes us think of an agent exerting some force to overcome resistance or overpower some other force. With the second ball-in-the-grass sentence, the force is literally a physical force. But with John, the force is a desire: a desire to go out which has been restrained. Similarly, the second Larry seems to house one psychic force impelling him to close the door and another that overpowers it.”<sup>12</sup><sup>12</sup>12 See pg. 354 of .
The sentences above can hardly seem unusual, and the pervasiveness of such metaphors is readily discovered. Indeed, “space and force are so basic to language that they are hardly metaphors at all.”<sup>13</sup><sup>13</sup>13 See pg. 357 of . Space and force are, however, integral parts of the description of physical reality. If, as the above quotations suggest, these concepts are building blocks of language itself, then what is the necessity of an external reality anyway? How may references to reality be distinguished from any other kind of reference if such physical terminology is (perhaps) always in use?
This discussion of Linguistics and Cognitive Science may be summed up as follows: There is reason to doubt that Cognitive Science has arrived at a clear distinction between language and the external reality which it, presumably, describes. This would be difficult to do, of course, but it is still somewhat unsatisfactory that it seems to be no help on this point. Nevertheless, it takes a realist position in speaking of a real world and reasoning that humans must have an inherent ability to recognize similarities in order to be able to deal with ambiguous data, as in the example with the glasses, which isn’t just in the mind. This pre-existing faculty for classification extends to the common possession of a Universal Grammar which operates on ordered strings of symbols, and which is the foundation for communication. This grammar may be given in terms of production rules, in a canonical form, for the manipulations of a finite discrete combinatorial system. In any event, it is clear that Linguistics and Cognitive Science have led to the conclusion that language may be understood in terms of a formal symbolic system. The consideration of formal systems in and of themselves, apart from consideration of any external physical “reality”, is the domain of mathematics, so this will be the next subject surveyed.
#### 1.2.2 Metamathematics
The modern study of formal systems within mathematics is referred to as metamathematics, the study of which originated with David Hilbert<sup>14</sup><sup>14</sup>14 For a leisurely exposition of metamathematics see .. “Metamathematics includes the description or definition of formal systems as well as the investigation of properties of formal systems. In dealing with a particular formal system, we may call the system the object theory, and the metamathematics relating to it the its metatheory.
From the standpoint of the metatheory, the object theory is not properly a theory at all … but a system of meaningless objects like the positions in a game of chess, subject to mechanical manipulations like the moves in chess. The object theory is described and studied as a system of symbols and of objects built up out of symbols. The symbols are regarded simply as various kinds of recognizable objects. To fix our ideas we may think of them concretely as marks on paper; or more accurately as abstracted from our experience with symbols as marks on paper. … The other objects of the system are only analyzed with regard to the manner of their composition out of the symbols. By definition, this is all that a formal system shall be as an object of study for metamathematics.
The metatheory belongs to intuitive and informal mathematics (unless the metatheory is itself formalized from a metatheory, which here we leave out of account). The metatheory will be expressed in ordinary language, with mathematical symbols, such as metamathematical variables, introduced according to need. The assertions of the metatheory must be understood. The deductions must carry conviction. They must proceed by intuitive inferences, and not, as the deductions in the formal theory, by applications of stated rules. Rules have been stated to formalize the object theory, but now we must understand without rules how these rules work. An intuitive mathematics is necessary even to define the formal mathematics.
… The methods used in the metatheory shall be restricted to methods, called finitary by the formalists, which employ only intuitively conceivable objects and performable processes. No infinite class may be regarded as a complete whole. Proofs of existence shall give, at least implicitly, a method for constructing the object which is being proved to exist.
This restriction is requisite for the purpose for which Hilbert introduces metamathematics. Propositions of a given mathematical theory may fail to have a clear meaning, and inferences in it may not carry indubitable evidence. By formalizing the theory, the development of the theory is reduced to form and rule. There is no longer ambiguity about what constitutes a statement of the theory, or what constitutes a proof in the theory. Then the question whether the methods which have been formalized in it lead to contradiction, and other questions about the effect of those methods, are to be investigated in the metatheory, by methods not subject to the same doubts as the methods of the original theory.”<sup>15</sup><sup>15</sup>15 This excellent and authoritative quote may be found on pp. 62-63 of .
According to the above characterization, metamathematics achieves clarity of meaning and definiteness of implication for the object theory, but the object theory is not all of mathematics. The undefined metatheory forever remains as an exterior criterion of the acceptability of the formalized object theory, and it is thus implicitly impossible, according to this standpoint, to ever arrive at a complete notion of mathematics which is both clear in its meaning and definite in its implications. It is clear, however, that this scheme allows a progressive clarification and development of the formal scheme. What, then, has been achieved along these lines?
It has been indicated that the object theory is judged from the perspective of the metatheory; Thus the metatheory must, if the meaning of statements in the object theory is to be definite, definitely classify each statement in the object theory as being either acceptable or not. This classification is achieved by a mapping of each statement in the object theory onto either “True” or “False”; This mapping is called the “interpretation” of the object theory. It is clear, then, that from the metamathematical perspective the notion of truth is contingent upon there being a separate metatheory which determines, via an interpretation, what is to be true.
Metamathematics is also supposed to make the acceptability of proofs of an assertion manifest, where a proof is a formal derivation of a true statement. Formal operations, being rules for the combination of meaningless symbols, can, in themselves, only lead from one combination of meaningless symbols to another and thus make one such combination depend upon another. Once an interpretation is provided, however, there is hope that certain formal operations may preserve the initial truth value in a progression of formal combinations leading to the statement it is desired to prove. Such a proof is not absolute, but rather expresses the truth of the conclusion as being contingent upon the truth of the initial statement in the derivation. This method cannot remove the ultimately contingent nature of such proofs, so that certain statements cannot be proven but may only be assigned the value “true” according to the interpretation. Such statements are called “axioms.” The truth-preserving formal rules would then be termed “correct”. Those statements which may be proven are termed “theorems.”
The axioms and the particular proof theory applied to an object theory are determined by the metatheory. Thus, from the metamathematical perspective, once a formal system is chosen, formal mathematics boils down to the selection of the axioms and the formal rules for carrying out proofs. Now it will be noticed that, once axioms and the formal rules for proofs are given, the formal system inherits an order leading from axioms to theorems. The direction and extent of this order will be determined by the axioms and proof system. There is thus a correspondence between the rules for proof accepted by the metatheory and some ordinal number. Such a selection cannot, in the above sense, be proven correct but may only be a matter of preference. The choice made in modern mathematics may, for the most part, be said to be for first order logic with proof based on the acceptance of arguments by induction, in the usual sense, on the integers, the integers having ordinal number $`\omega `$<sup>16</sup><sup>16</sup>16 For a critical discussion of this standpoint see ..
What is the result of this chosen perspective? Hilbert’s original motivation in proposing the metamathematical approach was to show, at least, that mathematics was both “consistent”, so that no false statement may be proven, and “complete”, so that any true statement may be proven. As is well known, Gödel<sup>17</sup><sup>17</sup>17 See either of or . dashed the hopes for the fulfillment of this program when he showed that, for any object theory which may model arithmetic and for which induction is defined to be over the integers, it may be proven that the object theory is both incomplete and cannot be proven to be consistent. It will be illuminating to proceed at this point to a discussion of the role of judgements based on the metatheory in the proof of the first of these assertions.
The key to the proof of these results lies in the notion of Gödel numbering. Gödel numbering encodes statements in first order logic as numbers in the arithmetic system which is formalized. Such numbers may then be arguments of predicates in the formal system so that it becomes possible, from the metamathematical perspective, to interpret some such statements as being self-referential. In particular, Gödel was able to construct a statement G such that from the metamathematical perspective G asserts that G is not a theorem in the object theory. It then is found to be impossible to prove G. This may be explained as follows: It is asked whether or not G is a theorem. If it is, then G, being a theorem, must assert a truth. G, however, according to the metatheory, asserts that G is not a theorem, so this is a contradiction. If the object theory is to be consistent then this alternative must be rejected. On the other hand, supposing G is not a theorem, it then follows, again according to the metatheory, that G asserts a truth. Then, since G is both true and not a theorem of the object theory, it follows that the object theory is incomplete.
Now G being true entails its negation being false, and, false statements being underivable, it follows that neither G nor its negation may be a theorem of the object theory. G is thus an “undecidable” proposition; It’s a statement the object theory can make no assertion about, one way or the other. In order to remedy this “hole” in the object theory it is necessary that the object theory be amended so that it adopts either G or its negation as a new axiom. The first option, usually referred to as being the standard one, is straight-forward enough. The second alternative, referred to as the non-standard one, it turns out, requires the augmenting of the object theory with a new class of symbols and a corresponding reinterpretation of predicates to admit the new class of symbols. This option points directly to an extension of the form of the object theory, as is illustrated by the reformulation of the Calculus within non-standard analysis<sup>18</sup><sup>18</sup>18 See .. The interesting thing about non-standard analysis is that it allows rigorous definition and manipulation of infinite and infinitesimal quantities. It is clear, in any event then, that Gödel’s Incompleteness Theorem, rather than merely being prohibitive, may be fruitful as well.
It should be carefully noted that the metamathematical basis for Gödel’s theorems was, in part, the acceptance of induction over the integers. While this is the standard approach, it is worth noting that Gödel’s theorem asserting the unprovability of consistency may be refuted if one is willing to accept transfinite induction over a larger ordinal. It has been shown that induction over an ordinal known as $`ϵ_0`$ allows a consistency proof to be carried through<sup>19</sup><sup>19</sup>19 See ..
In conclusion, the discussion of metamathematics has revealed a number of things. Metamathematics sets out, first of all, on the dual basis of an object theory and a metatheory. While the object theory is formalized so that acceptable statements may be definitely recognized, the metatheory, which justifies and interprets the object theory, is not necessarily formalized at all. The metatheory operates on an intuitive basis, which, in part, is simply a way of saying that the way that the metatheory operates is undefined. There is a parallel here between human language and physical reality on the one hand, and the object theory and the metatheory on the other. Proceeding nonetheless to consider the implications of this approach, it is found that the interpretive role played by the metatheory leads to the usual notions of truth and proof in an axiomatic mathematical system. Proof theory introduces a notion of order in the object theory. There is also a consequent parallel between the kinds of inductions accepted for proofs and the size of ordinals that must be admitted in the object theory. Gödel’s theorems then show that despite the definiteness of the notions of truth and proof which metamathematics utilizes, the presumably categorical power of these notions doesn’t obtain. It does, nevertheless, point to a limitless flexibility in the formation of the formal system. This flexibility is manifest in the extensibilty of the system of axioms and in the potential variety of variable types that may be constructed in non-standard theories. Infinities and infinitesimals, for example, naturally find their way into the formal system in this way. This fluid nature of the formal language of metamathematics suggests that if any comprehensive conclusions about the relationship between reality and language are to be arrived at, they might not derive from the given form of a language, but may be justified by a systematic referal to an exterior reality instead. For this reason attention now turns to an overview of Physics.
#### 1.2.3 Physics
Science is distinguished from other mental pursuits by its recourse to experimentation, this, it may be said, being a systematic procedure for determining how to use language to refer to “experience”. While “reality” is a scientifically indeterminate word, it is supposed that, via experimentation, Science refers to something other than a formal symbolic game<sup>20</sup><sup>20</sup>20 The notion of reality, from a quantum-mechanical perspective, is discussed in .. Thus the exploration of Physics will begin with a discussion of experimentation, the scientific method.
The scientific method may be described in a number of steps. First there is Observation. In this phase of the process it is simply asserted that something is noticed in the “outside world” and this something is described. It is important to note, right away, that science therefor deals with that which, though taken to come from outside of the linguistic apparatus, is, nevertheless, describable within this system.
After Observation comes a stage in which Questions may be formulated. Here the essential activity is that of trying to guess at any possible connections between different aspects of the physical situation and the particular phenomenon noticed in the previous step. In other words, an attempt is made to imagine ways in which the occurrence of the noticed phenomenon might depend on other physical events which then would precondition it. Loosely stated, this amounts to seeking causes for the observed effect. It should be noted that such questions, like the observed phenomenon, arise in an undefined way, as saying that such ideas are found by using one’s imagination does not define imagination itself. It is also worth noting that any presumed cause, like any other condition which may be distinguished, cannot always obtain; a term which always applies is not descriptive. Under these conditions it is at least conceivable that a supposed cause and the observed effect may be related.
The third stage of the scientific method is the formulation of a Hypothesis. The Hypothesis is a clearly stated trial explanation for the observed phenomenon chosen, according to best judgement, from the possibilities considered in the previous stage. This step is also undefined; it is simply a judgement call. It is clear that, from the scientific perspective, it cannot matter how questions or hypothesis are formed. The only thing that can matter is the usefulness of these questions and hypotheses once they are “found.”
The fourth step is to carry out a controlled Experiment. This means that a well defined series of pairs of physical situations is created such that in the first member of each pair the supposed cause obtains while the other half of the pair, called the control, is a situation which is identical to the first except that the supposed cause does not obtain. Having created these initial conditions, the experiment consists in waiting to see, and noting in each case, whether or not the effect to be explained occurs. It is presumed that any ”real” cause would be in effect during these trials and would show itself, at least sometimes, in the effect occurring in the non-controls and not occurring in the controls. Any hypothetical cause is accepted or rejected accordingly.
The final stage in this progression of the scientific method is that of Theorizing. Those relationships which pass the experimental step are taken, at least provisionally, to be factual. These relationships are then thought about (another undefined step) and an attempt is made to create a mathematical theory which comprises all of these relationships. This theory may make predictions as to what will happen in experiments not yet performed, and this becomes a test for the theory itself to pass. In any event, having cycled through these five stages, Science always returns to the first to look for new observations.
A number of questions may be raised concerning the adequacy of the scientific method outlined above. The world being tremendously complex, which phenomena are worth observing? Since phenomena, in order to be observed in the above sense, must be described, doesn’t this precondition what may be observed? Might it not be a consequence of this that only phenomena which are sufficiently simple may be observed? If the number of hypotheses that might explain an effect is not finite, doesn’t it follow that the experimental investigations might run themselves into a small corner in the sense of restricting the way the world is viewed by scientists? And might it not happen that experiment never checks all hypotheses which might explain a given effect? The scientific response to these issues may be inferred by examining Newton’s Rules of Reasoning in Philosophy<sup>21</sup><sup>21</sup>21 See Part III of .. While these rules do not constitute a formal oath of office for all scientists, it may be fairly asserted that they do reasonably reflect the practical attitude of most scientists.
Newton’s Rules of Reasoning in Philosophy are:
Rule I: We are to admit no more causes of natural things than such as are both true and sufficient to explain their appearances.
Rule II: Therefor to the same natural effects we must, as far as possible, assign the same causes.
Rule III: The qualities of bodies, which admit neither intensification nor remission of degrees, and which are found to belong to all bodies within the reach of our experiments, are to be esteemed the universal qualities of all bodies whatsoever.
Rule IV: In experimental philosophy we are to look upon propositions inferred by natural induction from phenomena as accurately or very nearly true, notwithstanding any contrary hypotheses that may be imagined, till such time as other phenomena occur, by which they may either be made more accurate, or liable to exceptions.
Rule IV refers to “natural induction”, and it should be noted that this term has been the subject of much debate. For the purposes of the argument to follow, this phrase will be replaced by “experimentation”. This seems reasonable as, after all, Newton’s rules are being interpreted and applied, rather than merely cited.
Newton’s first rule begins by subordinating causes to effects: Particular causes may be admitted only if they explain observed effects. However, it should be clearly noted that causes are otherwise unrestricted by Newton’s rules and it is nowhere indicated how one is to choose among many particular causes which might explain a particular effect. It is only asserted that if one of these causes is acceptable according to experiment, and is therefor “true”, then it may be adopted. This is in agreement with a minimal interpretation of the experimental procedure outlined above.
It is clear then that both the experimental method and Newton’s commentary on it leave the acts of observation, questioning, and hypothesizing without limits defined prior to experimentation. It follows from this that no limits, in principle, may be asserted to apply to the effective range of experimental method save any which may restrict the form of statements and hypotheses so that they are subject to controlled experiment. Thus the scientific method is not restricted to investigate only particular observations, and there is no bar to future effects necessitating the abandonment of previous causes, as is indicated explicitly in Rule IV, so that experimental investigation need not lead to a narrow outlook. It is not necessary that all hypothetical causes be checked.
What may be said about the form of statements about experimental observations and causes? It may be observed that the experimental method would be fruitless if it were merely the case that all observed effects were in one-to-one correspondence with causes, for then there would be no conceptual gain in speaking of causes: Note here that even in a deterministic physics, where prediction may in principle be perfect and invertible, an additional idea of the passage of time is involved. Thus, in general, effects are not equivalent to causes though they proceed from them. Furthermore, if the description of effects is not to be superfluous, then their number must be finite; In a description of an infinity of elements terms could be omitted and still one would retain a description which may correspond to the original. Appealing to Rule I then, it is clear that both causes and effects may and must be representable by a finite number of finite numbers, and this may be taken to be the form that experiment requires statements adopt. This conclusion is implicit in Newton’s Rules III and IV where the “qualities of bodies” are spoken of as being “accurately or very nearly true” and it is said that the “universal qualities” of bodies “admit neither intensification nor remission of degrees”. That physical quantities are to be represented by finite numbers is not explicit in the previous definition of the experimental method, though it is an unspoken canon of scientists.
Consider now Rule II: It is, as Newton asserts, a consequence of Rule I. It goes beyond Rule I in pointing out that physical explanation shall therefor be universal in that all phenomena must be viewed within a single framework. In fact, Rule II has the effect of imposing the requirement that no part of the universe, nor any aspect of its activity, may be considered forever isolated from the rest, as in that case the isolated part would be considered physically irrelevant through not having effects which experiments need address. This then also justifies Newton’s third rule, and Rule II furthermore serves as a constraint on theory, pushing physicists towards unification.
The discussion thus far has focused on experiment as a method of analyzing experience but it hasn’t yet been shown that this method, as opposed to the methods of linguistics and metamathematics, results in anything other than a formal symbolic game. Does the scientific method really access something separate from language itself, some external reality? In order to address this question the final aspect of the scientific method, that of theory, will be considered next.
There is a well-known understanding of the relationship between theory and experience: “Every theory can be divided into two separate parts, the formal part, and the interpretive part. The formal part consists of a purely logico-mathematical structure, i.e., a collection of symbols together with rules for their manipulation, while the interpretive part consists of a set of “associations”, which are rules which put some of the elements of the formal part into correspondence with the perceived world. The essential point of a theory, then, is that it is a mathematical model, together with an isomorphism between the model and the world of experience (i.e., the sense perceptions of the individual, or the “real world” - depending upon one’s choice of epistemology).”<sup>22</sup><sup>22</sup>22 See pg. 133 of .. Here again, as was the case with Linguistics and Metamathematics, there is a separation between experience and language, and where that separation is to be found, is, ultimately, intuited rather than defined.
If the separation between the real and the merely formal is not defined a priori, then perhaps, as a minimum, that separation is at least definitely stated in contemporary physics. Contemporary physics is, for the most part, quantum theory, and the mathematical model of quantum theory was given, essentially, by von Neumann<sup>23</sup><sup>23</sup>23 See .. Rather than discuss the intricacies of this model it is sufficient for current purposes to note here that calculations in this formalism are of two types: “Process 1” calculations which apply whenever new experimental data are acquired in “measurements”, and “Process 2” calculations which apply for those times in the interim between measurements<sup>24</sup><sup>24</sup>24See pg. 3 of .. This distinction in the formal operations of the theory clearly and formally distinguishes “real” data from merely formal variables. This formal separation would unquestionably draw a definite line between the “real” and the formal if it weren’t also the case that the interpretation of the model leads to difficulties<sup>25</sup><sup>25</sup>25 For a quick description of these difficulties see pp. 3-10 of . For a selection of some different ways of viewing the quantum formalism also consult pp. 181-195 of .. These difficulties, known as the “measurement problem”, have lead to a variety of elaborate attempts at resolution<sup>26</sup><sup>26</sup>26 See ., but, as yet, none are uniformly accepted.
Contemporary physics, having failed at clearly separating formalism and “reality”, may even be doubted to have adopted any such separation in the first place, instead almost reducing physics to formalism. To begin with, it has been shown, by those studying “Quantum Logic”, that the formal apparatus of quantum theory need not be motivated by experiment, but is, rather, a formal embodiment of the propositional calculus<sup>27</sup><sup>27</sup>27 For an overview of this formalism see . For an in depth treatment see .. The reduction of physics to formalism, it may be argued, is further supported by a founder of quantum theory, Niels Bohr, in the following comparison of quantum and classical physics: “In the case of quantum phenomena, the unlimited divisibility of events implied in such an (classical) account is, in principle, excluded by the requirement to specify the experimental conditions. Indeed, the feature of wholeness typical of proper quantum phenomena finds its logical expression in the circumstance that any attempt at a well-defined subdivision would demand a change in the experimental arrangement incompatible with the definition of the phenomena under investigation.”<sup>28</sup><sup>28</sup>28 See pg. 4 of .( italics mine). As to the concept of a physical condition, Bohr seems to be identifying it with, if not replacing it by, its formal description. Some may yet feel that these points are exaggerated and merely “philosophical”, and do not really apply to practical physics. In answer to this it is sufficient to turn consideration to the modern theory of quarks, for quarks are supposedly fundamental constituents of physical description motivated by experiment, and, at the same time, are subject to “confinement”. Thus, while quarks are believed to “exist” they cannot, in principle, ever be seen! Quarks can impact experimental results only indirectly, and are thus primarily products of theory and not observation.
Beyond formalization of the “existence” of external reality, science has even gone so far as to attempt the formalization of the individuals themselves who have experiences, so that none could speak of anything which is not formal. This program may be seen to be a consequence of Science’s previously mentioned drive towards universality. According to von Neumann: “… it is a fundamental requirement of the scientific viewpoint - the so-called principle of the psycho-physical parallelism - that it must be possible so to describe the extra-physical process of the subjective perception as if it were in reality in the physical world - i.e., to assign to its parts equivalent physical processes in the objective environment, in ordinary space.”<sup>29</sup><sup>29</sup>29 See pg. 418 of . The physical world being formalized, as indicated above, it then follows that individual experience ought to be as well. Along the lines of deterministic theory, work in Computer Science has produced Turing machines as well as other apparatus of Artificial Intelligence<sup>30</sup><sup>30</sup>30See for a basic coverage. For detailed theory consult .. von Neumann went so far, in his theory of self-reproducing automata, as to formalize deterministic machines that would carry on many of the formal operations of life-forms<sup>31</sup><sup>31</sup>31See .. von Neumann also addressed non-deterministic theory by showing that nearly deterministic behavior may be synthesized by the parallel organization of non-deterministic elements<sup>32</sup><sup>32</sup>32 See .. Consequently it is now common practice in physics to, in effect, replace observers by Turing machines<sup>33</sup><sup>33</sup>33See pg. 64 of ..
The above discussion has fairly well summarized the results of contemporary physics. A review of physics, then, indicates that Science attempts to answer the need for specifying a definite interpretation of its formalism through experimentation. The nature of the experimental procedure itself indicates that there are no necessary restrictions to the range of the interpretations thus derived, though such investigations are formally restricted to measurements expressible by finitely many finite numbers. Science is concerned with having an explanation of observations, but by no means need there be a unique explanation or theory. In fact, theories need not arise solely or uniquely from experiment either, as is illustrated in the case of quantum theory. This freeing of theory from strict correspondence with observation has resulted in modern physical theory being almost totally formal in nature. It has also left in question the degree to which anything besides a formalism is required or even specifiable.
### 1.3 A Departure from Prior Approaches
Having surveyed Linguistics and Cognitive Science, Metamathematics, and contemporary Physics, an attempt will now be made to glean any wisdom common to these disciplines. Some of those things which these disciplines have in common will be found to hold in the symbolic system developed here. Particular aspects specific to some of these areas will also feature in the pages to follow. And, as it will also be observed, there is a similarity in the difficulties facing these theories which will be an important indicator of a point of departure which distinguishes the theory developed here from previous developments.
It may be recalled that each of the noted areas does make use of written discrete combinatorial systems. Manipulations of these systems were, in all cases, specified by clearly stated formal grammars. In applying grammars rules were in effect according to the category of symbols being manipulated, there being no concern for any particular additional meaning attached to the symbols. This lends a certain rigor to these manipulations, but, in all cases, there was also an appeal made to an intuited level of interpretation and justification which, in the case of Physics at least, yields confusion if not contradiction. In no case was there an a priori clear division between what is to be merely formal and what is to be intuited.
The study of Linguistics additionally pointed to the particular relevance of the order of written symbols while also noting that the Universal Grammar may be taken to derive from an intuited notion of similarity or categorization. Consideration of Metamathematics taught that the notions Truth and Proof aren’t as categorical as might be formally desired, but that they are instead grounded, ultimately, in a choice of interpretation of the symbolism. Surprisingly, metamathematics, instead of demanding rigid and ultimate truths, requires an infinite flexibility in an allowed step-wise extension of any sufficiently complex formal system. This extension results in, among other things, a rigorous definition of infinities and infinitesimals. Experimentation, which was discussed in order to curtail this embarrassment of riches, demands, instead, operating with finite quantities, but still hasn’t categorically solved the problem of the interpretation of the formalism.
A concern common to all of the above considerations is, given that language is to play a role at all, the need for precisely determining the proper role and extent of the respective formalisms. In this connection Wittgenstein’s dictum<sup>34</sup><sup>34</sup>34For an overview of Wittgenstein’s philosophy see ., that that which cannot be said ought to be passed over in silence, seems pointedly relevant. It may, however, be pointed out that he who would wholely reject the relevance of language can hardly state his case! Furthermore, it is hardly possible to describe that which is outside of language. It seems, besides, that language ought to play some part in our experience as without such interpretation that experience is, arguably, a chaotic jumble. None of this, however, resolves the difficulties encountered above wherein language referred to something in addition to itself.
The notion that language might be able to represent and refer to something outside of itself might seem questionable. Perhaps such a conception is even self-contradictory, yet one might also conclude that empty formalism is the only alternative to such an approach and feel thereby driven to infer some external “physical” basis for, or content of, language. Re-examination of the nature of the axiomatic approach will suggest, however, that this is a false dichotomy and that one might instead adopt a third viewpoint which embraces none of the arbitrary formalisms such as have arisen in pure mathematics and yet also avoids such problematic combinations of intuition and formalism as have arisen in Physics. The considerations to follow outline such a re-examination and give an overview of the argument to be presented at length in the remainder of this thesis.
Recall that one lesson of Gödel’s Incompleteness Theorem is that any sufficiently powerful axiomatic system will have “holes” corresponding to undecidable propositions. There are no stated constraints on how such “holes” may be filled so that, in general, no axiom system is formally preferable to another. Usually one is thus led to recognize the limitations and arbitrary nature of the axioms of any particular metamathematical system. Not withstanding this result the axiomatic method is still employed by both mathematicians and physicists alike in order to draw definite conclusions from definite premises.
The formal symmetry in the admissibility of distinct axiomatic theories is usually broken either as a matter of taste, as in the case of pure mathematics, or via reference to experiment, as in the case of applied mathematics and Physics. That the relationship between empirical data and formal structures is indefinite has perhaps already been plausibly indicated and will be further advocated in Chapter 2. Any plausibility in these observations argues against the second means of identifying relevant formalisms though it cannot, by any means, prove that such a rejection is mandatory. If such a course is adopted, it yet remains to formally characterize the “physically relevant” structures within the formal theory. In this connection it may be anticipated that the presumed indefiniteness in the empirical relevance of particular formal structures, together with the previously discussed physical requirement that “measurements” be expressible as finite collections of finite (real) numbers, leads to the expectation that the finiteness of data might be the only possible formal requirement that could characterize an empirically relevant symbolism. This observation identifies an interesting class of formalisms to be investigated and gives some hope of success in replacing the empirical method with a strictly formal one.
The role of experiment being at least questionable, the axiomatic method yet remains arbitrary in that distinct and incompatible axiomatic systems are, in themselves, of equal standing. The selection of any particular collection of axioms being arbitrary, it would seem that any acceptable non-empirical symbolic method must be non-axiomatic. In light of this it may be asked how one is to determine a self-sufficient and constructive symbolic method. An answer: Rather than thinking of Gödel’s Incompleteness Theorem as being a constraint on particular axiom systems it is instead possible to interpret it as being a critique of the axiomatic method itself. Such a critique is a statement of what the formalisms of all such distinct axiomatic systems have in common.
The acceptability of such a statement suggests that the formal manipulations of a non-axiomatic system be taken to be those compatible with the adoption of any particular axiomatic system. This may also be thought of as requiring that the formal manipulations demanded by any given axiomatic system may be thought of as special restricted cases of the application of the formal rules of the non-axiomatic system. Such restrictions in the applicability of formal manipulations are determined by the axioms of the given axiomatic system so that such systems must be embeddable in the non-axiomatic system.
The axioms of such embedded formal systems may be thought of as summarizing the “truths” of the experiences of particular individuals so that, if this identification is made, it may be seen that the formal rules of the non-axiomatic system then ought to be consistent with the testimony of “observers”. If such an identification is to be made then it remains to at least plausibly illustrate that the “realistic” structure of everyday language may nevertheless be thought of in terms of, and justified by, the strictly formal considerations of a non-axiomatic theory. In this sense the formalism may be said to generate its own interpretation. Such a development, which coincides with an intuitive categorization of symbols such as Linguists infer in the operation of the Universal Grammar, is given in Chapter 3. Accepting such an identification then allows the derivation of the formal manipulations of the non-axiomatic system to be carried out, beginning with Chapter 4, in an intuitively plausible manner.
It should be noted that the non-axiomatic formal approach, by definition, does not presume that axioms may not be adopted, but, rather, attempts to investigate what is common to all such axiom systems. Definite conclusions correspond, as always, to definite axioms. While a non-axiomatic theory may make note of such connections it cannot endorse any conclusions as being categorically “true”. Conversely, any definite results must be those that may be derived on the basis of some axiomatic theory. If this is kept in mind it should be clear that this thesis can never make claims that one must necessarily accept, nor can it include formal manipulations that amount to proofs in the absense of reference to any axioms. In this sense a non-axiomatic theory may be said to be based upon the notion of contingency rather than truth. It may also be noted that, in accordance with the notion of a non-axiomatic theory, all non-axiomatic theories may be considered to be identical. Thus the formal structure of the non-axiomatic formal system will be identified in this thesis as being that of the General Symbolic System.
In order to make the above ideas clearer it may be helpful to make a detour in order to compare these ideas with some other historically important ideas not already discussed. It seems plausible, first of all, that if the non-axiomatic method is to be self-sufficient, then it might profitably be compared to the approach of various mathematical schools of thought. This comparison will then lead to a consideration of some of the philosophical ideas of Wittgenstein.
Historically, there have been three main schools of mathematical thought, namely: Formalism, Logicism, and Intuitionism. It is easy to quickly point out differences between the non-axiomatic method and the approach of each of these schools. Formalism has already been discredited by Gödel, but the symbolic method to be developed here manifestly differs from Formalism in the sense that there is to be no object theory separate from any intuitive metatheory. Logicism attempts to construct all of formal mathematics starting from a propositional system of assertions, assertions such as that a particular mathematical object has a particular property. The point of divergence between the approach developed here and that of Logicism is, as is explained in section 5 of Chapter 3, that the general symbolic system doesn’t adopt any particular kinds of mathematical objects as being such final categories. Finally, the Intuitionists are known for their rejection of the notions of mathematical infinities and of the Law of the Excluded Middle, which mandates that all statements are either true or false. The above approach, while evidently sharing the Intuitionistic attitude to truth, will distinguish itself from Intuitionism in fully embracing formal infinities.
Having distinguished the desired formalism from that of all standard mathematical schools of thought it would seem that the above ideas are on very philosophically shaky ground. In this vein it will be interesting to consider some of Wittgenstein’s ideas.
One of the foundational goals of Wittgenstein’s philosophy<sup>35</sup><sup>35</sup>35Refer again to for a discussion of Wittgenstein’s philosophy. is to prevent the imprecise use of language. Wittgenstein felt, in fact, that with a precise understanding of language there could be no questions, properly speaking, that could not be answered. In such a case the long-standing paradoxes of philosophy would have to therefor be rooted in an improper use of language. While his later philosophy attempted to explore language empirically<sup>36</sup><sup>36</sup>36These ideas were published posthumously in ., and is therefor apparently to be rejected here, his earlier philosophy, like the program to be pursued here, attempted to understand the nature of experience through an analysis of pure formalism<sup>37</sup><sup>37</sup>37Wittgenstein’s earlier philosophy is given by him in . and therefor might be akin to the ideas to be developed here.
Although Wittgenstein’s early ideas do proceed from a Logicist viewpoint, they will nevertheless be discussed. Wittgenstein espoused the idea that language gave a “picture” of reality, wherein the structure of language exactly mirrored the structure of an external reality. Wittgenstein thus postulates a neumenal world of ideas, similar to Plato’s idea of forms, parallel to a physical reality. Such an approach differs from the theory to be developed here in necessarily speaking of an external reality, and it seems suspicious that language and physical reality might coincidentally have identical structures, but the fact that the nature of language is to exactly determine the nature of experience is nevertheless quite similar to the idea that, in the general symbolic system, the formalism generates its own interpretation.
In further elaborating his idea of the “picture” Wittgenstein came to the conclusion that physical reality is, like ordinary written language, composed of atomic elements corresponding to those symbols representing ideas which are not further analyzed. In accordance with his Logicist views he also, therefor, came to the conclusion that language, and therefor physical reality also, is ultimately tautological in nature, there being no a priori way to distinguish atomic elements of reality. This led him to later reject the results of these investigations because he could not reconcile the complete symmetry between the atoms of his neumenal world and the apparent diversity of the physical world.
Wittgenstein’s early ideas may be compared to those of this thesis in the following way. While the structure of language is the only structure there will be that may be discussed in the general symbolic system, it is not presumed in the developments to follow that there will be atoms in the symbolic system into which all other expressions may be analyzed. Such an analysis amounts to a proof in a formal system, while such atoms correspond to axioms. The general symbolic system rejects, by definition, such a conclusion. Instead of this, as is explained in section 9 of Chapter 7, the considerations of the general symbolic system will lead to finite levels of description which may be further elaborated to an arbitrary degree of complexity. Additionally, acceptance of the methods of this thesis leads, unlike Wittgenstein’s early work, to the development, in Chapter 7, of a purely formal symbolic system to a point where the structure of language apparently does accurately reflect the diversity of experience. This is all that will be said of Wittgenstein’s ideas here, and discussion now returns to an overview of the mode of development of the general symbolic system.
Beyond the synthetic definition of, and presumed uniqueness of, the general symbolic system, it remains to indicate the particular formal manipulations which it should adopt. Here it suffices to generalize the formal rules of a particular simple axiomatic system as, after all, such structure must be embedded in the general symbolic system itself. Such a requirement is, as is indicated in Chapter 4 beginning with section 4, not merely restrictive, but decisive: The general symbolic system embeds the formal structure of Lattices and is thereby identified to be a division ring with a partial order.
The notions of conventional mathematics are developed with a bias in favor of finite quantities, this perhaps being associated with the atomism which corresponds to proofs in axiomatic theories. The methods of the Calculus, in their reliance on the $`ϵ`$-$`\delta `$ definition of a limit, are a case in point. Though finite quantities are perhaps favored by empiricists, no such preference is to be found in the basis for the development of the general symbolic system. Both for the sake of its own development, as well as for the sake of maintaining its comparability to the usual mathematics, Chapter 5 is devoted to the task of generalizing, within the general symbolic system, the formal operations of the Calculus. This generalization takes particular advantage of the non-commuting “multiplication” of the general symbolic system and thus echoes the stress which Linguistic analysis laid on the importance of the notion of order in language. Furthermore, Chapter 6 takes advantage of the generalized Calculus in order to arrive at some final conclusions as to the algebraic nature of the formal rules of the general symbolic system.
With all that has been said above, it should still be recalled that the purpose of the development of the general symbolic system is the investigation of the possibility that a purely formal, yet non-axiomatic, theory might nevertheless have a structure which yields predictions which agree with experiment. Such a program may be thought of as follows: If it fails, then this might be taken as a “proof” that experiment is needed in order to “cause” physical theory to correspond to experience. If it succeeds, however, then it is difficult to see how such a cause may be asserted to be necessary: Occam’s Razor argues against such a supposition.
As already noted above, the presumed formal characterization of empirically relevant symbolic structure is that of the finiteness of such symbolic structures. As the formal rules of the general symbolic system are taken to have been completed in Chapter 6, the first six sections of Chapter 7 are dedicated to exploring such special structures within the general symbolic system. This investigation leads to novel derivations of the formalisms of both General Relativity and Quantum Theory, as well as unsurprising characterizations of each as corresponding, respectively, to invertible and non-invertible manipulations of symbols. In so far as the empirical results of all of modern physics is conventionally taken to be encompassed within these formalisms it might perhaps also be concluded that the general symbolic system, in conjunction with the condition of structural finiteness, has succeeded to the same extent as empiricism in yielding an accurate determination of the nature of experience.
In the seventh section of Chapter 7 it is argued that the nature of the derivations of the Relativistic and Quantum-theoretical formalisms indicates that, if such is possible, all finite structure may be represented within either of these complementary schemes, though, as is characteristic of the non-axiomatic method, such a conclusion should not be thought to have been proven to be mandatory. The next section of Chapter 7 inclines, in fact, to the conclusion that such finite structures are formally incompatible. Resolving this conflict presumably requires either the consideration of non-finite formal structures or resorting to empiricism.
As the non-axiomatic method is conceptually consistent with adopting non-finite formal structures, the feasibility of adopting this first resolution of the above difficulty is investigated in the last section of Chapter 7. These considerations lead to the plausible conclusion that the Relativistic and Quantum-theoretical formalisms may be extended, in a mutually consistent manner, by the methods of Non-Standard Analysis. In particular, the formal manipulations of the general symbolic system are to be governed by both of the formalisms of Relativity and Quantum Theory, the only difference between this and the finitary case being that the symbols are no longer required to be finite in general. Established non-standard methods are then used to illustrate the relationship between this extended formal system and the usual finitary formalism which is thus embedded within it.
As a final observation, the role which Non-Standard Analysis plays in presumably reconciling Relativity and Quantum Theory as parts of the general symbolic system suggests that an identification may be made between the notion of the externality of the “physical reality” which finite empiricism forever appeals to and the presumed conceptual inadequacy of any finite structures within the general symbolic system. If this identification is made then, it is proposed, these two approaches may be thought of as being indistinct. While it would be inconsistent with the very notion of a non-axiomatic theory to demand such an identification, nevertheless it is hoped that such an identification might make rejecting serious consideration of such systems less likely.
## Chapter 2 A SCHEMATIC OVERVIEW OF PHYSICS
### 2.1 Introduction
The introductory chapter has provided a general orientation and motivation for the approach to be taken in this paper. The primary question raised by the discussion there is the necessity of language referring to an external “reality” even in such an objective discipline as Physics. Before proceeding to the formal development of the theory it will be helpful to re-examine, in a practical and more detailed way, the functioning of the scientific method, including the actual style of activity of both individual kinds of physicists and some of the historical development of Physics itself. It is hoped that then, after having more comprehensively illustrated the dubiousness of there having been a clear distinction between formalism and “realism” or empiricism, it will be possible to proceed more convincingly with the formal development to follow.
### 2.2 The Ideal Relationship Between Theory and Experiment
The growth and functioning of modern Physics can be understood in terms of a symbiotic relationship between two general activities of physicists: That of theory construction on the one hand, and that of observation on the other.
Theoretical physicists make predictions as to what will happen next in a given situation, and thus explain such events whenever they’re successful. Theoreticians also attempt to encompass as many conceivable situations as possible within as few distinct predictive schemes, or theories, as possible. Physics thus progresses whenever theoretical physicists are able to construct a theory to explain phenomena not previously understood, to predict new phenomena, or to “unify”<sup>1</sup><sup>1</sup>1 Hereafter the convention is adopted that the quotation of a word or phrase which is used, in context, in a definite sense will, if applicable, be presumed afterwards to indicate the definition of that word or phrase. old theories by creating a single theory which explains what was previously explained by two or more old ones. This activity yields a kind of unified understanding of the world.
Theoreticians can’t claim to have ever explained everything, nor can it reasonably be claimed that all of their predictions have been tested or confirmed. Apparently, the world is too uncooperative for that. There has thus been a definite role to be played by physicists who don’t merely predict, but who also observe.
Some of the terms used to describe these observations repeatedly occur in varying combinations. The recurrence of such “events”, which usually have relatively simple descriptions, apparently gives an opportunity for an exhaustive and methodical exploration of a part of experience. This exploration is undertaken by the experimental method.
By adopting conventions for the terms to be utilized in description it becomes possible to compare observations and thus to engage in “measurement”. The recurring conditions which define events also identify parts of description which refer to “measuring devices”. Measurement, by organizing descriptions in a conventional form, serves to help determine the extent to which theory succeeds, and even the failure of theory to explain a measurement still represents an advance for Physics, for it increases knowledge of the world and gives new challenges to the theoreticians.
It should be clear, then, that theory and observation reciprocally benefit the understanding which physicists aim for. Theory, as confirmed by experiment, serves as a basis for prediction and specifies and clarifies what is known, while observation identifies what is yet to be understood whenever unexplained measurements arise. In any event, theory cannot fruitfully deny the observations made in experimentation.
### 2.3 A Realistic Examination of the Behavior of Scientists
While, as just indicated, there are certainly mutually beneficial aspects to the relationship between theory and experiment which, in a manner of speaking, drive the engine of Science forward, there are also conflicts which arise in their interplay. This is, in part, because it is often difficult to recognize and accept when a previously successful system of concepts can no longer be relied upon when attempting to understand new experiences. It may even be said that theory preconditions what is noticed. There is a consequent on-going tension between groups of physicists, each comprising a school of thought, when their theories are faced with the extensive variety of experimental data. This tension may be seen in the debates between such schools of thought.
For this reason the sketch of the activities of physicists made in the last section, in stressing the mutually beneficial aspects of the relationship between theory and experiment, bears some further development. To this end some further general remarks about theory will be made, and it will be shown that physical debate has taken a predictable form.
In a physical theory prediction is a strictly mathematical operation: Given certain mathematical objects representing what is known, in the form of initial and boundary conditions, parameters, and physical constants, the theory then generates, in a well-defined way, other mathematical objects which are supposed to tell something of what is to happen later, in the future. The predicted future, in being derived from what is currently known, cannot be any more precisely defined than what is already known. Now the transformation by which the known thus yields predictions may or may not be invertible. If it is invertible, then prediction may be said to be a sort of relabeling, and physicists would, in this sense, know exactly as much about the future as they know about the present. It is usually otherwise, the future being somewhat less accessible than the present. When a theory’s predictive transformation is not invertible, when the future is somewhat indefinite, then prediction may be said to be “statistical”, by definition.
This consideration helps to explain the way in which experiments are, in fact, interpreted. If a theory yields, apparently, more information about the future than the past, then it may be concluded that some of the description of what is known is superfluous or physically irrelevant. If, alternately, prediction is less accurate than statements of the known, then, it being desirable to know the future as completely as possible, this deficiency may be attributed to the world being essentially statistical in nature, to a defect in the physical theory, or to an inadequacy in the determination of what is known by measurement. Exactly one of these factors must be the culprit, for it would be fruitless and ambiguous to maintain more than one independent cause for the same effect.
Physicists who don’t believe the world to be essentially statistical and who have faith in the physical theory would then assert that there “really” are more comprehensive initial data from which one could make better predictions. They would then advocate attempts to measure such data or assert that statistical methods are resorted to for practical reasons. Such has been the case for Classical Thermodynamics and Statistical Mechanics.
Physicists who accept such a physical theory and the initial data must assert that the world is essentially statistical in nature. This has been the position of adherents to Quantum Theory as exemplified by Niels Bohr’s Copenhagen Interpretation. Bohr, while not denigrating all of the results of the non-statistical General Relativity Theory, nevertheless rejected its methods and maintained the adequacy, or “completeness”, of the statistical Quantum Theory.
If a physicist accepts the adequacy of the initial data but rejects there being any essential indefiniteness in the world, then they are compelled to fault a physical theory which makes statistical predictions. This was Einstein’s reaction to Quantum Theory. Einstein couldn’t easily accept the idea that “God plays with dice”, and so he spent his last years searching for a generalization of his relativistic theory which could account for quantum phenomena without necessitating an essential randomness in the world. Such a theory would then replace Quantum Theory<sup>2</sup><sup>2</sup>2 For Einstein’s presentation of his physical beliefs see . For an independent and in-depth appraisal of Einstein’s realism, see ..
Most physicists today fall into one of these last two groups, but neither camp can claim a definitive victory so long as Quantum Theory and General Relativity can each claim unique successes and still resist unification. Each position was advocated in the famous Bohr - Einstein Debate. It is commonly believed that Bohr won this debate, and Quantum Theory enjoys a corresponding popular approval among physicists, but this is a position maintained by practical success and not, as indicated above, a position proven to be unassailable. In fact there still remains controversy over the interpretation of Quantum Theory, and so Relativists may remain “convicted, but not convinced”.
The notion of a statistical prediction has led to an overview of physical procedure and been useful in comparing the canonical schools of thought in contemporary physics. It has done even more: It has convincingly shown that each school of thought is offering a different solution to the same problem of the meaning of physical theory, and that the debate between these schools has taken a predictable form.
### 2.4 A Criticism of the Role of Experiment
The preceding analysis has highlighted the difficulty in uniquely interpreting the relationship between theory and experiment, and the consequent disputes that have arisen. Now experiment is merely a particular realm of experience to which language may address itself. In fact, according to Niels Bohr <sup>3</sup><sup>3</sup>3 See pg. 3 of reference .; “…by the word ‘experiment’ we can only mean a procedure regarding which we are able to communicate to others what we have done and what we have learnt.” Thus language must, in general, possess the same ambiguity of interpretation as Physics if experimentation is still uncritically referred to. Before carelessly accepting such an implication it would be best to critically examine the role of experiment in physics. This will be done in the following and it will lead to some interesting conclusions.
The defining virtue and test of Science is its ability to predict, where prediction, properly speaking, is a matter of telling what will happen in an experiment that has never been done before. Such a prediction can only be based on something which is presumed independently of the experiment actually being performed. Such prior knowledge comprises physical theory and may also be taken to define “reality” itself.
In order that Science may claim permanent relevance it is necessary that the supply of new experiments be inexhaustible, and thus infinite in number. It follows that reality may, in this sense, be said to be infinite as well. It would contradict the definition of reality, however, if observations were to be said to create reality: “Acts of observation”, in so far as such a term is used at all, must correspond to entirely unpredictable events.
Because a physical theory and the corresponding notion of reality must be presumed in order to make predictions, it is impossible, especially by experiment, to prove its correctness. It may, however, be “falsified” <sup>4</sup><sup>4</sup>4 For a detailed discussion of this notion see reference . by its inability to make acceptable predictions.
Consider, for example, the experimental results that are in accordance with Newton’s Theory of Gravitation. These same results are in accord with Einstein’s Theory of Gravitation as well, so that the same experimental data have received two quite distinct theoretical explanations. From an empirical standpoint, Einstein’s theory is distinguished from Newton’s in that it is able to naturally account for additional phenomena which Newton’s theory can address only with great difficulty if at all. Einstein’s theory is in this sense scientifically preferable to Newton’s, although there can be, as indicated above, no empirical guarantee that future observation won’t favor another theory over Einstein’s.
There cannot, therefor, be a direct path from experiment to theory, and consequently no way to empirically determine which individual theoretical assumptions are to be made or rejected in the construction of a theory: It is necessarily the theoretical system as a whole which must confront experiment. This indicates a kind of miracle, for any physical theory, including those that agree with experiment, must then be constructed in a way that has no well-defined connection with experimental results.
Considering the nebulous relationship between theory construction and the results of experimentation, it may well be wondered whether or not an attempt to construct physical theory without reference to experiment could be successful. In fact, if experiment is to be relied upon then it, in being inexhaustibly novel, can never allow for a final physical theory. In other words, a permanent role for experiment implies that physical theory forever be provisionally acceptable at best, and that there can be no place for necessity in Physics.
Conversely, if definite conclusions are desired, then it is ultimately necessary to ignore experiment in the construction of physical theory. This may seem like a heretical idea, but it is the only choice to make if, instead of rejecting the idea out of hand, it is desired to determine if physical theory can be constructed in this way. This is a worthwhile question to resolve, so this presumption will be made at this point.
Summing up then, it is usually assumed that there are two separate realms to be recognized: Physical theory and the reality to which it refers. Experience is then supposed to guide the construction of the physical theories which describe it, but it has been argued that there cannot be any defensible systematic procedure for incorporating the results of experiment into a physical theory. Accepting this conclusion leads to the realization that physical theories must therefor be miracles of independent creation which, nevertheless, do make predictions that agree with experiment, at least provisionally<sup>5</sup><sup>5</sup>5 For an extensive argumentation as to the freedom allowed in theory construction in Science, consult .. This being the case, it must then be possible to create physical theory without utilizing any experimental results whatsoever, and thus to abolish consideration of experimentation as a realm separate from physical theory itself.
## Chapter 3 THE UNIFICATION OF LANGUAGE AND EXPERIENCE
### 3.1 Introduction
In order to follow the course set out in the last chapter it will be necessary, first of all, to indicate how a language which makes reference to experiences need not, in fact, be referring to anything outside itself. This may be done by indicating how the “realistic” way in which language is used may be interpreted in a strictly formal sense. If this is to be done convincingly then a number of intuitively understood terms must be given explicit formal roles within language and the structure which is normally thought to be found within a reality external to language, as indicated in language itself, must be shown to have a strictly formal derivation. The discussion in this chapter will be carried out in conventional language and will be, in this sense, metamathematical. This should not, however, lead to the same difficulties indicated in the first chapter. This is because it is proposed that the formal system to be developed is not an “object language” which is separate from and subordinate to conventional language, but is, rather, equivalent to it. In other words, it is not presumed that the language developed is a formal language which may be intuitively judged in relation to a metalanguage, but instead the language itself is designed to incorporate and symbolize the intuitions of individuals. This is the condition which, as has been indicated, is required by the identification of the axiomatic systems embedded within the general symbolic system with the testimony of individuals.
### 3.2 Individuals and Language
Whatever applies to the relationship between theory and experiment will also affect the more general relationship between language and experience. In particular, the choice just made in the last chapter requires that experience not be considered separate from language itself. Nevertheless, language still contains statements referring to that which is experienced, and so such statements must remain intelligible. Thus statements about experience must be both accommodated and yet denied to refer to a separate category outside of language. This may seem paradoxical.
The unique solution to this difficulty which is adopted here is to identify the recognized significance of any statements made by individuals about their experiences with these statements themselves, so that such statements may be taken be a natural development within the language itself. Then all statements may be made in the form of references to experiences and yet no such separate realm need be recognized. Any “separation” will be a feature of the language itself. This identification has a number of consequences.
As all statements in language might be the assertions of some individual at some particular instant, it would seem to follow that all statements must have realizations in experience. However, because such statements, being merely strings of symbols, can only be judged according to their formal properties, it follows that while it may be intuitively appealing to base the considerations to follow on realistic explanations and examples, such an approach may be misleading, for it is easy to forget that such examples are, themselves, only to be judged according to their formal properties. The developments to follow will, accordingly, be kept, so far as seems reasonable, strictly formal and abstract.
It is evident that great care must be taken in dealing with statements about experience from an individual perspective. The delicate matter of constructing a strictly formal symbolic system which nevertheless accommodates the testimony of individuals is the task taken up next.
### 3.3 Perception and Meaning
The construction of the symbolic system starts, as indicated above, from the perspective of a single individual. It is assumed that the reader is an individual and therefor may understand the meaning of the terms “reader” and “individual” without a definition being provided.
It may be said that the symbolic system is to be utilized by this individual for the purpose of comprehending his experiences. In keeping with the individual perspective it may also be said that symbols may be utilized in the description of experiences at a particular instant. The developments to follow will be based on the consideration of the description of experiences at a particular instant, and the symbols will be said to identify what is experienced. The sense, if any, of these terms beyond their purely formal utilization is left to the reader’s insight.
Description will be provided by the combination of written symbols. The assignment, by an individual, of certain symbols to a particular instant will be said to comprise an act of “perception.” Such symbols will, in virtue of this assignment, be said to have been made “meaningful” to this individual. Because a given symbol either definitely does, or does not, appear in a given statement, it follows that a given symbol is, or is not, meaningful to an individual at a particular instant. Moreover, not all symbols need be meaningful.
Because written symbols may themselves be said to be experienced, it necessarily follows that perception may result in symbols being assigned, by an individual, to other symbols. This “meta-labeling” provides a means of inter-relating symbols because any symbols assigned to the same symbol are, in virtue of this very fact, grouped together. Such symbols may be said to be “associated.”
The association of symbols just mentioned raises the question of whether or not the development of the symbolic system can give anything distinct from set or class theory. In these theories the “existence” of sets or classes is required so long as they don’t violate any of the formal axioms of the theory. Such “existence” has not been mandated here, for the written strings of symbols are thus far presumed to correspond to testimony, and be merely given, and are not required by some formally stated rules to be formed. So long as this remains the case such an identification would be unjustified. This point will be returned to later, in the derivation of physical theories, where a condition for the existence of associations between symbols will be given.
### 3.4 An Important Illustrative Example
The meaning of a particular symbol is, at best, a necessarily private matter. It is not generally possible to say “what” a particular meaningful symbol means. It is only possible to note that certain symbols are consistently assigned to others. This situation may be made intuitively as well as formally clear with an example involving two hypothetical individuals, Joe and Tom, say, who may be supposed to refer to certain symbols common to both of them.
Let it be supposed that the symbols “stop signal” and “apple” may be referred to by both Joe and Tom, and that each indicates to the other that he takes these symbols to be meaningful. Now “stop signal” and “apple” may be, for the sake of argument, assigned in Joe’s mind to an experience identified as “A”, while “go signal” is assigned in his mind to the experience indicated by “B”. It is not generally presumed that an individual’s mind may be read like this, but this fiction is indulged in just to make a point.
Even supposing that Tom shares the same experiences to which “A” and “B” refer, it is perfectly admissible for Tom’s corresponding experiences to be the reverse of Joe’s. Joe and Tom could then agree in saying that the symbols “stop signal” and “apple” both refer to things which are “Red” and other things, such as the “go signal” may be agreed to be “Green”. However, Red would then be assigned to the experience A in Joe’s mind and to the experience B in Tom’s. The reverse would be the case for Green.
Any system of symbols corresponding to experience will accommodate such examples, and will therefor share in this indefiniteness as to the meaning of particular symbols. If, then, a symbol is supposed to have a definite meaning, then this meaning must be established solely within an individual’s mind. It is possible, however, to establish definite relations between symbols. In the above example, in fact, the convention was reached that the meta-label Red was to have the symbols “stop signal” and “apple” assigned to it.
The categories arrived at, such as Red and Green above, can be useful conventions which are stipulated according to some symbols having some meaning in common. In this case the red things have, according to both Joe and Tom, common meanings. Joe takes red things to share the meaning A, while Tom sees red things as sharing the meaning B. Again, this does not mean that red has the same meaning for Joe and Tom.
Such categorizations are always strictly provisional: there are no final categories and there can be no mention of necessary truths. In the above case it can never be asserted, for example, that there cannot be a symbol, such as “Fire Truck” which is not linked in Joe’s mind to the experience A. This may be because “Fire Truck” isn’t meaningful to Joe. It may be the case that Tom wouldn’t assign “Fire Truck” to the experience A either. It may also be the case that Tom may never have the experience A and so never be able to understand why Joe associates “apple” and “stop signal” as being “red” but doesn’t group “Fire Truck” under this metalabel. The usage of a particular symbol, such as red, can only be justified and stipulated solely by individuals, so that it can never be asserted that a particular symbol need apply to a particular experience. Distinct individuals may merely agree or disagree in their assignment of symbols.
The most precision with regard to meaning that may be expected of language, then, besides its including individual symbols of indefinite meaning, is that it also provide for the construction of provisionally relevant relations or groupings of symbols. The meta-labeling system discussed so far differs from the system of the usual formal logic in that it is not based on relations with a fixed and known number of “arguments”. Meta-labeling allows just such definite but flexibly chosen connections between symbols to be established in writing as to give partial exposition of the relations it represents. The specification of which symbols are associated by a given meta-label identifies a definite structure without, as is necessarily the case, making the meaning of any symbol evident. Language, in not being reliably able to convey a separate meaning, is thus seen to be strictly “formal” in nature.
### 3.5 The Development of Language and the Notion of Truth
How are languages to be judged? If the way in which symbols are combined in writing, and which combinations are to be allowed, is established once and for all, then it may be said that a particular well-defined language is in use. It otherwise is possible to arbitrarily amend the relations expressed by the language and, in this sense, the language will have been changed. Given various symbols which are connected, by meta-labeling, in definite ways, it is generally possible to give definite rules, called the “grammar”, for how such symbols may be combined in written expressions. Expressions which satisfy the rules may be said to be “correct.” It is thus necessary that a well-defined language have a grammar.
If given more than one language, how may they be compared? Certainly the comprehensiveness of a language is a primary consideration. Can a language be shown to be comprehensive?
If a language is developed to the point where several symbols may be consistently associated by a meta-label, then this meta-label may be said to identify a particular “thing”, and those symbols which it usually associates may be taken to be “properties” of that thing. Once this may be done multiple times, it may then be said, in a strictly formal sense, that these certain things “exist” to which the language refers.
Properties may be assigned to different things, although, as indicated before, such categorizations need not be permanent nor are they necessary. Assertions, relative to these adopted conventions, may then be made in the language, these assertions taking the form of indicating that a thing does or does not have a particular property assigned to it. The assertion, by an individual, that a thing has a particular property is, like all perceptions, an action which can have no necessary external justification. Even the individual himself may not provide a reason for a thing to have a property. The property is simply a part of the perception of the thing itself.
The discussion above, in showing that the notion of an assertion or “proposition” may be derived from the metalabeling which perception results in, should help to make clear the relationship between the theory herein developed and the approach of the school of thought known as “Logicism”. For, though the formal propositional apparatus from which Logicism proceeds has been shown to arise within the theory, nevertheless it must not be thought that the theory is equivalent to such an approach.
A difference between the two approaches may, in fact, be discerned by considering the old philosophical issue of the discussion of a particular object such as, for example, a chair. Whereas a logicist would, by the very nature of his approach, be constrained to assert that a chair does, “objectively” speaking, have certain properties, the approach espoused here starts from describing the perceptions of a given individual, and therefor does not dismiss, at the outset, the possibility that individual perceptions may differ. Considering a case of discussion of the color of the chair in which some individuals “are” color blind should make clear the primary role of convention, in which individuals “agree to agree”, as opposed to objective truth, and therefor also indicate an advantage of the present theory over Logicism. While, failing such agreement, effective discussion might be impossible, nevertheless undercutting the possibility of such disagreement can only be justified by arbitrarily dismissing some perceptions in favor of others.
Despite the above observations, if the descriptive system is greatly developed then it may contain many things and many assertions about them. Such a sufficiently developed system might exhaust all of the written statements any individual is likely to make in a lifetime. Such a system, in virtue of its complexity, might be deemed comprehensive, but, it may be asked, does it “tell the truth”?
Because experience, as a realm separate from language, has been abolished from consideration, there can be no way to judge a language apart from its formal properties. Thus there can be no external criterion on which to base an assertion of the truth-value of a statement. The only possible assertion of the adequacy of a language that may be sought, then, is in arguing that the language is, in its own nature, capable of representing any written relation. It is also worth noting that in rejecting the notion of truth there can consequently be no worry about consistency, nor can Gödel’s Incompleteness Theorem necessarily apply.
The way in which statements “about experience” can be made in a language without the necessity of an external realm of experience has been discussed. It has been found that the notions of things, properties, and assertions about things can be given purely formal characterizations. The notion of intrinsically true statements must, in the context of a non-axiomatic theory, be abandoned, but it is still reasonable to talk about the descriptive adequacy of a language.
## Chapter 4 BEGINNING CONSTRUCTION OF THE SYMBOLIC SYSTEM
### 4.1 Introduction
The last chapter addressed the way in which experiential statements may be given purely formal interpretations. There certain terms were introduced and certain arguments made, but it was not addressed how various formal manipulations were to be decided upon. Such a formal development of the general symbolic system is begun in this chapter.
### 4.2 Reserved Symbols and the Formal Nature of Language
It is not possible, as has been demonstrated, to in any way prejudge which symbols may play a role in individual perception in the general symbolic system. If the identification of meaningful symbols is left entirely to the individual, then the symbolic system must be such that, for any correctly written expression, each potentially meaningful symbol may be replaced by another to yield a new expression which may take on, according to individual perception, the same meaning as the last. This may be thought of as a recoding of the previous expression. This recoding will be the first aspect of the symbolic system to be specified.
In the written symbolic system symbols will be distinguished simply by their appearances as written; ’a’ is, for example, distinct from ’b’ in this sense. This will be indicated by writing a$``$b. It is important to note that no “equality” has been defined here.
The convention will be adopted that writing a$``$b will indicate that the symbol ’a’ is to be replaced by ’b’ whenever it appears in a written expression. Suppose, for example, that a$``$c and b$``$d. If this is done in such a way that if any pair of images, such as ’c’ and ’d’, are such that c$``$d then the corresponding pre-images, such as ’a’ and ’b’, are also such that a$``$b, then the replacement effected will be said to be a “relabeling”, it being presumed that relabelings are defined on all symbols.
It will never be the case that x$``$x for any symbol ’x’, and there is no requirement that an image and its corresponding preimage be distinct. Thus a given symbol may be retained under a given relabeling and, in fact, under all relabelings considered. The symbolic system thus allows that certain symbols may be used in the same way in all expressions in different relabelings. In this way it is possible to always retain particular symbols to play fixed roles in the symbolic system, but it has not been necessary to stipulate the “existence” of symbols as necessarily given: Such “reserved symbols” may be conventionally set aside for such use. Both ’$``$’ and ’$``$’ will be taken to be reserved symbols. Consequently, given a relabeling in which a$``$b and c$``$d, an expression such as a$``$c will be replaced by b$``$d, it being presumed that $``$ $``$ $``$ by convention.
It can be seen that having reserved symbols is the key to the adoption of a conventional language in which symbols are always used in the same way. Ironically, then, it is the very indefiniteness of the meaning of any symbol which will allow a language with rigid rules of combination, with a definite grammar, that is, to be constructed. In other words, the meaninglessness of language will be what allows a definite “formalism” to be used.
### 4.3 Partial Order and Formal Structure
In the last section it was found that the indefiniteness of the meaning of symbols provided a key to introducing reserved symbols which are adopted conventionally and allow for the possibility of a definite language being used. It will now be shown that the meta-labeling between symbols engenders a certain structure which may be explicitly indicated in the language through the use of reserved symbols indicating an order relation between symbols. The meta-labeling also will provide a definite means of combination of individual symbols, thus leading to the notion of an “expression”.
Consider the assignment of one symbol to another. In the example involving Joe and Tom both ’apple’ and ’stop signal’ were assigned to ’Red’. It may also have been the case that some symbol had been assigned to ’apple’, and some other symbol assigned to that, and so on. This chaining of the connections between symbols motivates a definition.
It will be said that ’x’ is a descendent of ’y’ iff
either ’x’ is assigned to ’y’
or ’x’ is assigned to a descendent of ’y’.
Obviously this is a recursive definition which is not constructive in the sense that a definite end to the recursion is not to be expected for any reason formally indicated. The assignment of symbols is not, however, something to be questioned: It may be said that it is contradictory to presume that individuals could stipulate such connections in an act of perception and not be able to identify any chain of assignments which would identify descendents of a given meta-label.
With this definition of descendents it is then possible to construct an order relation between meaningful symbols. It will be said that ’x’ is “less than or equal to” ’y’, written x$``$y, if ’x’ is a descendent of ’y’ or if it is ’y’ itself. ’x’ will be said to be “equal to” ’y’, written x = y, if both x$``$y and y$``$x hold. Clearly, if x = y then y = x. If x$``$y but x = y doesn’t hold, then x$`<`$y will be written. It will also be convenient to adopt the following conventions: x$``$y may also be written as y$``$x, and x$`<`$y may be written as y$`>`$x. The symbols $``$, $``$, $`<`$, and $`>`$ are taken to be reserved symbols.
Only the structural relationships indicated by meta-labelings can have any significance within the formal system, and these relationships may be conveyed by the reserved symbol $``$. Suppose, for example, that a$``$b, a$``$c, and b$``$d. Then (a$``$b)$``$(c$``$d), so that the meta-labeling structure between the pre-images ’a’ and ’b’ is inherited by the images ’c’ and ’d’. For this reason it may be said that the images are used “in the same way” as the corresponding pre-images.
Because, again, only the meta-labeling structure can be of formal significance, and because the symbol ’=’ indicates an “equality” based on meta-labeling, it follows that any two symbols ’x’ and ’y’ such that x = y may be substituted for one another anywhere within an expression. Thus meta-labeling is the source of “the principle of substitution.”
With these definitions it is clear that $``$ satisfies the standard definition of a “partial order”:
$$\text{x}\text{x for every x.}$$
(4.1)
$$\text{x}\text{y and y}\text{z implies that x}\text{z.}$$
(4.2)
$$\text{x}\text{y and y}\text{x implies that x = y.}$$
(4.3)
If, for meaningful x and y, either x$``$y or y$``$x, then x and y will be said to be “comparable”. Any system of meaningful symbols, all of which are comparable, will be said to form a “chain” and it may also be said that the order $``$ is then “total” for these symbols.
Along with the notion of a partial order there will also be an associated notion of a “bound”. Given two meaningful symbols x and y, any symbol w such that w$``$x and w$``$y will be said to be a “lower bound” of x and y. There need not be any lower bounds of x and y, but there may also be many. Among such w there may be at most one unique “greatest” one, denoted by x$``$y, such that e$``$x$``$ y for any lower bound e.
It will be noticed that the string of symbols x$``$y stands for a single symbol determined by meta-labeling with respect to the two symbols x and y. Thus x$``$y is to be considered as a kind of unit. In order to make this clear it will be convenient to introduce the reserved symbols ’(’ and ’)’ and write (x$``$y).
There will be notions of an “upper bound” and for a “least upper bound” for meaningful symbols which are defined in correspondence with the analogous notions above. A unique least upper bound of x and y will be denoted by (x$``$y), where $``$ is taken to be a reserved symbol.
Both of $``$ and $``$ act as a means of combining, according to meta-labeling, several individual symbols into one unit. Because, for example, (a$``$b) is considered as a unit, it makes perfect sense to consider a compound such as (a$``$b)$``$c. Such a compound, combining several individual symbols, will be called an “expression”. Expressions will be primary constituents of the written symbolic system.
The above considerations have introduced $``$ and $``$ and found their significance as a means of the formal combination of individual symbols into expressions. It has also been shown how the meta-labeling involved in perception may be explicitly indicated by using the reserved symbol $``$. These are fundamental steps in the development of the formal system, but there remain further developments to consider.
### 4.4 A Preliminary Consideration of the Formal System
Prior discussion has indicated that the general symbolic system to be developed here cannot admit of the notion of truth nor proceed on an axiomatic basis: Its most concrete realizations of relationships can only be expressions of contingency. While no axioms or truths may be embraced, it would come to the same thing if any were explicitly rejected either. Thus this system must be formally compatible with the structure of any axiomatic system. This observation justifies summarily citing and applying the results of a particular axiomatic theory in so far as the general symbolic system cannot have formal rules for the manipulation of symbols which are incompatible with it. This state of affairs will be taken advantage of in what follows.
A symbolic system with an associated partial order $``$ and the operations $``$ and $``$ may be viewed on the basis of the usual Set theory, and then it comprises what is commonly known as a LATTICE . If attention is restricted to only meaningful symbols then the standard developments of Lattice Theory will provide necessary specifications of the symbolic system. For such a structure it can be shown<sup>1</sup><sup>1</sup>1For the Lattice-Theoretical results cited in this section see pp.1-17 of reference . that the following results, which, for convenience, are cited as they appear in the above reference, must hold:
L2. x$``$y=y$``$x and x$``$y=y$``$x. (COMMUTATIVITY)
L3. x$``$(y$``$z)=(x$``$y)$``$z and x$``$(y$``$z)=(x$``$y)$``$z. (ASSOCIATIVITY)
CONSISTENCY: x$``$y is equivalent to imposing either of x$``$y=x and x$``$y=y.
THEOREM 8: Any system with two binary operations which satisfies L2, L3, and CONSISTENCY is a Lattice, and conversely.
Thus L2, L3, and CONSISTENCY completely characterize Lattices. It must be the case, then, that at least one of these conditions must be violated in order that a symbolic system not comprise a Lattice.
Further results for Lattices:
Lemma 3: If y$``$z then (x$``$y)$``$(x$``$z) and (x$``$y)$``$(x$``$z).
Lemma 4: (x$``$y)$``$(x$``$z)$``$x$``$(y$``$z).
Suppose that O is such that, for every meaningful symbol x, O$``$x holds, and that O is defined identically. It follows that O$``$O$`{}_{}{}^{}`$O, so that O=O. Thus any such symbol O, and likewise any symbol I such that x$``$I for any meaningful x, must be unique. Such O and I will be universal bounds for meaningful symbols.
O and I have been introduced, in part, in order to show that the $``$ in Lemma 4 cannot be strengthened to an equality. This can be shown by considering the system of inequalities O$``$x, O$``$y, O$``$z, x$``$I, y$``$I, and z$``$I. Then x$``$(y$``$z)=x$``$I=x and (x$``$y)$``$(x$``$z)=O$``$O=I, where it’s not the case that x=I. The equality in Lemma 4 may hold in some Lattices, however. Such Lattices will be called “DISTRIBUTIVE”.
THEOREM 10: In a distributive lattice, if c$``$x=c$``$y and c$``$x=c$``$y, then x=y.
Theorem 10 provides something like a cancellation law for distributive lattices.
Consider a distributive lattice with O and I. From the given lattice a new distributive one may be constructed which extends the original lattice by including, for each x in the original lattice, a corresponding new symbol x. This new symbol is to be defined solely by the requirement that O$``$x$``$I and O$``$x$`{}_{}{}^{}`$I. In such a lattice it will be the case that x$``$x=O and x$``$x=I. Such an x will be called the “complement” of x. According to Theorem 10, a symbol x can have at most one complement in a distributive lattice. A lattice in which each symbol has a complement will be called a “complemented” lattice. A complemented distributive lattice will be said to be “BOOLEAN”.
These considerations have been preparation for
THEOREM 16: In a Boolean Lattice:
L9. (x)=x, and
L10. (x$``$y)=x$`{}_{}{}^{}`$y.
A final lattice theoretical result will be of central importance here:
THEOREM 7: Any chain is a distributive lattice.
The above survey of Lattice Theory describes a formal structure which may be generated according to an individuals metalabeling, so the general symbolic system must therefor be able to somehow embody such a structure. It must be stressed, however, that the general symbolic system has not been shown to be a lattice. The reason for this distinction will be given in the next section.
### 4.5 The Necessity of Generalizing Lattices
With the last result cited in the last section it is appropriate to ask whether or not lattices, generated according to the order $``$ and the operations $``$ and $``$, are suitable structures for the general symbolic system. It is to be expected that they are not, for, in addition to being based upon metalabeling structure, lattices stipulate that such metalabeling structures must have least upper and greatest lower bounds while such requirements cannot be part of a non-axiomatic theory.
The selection of meaningful symbols, and of the meta-labeling between them, is almost entirely unrestricted and left up to the individual. From this it follows that any formal rules for symbols which may be adopted for the symbolic system must apply whatever the meaningfulness or order relations of the symbols may be. Now it is possible for any three particular symbols to form, according to individual perception, a chain, so that any acceptable formal rule must apply as if the whole symbolic system is a lattice and, in fact, a chain. Theorem 7 then indicates that such a lattice must be distributive.
The lattice structure is the minimal embodiment, as demonstrated by Theorem 8, of the partial order $``$ and the associated operations $``$ and $``$ which are generated by unrestricted meta-labeling with least upper and greatest lower bounds. Because, as has been shown, free meta-labeling is precisely what may and must be conveyed by language, it is unacceptable to have a further restriction placed upon such a lattice structure. Thus the general symbolic system must be able to accommodate lattices in general, and non-distributive lattices in particular. It may be noted, in passing, that the necessity of not restricting the general symbolic system to Boolean lattices decisively distinguishes it from the mathematical methods acceptable to the Intuitionists.
The above considerations lead to the conclusion that the general symbolic system must be distinct from, and yet embody, lattices formed by meaningful symbols. This may only be accomplished by embedding such lattices in a distinct formal system. It will then be necessary that this formal system itself have operations which are analogues of the operations $``$ and $``$, and that these analogous operations satisfy a distributive law.
In order to justify the program pursued in this thesis it is necessary that the formal system may consistently reflect the “intuited” combinations of meaningful symbols perceived by individuals and yet also obey formal rules which apply generally. Thus, for any manipulations in the formal system, it will still be necessary to be able to identify corresponding meaningful expressions and, in this way, be able to assert that the system has yielded its own interpretation. Now there will generally be more than one expression in the formal system corresponding to a given meaningful one, for since a lattice of meaningful symbols may be embedded within the formal system, it would otherwise be formally identical to such a lattice. It necessarily follows that the general symbolic system has a “formal redundancy”: The process of “interpretation” whereby meaningful expressions are extracted from expressions in the formal system is generally a many-to-one transcription.
In conclusion, the only apparent shortcoming of the choice of the lattice structure as the general symbolic system lies in its inability to maintain both the freedom necessary in meta-labeling and obedience to generally applicable formal rules. Nevertheless, because the general symbolic system must embed the lattice structure, a unique and definite means of specifying the structure of the formal symbolic system has been indicated. This consists in generalizing the formal rules for a distributive lattice while explicitly retaining those lattice-theoretical results which don’t depend on which symbols are meaningful and therefor fall into the order $``$. Other results may not be taken over so directly, but may still suggest extensions. The task of this generalization will be taken up in the next section.
### 4.6 The Symbolic System as a Division Ring
The general symbolic system will be constructed in a very nearly direct analogy with the structure of lattices found above, though it must be stressed that what is pursued here is merely the specification of formal rules for symbolic manipulation which are not inconsistent with and yet are sufficient for the embodiment of the lattice structures discussed above. In particular, the general symbolic system will necessarily begin with two reserved symbols, referred to as operations, denoted by +, called “addition” and corresponding to $``$, and $``$, called “multiplication” and corresponding to $``$. The symbols $``$ and $``$ will still be retained in expressions whenever all non-reserved symbols are those that would result from the interpretation of an expression in the general symbolic system, and their appearance will indicate this state of affairs.
In the following, certain formal operations will be defined in terms of allowed substitutions. Recall that the “equality” of symbols, indicated by =, justified the principle of substitution. The formal symbol = will, correspondingly, be taken to be an extension of the equality previously defined by metalabeling. Thus = will always justify substitution. Another note: The presentation will parallel that of a typical axiomatization of a mathematical structure, but it is important to remember that no notion of “existence” is to be inferred. This point will soon be stressed with an observation relating to an inverse operation for $``$.
Recall that x$``$y was considered to be a symbol which corresponded in a particular way to the symbols x and y. + and $``$ will be defined to have the same formal property, so that the “combinations” that + and $``$ yield may be represented by a new symbol.
Given the combination of symbols a and b indicated by a$``$b, this same combination may be written as a new symbol c, where
c = a$``$b. (4.4)
Given the combination of symbols a and b indicated by a + b, this same combination may be written as a new symbol c, where
a+b=c. (4.5)
The requirements (4.4) and (4.5) indicate that the operations + and $``$ are, in the usual terminology, “closed”.
The associative property, specified above in L3, applies in all lattices and must therefor be taken over directly. Thus the following substitutions are permitted:
$$(ab)c=a(bc)$$
(4.6)
$$(a+b)+c=a+(b+c)$$
(4.7)
The distributive property will, as indicated in the last section, have to be imposed. Thus the substitutions:
$$a(b+c)=(ab)+(ac)$$
(4.8)
$$(a+b)c=(ac)+(bc)$$
(4.9)
It is necessary to give the second rule, rule (4.9), because a$``$b=b$``$a need not hold in general.
It will be convenient, for the sake of referring to some ideas in the usual mathematics, to revert to language using the notion of existence for a moment, this being done solely for the sake of motivation of a formal substitution rule. In this spirit the cancellation property for distributive lattices, which was cited in Theorem 10, leads, as a trial generalization, to the investigation of the following standard group-theoretic rules:
For any symbols a and b there are symbols x and y such that
a$``$x=b and y$``$a=b. (4.10)
For any symbols a and b there are symbols x and y such that
a+x=b and y+a=b. (4.11)
It is easily shown that rules (4.4),(4.6), and (4.10) lead to the availability of a unique symbol “1”<sup>2</sup><sup>2</sup>2 See pg. 201 of . , the “multiplicative identity”, such that
For every symbol a, a$``$1=1$``$a=a. (4.12)
In a similar fashion, rules (4.5),(4.7), and (4.11) lead to a unique symbol “0”, the “additive identity”, such that
For every symbol a, a+0=0+a=a. (4.13)
It can, furthermore, be shown that 0$``$a=a$``$0=0 for any symbol a. With this in mind it is clear that rule (4.10) cannot be maintained.
Returning now to the consideration of substitution rules, it is clear that the symbols “0” and “1” may be reserved, and the rules (4.10) and (4.11) are to be replaced by the following rules:
For any symbol symbol a which is not such that a=0, another symbol, a<sup>-1</sup>, may be written such that
a$``$a<sup>-1</sup>=a$`{}_{}{}^{1}`$a=1. (4.14)
For any symbol a another symbol, -a, may be written such that
a+(-a)=(-a)+a=0. (4.15)
The expression a+(-b) is conveniently written as a-b. It is easily shown that (-1)$``$a= -a. a<sup>-1</sup> will be referred to as the “multiplicative inverse” of a, while -a will be called the “additive inverse” of a or the “negative” of a. Evidently, the inversion of either multiplication or addition serves to erase symbols in written expressions. The operation of inversion gives a good opportunity to emphasize that the concept of existence is not to be applied. If, in fact, a were to refer to, for example, a “dog”, then what, if anything, might a<sup>-1</sup> refer to ? The inverse of a dog ?
Finally, it can be shown <sup>3</sup><sup>3</sup>3 See pg. 201 of reference . that a certain substitution rule for addition holds:
a+b=b+a. (4.16)
This property of addition is referred to by saying that addition is “commutative”. It turns out that multiplication is not necessarily commutative<sup>4</sup><sup>4</sup>4 See pp. 92-94 of ..
The rules (4.4), (4.5), (4.6), (4.7), (4.8), (4.9), (4.12), (4.13), (4.14), (4.15), and (4.16), taken together, indicate that the general symbolic system comprises what is generally known as a “division ring” or a “skew field”. It should be recalled that these rules merely specify the way in which symbols are to be formally manipulated, so that, for example, the “existence” of an inverse of a symbol is strictly equivalent to erasure of that particular symbol being allowed. In this sense algebra is here regarded to be a branch of orthography.
This is the extent of the exploration of the formal algebraic properties of the symbolic system that will be undertaken at this point. It remains to investigate the role of the order $``$ in the symbolic system, and this task will be taken up in the next section.
### 4.7 Order in the Symbolic System
The order in a lattice is generated and defined by the meta-labeling inter-relationships between meaningful symbols, and its sole purpose is the conveyance of this information. The general symbolic system is to be a proper extension of a lattice in which not all symbols may correspond to meaningful ones. Since not all symbols, even after interpretation, may be meaningful, it follows that no order can be defined, according to meta-labeling, on all of the symbolic system. This has the obvious consequence that Induction, and with it Proof Theory, cannot, in any sense, be carried out in the general symbolic system. It also forecloses any paradoxes associated with the Axiom of Choice. It will be the case, however, that part of the symbolic system will have an order defined on it: This will be a partial order such as is induced under the embedding of a lattice of meaningful symbols. There is a definite necessity, then, for extending the notion of order to the formal symbolic system.
Those symbols which correspond to meaningful symbols will still obey the laws found before for partial orders:
x$``$x for every x. (4.17)
x$``$y and y$``$z implies that x$``$z. (4.18)
x$``$y and y$``$x implies that x = y. (4.19)
The uncritical transcription of Lemma 3 would read:
If y$``$z then x$``$y$``$x$``$z and x+y$``$x+z for any x. (4.20)
Can this rule be accepted as it is? Consider a symbol y$``$0. From the second half of (4.20) it follows that 0$``$ -y. Thus, for any meaningful symbol y, exactly one of y = 0, y $`<`$ 0, or -y $`<`$ 0 must hold. Now consider y $`<`$ 0 and x = -1 and apply the first half of (4.20). Then the conclusion would be reached that both y $`<`$ 0 and -y $`<`$ 0, but this definitely cannot be accepted.
It is thus required that x be restricted in the first half of (4.20), so that the condition x $`>`$ 0 is adopted. This results in the adoption of the two rules
If y $`<`$ z and x $`>`$ 0 then x$``$y $`<`$ x$``$z. (4.21)
If y $`<`$ z then x+y $`<`$ x+z for any x. (4.22)
It can now be concluded that the symbolic system is a division ring on part of which an order $``$ obeying the requirements (4.17) through (4.19), and (4.21) and (4.22), is defined. The nature of the symbolic system is still very abstract, but a more concrete picture will start to emerge in the next section.
### 4.8 Finiteness and Some Concrete Examples
Now that the the abstract algebraic and ordinal properties of the general symbolic system have been somewhat developed, it is possible to exhibit some of the usual mathematical objects that should be expected to find a place within the system. These considerations will start with the fundamental notions of finiteness and infiniteness, and will then proceed to the construction of symbols that multiplicatively commute with all other symbols, as well as to the construction of symbols with tailor-made commutation properties. These developments will illustrate the non-trivial nature of the rules thus far adopted.
Consider a symbol s to which various other symbols are assigned. Suppose that x is one such assign, so that x is a descendent of s but not of any other descendent of s. Construct a new symbol s/x such that s/x has the same assigns as s except it omits x. If there is not a one-to-one relabeling from the assigns of s to those of s/x, then s will be said to have “finitely many” assigns. Otherwise, s will be said to have “infinitely many” assigns.
If a symbol s has finitely many assigns, then it will be possible to explicitly arrange them in a terminating list. Given such a list a corresponding sum
m = $`\underset{\text{m terms}}{\underset{}{1+1+1+\mathrm{}+1}}`$ (4.23)
may be derived, where the summands 1 correspond to the distinct assigns of s. s will then be said to have m assigns and 1 will have been added “m times”. Such m will themselves be said to be “finite”, and also to be “Natural Numbers”. Natural numbers serve to “count” the assigns of a given finite symbol.
Besides this elementary construction of the natural numbers, other constructions starting from 1 and using finitely many applications of the operations +, -, $``$, and ($``$)<sup>-1</sup>, but not requiring (0)<sup>-1</sup>, result, in the obvious way, in the generation of the Integers and Rational numbers.
Consider a natural number m and any other symbol q. Then
m$``$q = (1+1+1+$`\mathrm{}`$+1)$``$q = (q+q+q+$`\mathrm{}`$+q)
and
$$\text{q}\text{m = q}\text{(1+1+1+}\mathrm{}\text{+1) = (q+q+q+}\mathrm{}\text{+q)}.$$
so that m$``$q = q$``$m for any q. Thus the naturals commute with all symbols. So too will the integers and rationals.
Consider now those symbols which are generated by finitely many operations +, -, $``$, and ($``$)<sup>-1</sup> and finitely many other symbols x<sub>1</sub>, x<sub>2</sub>,$`\mathrm{}`$,x<sub>m</sub>. The resulting symbols will be rational polynomials in x<sub>1</sub>, x<sub>2</sub>,$`\mathrm{}`$,x<sub>m</sub>. If, as the x<sub>1</sub>, x<sub>2</sub>,$`\mathrm{}`$,x<sub>m</sub> are varied, the same sequence of operations is applied to the new x<sub>i</sub>, then a function will have been established. Thus all of elementary algebraic mathematics is available in the symbolic system.
Besides the development of this algebra, it is also of basic importance to investigate the constructive aspects of the commutativity of the $``$ operation. In this case it has already been constructively established by Hilbert <sup>5</sup><sup>5</sup>5 See pp. 92-94 of reference . that the rules for a division ring with order allow for symbols with tailor-made commutation properties, so that this detail need not be further addressed.
In this section some finite constructions within the symbolic system have been investigated. It has been found that the usual algebraic objects of elementary mathematics are to be found in the symbolic system. The integers and rationals, in particular, provide such examples of symbols which, in addition, commute with every other symbol. It has also been found that non-commuting symbols may, as is necessarily the case, be formally constructed. These considerations have focused on finite manipulations of symbols, but it remains to investigate infinite manipulations. This will be done in the next chapter.
## Chapter 5 THE CALCULUS GENERALIZED
### 5.1 Introduction
The usual presentation of mathematics is developed in terms of finite quantities: Functions and numbers are, for the most part, finite. The Calculus is developed in terms of the concept of a limit which, again, relies on finite quantities. It will be important for the sake of later considerations that this prejudice not be maintained. Then, whenever quantities turn out to be, or are assumed to be, finite, this state of affairs can be explicitly recognized rather than tacitly assumed.
In the sections to follow a general theory of functions, differentiation, and integration will be developed within the general symbolic system which doesn’t presume that the quantities involved are necessarily finite: The formal methods developed can be applied independently of this condition and yet extend the results given by the standard finite version of the Calculus.
### 5.2 Functions
Polynomials, as has been seen, give an example of functions, in the usual sense, in the general symbolic system. It will be necessary to give, however, a general characterization of functions in the general symbolic system which doesn’t rely on their having been finitely generated. The necessity of maintaining the formal structure of statements under relabeling will show that a particular modified notion of a function is required.
Suppose A = a$`{}_{1}{}^{}`$a$`{}_{2}{}^{}`$a$`{}_{3}{}^{}\mathrm{}`$a$`{}_{n}{}^{}_{j=1}^n`$a<sub>j</sub>. If the a<sub>j</sub> are allowed to vary, that is, a relabeling is applied where each a<sub>i</sub> is replaced by a new “argument”, then this product will give a simple example of a function of many variables. Whatever general definition of a function is adopted, it should recognize this product as giving a function.
Let the a<sub>j</sub> vary through, in particular, a relabeling a$`{}_{j}{}^{}`$q<sub>j</sub>. The resulting product must be formally identical to the original and have the same interpretation. Thus each a<sub>j</sub> must vary in the same way, as given by q<sub>k</sub>=w(a<sub>k</sub>), where w is some well-defined rule for relabeling individual variables. It must then be the case that w(A) = w($`_{k=1}^n`$a<sub>k</sub>) = $`_{k=1}^n`$w(a<sub>k</sub>) = $`_{k=1}^n`$q<sub>k</sub>. As a special case this gives w(x<sup>2</sup>) = w(x)<sup>2</sup>; This will determine w.
Choose any v$``$0. Then an equation may be written involving a symbol u, u not corresponding solely to x, such that w(x)$``$v = u$``$x. Now w(x<sup>2</sup>) = u$``$x$`{}_{}{}^{2}`$v<sup>-1</sup> = w(x)<sup>2</sup> = u$``$x$``$v$`{}_{}{}^{1}`$u$``$x$``$v<sup>-1</sup>, so that u = v. Thus relabelings w<sub>u</sub> are given, in general by w<sub>u</sub>(x) = u$``$x$``$u<sup>-1</sup> for u$``$0. In order that these relabelings may be non-trivial it is only necessary that u not commute with all of the arguments of the function to be defined. This indicates the important role to be played by the availability of symbols which don’t commute with given symbols, for if no such symbols are available then even a simple product would be incompatible with relabeling, let alone a natural definition of a function motivated by it.
It will be noticed, in the case of the product A = $`_{j=1}^n`$a<sub>j</sub>, that any symbol which commutes with all of the a<sub>j</sub> also commutes with A. This will serve as a defining characteristic of a function.
The considerations to follow in the development of the theory of functions and differentiation, while proceeding from different motivation, nevertheless formally parallel those already given by Dirac <sup>1</sup><sup>1</sup>1 See reference . in a paper addressing algebra in quantum theory. In particular, Dirac formalized operations already found in Quantum Theory and suggested formal axioms in order to carry out these desired operations. He also worked within the framework of standard mathematics with finite quantities and accepted the notion of mathematical truth. His proposals for the definition of functions and differentiation were made for application to the Quantum Theory in particular, but were not proposed as replacements for the corresponding notions in mathematics itself. Here the developments proceed from a non-physical motivation and basis, and mathematical quantities are not necessarily finite. With this difference between Dirac’s approach and the approach of this thesis in mind, the results for these subjects to be given below will be merely cited, but not formally derived here, unless they represent developments not formally parallel to Dirac’s presentation.
DEFINITION: Given certain symbols q<sub>α</sub> ordered by their index $`\alpha `$, a rule, called $`\varphi `$, for generating values $`\varphi `$(q<sub>α</sub>) from the ”arguments” q<sub>α</sub>, will be called a FUNCTION iff
a) Any $`\beta `$ which commutes with all of the q<sub>α</sub> also commutes with $`\varphi `$(q<sub>α</sub>), and
b) Whenever q$`{}_{}{}^{}{}_{\alpha }{}^{}`$= u$``$q$`{}_{\alpha }{}^{}`$u<sup>-1</sup> for any u $``$0, it follows that $`\varphi `$(q$`{}_{}{}^{}{}_{\alpha }{}^{}`$) = u$`\varphi `$(q<sub>α</sub>)$``$u<sup>-1</sup>.
Some results for functions follow.
1. Functions are well-defined, which means that any u, such as that above, will give the same value $`\varphi `$(q<sub>α</sub>).
2. $`\varphi `$(q$`{}_{}{}^{}{}_{\alpha }{}^{}`$) commutes with whatever commutes with q$`{}_{}{}^{}{}_{\alpha }{}^{}`$.
3. If $`\varphi `$ and $`\psi `$ are functions, then so are $`\varphi \pm \psi `$ and $`\varphi \psi `$.
4. V(x) $``$ x<sup>-1</sup> is a function.
5. Composition of functions gives a function.
6. Functions of the same variables commute.
Thus all arguments of a function commute.
It is to be noted that polynomials are still, according to the new definition, functions.
It will be noticed that a function whose arguments commute with all symbols will have a function value that is invariant under relabelings. Such a function will be called a “constant”. It is also the case that if the symbol u which defines the relabeling commutes with an argument q<sub>α</sub> then it will not change and the function will be constant in q<sub>α</sub>. If it is arranged that only one argument varies and yet the function value does not change, then the function is constant with respect to that variable. In fact, the function will not depend upon that variable, and so u will commute with all of the arguments that the function does depend upon, and then the function value won’t change under this relabeling.
### 5.3 Differentiation
In this section a generalization of the usual notion of differentiation will be given. The notion of differentiation given here will not depend upon the notion of a limit, but will still give the usual result for the differentiation of sums and products of functions, and will therefor agree with the results of the usual differentiation on analytic functions. It will also turn out that, for the new notion of differentiation, every function is differentiable.
Consider the function $`\varphi `$(q<sub>α</sub>). Then \[q<sub>α</sub>, q<sub>β</sub>\]$``$ q$`{}_{\alpha }{}^{}`$q<sub>β</sub>-q$`{}_{\beta }{}^{}`$q<sub>α</sub> $``$ 0. If q<sub>β</sub> doesn’t commute with everything, then another symbol p<sub>β</sub> may be found such that
$$\text{[q}\text{α}\text{,p}\text{β}\text{] =}\{\begin{array}{cc}\hfill 0& \text{if }\alpha \beta \hfill \\ \hfill 1& \text{if }\alpha \text{=}\beta \hfill \end{array}$$
Suppose Q=$`\varphi `$(q<sub>σ</sub>). Define $`_\alpha \varphi `$(q<sub>σ</sub>)= Q$``$p<sub>α</sub>-p$`{}_{\alpha }{}^{}`$Q.
This is the “partial derivative of $`\varphi `$ with respect to q<sub>α</sub>”. Note that if $`\varphi `$ is constant with respect to q<sub>α</sub> then $`_\alpha \varphi `$=0. Other properties of differentiation will now be given:
1. All p$`{}_{}{}^{}{}_{\alpha }{}^{}`$ satisfying \[q<sub>α</sub>,p$`{}_{}{}^{}{}_{\alpha }{}^{}`$\]=1 give the same value for $`_\alpha \varphi `$.
2. $`_\alpha \varphi `$ is also a function.
3. $`_\alpha `$($`\varphi `$+$`\psi `$)= $`_\alpha \varphi `$+$`_\alpha \psi `$ and $`_\alpha `$($`\varphi \psi `$)=($`_\alpha \varphi `$)$`\psi `$+$`\varphi `$($`_\alpha \psi `$).
4. The Chain Rule holds: $`_s\varphi `$(Q<sub>γ</sub>((s))=$`_{k=1}^n_k\varphi `$(Q<sub>γ</sub>(s))$`_s`$Q<sub>k</sub>(s).
5. $`_\alpha `$\[$`_\beta \varphi `$(Q<sub>γ</sub>)\]$`_{\alpha \beta }\varphi `$(Q<sub>γ</sub>) is such that $`_{\alpha \beta }`$=$`_{\beta \alpha }`$. In other words, the order of partial differentiation doesn’t matter.
Item number 3 indicates that $``$ is a ”derivation”, and so differentiation, in the generalized sense, treats sums and products of functions in the same way as ordinary differentiation does.
### 5.4 Integration
Differentiation begins with one function and returns another. It may be asked whether or not this process is invertible, so that a converse process of determining a function whose derivative is a given function can be performed. It will be shown that this problem can be solved in general and, with the condition that only constants have derivatives equal to zero, for a function which is unique up to a constant. Then some results about integration will be shown and some particular integrations, and functions arrived at through integration, will be defined. These results will be useful in the further development of the theory.
Suppose $`\mathrm{\Psi }`$ is a function of the variable s and other variables q<sub>α</sub>. Suppose also that a function $`\mathrm{\Phi }`$ is such that $`_s\mathrm{\Phi }`$(s,q<sub>α</sub>) = $`\mathrm{\Psi }`$(s,q<sub>α</sub>).
Then, for p defined as above,
$`_s\mathrm{\Phi }`$(s,q<sub>α</sub>)=$`\mathrm{\Phi }`$(s,q<sub>α</sub>)$``$p-p$`\mathrm{\Phi }`$(s,q<sub>α</sub>)=$`\mathrm{\Psi }`$(s,q<sub>α</sub>) (5.1)
where s$``$p-p$``$s=1 and q$`{}_{\alpha }{}^{}`$p-p$``$q<sub>α</sub>=0 for each q<sub>α</sub>, (5.2)
while s$``$q<sub>α</sub>-q$`{}_{\alpha }{}^{}`$s = q$`{}_{\alpha }{}^{}`$q<sub>β</sub>-q$`{}_{\beta }{}^{}`$q<sub>α</sub>=0 also holds. (5.3)
It can be seen that, for a symbol p defined to satisfy the conditions (5.2), and given $`\mathrm{\Psi }`$(s,q<sub>α</sub>), the symbol $`\mathrm{\Phi }`$(s,q<sub>α</sub>) may be defined to satisfy (5.1). This may be done because all that is under consideration here is a symbolic method that needs to be consistent with algebra, but is not such that each symbol needs to have an algebraic derivation. There is thus always at least one such symbol $`\mathrm{\Phi }`$, so all functions $`\mathrm{\Psi }`$ have “antiderivatives”.
Given such a $`\mathrm{\Phi }`$, consider now the function $`\mathrm{\Phi }`$+c, where c is a constant with respect to s. Then $`_s`$($`\mathrm{\Phi }`$+c)=$`_s\mathrm{\Phi }`$=$`\mathrm{\Psi }`$. Thus antiderivatives can, at best, be defined up to a constant in this way. This is as close as antiderivatives can be to being unique. It will thus now be assumed that only constants have zero derivatives.
Under these conditions, if both $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$ are such that $`_s\mathrm{\Phi }_1`$=$`_s\mathrm{\Phi }_2`$=$`\mathrm{\Psi }`$ then $`_s`$($`\mathrm{\Phi }_1`$-$`\mathrm{\Phi }_2`$)=0, so that $`\mathrm{\Phi }_1`$=$`\mathrm{\Phi }_2`$+c, for some constant c. Then $`\mathrm{\Phi }_2`$(b,q<sub>α</sub>)-$`\mathrm{\Phi }_2`$(a,q<sub>α</sub>) = $`\mathrm{\Phi }_1`$(b,q<sub>α</sub>)-$`\mathrm{\Phi }_1`$(a,q<sub>α</sub>) for such $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$ and particular a and b. It is then possible to uniquely define
$`_a^b\mathrm{\Psi }`$($`\tau `$,q<sub>α</sub>)d$`\tau `$ $``$ $`\mathrm{\Phi }`$(b,q<sub>α</sub>)-$`\mathrm{\Phi }`$(a,q<sub>α</sub>) for any $`\mathrm{\Phi }`$ such that $`_s\mathrm{\Phi }`$(s,q<sub>α</sub>)=$`\mathrm{\Psi }`$(s,q<sub>α</sub>). (5.4)
This condition gives $`_a^b\mathrm{\Psi }`$($`\tau `$,q<sub>α</sub>)d$`\tau `$ as a well-defined function of a, b, and q<sub>α</sub>, and also gives a relationship formally identical to the usual Fundamental Theorem of Calculus. For this reason it is reasonable to call $`_a^b\mathrm{\Psi }`$($`\tau `$,q<sub>α</sub>)d$`\tau `$ the “definite integral” of $`\mathrm{\Psi }`$ from a to b. With this general definition it is easy to deduce the usual formal properties of integrals:
1. $`_x_a^x\mathrm{\Psi }`$(s,q<sub>α</sub>)ds = $`_x`$\[$`\mathrm{\Phi }`$(x,q<sub>α</sub>)-$`\mathrm{\Phi }`$(a,q<sub>α</sub>)\] = $`_x\mathrm{\Phi }`$(x,q<sub>α</sub>) = $`\mathrm{\Psi }`$(x,q<sub>α</sub>).
2. Suppose $`\varphi `$(s)=$`\stackrel{~}{\varphi }`$(u(s)). Then $`_s\varphi `$ = $`_u\stackrel{~}{\varphi }_s`$u, so $`_\alpha ^\beta _u\stackrel{~}{\varphi }_s`$u ds = $`_a^b_s\varphi `$ds when u(a)=$`\alpha `$ and u(b)=$`\beta `$. Thus the Chain Rule leads, as usual, to the rule for changing the variables of integration.
3. $`_a^b\varphi `$d$`\tau `$+$`_b^c\varphi `$d$`\tau `$=$`_a^c\varphi `$d$`\tau `$. Thus integrals may be considered, as usual, to be given by sums. It also follows that $`_a^b\varphi `$d$`\tau `$ = - $`_b^a\varphi `$d$`\tau `$.
4. $`_a^b`$($`\alpha \varphi `$+$`\beta \psi `$)d$`\tau `$ = $`\alpha _a^b\varphi `$d$`\tau `$ \+ $`\beta _a^b\psi `$d$`\tau `$ for constant $`\alpha `$ and $`\beta `$.
5. $`_x`$($`\varphi \psi `$) = ($`_x\varphi `$)$`\psi `$ \+ $`\varphi `$($`_x\psi `$) so $`_a^b\varphi `$($`_x\psi `$)dx = $`_a^b_x`$($`\varphi \psi `$)dx - $`_a^b`$($`_x\varphi `$)$`\psi `$dx = $`\varphi \psi |_a^b`$ \- $`_a^b`$($`_x\varphi `$)$`\psi `$dx. This is the usual rule for integration by parts.
The differentiation of polynomials follows the usual rules, so it follows that integration of polynomials will also follow the usual rules. It can thus be seen that the theory of integration given here gives a non-trivial extension of the usual theory.
The tools just developed will now be applied to yield generalized versions of the logarithm and exponential functions.
As every function is integrable, the integral $`_a^b\tau ^1`$d$`\tau `$ may be considered without regard to the behavior, finite or otherwise, of the integrand.
Define now $`\mathrm{ln}`$(x)$`_a^x\tau ^1`$d$`\tau `$, with a as yet unspecified. Then $`\mathrm{ln}`$(x$``$y) = $`_a^{xy}\tau ^1`$d$`\tau `$ = $`_a^x\tau ^1`$d$`\tau `$ \+ $`_x^{xy}\tau ^1`$d$`\tau `$ = $`\mathrm{ln}`$(x) + $`_x^{xy}\tau ^1`$d$`\tau `$.
Let u($`\tau `$) = x$`{}_{}{}^{1}\tau `$. Then $`\mathrm{ln}`$(x$``$y) = $`\mathrm{ln}`$(x) + $`_1^y`$u<sup>-1</sup>du = $`\mathrm{ln}`$(x) + $`_a^y\tau ^1`$d$`\tau `$ \- $`_1^a\tau ^1`$d$`\tau `$ = $`\mathrm{ln}`$(x) + $`\mathrm{ln}`$(y) - $`_1^a\tau ^1`$d$`\tau `$. Now if a = 1 is chosen, then
$`\mathrm{ln}`$(x$``$y) = $`\mathrm{ln}`$(x) + $`\mathrm{ln}`$(y). (5.5)
This choice is made, so that, in general:
$`\mathrm{ln}`$(x) $``$ $`_1^x\tau ^1`$d$`\tau `$. (5.6)
Now $`_x\mathrm{ln}`$(x)$`|_{x=y}`$ = $`\frac{1}{y}`$, so that $`\mathrm{ln}`$(w) = $`\mathrm{ln}`$(v) implies that w = v, and so $`\mathrm{ln}`$ is a one-to-one function and may be inverted. Let g be the inverse of $`\mathrm{ln}`$. Then, according to (5.5), g must satisfy
g(x+y) = g(x)$``$g(y) with g(0)=1. (5.7)
Now, for rational p, it follows that $`\mathrm{ln}`$(x<sup>p</sup>) = p$`\mathrm{ln}`$(x), so then
x<sup>p</sup> = g($`\mathrm{ln}`$(x<sup>p</sup>)) = g(p$`\mathrm{ln}`$(x)). (5.8)
This provides a general definition of x<sup>p</sup> regardless of whether or not p is rational. In fact, g defines irrational symbols, so that g definitely extends the known range of the symbolic system. Such irrationals correspond to infinite sequences of integers in their decimal expansions, so that the inclusion of irrationals in the symbolic system requires the admittance of the results of limits and, in particular, requires the inclusion of infinite symbols and of reals. If the symbol “e” is defined by the requirement that it be such that $`\mathrm{ln}`$(e) = 1, then e<sup>x</sup>=g(x$`\mathrm{ln}`$(e)) = g(x). Thus g clearly generalizes the exponential function.
### 5.5 The Relationship to the Usual Calculus
An abstract theory of functions, differentiation, and integration has been given which relies on the constructibility of symbols with tailor-made commutation properties but does not rely on the use of limits nor assume that the quantities involved are finite. Moreover, differentiation obeys the usual rules when operating upon sums and products of functions and so the generalized theory is an extension of the usual Calculus. In fact, as a basis for comparison between the usual and the generalized Calculus, it may be helpful to think of implementing the given rules for the generalized theory in a symbolic math program. The computations performed would then be identical to those performed in a version of the program based upon the usual Calculus whenever exact results are required and recourse to numerical methods for approximating limits is not allowed.
Conversely, if it is assumed that all quantities involved in computation are finite, then the generalized functions, differentiation, and integration given here reduce to those of the usual Calculus. It is thus possible to apply the formal methods of the Calculus in a domain where the finiteness of quantities is not presumed and to also, then, investigate explicitly where and how the finiteness of quantities affects the results.
## Chapter 6 COMPLETION OF THE SYMBOLIC SYSTEM
### 6.1 Functional Equations and Their Consequences
With the development of the Calculus it is now possible to derive some further results on the algebraic nature of the general symbolic system. This will be done by solving, by the methods of the calculus just developed, certain functional equations which are constructed based on the relationship between the formal system and an embedded lattice of meaningful symbols that results from it upon interpretation. These considerations will indicate that $``$ may, in fact, be taken as the extension of $``$ to the general symbolic system, and that complex numbers must necessarily be included in the general symbolic system.
Consider the meaningful symbols a, b, and c. In considering meaningful symbols, it follows that there is a function $`\varphi `$ which maps these symbols into corresponding symbols in the symbolic system. Thus let $`\varphi `$(a) = x, $`\varphi `$(b) = y, and $`\varphi `$(c) = z. Now $`\varphi `$ must be such that, for some function F,
$`\varphi `$(a$``$b) = F\[$`\varphi `$(a), $`\varphi `$(b)\] = F\[x, y\] (6.1)
The associative property of $``$ then immediately results in
F\[x, F\[y, z\]\] = F\[F\[x, y\], z\] (6.2)
This will be referred to as the ASSOCIATIVITY EQUATION.
The Associativity equation is a functional equation for F, and the solution of functional equations is, in general, a difficult problem<sup>1</sup><sup>1</sup>1 See and . the solution of which generally requires special additional assumptions about the function to be determined. In the general symbolic system, however, this equation can be solved in all generality without restrictive assumptions about F<sup>2</sup><sup>2</sup>2In this connection also consult .. The solution process is intricate, however, and the details are left to the appendix. The solution is here simply given as
F\[x, y\] = $`\mathrm{\Phi }^1`$\[$`\alpha \mathrm{\Phi }`$(x)$`\mathrm{\Phi }`$(y)\] (6.3)
where $`\alpha `$ is a fixed constant and the form of $`\mathrm{\Phi }`$ is specified in the the appendix. The form of the equation (6.3) indicates that the operation $``$ may definitely be taken to correspond to the operation $``$. This is significant because the operation $``$ was merely motivated by the operation $``$ without there being a definite connection between them assumed.
Further results about the symbolic system may be derived by considering the relationship between interpretation and operations between meaningful symbols. Consider, in particular, a Boolean Lattice of meaningful symbols, where a complement a is given for every meaningful symbol a. Recall that (a) = a. Now, for some function f, it must be that
$`\varphi `$(a) = f($`\varphi `$(a)). (6.4)
From this it follows that f satisfies
f(f(x)) = x. (6.5)
This is another functional equation; it shall be referred to as the CATEGORICITY EQUATION.
As was the case with the Associativity equation, the solution of the Categoricity equation may be found in the appendix. The solution is given, in general, by
f(x) = \[ c<sup>r</sup> \- x<sup>r</sup> \]$`^{\frac{1}{r}}`$ (6.6)
where c and r are arbitrary fixed constants.
So far, the consideration of finite constructions within the symbolic system has resulted in rationals and polynomials. The discussion of the Calculus has even extended the symbolic system to irrational numbers and reals. The equation (6.6) must, however, be solvable for any and all constants c, r. This necessarily results in the admission of further kinds of symbols. Consider, in particular, r = 2, with c and x rationals such that x<sup>2</sup> $`>`$ c<sup>2</sup> \+ 1. Then, in order that f(x) may be solved for, it is necessary that (-1)$`^{\frac{1}{2}}`$ be a symbol in the general symbolic system. Thus the general symbolic system must include complex numbers. This is quite an ironic result in the sense that the interpretability of a Boolean Lattice has resulted in the necessity of the symbolic system being complex, at least, whereas such complex quantities only entered physical theory with the advent of Quantum theory and are taken to reflect a non-Boolean logic.
The considerations of this section have shown a definite connection between the algebraic operations and algebraic types of the symbols of the general symbolic system on the one hand, and the corresponding operations in an embedded lattice of meaningful symbols on the other. With these results it is possible to proceed, in the next section, to a complete algebraic specification of the general symbolic system.
### 6.2 Frobenius’s Theorem and the Completion of the System
The specification of the general symbolic system is almost complete. It only remains to consider some consequences of the way in which finitely many symbols may be combined and encoded. This will have consequences for the algebraic form of the system.
Consider the symbols that may be finitely generated by the algebraic operations and finitely many variables and complex numbers. The result will be a polynomial p of finite order, over the complex numbers, and in the given variables. Suppose attention is restricted to only the consideration of finite symbols. In this context it is a standard result that any such given finite order polynomial corresponds to a definite list of its distinct zeros, and conversely that a polynomial may be constructed which has precisely a given list of zeroes. Thus a correspondence may be given between finite selections of meaningful symbols and polynomials which represent them. It must certainly be possible to construct a restricted form of the formal system which is comprised only of symbols which are generated in this way. The general symbolic system must then be an extension of this system.
The restricted finite form of the symbolic system would therefor have to be algebraic over the reals at least, and the general system must be an extension of this restricted system. In order that uniform rules as to the algebraic form of the symbols in the general symbolic system may be adopted, it is necessary that the general system have the same algebraic form as the finite restricted system. The general form of division rings of finite symbols algebraic over the reals has already been established in a theorem of Frobenius <sup>3</sup><sup>3</sup>3 See pp. 327-329 of reference .. This theorem indicates that such a division ring must be isomorphic to one of the real field, the complex numbers, or the real quaternions.
In addition to the information given by Frobenius’s Theorem, it may also be concluded that neither the form of the reals nor complex numbers may be adopted for the general symbolic system because these are commutative fields, and the formal system must include non-commuting variables. The unequivocal conclusion is therefor that in the general symbolic system the symbols must have the algebraic form of the real quaternions, and in fact be real quaternions in the case that they are to be finitely generated. This simple algebraic result will later be shown to be of importance in accounting for the four-dimensionality of space-time in the derivation of the general relativistic formalism. This requirement will also distinguish the quantum theoretical derivation to follow in that while the usual theory operates over the complex field, there is no clear justification for or understanding of the significance of this particular choice, whereas such understanding is supplied here.
To sum up, it has been found that every symbol $`\omega `$ in the general symbolic system will have the form
$`\omega `$ = $`\omega _0`$ \+ $`\omega _1`$i +$`\omega _2`$j +$`\omega _3`$k, where (6.7)
i<sup>2</sup>=j<sup>2</sup>=k<sup>2</sup>= -1, i$``$j=-j$``$i=k, j$``$k=-k$``$j=i, and k$``$i=-i$``$k=j. (6.8)
In general the $`\omega _k`$ will commute with one another, will satisfy the rules of a division ring, and will comprise the components of such four-tuples, but they will not generally be real numbers, for the finiteness of the $`\omega _k`$ is not presumed for the general symbolic system. It may also be observed that such $`\omega `$ will not be generally suitable as arguments for functions, as defined here, for only restricted classes of such symbols will commute with one another.
One final note. Given the expression (6.7) for $`\omega `$ it is possible to define the CONJUGATE of $`\omega `$, denoted by $`\overline{\omega }`$, as
$`\overline{\omega }`$ = $`\omega _0`$ \- $`\omega _1`$i -$`\omega _2`$j -$`\omega _3`$k. (6.9)
The introduction of the conjugate also allows the INNER PRODUCT, between two symbols a and b, to be defined by
(a, b) = a$`\overline{\text{b}}`$. (6.10)
Then, in the usual way, the NORM may be defined by
$`||`$a$`||`$ = $`\sqrt{\text{ (a, a)}}`$. (6.11)
### 6.3 Conclusions on the Algebraic System
The time has finally come to take stock of the net effect of all these abstract considerations about the general symbolic system. The following conclusions have been reached:
1. The general symbolic system is a division ring with the algebraic form of the real quaternions.
2. The formal system has a partial order defined on part of it.
3. Symbols have conjugates and norms, and an inner product is defined.
4. Symbols are not necessarily finite, and infinite symbols must be included in the general symbolic system.
5. A generalized version of the Calculus has been defined on the symbolic system without reference to finite quantities.
6. The generalized Calculus gives an extension of the usual limit-defined notions, so that all functions are both differentiable and integrable.
With these observations it will now be possible to proceed to the derivation of physical theories without the utilization of any experimental results whatsoever!
## Chapter 7 THE DERIVATION OF PHYSICAL THEORY
### 7.1 Introduction
The considerations of the previous chapters have been preliminary to this chapter, for to substantiate a claim that physical theory can be derived in a purely formal manner, without recourse to experiment, it is necessary to first develop a formal system on its own terms. Having done so, the task of constructing physical theory will be taken up in the following manner.
It has already been argued that the very nature of experimentation is such that it must be possible to state its results and carry out its predictions in terms of finitely many finite symbols. It is therefor of particular interest to investigate, as a preliminary special case, the restriction of the general symbolic system that results if it is first assumed that the variables in the symbolic system are all finite. It will then be argued that the associated finite relations comprising physical laws will be expressible in a particular canonical form. The analysis of this form may be broken into two cases.
First it will be shown what form of predictive theory results if it is assumed that the finite relations are to be determinate. This will yield a novel derivation and justification for General Relativity Theory<sup>1</sup><sup>1</sup>1 There are many presentations of Relativity. For an intuitive motivation and overview, see . For a more rigorous introduction in the geometric style, see . For an exposition given in the language of forms, see . A non-geometrical route is taken in . Finally, a concise mathematical presentation may be found in . which is inevitable, on these assumptions, and which is free from experimental or merely plausible theoretical justification. It will also follow that Relativity is the unique form of finite deterministic theory. From Einstein’s Field Equation, serving as the basis of General Relativity Theory, the dependent notion of Turing Computability will also arise.
Finite laws that are statistical are discussed next. In order that these laws give information that is independent of Einstein’s theory, it is necessary that the quantities they provide for prediction must not, at the very least, be determinable by Turing Machines. This will lead to a unique set of requirements for statistical predictions, and consequently to a propensity theory of probability which is similar, in its foundations, to Von Mises’s Frequency Theory of Probability <sup>2</sup><sup>2</sup>2 See references and .. This theory will apply to the cases of both continuous and discrete representation of data, and so is perfectly general.
Combining these interpretive aspects of finite statistical theory with the previously determined form of finite relations will yield Quantum Theory as the unique embodiment of finite statistical theory<sup>3</sup><sup>3</sup>3 For an axiomatic presentation of the quantum formalism see . Note especially the comment on pg. 75 in regard to the finiteness of experimental data.. The strength of this approach is shown in that it dismisses the controversies that have arisen in deciding the proper way to interpret the quantum formalism.
Such discussion completely analyzes finitary theory but it is found, in fact, that these two canonical finite physical theories apparently result in a mathematical dichotomy. This difficulty may be understood to stem from one of two sources and lead to correspondingly definite means of resolution. One solution is to maintain the restriction to finite mathematical theories and resign one’s self to forever deferring to experiment. Alternately, the assumption of finiteness could be dropped, and this, it will be argued, results in a unique non-finitary formalism. Interestingly enough, these two solutions, while embodying the contending philosophical approaches of pragmatic empiricism and of idealistic metaphysics, are, in a definite sense, not at all distinct.
### 7.2 The Form of Finite Relations
Now that the abstract algebraic properties of the general symbolic system have been specified it is possible to arrive at some concrete conclusions about the kinds of descriptions this system can provide. Starting from the formal assumption that all descriptive symbols are finite it will be found, to begin with, that such finite relations must be expressible in a particular canonical form. The necessity of this form will have implications for the construction of physical theory.
Consider any function in the general symbolic system. Such a function, and indeed any term in any equation, will, in general, have a value and arguments of the form of quaternions. Thus any equations can be expressed as a system of four functions set equal to zero, each function having values of the form of either real or complex numbers. For each such function, the corresponding arguments must be commuting symbols.
In the general case the constructibility of symbols which didn’t commute with given arguments was of central importance in the derivation of the generalized Calculus. In the finite case, however, it has been observed that the operations of the usual Calculus, with their associated limit dependence, may be substituted and do not require that the arguments not commute with certain symbols. For this reason it is possible to take the arguments of all functions to be complex when it is assumed that all symbols are finite.
It should be pointed out that expressions in the finite case therefor, while yet in correspondence with expressions in terms of real quaternions, do not directly express the non-commutative nature of the division ring of the general symbolic system. With this understanding any finite relations will, in fact, be expressed in terms of complex-valued functions of complex arguments.
Consider any finite collection of distinct symbols each of which has been put into correspondence with, and therefor represented by, a distinct complex number q<sub>α</sub>. It is a standard result of complex variable theory that it is then possible to construct an analytic function $`\varphi `$(z) which has zeroes at precisely the points z=q<sub>α</sub>. Now the condition $`\varphi `$(z)=0 is equivalent to the condition $`|\varphi `$(z)$`|`$=0, so that given a known collection of data represented by the complex numbers q<sub>α</sub>, the same collection is represented by a certain analytic function $`\varphi `$(z) and is recovered by solving the equation $`|\varphi `$(z)$`|`$=0. In this way it is seen that the analytic function $`\varphi `$(z) associates the zeroes q<sub>α</sub> and thus expresses a relation between them. It is thus the case that any quantity depending on the collection of data q<sub>α</sub> will, in fact, be a function of $`|\varphi `$(z)$`|`$. The general conclusion is reached that all statements and quantities upon which finite physical theory may be based are to be found in this form.
The considerations above were made under the assumption that the data q<sub>α</sub> were already given and then $`\varphi `$(z) was to be constructed to represent them. However, the converse procedure is completely justified, so that, starting with an analytic function selected according to some criterion, $`|\varphi `$(z)$`|`$ will have to be the quantity upon which prediction depends and it will, in general, automatically engender an associated collection of data. It may be noted in passing that the consideration of analytic functions generally might also be used to extend the theory to the case of countably many data, as analytic functions with that many zeroes may be constructed. In the same way as above, the formal criteria which select $`\varphi `$(z) also correspond to certain “phenomena” so that the usual intuitive process of looking for quantities to represent what is observed may be reversed. This observation helps in understanding how a completely formal theory might reproduce the physical theories which have been motivated by experiments and observation in the past, although the difficulty remains of identifying such constructive and yet abstract criteria which $`\varphi `$(z) is to satisfy.
It has been found that all finite predictive quantities must be expressible in a particular canonical form, and that this form also indicates a way to formally bypass experiment. The next section will move beyond these abstractions to explore an application of this idea and to indicate a particular success of this approach.
### 7.3 Determinism and Relativity
#### 7.3.1 Introduction
In the sections to follow it will first be demonstrated that a particular form of $`|\varphi `$(z)$`|`$, called the interval, must be derivable which gives invertible relations between pairs of symbols. Next it will be shown how the interval may be put into a canonical form at a point by a change of coordinates, and that this form has implications for the corresponding physical theory. In particular, relations between events may be expressed in terms of the relationship between locally preferred coordinate systems, each of which give the interval the canonical form at a particular point, associated with each event. The sequence of events determined by the theory will then be defined by the allowable evolution of the preferred coordinate systems associated with these events. In order to determine this evolution it is next shown how the evolution of the preferred coordinate system may be expressed, via the geodesic equation, in any system of coordinates. Examination of these equations will indicate the quantities in the geodesic equation which will govern the evolution of the preferred coordinate system, and then it will be shown that there is a unique way that the behavior of the preferred coordinate systems may be determined in a manner independent of the coordinate system in which the definition is expressed. This will yield Einstein’s Field Equations. Finally, it will be shown that Einstein’s equations are, as they must be, equivalent to the formal requirement of the general symbolic system which started these considerations.
#### 7.3.2 Derivation of the Interval
In this section it will be demonstrated that a particular form of $`|\varphi `$(z)$`|`$ must be derivable which gives invertible relations between pairs of symbols. This will give a novel derivation and interpretation of General Relativity Theory without any further assumptions or recourse to experiment. This derivation will thus explain all relativistic phenomena and concepts, such as gravitation, the behavior of light signals and space-time, and the unification of the concepts of mass, momentum, and energy, as the natural implications of a formal necessity found in the general symbolic system.
In this section it will, in fact, be demonstrated that the necessity of allowing one-to-one relabeling in the general symbolic system requires that a particular form of finite relation, determined by $`|\varphi `$(z)$`|`$, must be derivable which gives invertible relations between pairs of symbols.
Recall the formal process of relabeling in the general symbolic system. In order to investigate a special restricted case ,require, in addition, that this procedure be one-to-one: Then x$``$f(x) is such that a = b exactly when f(a) = f(b). Thus such one-to-one relabeling maintains the distinctness of symbols: They may be said to be “separated”, so that if a$``$b then f(a)$``$f(b). This is a definite symmetric relationship between symbols which is maintained under invertible relabeling. As such manipulations are an inherent part of the general symbolic system, this precise relationship must find its expression for finite relations as well.
Suppose $`\varphi `$(z)=m$``$e<sup>ıθ</sup>, where m$``$0, is to express this relationship for pairs of symbols a and b. Definite values for a and b must then give definite values for $`\varphi `$(z). However, m=$`|\varphi `$(z)$`|`$, upon which expression of this relationship depends, does not generally identify a unique function $`\varphi `$(z): $`\theta `$ is as yet undetermined. Such a specification requires a conventional value of $`\theta `$ be chosen, and the most convenient is $`\theta `$=0. Then $`|\varphi `$(z)$`|`$=$`\varphi `$(z)= m $``$0. In other words, the realization of the relationship maintained by invertible relabeling requires, for its finite expression, utilizing an analytic $`\varphi `$-function which is real-valued and non-negative.
The labels a and b referred to above are representatives of the general symbolic system and, as such, are given by real quaternions in the finite case. Each label can, therefor, be represented by 4-tuples of real numbers: x<sup>μ</sup>, where $`\mu `$=0,$`\mathrm{}`$,3. Thus $`\varphi `$(z) may be rewritten as $`\delta `$(x$`{}_{1}{}^{}{}_{}{}^{\mu }`$,x$`{}_{2}{}^{}{}_{}{}^{\nu }`$). Let dx<sup>μ</sup>=x$`{}_{2}{}^{}{}_{}{}^{\mu }`$-x$`{}_{1}{}^{}{}_{}{}^{\mu }`$. Then, with some abuse of notation, $`\delta `$(x$`{}_{1}{}^{}{}_{}{}^{\mu }`$,dx<sup>μ</sup>)$`\delta `$(x$`{}_{1}{}^{}{}_{}{}^{\mu }`$,x$`{}_{2}{}^{}{}_{}{}^{\mu }`$) may be defined. $`\delta `$, being analytic, may be expanded in a Taylor series for x$`{}_{1}{}^{}{}_{}{}^{\mu }`$ fixed and dx<sup>μ</sup>=O($`ϵ`$):
$`\delta `$(x$`{}_{1}{}^{}{}_{}{}^{\mu }`$, dx<sup>μ</sup>)=g(x$`{}_{1}{}^{}{}_{}{}^{\mu }`$)+g<sub>α</sub>(x$`{}_{1}{}^{}{}_{}{}^{\mu }`$)dx<sup>α</sup>+g<sub>αβ</sub>(x$`{}_{1}{}^{}{}_{}{}^{\mu }`$)dx<sup>α</sup>dx<sup>β</sup>+g<sub>αβγ</sub>(x$`{}_{1}{}^{}{}_{}{}^{\mu }`$)dx<sup>α</sup>dx<sup>β</sup>dx<sup>γ</sup>+O($`ϵ^4`$). (7.1)
Note here that Einstein’s convention, where repeated indices are summed over, is utilized in (7.1). Also note that all g-functions are symmetric in their indices as $`\delta `$ is analytic. Now it may be seen, at O(1), that:
$`\delta `$(x$`{}_{1}{}^{}{}_{}{}^{\mu }`$, x$`{}_{2}{}^{}{}_{}{}^{\mu }`$)-$`\delta `$(x$`{}_{2}{}^{}{}_{}{}^{\mu }`$, x$`{}_{1}{}^{}{}_{}{}^{\mu }`$)= g(x$`{}_{1}{}^{}{}_{}{}^{\mu }`$)-g(x$`{}_{2}{}^{}{}_{}{}^{\mu }`$). (7.2)
But $`\delta `$ is symmetric in its arguments, so that:
g(x$`{}_{1}{}^{}{}_{}{}^{\mu }`$)=g(x$`{}_{2}{}^{}{}_{}{}^{\mu }`$). (7.3)
Thus g is locally constant, and this result may be continued so that g=0 may be stipulated globally. This amounts to choosing $`\delta `$(x<sup>μ</sup>, x<sup>μ</sup>)=0 rather than some other fixed number. Consider now $`\delta `$(x$`{}_{1}{}^{}{}_{}{}^{\mu }`$, -dx<sup>μ</sup>). As $`\delta `$0, it follows that all odd-ordered terms must be zero, so that:
$`\delta `$(x$`{}_{1}{}^{}{}_{}{}^{\mu }`$, dx<sup>μ</sup>) = g<sub>αβ</sub>(x$`{}_{1}{}^{}{}_{}{}^{\mu }`$)dx<sup>α</sup>dx<sup>β</sup>+O($`ϵ^4`$). (7.4)
It may be recalled that $`\delta `$ was to fulfill the exact requirement of expressing the separateness of symbols and nothing more. Thus, because the second order term in (7.4) already suffices to make this distinction in general, the higher order terms must be dropped. The form of the function $`\delta `$ is thus uniquely given by
$`\delta `$(x$`{}_{1}{}^{}{}_{}{}^{\mu }`$, dx<sup>μ</sup>) = g<sub>αβ</sub>(x$`{}_{1}{}^{}{}_{}{}^{\mu }`$)dx<sup>α</sup>dx<sup>β</sup>. (7.5)
Because $`\delta `$0 and is real, it follows that $`\delta `$ may be rewritten as $`\delta `$=ds<sup>2</sup>, where ds$``$0 is real-valued. The value ds corresponding to the labels x$`{}_{1}{}^{}{}_{}{}^{\mu }`$ and x$`{}_{2}{}^{}{}_{}{}^{\mu }`$ is called the “interval” between them. This allows (7.5) to be put in the conventional form:
ds<sup>2</sup>= g<sub>αβ</sub>(x$`{}_{1}{}^{}{}_{}{}^{\mu }`$)dx<sup>α</sup>dx<sup>β</sup>. (7.6)
This equation now defines, for given points x$`{}_{1}{}^{}{}_{}{}^{\mu }`$ and x$`{}_{2}{}^{}{}_{}{}^{\mu }`$, the interval ds once the functions g<sub>αβ</sub> may be stipulated. It must be remembered that conditions specifying the g<sub>αβ</sub> are required before this definition is complete. It should also be noted that, if ds<sup>2</sup> is to express only the separateness of symbols, then the g<sub>αβ</sub> must be determined by a self-consistent equation with no other data required. Otherwise, the g<sub>αβ</sub> will reflect both the separateness of symbols and something else.
#### 7.3.3 Changes of Coordinates and the Canonical Form
Now, in this section, it will be shown how the interval just derived may be, and must be able to be, put into a canonical form at each point by a change of coordinates. The preferred coordinate system which is thus associated with each point then serves as a definite means of speaking of a relationship between points: The canonical form at a point will identify a locally preferred coordinate system, so that events may be characterized in a definite way and it thus becomes possible to speak of definite relations between them. This will lead, in short order, to a logical derivation of the usual relativistic space-time picture of physical reality.
The coordinates x<sup>μ</sup> of labels, being subject to relabeling, cannot have a direct significance. It is possible, however, to indicate at this point how relabeling is effected and how this affects the terms of (7.6).
Suppose x<sup>μ</sup>=x<sup>μ</sup>($`\overline{x}^\alpha `$) gives a relabeling, or change of coordinates. It is only required that this substitution leave ds<sup>2</sup> invariant in (7.6). Thus the change of coordinates need only be continuous, and substitution of
dx<sup>μ</sup>=$`\frac{\text{x}^\mu }{\overline{x}^\alpha }`$d$`\overline{x}^\alpha `$. (7.7)
into (7.6) leads to
$`\overline{\text{g}}_{\mu \nu }`$=g$`{}_{\alpha \beta }{}^{}\frac{\text{x}^\alpha }{\overline{x}^\mu }\frac{\text{x}^\beta }{\overline{x}^\nu }`$. (7.8)
Now because ds<sup>2</sup> is a real quadratic form it is, according to a standard theorem<sup>4</sup><sup>4</sup>4 See pg. 207 of ., always possible, at any point x<sup>μ</sup>, to change coordinates continuously so as to obtain the canonical form
$$\text{ds}^2\text{=}\eta _{\alpha \beta }\text{d}\xi ^\alpha \text{d}\xi ^\beta \text{=}\eta _{\alpha \alpha }\text{(d}\xi ^\alpha \text{)}\text{2}\text{ where }|\eta _{\alpha \beta }|\text{=}\delta _{\alpha \beta }\text{.}$$
(7.9)
It is clear that a uniform definition of ds cannot be arrived at unless the signs of $`\eta _{\alpha \alpha }`$ are determined in the same way for all points x<sup>μ</sup>. Now each x<sup>μ</sup> is to correspond to some event which an individual perceives at a particular instant, so that one of the variables x<sup>μ</sup> must be specially identified, by a sign for $`\eta _{\alpha \alpha }`$ distinct from that of the others, to serve as an indicator of the instant in question. The arbitrary choice that x<sup>0</sup> be identified by $`\eta _{00}`$= 1, and thus also that $`\eta _{mm}`$= -1 for m=1, 2, and 3, is made.
Once the canonical form of ds<sup>2</sup> has been decided upon, the changes of coordinates that may lead to it at a particular point take on a definite significance. It is then possible to speak of definite relationships between different points based on the relationship between the preferred coordinate systems associated with these points.
This has the consequence, first of all, that the labels x<sup>μ</sup> will now be formally identical to the usual space-time which comprises a pseudo-Riemannian manifold. In such a scheme there is an associated notion that the relation between events which ds specifies is given by a sort of signaling. In the standard empirical justification of Relativity, this signaling is taken to be a property of light. Here, no such recourse to experiment is implied, but it is instead asserted that the formal properties of the symbolic system itself require that such a relation between events obtain, and that, if it is desired that this relation be interpreted in in terms of signaling, there be “something” which travels between events and acts as such a signal. This something, being always associated with the events which it acts as a signal for, can be thought of as being made up of whatever causes the perception of these events. It follows that if the identification of events is ultimately made visually, for example, then the signal between events must be realized by light.
The necessity of a uniform definition of ds also has the consequence that instants fall, from a finitely justifiable point of view, into a linear order as does the real-valued variable x<sup>0</sup>, and thus introduces the qualitative temporal notions of “earlier” and “later” to the instants. It is to be noticed that in this construction of the notion of time, the linear and continuous aspects of time are not assumed, based on experience, but are instead derived from the formal nature of the symbolic system as a whole.
It must also be stressed that the three remaining non-temporal labels x<sup>m</sup> are not presumed in order that the usual notion of space may find a place within the formalism. Instead, this three-dimensional aspect of finite description rather justifies the notion of space, for the relations found in “spatial perceptions” are never argued for on the basis of infinitely many examples. Because the notion of space may be derived simply from the assumption of a finite check of its relations, and always includes this assumption, it would be superfluous to assert any further source of the formal nature of spatial perceptions. Both space and time may thus be seen to be moulds into which experiences are poured when construed finitely.
Once the functions g<sub>αβ</sub> are determined the differential relation (7.6) for ds<sup>2</sup> may, in most cases, be integrated along a path for which ds$``$0. Paths which give the least value may then uniquely define the interval between distant events. It will then be possible, by a chain of intermediate transformations, to relate such events by changes of variables.
This process will not, however, necessarily relate all pairs of points x<sup>μ</sup>. It may happen that the extension of a path along which ds$``$0 is blocked by the unsolvability of this condition. This will be the case where attempting to cross regions in which ds$``$0 cannot be satisfied or in which ds simply vanishes. Such pairs of events will then not be finitely related in this direct way, and Relativity is, in this sense, a “local” theory. Such events cannot “communicate” by (light) signaling, but the pairs of points x<sup>μ</sup> will still be formal parts of the descriptive apparatus of the theory and serve to mathematically complete the manifold. A new notion, that of relatively isolated parts of “reality”, even parts which are never experienced, is thus introduced.
Another interesting remark may be made at this point. The functions specifying changes of coordinates have thus far been interpreted “passively” as giving a mere renaming of the same events. It is unavoidable, however, that these functions will also relate pairs of events identified by the arguments and the values of these transformations. Because these events may both fall within a particular individual’s experiences, the change of coordinates may be interpreted “actively” as referring to transitions between that individual’s experiences. The change of coordinates may thus show, at a particular instant, how that individual’s experiences are related, or, if the transformation is between distinct instants for the individual, how “earlier” events are related to “later” ones.
If these events are considered to be experienced by distinct individuals, then the change of coordinates specifies a “change of reference frame.” Once the interval is determined for all events, it follows that it will inter-relate the experiences of all distinct observers. The equal admissibility of all continuous changes of variables then gives all observers, which may be related in this way by changes in reference frame, an equal observational status.
#### 7.3.4 The Geodesic Equation
The canonical form discussed in the last section identifies a locally preferred coordinate system that is associated with each point in space-time. With these coordinate systems characterizing each point it now has become possible, in principle at least, to speak in some definite way about the relationship between events. Progress on this count will be made in this section, where the form of the equations describing the evolution of the locally preferred coordinate systems will be determined. It will then be apparent which quantities in these equations must be determinative in any law for this evolution.
It has been found that the finite-relatedness, generated by invertible relabeling, of pairs of events may be expressed in terms of the real quadratic form ds<sup>2</sup>, and that ds<sup>2</sup> depends, in turn, on the functions g<sub>αβ</sub>. The first step in the determination of the g<sub>αβ</sub> will now be undertaken.
For ds$`>`$0 it follows from (7.9) that:
$$\text{ 1 = }\eta _{\alpha \alpha }\left(\frac{\text{d}\xi ^\alpha }{\text{ds}}\right)^2.$$
(7.10)
From this it follows that
$$0\text{ = }\eta _{\alpha \alpha }\left(\frac{\text{d}\xi ^\alpha }{\text{ds}}\right)\left(\frac{\text{d}\text{2}\xi ^\alpha }{\text{ds}\text{2}}\right).$$
(7.11)
In order that (7.11) hold generally it is necessary that
$$\frac{\text{d}\text{2}\xi ^\alpha }{\text{ds}\text{2}}\text{ = 0.}$$
(7.12)
For any parameter $`\omega `$ and new coordinate system $`\xi ^\alpha `$=$`\xi ^\alpha `$(x<sup>μ</sup>) the equation (7.12) takes the form
$$\frac{\text{d}\text{2}\text{x}^\lambda }{\text{d}\omega ^2}\text{ + }\mathrm{\Gamma }_{\mu \nu }^\lambda \frac{\text{dx}\text{μ}}{\text{d}\omega }\frac{\text{dx}\text{ν}}{\text{d}\omega }\text{ =0}$$
(7.13)
where $`\mathrm{\Gamma }_{\mu \nu }^\lambda \frac{\text{x}^\lambda }{\xi ^\alpha }\frac{^2\xi ^\alpha }{\text{x}^\mu \text{x}^\nu }`$. (7.14)
(7.13) are usually known as the “geodesic equations”, while the $`\mathrm{\Gamma }_{\mu \nu }^\lambda `$ are known as the “Christoffel” symbols of the first kind. These equations indicate the form taken by the conditions (7.12) in any coordinate system. Note that any differences between points x<sup>μ</sup> must be determined by $`\mathrm{\Gamma }_{\mu \nu }^\lambda `$ at the points, but that it is always possible to choose a coordinate system so that $`\mathrm{\Gamma }_{\mu \nu }^\lambda `$ vanishes at a given point. Thus the $`\mathrm{\Gamma }_{\mu \nu }^\lambda `$ determine the relationship between points under relabeling, but not in a manner independent of the associated change in coordinates.
#### 7.3.5 The Metric Field Equation
It will now be shown how the metric field g<sub>αβ</sub> must be determined in a manner independent of the coordinates. With the determination of g<sub>αβ</sub> the interval will also be defined.
Because both g<sub>αβ</sub> and $`\mathrm{\Gamma }_{\mu \nu }^\lambda `$ determine ds, it must be the case that they are related to each other. In fact, it may be shown that they satisfy the relations
$$\frac{\text{g}_{\mu \nu }}{\text{x}^\lambda }\text{ = }\mathrm{\Gamma }_{\lambda \mu }^\rho \text{g}_{\rho \nu }+\mathrm{\Gamma }_{\lambda \nu }^\rho \text{g}_{\rho \mu }\text{ , and}$$
(7.15)
$$\mathrm{\Gamma }_{\lambda \mu }^\sigma \text{ = }\frac{1}{2}\text{g}^{\nu \sigma }\left[\frac{\text{g}_{\mu \nu }}{\text{x}^\lambda }\text{ + }\frac{\text{g}_{\lambda \nu }}{\text{x}^\mu }\text{ - }\frac{\text{g}_{\mu \lambda }}{\text{x}^\nu }\right]$$
(7.16)
$$\text{where g}^{\nu \sigma }\text{g}_{\rho \nu }\text{ = }\delta _\rho ^\sigma \text{ is the delta function.}$$
(7.17)
Note that g<sup>νσ</sup> may be defined because the coordinate transformations considered are invertible. Because all $`\mathrm{\Gamma }_{\lambda \nu }^\rho `$ may be made to vanish at a point, it follows from (7.15) that the same may be done for all $`\frac{\text{g}_{\mu \nu }}{\text{x}^\lambda }`$. Then the equations (7.15) and (7.16) give no further information. A non-vanishing relationship may, however, be derived. This relationship will be expressed in a coordinate system in which $`\mathrm{\Gamma }_{\lambda \mu }^\sigma `$ vanishes, and may be derived by differentiating (7.16) with respect to x<sup>β</sup>. Then
$$\frac{\mathrm{\Gamma }_{\lambda \mu }^\sigma }{\text{x}^\beta }\text{ = }\frac{1}{2}\text{g}^{\nu \sigma }\left[\frac{^2\text{g}_{\mu \nu }}{\text{x}^\beta \text{x}^\lambda }\text{ + }\frac{^2\text{g}_{\lambda \nu }}{\text{x}^\beta \text{x}^\mu }\text{ - }\frac{^2\text{g}_{\mu \lambda }}{\text{x}^\beta \text{x}^\nu }\right].$$
(7.18)
It should be noted that this relationship is linear in the second derivatives of g<sub>μν</sub> with coefficients depending upon g<sub>μν</sub>. Such a differential equation is amenable to integration by parts, so that an equivalent condition may be derived from a variational principle. This will have the advantage of guaranteeing that the equation thus derived will determine g<sub>μν</sub> in a way that is independent of the coordinate system chosen. It is, in fact, necessary that such an equation may be found if invertible relabeling is to be realized in the finite system.
Consider the variational principle
$$\text{ I = }\text{R}\sqrt{\text{g}}\text{d}\text{4}\text{x}\text{, where }\delta \text{I = 0.}$$
(7.19)
Here g is the Jacobian of the transformation from the locally canonical coordinate system, which is given by the determinant of g<sub>αβ</sub>, and R is a scalar function of g<sub>αβ</sub> and its first and second derivatives.
Now equation (7.18) is a relation in symbols with four indices, so it will be advantageous to think of R as
$$\text{R = g}^{\mu \nu }\text{R}_{\mu \nu }\text{ = g}^{\mu \nu }\text{R}_{\text{ }\mu \lambda \nu }^\lambda .$$
(7.20)
If R is to be a scalar, then R$`{}_{\text{ }\mu \omega \nu }{}^{}{}_{}{}^{\lambda }`$ must be a tensor, and if $`\delta `$I = 0 is to yield a relation of the form of equation(7.18), then this tensor must be linear in the second derivatives of g<sub>αβ</sub>. Adventitiously, it’s the case that there is, up to constant multiples, exactly one tensor with these properties<sup>5</sup><sup>5</sup>5Proof of this assertion may be found on pp. 131-134 of .. It is the Riemann-Christoffel Curvature Tensor. Then R must be the Riemann Curvature Scalar.
Applying this information to (7.19) yields
$$\text{R}^{\mu \nu }\text{ - }\frac{1}{2}\text{g}^{\mu \nu }\text{R = 0.}$$
(7.21)
This is Einstein’s Field Equation (in free space). It is, according to the previous analysis, the unique equation defining a self-determined g<sub>αβ</sub> field which may merely express a finite local relation within the general symbolic system, this relation being determined solely by the requirement that relabeling be carried out in a one-to-one fashion.
#### 7.3.6 The Significance of the Field Equation
The equation (7.21) must be satisfied by the metric in order to result in a self-consistent finite theory which embodies simply the notion of relabeling. It should be possible, according to the reversible way in which the theory has been constructed, to make the equation (7.21) the sole basis of the entire theory. This is, as will be indicated, in fact possible. This approach may also be extended to provide the basis of all theories which are merely consistent with relabeling.
The theory presented so far relating to the interpretation of equation (7.21) is somewhat unsatisfactory in that a manifold on which (7.21) may be solved anywhere does not in any way indicate any finite number of points which may be distinguished from the others. There has to be some way in which a particular finite collection of points may be singled out so that equation (7.21) describes their behavior. The very construction of these equations indicates that such a distinction must be provided by a topological invariant. The curvature, in itself, is not eligible to make such a finitely defined distinction, so the only possibility is that points may be distinguished by there being singularities in the metric.
This representation of the interaction of a finite number of distinguished points was achieved by Einstein and Infeld <sup>6</sup><sup>6</sup>6 See reference .. Their conclusions were that the equations (7.21) were, in themselves, sufficient to completely determine the behavior of the singularities representing distinguished points. Moreover, it was found that the condition that the singularities of the metric field be restricted to the distinguished points resulted in these points moving according to the geodesic equation. Thus the Field Equation implies, in this sense, the geodesic equation and consequently that the intervals between points be invariant under one-to-one relabeling.
The preceding analysis also indicates how to determine all finite relations which are consistent with relabeling, for such relations must be determined by solutions of an extension of equation (7.21). Extensions of equation (7.21) will have to be of the form
$$\text{R}^{\mu \nu }\text{ - }\frac{1}{2}\text{g}^{\mu \nu }\text{R = kT}^{\mu \nu }.$$
(7.22)
An identity of Bianchi<sup>7</sup><sup>7</sup>7 For this and other mathematical details in this section, see . then indicates, as a consequence of (7.22), that
$$\text{T}_{\text{ :}\nu }^{\mu \nu }\text{ =0 , where“:” indicates covariant differentiation.}$$
(7.23)
This is a conservation equation, and T<sup>μν</sup> will give what will be interpreted as an Energy-Momentum Tensor. In fact this interpretation of T<sup>μν</sup> can be justified fully. The phenomena characterized by (7.22) will, because of the maintenance of separation which the construction of (7.22) uniquely reflects, be localizable, and localizable phenomena may be described by corresponding proper velocities
v<sup>μ</sup> = $`\frac{\text{dx}^\mu }{\text{ds}}`$. (7.24)
Now a general second order tensor may be expressed as a sum of terms which are the product of a scalar and two vectors. Thus, for a given scalar field $`\rho `$,
$$\text{T}^{\mu \nu }\text{ = }\rho \text{v}^\mu \text{v}^\nu $$
(7.25)
gives the most general expression for a source contribution for (7.22).
Applying (7.23) to this source results in
v<sup>μ</sup>($`\rho `$v<sup>ν</sup>)$`_{\text{:}\nu }`$ \+ $`\rho `$v<sup>ν</sup>(v<sup>μ</sup>)$`_{\text{:}\nu }`$ =0. (7.26)
Multiplying (7.26) by v<sub>μ</sub> yields
v<sub>μ</sub>v<sup>μ</sup>($`\rho `$v<sup>ν</sup>)$`_{\text{:}\nu }`$ \+ $`\rho `$v<sup>ν</sup>(v<sub>μ</sub>v$`{}_{\text{:}\nu }{}^{}{}_{}{}^{\mu }`$)=0. (7.27)
Now g<sub>μν</sub>v<sup>μ</sup>v<sup>ν</sup> = 1 implies that 0 = (g<sub>μν</sub>v<sup>μ</sup>v<sup>ν</sup>)$`_{\text{:}\sigma }`$=2g<sub>μν</sub>v<sup>μ</sup>v$`{}_{\text{:}\sigma }{}^{}{}_{}{}^{\nu }`$, and this yields
v<sub>ν</sub>v$`{}_{\text{:}\sigma }{}^{}{}_{}{}^{\nu }`$ = 0. (7.28)
With this the last term of (7.27) drops out, leaving
($`\rho `$v<sup>ν</sup>)$`_{\text{:}\nu }`$=0. (7.29)
This is the equation of conservation of Energy-Momentum, so that imposing the field equation has the automatic consequence of requiring the conservation of energy-momentum.
Returning attention to (7.26), it is seen that the first term vanishes, so that, for $`\rho `$0, this leaves
v<sup>ν</sup>v$`{}_{\text{:}\nu }{}^{}{}_{}{}^{\mu }`$=0. (7.30)
Consider the geodesic equation when v<sup>μ</sup>=v<sup>μ</sup>(x<sup>α</sup>(s)), so that multiple paths may be explored. It may be confirmed that in this case the geodesic equation is given by (7.30). Thus imposing the field equation also requires that phenomena follow the geodesic equation.
Thus, in both the free space case, given by (7.21), and in the presence of “matter”, as is the case for (7.22), the Field Equation is, in itself, a complete basis for the entire Relativity Theory. This is of particular importance in that it finally justifies, in a restricted case, the faith that a purely formal symbolic theory of experience would not necessarily degenerate to mere tautologies, but rather could lead to a constructive theory which accounts for some of the structure in the diversity of experience. Relativity does, after all, describe gravitation and therefor leads to an agglomeration and organization of the various labels which correspond to “matter”. This leads to the conclusion that, for a given finite level of description which is inter-related in a one-to-one manner, the general symbolic system has thus far succeeded in formally generating structure and the corresponding metalabeling.
### 7.4 On Turing Computability
Einstein’s Field Equation leads, as demonstrated in the last section, to the satisfaction of the geodesic equation for the description of traceable phenomena in space-time. An important limiting case for the solution of the geodesic equation is that of arbitrarily slow processes. In such cases the geodesic equation reduces<sup>8</sup><sup>8</sup>8 See, again, reference . to the Newtonian form
$$\frac{\text{dv}^m}{\text{dx}^0}\text{ = }(\text{g}_{00}^{}{}_{}{}^{\frac{1}{2}})_{\text{, m}}\text{ for m = 1,2,3.}$$
(7.31)
Here “,” indicates partial differentiation with respect to the variable with the indicated index. The form of these equations indicates that the particle motions described are those of motion according to Newton’s laws in a conservative potential. Such potentials may describe, among other things, elastic collisions. Thus the notion of a system of particles that interact by elastic collisions has a strictly theoretical derivation.
From this intuitively appealing special case of phenomena that may be generated from Einstein’s Field Equation, important conclusions follow. In fact, such elastic collisions define the ideal behavior of billiard balls, and, as a purely conceptual matter, it has been demonstrated <sup>9</sup><sup>9</sup>9 See reference . that such a “physical system” may perform all of the computational steps necessary in order to realize a Turing Machine. Thus Relativity, as a special case, contains all of Recursive Function Theory. If Church’s Thesis, that all computations may be carried out on a Turing Machine, is accepted then this demonstrates that that Relativity entails, in fact, all conceivable symbolic computations. Moreover, as all Turing Machines may be encoded as integers<sup>10</sup><sup>10</sup>10Kleene gave such a standard form. See ., it would follow that any symbolic computation would be specified by a particular integer, and the effects of Relativity, in restricted special cases, could be given by this same number.
It should also be observed that both Set and Class Theory are given in the context of a formal logic. Thus all derivations of statements in Set and Class Theory follow steps that may be directed by Turing Machines. It follows that the General Symbolic System cannot be any less powerful than either Set or Class Theory, as it, in fact, includes them.
### 7.5 Statistical Prediction and Probability
Relativity Theory, its derivation, and some of its consequences have now been discussed. It has been found that Relativity Theory may be derived on a strictly formal basis and includes, as a special case, all of Recursive Function theory and Set and Class Theory as well. Given the breadth of these results it may well be wondered whether or not there can be any other finite theory which hasn’t already been encompassed in Relativity Theory. The investigation of this possibility will be taken up next.
Relativity Theory has been shown to be the canonical finitary formalism in the case when invertible operations are performed on individual symbols. It will be recalled, however, that the general symbolic system, besides such individual symbols, also admits of metalabels or, in other terminology, “categories” as fundamental concepts. If the general symbolic system were exhausted by Relativity, then this would amount to saying that no formal laws may be stated in terms of categories. It is to be expected, and it will be shown in the following, therefor, that there are some finitary laws which take metalabels as such as their arguments. These relations will go beyond Relativity Theory.
Any finite relations which are to go beyond the predictions which Relativity already provides must, at least, not be computable by Turing Machines. Thus the considerations to follow will begin by considering what role Turing Machines may play in the theory of finite relations.
Consider a general function $`\rho `$ which is to represent a finite relation. As indicated in the last section, any Turing computable relations may be represented by an integer m which encodes an appropriate Turing Machine. Thus the function $`\rho `$ must have the dependence
$$\rho =\rho \text{(}\varphi \text{, m)}$$
(7.32)
where $`\varphi `$ is itself an undetermined function which gives a finite relation. $`\varphi `$ will therefor be, in general, a complex analytic function. Such functions correspond, as has been noted, to at most countably many zeros, and so $`\varphi `$ can give a relation between at most countably many data. The domain, and thus the range, of $`\varphi `$ also may thus, without loss in generality, be taken to be the Natural numbers. The distinct natural numbers making up the range of $`\varphi `$ will be referred to as “attributes”. These attributes naturally classify the action of $`\varphi `$ in the sense that $`\varphi `$(n) = $`\varphi `$(m) indicates that n and m have the same attributes.
It may thus be said that a finite relation $`\rho `$ corresponds to a particular at most countable variety of integral attributes which ”occur” in some sequence, together with the specification of a Turing Machine. The sequence itself corresponds to some finite relation which may be specified by an analytic function $`\varphi `$.
Now, for a fixed variety of attributes, any such function $`\rho `$($`\varphi `$, m) may be recoded as
$`\rho `$($`\varphi `$, m) = $`\omega `$($`\stackrel{~}{\varphi }_m`$) where (7.33)
$`\stackrel{~}{\varphi }_m`$(n)$`\varphi `$(T$`_\text{m}`$(n)) (7.34)
If there is to be such a thing as a finite relation $`\rho `$ which is independent of Relativity, and thus also of Turing Computability, then it must be the case that
$$\rho \text{(}\varphi \text{, m) = }\rho \text{(}\varphi \text{, s) for all naturals m,s .}$$
(7.35)
Such $`\rho `$ will be said to be “T-independent”. Some conclusions about T-independent functions will now be derived.
Consider those Turing Machines T<sub>m</sub> which permute initial segments of the sequence corresponding to $`\varphi `$. If $`\rho `$ is to be invariant under the action of all such Turing Machines then it must, first of all, be invariant under the action of all T<sub>m</sub> which permute the same initial segment of the sequence of attributes of $`\varphi `$. Then the dependence of $`\rho `$ on this initial segment must be expressible in terms of the relative frequencies of each attribute within this segment. If this is to apply, for a general finite function $`\rho `$, for all initial segments, then, first of all, the relative frequencies of each attribute must approach limits. $`\rho `$ will then have to be a function of these limiting relative frequencies. In particular, $`\rho `$ may be the limiting relative frequency of one of the attributes of $`\varphi `$. The limiting value, p(a), of the relative frequency of an attribute “a” for a T-independent function will be called its “probability”.
Now suppose T<sub>m</sub> yields an increasing function of n. Then, in general, $`\stackrel{~}{\varphi }_m`$(n) = $`\varphi `$(T<sub>m</sub>(n)) will have the effect of generating a new sequence in which certain of the original sequence will have been omitted based solely on the place, in terms of n, they occupied. Such “place selections” must not, therefor, alter the limiting relative frequency of an attribute if $`\rho `$ is to give a finite relationship independent of what Turing Machines may yield. Such sequences, being “insensitive” to place selections, will be said to be “Random”. It is of particular significance that ”random” sequences must be convergent as any particular segments of such sequences which might prevent convergence of the associated relative frequencies may be removed by place selections.
It should be pointed out that, for a given sequence to be random, it can in no sense be “known” term-wise, for the sequence is only to be acknowledged in the role which the theory assigns to it, and the theory considers a random sequence and any other sequence derived from it by place selections to be interchangeable. With this observation it becomes clear that while the argument of a probabilistic function may be described as a random sequence, it may more simply be described as a proper category, where a proper category is a metalabel which cannot be defined in terms of its corresponding assigns. To be clear, the unknowability of the random sequence entails the impossibility of uniquely determining $`\varphi `$.
From the above considerations it may be concluded that the quantities which T-independent finite functions may give as description beyond those obtainable from Turing machines may be thought of as being limiting relative frequencies of sequences of attributes. Furthermore, the sequences of attributes which these limits apply to must be random, that is, they must be insensitive to place selection. These two properties correspond, in an obvious way, to the axiomatic basis for Von Mises’s frequency theory of probability, although the theory presented here has a strictly conceptual basis and is not asserted to have an empirical basis as is that of Von Mises. It should also be stressed that this theory is distinct from the usual measure-theoretic formulation of probability<sup>11</sup><sup>11</sup>11 See . or the classical theory of probability<sup>12</sup><sup>12</sup>12 For a comprehensive survey of distinct theories of probability see . determines features which all theories of probability ought to share..
Recall that the question at issue here is the possibility of there being a finite theory that is independent of Turing computability. With this in mind it is clear that something must be said about the acceptability of infinite sequences necessarily playing a role in the theory. In this connection it must be recalled that it is the limiting relative frequency itself, and the randomness of the sequence, which are demanded by the theory, and not any particular infinite sequence of occurrences of attributes nor any calculations involving an infinite quantity; The calculation of this limiting frequency, which depends solely on $`\rho `$, need not require any such limiting processes to be carried out.
In fact, a finite T-independent theory treats probability as a “propensity”. This propensity, as defined by $`\rho `$ and not requiring the elaboration of a sequence of attributes, operates somewhat like a force field in classical physics: Just as a force causes a definite dynamical evolution, so a probability causes an evolution of the relative frequency of an attribute towards a definite limit. This action is unlike that of a classical force, however, in so far as the propensity governs the collective behavior of the whole sequence and may therefor be considered to act non-locally<sup>13</sup><sup>13</sup>13 Quantum theory may, in fact, be formulated as a strictly causal though non-local theory as was done by deBroglie and Bohm. See .. In this sense probability may be thought of as a non-local relation. It is only after viewing probability in this finitely-oriented way that a finite sequence of attributes may be explicitly introduced which is supposed to be governed by a law of probability. Then it may be argued, based on this, that as the sequence is extended under the given probability conditions its relative frequencies will approach particular limits. There can, in fact, be no cause to doubt that the sequence is random as the sequence is, after all, an artifact of the probabilistic theory itself.
One last remark on statistical theory in general. It will be noticed that the above theory, being based on the notion of T-independence, treats of discrete attributes. This is not, however, an essential restriction, for this probability theory may be extended, in the usual way, to probabilities over continuous attributes. Randomness may still play a role in the continuous theory as continuous distributions may be approximated by discrete ones on which the above notions of Randomness make sense. The discussion to follow will therefor retain the discrete perspective for the sake of ease of presentation and derivation.
### 7.6 Quantum Theory
The possibility of extending the theory of finite relations beyond what is given in Relativity was investigated in the last section. These considerations led to the conclusion that if such were possible, then the corresponding theory would have to be a probabilistic one: The foundations of the necessary form of such a theory were in fact given. Of particular interest was the conclusion that a finite theory of probability is necessarily a “realistic” one in which probability is a certain propensity for relative frequencies to approach particular limits. These were important but also very general conclusions. The goal of this section is to continue on to derive a more precise and quantitative theory. The natural development of the foregoing theory will lead uniquely to Quantum Theory, thus giving it a definite theoretical significance quite apart from any appeal to experimental justification.
The probability of an attribute has been defined to be its limiting relative frequency in a random sequence. Suppose that certain attributes a<sub>α</sub> are grouped together and all considered instances of the attribute a. Calculation of p(a) must then be such that the appearance of a<sub>α</sub> in the random sequence adds one to the frequency of occurrence of the a’s, and thus contributes to the relative frequency also. Consider an initial segment of length n. Suppose a appears n<sub>a</sub> times, while each a<sub>α</sub> appears n<sub>α</sub> times. Then n<sub>a</sub> = $`_\alpha `$ n<sub>α</sub>. Because all terms are positive, limits may be interchanged to derive
$$\text{p(a) = }\underset{n\mathrm{}}{lim}\frac{\text{n}\text{a}}{\text{n}}\text{ = }\underset{\alpha }{}\underset{n\mathrm{}}{lim}\frac{\text{n}_\alpha }{\text{n}}\text{ = }\underset{\alpha }{}\text{p(a}\text{α}\text{)}$$
(7.36)
As a special case of this suppose all attributes $`\alpha `$ are considered to be instances of the attribute u. Then
$$\underset{\alpha }{}\text{p(}\alpha \text{) = p(u) = 1.}$$
(7.37)
Suppose an attribute a is also considered to be a case of an attribute b, though not necessarily the other way around. Suppose also that p(b) $``$0. Consider, again, an initial segment of length n. Let n$`_{\text{a}|\text{b}}`$ be the number of attributes a found among the b’s that were found in the initial segment. Also let n$`_\text{a}`$ be the number of a’s found and n$`_\text{b}`$ the number of b’s found. Note that then n$`_{\text{a}|\text{b}}`$ = n$`_\text{a}`$, and if n$`{}_{\text{b}}{}^{}\mathrm{}`$ then it follows that n$`\mathrm{}`$. Then
$$\text{p(a}|\text{b) }\underset{\text{n}_\text{b}\mathrm{}}{lim}\frac{\text{n}_{\text{a}|\text{b}}}{\text{n}_\text{b}}\text{ = }\underset{\text{n}\mathrm{}}{lim}\left(\frac{\text{n}_\text{a}}{\text{n}}\right)\left(\underset{\text{n}\mathrm{}}{lim}\frac{\text{n}_\text{b}}{\text{n}}\right)^1\text{ = p(a)[p(b)]}\text{-1}\text{.}$$
(7.38)
so that it follows that
p(a) = p(a$`|`$b)$``$p(b). (7.39)
It should be stressed that in order to apply this rule to factor a probability it must be the case that the occurrence of a implies the occurrence of b.
The considerations above have indicated some relations, holding between the probabilities determined by a finite relation, that must be satisfied, but didn’t tell anything of how these probabilities were, themselves, to be determined. This will be addressed next.
In the determination of an appropriate function to represent a finite relation there will generally be certain data which is considered to be dependent on other data and primarily varies and other data which is presumed, and not considered to vary. Such is the case, in standard physical theories, with dependent variables on the one hand, and with independent variables, parameters, and “physical constants” on the other. Such a segregation of the data is not proven to be absolute, but, rather, is a matter of convenience in stating the relations which are utilized in the theory. For this reason, a means of indicating a segregation of data will also be adopted here.
It has already been shown above that finite relations must take a particular form. Probabilities, being determined by finite relations, must also take this form, and so must be given as functions of the magnitude, or, equivalently, of the square of the magnitude, of some complex analytic function. The particular way in which a probability is to depend upon the square of the magnitude of an analytic function must be well-defined and the same for all probability calculations but is otherwise unrestricted and so may be determined simply as a matter of convenience. It turns out to be convenient to choose the convention that probability be given, in fact, by the square of the magnitude of some complex analytic function. This choice results in the probability function being a measure<sup>14</sup><sup>14</sup>14 See pp. 70-72 of . and also avails the theory of a geometric analysis on a Hilbert Space, the given magnitude easily being constructed by an inner product<sup>15</sup><sup>15</sup>15 For the relevant functional analysis see, e.g., .. Combining the above observations results in the following notation and results:
In the context of certain presumed data which are symbolized by “i”, the probability, p(a), of a particular attribute a is given by the square of the magnitude of a certain analytic function, referred to from now on as the corresponding “probability amplitude”<sup>16</sup><sup>16</sup>16 For an exposition of Quantum Theory on this basis, see ., determined by the attributes under consideration and evaluated at a. This amplitude is denoted, as was done by Dirac<sup>17</sup><sup>17</sup>17 See ., by $`<`$a$`|`$i$`>`$. Thus
$$\text{ p(a) = }|<\text{a}|\text{i}>|^2$$
(7.40)
Now suppose that p(a)=$`|<`$a$`|`$i$`>|^2`$, p(b)=$`|<`$b$`|`$i$`>|^2`$, and p(a$`|`$b)=$`|<`$a$`|`$b$`>|^2`$. Applying equation(7.39) results in
$$<\text{a}|\text{i}>\text{ = }\text{e}^{ı\theta }<\text{a}|\text{b}><\text{b}|\text{i}>\text{.}$$
(7.41)
Here it is to be noted that $`\theta `$, if it need be determined, must yet be determined by other considerations. In fact, it is a standard insight of Quantum Theory that such exponential factors allow for, and demand, the conservation of such quantities as energy, charge, angular momentum, etc.<sup>18</sup><sup>18</sup>18 See vol. III of .. Also, it should be reiterated that this result applies, as does (7.39), only in the case that the occurrence of a implies the occurrence of b. Conversely, whenever an amplitude may be factored into a product of amplitudes it will have to be the case that any corresponding probabilities will be for attributes which are related in this way.
The only other relation from probability theory is equation(7.36). It applies in the case when a variety of attributes a<sub>α</sub> are to be considered to be instances of another attribute a. This leads to the relation
$$|<\text{a}|\text{i}>|^2\text{ = }\underset{\alpha }{}|<\text{a}_\alpha |\text{i}>|^2\text{.}$$
(7.42)
This rule is evidently not very helpful in determining $`<`$a$`|`$i$`>`$ itself, although it does lead to the conclusion, following the case in (7.37), that $`|<`$u$`|`$i$`>|^2`$ =1. In other words, the “total” amplitude must be normalized. Thus the probability rules offer little help, but also little hindrance, in the determination of the amplitudes that will play a role in the theory. The only possible approach that remains in determining a constructive theory is thus to analyze amplitudes according to the simple fact that they are given by analytic functions.
A defining characteristic of analytic functions is that they may be expressed as power series sums, and such series may be manipulated term-wise, and analyzed in this way. Consider, then, an amplitude and a power series expansion of it. Each of the summands, in order that it may be interpreted in the theory, must be somehow related to an amplitude. Each term may, in fact, be formally analyzed according to the rule in equation (7.41), so that, in general, an amplitude $`<`$a$`|`$i$`>`$ may be analyzed as
$$<\text{a}|\text{i}>\text{ = }\underset{\beta }{}\text{e}^{ı\theta _\beta }<\text{a}|\beta ><\beta |\text{i}>\text{.}$$
(7.43)
It must be remembered that the $`\beta `$ are introduced here in a strictly formal sense, for if the occurrence of a were to imply the occurrence of any such $`\beta `$ then application of equation (7.41) would yield
$$<\text{a}|\text{i}>\text{ = }\text{e}^{ı\delta }<\text{a}|\beta ><\beta |\text{i}>\text{.}$$
(7.44)
and thus the sum in equation (7.43) would have to either degenerate into this single term or else lead to contradictions if more than one distinct such $`\beta `$ were to be found. On the other hand, if, starting from a situation described by the amplitude $`<`$a$`|`$i$`>`$ given by (7.43), an event described by a particular $`\beta `$ is considered to occur, then the amplitude for a would be of the form $`<`$a$`|\beta >`$ and the transition between these situations could, to needlessly indulge in a “physical” description of what may be considered a logical process, be said to be a case of a “collapse” of the “wave function” $`<`$a$`|`$i$`>`$. It should be remembered that the theory is derived here under the assumption that there is no separate realm of experience which may affect the developments within the theory.
If the expansion (7.43) is not to refer to any necessary intermediate process $`\beta `$ in the description of the attribute a, then there must be a great deal of freedom in the choice of such formally chosen “virtual” attributes. The selection of any particular expansion must, therefor, be justified on the basis of calculational advantages. The benefit of this analysis consists in two effects. First of all, each summand may be simple and so much so, in fact, that the values of such amplitudes may be assumed in the construction of a calculational theory. Such is, in fact, done in the case of the so-called ”bare” mass and ”bare” charge of the electron in Quantum Theory. It may also happen, as a second benefit, that the selection of a particular expansion may result in rapid convergence of the sum (7.43) so that it might even be terminated after a few terms.
At this point it is clear that the above considerations have completely reproduced Feynman’s formalization of Quantum Theory, including a justification of the central role played by the Feynman diagram approximation technique<sup>19</sup><sup>19</sup>19 See for the quantum formalism, and consult for a discussion of the Feynman diagram technique. but with the advantage of showing that corresponding graph-theoretical concepts naturally arise in the algebraic analysis of $`\varphi `$. It should be noted that this approach, in basing all calculations on a single history, is unsuitable for addressing the intrinsic compatibility of quantum theory with “multiple observers”. However, because this approach is avowedly applied as an approximation, it is also not necessary to do so, for any effects associated with this kind of situation may always be relegated to the uncalculated terms in the approximation.
Another formal consequence of the analytic nature of the amplitudes is of paramount importance when it is desired that exact calculations may, in principle, be made. Because expansions of analytic functions, such as in (7.43), may be manipulated term-wise in the case of many mathematical operations, such operations will be linear in the amplitudes. Thus the theory given here may also, for those that prefer to do so, be translated into the coordinate-free operator formalism which Von Neumann gave to Quantum Theory<sup>20</sup><sup>20</sup>20 Again, see ., with one important difference: All transformations of state must now be represented by application of a (Hermitian) linear operator as there can no longer be special procedures for representing “measurements”. This then leads uniquely to a quantum theory which is formally identical to that of Everett’s Many Worlds approach<sup>21</sup><sup>21</sup>21 See pp. 3-149 of for an exposition of this theory. although it must be stressed that in the context of this thesis there is no implication of, nor can there be, histories and/or universes which are anything but formally separate. As all changes of state are of the same formal nature, it would be unacceptable, besides merely troubling or unbelievable, for any to be especially felt as a “split of the universe” or a “branch in one’s history”.
Everett’s derivation of his theory assumed that a Schrödinger equation was provided for the time evolution of the state, but the theory may be derived without this assumption. As Everett has shown, there is, corresponding to a sequence of attributes which defines a history, a corresponding sequence of applications of corresponding operators. Each history corresponds, exactly as in Feynman’s sum over histories approach to Quantum Mechanics<sup>22</sup><sup>22</sup>22See for an exposition of this theory., to a summand in the superposition (7.43) which formally represents an amplitude. In this way the superposition may be thought of as a net of histories.
Begin with some operator H which is to be a “Hamiltonian” which gives the evolution of what is to be considered, by definition of this choice, an “isolated” system. In that case all other operators must represent, by definition, “interactions”. In examining any history one may then refer to “change” whenever H alternates with some other operator. One may then construct a corresponding notion of time for each individual history where distinct points in time are identified by such transitions while interactions shared with other histories identify related moments.
Everett has shown<sup>23</sup><sup>23</sup>23 Refer, once more, to pp. 70-72 of . that there is a measure on histories which agrees with the probability quantum mechanics assigns to them. Thus, it may be observed here that a direction of a history which progresses towards the appropriate limiting frequencies may be identified with the future along that history. This then accounts for the ordinal nature of time. If it is desired that time also possess a measure, then this may be achieved by selecting some operator, T say, distinct from H and identifying any successive interactions involving T with unit intervals in time along that history. The rate of change will then correspond to the likelihood of such interactions.
While all formal histories which correspond to terms in (7.43) play a role, they weigh in according to the measure assigned to them - in other words, according to their probability: This is the quantal basis for the variational form of the laws of physics. It is thus possible to dispense with relatively unlikely histories and, conversely, any initial conditions which may play a role in theory must be sufficiently likely. In this sense initial conditions no longer enter into theory as independent data. This makes sense as, after all, the distant past, in not being accessible to personal experience, has always been reconstructed in light of current understanding or theory.
The whole character of the theory is determined, in fact, by H, because identifying H is a matter of identifying when, by contrast, something may happen. If an operator has a particular symmetry, then the history corresponding to it will appear with a corresponding degeneracy, or repeated appearance, in the superposition (7.43) so that symmetries correspond to prevalent terms and thus to significant histories. This explains why the search for the laws governing quantum-scale interactions has been successfully guided by considerations of symmetry. Although this account of quantum theory is by no means a usual one, the formal rigor of all mathematical steps has already been born out in the references cited and the work done above, with additional confirmation to be found in the history of the development of quantum theory.
The above analysis is self-contained and fully determinate excepting that it remains to explain how the Hamiltonian operator arises as, surely, not just any one will do! Such an explanation will, however, proceed from the following discussion which links the above ideas with formally thermodynamic considerations<sup>24</sup><sup>24</sup>24 For an exposition of the required Thermodynamic theory see, for example, vol. I of ..
First of all, all component histories will, by definition, converge in their futures to the respective probable states corresponding to the analytic state function $`\varphi `$ which finitely defines the system. This being true individually, it also follows that all histories proceeding from a single state $`\varphi `$, and thus forming a system, will also converge towards mutual equilibrium. This limiting condition is that which is guaranteed by the so-called Zeroth Law of Thermodynamics, wherein a temperature function is introduced to express the conditions of relative equilibrium holding between any collection of systems. That the progress towards the limit identifies a direction in which time “flows” in a particular history is the essential content of the Second Law of Thermodynamics.
The First Law of Thermodynamics is that of Energy Conservation. It has, however, already been noted, in the discussion following equation (7.41), that energy is conserved in quantum theory. It is also well-known that the boundedness, and general smallness, of the ground state energy of quantal systems leads to the Third Law of Thermodynamics. It may be therefor concluded from the above analysis that all of the laws of Thermodynamics follow within this quantum formalism.
Quantal systems being therefor thermodynamic systems, it follows that they may be analyzed thermodynamically. Since the theory presented here is strictly formal, it necessarily follows that the ensembles used in the analysis must be (virtual) Gibbs Ensembles. As is always the case in the thermodynamic analysis of a system, the properties of the thermodynamic system are completely defined by the (Grand) Partition Function which corresponds to the Gibbs Ensemble. Such “bulk” properties as the volume and pressure exerted by the system are, in fact, expressible, in a standard way, in terms of the partition function<sup>25</sup><sup>25</sup>25Refer, again, to .. Now note: Both the partition function and $`\varphi `$ completely determine precisely the behavior of the system so that these functions are interchangeable! Thus $`\varphi `$ will be taken to be the partition function.
Identifying $`\varphi `$ to be the partition function has immediate consequences. First of all, as $`\varphi `$ is necessarily a complex analytic function, it naturally follows that this is also the proper way of thinking about the partition function. In this case one no longer introduces the complex “fugacity”<sup>26</sup><sup>26</sup>26 Refer to for an explanation of this concept. into the partition function as a formal trick but, instead, the question rather becomes how one starts with a complex-valued partition function with complex arguments and ends up only concerned with a real-valued partition function with real-valued arguments. In this connection it should be recalled that $`\varphi `$ is constructed so as to encode certain corresponding complex zeroes. The nature of the events which may occur in the given thermodynamic system is therefor determined by the location of these zeroes: There must be certain criteria as to the situation of these zeroes the satisfaction of which results in singular phenomena.
At this point it is appropriate to introduce, as a formally necessary development of the finite case rather than as a merely plausible addition, a classic result of Yang and Lee<sup>27</sup><sup>27</sup>27 See . which states that any ambiguity in the value of certain “bulk” variables, such as density, and the corresponding changes of “phase” of the system, such as from liquid to solid, must occur precisely at points on the real axis which are not separated, by any neighborhood in the complex plane, from the zeros of $`\varphi `$ (the partition function). This shows, to begin with, that the variables which are to represent points of transition between phases must be real: This justifies requiring all variables, in the end, to be real-valued.
According to this result it also follows that, as the positioning of the zeroes of $`\varphi `$ corresponds to condensation phenomena, these zeroes also define “material structure” in concrete and tactile terms such as pressure, volume, and temperature. Such relatively stable material structure is a candidate for treatment as an isolated system. The symmetries of such systems then will identify corresponding Hamiltonians.
This self-contained Quantum-Theoretical analysis, like that of the Relativistic analysis before it, reduces the problem of the variety of concrete experience to that of the problem of the “existence” of any phenomena whatsoever. In other words, if $`\varphi `$ has any zeroes at all then this automatically engenders a corresponding physical picture, operating according to the quantal laws, comprised of interacting phases of matter, such as solids, liquids, and gases, each occupying certain volumes and “influencing” one another barometrically and thermally or in terms of other recognizable thermodynamic variables (Such as, for example, the chemical potential).
In reviewing the considerations above it is clear that the attempt to construct a finite theory beyond that which is determined by Relativity has resulted in an essentially unequivocal derivation of the mathematical formalism of Quantum Theory. It has done so without appeal to experiment and without accepting the notion of an external realm of experience to which the formalism need refer. In doing this it has been demonstrated that the great variety of physical interpretations of the quantum formalism are, presumably, unnecessary, for from a finite strictly mathematical standpoint a unique propensity interpretation of probability and of Quantum Theory arises.
### 7.7 Conclusions on the Derived Theories
In the preceding it has been found possible to derive both Relativity and Quantum Theory on a strictly conceptual basis. That these derivations may have been performed without appeal to experiment obviously undercuts the notion that these theories must be thought of as describing or being associated with an external physical reality separate from language.
Again, no appeal to experiment has been made, but, rather, the principle aspects of the phenomena which these theories address have been shown to follow from strictly formal principles; It should be understood that, instead of experiment justifying theory, it has been found that the scope of thinkable finitary theory has shown why such such phenomena presumably must occur. In particular, Relativity Theory has been argued to uniquely reflect relations which maintain one-to-one relabeling within the finite general symbolic system. Quantum Theory was subsequently argued to reflect every and anything else that may be called a finite relation which is not justified on a Relativistic basis. In particular, it has been indicated that Quantum Theory, instead of treating of relationships between individual symbols, yields rather all finite relations which govern proper categories or, in other words, are based upon metalabeling relationships as such.
As it has been argued that consistency in the utilization of individual symbols, together with the conveyance of metalabeling relationships between symbols, exhausts all that may be required of language, it is to be expected that Quantum Theory, in virtue of the nature of its construction, should, together with Relativity Theory, provide a complete and consistent description, if such is possible, of all finite relations within the general symbolic system. This possibility is precisely what is investigated in the next section.
### 7.8 Einstein’s Alternatives and Finite Reality
Relativity and Quantum Theory should together, if such is possible, comprise a complete finite theory. It will be found, however, that they are not merely distinct but also, in the context of a finite restriction of the general symbolic system, incompatible. This incompatibility will be shown by an argument similar to one that Einstein gave<sup>28</sup><sup>28</sup>28 See pp. 168-173 of .. The possible ways of resolving this incompatibility will be investigated and this will lead, on the basis of the general symbolic system, to plausible far-reaching conclusions about physical theory.
The general symbolic system now has, in Quantum Theory and Relativity, two distinct formal procedures which may be applied to the data which comprise any given “physical situation”. The question now arises as to which parts of the data each of the theories is to apply. Is it possible, in other words, to recognize distinctly “classical” data to which Relativity alone applies?
Relativistic phenomena are characterized by the energy-momentum tensor, which is, it may be recalled, the right-hand side of equation (7.22). Conversely, equation (7.22) indicates that this tensor is a function of the space-time itself, so that a flow of energy must correspond to changes in the space-time manifold. As quantum processes are not energetically isolated it follows, therefor, that they too must be dependent, at least in part, on the space-time variables. This influence may also be seen to go the other way, as it is well known that, in accordance with the Third Law of Thermodynamics, Quantum Theory results, in a number of cases, in predictions that agree with those of Classical Physics and, in fact with, Relativity. This is also the content of Bohr’s Correspondence Principle.
Because of this limiting behavior there can be no finitary conceptual distinction in general between a result produced by quantal laws and a result which is the outcome of local deterministic processes. In short, as is usually taken as a physical necessity, Quantum and Relativistic processes “interact”: They may not be considered to be disjoint physical theories which don’t share common variables upon which they depend. Thus quantum-theoretical amplitudes must generally depend, in part, on space-time variables.
Let a be an attribute which indicates something about events at a particular space-time point P<sub>1</sub>. Also let $`<`$a$`|`$i$`>`$ be the amplitude associated with this attribute. Recall that, in a finite probabilistic theory, this amplitude must be considered to be a kind of propensity, changes of which therefor have an immediate space-time significance.
Now consider some attributes $`\beta `$ which aren’t implied by a so that equation (7.43) may be applied. Then
$$<\text{a}|\text{i}>\text{ = }\underset{\beta }{}\text{e}^{ı\theta _\beta }<\text{a}|\beta ><\beta |\text{i}>$$
(7.45)
holds. Note that $`\beta `$ could represent almost anything.
Consider an attribute $`\beta `$ which indicates something about events at a space-time point P<sub>2</sub>. Further suppose that the interval from P<sub>1</sub> to P<sub>2</sub> is positive. Now the occurrence of $`\beta `$, according to (7.43), requires the formal transition $`<`$a$`|`$i$`><`$a$`|\beta >`$. If thought about in physical terms, such a change in the amplitude is an immediate “reality”, so that a change at P<sub>2</sub> entails an instantaneous change at P<sub>1</sub>. But, even if viewed in a strictly formal sense, this is obviously a violation of the restriction to local interactions in Relativity Theory. Thus, while Relativity and Quantum Theory cannot be kept apart, they are nevertheless incompatible.
There are exactly two assumptions that have been made in the construction of these finitary physical theories; That there is no necessity for a separate realm of experience apart from language itself, and that all descriptive quantities in the theory must be finite. At least one of these assumptions must be dropped.
If the first assumption is dropped, then, as noted before, the permanent role of experiment entails that there can never be a final understanding or ultimate physical theory. If the second assumption is dropped then it will be necessary to have a mathematics that can naturally handle infinite quantities, as was attempted in the physical programme known as “renormalization”<sup>29</sup><sup>29</sup>29 See ., and interpret such quantities in a sensible fashion. In order to maintain the derivation of the theory of the general symbolic system without recourse to experiment the second option will next be explored.
### 7.9 On Generalizing Finite Theories
It may have seemed at the time that consideration of infinite quantities and the development of a Calculus which doesn’t depend upon finitely defined limits was an extravagance, but it is clear now that there was a definite justification for these steps. In fact, had the Calculus been restricted to the usual finitely defined operations it would have been impossible to identify the cause of the inconsistency, in a non-empirical setting, with simply the finiteness of the variables, for the operations of the Calculus also went into the derivation of the theory and might have been connected with the difficulties.
Having already laid the groundwork, the extension of the previous theory to the infinite realm is relatively simple. It has already been noted that the operations of the generalized Calculus extend those of the usual Calculus to the infinite case, so the Calculus may still be applied in the derivations of exact theory. Suppose now that some equation is to hold in the potentially infinite case. This same equation must still apply even if all of its arguments are finite, but in that case the given equation would have to be part of the theory for the finite case. It thus follows that the general theory for the manipulation of symbols is formally identical to, in fact it is determined by the same equations (7.22), (7.40), (7.41), and (7.43) as govern the finitely-based theory which has already been derived above. The only difference between the finite and non-finite cases is, therefor, merely the potentially non-finite nature of the variables employed in the theory.
While an accurate conclusion, it is still not very informative to merely imply that the variables in the general symbolic system are (generally) non-finite analogues of the finite space-time coordinates and probabilities which characterize the usual Relativistic and Quantum-theoretical formalisms. In order to more definitely relate the general symbolic system to its finite restriction it is therefor desirable to introduce the methods of Non-Standard Analysis. These methods allow for the rigorous definition and algebraic manipulation of non-finite quantities in accordance with the rules of a division ring, as well as for a canonical representation of all such quantities as a unique sum of finite and non-finite parts.
Non-Standard Analysis, as previously noted, has been rigorously developed elsewhere, so it is neither necessary nor desirable to go into much detail in explaining or justifying these methods. It should be helpful, however, if an intuitively plausible exposition is provided which, moreover, shows that the definition of the extension of the finite real number system employed in non-standard analysis may reasonably be thought of as arising from a generalization, to a non-finite case, of the notions already introduced in the analysis of the finite case.
The most general kind of relation in the general symbolic system is that which allows for the non-invertible manipulation of symbols. As was previously indicated, such relations were, in the finite case, determined by the formal rules of Quantum Theory. These rules still apply and will determine the relations which define the non-standard extension of the reals.
Recall the Quantum-theoretical rule (7.36) and its special case (7.37): These rules, considered respectively, require the finite additivity of probability and that probabilities should be real numbers in the unit interval. Considered without regard to the notion of randomness, the development of the theory of finite relations, in accordance with these rules, leads to the construction of a finitely additive measure on sequences. Given such a measure, m, and any two sequences, {a<sub>n</sub>} and {b<sub>n</sub>}, the measure may be used to quantify the degree of term-wise agreement between any two such sequences: m{n$`|`$ a<sub>n</sub> = b<sub>n</sub>} gives this measure.
In the finite case the desired T-independence of this measure requires randomness of such sequences. It may be recalled that the relative frequencies of random sequences necessarily converge to a limit, which is the probability, and that any two random sequences which correspond to the same measure are interchangable and correspond to the same probability. In this sense sequences are, as a natural development within the finite general symbolic system, “identified” according to their having the same limit, so that limits play the role of measure in this case.
The identification of sequences with common limits suggests the construction of equivalence classes $``$ a<sub>n</sub> $``$ defined by $``$ a<sub>n</sub> $``$ $``$ {{b<sub>n</sub>} $`|`$ lim<sub>n→∞</sub> (a<sub>n</sub> \- b<sub>n</sub>)=0 }. The equivalence between sequences which is thus defined with respect to limits is a very liberal one, there being very many sequences which converge to a given limit. In other words, probability, thought of as a measure, distinguishes relatively few sequences because it neglects the mode of convergence of such sequences.
Algebraic operations may be defined upon sequences, though if such operations are to be relevant to, and reliably reflected in, the probabilities associated with such sequences, then these operations must be defined term-wise. Therefor the definitions {a<sub>n</sub>}$``${b<sub>n</sub>} $``$ {a<sub>n</sub>$``$b<sub>n</sub>} and {a<sub>n</sub>}+{b<sub>n</sub>} $``$ {a<sub>n</sub>+b<sub>n</sub>} are made. With these definitions it may be confirmed that the equivalence classes corresponding to such sequences also satisfy these equations.
The above construction of limit-based equivalence classes is reminiscent of the construction of the real number system as a collection of equivalence classes of Cauchy sequences of rational numbers. In this derivation algebraic operations on the equivalence classes are like-wise well-defined in terms of the corresponding term-wise operations on the elements of representative sequences from each of the classes. In order to generalize the real numbers to the non-finite case it therefor is natural to attempt a generalization of this procedure by replacing the “limit measure” with the least restrictive measure conceivable.
As indicated above, a measure categorizes a sequence according to the value it assigns it. In order to make the greatest number of sequences distinct it is necessary to allow the measure to take on the least range of values. This identifies the measure to be employed in the extension of the reals as the discrete measure, where this measure only takes on either of the two values 0 or 1 for any subset of the integers. In order that this measure lead to a number system which extends the reals it is necessary that it assign the value 0 to all finite subsets of the integers. This corresponds to the indifference of the “limit measure” to any finite portion of a sequence. This requirement also leads to the conclusion that infinite subsets of the integers have discrete measure 1.
Given such a discrete measure m defined on the integers, a new definition of equivalence classes of sequences may be specified in analogy with that already determined by limits. In particular: $``$ a<sub>n</sub> $``$ $``$ {{b<sub>n</sub>} $`|`$ m{ n $`|`$ a<sub>n</sub>=b<sub>n</sub>}=1 }. When this is done it is easily confirmed that, as before, the addition and multiplication of equivalence classes is well-defined. It is also clear that the real numbers may be embedded in this scheme with the mapping a$``$a$``$, which takes the real number a to the equivalence class of sequences which contains a sequence all of the terms of which are identically a. The above construction therefor extends the real number system. This extension is referred to as the hyper-real number system.
Unlike before, however, the discrete measure takes into account the whole sequence, and not just its limit, so that it categorizes convergent sequences also according to their modes of convergence. To see this in detail it suffices to extend the notion of order to the new equivalence classes: Take $``$a<sub>n</sub>$`<`$b<sub>n</sub>$``$ exactly wherever m{n $`|`$ a$`{}_{n}{}^{}<`$b<sub>n</sub>}=1. Then it may be seen that, according to this definition, $`\frac{1}{n^2}<\frac{1}{n}`$ and $``$n$`<`$n$`{}_{}{}^{2}`$ are examples of inequalities which hold in the hyper-real number system. As may be expected, these inequalities are between what may be referred to, in a precise way, as infinitesimal and infinite numbers, respectively. Indeed, a finite number x may be defined to be one which is bounded by some positive real a so that -a$`<`$x$`<`$a. Numbers which are not finite are termed infinite. Infinitesimals are taken to be numbers bounded by all positive real numbers.
The foregoing makes it clear that the hyper-reals constitute a proper extension of the real number system. Beyond this, the formal relationship between the hyper-reals and the embedded finite real number system may readily be made clear<sup>30</sup><sup>30</sup>30Rigorous proof of the following and other results of Non-Standard Analysis may be found in . See, in particular, pp. 1-105 there for an introduction to non-standard methods along the lines of the presentation here. Consider a finite hyper-real x. Let a be the least upper bound of the real numbers less than x. Then $`ϵ`$=x-a is infinitesimal. The decomposition x = a+$`ϵ`$ is unique. Thus all hyper-reals may be expressed as a unique sum of finite and non-finite numbers.
With this result it is possible to have an at least intuitive understanding of the hyper-real number system and how it relates to the previously developed finitary theory. It now remains to apply this general understanding to see how working in the context of the hyper-real number system might resolve the previously mentioned incompatibility between the finite Relativistic and Quantum-theoretical formalisms.
Consider any combination of Relativistic and Quantum-Theoretic data in the non-finitary general symbolic system. Express this data as canonical sums of finite and non-finite numbers. The standard (finite) parts of the data may then be identified with the data upon which the finitary theory holds. As previously discovered, such restricted data cannot form the basis of a satisfactory physical theory in the general symbolic system.
Such data, while giving a finite level of description, do not, however, give the “whole picture” in the non-finitary case. The residual parts of the data, due to their infinitesimal or infinite nature, while either empirically entirely inaccessible or irrelevant in principle in the context of a finitary theory, are, nevertheless, critical in resolving the apparent incompatibility of Quantum Theory and Relativity. Two numbers which differ, even if only by an infinitesimal amount, are different, and the differences to be found in the residual parts of data are presumably the only means of formally rejecting the arguments of the last section.
Put more concretely, the non-local correlations which Quantum Theory formally demands but finitary Relativity forbids are expressed by quantities neglected in the finitary theory. It may be said, in this sense, that the elaboration of finitary theory, while not adopting any axioms, has nevertheless resulted in a level of description. Changes in such neglected quantities amount to signals which may travel at infinite speed and yet are not incompatible with the finitary theory. This is apparently the unique way, String Theory not withstanding, in which the Quantum Theoretical and Relativistic formalisms may be unified in the general symbolic system. This unification also suggests a resolution between the corresponding conflict between a preference for a discrete or a continuous mode of description of experience: The non-finitary general theory comprises a third kind of description which chooses neither over the other, and in this context neither probabilistic nor deterministic theory takes precedence over the other or is wholely adequate.
As a final observation, the role which Non-Standard Analysis plays in presumably reconciling Relativity and Quantum Theory as parts of the general symbolic system suggests that an identification may be made between the notion of the externality of the “physical reality” which finite empiricism forever appeals to and the presumed conceptual inadequacy of any finite structures within the general symbolic system. In particular, there is a similarity between these two approaches in that finitely based empirical theories are subject to arbitrarily many experimental revisions of their data, while the finite structures of the general symbolic system are themselves subject to infinite elaboration. If this identification is made then, it is proposed, these two approaches may be thought of as being indistinct. While it would be inconsistent with the very notion of a non-axiomatic theory to demand such an identification, nevertheless it is hoped that such an identification might make rejecting serious consideration of such systems less likely.
## Chapter 8 DISCUSSION AND CONCLUSIONS
The object of this thesis has been the development of a written discrete combinatorial symbolic system based on formal rules for the manipulation of symbols which apply without regard to the particular meaning of the symbols involved. Motivation for the particular rules to be adopted was sought through examination of some contemporary scientific theories. Beginning with discussions of Linguistics, Metamathematics, and Physics it was found that, although each proceeded, in part, from a formal axiomatic basis, there were still a number of difficulties associated with the incorporation of an additional intuited meaningful aspect of experience. The discussion of Metamathematics, however, led to the suggestion that the notions of the truth and the formal provability of any particular assertion were, in part, determined by the formally unrestricted choice of an interpretation. Consequently it was suggested that perhaps it is unnecessary and undesirable to make such an apparently arbitrary choice. Perhaps, instead, the formalism could be shown to yield its own interpretation and thus avoid the possibility of the formal rules conflicting with intuition.
In order to further motivate and justify this approach a schematic overview of physics in which, beyond the idealistic presentation of the introduction, a practical description of the activities of physicists and their implicit philosophies, as well as a discussion of the history of the development of Physics itself, was given. This presentation clarified the nature of the conflict between formalism and intuition and indicated the way in which this difficulty has been addressed, though by necessity impermanently, through the adoption of either statistical or non-statistical theories. As one is always free to adopt theories of either character, it was next attempted to take a step back to a stage prior to theoretical analysis in order to seek a basis for a formal theory.
The first stage of the scientific procedure, prior to that of theory construction, is that of observation and description. At this point it was argued, aiming at plausibility rather than proof, that a number of the features of experience which seem to be merely intuited may nevertheless be shown to derive from a merely formal basis. It was argued, in fact, that while any particular grouping, or categorization, of symbols( such a process being a necessary prerequisite for any formal theory ) cannot be defended a priori, nevertheless the mere faculty of being able to adopt some such description leads to the utilization of realistic structures within language. It was therefor decided that, even on an intuitive level, the formalism should be based not on the expression of “truths” or “facts”, but rather constructed in order to express contingency, in which case a theory may be judged completely satisfactory if its elements are convincingly inter-related and the structure of the language is rich enough to be descriptive.
Such a formal theory was then presented. It was shown that in such a theory the freedom of labeling consequent to the absence of facts leads to groupings of symbols, here called metalabelings, and to a further facility for relabeling which, almost paradoxically, enables one to conventionally reserve symbols to be used in particular ways. Such considerations were then applied to the development of the so-called “general symbolic system”. This system takes the algebraic form of the real quaternions, the system inheriting a partial order based on metalabeling, though it also naturally admits of infinite and infinitesimal quantities. In order to demonstrate the feasibility of such a system it was necessary to next develop a generalization of the Calculus which is not based on $`ϵ`$-$`\delta `$ arguments. With this it was possible to complete the development of the formal system, although it remained to demonstrate its relevance by showing its utility in application to Physics.
Physical theory was developed in accordance with the requirements of empiricism which dictate that physical description may be given in terms of finitely many finite variables. Such a starting point was invoked in basing all theory on the location of the zeroes of a complex analytic state function $`\varphi `$. Theory was then developed along two lines. First it was shown that Relativity Theory, including the space-time structure to which it applies, resulted if it was desired that all operations in the theory should express invertible relations. As was also indicated, it turns out that Turing computability is a notion which may be thought of as being subsidiary to Relativity Theory.
Relativity presumably being the complete expression of determinism, extension of the theory led to probabilistic developments. It was then argued that the propensity-based probability theory arrived at, expressing relations between proper categories, was, in fact, Quantum Theory. As was to be expected, this derivation also serves to justify the assertion that Quantum Theory is an intrinsically probabilistic theory. The strictly formal basis of this derivation bypassed the notorious “measurement problem” of quantum theory and resulted in a formalism and descriptive apparatus formally similar to that of Feynman’s Sum Over Histories and Everett’s Many Worlds interpretations of quantum mechanics. The resulting perspective is unique, however, in indicating that all features of the dynamical and thermodynamical laws and description, including those materialistic phenomenological aspects of individual experience described by pressure, volume, and temperature, may be thought of as simply being consequences of basing theory on the finitary state function $`\varphi `$.
The finite application of the general symbolic system had, at this point, reasonably justified its relevance in comparison with the usual empirical epistemology. Investigation of its internal coherence led, however, to the interesting conclusion that retaining a formally finite basis for theory required that the formalism be supplemented by something outside itself, such as is done in experiment, and yet does not result in a conceptually coherent description of experience. Conversely, retaining a strictly formal perspective apparently requires rejecting finitism and therefor utilizing symbols which cannot be empirically verified if it is assumed that measurements are represented by finite collections of finite numbers. Extension of the formal theory in such an infinitary approach turns out to be surprisingly easy and is perhaps the only way to achieve a desired unification of the Relativistic and Quantum-theoretical formalisms.
## Appendix A THE SOLUTION OF FUNCTIONAL EQUATIONS
### A.1 The Associativity Equation
The equation to be solved is
F\[x,F\[y,z\]\] = F\[F\[x,y\],z\]. (A.1)
From (A.1) may be derived
$`_1`$F\[x,F\[y,z\]\] = $`_1`$F\[F\[x,y\],z\]$`_1`$F\[x,y\] (A.2)
and
$`_2`$F\[x,F\[y,z\]\]$`_1`$F\[y,z\] = $`_1`$F\[F\[x,y\],z\]$`_2`$F\[x,y\]. (A.3)
Note that the order of the factors is unimportant because both terms in each product are functions of the same variables and therefor commute. For F not constant in any of x, y, or z it follows that $`_1`$F\[x,y\] $``$0 and $`_2`$F\[x,y\] $``$0. Now (A.2) and (A.3) may be combined in the form
$`_1`$F<sup>-1</sup>\[x,y\]$`_2`$F\[x,y\] = $`_1`$F<sup>-1</sup>\[x,F\[y,z\]\]$`_2`$F\[x,F\[y,z\]\]$`_1`$F\[y,z\]. (A.4)
This expression may be simplified by letting
G\[x,y\] = $`_1`$F<sup>-1</sup>\[x,y\]$`_2`$F\[x,y\]. (A.5)
Then follows
G\[x,y\] = G\[x,F\[y,z\]\]$`_1`$F\[y,z\] (A.6)
and
G\[x,y\]$``$G\[y,z\] = G\[x,F\[y,z\]\]$`_2`$F\[y,z\]. (A.7)
Now
$`_z`$G\[x,y\] = 0 =$`_2`$G\[x,F\[y,z\]\]$`_2`$F\[y,z\]$`_1`$F\[y,z\] + G\[x,F\[y,z\]\]$`_{21}`$F\[y,z\] (A.8)
and
$`_y`$\[G\[x,y\]$``$G\[y,z\]\] = $`_2`$G\[x,F\[y,z\]\]$`_1`$F\[y,z\]$`_2`$F\[y,z\] + G\[x,F\[y,z\]\]$`_{12}`$F\[y,z\]. (A.9)
Subtracting (A.8) from (A.9) results in
$`_y`$\[G\[x,y\]$``$G\[y,z\]\] = G\[x,F\[y,z\]\]$``$\[$`_{12}`$F\[y,z\] - $`_{21}`$F\[y,z\]\]. (A.10)
Now $`_{12}`$F = $`_{21}`$F so that
$`_y`$\[G\[x,y\]$``$G\[y,z\]\] = 0. (A.11)
This may be solved by
G\[x,y\] = $`\beta `$H(x)$``$H<sup>-1</sup>(y) (A.12)
where H is any function such that H(y) $``$ 0 and $`\beta `$ is a constant. Now from this and (A.6) it follows that
$`_1`$F\[y,z\] = H(F\[y,z\])$``$H<sup>-1</sup>(y). (A.13)
Similarly, (A.7) and (A.12) yield
$`_2`$F\[y,z\] = $`\beta `$H(F\[y,z\])$``$H<sup>-1</sup>(z). (A.14)
Now define the function $`\mathrm{\Phi }`$ by the requirement that
$`\mathrm{ln}`$($`\mathrm{\Phi }`$($`\omega `$)) = $`_1^\omega `$H<sup>-1</sup>($`\tau `$)d$`\tau `$. (A.15)
Then
$`\mathrm{ln}`$($`\mathrm{\Phi }`$(F\[y,z\])) = $`_1^{\text{F[y,z]}}`$H<sup>-1</sup>($`\tau `$)d$`\tau `$ and $`\mathrm{ln}`$($`\mathrm{\Phi }`$(y)) = $`_1^\text{y}`$H<sup>-1</sup>($`\tau `$)d$`\tau `$. (A.16)
Thus $`_y\mathrm{ln}`$($`\mathrm{\Phi }`$(y)) = H<sup>-1</sup>(y) and, by (A.13), it follows that
$`_y\mathrm{ln}`$($`\mathrm{\Phi }`$(F\[y,z\])) = H<sup>-1</sup>(F\[y,z\])$`_1`$F\[y,z\] = H<sup>-1</sup>(y). (A.17)
Thus
$`\mathrm{ln}\left(\frac{\mathrm{\Phi }\text{(F[y,z])}}{\mathrm{\Phi }\text{(y)}}\right)`$ = $`\mathrm{ln}`$(K(z)) for some function K. (A.18)
Now (A.14) indicates that
$`_z\mathrm{ln}`$($`\mathrm{\Phi }`$(F\[y,z\])) = H<sup>-1</sup>(F\[y,z\])$`_2`$F\[y,z\] = $`\beta `$H<sup>-1</sup>(z) (A.19)
so that
$`_z\mathrm{ln}`$($`\mathrm{\Phi }`$(F\[y,z\])) = $`\beta _z\mathrm{ln}`$($`\mathrm{\Phi }`$(z)) = $`_z\mathrm{ln}`$($`\mathrm{\Phi }^\beta `$(z)). (A.20)
Thus
$`_z\mathrm{ln}\left(\frac{\mathrm{\Phi }\text{(F[y,z])}}{\mathrm{\Phi }^\beta \text{(z)}}\right)`$ = 0. (A.21)
Now considering (A.18) also this leads to
$`_z\mathrm{ln}`$(K(z)) = $`_z\mathrm{ln}`$($`\mathrm{\Phi }`$(F\[y,z\])) = $`\beta _z\mathrm{ln}`$($`\mathrm{\Phi }`$(z)). (A.22)
Thus, for a constant c, and by (A.18), it follows that
$`\mathrm{ln}`$(K(z)) = $`\mathrm{ln}`$($`\mathrm{\Phi }^\beta `$(z)) - $`\mathrm{ln}`$(c) = $`\mathrm{ln}`$($`\mathrm{\Phi }`$(F\[y,z\])) - $`\mathrm{ln}`$($`\mathrm{\Phi }`$(y)). (A.23)
This last equation is equivalent to
c$`\mathrm{\Phi }`$(F\[y,z\]) = $`\mathrm{\Phi }`$(y)$`\mathrm{\Phi }^\beta `$(z). (A.24)
Note that c must be non-zero. Now it remains to determine $`\beta `$. Starting from (A.1) it is evident that
c$`\mathrm{\Phi }`$(F\[x,F\[y,z\]\]) = c$`\mathrm{\Phi }`$(F\[F\[x,y\],z\]). (A.25)
Applying (A.24) to (A.25) twice results in
$`\mathrm{\Phi }`$(x)$`\text{[ c}{}_{}{}^{1}\mathrm{\Phi }\text{(y)}\mathrm{\Phi }^\beta \text{(z)]}^\beta `$ = c$`{}_{}{}^{1}\mathrm{\Phi }`$(x)$`\mathrm{\Phi }^\beta `$(y)$`\mathrm{\Phi }^\beta `$(z). (A.26)
and this is equivalent to
c<sup>β-1</sup> = $`\mathrm{\Phi }^{\beta (\beta 1)}`$(z). (A.27)
Thus it must be that $`\beta `$ = 1 if $`\mathrm{\Phi }`$ is not to be constant. The solution of (A.1) is therefor
$`\mathrm{\Phi }`$(F\[x,y\]) = c$`{}_{}{}^{1}\mathrm{\Phi }`$(x)$`\mathrm{\Phi }`$(y). (A.28)
### A.2 The Categoricity Equation
The equation to be solved is
$`\varphi `$(a) = f($`\varphi `$(a)). (A.29)
This correspondence must hold in all cases, including when attention is restricted to a Boolean Lattice. This is the case under which (A.29) will be solved.
Let h = $`\mathrm{\Phi }\varphi `$. Then, by (A.28), it follows that
c$``$h(a$``$b) = h(a)$``$h(b). (A.30)
Let h(a) = g(h(a)). From (A.29) it follows that
$`\varphi `$(a) = \[$`\mathrm{\Phi }^1`$g$`\mathrm{\Phi }`$\]$`\varphi `$(a). (A.31)
Thus f = $`\mathrm{\Phi }^1`$g$`\mathrm{\Phi }`$ so that finding f is reduced to finding g.
Now, according to (A.29), (A.30), and DeMorgan’s Law, it follows that
c$``$g(h(a$``$b)) = g(h(a))$``$g(h(b)) (A.32)
so that
h(b) = g(g(h(b))) = g$`\left(\frac{\text{c}\text{g(h(a}\text{b))}}{\text{g(h(a))}}\right)`$. (A.33)
Similarly, with the help of (A.33), it follows that
c$``$h(a$`{}_{}{}^{}`$b) = g(h(a))$``$g$`\left(\frac{\text{c}\text{g(h(a}\text{b))}}{\text{g(h(a))}}\right)`$. (A.34)
Now (A.30) may be used to evaluate this same quantity in another way;
c$``$h(a$`{}_{}{}^{}`$b) = h(b)$``$g$`\left(\frac{\text{c}\text{h(a}\text{b)}}{\text{g(h(a))}}\right)`$. (A.35)
Then (A.34) and (A.35) combine to yield
g(h(a))$``$g$`\left(\frac{\text{c}\text{g(h(a}\text{b))}}{\text{g(h(a))}}\right)`$ = h(b)$``$g$`\left(\frac{\text{c}\text{h(a}\text{b)}}{\text{g(h(a))}}\right)`$. (A.36)
Consider a special case where a = i$``$j and b = i$``$j. Then a$``$b = a and a$``$b = b. Applying this to (A.36) yields
g(h(a))$``$g$`\left(\frac{\text{c}\text{g(h(b))}}{\text{g(h(a))}}\right)`$ = h(b)$``$g$`\left(\frac{\text{c}\text{h(a)}}{\text{h(b)}}\right)`$. (A.37)
Let y = g(h(a)), z = h(b), u = $`\frac{\text{c}\text{g(y)}}{\text{z}}`$, and v = $`\frac{\text{c}\text{g(z)}}{\text{y}}`$. Then (A.37) becomes simply
z$``$g(u) = y$``$g(v). (A.38)
Operating with $`_y`$ on (A.38) results in
c$``$g(u)$``$g(y) = g(v) - v$``$g(v). (A.39)
Similarly, taking $`_z`$ of (A.38) leads to
c$``$g(v)$``$g(z) = g(u) - u$``$g(u). (A.40)
Operating on (A.40) with $`_y`$ now gives
u$``$y$``$g<sup>′′</sup>(u)$``$g(y) = v$``$z$``$g<sup>′′</sup>(v)$``$g(z). (A.41)
By using (A.38) this may be reduced to
u$``$g<sup>′′</sup>(u)$``$g(u)$``$g(y) = v$``$g<sup>′′</sup>(v)$``$g(v)$``$g(z). (A.42)
Equations (A.39) and (A.40) may be solved for g(y) and g(z) and the results substituted into (A.42). This results in
$$\frac{\text{u}\text{g}\text{′′}\text{(u)}\text{g(u)}}{\text{g}\text{}\text{(u)}\text{(u}\text{g}\text{}\text{(u) - g(u))}}\text{ = }\frac{\text{v}\text{g}\text{′′}\text{(v)}\text{g(v)}}{\text{g}\text{}\text{(v)}\text{(v}\text{g}\text{}\text{(v) - g(v))}}\text{ = }\lambda \text{ ( a constant ).}$$
(A.43)
Evidently, then, the variables u and v may be separated in this differential equation. The first integral of the resulting equation is given by
g(x)$``$g(x) = $`\alpha `$x (A.44)
where $`\alpha `$ is a constant of integration.
The solution of (A.44) falls into two cases corresponding to $`\lambda `$ = 1 and $`\lambda `$ $``$ 1.
For $`\lambda `$ = 1 the solution of (A.44) is of the form g(x) = A$``$x<sup>α</sup>. Then
x = g(g(x)) = g(A$``$x<sup>α</sup>) = A$`{}_{}{}^{\alpha \text{+}1}`$x$`^{\alpha ^2}`$ (A.45)
This amounts to
A$`^{\text{-(}\alpha \text{+}1\text{)}}`$ = x$`^{\text{(}\alpha \text{+}1\text{)(}\alpha \text{-}1\text{)}}`$. (A.46)
This implies that $`\alpha `$ = $`\pm `$1. For $`\alpha `$ = 1, A = $`\pm `$1. For $`\alpha `$ = - 1, A need merely be non-zero. Thus the $`\lambda `$ = 1 case solutions are of the form
g(x) = $`\pm `$x, or g(x) = A$``$x<sup>-1</sup>. (A.47)
According to the previous section of the appendix, it was found that $``$ formally corresponds to $``$. It is necessary, then, that $``$-products may be factored so that an interpretation may be uniquely determined. It may therefor be seen that all of the solutions for g(x), for the $`\lambda `$ = 1 case, must be rejected. In particular, g(x) = x obviously doesn’t distinguish x from g(x) in the first place. g(x) = -x must be rejected because products of negatives don’t retain specific signs for the factors, so that it would be generally impossible to determine which factors represent g(x) for some x. Finally, g(x) = A$``$x<sup>-1</sup> must also be rejected. In this case it is because the operation of inversion with respect to $``$ has already been reserved, by design, for the task of indicating erasures. It would be impossible to necessarily determine, then, whether a factor x<sup>-1</sup> would indicate that x is to be erased at some point or that it indicates that x is not a meaningful symbol. There therefor only remains one possible class of solution.
For $`\lambda `$ 1 the solution of (A.44) is defined by
g$`^\text{r}`$(x) = $`\alpha `$x$`^\text{r}`$ \+ B, where r = (1-$`\lambda `$). (A.48)
Applying (A.48) twice to (A.38) results in
($`\alpha ^2`$c$`^\text{r}`$ \- B)$``$(y$`^\text{r}`$ \- z$`^\text{r}`$) = 0, so B = $`\alpha ^2`$c$`^\text{r}`$. (A.49)
In light of (A.49), it follows that substituting (A.48) into x = g(g(x)) results in
(1 - $`\alpha ^2`$)$``$x$`^\text{r}`$ = $`\alpha ^2`$($`\alpha `$ \+ 1)$``$c$`^\text{r}`$, so $`\alpha `$ = $`\pm `$1. (A.50)
The only solution for the case $`\lambda `$1 which, at least, has not already been found in the previous case is then
g$`^\text{r}`$(x) = c$`^\text{r}`$ \- x$`^\text{r}`$. (A.51)
This general solution of this equation requires that complex numbers be considered, and thus such numbers must be admitted to the general symbolic system. Such a number system may be used to distinguish the factor x from g(x). Thus the solvability of (A.29) necessitates the inclusion of complex numbers in the symbolic system, and, conversely, gives a definite rationale for using such symbols.
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# The Quest for the Neutrino Mass Spectrum
## 1 Introduction
Recently the particle physics community was shocked with breathtaking news from the neutrino sector: Neutrino oscillations have been confirmed finally in the Super-Kamiokande experiment. Now for the first time, ongoing and future experiments in neutrino oscillations (Super-Kamiokande, Borexino, SNO, MINOS, KAMLAND, MINIBOONE,…) and double beta decay (Heidelberg–Moscow, GENIUS,…) together can aim to solve the neutrino mass puzzle.
It was in 1930, when Wolfgang Pauli (fig. 2) wrote his famous letter adressed as Liebe radioaktive Damen und Herren (Dear radioactive Ladies and Gentlemen), where he informed the participants of a nuclear physics workshop in Tübingen about his absence (he preferred to participate in a dance party) and postulated the neutrino to solve the problem of energy nonconservation in the nuclear beta decay. In 1956 the neutrino was observed for the first time by Clyde Cowan and Fred Reines in Los Alamos, who originally planned to explode a nuclear bomb for their experiment . Finally, two years ago, the Super-Kamiokande experiment, a 50,000 ton water tank viewed by more than 11,000 photo multipliers 1,000 meter underground below a holy mountain in Japan, announced a significant signal for neutrino oscillations and established a non-vanishing mass of the neutrino as the first experimental signal of physics beyond the standard model. However, in spite of these successes, entering a new millenium the neutrino is still the most mysterious of the known particles. Alternatingly compared with spaceships travelling through the universe, ghosts penetrating solid rocks and vampires missing a mirror image , it still inspires the phantasy of hundreds of adventurous particle, nuclear and astro physicists being motivated by the hope, the neutrino could act as a key to the old human dream of a final theory, describing all particles and forces in a unified framework, and to a deeper understanding of the fate of the universe.
The attributes, making the neutrino this kind of outlaw among the known particles, are the following:
* The neutrino seems to possess an at least million times smaller mass than the lightest of the remaining particles, the electron. While in the standard model the neutrino was introduced as massless “by hand”, this feature is especially problematic in unified theories, where the common treatment of neutrinos and charged fermions in extended multiplets implies them to have (Dirac) mass terms of the same order of magnitude as the other fermions.
* Among all fundamental fermions the neutrino is the only one being electrically uncharged. Thus the neutrino interacts a billion times less often than an electron and may penetrate the entire earth without even be deviated. This is the reason why neutrinos, in spite of their tiny masses, may be that abundant that they contribute substantially to the mass of the universe, about twenty times more than the mass of all visible stars in the sky, and may influence the evolution of the universe, e.g. the growth of structures, in a significant way.
The basis for an understanding of these features relates them to each other and was proposed in 1933 by the Italian theoretician Ettore Majorana , one year before he dissappeared under mysterious circumstances. Majorana found out that neutrinos, due to their neutral charge, can be identical with their antiparticles, triggered by a new, so-called Majorana mass term <sup>*</sup><sup>*</sup>*In fact also pure usual “Dirac” mass terms for the neutrino are possible but are disfavored in most fundamental theories.. In 1979 T. Yanagida and independently Murray Gell-Mann (nobel prize 1969), P. Ramond and R. Slansky found out that these additional Majorana mass terms may cancel almost totally the usual Dirac mass terms in the so-called “see-saw mechanism” , yielding a natural explanation of the tiny neutrino masses. This would require the existence of right-handed heavy neutrinos as they are naturally predicted in “left-right-symmetric” unified models. Alternative mechanisms motivate neutrino masses at the weak scale, a famous example is R-parity violating supersymmetry, see e.g. , where neutrino masses provide a window into deep relations of particles and forces. Also gravity induced non-renormalizable mass terms can play a role in string-motivated scenarios, see e.g. . The exact value of the mass then is correlated with a higher energy scale predicted by the underlying unified gauge group, and offers one of the rare possibilities to test these theories, since most of the predictions are observable only at very high energies, which are lying beyond the reach of present and future accelerators (see fig. 4). The question to experimentalists thus remains: What is the mass of the neutrino? The following review will outline the way to answer this question, concentrating on two experimental approaches, yielding the complementary pieces to solve the puzzle: Only both neutrino oscillations and neutrinoless double beta decay together could solve this absolute neutrino mass problem.
## 2 Neutrino oscillations
The fact that neutrinos are massive has finally been established by neutrino oscillation experiments. Neutrino oscillations are a quantum mechanical process based on mixing between the three neutrino flavors, which is possible if the flavor (interaction) eigenstates $`\nu _\alpha `$ do not coincide with the mass eigenstates $`\nu _i`$. The flavor eigenstates are thus given by a superposition of the mass eigenstates:
$$\nu _\alpha =\underset{i=1}{\overset{3}{}}U_{\alpha i}\nu _i$$
(1)
In that case a neutrino, which is emitted as a flavor eigenstate $`\nu _\alpha `$ in a weak reaction, propagates as a superposition of the three mass eigenstates. If these mass eigenstates are non-degenerate, they travel with different velocities and the composition in eq. (1) is getting out of phase. With a probability, which is a function of the mass squared differences $`\mathrm{\Delta }m^2=m_i^2m_j^2`$ and the mixing $`U_{\alpha i}`$, after a certain distance the neutrino interacts as a different mass eigenstate $`\nu _{\beta \alpha }`$ (see fig. 6). Obviously neutrino oscillation experiments cannot give any information about the absolute mass scale in the neutrino sector, but yield informations about mass (squared) differences, only. Since the probability oscillates with the propagation distance, this phenomenon, which was predicted by Bruno Pontecorvo, after he disappeared in 1950 from England and later showed up again in Russia, is called neutrino oscillations .
Up to now, hints for neutrino oscillations have been observed in solar and atmospheric neutrinos as well as the accelerator experiment LSND (for an overview see fig. 8). It should be stressed that besides neutrino oscillations also new interactions beyond the standard model may provide solutions to some of the neutrino anomalies, see
* A deficit of the number of solar neutrinos being expected has been confirmed in many experiments after the pioneering Chlor experiment of Ray Davis in the Homestake mine. The oscillation mechanism of the solar $`\nu _e`$ in (as normally assumed) $`\nu _\mu `$ <sup>§</sup><sup>§</sup>§An alternative would be a fourth sterile $`\nu _s`$, see section 7 may be induced via two different mechanisms. The usual neutrino oscillation mechanism requires maximal mixing and suffers from the fact, that for this case the distance earth-sun has to be finetuned (vacuum oscillations). An alternative solution has been suggested by works of S. Mikheyev, Alexei Smirnov and L. Wolfenstein : Resonant conversions, which are triggered by matter effects in the solar interior implying a level crossing of mass eigenstates, can cause the neutrino deficit. In this case both small as well as large mixing are allowed. The different solutions of the solar neutrino experiments correspond to different combinations of mass squared differences $`\mathrm{\Delta }m_{12}^2`$ and mixing matrix elements $`U_{12}^2`$. They will be tested by ongoing and future experiments such as Super-Kamiokande, SNO and BOREXINO in the next years If one allows for larger confidence belts a third MSW “LOW” solution appears, which can be tested via its strong day-night effect at low neutrino energies, observable at BOREXINO, LENS and the double beta and dark matter detector GENIUS (see below) .. Vacuum oscillations should lead to seasonal variations, the small mixing MSW solutions should imply distortions of the energy spectrum and the large mixing angle solution should show a small spectral distortion, a day-night effect of the total rate and a disappearance signal in the long baseline reactor experiment KAMLAND just under construction.
* A similar effect has been observed in atmospheric neutrinos , which stem from the decay of the pions produced from cosmic ray interactions in the upper atmosphere and the following-up decays. Here Super-Kamiokande obtained a high precision result of a deficit of muon neutrinos compared to electron neutrinos. Even more convincing is the distortion observed for the zenith angle dependence of the muon neutrino flux, which provides a strong hint for $`\nu _\mu \nu _\tau `$ oscillations with maximal mixing and information about $`\mathrm{\Delta }m_{23}^2`$ and $`U_{23}^2`$ . Future long baseline experiments, K2K (already running), MINOS, and CERN-Gran Sasso , looking for oscillations in accelerator produced neutrino beams over distances of several hundred kilometers will provide a check of this result by directly looking for $`\nu _\tau `$ appearance and have the possibility to search for small contributions of $`\nu _e\nu _\tau `$ oscillations.
* Also an accelerator experiment, LSND, has reported evidence for $`\nu _e\nu _\mu `$ neutrino oscillations. However, this evidence is generally understood as the most ambiguous. The KARMEN experiment has excluded a large part of the favored region of LSND. Since only two experimental evidences may be fitted with only three neutrinos. the LSND result would require the existence of a fourth, sterile (i.e. not weakly interacting) neutrino (see section 7). A decisive test will be obtained from the MINIBOONE experiment .
## 3 Neutrinoless double beta decay
Double beta decay ($`0\nu \beta \beta `$) corresponds to two single beta decays occurring in one nucleus and converts a nucleus (Z,A) into a nucleus (Z+2,A) (see fig. 6). While the standard model (SM) allowed process emitting two antineutrinos
$$^A_ZX_{Z+2}^AX+2e^{}+2\overline{\nu }_e$$
(2)
is the rarest process observed in nature with half lives in the region of $`10^{2124}`$ years, more interesting is the search for the lepton number violating and thus SM forbidden neutrinoless mode,
$$^A_ZX_{Z+2}^AX+2e^{}$$
(3)
which has been proposed by W.H. Furry in 1939 . In this case the neutrino is exchanged between the vertices (see fig. 6), a process being only allowed if the intermediate neutrino has a Majorana mass. Neutrinoless double beta decay, when observed, also does not measure directly the neutrino mass. Since the neutrino in the propagator is only virtual, it does not have a definite mass. Propagating in the nucleus is the flavor eigenstate with the so-called effective neutrino Majorana mass
$$m=\left|\underset{j}{}|U_{ej}|^2e^{i\varphi _j}m_j\right|,$$
(4)
which is a function of the mixing angles $`U_{ej}`$, complex phases $`\varphi _j`$, which allow for cancellations of the entering masses, and the neutrino mass eigenvalues. This quantity has exciting connections to the observables in neutrino oscillation experiments. The most stringent limit on this quantity, $`m<0.35`$ eV, is obtained by the Heidelberg–Moscow experiment , which was initiated by one of the authors and is running since 10 years in the Gran Sasso underground laboratory in Italy. An impressive breakthrough to $`10^210^3`$ eV could be obtained realizing the GENIUS project proposed in 1997 , a further proposal of H.V. Klapdor-Kleingrothaus, operating 1-10 tons of enriched Germanium directly in a tank of 12 m diameter and height filled with liquid nitrogen.
How are the results in double beta decay and neutrino oscillations related? In a recent work the authors of this article in collaboration with Alexei Smirnov from the ICTP Trieste were studying the relations of the neutrino oscillation parameters and the effective Majorana mass in the several possible neutrino mass scenarios and settled the conditions under which the neutrino mass spectrum can be reconstructed with future projects (see fig. 10). In the following we will concentrate on three extreme cases as examples, the hierarchical spectrum, the degenerate scheme and the inverse hierarchical scheme.
## 4 Hierarchical schemes
Hierarchical spectra (fig. 12)
$$m_1m_2m_3$$
(5)
can be motivated by analogies with the quark sector and the simplest see-saw models. In these models the contribution of $`m_1`$ to the double beta decay observable $`m`$ is small. The main contribution is obtained from $`m_2`$ or $`m_3`$, depending on the solution of the solar neutrino deficit.
If the small mixing angle solution is realized in solar neutrinos (i.e. small $`\nu _e\nu _\mu `$ mixing), the contribution of $`m_2`$ is small due to the small admixture $`U_{e2}`$. The same is true for vacuum oscillations, where $`U_{e2}`$ is maximal but the mass of the second state is tiny. In these cases the main contribution to $`m`$ comes from $`m_3`$. The contribution of the latter is shown in fig. 14. Here lines of constant $`m`$ are shown as functions of the oscillation parameters $`\mathrm{\Delta }m_{13}^2`$ and $`U_{13}`$, parametrized by $`\mathrm{sin}^22\theta _{13}`$. The shaded areas show the mass $`m_3\sqrt{\mathrm{\Delta }m_{13}^2}`$ favored by atmospheric neutrinos with the horizontal line indicating the best fit value. The region to the upper right is excluded by the nuclear reactor experiment CHOOZ , implying $`m<210^3`$ eV in the range favored by atmospheric neutrinos. Obviously in this case only the 10 ton GENIUS experiment could observe a positive $`0\nu \beta \beta `$ decay signal. A coincidence of such a measurement with a signal of $`\nu _e\nu _\tau `$ oscillations at MINOS and a confirmation of solar vacuum or small mixing MSW oscillations by solar neutrino experiments would be a strong hint for this scheme.
If the large mixing solution of the solar neutrino deficit is realized, the contribution of $`m_2`$ becomes large due to the almost maximal $`U_{e2}`$, now. Fig. 16 shows values of $`m`$ in the range of the large mixing angle solution (closed line). The almost horizontal lines correspond to constant day-night asymmetries. A coincident measurement of $`m10^3`$ eV, a day-night asymmetry of 0.07 at future oscillation experiments and a confirmation of the large mixing angle solution by KAMLAND would identify a single point in the large mixing angle MSW solution (in this example near the present best-fit point) and provide a strong hint for this scheme.
## 5 Degenerate schemes
Degenerate schemes (fig. 18)
$$m_1m_2m_3m_0$$
(6)
require a more general (and more complicated) form of the see-saw mechanism . One of their motivations is also, that a large overall mass scale allows neutrinos to be cosmologically significant. Neutrinos with an overall mass scale of a few eV could play an important role as “hot dark matter” component of the universe. When structures were formed in the early universe, overdense regions of (cold) dark matter provide the seeds of the large scale structure, which later formed galaxies and clusters. A small “hot” (relativistic) component could prevent an overproduction of structure at small scales. Since structures redshift photons, this should imply also imprints on the cosmic microwave background (CMB), which could be measured by the future satellite experiments MAP and Planck .
In degenerate schemes the mass differences are not significant. Since the contribution of $`m_3`$ is strongly bounded by CHOOZ again, the main contributions to $`m`$ come from $`m_1`$ and $`m_2`$. The relative contributions of these states depend on their admixture of the electron flavor, which is determined by the solution of the solar neutrino deficit.
In fig. 20 lines of constant double beta decay observables (solid curved lines) are shown together with information from cosmological observations about the overall mass scale (horizontal lines). Shown are best fits to the CMB and the large scale structure of Galaxy surveys in different cosmological models as well as the sensitivity of MAP and Planck. E.g., a $`\mathrm{\Lambda }`$CHDM model with a total $`\mathrm{\Omega }_m=0.5`$ of both cold and hot dark matter as well as a cosmological constant, and a Hubble constant of $`h=0.6`$ would imply an overall mass scale of about 0.5 eV. However, the contributions of different mass eigenstates are in the same order of magnitude and may cancel, now. Assuming a mixing corresponding to the best fit of solar large mixing MSW or vacuum oscillations this yields $`m=0.20.5`$ eV, just in the range of the recent half life limit of the Heidelberg–Moscow experiment. If even larger mixing turns out to be realized in the solution of the solar neutrino deficit, this allows for a larger cancellation. A coincidence of the absolute mass scale reconstructed from double beta decay and neutrino oscillations with a direct measurement of the neutrino mass in tritium beta decay spectra or its derivation from cosmological parameters determined from the CMB in the satellite experiments MAP and Planck would prove this scheme to be realized in nature. To establish this triple evidence however is difficult due to the restricted sensitivity of the latter approaches. Future tritium experiments aim at a sensitivity down to $`𝒪`$(0.1 eV) and MAP and Planck have been estimated to be sensitive to $`m_\nu =0.50.25`$ eV. Thus for neutrino mass scales below $`m_0<0.1`$ eV only a range for the absolute mass scale can be fixed by solar neutrino experiments and double beta decay.
## 6 Inverse Hierarchy
A further possibility is an inverse hierarchical spectrum (fig. 22)
$$m_3m_2m_1$$
(7)
where the heaviest state with mass $`m_3`$ is mainly the electron neutrino, now.
Its mass is mainly determined by the atmospheric neutrinos, $`m_3\mathrm{\Delta }m_{23}`$. Thus for the case of the small mixing angle solution one gets a unique prediction of $`m=(58)10^2`$ eV, which could be tested by the 1 ton version of GENIUS. For the vacuum or large mixing MSW solution cancellations of the two heavy states become possible and $`m<810^2`$ eV. A test of the inverse hierarchy is possible in matter effects of neutrino oscillations. For this case the MSW level crossing happens for antiparticles rather than for particles. Effects could be observable in long baseline experiments and in the neutrino spectra of supernovae .
## 7 Four neutrinos
Sterile neutrinos, which do not couple to the weak interactions, can easily be motivated in superstring inspired models: multidimensional candidates for a final “Theory of Everything”, in which the fundamental constituents of matter have a string rather than a particle character. Such theories could accomodate for additional neutrinos in different ways. Examples are extended gauge groups, fermions living in extra (compactified) dimensions as well as a mirror world, which contains a complete duplicate of matter and forces building the universe, interacting only via gravity. In the latter case $`m=0.002`$ eV is predicted .
If the four neutrinos are arranged as two pairs of degenerate states (mainly $`\nu _e\nu _s`$ for solar and $`\nu _\mu \nu _\tau `$ for atmospheric neutrinos) separated by a LSND gap, all three neutrino anomalies can be solved and the two heavy states can account for the hot dark matter. The main contribution to $`m`$ comes from the heavy states, then, and can be derived from the LSND result. Depending on the phase of these two contributions $`m`$ can be as large as $`𝒪(10^3`$ eV). A strong hint for the scheme would be a coincidence of the $`\mathrm{\Delta }m^2`$ favored in LSND and possibly MINIBOONE, cosmological observations and double beta decay, together with the discovery of sterile neutrinos in solar neutrino oscillations by SNO.
## 8 Summary
The recent years brought exciting developments in neutrino physics. Neutrino oscillations have finally been confirmed in atmospheric neutrinos and at the same time double beta decay experiments realized for the first time a sensitivity, leading to strong implications on the neutrino mass spectrum and cosmological parameters. After this particle physics now seems to enter its “neutrino epoche”: The neutrino mass spectrum and its absolute mass scale offer unique possibilities to provide crucial information for cosmology and theories beyond the standard model. Only both neutrino oscillations and neutrinoless double beta decay together have the chance to solve this neutrino mass problem (see also, e.g. ) and to set the absolute scale in the neutrino sector: If the solution of the solar neutrino deficit and the character of hierarchy (direct or inverse) is determined in neutrino oscillation experiments, the following informations will be obtained from a future double beta decay project:
For the case of direct/normal hierarchy, a confirmation of the small mixing MSW solution would mean: If double beta decay would be measured with $`m>0.1`$ eV this would establish a degenerate spectrum with a fixed mass scale. If $`m`$ is measured in the range $`(0.53)10^2`$ eV a partially degenerate spectrum, $`m_1m_2m_3`$, with fixed mass scale is realized in nature. For $`m<210^3`$ eV a hierarchical spectrum exists in nature. For the large mixing MSW solution a value of $`m>310^2`$ eV implies a degenerate spectrum with a region for the mass scale determined by the solar mixing angle. For $`m<210^2`$ eV a partially degenerate or hierarchical spectrum is realized in nature and a region for the mass scale is set by the solar mixing angle. If $`m<210^3`$ eV is measured the spectrum is hierarchical. If vacuum oscillations are the correct solution for the solar neutrino deficit a value of $`m>310^2`$ eV implies degeneracy, $`m>210^3`$ eV partial degeneracy and $`m<210^3`$ eV hierarchy, but no information about the absolute mass scale is obtained.
For the case of inverse hierarchy the situation is more predictive. For the small mixing angle MSW solution $`m(58)10^2`$ eV is expected. For large mixing angle MSW or vacuum oscillations one awaits $`m<810^2`$, above this value the scheme approaches the degenerate case.
In four neutrino schemes $`m`$ can be as large as $`𝒪(10^3)`$ eV. A conincidence of a double beta decay signal with the $`\mathrm{\Delta }m^2`$ favored in LSND and possibly in MINIBOONE, an imprint of neutrinos as hot dark matter in the CMB as well as the discovery of sterile neutrinos in SNO would prove the scheme and fix the mass scale.
This outcome will be a large step both towards the understanding of the evolution of the universe and towards the dream of a unified theoretical description of nature. We are entering an exciting decade!
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# Rigged Hilbert Spaces associated with Misra-Prigogine-Courbage Theory of Irreversibility.
## I Introduction.
The aim of this paper is to demonstrate unsuspected mathematical implications of the Misra-Prigogine-Courbage theory of irreversibility, one of the two theories of ”intrinsic irreversibility” developed by the group of Brussels (Belgium).
More precisely, it will be proved that MPC-theory is strongly connected with the theory of Rigged Hilbert Spaces (RHS). This is important for two reasons. First, because it increases the mathematical meaning of the $`\mathrm{\Lambda }`$ transformation, and relates it with a well known and successful subject (specially in Quantum Physics ) as is the Theory of Riggings. And second, for it makes possible a relation and a comparison with the other version of irreversibility, namely the ”Rigged Hilbert Space Extension through the Spectral Decomposition” .
The Misra-Prigogine-Courbage theory is based on the Internal Time superoperator and its associated the $`\mathrm{\Lambda }`$ transformations . Let us briefly explain this formalism considering a hamiltonian system: the motion of the dice when it is thrown. This system is theoretically deterministic and reversible, but in reality it is impossible to predict the result of one bet, by solving the equations of motion. It is so because these equations are dynamically unstable, i.e. any initial condition C is surrounded by many others C’ almost identical to C, but yielding completely different results. Therefore this kind of dynamics can be considered, ”for all practical purposes”, as a stochastic process, and solved using the theory of probabilities instead of newtonian mechanics.
The Brussels School has proved that any reversible but unstable enough dynamic, determines a class of $`\mathrm{\Lambda }`$ transformations. Each member of the class can be considered as an equivalence ”for all practical purposes” between the dynamic and a stochastic Markov process, irreversibly convergent towards an equilibrium density. Moreover, non-isomorphic dynamical processes are transformed by the $`\mathrm{\Lambda }`$ in non-isomorphic processes .
Therefore it is not necessary to use an arbitrary, observer-dependent, or extrinsic ”coarse graining” independent of the dynamics to transform a reversible evolution into an irreversible one going to equilibrium. The system itself, if it is sufficiently unstable, defines its own class of $`\mathrm{\Lambda }`$ superoperators, transforming its dynamic uncertainties (dues to unstability) into probabilistic estimations. (Actually, there are systems that also define a class of conditional expectations that constitute non arbitrary and dynamic-dependent, or intrinsic ”coarse grainings” projections, wich also yields an irreversible Markov process )
## II The $`\mathrm{\Lambda }(T)`$ formalism of the MPC-theory.
In this section we will introduce a notation that can be used both for classical and quantum systems
### A The classical case.
Let us consider an abstract dynamical system . Let $`\mathrm{\Omega }`$ be the states space (for example, the phase space), $``$ the $`\sigma `$-algebra of measurable sets of $`\mathrm{\Omega },`$ and $`\mu `$ the corresponding measure (e.g. Liouville measure). Let $`S_t:\mathrm{\Omega }\mathrm{\Omega }`$ be the time-evolution operator on phase space, with $`t𝔾`$, where $`𝔾`$ will be $``$ for the flows (i.e. continuous dynamical systems)and $``$ for the cascades (discrete dynamical systems). $`𝔏`$:$`=L^2(\mathrm{\Omega },,\mu )`$ will denote the Hilbert space of the equivalent classes ”a.e.” (almost everywhere) of measurable functions of $`\mathrm{\Omega }`$ in $``$ of square integrable modulus with respect to $`\mu .`$ Then, $`S_t`$ induce an unitary evolution $`U_t`$ over $``$, i.e. a unitary representation of the group $`(𝔾,+)`$ over $``$, defined as:
$$(U_t\rho )(\omega )=\rho (S_t(\omega ))\text{ , provided }\rho 𝔏\text{, and }\omega \mathrm{\Omega }$$
(1)
where $`U_t`$ is known as the Koopman operator.
$`𝒟`$ will denote the subspace of dimension one generated by the constant function 1:
$`𝒟:=\{\alpha 1/\alpha ;1:\mathrm{\Omega }1(\omega )=1,\text{if}\omega \mathrm{\Omega }\}`$
and we will write $`=𝒟^{}`$. Then:
$$𝔏=𝒟$$
(2)
The positive $`\rho 𝔏`$ (i.e.$`\rho (\omega )0`$ for every $`\omega \mathrm{\Omega }`$), which are also normal ( in the sense of the $`L^1`$ norm, i.e. $`_\mathrm{\Omega }\rho (\omega )𝑑\omega =1`$), will be the ”probability density functions” or the Gibbs ”ensembles” of the system. If $`\mu `$ is normalized, in such a way that $`_\mathrm{\Omega }𝑑\mu =\mu (\mathrm{\Omega })=1,`$ then the constant function equal to one is an invariant density under $`U_t`$ as a consequence of eq. (2.1). Furthermore it can be demonstrated that if the dynamical system is mixing $`U_t\rho 1`$ in a weak sense. Therefore 1 is called the equilibrium density and it is symbolized as $`1=\rho _{eq}`$. Also $`U_t|_{}0`$ in a weak sense .
If $`𝔾=R,`$ and being $`U_t`$ unitary in $`𝔏`$, there is a self-adjoint generator $`L`$ such that:
$$U_t=e^{iLt}$$
(3)
If $`S_t`$ is also a hamiltonian flux, with a hamiltonian function $`H`$ , then $`L`$ is call the Liouvillian, and (2.1) is equivalent to the Liouville differential equation:
$$L\rho =i_t\rho $$
(4)
where $`L=i\{H,.\}`$ and $`\{`$,$`\}`$ is the Poisson bracket.
### B The quantum case.
Non-trivial quantum systems have a continuous spectrum. In this case the equilibrium state is not an ordinary state but a ”singular diagonal” state . These facts force us to use an extension of the usual quantum mechanics formalism. Following the line of thought of the cited papers, and taking into account eq. (2.2) we postulate that our state space is:
$`𝔏=𝒟=\{\rho =\rho ^d+\rho ^c:\rho ^d𝒟\text{ y }\rho ^c\}`$
Then the $`\rho ^{}s`$ will evolve under a generalized Liouville equation . Space $`𝒟`$ contains the information about the probability density of the states. Space $``$ contains the information about correlations, coherent state superposition, and covariance between observables . Then, if the dynamics is mixing we will have $`\rho _t\rho _{eq},`$ and $`\rho _t^c0,`$ in a weak sense .
Let us now consider a quantum system defined in a Hilbert states space $``$ and its ”complete set of commuting observables”. Let $`𝒟_𝔄`$ be the maximal abelian von Neumann algebra that contains this set. We will call $`I`$ to its unit element. If some observables are essentially selfadjoint unbounded operators, we shall consider the algebra generated by their spectral projections, which are bounded. Let $`𝒜==\{`$Hilbert-Schmidt operators over $`\}`$, with respect to the scalar product $`\rho |\sigma _{}=Tr(\rho ^{}\sigma ).`$ As it is well known, this space is a Hilbert space. Let us now consider the algebra of observables of the system: $`𝔄=𝒟_𝔄𝒜`$, where $`𝒟_𝔄`$ is the diagonal part of the algebra, and $`𝒜`$ the non diagonal part of it. Let us define $`𝔏`$ as the dual space of $`𝔄`$, namely: $`𝔏=(A)^{}=𝒟`$, where $`𝒟=𝒟_𝔄^{}`$ $`=\{\rho _d:\rho _d`$ linear and continuous functional over $`𝒟_𝔄\},`$ and where $``$ has been identified with $`^{}.`$
The $`\rho =\rho ^d+\rho ^c`$ are non-negative, in the sense that for any $`A𝔄`$ , we have: $`\rho (A^{}A)0,`$ where $`A^{}=(A^d)^{}+(A^c)^{}`$ and they are normal in the sense that $`\rho (I)=1`$ (as $`I𝒟_𝔄,`$ $`\rho (I)=\rho ^d(I)).`$ These $`\rho `$ will be considered as the possible states of the system. The $`\rho ^d`$ which are non-negative as linear functional over the von Neumann algebra $`𝒟_𝔄`$ and normal in the sense of $`\rho ^d(I)=1,`$ will be the diagonal states, as are, e.g., the equilibrium states.
Let us observe that if we would take $`𝒜`$ as the set of compact operators over $``$, as in ref. , its dual space would be the space of nuclear operators, which is a subset of the Hilbert-Schmidt operators , namely $`𝒜`$.
### C The $`\mathrm{\Lambda }`$ transformation.
Let us consider a continuous linear operator $`𝚲:LL`$ , such that:
i) $`𝚲`$ preserves probabilities, in the strong sense that $`𝚲|_𝒟=I_𝒟`$ which is the identity in $`𝒟`$. I.e., $`𝚲=I_𝒟\mathrm{\Lambda },`$ where $`\mathrm{\Lambda }:`$ is linear and continuous in the Hilbert space $``$. In particular $`𝚲\rho _{eq}=\rho _{eq}.`$
ii) $`𝚲`$ transforms ensembles into ensembles, namely $`𝚲`$ preserves the positivity and the normalization. Therefore $`\mathrm{\Lambda }`$ must be non negative and symmetric. As the domain of $`\mathrm{\Lambda }`$ is the whole $``$, $`\mathrm{\Lambda }`$ must be self adjoint .
iii)$`𝚲`$ is not a ”coarse-graining”, namely it doesn’t neglect information as a ”coarse-graining”-projector. It is only a ”change of representation” that ”reorganizes”, or ”redefines” the information content of the densities, in such a way that the resulting theory is closer to actual experimental possibilities and to physical reality. This last requirement is attained by making $`𝚲`$ an injective and dense range application (states with ”infinite information content” are not in the range of $`𝚲)`$.
In fact, properties i) and ii) above, plus the injectivity, assure that the range of $`\mathrm{\Lambda }`$ must be either $``$ or dense in $``$ . If the range of $`\mathrm{\Lambda }`$ is $``$, then $`\mathrm{\Lambda }^{1\text{ }}`$is continuous and therefore it is an isomorphism and a homeomorphism, and then $`\mathrm{\Lambda }U_t\mathrm{\Lambda }^1`$ is a dynamical system equivalent to $`U_t.`$ On the contrary, if the range of $`\mathrm{\Lambda }`$ is dense in $``$ then $`\mathrm{\Lambda }^1`$ is unbounded and this singularity of $`\mathrm{\Lambda }^1`$ is essential because it gives new properties to $`\mathrm{\Lambda }`$ that can be considered as ”catastrophic” (i.e. with strong ”qualitative changes” ). Precisely the hamiltonian system $`U_{t\text{ }}`$ is transformed by the $`\mathrm{\Lambda }`$ into a stochastic process $`W_t=\mathrm{\Lambda }U_t\mathrm{\Lambda }^1.`$ As now $`\mathrm{\Lambda }^1`$ is unbounded its domain can be extended beyond the range of $`\mathrm{\Lambda }.`$
Nevertheless, there is not reason for the positivity of $`W_t`$ (and therefore for its markovian character), for any $`tG`$ beyond the range of $`\mathrm{\Lambda }.`$ So the unboundedness of $`\mathrm{\Lambda }^1`$ is the crucial ”detail” that makes that the $`W_t`$ do not form a group and breakes the time-symmetry .
iv) $`W_t=\mathrm{\Lambda }U_t\mathrm{\Lambda }^1,`$ $`t0`$ is the evolution operator of a strong Markov process, namely a monotonously convergent process to the null vector in the Hilbert topology of $``$ (and not only in a weak sense as in the mixing dynamics). I.e.: $`W_t\rho _{}0`$ if $`t\mathrm{},`$ for any $`\rho Dom(\mathrm{\Lambda }^1).`$ This property is similar to a Markov exact process , but in space $`L^2`$ instead of $`L^1.`$
Then, $`\mathrm{\Lambda }^2`$ is a decreasing Liapounov variable of the considered dynamics in the following sense:
$$W_t\rho _{}^2=\mathrm{\Lambda }U_t\mathrm{\Lambda }^1\rho \mathrm{\Lambda }U_t\mathrm{\Lambda }^1\rho _{}=\rho _t\mathrm{\Lambda }^2\rho _t_{}0$$
(5)
where $`\rho _t=U_t\mathrm{\Lambda }^1\rho .`$
Accordin to the Brussels group a dynamical system is intrinsically or essentially random if there exists a $`𝚲:𝔏L`$ with the properties above. In order that this happens it is necessary the mixing character of dynamic, and it is sufficient the existence of an age or internal time operator $`T`$ . For the flows $`T`$ is a kind of ”time-position operator”, similar to the ”space-position” operator $`Q`$ of quantum mechanics, but it acts in space $``$ instead of space $``$ .
The Liouville operator $`L`$ is the ”canonical conjugate momentum” of $`T`$ :
$$[T,L]=iU_t^{}TU_t=T+t,t$$
(6)
This last equation between the ”internal time” and the ”external time” $`t`$, which can also be used for cascades, can be considered as a general definition of $`T:`$
$$U_t^{}TU_t=T+t,t𝔾$$
(7)
The construction of operator $`T`$ is completely similar to that of operator $`Q.`$ $`T`$ exist iff the unitary representation $`U_t`$ of ($`𝔾`$,$`+)`$ in $``$ is imprimitive with respect to $`𝔾`$ . This means that there is a spectral measure $`E`$, defined over a $`\sigma `$-algebra $``$ of $`𝔾`$, and whose values are orthogonal projectors of $``$ , such that:
$$U_t^{}E(\mathrm{\Delta })U_t=E(\mathrm{\Delta }+t),\text{ for }t𝔾,\text{ and }\mathrm{\Delta }$$
(8)
In such a case:
$$T=\underset{𝔾}{}s\text{ }𝑑E$$
(9)
From eq. (2.8), for all $`\rho `$, the numerical measure $`\mathrm{\Delta }\rho |E(\mathrm{\Delta })\rho _{}`$, is translational invariant. Then, if $`𝔾=R,`$ it is equivalent to the Lebesgue measure . In other words, for flows, the spectrum of $`L`$ must be absolutely continuous and uniform . This condition is fulfilled for classical and quantum K-flows . Going back to the general case, any $``$-measurable function $`\lambda :𝔾[0,1],`$ such that:
i) $`\lambda `$ is decreasing, i.e.: $`r<s\lambda (r)\lambda (s).`$
ii)$`\lambda (t)1`$ if $`t\mathrm{}`$ and $`\lambda (t)0`$ if $`t\mathrm{}.`$
iii) If $`t0:\frac{\lambda (s+t)}{\lambda (s)}0,`$ i.e.: $`r<s\frac{\lambda (r+t)}{\lambda (r)}\frac{\lambda (s+t)}{\lambda (s)}`$ and $`\frac{\lambda (s+t)}{\lambda (s)}0`$ if $`s\mathrm{}.`$
defines a $`𝚲`$ as:
$$\mathrm{\Lambda }=\lambda (T)=\underset{𝔾}{}\lambda (s)\text{ }𝑑E$$
(10)
$$𝚲=I_𝒟\mathrm{\Lambda }$$
(11)
Since operator $`T`$ fulfills eq. (2.7), we have:
$`U_t^{}\lambda (T)U_t`$ $`=`$ $`\lambda (T+t),\text{for }t𝔾`$ (12)
$`U_t^{}\lambda ^2(T)U_t`$ $`=`$ $`\lambda ^2(T+t),\text{for }t𝔾`$ (13)
where $`\lambda ^2(T)=\mathrm{\Lambda }^2`$ is the decreasing Liapounov variable.
## III The rigged Hilbert spaces.
Let $``$ be a separable Hilbert space (e.g. a Liouville space, but here it will be considered in a general sense). Let $`\mathrm{\Psi }`$ be a proper vector subspace of $``$. Let us suppose that in $`\mathrm{\Psi }`$ it is defined a countable family of Hilbert norms $`\{||.||_n\}_{nN}`$ , where $`||.||_n=.|._n^{\frac{1}{2}},`$ $`N^+`$ (the set of rational non negative numbers, therefore the set $`N`$ is countable), such that:
(i) $`n_1n_2||.||_{n_1}||.||_{n_2}`$ and both norms are compatible, meaning that if $`\{\rho _n\}`$ is a Cauchy sequence in both norms, and if $`\rho _m_{n_1}0`$ then $`\rho _m_{n_2}0.`$
(ii) $`N`$ has a minimum element, that wil be assumed to be zero (for simplicity), $`||.||_0=||.||_{}`$, and $`\mathrm{\Psi }`$ is dense in $``$ (meaning that the completion of $`(\mathrm{\Psi },||.||_0)`$ is $``$)
In such a case, the completion of $`\mathrm{\Psi }`$ with $`||.||_{n_i}`$ will be denoted $`\mathrm{\Phi }_{n_i}`$, whose elements are the equivalence classes $`[\{\rho _m\}]_{n_i}`$ of the Cauchy sequences $`\{\rho _m\}`$ of $`\mathrm{\Psi },`$ where:
$`\{\rho _m\}\{\sigma _m\}\rho _m\sigma _m_{n_i}0\text{ if }m\mathrm{}`$
$`\mathrm{\Psi }\mathrm{\Phi }_{n_2}\mathrm{\Phi }_{n_1}\mathrm{\Phi }_0=`$
condition (i) assures the injectivity and continuity of the canonical applications $`i_{n_2,n_1}:\mathrm{\Phi }_{n_2}\mathrm{\Phi }_{n_1}`$ defined by $`[\{\rho _m\}]_{n_2}[\{\rho _m\}]_{n_1},`$ and condition (ii) that they have a dense range.
Let us now consider the local convex topology defined by the family of norm on $`\mathrm{\Psi }.`$ In other words the topology such that:
$`\{\rho _m\}\rho \text{for every }nN:\rho _m\rho _n0`$
There are three possibilities:
A) $`N`$ has a maximum ñ.
In this case, from condition (i) we have:
$`\{\rho _m\}\rho \rho _m\rho _{\stackrel{~}{n}}0`$
But this is the norm topology $`||.||_{\stackrel{~}{n}}.`$ Then, by completion of $`(\mathrm{\Psi },||.||_{\stackrel{~}{n}})`$ we get a Hilbert space $`\mathrm{\Phi }_H,`$ such that for every $`nN:\mathrm{\Phi }_H\mathrm{\Phi }_n.`$
B) $`N`$ has a supreme ñ, but ñ $`N`$.
In this case we do not get a Hilbert space, but a $`\sigma `$Hilbertian space , that we shall call $`\mathrm{\Phi }_K`$ (because, as we shall see later, this is a Köthe space, if $`\mathrm{\Lambda }`$ is compact). In particular, if $`N`$ has a maximum $`\stackrel{~}{n},`$ but we only consider the family of norms $`\{||.||_n\}_{nN\{\stackrel{~}{n}\}}`$, we get a space $`\mathrm{\Phi }_K=\{\mathrm{\Phi }_n:nN\{\stackrel{~}{n}\}\}`$ such that $`\mathrm{\Phi }_H\mathrm{\Phi }_K`$.
C) $`N`$ is not bounded from above.
In this case we get a smallest $`\sigma `$Hilbertian space with a stronger topology than $`\mathrm{\Phi }_K`$ (and therefore easier to transform in a nuclear topology, by endowing $`\mathrm{\Lambda }`$ with more properties$`).`$ We shall call this space $`\mathrm{\Phi }.`$ Precisely, there is a sequence of subsets of $`N`$, $`\{N_p\}`$, such that: $`N_1N_2\mathrm{}`$, and $`_pN_p=N.`$ In this way we can obtain a sequence of spaces $`\mathrm{\Phi }_{H_p}`$ as in paragraph (A). Now, condition (i) assures that the canonical mappings $`i_{H_p}:\mathrm{\Phi }_{H_p}`$ $``$, defined as the $`i_{n_2,n_1}`$ when $`n_2=1,`$ $`n_1=0,`$ and the mapping $`i:\mathrm{\Phi }`$, defined as: $`i(\stackrel{~}{\rho })=i_{H_p}(\stackrel{~}{\rho }),`$ for every $`\stackrel{~}{\rho }\mathrm{\Phi }`$ and every $`p`$, are all of them injectives and continuous, and condition (ii) assures that they have a dense range.
In any of these cases it is usual to say that we have rigged the Hilbert space $``$ with another Hilbert space $`\mathrm{\Phi }_H`$ or with a $`\sigma `$Hilbertian space, either $`\mathrm{\Phi }_K`$ or $`\mathrm{\Phi }.`$ Really we must also consider the corresponding antidual spaces (of continuous antilinear functions) that we shall call $`\mathrm{\Phi }_H^\times ,`$ $`\mathrm{\Phi }_K^\times ,`$ $`\mathrm{\Phi }^\times ,`$ $`^\times =`$. In fact, as the topologies of $`\mathrm{\Phi }_H,`$ $`\mathrm{\Phi }_K,`$ and $`\mathrm{\Phi }`$ are stronger than that of $``$, they make possible the existence of larger sets of continuous antilinear functionals. Therefore we have:
$`\mathrm{\Phi }_H\mathrm{\Phi }_H^\times \text{ };\text{ }\mathrm{\Phi }_K\mathrm{\Phi }_K^\times \text{ };\text{ }\mathrm{\Phi }\mathrm{\Phi }^\times `$
where the corresponding inclusions are continuous and their images are dense.
Let us consider a rigging of type (A). Let $`:\mathrm{\Phi }_H^\times \mathrm{\Phi }_H`$ be the Riesz representation: to every antilinear continuous functional $`F`$ it associates the vector $`\rho _F`$ such that:
$$\sigma \rho _F_{\mathrm{\Phi }_H}=F(\sigma )\text{ , for every }\sigma \mathrm{\Phi }_H$$
(14)
It is known that $``$ is an isometric isomorphism. Nevertheless, $`\mathrm{\Phi }_H\mathrm{\Phi }_H^\times `$ if we consider these spaces just like sets. Then:
$`\mathrm{\Phi }_H\mathrm{\Phi }_H^\times `$
It is easy to see that $`|_{}=R`$ is a non negative operator. In fact, from Riesz representation of $`^\times `$ in $``$, all $`\rho `$ can be considered as an antilinear continuous functional on $``$,
$$\sigma \sigma \rho _{}$$
(15)
On the other hand, as $`\mathrm{\Phi }_H^\times `$ the same $`\rho `$ can be thought as a functional:
$$\sigma \sigma R(\rho )_{\mathrm{\Phi }_H}$$
(16)
Then:
$$\sigma \rho _{}=\sigma R(\rho )_{\mathrm{\Phi }_H}$$
(17)
If, in particular, $`\sigma =R(\rho ),`$ then:
$$R(\rho )\rho _{}=R(\rho )R(\rho )_{\mathrm{\Phi }_H}0$$
(18)
Therefore $`R`$ is non negative and thus it has a square root $`J=R^{\frac{1}{2}},`$ which also is non negative, continuous and self-adjoint in $`,`$ injective and with dense range. Furthermore it is proven in ref. that:
$$J\sigma J\rho _{\mathrm{\Phi }_H}=\sigma \rho _{}$$
(19)
in such a way that $`J`$ turns out to be an isometry.
Viceversa., if we have an operator $`J:`$ with the same properties as above, then the relation:
$$J\sigma J\rho _\mathrm{\Psi }=\sigma \rho _{}$$
(20)
defines a scalar product on $`\mathrm{\Psi }=Ran(J),`$ whose completion is a Hilbert space $`\mathrm{\Phi }_H`$ which riggs $``$ in a canonical way. Let us remark that giving a rigging type (A) is equivalent to giving an operator $`J`$ with the properties listed above, that we shall call the associated operator to the rigging. Let us also observe that, as the range of $`J`$ is dense, $`J^1`$ is unbounded.
Furthermore, an operator $`J`$ defines a canonical rigging of type (B) and another one of type (C). In fact, let $`N_B=\{n=p/p+1:p=0,1,2,\mathrm{}\}`$ and let $`N_C=\{0,1,2,\mathrm{}\}.`$ Let us also define on $`\mathrm{\Psi }`$ the Hilbert norms:
$$J\rho _n^2=J^n\rho J^n\rho _n=\rho \rho _{}$$
(21)
where $`nN_{B\text{ }}`$ in the first case, and $`nN_C`$ in the second one. (If we had chosen as $`N_B`$ any other subset of rational, non-negative numbers with a supreme equal to 1, we would have had another equivalent sequence of norms, yielding the same $`\mathrm{\Phi }_K.`$ If we had chosen as $`N_C`$ another growing sequence of rational numbers we would have had the same space $`\mathrm{\Phi }).`$
Even if $`J:\mathrm{\Phi }_H`$ is an isometry, this would not be the case for $`J,`$ considered as an operator $`J:`$. In fact, the difference between the second operator and the first one, is that the latter establishes the ”deformation power” of the former. For instance, the first one maps the unit sphere $`S`$ of $``$ in the unit sphere $`S_H`$ of $`\mathrm{\Phi }_{H.}`$ On the other hand, if the second one is bounded or continuous, $`S_H`$ is only a bounded set of $``$. But if it is compact , $`S_H`$ will be an ellipsoid whose semiaxis go to zero, i.e. the action of $`J`$ is more ”drastic”. If it is Hilbert-Schmidt or nuclear , $`S_H`$ will be an ellipsoid with semiaxis going to zero in $`l^2`$ or in $`l^1`$, respectively.
Example 1.
Let us suppose that $`J:`$ is a compact operator. Then we will show that $`\mathrm{\Phi }_K`$ is a Köthe space . As $`J`$ is compact, there is an orthonormal basis $`\{\rho _k\}`$ of $``$ and a sequence of numbers $`\lambda _k0,`$ $`\lambda _k0`$, such that:
$`\text{for every }\sigma \text{, it is: }\sigma ={\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}a_k\rho _k\text{ , and }J\sigma ={\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\lambda _k\sigma \rho _k_{}\rho _k`$
Then for any $`nN_B`$, we have:
$`J^n\sigma J^n\rho _n`$ $`=`$ $`\sigma \rho _{}={\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}a_k^{}b_k={\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\sigma \rho _k_{}^{}\rho \rho _k_{}=`$
$`=`$ $`{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}J^n\sigma \rho _k_{}^{}J^n\rho \rho _k_{}\lambda _k^{2n}`$
In other words:
$`\mathrm{\Phi }_n=l^2(\lambda _k^n)=\left\{\left\{a_k\right\}/{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\left|a_k\right|^2\lambda _k^{2n}<\mathrm{}\right\}`$
and $`\mathrm{\Phi }_K=_{nN_B}\mathrm{\Phi }_n`$ , endowed with the sequence of Hilbert norms of all these spaces $`\mathrm{\Phi }_n,`$ is by definition, a Köthe-Toeplitz space.
Furthermore if $`J`$ is not only compact but also satisfies the condition:
$$\overline{l\stackrel{´}{ı}m}\frac{\lambda _{k+1}}{\lambda _k}<1$$
(22)
(which in particular, using the quotient theorem for series, implies that $`_{k=1}^{\mathrm{}}\lambda _k<\mathrm{}`$, and therefore that $`J`$ is nuclear) then it is attained a necessary and sufficient condition for $`\mathrm{\Phi }_K`$ being a Köthe nuclear space, namely: for every $`n_1N_B`$ there exists a $`n_2N_B,`$ $`0n_1<n_2<1,`$ such that:
$$\underset{k=1}{\overset{\mathrm{}}{}}\frac{\lambda _k^{2n_1}}{\lambda _k^{2n_2}}<\mathrm{}$$
(23)
In fact, using the same quotient theorem, but now in the serie (3.10), we have:
$$\frac{\frac{\lambda _{k+1}^{2n_1}}{\lambda _{k+1}^{2n_2}}}{\frac{\lambda _k^{2n_1}}{\lambda _k^{2n_2}}}=\left(\frac{\lambda _{k+1}}{\lambda _k}\right)^{2(n_2n_1)}$$
(24)
But from eq. (3.9) and considering that $`n_2n_1>0`$ we have:
$$\overline{l\stackrel{´}{ı}m}\left(\frac{\lambda _{k+1}}{\lambda _k}\right)^{2(n_2n_1)}<1$$
(25)
Example 2.
Let us now consider the most typical quantum mechanics rigging, as explained in . The physical quantum system is represented by an algebra of observables, acting on a vector space with an inner product $`.|..`$ $`R_n`$ will denote the eigenspace of the eigenvalue $`n`$ of the hamiltonian of the system. $`N`$ will be the ”number of modes” operator of the system. If $`N_C=\{0,1,2,\mathrm{}\},`$ let $`\mathrm{\Psi }=\underset{nN_C}{}R_n,`$ be the set of all those states that are finite linear combinations of the energy eigenstates of the system. For each $`nN_C,`$ let us define a scalar product $`.|._n`$ on $`\mathrm{\Psi }`$ as:
$$\varphi \psi _n=\varphi (N+I)^n\psi \text{ for every pair }\varphi \text{ and }\psi \text{ of }\mathrm{\Psi }$$
(26)
or which is equivalent, if $`J=(N+I)^{\frac{1}{2}},\varphi =J^n\sigma ,\psi =J^n\rho :`$
$$J^n\sigma J^n\rho _n=(N+I)^{n/2}\sigma (N+I)^{n/2}\rho _n=\sigma \rho $$
(27)
In this way we have defined a family of Hilbert norms:
$$\psi _n=\sqrt{\psi \psi _n}$$
(28)
and the corresponding rigging of the Hilbert space $``$ obtained by the completion of $`\mathrm{\Psi }`$ with the norm (3.15) for $`n=0`$. In this case, the associated mapping $`J`$ , has the spectrum:
$$\{\lambda _k:k=1,2,\mathrm{}\}=\{\frac{1}{\sqrt{k+1}}\text{ : }k=1,2,..\}$$
(29)
Then, $`J^1`$ has spectrum $`\{\sqrt{\overline{k+1}}:k=1,2,\mathrm{}\}`$, which is an unbounded set, and therefore $`J^1`$ is an unbounded operator. Furthermore, it is obvious that, even if $`J`$ is neither a nuclear operator (since $`_{k=1}^{\mathrm{}}\lambda _k=\mathrm{})`$ nor a Hilbert-Schmidt one (since $`_{k=1}^{\mathrm{}}`$ $`\lambda _k^2=\mathrm{}),`$ the powers of $`J`$ are nuclear operators. In fact, as $`J`$ <sup>2</sup> is a Hilbert-Schmidt operator, it follows that $`J^n`$ is nuclear for every $`n4.`$ That’s why, in this case, we have rigged $``$ with a nuclear space. It’s important to remark that this is a rigging of $``$, not of $``$, and that this system is dynamically stable. So, this $`J`$ is not a $`\mathrm{\Lambda }`$ in the sense of the Brussels group.
Finally let us observe that there exists riggings of more general types , i.e. using nonmetrizable spaces. But they are not relevant for this paper.
## IV Equivalence $`\mathrm{\Lambda }(T)`$-coherent rigging.
From all we have said it is obvious that, given a dynamical system with an internal time operator $`T`$ and a $`𝚲=I_D\mathrm{\Lambda }`$ such that $`\mathrm{\Lambda }=\lambda (T),`$ then there is a canonic operator $`J`$ defined as $`J=\mathrm{\Lambda }`$, endowed with the necessary properties to define a rigging of $``$ of each one of the three types (A), (B), and (C). This rigging is deeply related with the dynamics, since it is defined through a $`\mathrm{\Lambda }`$, and eq. (2.8) is valid. This means that a relation exists between that family of growing norms, defined by $`\mathrm{\Lambda }`$ via eqs. (3.6) or (3.8) and the time evolution of the dynamics. In fact, both $`\mathrm{\Lambda }`$ and the Liapunov variable $`\mathrm{\Lambda }^2=R`$, are ”decreasing functions” of $`T,`$ in the sense that they are respectively equal to $`\lambda (T)`$ and $`\lambda ^2(T),`$ being $`\lambda (t)0`$ for $`t\mathrm{}.`$ In these cases we will say that the rigging is coherent with the dynamics.
Since a $`\mathrm{\Lambda }`$ defines a type (A) rigging, we have a base to say that the stochastic process whose semigroup of contracting operators is $`W_t=\mathrm{\Lambda }U_t\mathrm{\Lambda }^1`$ in $``$ is, in some sense, ”equivalent” (for $`t>0`$) to the dynamical system whose group of unitary evolution is $`U_t.`$ In fact, as we have said $`\mathrm{\Lambda }:\mathrm{\Phi }_H`$ is an isometry and therefore, in $`\mathrm{\Phi }_H`$, the above process is the isometric image, for $`t`$ $`0,`$ of the $`U_t`$ dynamics. Something very similar, but not so ”perfect”, happens if we use the type (B) rigging. In this case we have a sequence of isometries $`\{\mathrm{\Lambda }^n:\mathrm{\Phi }_n\}`$ with $`nN_B,`$ being $`\mathrm{\Lambda }^n\mathrm{\Lambda }`$ if $`n1`$ ($`n=p/p+1`$ if $`p\mathrm{}).`$ Thus $`W_t`$ turns out to be the limit of a sequence of isometric images of $`U_t.`$
Let us now consider any rigging of $`,`$ which is coherent with a dynamic that has an internal time $`T`$. This rigging may be of any of the three types (A), (B), or (C), and it must be defined by a unique $`J=\mathrm{\Lambda }=\lambda (T),`$ with $`\lambda :[0,1]`$ endowed with the properties listed above the eq. (2.6). Then let us define : $`𝚲=I_𝒟\mathrm{\Lambda }.`$ In this way we have all the properties of $`𝚲`$ with the exception of the normalization and the monotonous convergence of $`W_t\rho `$ to zero.
The normalization turns out to be trivial in the quantum case, since it is defined by the diagonal part: if $`\rho `$ is normal then:
$$\left(𝚲\rho \right)(I)=(I_𝒟\rho ^d)(I)=\rho ^d(I)=1$$
(30)
In the classical case, if $`\rho `$ is normal, namely if $`_\mathrm{\Omega }\rho 𝑑\mu =1,`$ then:
$$_\mathrm{\Omega }𝚲(\rho )\text{ }𝑑\mu =_\mathrm{\Omega }[\rho ^d+\mathrm{\Lambda }(\rho ^c)]\text{ }𝑑\mu =_\mathrm{\Omega }\rho ^d\text{ }𝑑\mu +_\mathrm{\Omega }\mathrm{\Lambda }(\rho ^c)\text{ }𝑑\mu =_\mathrm{\Omega }\rho ^d\text{ }𝑑\mu =1$$
(31)
since in this case $`=𝒟^{}`$ and $`1𝒟`$, while $`\mathrm{\Lambda }(\rho ^c)`$, then:
$$_\mathrm{\Omega }\mathrm{\Lambda }(\rho ^c)\text{ }𝑑\mu =_\mathrm{\Omega }\mathrm{\Lambda }(\rho ^c).1\text{ }𝑑\mu =\mathrm{\Lambda }(\rho ^c)1=0$$
(32)
(for the same reason $`_\mathrm{\Omega }\rho ^cd\mu =0).`$
Let us now consider $`W_t=\mathrm{\Lambda }U_t\mathrm{\Lambda }^1,`$ where $`\mathrm{\Lambda }^1=\lambda ^1(T)=_{}\frac{1}{\lambda (s)}𝑑E.`$ Then, for any $`\rho `$ in the domain of $`\mathrm{\Lambda }^1`$, we have:
$`W_t\rho _{}^2`$ $`=`$ $`\mathrm{\Lambda }U_t\mathrm{\Lambda }^1\rho _{}^2=(U_t^{}\mathrm{\Lambda }U_t)\mathrm{\Lambda }^1\rho _{}^2=`$ (33)
$`=`$ $`\lambda (T+t)\lambda ^1(T)\rho _{}^2={\displaystyle _{}}\left[{\displaystyle \frac{\lambda (s+t)}{\lambda (s)}}\right]^2dE\rho _{}^2`$ (34)
being the function $`s\left[\frac{\lambda (s+t)}{\lambda (s)}\right]^2`$ non negative and bounded by the integrable function $`1`$. As $`\frac{\lambda (s+t)}{\lambda (s)}`$ goes monotonously to $`0`$ (see above eq. (2.6)), from the Lebesque dominated convergence theorem we have that $`W_t\rho 0.`$
## V Synthesis of both formalisms.
In this section, we will relate the formalism of the $`\mathrm{\Lambda }`$ with the formalism of a rigging with $`\mathrm{\Phi }_H.`$
Let us consider the Koopman operator of the dynamic $`U_t:`$, with certain $`\mathrm{\Lambda }:`$ We have proved that this is equivalent to a rigging of $``$ with a Hilbert space $`\mathrm{\Phi }_H`$ with an inner product:
$$\mathrm{\Lambda }\sigma \mathrm{\Lambda }\rho _{\mathrm{\Phi }_H}=\sigma \rho _{}$$
(35)
as well as with its antidual $`\mathrm{\Phi }_H^\times `$. This rigging defines the following operators:
1) $`\mathrm{\Lambda }:\mathrm{\Phi }_H,`$ namely function $`\mathrm{\Lambda }`$, but with a restricted range.
2) $`\mathrm{\Lambda }^\times :\mathrm{\Phi }_H^\times `$, namely the antitransposed former operator defined as:
$$\rho \mathrm{\Lambda }^\times (F)_{}:=F\left(\mathrm{\Lambda }^{}\rho \right)=F\left(\mathrm{\Lambda }\rho \right)$$
(36)
where, for simplicity, only for $``$ we have made the identification: $`=^\times .`$
3) $`:\mathrm{\Phi }_H^\times \mathrm{\Phi }_H,`$ namely the Riesz representation (already defined in eq. (3.1)), which is related with the former operator by:
$$=\mathrm{\Lambda }\mathrm{\Lambda }^\times $$
(37)
and its inverse $`^1:\mathrm{\Phi }_H\mathrm{\Phi }_H^\times .`$
4) $`|_{}=R0`$, such that:
$$R=\mathrm{\Lambda }^2\text{ , or which is equivalent, }\mathrm{\Lambda }=\sqrt{R}$$
(38)
5) $`\mathrm{\Lambda }^1`$as an extension of the $`\mathrm{\Lambda }`$ inverted operator, or which is the same thing:
$$\mathrm{\Lambda }^1=\sqrt{R^1}$$
(39)
We also obtain some important operators combining the rigging with the dynamics:
6) $`W_t=\mathrm{\Lambda }U_t\mathrm{\Lambda }^1,`$ $`tG^+`$, namely the evolution operator of the Markov semigroup that we have already considered and on which is based the $`\mathrm{\Lambda }`$ formalism.
7) $`\overline{U}_t:\mathrm{\Phi }_H^\times \mathrm{\Phi }_H^\times ,`$ $`tG^+,`$ namely the extension of a semigroup of $`U_t`$ to $`\mathrm{\Phi }_H^\times `$, defined by:
$$\left(\overline{U_t}(F)\right)(\mathrm{\Lambda }\rho )=F\left(U_t\mathrm{\Lambda }\rho \right)$$
(40)
which is the base of the rigging formalism.
8) $`Y_t:,`$ $`tG^+`$ defined as:
$$Y_t=\mathrm{\Lambda }^\times \overline{U_t}\left(\mathrm{\Lambda }^\times \right)^1$$
(41)
9) $`V_t:\mathrm{\Phi }_H^\times \mathrm{\Phi }_H^\times ,`$ $`tG^+`$ defined as:
$$V_t=\left(\mathrm{\Lambda }^\times \right)^1U_t\mathrm{\Lambda }^\times $$
(42)
10) $`Z_t:\mathrm{\Phi }_H\mathrm{\Phi }_{H,}`$ $`tG^+`$ defined as:
$$Z_t=\overline{U_t}^1$$
(43)
11) $`X_t:\mathrm{\Phi }_H^\times \mathrm{\Phi }_H^\times ,`$ $`tG^+`$ defined as:
$$X_t=^1W_t$$
(44)
The following propositions make clear the deep relation among all these operators.
Theorem:
For any $`tG^+,`$ we have:
i) $`V_t=X_t.`$
ii) $`V_t`$ defines a strong Markov process.
iii) $`V_t\overline{U}_t.`$
iv) $`Y_tU_t.`$
v) $`W_tZ_t.`$
vi) $`Z_t`$ defines a dynamic which is equivalent to $`\overline{U}_t.`$
Demonstration:
Let $`tG^{+,}F\mathrm{\Phi }_H^\times ,`$ and $`\rho \mathrm{\Phi }_H.`$ Then we have:
$$X_t=(\mathrm{\Lambda }\mathrm{\Lambda }^\times )^1W_t(\mathrm{\Lambda }\mathrm{\Lambda }^\times )=\left(\mathrm{\Lambda }^\times \right)^1(\mathrm{\Lambda }^1W_t\mathrm{\Lambda })\mathrm{\Lambda }^\times =\left(\mathrm{\Lambda }^\times \right)^1U_t\mathrm{\Lambda }^\times =V_t$$
(45)
and so (i) is demonstrated.
As Riesz representation is an isometric isomorphism, then $`X_t`$ in $`\mathrm{\Phi }_H^\times `$ is equivalent to $`W_t`$ in $`\mathrm{\Phi }_H.`$ Now, we have just demonstrated that $`V_t=X_t`$, so (ii) is also demonstrated.
In order to demonstrate (iii), it is enough to show that:
$$\mathrm{\Lambda }^\times \overline{U_t}U_t\mathrm{\Lambda }^\times $$
(46)
Now, it is:
$$\rho \mathrm{\Lambda }^\times \overline{U_t}(F)=\left[\overline{U_t}(F)\right](\mathrm{\Lambda }\rho )=F(U_t\mathrm{\Lambda }\rho )$$
(47)
while:
$$\rho U_t\mathrm{\Lambda }^\times (F)=U_t\rho \mathrm{\Lambda }^\times (F)=F(\mathrm{\Lambda }U_t\rho )$$
(48)
As we know that $`\mathrm{\Lambda }`$ do not commute with $`U_t`$ (see (2.12)), it turns out that the r.h.s. of the two last equations are not equal for all $`F`$ and for all $`\rho `$. Therefore eq. (5.12) is proved. So (iii) is demonstrated.
As $`(\mathrm{\Lambda }^\times )^1`$ is a bijection , we have:
$$Y_t=\mathrm{\Lambda }^\times \overline{U_t}\left(\mathrm{\Lambda }^\times \right)^1\mathrm{\Lambda }^\times V_t\left(\mathrm{\Lambda }^\times \right)^1=U_t$$
(49)
which proves (iv).
Finally, if we take into account (i), plus the relation that can be obtained from eq. (5.10), and the fact that $``$ is bijective:
$$Z_t=\overline{U_t}^1V_t^1=W_t$$
(50)
The last part is similar to the proof of (ii).
## VI Acknowledgement
The author wishes to thank Professor Mario A.Castagnino and Dr. Roberto Laura for many teachings, important dialogs and continuous encouragement, and also to the ”unknown referee”, whose deep remarks helped me to understand some mistakes.
## VII References
B.Misra, I.Prigogine, M.Courbage. Physica 98 A (1979) 1-26
Idem Proc. Nat. Acad. of Sci. USA 76 (1979) 4768-72
I.E.Antoniou, These, presentée dans l’Université Libre de Bruxelles.
S.Goldstein, B.Misra, M.Courbage. J.Stat.Phys. 25 (1981) 111-126
I.E.Antoniou, I.Prigogine. Physica A 192 (1993) 443-464
I.E.Antoniou, S.Tasaki. Int.Quantum Chemistry, 46 (1993), 425-474.
Idem. Physica A, 190 (1992),303.
I.Prigogine, ”Time, Chaos and the Laws of Nature”, Conference, July 27 th, 1994.
M.A.Castagnino, F.Gaioli, E.Gunzig. Found. of Cosmic Physic. In press.
A.Böhm, M.Gadella, ”Dirac Kets, Gamow vectors and Gel’fand Triplets”,
Springer Verlag (1989).
M.Cotlar, ”Equipación con Espacios de Hilbert”, Cursos y Seminarios de Matemática (Número 15, 1968), Univ. Buenos Aires.
I.M.Gelfand, G.E.Chilov, ”Les Distributions”, vol. 2, Dunod (1964)
I.M.Gelfand, N.Y.Vilenkin, ”Les Distributions”, vol 4, Dunod (1967)
A.Böhm, ”The Rigged Hilbert Space and Quantum Mechanics”, Lecture Notes in Physics, Número 78.
I.Antoniou, Z.Suchanecki, in ’Nonlinear, deformed and irreversible quantum systems’. H.D. Doebner et al editors, World Scientific (1995)
I.Antoniou, Z.Suchanecki, R.Laura, S.Tasaki, ’Intrinsic irreversibility of quantum systems with diagonal singularity’, Physica A, 241 (1997), 737-772.
M.Reed, B.Simon, ”Methods of Modern Mathematical Physics”, Acad.Press, vol. 1 (1979)
M.Mackey,”Times Arrow: The Origins of Thermodynamic Behaviour”Springer Verlag (1992)
V.I.Arnold, A.Avez, ”Ergodic Problems of Classical Mechanics”, Benjamin (1968)
W.Rudin, ”Functional Analysis”, Mc. Graw Hill (1973)
R.Thom, ”Structural Stability and Morphogenesis”,Benjamin (1975)
B.Misra, Proc.Natl.Acad.Sci.USA, vol 75, N4 (April 1978) 1627-31
K.Maurin, ”General eigenfuction expansions and unitary representations of topological groups”, Warzawa (1968)
N.N.Bogolubov, A.A.Logunov and I.T.Todorov, ”Introduction to Axiomatic Quantum Field Theory”, Benjamin (1975)
R.Laura and A.R.Ordóñez, ”Internal Time Superoperator for Quantum Systems with Diagonal Singularity”, to appear in: Inter.Jour.Theor.Phys. (1997)
G.G.Emch, Commun. Math. Phys. 49, 191-215, (1976)
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# 1 Introduction
## 1 Introduction
We address the Lagrangian antifield BRST formalism of , phrased in terms of exterior forms of finite jet order on the infinite order jet space of physical fields, ghosts and antifields on a base manifold $`X`$. Horizontal (semibasic) exterior forms constitute a bicomplex with respect to the BRST operator $`𝐬`$ and the horizontal (total) differential $`d_H`$. It is graded by the ghost number $`k`$ and the form degree $`p`$. We aim to study the iterated $`𝐬`$-cohomology $`H^{k,p}(𝐬|d_H)`$ of the $`d_H`$-cohomology groups of this bicomplex (i.e., the term $`E_2^,`$ of its spectral sequence ). In terms of form degree $`p=n=dimX`$, this cohomology coincides with the well-known local BRST cohomology (i.e., $`𝐬`$-cohomology modulo $`d_H`$). If $`p<n`$, iterated BRST cohomology unlike local BRST cohomology is defined only for $`d_H`$-closed forms.
The above mentioned BRST formalism is usually formulated on a contractible base $`X`$ when the $`d_H`$-cohomology of form degree $`0<p<n`$ are trivial in accordance with the algebraic Poincaré lemma (see Lemma 4 below). If $`X=R^n`$, there is an isomorphism of the iterated (and local) BRST cohomology groups $`H^{k,n}(𝐬|d_H)`$, $`kn`$, to the cohomology $`H_{\mathrm{tot}}^{k+n}`$ groups of the total BRST operator $`\stackrel{~}{𝐬}=𝐬+d_H`$ on horizontal exterior forms of total ghost number $`k+n`$ . We generalize this result to the case of iterated BRST cohomology of form degree $`p<n`$ and an arbitrary connected manifold $`X`$ (see Corollary 1 below). To construct the corresponding (global) descent equations, one needs the $`d_H`$-cohomology of exterior forms on the infinite order jet space. The study of this cohomology is the key point of our consideration.
Note that the descent equations for representativers of local BRST cohomology groups of form degree $`p<n`$ are also constructed, but this cohomology fails to be related to cohomology of the total BRST operator (see also where BRST cohomology modulo the exterior differential $`d`$ in the Yang–Mills gauge theory are considered).
To avoid the sophisticate techniques of ghosts and antifields , we here study iterated cohomology of the graded differential algebra $`𝒪_{\mathrm{}}^{}`$ of exterior forms of finite jet order on the infinite order jet space $`J^{\mathrm{}}Y`$ of an affine bundle $`YX`$. Note that affine bundles provide a standard framework in quantum field theory. Let $`𝒪_{\mathrm{}}^{}`$ be endowed with a nilpotent form degree preserving endomorphism $`𝐬`$ such that horizontal elements of $`𝒪_{\mathrm{}}^{}`$ constitute a bicomplex with respect to $`𝐬`$ and $`d_H`$, i.e., $`𝐬d_H+d_H𝐬=0`$. We agree to call a gradation degree $`k`$ with respect to $`𝐬`$ the ghost number. Suppose that $`𝐬`$ vanishes on exterior forms on $`X`$ and that these forms are of zero ghost number and are not $`𝐬`$-exact. The goal is the following.
THEOREM 1. If $`YX`$ is an affine bundle, the iterated $`𝐬|d_H`$-cohomology is the following.
(i) The $`H^{0,p<n}(𝐬|d_H)`$ is trivial.
(ii) $`H^{0,p<n}(𝐬|d_H)=H^p(X)`$ where $`H^{}(X)`$ is de Rham cohomology of $`X`$.
(iii) $`H^{k,n}(𝐬|d_H)=H_{\mathrm{tot}}^{k+n}`$, $`k<n`$ or $`k>1`$.
(iv) Let $`\gamma _p:H^p(X)H_{\mathrm{tot}}^p`$, $`0p<n`$, be a natural homomorphism corresponding to the monomorphism of the algebra $`𝒪^{}(X)`$ of exterior forms on $`X`$ to $`𝒪_{\mathrm{}}^{}`$. Put $`\overline{H}^p=H_{\mathrm{tot}}^p/\mathrm{Im}\gamma _p`$. If $`nk<1`$, there is a monomorphism $`\overline{H}^{k+n}H^{k,n}(𝐬|d_H)`$ such that $`H^{k,n}(𝐬|d_H)/\overline{H}^{k+n}=\mathrm{Ker}\gamma _{k+n+1}`$. In particular, $`\mathrm{Ker}\gamma _0=0`$ and $`\overline{H}^0=H_{\mathrm{tot}}^0/H^0(X)`$.
(v) $`H^{1,n}(𝐬|d_H)=\overline{H}^{n1}`$.
Note that the operator $`𝐬`$ in Theorem 1 may have different physical origins. The following corollary of Theorem 1 corresponds to the case of iterated BRST cohomology.
COROLLARY 2. Let $`YX`$ be a vector bundle and $`P_{\mathrm{}}^{}𝒪_{\mathrm{}}^{}`$ a subalgebra of exterior forms which are polynomial in fibre coordinates of $`J^{\mathrm{}}YX`$. There is the decomposition of $`C^{\mathrm{}}(X)`$-modules $`P_{\mathrm{}}^{}=𝒪^{}(X)(P_{\mathrm{}}^{})_{>0}`$. Let $`\stackrel{~}{𝐬}(P_{\mathrm{}}^{})(P_{\mathrm{}}^{})_{>0}`$. Then, $`\gamma _p`$, $`0p<n`$, are monomorphisms, and the items (iv), (v) of Theorem 1 state isomorphisms $`H^{k,n}(𝐬|d_H)=H_{\mathrm{tot}}^{k+n}/H^{k+n}(X)`$, $`nk1`$.
In particular, if $`X=R^n`$, the iterated BRST cohomology of form degree $`p<n`$ (except $`H^{0,0}(𝐬|d_H)=R`$) is always trivial in contrast with the local BRST cohomology.
In fact, to prove Theorem 1, we will obtain $`d_H`$\- and $`\delta `$-cohomology of the variational complex of the graded differential algebra $`𝒪_{\mathrm{}}^{}`$ in the case of an arbitrary smooth bundle $`YX`$.
The $`𝒪_{\mathrm{}}^{}`$ is the direct limit of graded differential algebras of exterior forms on finite order jet manifolds. It consists of exterior forms on finite order jet manifolds modulo the pull-back identification. Passing to the direct limit of the de Rham complexes of exterior forms on finite order jet manifolds, de Rham cohomology of $`𝒪_{\mathrm{}}^{}`$ has been found to coincide with de Rham cohomology of the bundle $`Y`$ . However, this is not a way of studying other cohomology groups of $`𝒪_{\mathrm{}}^{}`$. Therefore, we enlarge $`𝒪_{\mathrm{}}^{}`$ to the structure algebra $`𝒯_{\mathrm{}}^{}`$ of the sheaf of germs of exterior forms on finite order jet manifolds. One can say that $`𝒯_{\mathrm{}}^{}`$ consists of exterior forms of locally finite jet order on $`J^{\mathrm{}}Y`$. The $`d_H`$\- and $`\delta `$-cohomology of $`𝒯^{\mathrm{}}`$ has been investigated in . We simplify this investigation due to Lemma 4 below and prove that $`𝒪_{\mathrm{}}^{}`$ and $`𝒯_{\mathrm{}}^{}`$ have the same $`d_H`$\- and $`\delta `$-cohomology (see Theorem 6 below). In particular, this provides a solution of the global inverse problem of the calculus of variations in the class of exterior forms of finite jet order.
For the proof of Theorem 1, it is quite important that, if $`YX`$ is an affine bundle, $`d_H`$-cohomology $`H^{<n}(d_H;𝒪_{\mathrm{}}^{})`$ of $`𝒪_{\mathrm{}}^{}`$ coincides with de Rham cohomology of the base $`X`$. It follows that every $`d_H`$-closed $`(k<n)`$-form $`\varphi 𝒪_{\mathrm{}}^{}`$ splits into the sum $`\varphi =\phi +d_H\xi `$ where $`\xi 𝒪_{\mathrm{}}^{}`$ and $`\phi `$ is a closed form on $`X`$. Since the operator $`𝐬`$ annihilates these forms, the system of global descent equations can be constructed though its right-hand side is not zero. We come to the same result for the polynomial algebra $`P_{\mathrm{}}^{}`$ and its subalgebra $`\overline{P}_{\mathrm{}}^{}`$ of $`x`$-independent forms.
## 2 The differential calculus on $`J^{\mathrm{}}Y`$
Smooth manifolds throughout are assumed to be real, finite-dimensional, Hausdorff, paracompact, and connected. Put further dim$`X=n`$. The standard notation of jet formalism is utilized. Following the terminology of , by a sheaf $`S`$ on a topological space $`Z`$ is meant a sheaf bundle $`SZ`$. Accordingly, $`\mathrm{\Gamma }(S)`$ denotes the canonical presheaf of sections of the sheaf $`S`$, and $`\mathrm{\Gamma }(Z,S)`$ is the set of global sections of $`S`$.
Recall that the infinite order jet space $`J^{\mathrm{}}Y`$ of a smooth bundle $`YX`$ is defined as a projective limit $`(J^{\mathrm{}}Y,\pi _r^{\mathrm{}})`$ of the inverse system
$$X\stackrel{\pi }{}Y\stackrel{\pi _0^1}{}\mathrm{}J^{r1}Y\stackrel{\pi _{r1}^r}{}J^rY\mathrm{}$$
(1)
of finite order jet manifolds $`J^rY`$ of $`YX`$, where $`\pi _{r1}^r`$ are affine bundles. Bearing in mind Borel’s theorem, one can say that $`J^{\mathrm{}}Y`$ consists of the equivalence classes of sections of $`YX`$ identified by their Taylor series at points of $`X`$. Endowed with the projective limit topology, $`J^{\mathrm{}}Y`$ is a paracompact Fréchet manifold . A bundle coordinate atlas $`\{U_Y,(x^\lambda ,y^i)\}`$ of $`Y`$ yields the manifold coordinate atlas $`\{(\pi _0^{\mathrm{}})^1(U_Y),(x^\lambda ,y_\mathrm{\Lambda }^i)\}`$, $`0|\mathrm{\Lambda }|`$, of $`J^{\mathrm{}}Y`$, together with the transition functions
$$y_{}^{}{}_{\lambda +\mathrm{\Lambda }}{}^{i}=\frac{x^\mu }{x^\lambda }d_\mu y_\mathrm{\Lambda }^i,$$
(2)
where $`d_\mu `$ denotes the total derivative $`d_\mu =_\mu +\underset{|\mathrm{\Lambda }|0}{}y_{\mu +\mathrm{\Lambda }}^i_i^\mathrm{\Lambda }`$.
With the inverse system (1), one has the direct system
$$𝒪^{}(X)\stackrel{\pi ^{}}{}𝒪_0^{}\stackrel{\pi _0^1^{}}{}𝒪_1^{}\stackrel{\pi _1^2^{}}{}\mathrm{}\stackrel{\pi _{r1}^r^{}}{}𝒪_r^{}\mathrm{}$$
(3)
of $`R`$-algebras $`𝒪_r^{}`$ of exterior forms on finite order jet manifolds, where $`\pi _{r1}^r^{}`$ are pull-back monomorphisms. Its direct limit is the above mentioned graded differential $`R`$-algebra $`𝒪_{\mathrm{}}^{}`$. It is a differential calculus over the $`R`$-ring $`𝒪_{\mathrm{}}^0`$ of continuous real functions on $`J^{\mathrm{}}Y`$ which are the pull-back of smooth functions on finite order jet manifolds. Let us enlarge the ring $`𝒪_{\mathrm{}}^0`$ to the $`R`$-ring $`𝒯_{\mathrm{}}^0`$ of continuous real functions on $`J^{\mathrm{}}Y`$ such that, given $`f𝒯_{\mathrm{}}^0`$ and any point $`qJ^{\mathrm{}}Y`$, there exists a neighborhood of $`q`$ where $`f`$ coincides with the pull-back of a smooth function on some finite order jet manifold. Let $`𝔒_r^{}`$ be a sheaf of germs of exterior forms on the $`r`$-order jet manifold $`J^rY`$ and $`\mathrm{\Gamma }(𝔒_r^{})`$ its canonical presheaf. There is the direct system of canonical presheaves
$`\mathrm{\Gamma }(𝔒_X^{})\stackrel{\pi ^{}}{}\mathrm{\Gamma }(𝔒_0^{})\stackrel{\pi _0^1^{}}{}\mathrm{\Gamma }(𝔒_1^{})\stackrel{\pi _1^2^{}}{}\mathrm{}\stackrel{\pi _{r1}^r^{}}{}\mathrm{\Gamma }(𝔒_r^{})\mathrm{}.`$
Its direct limit $`𝔒_{\mathrm{}}^{}`$ is a presheaf of graded differential $`R`$-algebras on $`J^{\mathrm{}}Y`$. Let $`𝔗_{\mathrm{}}^{}`$ be the sheaf constructed from $`𝔒_{\mathrm{}}^{}`$ and $`\mathrm{\Gamma }(𝔗_{\mathrm{}}^{})`$ its canonical presheaf. The structure algebra $`𝒯_{\mathrm{}}^{}=\mathrm{\Gamma }(J^{\mathrm{}}Y,𝔗_{\mathrm{}}^{})`$ of the sheaf $`𝔗_{\mathrm{}}^{}`$ is a differential calculus over the $`R`$-ring $`𝒯_{\mathrm{}}^{}`$. There are the $`R`$-algebra monomorphisms $`𝔒_{\mathrm{}}^{}\mathrm{\Gamma }(𝔗_{\mathrm{}}^{})`$ and $`𝒪_{\mathrm{}}^{}\mathrm{\Gamma }(𝒯_{\mathrm{}}^{})`$. Since the paracompact space $`J^{\mathrm{}}Y`$ admits a partition of unity by elements of $`𝒯_{\mathrm{}}^0`$ , sheaves of $`𝒯_{\mathrm{}}^0`$-modules on $`J^{\mathrm{}}Y`$ are fine and acyclic. Then, the abstract de Rham theorem on cohomology of a sheaf resolution can be called into play in order to obtain cohomology of the algebra $`𝒯_{\mathrm{}}^{}`$.
For short, we agree to call elements of $`𝒯_{\mathrm{}}^{}`$ the exterior forms on $`J^{\mathrm{}}Y`$. Restricted to a coordinate chart $`(\pi _0^{\mathrm{}})^1(U_Y)`$ of $`J^{\mathrm{}}Y`$, they can be written in a coordinate form, where horizontal and contact forms $`\{dx^\lambda ;\theta _\mathrm{\Lambda }^i=dy_\mathrm{\Lambda }^iy_{\lambda +\mathrm{\Lambda }}^idx^\lambda \}`$ provide generators of the algebra $`𝒯_{\mathrm{}}^{}`$. There is the canonical decomposition of $`𝒯_{\mathrm{}}^{}`$ into $`𝒯_{\mathrm{}}^0`$-modules $`𝒯_{\mathrm{}}^{k,s}`$ of $`k`$-contact and $`s`$-horizontal forms:
$`𝒯_{\mathrm{}}^{}=\underset{k,s}{}𝒯_{\mathrm{}}^{k,s},h_k:𝒯_{\mathrm{}}^{}𝒯_{\mathrm{}}^{k,},h^s:𝒯_{\mathrm{}}^{}𝒯_{\mathrm{}}^{,s},0k,0sn.`$
Accordingly, the exterior differential on $`𝒯_{\mathrm{}}^{}`$ splits into the sum $`d=d_H+d_V`$ of horizontal and vertical differentials such that
$`d_Hh_k=h_kdh_k,d_H(\varphi )=dx^\lambda d_\lambda (\varphi ),\varphi 𝒯_{\mathrm{}}^{},`$
$`d_Vh^s=h^sdh^s,d_V(\varphi )=\theta _\mathrm{\Lambda }^i_i^\mathrm{\Lambda }\varphi .`$
## 3 The horizontal complex
Being nilpotent, the differentials $`d_V`$ and $`d_H`$ provide the natural bicomplex $`\{𝒯_{\mathrm{}}^{k,m}\}`$ of the graded differential algebra $`𝒯_{\mathrm{}}^{}`$. Let us consider its row
$$0R𝒯_{\mathrm{}}^0\stackrel{d_H}{}𝒯_{\mathrm{}}^{0,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}𝒯_{\mathrm{}}^{0,n}\stackrel{d_H}{}0$$
(4)
called the horizontal complex. The corresponding complex of sheaves
$$0R𝔗_{\mathrm{}}^0\stackrel{d_H}{}𝔗_{\mathrm{}}^{0,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}𝔗_{\mathrm{}}^{0,n}\stackrel{d_H}{}0$$
(5)
except the last term is exact (see Lemma 4 below). Then, since the sheaves $`𝔗^{0,<n}`$ of $`𝒯_{\mathrm{}}^0`$-modules on $`J^{\mathrm{}}Y`$ are fine, we obtain from the abstract de Rham theorem and Lemma (5) below that $`d_H`$-cohomology $`H^{<n}(d_H;𝒯_{\mathrm{}}^{})`$ of the horizontal complex (4) is equal to de Rham cohomology $`H^{}(Y)`$ of the bundle $`Y`$. Theorem 6 below shows that $`d_H`$-cohomology $`H^{<n}(d_H;𝒪_{\mathrm{}}^{})`$ of the horizontal complex
$$0R𝒪_{\mathrm{}}^0\stackrel{d_H}{}𝒪_{\mathrm{}}^{0,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}𝒪_{\mathrm{}}^{0,n}\stackrel{d_H}{}0$$
(6)
of the algebra $`𝒪_{\mathrm{}}^{}`$ is the same. However, one should complete the horizontal complex (4) in the variational one in order to say something on the $`n`$th cohomology group of $`d_H`$ (see the relation (19) below).
## 4 The variational complex
Let us consider the variational operator $`\delta =\tau d`$ on $`𝒯_{\mathrm{}}^{,n}`$ where
$`\tau ={\displaystyle \underset{k>0}{}}{\displaystyle \frac{1}{k}}\overline{\tau }h_kh^n,\overline{\tau }(\varphi )={\displaystyle \underset{0\mathrm{\Lambda }}{}}(1)^\mathrm{\Lambda }\theta ^i[d_\mathrm{\Lambda }(_i^\mathrm{\Lambda }\varphi )],\varphi 𝒯^{>0,n}_{\mathrm{}},`$
is the projection $`R`$-module endomorphism of $`𝒯_{\mathrm{}}^{}`$ such that $`\tau d_H=0`$ (see, e.g., ). The $`\delta `$ is nilpotent, and obeys the relation
$$\delta \tau \tau d=0.$$
(7)
Put $`𝔈_k=\tau (𝔗_{\mathrm{}}^{k,n})`$, $`E_k=\tau (𝒯_{\mathrm{}}^{k,n})`$, $`k>0`$. Since $`\tau `$ is a projector, there are isomorphisms
$`\mathrm{\Gamma }(𝔈_k)=\tau (\mathrm{\Gamma }(𝔗_{\mathrm{}}^{k,n})),E_k=\mathrm{\Gamma }(J^{\mathrm{}}Y,𝔈_k).`$
With operators $`d_H`$ and $`\delta `$, we have the variational complex
$$0R𝔗_{\mathrm{}}^0\stackrel{d_H}{}𝔗_{\mathrm{}}^{0,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}𝔗_{\mathrm{}}^{0,n}\stackrel{\delta }{}𝔈_1\stackrel{\delta }{}𝔈_2\mathrm{}$$
(8)
of the sheaf $`𝔗_{\mathrm{}}^{}`$ and that
$$0R𝒯_{\mathrm{}}^0\stackrel{d_H}{}𝒯_{\mathrm{}}^{0,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}𝒯_{\mathrm{}}^{0,n}\stackrel{\delta }{}E_1\stackrel{\delta }{}E_2\mathrm{}$$
(9)
of its structure algebra $`𝒯_{\mathrm{}}^{}`$. The similar variational complex $`\{𝒪_{\mathrm{}}^{},\overline{E}_k\}`$ of the algebra $`𝒪_{\mathrm{}}^{}`$ takes place. There are the well-known statements summarized usually as the algebraic Poincaré lemma (see, e.g., ).
LEMMA 3. If $`Y`$ is a contractible bundle $`R^{n+p}R^n`$, the variational complex $`\{𝒪_{\mathrm{}}^{},\overline{E}_k\}`$ of the graded differential algebra $`𝒪_{\mathrm{}}^{}`$ is exact.
It follows that the variational complex of sheaves (8) is exact for any smooth bundle $`YX`$. Moreover, the sheaves $`𝔗_{\mathrm{}}^{,n}`$ in this complex are fine, and so are the sheaves $`𝔈_k`$ in accordance with Lemma 4 below. Hence, the variational complex (8) is a resolution of the constant sheaf $`R`$ on $`J^{\mathrm{}}Y`$.
LEMMA 4. Sheaves $`𝔈_k`$, $`k>0`$, are fine.
Proof. Though the $`R`$-modules $`E_{k>1}`$ fail to be $`𝒯_{\mathrm{}}^0`$-modules , one can use the fact that the sheaves $`𝔈_{k>0}`$ are projections $`\tau (𝔗_{\mathrm{}}^{k,n})`$ of sheaves of $`𝒯_{\mathrm{}}^0`$-modules. Let $`𝔘=\{U_i\}_{iI}`$ be a locally finite open covering of $`J^{\mathrm{}}Y`$ and $`\{f_i𝒯_{\mathrm{}}^0\}`$ the associated partition of unity. For any open subset $`UJ^{\mathrm{}}Y`$ and any section $`\phi `$ of the sheaf $`𝔗_{\mathrm{}}^{k,n}`$ over $`U`$, let us put $`h_i(\phi )=f_i\phi `$. The endomorphisms $`h_i`$ of $`𝔗_{\mathrm{}}^{k,n}`$ yield the $`R`$-module endomorphisms
$`\overline{h}_i=\tau h_i:𝔈_k\stackrel{\mathrm{in}}{}𝔗_{\mathrm{}}^{k,n}\stackrel{h_i}{}𝔗_{\mathrm{}}^{k,n}\stackrel{\tau }{}𝔈_k`$
of the sheaves $`𝔈_k`$. They possess the properties required for $`𝔈_k`$ to be a fine sheaf. Indeed, for each $`iI`$, $`\mathrm{supp}f_iU_i`$ provides a closed set such that $`\overline{h}_i`$ is zero outside this set, while the sum $`\underset{iI}{}\overline{h}_i`$ is the identity morphism. $`\mathrm{}`$
## 5 Cohomology of $`𝒯_{\mathrm{}}^{}`$
LEMMA 5. There is an isomorphism
$$H^{}(J^{\mathrm{}}Y,R)=H^{}(Y,R)=H^{}(Y)$$
(10)
between cohomology $`H^{}(J^{\mathrm{}}Y,R)`$ of $`J^{\mathrm{}}Y`$ with coefficients in the constant sheaf $`R`$, that $`H^{}(Y,R)`$ of $`Y`$, and de Rham cohomology $`H^{}(Y)`$ of $`Y`$.
Proof. Since $`Y`$ is a strong deformation retract of $`J^{\mathrm{}}Y`$ (see, e.g., ), the first isomorphism in (10) follows from the Vietoris–Begle theorem , while the second one results from the well-known de Rham theorem. $`\mathrm{}`$
Let us consider the de Rham complex of sheaves
$$0R𝔗_{\mathrm{}}^0\stackrel{d}{}𝔗_{\mathrm{}}^1\stackrel{d}{}\mathrm{}$$
(11)
on $`J^{\mathrm{}}Y`$ and the corresponding de Rham complex of their structure algebras
$$0R𝒯_{\mathrm{}}^0\stackrel{d}{}𝒯_{\mathrm{}}^1\stackrel{d}{}\mathrm{}.$$
(12)
The complex (11) is exact due to the Poincaré lemma, and is a resolution of the constant sheaf $`R`$ on $`J^{\mathrm{}}Y`$ since sheaves $`𝔗_{\mathrm{}}^r`$ are fine. Then, the abstract de Rham theorem and Lemma 5 lead to the following.
PROPOSITION 6. De Rham cohomology $`H^{}(𝒯_{\mathrm{}}^{})`$ of the graded differential algebra $`𝒯_{\mathrm{}}^{}`$ is isomorphic to that $`H^{}(Y)`$ of the bundle $`Y`$.
It follows that every closed form $`\varphi 𝒯_{\mathrm{}}^{}`$ splits into the sum
$$\varphi =\phi +d\xi ,\xi 𝒯_{\mathrm{}}^{},$$
(13)
where $`\phi `$ is a closed form on the bundle $`Y`$.
Similarly, from the abstract de Rham theorem and Lemma 5, we obtain the following.
PROPOSITION 7. There is an isomorphism between $`d_H`$\- and $`\delta `$-cohomology of the variational complex (9) and de Rham cohomology of the bundle $`Y`$, namely,
$`H^{k<n}(d_H;𝒯_{\mathrm{}}^{})=H^{k<n}(Y),H^{kn}(\delta ;𝒯_{\mathrm{}}^{})=H^{kn}(Y).`$
This isomorphism recovers the results of , but note also the following. The relation (7) for $`\tau `$ and the relation $`h_0d=d_Hh_0`$ for $`h_0`$ define a homomorphisms of the de Rham complex (12) of the algebra $`𝒯_{\mathrm{}}^{}`$ to its variational complex (9). The corresponding homomorphism of their cohomology groups is an isomorphism by virtue of Proposition 5 and Proposition 5. Then, the splitting (13) leads to the following decompositions.
PROPOSITION 8. Any $`d_H`$-closed form $`\sigma 𝒯^{0,m}`$, $`m<n`$, is represented by the sum
$$\sigma =h_0\phi +d_H\xi ,\xi 𝒯_{\mathrm{}}^{m1},$$
(14)
where $`\phi `$ is a closed $`m`$-form on $`Y`$. Any $`\delta `$-closed form $`\sigma 𝒯^{k,n}`$, $`k0`$, splits into
$`\sigma =h_0\phi +d_H\xi ,k=0,\xi 𝒯_{\mathrm{}}^{0,n1},`$ (15)
$`\sigma =\tau (\phi )+\delta (\xi ),k=1,\xi 𝒯_{\mathrm{}}^{0,n},`$ (16)
$`\sigma =\tau (\phi )+\delta (\xi ),k>1,\xi E_{k1},`$ (17)
where $`\phi `$ is a closed $`(n+k)`$-form on $`Y`$.
## 6 Cohomology of $`𝒪_{\mathrm{}}^{}`$
THEOREM 9. Graded differential algebra $`𝒪_{\mathrm{}}^{}`$ has the same $`d_H`$\- and $`\delta `$-cohomology of the variational complex as $`𝒯_{\mathrm{}}^{}`$.
Proof. Let the common symbol $`D`$ stand for $`d_H`$ and $`\delta `$. Bearing in mind decompositions (14) – (17), it suffices to show that, if an element $`\varphi 𝒪_{\mathrm{}}^{}`$ is $`D`$-exact in the algebra $`𝒯_{\mathrm{}}^{}`$, then it is so in the algebra $`𝒪_{\mathrm{}}^{}`$. Lemma 4 states that, if $`Y`$ is a contractible bundle and a $`D`$-exact form $`\varphi `$ on $`J^{\mathrm{}}Y`$ is of finite jet order $`[\varphi ]`$ (i.e., $`\varphi 𝒪_{\mathrm{}}^{}`$), there exists an exterior form $`\phi 𝒪_{\mathrm{}}^{}`$ on $`J^{\mathrm{}}Y`$ such that $`\varphi =D\phi `$. Moreover, a glance at the homotopy operators for $`d_H`$ and $`\delta `$ shows that the jet order $`[\phi ]`$ of $`\phi `$ is bounded for all exterior forms $`\varphi `$ of fixed jet order. Let us call this fact the finite exactness of the operator $`D`$. Given an arbitrary bundle $`Y`$, the finite exactness takes place on $`J^{\mathrm{}}Y|_U`$ over any open subset $`U`$ of $`Y`$ which is homeomorphic to a convex open subset of $`R^{\mathrm{dim}Y}`$. Let us prove the following.
(i) Suppose that the finite exactness of the operator $`D`$ takes place on $`J^{\mathrm{}}Y`$ over open subsets $`U`$, $`V`$ of $`Y`$ and their non-empty overlap $`UV`$. Then, it is also true on $`J^{\mathrm{}}Y|_{UV}`$.
(ii) Given a family $`\{U_\alpha \}`$ of disjoint open subsets of $`Y`$, let us suppose that the finite exactness takes place on $`J^{\mathrm{}}Y|_{U_\alpha }`$ over every subset $`U_\alpha `$ from this family. Then, it is true on $`J^{\mathrm{}}Y`$ over the union $`\underset{\alpha }{}U_\alpha `$ of these subsets.
If these assertions hold, the finite exactness of $`D`$ on $`J^{\mathrm{}}Y`$ takes place because one can construct the corresponding covering of the manifold $`Y`$ (, Lemma 9.5).
Proof of (i). Let $`\varphi =D\phi 𝒪_{\mathrm{}}^{}`$ be a $`D`$-exact form on $`J^{\mathrm{}}Y`$. By assumption, it can be brought into the form $`D\phi _U`$ on $`(\pi _0^{\mathrm{}})^1(U)`$ and $`D\phi _V`$ on $`(\pi _0^{\mathrm{}})^1(V)`$, where $`\phi _U`$ and $`\phi _V`$ are exterior forms of finite jet order. Let us consider their difference $`\phi _U\phi _V`$ on $`(\pi _0^{\mathrm{}})^1(UV)`$. It is a $`D`$-exact form of finite jet order which, by assumption, can be written as $`\phi _U\phi _V=D\sigma `$ where $`\sigma `$ is also of finite jet order. Lemma 6 below shows that $`\sigma =\sigma _U+\sigma _V`$ where $`\sigma _U`$ and $`\sigma _V`$ are exterior forms of finite jet order on $`(\pi _0^{\mathrm{}})^1(U)`$ and $`(\pi _0^{\mathrm{}})^1(V)`$, respectively. Then, putting
$`\phi ^{}|_U=\phi _UD\sigma _U,\phi ^{}|_V=\phi _V+D\sigma _V,`$
we have $`\varphi =D\phi ^{}`$ on $`(\pi _0^{\mathrm{}})^1(UV)`$ where $`\phi ^{}`$ is of finite jet order.
Proof of (ii). Let $`\varphi 𝒪_{\mathrm{}}^{}`$ be a $`D`$-exact form on $`J^{\mathrm{}}Y`$. The finite exactness on $`(\pi _0^{\mathrm{}})^1(U_\alpha )`$ holds since $`\varphi =D\phi _\alpha `$ on every $`(\pi _0^{\mathrm{}})^1(U_\alpha )`$ and, as was mentioned above, the jet order $`[\phi _\alpha ]`$ is bounded on the set of exterior forms $`D\phi _\alpha `$ of fixed jet order $`[\varphi ]`$. $`\mathrm{}`$
LEMMA 10. Let $`U`$ and $`V`$ be open subsets of a bundle $`Y`$ and $`\sigma 𝔒_{\mathrm{}}^{}`$ an exterior form of finite jet order on $`(\pi _0^{\mathrm{}})^1(UV)J^{\mathrm{}}Y`$. Then, $`\sigma `$ splits into a sum $`\sigma _U+\sigma _V`$ of exterior forms $`\sigma _U`$ and $`\sigma _V`$ of finite jet order on $`(\pi _0^{\mathrm{}})^1(U)`$ and $`(\pi _0^{\mathrm{}})^1(V)`$, respectively.
Proof. By taking a smooth partition of unity on $`UV`$ subordinate to the cover $`\{U,V\}`$ and passing to the function with support in $`V`$, one gets a smooth real function $`f`$ on $`UV`$ which is 0 on a neighborhood of $`UV`$ and 1 on a neighborhood of $`VU`$ in $`UV`$. Let $`(\pi _0^{\mathrm{}})^{}f`$ be the pull-back of $`f`$ onto $`(\pi _0^{\mathrm{}})^1(UV)`$. The exterior form $`((\pi _0^{\mathrm{}})^{}f)\sigma `$ is zero on a neighborhood of $`(\pi _0^{\mathrm{}})^1(U)`$ and, therefore, can be extended by 0 to $`(\pi _0^{\mathrm{}})^1(U)`$. Let us denote it $`\sigma _U`$. Accordingly, the exterior form $`(1(\pi _0^{\mathrm{}})^{}f)\sigma `$ has an extension $`\sigma _V`$ by 0 to $`(\pi _0^{\mathrm{}})^1(V)`$. Then, $`\sigma =\sigma _U+\sigma _V`$ is a desired decomposition because $`\sigma _U`$ and $`\sigma _V`$ are of finite jet order which does not exceed that of $`\sigma `$. $`\mathrm{}`$
## 7 The global inverse problem
The expressions (15) – (16) in Proposition 5 provide a solution of the global inverse problem of the calculus of variations on fibre bundles in the class of Lagrangians $`L𝒯_{\mathrm{}}^{0,n}`$ of locally finite order (which is not so interesting for physical applications). These expressions together with Theorem 6 give a solution of the global inverse problem of the finite order calculus of variations.
PROPOSITION 11. (i) A finite order Lagrangian $`L𝒪_{\mathrm{}}^{0,n}`$ is variationally trivial, i.e., $`\delta (L)=0`$ iff
$$L=h_0\phi +d_H\xi ,\xi 𝒪_{\mathrm{}}^{0,n1},$$
(18)
where $`\phi `$ is a closed $`n`$-form on $`Y`$. (ii) A finite order Euler–Lagrange-type operator satisfies the Helmholtz condition $`\delta ()=0`$ iff
$`=\delta (L)+\tau (\varphi ),L𝒪_{\mathrm{}}^{0,n},`$
where $`\varphi `$ is a closed $`(n+1)`$-form on $`Y`$ (see also ).
A solution of the global inverse problem of the fixed order calculus of variations has been suggested in by computing cohomology of the fixed order variational sequence. However, the proof of the local exactness of this variational sequence requires rather sophisticated ad hoc techniques in order to be reproduced (see also ). The first thesis of agrees with that of Proposition 7i, but says that the jet order of the form $`\xi `$ in the expression (18) is $`k1`$ if $`L`$ is a $`k`$-order variationally trivial Lagrangian. The second one states that a $`2k`$-order Euler–Lagrange operator can be always associated with a $`k`$-order Lagrangian.
One obtains from Proposition 7(i) that the cohomology group $`H^n(d_H;𝒪_{\mathrm{}}^{})`$ of the complex (5) obeys the relation
$$H^n(d_H;𝒪_{\mathrm{}}^{})/H^n(Y)=\delta (𝒪_{\mathrm{}}^{0,n}),$$
(19)
where $`\delta (𝒪_{\mathrm{}}^{0,n})`$ is the $`R`$-module of Euler–Lagrange forms on $`J^{\mathrm{}}Y`$.
## 8 The case of an affine bundle
Let $`YX`$ be an affine bundle. Since $`X`$ is a strong deformation retract of $`Y`$, de Rham cohomology of $`Y`$ is equal to that of $`X`$. It leads to the cohomology isomorphisms
$`H^{<n}(d_H;𝒪_{\mathrm{}}^{})=H^{<n}(X),H^0(\delta ;𝒪_{\mathrm{}}^{})=H^n(X),H^k(\delta ;𝒪_{\mathrm{}}^{})=0.`$
Hence, every $`d_H`$-closed form $`\varphi 𝒪_{\mathrm{}}^{0,m<n}`$ splits into the sum
$$\varphi =\phi +d_H\xi ,\xi 𝒪_{\mathrm{}}^{0,m1},$$
(20)
where $`\phi `$ is a closed form on $`X`$.
In the case of an affine bundle $`YX`$, horizontal complexes (4) – (5) induce similar complexes on the base $`X`$ as follows. It is quite important for the cohomology calculation of polynomial complexes in next Section.
Let us consider the open surjection $`\pi ^{\mathrm{}}:J^{\mathrm{}}YX`$ and the direct image $`\{\pi _{}^{\mathrm{}}𝔗_{\mathrm{}}^{}\}`$ on $`X`$ of the sheaf $`𝔗_{\mathrm{}}^{}`$. Its stalk at a point $`xX`$ consists of the equivalence classes of sections of the sheaf $`𝔗_{\mathrm{}}^{}`$ which coincide on the inverse images $`(\pi ^{\mathrm{}})^1(U_x)`$ of neighbourhoods $`U_x`$ of $`x`$. Put further the notation $`𝔗𝔛_{\mathrm{}}^{}=\pi _{}^{\mathrm{}}𝔗_{\mathrm{}}^{}`$. Since $`\pi _{}^{\mathrm{}}R=R`$, we have the following complex of sheaves on $`X`$:
$$0R𝔗𝔛_{\mathrm{}}^0\stackrel{d_H}{}𝔗𝔛_{\mathrm{}}^{0,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}𝔗𝔛_{\mathrm{}}^{0,n}\stackrel{d_H}{}0.$$
(21)
Every point $`xX`$ has a base of open contractible neighbourhoods $`\{U_x\}`$ such that the sheaves $`𝔗_{\mathrm{}}^{0,}`$ of $`𝒯_{\mathrm{}}^{}`$-modules are acyclic on the inverse images $`(\pi ^{\mathrm{}})^1(U_x)`$ of these neighbourhoods. Then, in accordance with the Leray theorem , cohomology of $`J^{\mathrm{}}Y`$ with coefficients in the sheaves $`𝔗_{\mathrm{}}^{0,}`$ are isomorphic to that of $`X`$ with coefficients in their direct images $`𝔗𝔛_{\mathrm{}}^{0,}`$, i.e., the sheaves $`𝔗𝔛_{\mathrm{}}^{0,}`$ on $`X`$ are acyclic. Furthermore, Lemma 4 shows that the complexes of sections of sheaves $`𝔗_{\mathrm{}}^{0,<n}`$ over $`(\pi _0^{\mathrm{}})^1(U_x)`$ are exact. It follows that the horizontal complex (21), except the last term, is also exact. Due to the $`R`$-algebra isomorphism $`𝒯_{\mathrm{}}^{}=\mathrm{\Gamma }(X,𝔗𝔛_{\mathrm{}}^{})`$, one can think of the horizontal complex (4) as being the complex of the structure algebras of sheaves of the horizontal complex (21) on $`X`$.
## 9 Cohomology of polynomial complexes
Given the sheaf $`𝔗𝔛_{\mathrm{}}^{}`$ on $`X`$, let us consider its subsheaf $`𝔓_{\mathrm{}}^{}`$ of germs of exterior forms which are polynomials in the fiber coordinates $`y_\mathrm{\Lambda }^i`$, $`|\mathrm{\Lambda }|0`$, of the topological fiber bundle $`J^{\mathrm{}}YX`$. This property is coordinate-independent due to the transition functions (2). The $`𝔓_{\mathrm{}}^{}`$ is a sheaf of $`C^{\mathrm{}}(X)`$-modules. Its structure algebra $`𝒫_{\mathrm{}}^{}`$ is a $`C^{\mathrm{}}(X)`$-subalgebra of $`𝒯_{\mathrm{}}^{}`$. For short, one can say that $`𝒫_{\mathrm{}}^{}`$ consists of exterior forms on $`J^{\mathrm{}}Y`$ which are locally polynomials in fiber coordinates $`y_\mathrm{\Lambda }^i`$.
We have the subcomplex
$$0R𝔓_{\mathrm{}}^0\stackrel{d_H}{}𝔓_{\mathrm{}}^{0,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}𝔓_{\mathrm{}}^{0,n}$$
(22)
of the horizontal complex (21) on $`X`$. As a particular variant of the algebraic Poincaré lemma, the exactness of the complex (22) has been repeatedly proved (see, e.g., ). Since the sheaves $`𝔓_{\mathrm{}}^{0,}`$ of $`C^{\mathrm{}}(X)`$-modules on $`X`$ are acyclic, the complex (22) is a resolution of the constant sheaf $`R`$ on $`X`$. Hence, cohomology of the complex
$$0R𝒫_{\mathrm{}}^0\stackrel{d_H}{}𝒫_{\mathrm{}}^{0,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}𝒫_{\mathrm{}}^{0,n}$$
(23)
of the structure algebras $`𝒫_{\mathrm{}}^{0,<n}`$ of sheaves $`𝔓_{\mathrm{}}^{0,<n}`$ is equal to de Rham cohomology of $`X`$. It follows that every $`d_H`$-closed polynomial form $`\varphi 𝒫_{\mathrm{}}^{0,m<n}`$ splits into the sum
$$\varphi =\phi +d_H\xi ,\xi 𝒫_{\mathrm{}}^{0,m1},$$
(24)
where $`\phi `$ is a closed form on $`X`$. Let $`P_{\mathrm{}}^{}`$ be $`C^{\mathrm{}}(X)`$-subalgebra of the polynomial algebra $`𝒫_{\mathrm{}}^{}`$ which consists of exterior forms which are polynomials in the fiber coordinates $`y_\mathrm{\Lambda }^i`$. Obviously, $`P_{\mathrm{}}^{}`$ is a subalgebra of $`𝒪_{\mathrm{}}^{}`$. As a repetition of Theorem 6, one can show that $`P_{\mathrm{}}^{}`$ have the same cohomology as $`𝒫_{\mathrm{}}^{}`$, i.e., if $`\varphi `$ in the decomposition (24) is an element of $`P_{\mathrm{}}^{0,}`$, then $`\xi `$ is so.
Let us consider the subsheaf $`\overline{𝔓}_{\mathrm{}}^{}`$ of the sheaf $`𝔓_{\mathrm{}}^{}`$ which consists of germs of $`x`$-independent polynomial forms. Its structure algebra $`\overline{P}_{\mathrm{}}^{}`$ is a subalgebra of the algebra $`P_{\mathrm{}}^{}`$. We have the complex of sheaves
$`0R\overline{𝔓}_{\mathrm{}}^0\stackrel{d_H}{}\overline{𝔓}_{\mathrm{}}^{0,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}\overline{𝔓}_{\mathrm{}}^{0,n}`$
which fails to be exact. The obstruction to its exactness at the term $`\overline{𝔓}_{\mathrm{}}^{0,k}`$ is provided by the germs of $`k`$-forms on $`X`$ with constant coefficients . Let us denote the sheaf of such germs by $`S_X^k`$. For any $`0<k<n`$, we have the short exact sequence of sheaves
$`0\mathrm{Im}d_H\mathrm{Ker}d_HS_X^k0`$
and the sequence of their structure modules
$`0\mathrm{\Gamma }(X,\mathrm{Im}d_H)\mathrm{\Gamma }(X,\mathrm{Ker}d_H)\mathrm{\Gamma }(X,S_X^k)0`$
which is exact because $`S_X^k`$ is a subsheaf of $`R`$-modules of the sheaf $`\mathrm{Ker}d_H`$. Then, the $`k`$th cohomology group of the horizontal complex
$`0R\overline{P}_{\mathrm{}}^0\stackrel{d_H}{}\overline{P}_{\mathrm{}}^{0,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}\overline{P}_{\mathrm{}}^{0,n}`$
of the algebra $`\overline{P}_{\mathrm{}}^{}`$ is isomorphic to the $`R`$-module $`\mathrm{\Gamma }(X,S_X^k)`$ of global constant $`k`$-forms on the manifold $`X`$. If a manifold $`X`$ does not admit an affine coordinate atlas, the module $`\mathrm{\Gamma }(X,S_{\mathrm{}}^{0,k})`$ is empty and, consequently, the differential $`d_H`$ is exact on the algebra $`\overline{𝒫}^{0,<n}`$. Otherwise, any $`d_H`$-closed element $`\varphi \overline{P}_{\mathrm{}}^{0,k}`$, $`0<k<n`$, splits into the sum $`\varphi =\phi +d_H\xi `$, $`\phi \mathrm{\Gamma }(X,S_X^k)`$, $`\xi \overline{P}_{\mathrm{}}^{0,k1}`$.
## 10 Iterated cohomology
Turn to the proof of Theorem 1. By assumption, the horizontal complex (6) is a $`𝐬|d_H`$-bicomplex. Its iterated cohomology group $`H^{k,m}(𝐬|d_H)=E_2^{k,m}`$, $`kZ`$, $`0mn`$, consists of $`d_H`$-closed horizontal $`m`$-forms $`\omega 𝒪^{0,m}`$ of ghost number $`k`$ such that $`𝐬\omega `$ is $`d_H`$-exact, which are taken modulo exterior forms $`𝐬\psi +d_H\xi `$ where $`\psi `$ is a $`d_H`$-closed form. Then, the assertions (i) and (ii) of Theorem 1 follows immediately from the decomposition (20) and the assumption that forms on $`X`$ are of zero ghost number and are not $`𝐬`$-exact. The proof of assertions (iii) – (v) is based on the analysis of descent equations. Since the operator $`𝐬`$ annihilates exterior forms on $`X`$, descent equations can be constructed, but their right-hand side is not necessarily zero. The key point lies in the existence of closed forms on $`X`$ which are exact with respect to the total operator $`\stackrel{~}{𝐬}=𝐬+d_H`$.
Proof of (iii). Let $`\omega _n`$ be a representative of the iterated cohomology group $`H^{k,n}(𝐬|d_H)`$, $`k<n`$ or $`k1`$. It is a horizontal $`n`$-form of ghost number $`k`$ such that $`𝐬\omega _n`$ is $`d_H`$-exact, i.e.,
$$𝐬\omega _n+d_H\omega _{n1}=0.$$
(25)
Acting on this equality by $`𝐬`$, we observe that $`𝐬\omega _{n1}`$ is a $`d_H`$-closed form of non-vanishing ghost number $`k+2`$. Therefore, it is $`d_H`$-exact, i.e.,
$`𝐬\omega _{n1}+d_H\omega _{n2}=0.`$
Iterating the arguments, one concludes the existence of a family $`\{\omega _{np}\}`$, $`0pn`$, of horizontal $`(np)`$-forms $`\omega _{np}`$ of non-vanishing ghost numbers $`k+p`$ which obey the descent equations
$$𝐬\omega _{np}+d_H\omega _{np1}=0,𝐬\omega _0=0,0p<n.$$
(26)
It may happen that $`\omega _{np}=0`$, $`pp_0`$, for some $`p_0`$. Put
$$\stackrel{~}{\omega }_n=\underset{p=0}{\overset{n}{}}\omega _{np}.$$
(27)
It is an $`\stackrel{~}{𝐬}`$-closed form of total ghost number $`k+n`$. Let $`\{\omega _{np}^{}\}`$ be another solution of the descent equations (26) for given $`\omega _n`$. It is readily observed that $`\stackrel{~}{\omega }_n\stackrel{~}{\omega }_n^{}`$ is an $`\stackrel{~}{𝐬}`$-exact form. Let the iterated cohomology class of $`\omega _n`$ be zero, i.e. $`\omega _n=𝐬\xi _n+d_H\xi _{n1}`$ where $`\xi _n`$ is $`d_H`$-closed. Then, $`\{\omega _n,𝐬\xi _{n1},0,\mathrm{},0\}`$ is a solution of the descent equations such that $`\stackrel{~}{\omega }_n`$ (27) is an $`\stackrel{~}{𝐬}`$-exact form. Conversely, let a horizontal exterior form $`\stackrel{~}{\omega }`$ of total ghost number $`n+k`$ be $`\stackrel{~}{𝐬}`$-closed. It splits into the sum (27) whose summands obey the descent equations (26). The higher term $`\omega _n`$ of this sum fulfills the relation (25), i.e., is a representative of iterated cohomology. Since all $`d_H`$-closed forms of non-vanishing ghost number are $`d_H`$-exact, the descent equations (26) show that: (i) if $`\omega _n=0`$, then $`\stackrel{~}{\omega }`$ is $`\stackrel{~}{𝐬}`$-exact, and (ii) if $`\stackrel{~}{\omega }=\stackrel{~}{𝐬}\xi `$ is $`\stackrel{~}{𝐬}`$-exact, then $`\omega _n=𝐬\xi _n+d_H\xi _{n1}`$ is of zero cohomology class. Thus, we come to a desired isomorphism.
Proof of (iv). In contrast with the previous case, $`d_H`$-closed forms now are not necessarily $`d_H`$-exact. Therefore, a representative $`\omega _n`$ of iterated cohomology defines a system of descent equations where the descent equation of vanishing ghost number
$$𝐬\omega _{n+k+1}+d_H\omega _{n+k}=\phi $$
(28)
has a closed $`(n+k+1)`$-form $`\phi `$ on $`X`$ in its right-hand side. Accordingly, the form $`\stackrel{~}{\omega }_n`$ (27) fulfills the equality $`\stackrel{~}{𝐬}\stackrel{~}{\omega }_n=\phi `$. It follows that the system of descent equations
$$\stackrel{~}{𝐬}\stackrel{~}{\omega }^\phi =\phi ,$$
(29)
which contains the equation (28), admits a solution $`\stackrel{~}{\omega }^\phi `$ iff the cohomology class of $`\phi `$ belongs to $`\mathrm{Ker}\gamma _{n+k+1}`$. A higher term $`\omega _n^\phi `$ of every such solution obeys the relation (25) and, consequently, is a representative of iterated cohomology. If $`\omega _n^\phi =\omega _n^\phi +𝐬\xi _n+d_H\xi _{n1}`$ is another representative of the same cohomology class, then $`\stackrel{~}{\omega }+𝐬\xi _{n1}`$ is also a solution of the same descent equations (29). Consequently, any closed $`(n+k+1)`$-form $`\phi `$ on $`X`$ whose cohomology class belongs to $`\mathrm{Ker}\gamma _{n+k+1}`$ defines a subset $`A_\phi `$ of the iterated cohomology group $`H^{k,n}(𝐬|d_H)`$, given by the higher terms of solutions of the descent equations (29). In particular, let us consider $`A_{\phi =0}`$. The difference from the proof of item (iii) lies in the fact that, if the higher term $`\omega _n`$ of an $`\stackrel{~}{𝐬}`$-closed form $`\stackrel{~}{\omega }`$ vanishes, then $`\stackrel{~}{\omega }=\stackrel{~}{𝐬}\xi +\psi `$ where $`\psi `$ is a $`(n+k)`$-closed form on $`X`$. It follows that $`A_{\phi =0}=\overline{H}^{k+n}`$. One can justify easily that: (i) $`\omega _n^\phi A_\phi `$ and $`\omega _n^0A_{\phi =0}`$ implies $`\omega _n^\phi +\omega _n^0A_\phi `$, and (ii) $`\omega _n^\phi ,\omega _n^\phi A_\phi `$ implies $`\omega _n^\phi \omega _n^\phi A_{\phi =0}`$. It follows that $`A_\phi `$ is an affine space modelled over the linear space $`\overline{H}^{k+n}`$. Let $`\phi ^{}=\phi +d\sigma `$. Then, a solution $`\stackrel{~}{\omega }^\phi `$ of the descent equations (29) defines a solution $`\stackrel{~}{\omega }^\phi +\sigma `$ of the descent equations (29) where $`\phi `$ is replaced with $`\phi ^{}`$. These solutions have the same higher term $`\omega _n^\phi =\omega _n^\phi ^{}`$ and, consequently define the same representative of iterated cohomology. It follows that $`A_\phi =A_{\phi +d_H\sigma }`$, i.e., $`A_\phi `$ is set by the cohomology class $`[\phi ]\mathrm{Ker}\gamma _{k+n+1}`$ of $`\phi `$. It remains to show that $`H^{k,n}(𝐬|d_H)`$ is a disjoint union of sets $`A_{[\phi ]}`$, $`[\phi ]\mathrm{Ker}\gamma _{k+n+1}`$. Indeed, let a representative of iterated cohomology defines different systems of descent equations (29) with $`\phi `$ and $`\phi ^{}`$ in the right-hand side. Then, one can easily justify that $`\phi ^{}\phi `$ is an exact form. Turn to the particular case $`k=n`$. If a constant function on $`X`$ is $`\stackrel{~}{𝐬}`$-exact, it is $`𝐬`$-exact. Therefore, $`\mathrm{Ker}\gamma =0`$.
Proof of (v). In this case, we have the descent equations (26) with the zero right-hand side, but $`H^{1,n}(𝐬|d_H)=H_{\mathrm{tot}}^1/\mathrm{Im}\gamma _{n1}`$.
Note that, in the case of BRST cohomology (see Corollary 1), the right-hand side of the global descent equations for any total ghost number remains zero.
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# CONSTRAINTS ON CLUSTER PROPERTIES USING X-RAY & OPTICAL DATA
## 1 Introduction
Our main purpose when we started this work was to perform realistic simulations of the Sunyaev-Zel’dovich effect (SZE) in order to check different substraction techniques applied to this effect. By realistic we understand that our simulations must take into account first, the cluster population; that is, how is the cluster distribution in the mass-redshift space. Second, we must be able to describe cluster properties such as the scalings temperature-mass ($`TM`$) and X-ray luminosity-mass ($`L_xM`$). These descriptions of the cluster population and cluster scalings must be realistic in the sense that both descriptions must be in agreement with recent cluster data.
The idea of this work is to find such a realistic description of the clusters (population + intrinsic properties) by comparing our model with recent X-ray and optical data and fitting the free parameters of the model to the data in order to search for the best fitting values.
Several authors have used different cluster data sets in an attempt to constrain the cosmology. The usuall procedure is, starting from the Press-Schechter (PS) mass function, fitting the experimental mass function or by using a given (that is fixed or non-free parameter) $`TM`$ or $`L_xM`$ relation construct some of the other cluster functions (temperature, X-ray luminosity or flux functions) and then compare with the corresponding data sets. This can be a dangerous process for two reasons. First, when considering just one data set one ensures that his best fitting model is compatible with that data set but it can be inconsistent with others. We have observed this behaviour in many cases. For instance some of the models inside the $`68\%`$ confidence level region in the $`\mathrm{\Omega }\sigma _8`$ space in Bahcall & Fan (1998) are inconsistent with, for example, the luminosity function of Ebeling et al. (1997). This point suggests that a consistent analysis should take into account the information coming from different experiments to avoid these incompatibilities. The second problem comes when some authors fix the $`TM`$ or $`L_xM`$ relations. The scatter in this correlations is large enough to introduce important uncertainties in the final result (Voit & Donahue (1998)). When these authors fix some of these relations, they probably are introducing a sistematic bias in their conclusions due to the fact that they suppose that the virial relation, for instance, is a good enough aproximation everywhere. We want to avoid all those uncertainties by fitting different data sets simultaneously without doing any assumption about the cosmology and the $`TM`$ and $`L_xM`$ relations.
## 2 Data
In this section we present the five different data sets we have used in our fit. The first data set is the cluster mass function given in Bahcall & Cen. (1993). To take into account evolutionary effects of the cluster mass function we use the evolution of the mass function given by Bahcall & Fan 1998. With the $`TM`$ and $`L_xM`$ relations we can build the cluster temperature, X-ray luminosity and flux functions. For the temperature function we use the one given by Henry & Arnaud (1991). For the X-ray luminosity function, that of Ebeling et al. 1997 and finally for the X-ray cluster flux function, the one determined by Rosati et al. 1998 for low-flux clusters, and the one obtained by De Grandi et al. (1999) for high-flux clusters. With all these data curves and with our model we are now able to fit the model and say something about its free parameters. The results will be robust and consistent in the sense that we did not make any assumption about the cosmology (we did not fix any cosmological parameter) nor about the intrinsic cluster properties ($`TM`$ and $`L_xM`$).
## 3 The model
Our model consists of two parts. First the description of the cluster population and second the intrinsic cluster properties. For the cluster population we assume the standard Press-Schechter (PS) formalism (Press & Schechter 1974). This formalism depends on three parameters: the density of the universe $`\mathrm{\Omega }`$, the amplitude of the power spectrum in units of $`\sigma _8`$ and finally the shape parameter of the power spectrum $`\mathrm{\Gamma }`$. In this work we have only considered low density models with $`\mathrm{\Lambda }=0`$. In a subsequent paper we will also include in our analisis flat $`\mathrm{\Lambda }`$CDM models.
The PS approach is supported by N-body numerical simulations which do show a good agreement with the PS parametrization (Lacey & Cole 1994, White et al. 1993, Efstathiou et al. 1988, etc).
With the PS formalism we know the comoving number density of clusters with $`M[M,M+dM]`$ and at a given redshift. This will allow us to distribute the clusters in the $`Mz`$ space for a given solid angle and cosmology ($`\mathrm{\Omega },\sigma _8,\mathrm{\Gamma }`$). What we need now is the second part of the model, the intrinsic cluster properties. Basically what we need is the $`TM`$ relation for which the virial relation is usually assumed (Navarro et al. 1995, Bryan & Norman 1998). One problem is that it is not clear to what extend the virial relation is true for high or even intermediate redshift. In this work we have decided to consider this scaling as a free parameter relation:
$$T_{gas}=T_0M_{15}^\alpha (1+z)^\psi $$
(1)
where $`T_0,\alpha `$ and $`\psi `$ are our three free parameters. $`M_{15}`$ is the cluster mass in $`h^110^{15}M_{}`$ units. With the $`TM`$ relation and the PS mass function we can build the temperature function, simply by doing:
$$\frac{dN(T,z)}{dT}=\frac{dN(M,z)}{dM}\frac{dM}{dT}$$
(2)
Our aim is to use as many data sets as possible and so, we also want to include in our analysis other relevant data sets like the X-ray luminosity and flux functions. To do that we need another scaling law, the $`L_xM`$ relation in order to relate the mass function with the X-ray luminosity and flux functions similarly as we did with the temperature function in eq. (2). As in the $`TM`$ relation, the $`L_xM`$ is not well established yet and for this reson we have adopted a parametrization similar to that given in eq. (1):
$$L_x^{Bol}=L_0M_{15}^\beta (1+z)^\varphi $$
(3)
where we added three additional free parameters: $`L_0,\beta `$ and $`\varphi `$. From this relation is possible to build the $`S_xM`$ relation. Just taking,
$$S_x^{Bol}=\frac{L_x^{Bol}}{4\pi D_l(z)^2}$$
(4)
Summarising, the final number of free parameters are 9 : $`\mathrm{\Omega },\sigma _8,\mathrm{\Gamma },T_0,\alpha ,\psi ,L_0,\beta `$ and $`\varphi `$.
We now want to play with these parameters and look for the best fitting model to the different data sets.
## 4 Best fit
To fit the five data sets we must decide which estimator will pick up our best fitting model. Because of the scaling relations assumed between the mass and temperature ($`TM`$) and the mass and luminosity in the X-ray band ($`L_xM`$), then there must be some correlations among the five data sets predicted by the model. Therefore we should start by considering an estimator like the standard likelihood estimator.
In our case, the model depends on 9 free parameters and if we consider a grid of, let’s say 5 values per parameter, then we should compute the correlation matrix for $`5^9`$ 1 million different models. This process would take many years. A faster technique would require a search method that avoids exploring all the parameter space. This can be the solution if we are interested just in the best model but we want also the error bars of our model, or in other words the marginalized probability distribution of the parameters. In order to do that, we need to know the probability in a given regular grid.
To simplify the problem, the most simple approach is to consider as our estimator the standard $`\chi _{joint}^2`$:
$$\chi _{joint}^2=\chi _M^2+\chi _{M(z)}^2++\chi _T^2++\chi _{L_x}^2+\chi _{S_x}^2$$
(5)
where $`\chi _x^2`$ represents the corresponding ordinary $`\chi ^2`$ for the five different data sets and we are supposing that the correlation matrix is in this case diagonal.
By doing this, we know that we are forgetting the correlations between the curves and that there will be some bias in our estimation. For this reason we want to check other more ellaborated estimators.
We have considered as a second estimator of the best model the next one based on bayesian theory (Lahav et al, 2000);
$$2lnP_L=\chi _L^2$$
(6)
where,
$$\chi _L^2=\underset{i}{\overset{5}{}}N_iln(\chi _i^2)$$
(7)
In this estimator $`\chi ^2`$ is the same as before and $`N_i`$ represents the number of data points for the data set $`i`$. The authors have shown that this estimator is apropiate for the case when different data sets are combined together, as is our case. The factor $`N_i`$ plays the role of a weight factor. Those data sets with more measures (more data points) are considered as more realistic.
We have checked both estimators by doing a bias test. In this test we have simulated the data for a known model with the corresponding error bars. The input model was selected according to the criterium that it would be as close as possible to the data (for instance the model which minimises $`\chi _{joint}^2`$). In the simulations we have taken into account all the characteristics of the data, that is, sky coverage, limiting flux, maximun redshift etc. Then we compare each one of these realizations with the mean value of the different models and for each realization we get the best model using both estimators.
We have concluded from this test that the second estimator works better than the standard $`\chi _{joint}^2`$. There is still some bias with the second estimator but the agreement between the input model and the recovered one is very good.
Using the second estimator we have computed the probability distribution in our 9 parameter space. We have used a grid with about 1 million different models and for each of them we have computed its $`P_L`$ (eq. 6). Knowing that, we can obtain the best model (maximum probability, see table 1).
## 5 Conclusions
In this work we have shown that combining different data sets and using the Lahav’s et al. estimator is a powerful tool to constrain the cosmology and cluster parameters. In a future paper (Diego et al, in preparation) we will present the full analysis taking into account both, open ($`\mathrm{\Lambda }=0`$) and flat ($`\mathrm{\Lambda }>0`$) CDM models. The marginalized probability distributions for the parameters of the models will be also included.
Additional data coming from high redshift clusters (CHANDRA, XMM-Newton, PLANCK) will improve this result.
Particularly interesting is the work that can be done with future CMB surveys. The PLANCK satellite will explore the whole sky at nine different frecuencies (from 30 Ghz to 800 Ghz) and with resolutions between 5 arcmin and 30 arcmin. At these frecuencies and with those resolutions we have shown (Diego et al. in preparation) that several clusters are expected to be observed at high redshift ($`z>2`$) through the Sunyaev-Zel’dovich effect. The information that these clusters will provide will be decisive to definitely exclude many models.
## Acknowledgments
We would like to thank to Piero Rosati for kindly providing his data for the differential flux function. JMD acknowledges the DGES for a fellowship. JMD, EM, JLS, & LC thanks CFPA/Astronomy Dept. Berkeley for the facilities given during this work.
## References
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# The Curvature of a Single Contraction Operator on a Hilbert Space AMS Subject Classification: 47A13 (Primary); 47A20 (Secondary). Keywords: operator, curvature
## 1 The curvature of a single operator
This note studies Arveson’s curvature invariant for $`d`$-contractions $`T=(T_1,T_2,\mathrm{},T_d)`$ for the special case $`d=1`$, referring to a single contraction operator $`T`$ on a Hilbert space. It establishes a formula which gives an easy-to-understand meaning for the curvature of a single contraction. The formula is applied to give an example of an operator with nonintegral curvature. Under the additional hypothesis that the single contraction $`T`$ be “pure”, we show that its curvature $`K(T)`$ (defined below) is given by $`K(T)=\text{index }(T):=(dim\text{ker }(T)dim\text{coker }(T))`$.
Let $`T`$ be a contraction operator on a Hilbert space $`H`$, and $`\mathrm{\Delta }_T:=\sqrt{1TT^{}}`$. Assume that $`\mathrm{\Delta }_T`$ has finite rank. Then the curvature $`K(T)`$ of $`T`$ (our shorthand for what should properly be called the curvature of the Hilbert module associated with $`T`$) is defined in as:
$$K(T):=_{|z|=1}𝑑z\underset{r1}{lim}(1r^2)\mathrm{tr}(\mathrm{\Delta }_T(1rzT^{})^1(1r\overline{z}T)^1\mathrm{\Delta }_T).$$
(1)
This is a specialization to the case of a single operator of Arveson’s more general theory of $`d`$-contractions, which are finite sets of $`d`$ commuting operators satisfying an auxiliary condition analogous to contractiveness of our $`T`$.
We refer the reader to , , and for the definition and basic properties of $`d`$-contractions. However, we consider $`d`$-contractions for $`d>1`$ solely for purposes of placing our results within the framework of the more general theory, and essentially no knowledge of $`d`$-contractions is necessary to follow our proofs. The only reliance on the general theory is that Arveson’s Stability of Curvature result, , Section 3, Corollary 1, is used in the proof of Proposition 1. However, as noted there, the reader can easily establish this result directly for the special case of a 1-contraction, which is all that we need.
The definition of curvature implicitly assumes the existence of the limit in (1). A theorem stated in , and proved in (Theorem A), guarantees the existence of the limit for almost all $`z`$, and moreover bounds it above by the rank of $`\mathrm{\Delta }_T`$. For the case of a single operator, this also follows from the discussion of , Chapter VI, Section 1, particularly, page 238, equation (1.5).
Let $`T:HH`$ be a contraction on a Hilbert space $`H`$ with $`\text{rank }\sqrt{1TT^{}}`$ finite. Note that this implies that $`\text{range }\sqrt{1TT^{}}=\text{range }(1TT^{})`$, a fact which will be used frequently without comment.
First we associate with $`T`$ a partial isometry $`Q`$ with the same curvature, so that for most purposes of computing the curvature, we may assume that $`T`$ is itself a partial isometry. This is not always necessary, but it makes many problems easier to think about.
###### Proposition 1
With $`T`$ as just described, set
$$Q:=\left[\begin{array}{cc}T& \sqrt{1TT^{}}\\ 0& 0\end{array}\right],$$
considered as an operator on $`H\text{range }(1TT^{})`$.
Then $`Q`$ is a partial isometry with $`K(Q)=K(T)`$.
Moreover, $`\text{rank }(1QQ^{})=\text{rank }(1TT^{})`$,
and, $`\text{rank }(1Q^{}Q)=\text{rank }(1T^{}T)`$.
Proof: That $`K(Q)=K(T)`$ follows from one of Arveson’s key results for $`d`$-contractions, Stability of Curvature, , Section 3, Corollary 1. For our case of a 1-contraction, a proof can alternatively be obtained by a straightforward calculation of $`K(Q)`$, based on its definition (1).
Since
$$1QQ^{}=\left[\begin{array}{cc}0\hfill & 0\hfill \\ 0\hfill & 1_{\text{Range }(1TT^{})}\hfill \end{array}\right],$$
(2)
it is obvious that $`Q`$ is a partial isometry with $`\text{rank }(1QQ^{})=\text{rank }(1TT^{})`$.
Next we show that $`\text{rank }(1Q^{}Q)=\text{rank }(1T^{}T)`$. We have
$$1Q^{}Q=\left[\begin{array}{cc}1T^{}T& T^{}\sqrt{1TT^{}}\\ \sqrt{1TT^{}}T& TT^{}\end{array}\right].$$
Since the rank of an operator matrix is at least as large as the rank of any entry, if $`\text{rank }(1T^{}T)`$ is infinite, so is $`\text{rank }(1Q^{}Q)`$. Thus we may assume that $`\text{rank }(1T^{}T)`$ is finite.
Let $`C_i,i=1,2`$, denote the $`i`$’th column of the matrix for $`1Q^{}Q`$, considered in the obvious way as operators, e.g., $`C_1:H\text{range }(1TT^{})`$. Then
$$C_2\sqrt{1TT^{}}=\left[\begin{array}{c}T^{}(1TT^{})\\ TT^{}\sqrt{1TT^{}}\end{array}\right]=\left[\begin{array}{c}(1T^{}T)T^{}\\ \sqrt{1TT^{}}TT^{}\end{array}\right]=C_1T^{}$$
Since the domain of $`C_2`$ is $`\text{range }\sqrt{1TT^{}}`$, this implies that $`\text{range }C_2\text{range }C_1`$, and hence $`\text{range }(1Q^{}Q)=\text{range }C_1`$.
It is well known (e.g., , Section 147) that $`\sqrt{1TT^{}}T=T\sqrt{1T^{}T}`$, so
$$C_1=\left[\begin{array}{c}(1T^{}T)\\ T\sqrt{1T^{}T}\end{array}\right]=\left[\begin{array}{c}\sqrt{1T^{}T}\\ T\end{array}\right]\sqrt{1T^{}T}$$
Since
$$\left[\begin{array}{c}\sqrt{1T^{}T}\\ T\end{array}\right]$$
is an isometry, the map $`\sqrt{1T^{}T}xC_1x,xH`$, defines an isometric bijection between $`\text{range }\sqrt{1T^{}T}`$ and $`\text{range }C_1=\text{range }(1Q^{}Q)`$. Hence $`\text{rank }(1T^{}T)=\text{rank }(1Q^{}Q)`$.
Next we derive a simple formula for $`K(Q)`$, along with a variant formula for $`K(T)`$ which does not mention $`Q`$. The formula for $`K(Q)`$ seems particularly helpful in thinking about these problems.
###### Theorem 2
Let $`Q`$ be a partial isometry such that $`\mathrm{\Delta }_Q:=\sqrt{1QQ^{}}`$ has finite rank, and let $`e_1,e_2,\mathrm{},e_q`$ be an orthonormal basis for $`\text{range }(\mathrm{\Delta }_Q)`$. Then
$$K(Q)=\underset{k=1}{\overset{q}{}}\underset{n\mathrm{}}{lim}Q^ne_k^2.$$
Moreover, for any contraction $`T`$ for which $`\mathrm{\Delta }_T`$ has finite rank,
$$K(T)=\underset{n\mathrm{}}{lim}\mathrm{tr}(T_{}^{}{}_{}{}^{n}T^n(1TT^{})).$$
Proof:
The boundedness of the integrand of the curvature justifies application of the Lebesgue Dominated Convergence Theorem to interchange limit and integral in the definition (1) of curvature:
$`K(Q)`$ $`:=`$ $`{\displaystyle _{|z|=1}}\underset{r1}{lim}(1r^2){\displaystyle \underset{k=1}{\overset{q}{}}}(1rzQ^{})^1e_k,(1r\overline{z}Q)^1e_kdz`$ (3)
$`=`$ $`{\displaystyle \underset{k=1}{\overset{q}{}}}\underset{r1}{lim}(1r^2){\displaystyle _{|z|=1}}(1r\overline{z}Q)^1e_k,(1r\overline{z}Q)^1e_k𝑑z`$ (4)
$`=`$ $`{\displaystyle \underset{k=1}{\overset{q}{}}}\underset{r1}{lim}(1r^2){\displaystyle _{|z|=1}}{\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}(r\overline{z}Q)^ie_k,{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}(r\overline{z}Q)^je_k𝑑z`$ (5)
$`=`$ $`{\displaystyle \underset{k=1}{\overset{q}{}}}\underset{r1}{lim}(1r^2){\displaystyle _{|z|=1}}{\displaystyle \underset{i,j=0}{\overset{\mathrm{}}{}}}r^{i+j}z^{ji}Q^ie_k,Q^je_kdz`$ (6)
$`=`$ $`{\displaystyle \underset{k=1}{\overset{q}{}}}\underset{r1}{lim}(1r^2){\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}r^{2i}Q^ie_k^2`$ (7)
$`=`$ $`{\displaystyle \underset{k=1}{\overset{q}{}}}\underset{i\mathrm{}}{lim}Q^ie_k^2`$ (8)
Equation (7) was obtained by interchanging the infinite sum and integration. This is justified because for fixed $`r`$, the infinite sum converges absolutely with sum of absolute values bounded above by $`(1r^2)^2`$.
Equation (8) is justified as follows. For fixed $`k`$, consider the decreasing sequence
$$1Qe_kQ^2e_k\mathrm{}\underset{i\mathrm{}}{lim}Q^ie_k,$$
and set $`L:=lim_i\mathrm{}Q^ie_k^2`$. Then for any positive integer $`m`$,
$`L`$ $`=`$ $`L\underset{r1}{lim}(1r^2){\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}r^{2i}`$
$``$ $`\underset{r1}{lim}(1r^2){\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}r^{2i}Q^ie_k^2`$
$`=`$ $`\underset{r1}{lim}(1r^2){\displaystyle \underset{i=m}{\overset{\mathrm{}}{}}}r^{2i}Q^ie_k^2`$
$``$ $`Q^me_k^2\underset{r1}{lim}(1r^2){\displaystyle \underset{i=m}{\overset{\mathrm{}}{}}}r^{2i}`$
$`=`$ $`Q^me_k^2.`$
For sufficiently large $`m`$, the right side is arbitrarily close to $`L`$, showing that
$$\underset{r1}{lim}(1r^2)\underset{i=0}{\overset{\mathrm{}}{}}r^{2i}Q^ie_k^2=\underset{i\mathrm{}}{lim}Q^ie_k^2,$$
thus proving (8).
This proves the asserted formula for $`K(Q)`$. To prove the alternative formula for $`K(T)`$, define $`Q`$ to be the partial isometry of Proposition 1 with $`K(Q)=K(T)`$. Recall from (2) that
$$\mathrm{\Delta }_Q=\left[\begin{array}{cc}0& 0\\ 0& 1\end{array}\right],$$
and check that
$$Q^n=\left[\begin{array}{cc}T^n& T^{n1}\mathrm{\Delta }_T\\ 0& 0\end{array}\right].$$
(9)
The formula for $`K(T)`$ follows immediately upon combining these observations, Proposition 1, the formula just proved for $`K(Q)`$, and the cyclic property of the trace:
$`K(T)`$ $`=`$ $`K(Q)={\displaystyle \underset{k=1}{\overset{q}{}}}\underset{n\mathrm{}}{lim}Q^ne_k^2`$
$`=`$ $`{\displaystyle \underset{k=1}{\overset{q}{}}}\underset{n\mathrm{}}{lim}\mathrm{\Delta }_TT_{}^{}{}_{}{}^{n}T^n\mathrm{\Delta }_Te_k,e_k`$
$`=`$ $`\underset{n\mathrm{}}{lim}\mathrm{tr}(\mathrm{\Delta }_TT_{}^{}{}_{}{}^{n}T^n\mathrm{\Delta }_T)`$
$`=`$ $`\underset{n\mathrm{}}{lim}\mathrm{tr}(T_{}^{}{}_{}{}^{n}T^n(1TT^{}))`$
A simple sufficient condition for the curvature to vanish is an immediate corollary:
###### Corollary 3
Any contraction $`T`$ whose positive powers $`T^n`$ converge strongly to 0 has vanishing curvature: $`K(T)=0`$.
## 2 Relation to Arveson’s curvature formula
Arveson established a different formula for the curvature of a $`d`$-contraction $`T`$. Specialized to the case $`d=1`$, it reads:
$$K(T)=\underset{n\mathrm{}}{lim}\frac{\mathrm{tr}(1T^nT_{}^{}{}_{}{}^{n})}{n}.$$
(10)
In order to make clear how our formula fits into Arveson’s framework, we now derive ours assuming his. However, the resulting proof is not notably simpler than the direct proof above, and Arveson’s proof is even more involved, corresponding to the fact that the case $`d>1`$ is probably fundamentally more difficult than $`d=1`$. For the single operator case $`d=1`$, Arveson’s formula follows similarly from ours.
Let $`T`$ be a contraction with $`\mathrm{\Delta }_T`$ of finite rank, and $`e_1,\mathrm{},e_q`$ an orthonormal basis for $`\text{range }(\mathrm{\Delta }_T)=\text{range }(1TT^{})`$. First note the collapsing sum:
$$1T^nT_{}^{}{}_{}{}^{n}=\underset{i=0}{\overset{n1}{}}T^i(1TT^{})T_{}^{}{}_{}{}^{i}.$$
Hence
$`\underset{n\mathrm{}}{lim}{\displaystyle \frac{\mathrm{tr}(1T^nT_{}^{}{}_{}{}^{n})}{n}}`$ $`=`$ $`\underset{n\mathrm{}}{lim}{\displaystyle \frac{1}{n}}{\displaystyle \underset{i=0}{\overset{n1}{}}}\mathrm{tr}(T_{}^{}{}_{}{}^{i}T^i(1TT^{}))`$
$`=`$ $`\underset{i\mathrm{}}{lim}\mathrm{tr}(T_{}^{}{}_{}{}^{i}T^i(1TT^{})).`$
The last equality was obtained as follows. Consider the sequence
$`a_i`$ $`:=`$ $`\mathrm{tr}(T_{}^{}{}_{}{}^{i}T^i(1TT^{}))`$
$`=`$ $`\mathrm{tr}((T^i\mathrm{\Delta }_T)^{}(T^i\mathrm{\Delta }_T))`$
$`=`$ $`{\displaystyle \underset{k=1}{\overset{q}{}}}(T^i\mathrm{\Delta }_T)^{}(T^i\mathrm{\Delta }_T))e_k,e_k`$
$`=`$ $`{\displaystyle \underset{k=1}{\overset{q}{}}}T^i\mathrm{\Delta }_Te_k^2.`$
The last expression makes clear that $`a_1a_2\mathrm{}0`$, so that the sequence has a limit $`L=lim_i\mathrm{}a_i`$. We shall show that for any such sequence $`a_i`$,
$$\underset{n\mathrm{}}{lim}(1/n)\underset{i=0}{\overset{n1}{}}a_i=L.$$
Since the sequence $`\left\{\frac{1}{n}_{i=0}^{n1}a_i\right\}`$ is bounded above by $`a_0`$, it is enough to show that its only possible accumulation point is $`L`$.
For any fixed $`m`$ and all $`nm`$,
$$L\frac{1}{n}\underset{i=0}{\overset{n1}{}}a_i\frac{1}{n}\underset{i0}{\overset{m1}{}}a_i+\frac{nm}{n}a_m.$$
Letting $`n`$ tend to infinity with $`m`$ fixed, we see that any accumulation point of the sequence $`\left\{\frac{1}{n}_{i=0}^{n1}a_i\right\}`$ must lie between $`L`$ and $`a_m`$. Finally, letting $`m`$ tend to infinity shows that $`L`$ is the only accumulation point. .
## 3 A simple formula for the curvature of a single, pure, contraction.
A single contraction $`T`$ on a Hilbert space $`H`$ will be called pure if for all $`hH`$, $`lim_n\mathrm{}T_{}^{}{}_{}{}^{n}h=0`$; i.e., if the adjoint powers $`T_{}^{}{}_{}{}^{n}`$ converge strongly to 0. This is the specialization to the case $`d=1`$ of Arveson’s more complicated definition of a pure $`d`$-contraction.
Arveson remarked that it is generally difficult to determine the curvature of a $`d`$-contraction, but that “in the few cases where the computations can be explicitly carried out, the curvature turns out to be an integer.” This led him to ask if the curvature of a pure $`d`$-contraction need always be an integer.
This was a surprising suggestion, because nothing in the definition of curvature suggests that it should be an integer. Subsequently, D. Greene, S. Richter, and C. Sundberg proved that indeed the curvature of any pure $`d`$-contraction is an integer. However, their function-theoretic methods do not seem to give an effective procedure for calculating this integer in particular cases, and a geometric, operator-theoretic interpretation of the curvature of a general $`d`$-contraction remains elusive as of this writing.
Our contribution toward understanding the meaning of the curvature invariant is a simple, usually easily computable, formula for the curvature of single, pure contraction; i.e., the special case $`d=1`$. It states that the curvature is the difference of the dimensions of two subspaces, and hence is obviously integral. The methods of proof are operator-theoretic, based on unitary dilation theory as set forth in . It uses neither the Greene/Richter/Sundberg result nor their function-theoretic methods, and thus gives an independent proof of their result for the special case $`d=1`$.
Our characterization of the curvature of a single pure contraction is:
###### Theorem 4
Let $`T`$ be a pure contraction operator such that $`\mathrm{\Delta }_T:=\sqrt{1TT^{}}`$ has finite rank. Then its curvature $`K(T)`$ is the integer
$$K(T)=dim\text{range }(1TT^{})dim\text{range }(1T^{}T).$$
(11)
A counterexample in the next section uses Theorem 2 to show that the hypothesis that $`T`$ be pure is essential.
Before proving the theorem, we review some standard facts about unitary dilations. Proofs can be found in , particularly Chapters 1, 2, and 6. We give specific references from this work for key facts required by the proof.
Let $`T`$ be a contraction on a Hilbert space $`H`$, and $`U`$ its minimal unitary dilation to a larger Hilbert space $`KH`$. This means that $`P_HU^n|H=T^n`$ for all $`n0`$, where $`P_H`$ denotes the projection to $`H`$, and minimality means that $`K=_{n=\mathrm{}}^{\mathrm{}}U^nH`$.
1. The minimal unitary dilation $`U`$ for $`T`$ may be constructed as follows. Define
$$K:=\mathrm{}\overline{\mathrm{\Delta }_TH}\overline{\mathrm{\Delta }_TH}H\overline{\mathrm{\Delta }_T^{}H}\overline{\mathrm{\Delta }_T^{}H}\mathrm{},$$
(12)
where the overscore denotes closure. (The closures turn out to be unnecessary in our context, but that only becomes apparent later.) Consider $`H`$ as embedded in $`K`$ in the obvious way. Then $`U`$ is defined on $`K`$ by:
$`U(\mathrm{},b_2,b_1,b_0,h,a_0,a_1,a_2,\mathrm{}):=`$ (13)
$`(\mathrm{},b_2,b_1,\overline{)Th+\mathrm{\Delta }_Tb_0},T^{}b_0+\mathrm{\Delta }_T^{}h,a_0,a_1,\mathrm{}).`$
Here $`b_i\overline{\mathrm{\Delta }_T}H`$, $`a_i\overline{\mathrm{\Delta }_T^{}}H`$, and zero’th components (vectors in $`H`$) are distinguished by boxes.
The realization of $`U`$ just given is best for some purposes, but a change of notation will bring out more clearly the features which will be important to us. Set $`:=(UT)H`$ and $`_{}:=(U^{}T^{})H`$. Informally, $``$ is the leftmost $`\overline{\mathrm{\Delta }_T^{}H}`$ factor in (12). The other $`\overline{\mathrm{\Delta }_T^{}H}`$ factors are images of the leftmost under positive powers of $`U`$. Similarly, $`_{}`$ is the rightmost $`\overline{\mathrm{\Delta }_TH}`$ factor, and the other $`\overline{\mathrm{\Delta }_TH}`$ factors are images of it under negative powers of $`U`$.
To reflect these insights, instead of realizing $`K`$ as above, think of it as follows:
$$K\mathrm{}U^2_{}U^1_{}_{}HUU^2\mathrm{}$$
(14)
Here $``$ stands for unitary equivalence.
The conceptual advantage of (14) is that it makes clear at a glance much of the action of $`U`$ on $`K`$. Unfortunately, it is awkward for the purpose of defining $`U`$ due to logical circularity.
Embedded in $`U`$ are two bilateral shifts which interact in a complicated way. One shifts $``$, and the other shifts $`_{}`$. One half of each shift is transparently visible in (14). For example, $`U`$ obviously acts as a unilateral shift (with multiplicity $`\text{dim}`$) on the invariant subspace
$$M()^+:=\underset{n=0}{\overset{\mathrm{}}{}}U^n.$$
(15)
Since all iterates $`U^n,\mathrm{}n\mathrm{}`$ are easily seen to be pairwise orthogonal, also $`U`$ acts as a bilateral shift on the invariant subspace
$$M():=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}U^n,$$
(16)
but the left half of this subspace, $`M()M()^+`$, is embedded in a non-transparent way in $`K`$.
A subspace $`𝒮`$ such that the subspaces $`U^n𝒮`$ are pairwise orthogonal, $`\mathrm{}<n<\mathrm{}`$, is called a wandering subspace for $`U`$. Thus $``$ is a wandering subspace, and so is $`_{}`$. For any wandering subspace $`𝒮`$, we’ll use the notation $`M(𝒮)`$ as defined in (16) with $``$ replaced by $`𝒮`$.
2. The contraction $`T`$ is pure, i.e., $`T_{}^{}{}_{}{}^{n}0`$ strongly, if and only if
$$K=M(_{}):=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}U^n_{}$$
(Chap. 2, Thm. 1.1, p. 57).
3. If $`T`$ is pure, then $`dimdim_{}`$; equivalently, $`\text{rank }\mathrm{\Delta }_T^{}\text{rank }\mathrm{\Delta }_T`$. In particular, under our hypotheses that $`T`$ is pure with $`\mathrm{\Delta }_T`$ of finite rank, also $`\mathrm{\Delta }_T^{}`$ has finite rank, and both $``$ and $`_{}`$ are finite dimensional.
This follows from item 2 above combined with , Chap. 1, Prop. 2.1, p. 4. Alternatively, it can be obtained for the case that we’ll need, $`dim_{}<\mathrm{}`$, from the Reciprocity Lemma 5 below with $`^{}:=_{}`$. Assuming temporarily that $`dim`$ is known to be finite, the Reciprocity Lemma applies as follows:
$$dim=\mathrm{tr}(P_{})=\mathrm{tr}(P_{}P_{M(_{})})=\mathrm{tr}(P_{_{}}P_{M()})\mathrm{tr}(P_{_{}})=dim_{}.$$
The case of an infinite-dimensional $``$ can be ruled out by applying the same reasoning with $``$ replaced by finite-dimensional subspaces of $``$.
4. When $`T`$ is a partial isometry,
$$U_{}=\overline{\mathrm{\Delta }_TH}.$$
In particular, $`U_{}H`$.
This is immediate from (13) after recalling that a partial isometry $`T`$ satisfies $`T^{}(1TT^{})H=\{0\}`$.
Let $`E`$ and $`F`$ be projections on a Hilbert space, at least one of which has finite rank. Then $`\mathrm{tr}(EF)=\mathrm{tr}(E^2F)=\mathrm{tr}(EFE)`$, so $`\mathrm{tr}(EF)`$ is always non-negative, is zero if and only if $`E`$ and $`F`$ have orthogonal ranges, and takes on its maximum value dim ($`E`$) or dim ($`F`$) only when $`EF`$ or $`FE`$. Thus $`\mathrm{tr}(EF)`$ serves as a measure of how nearly the ranges of $`E`$ and $`F`$ coincide. For lack of a standard term, call $`\mathrm{tr}(EF)`$ the affinity between the ranges of $`E`$ and $`F`$.
The following lemma, which we call the Reciprocity Lemma, may have some interest in its own right. It states that for wandering subspaces $``$ and $`^{}`$ for a unitary operator $`U`$, the affinity between $``$ and the closed span of the iterates $`U^n^{},\mathrm{}n\mathrm{}`$ is invariant under interchange of $``$ and $`^{}`$.
###### Lemma 5
\[Reciprocity Lemma\] Let $``$ and $`^{}`$ be finite dimensional wandering subspaces for a unitary operator $`U`$ on a Hilbert space $`K`$, and set
$$M():=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}U^n\text{and}M(^{}):=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}U^n^{}.$$
Then, denoting by $`P_S`$ the projection on an arbitrary subspace $`S`$ of $`K`$,
$$\mathrm{tr}(P_{}P_{M(^{})})=\mathrm{tr}(P_{^{}}P_{M()}).$$
Proof:
Note that
$$P_{M()}=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}P_{U^n}=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}U^nP_{}U^n,$$
(17)
the sums converging in the strong operator topology.
Multiply (17) by $`P_{^{}}`$ on the left, take the trace of both sides, and suppose we can justify an interchange of sum and trace, obtaining
$`\mathrm{tr}(P_{^{}}P_M_{})`$ $`=`$ $`\mathrm{tr}({\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}P_{^{}}U^nP_{}U^n)`$ (18)
$`=`$ $`{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\mathrm{tr}(P_{^{}}U^nP_{}U^n).`$
Then the following simple calculation establishes the lemma:
$`\mathrm{tr}(P_{^{}}P_{M()})`$ $`=`$ $`{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\mathrm{tr}(P_{^{}}U^nP_{}U^n)`$ (19)
$`=`$ $`{\displaystyle \underset{n}{}}\mathrm{tr}(P_{}U^nP_{^{}}U^n)`$
$`=`$ $`\mathrm{tr}(P_{}P_{M(^{})}),`$
where the last line was obtained from (18) with $``$ and $`^{}`$ interchanged and the summation index $`n`$ replaced by $`n`$.
The interchange of sum and trace required to justify the above calculation is not immediate because the trace is not continuous in the strong operator topology. However, the trace is well-known to be a normal linear functional, which implies that for any increasing sequence of trace class positive operators
$$A_1A_2\mathrm{}A_m\mathrm{}$$
converging in the strong operator topology to a trace class operator $`A`$, we have
$$\underset{m\mathrm{}}{lim}\mathrm{tr}(A_m)=\mathrm{tr}(A).$$
(20)
This property is a slight specialization of the definition of normality. It follows routinely from the definition $`\mathrm{tr}A:=_{i=1}^{\mathrm{}}Ae_i,e_i`$, with $`\{e_i\}`$ an orthonormal basis.
Noting that
$$\mathrm{tr}(P_{^{}}U^nP_{}U^n)=\mathrm{tr}(P_{^{}}U^nP_{}U^nP_{^{}})$$
and applying (20) with $`A_m:=_{n=m}^m\mathrm{tr}(P_{^{}}U^nP_{}U^nP_{^{}})`$ proves (18).
I thank W. Arveson for suggesting the above proof to replace the unattractive direct calculation of an earlier draft.
Proof of Theorem 4:
Proposition 1 shows that we may assume that $`T`$ is a partial isometry. We are going to use Theorem 2 to calculate $`K(T)`$ by calculating
$$\underset{n\mathrm{}}{lim}T^ne^2$$
for $`e\mathrm{\Delta }_TH`$. For any $`hH`$,
$$T^nh^2=P_HU^nh^2.$$
Since $`U^nhHM()^+=H(_{k=0}^{\mathrm{}}U^k`$),
$`T^nh^2`$ $`=`$ $`U^nh^2P_{M()^+}U^nh^2`$ (21)
$`=`$ $`h^2P_{M()^+}U^nh^2.`$
Write $`K=M()R`$, where (as always), the direct sum denotes an orthogonal direct sum, so this defines the subspace $`R`$, which reduces $`K`$ because $`M()`$ does. Then any $`kK`$ can be written
$$k=\underset{i=\mathrm{}}{\overset{\mathrm{}}{}}U^if_i+r,$$
with $`f_i`$ and $`rR`$. And, for any $`hH`$,
$$h=\underset{i=\mathrm{}}{\overset{1}{}}U^if_i+r.$$
(22)
Substituting (22) in (21) gives:
$`\underset{n\mathrm{}}{lim}T^nh^2`$ $`=`$ $`h^2\underset{n\mathrm{}}{lim}{\displaystyle \underset{i=n}{\overset{1}{}}}f_i^2`$ (23)
$`=`$ $`h^2P_{M()}h^2`$
$`=`$ $`h^2P_{M()}h,h.`$
Choose an orthonormal basis $`e_1,\mathrm{},e_q`$ for $`\text{range }\mathrm{\Delta }_TH`$. Then substituting (23) in Theorem 2 gives:
$`K(T)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{q}{}}}\underset{n\mathrm{}}{lim}T^ne_i^2`$ (24)
$`=`$ $`q{\displaystyle \underset{i=1}{\overset{q}{}}}P_{M()}e_i,e_i`$
$`=`$ $`q\mathrm{tr}(P_{M()}P_{\mathrm{\Delta }_TH}).`$
Item 2 remarked that for a partial isometry $`T`$, $`\mathrm{\Delta }_TH=U_{}`$, and substituting this in (24) gives:
$$K(T)=q\mathrm{tr}(P_{M()}P_U_{}).$$
(25)
By the Reciprocity Lemma 5,
$$\mathrm{tr}(P_{M()}P_U_{})=\mathrm{tr}(P_{M(U_{})}P_{}).$$
But obviously,
$$M(U_{}):=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}U^nU_{}=\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}U^k_{}=M(_{}).$$
Combining these facts gives the desired conclusion:
$`K(T)`$ $`=`$ $`q\mathrm{tr}(P_{M()})P_U_{})=q\mathrm{tr}(P_{M(U_{})}P_{})`$
$`=`$ $`q\mathrm{tr}(P_{M(_{})}P_{})=q\mathrm{tr}(P_{})`$
$`=`$ $`dim\text{range }\mathrm{\Delta }_Tdim\text{range }\mathrm{\Delta }_T^{}`$
$`=`$ $`dim\text{range }(1TT^{})dim\text{range }(1T^{}T).`$
The second line follows from item 2’s observation that the hypothesis that $`T`$ be pure is equivalent to $`M(_{})=K`$.
Recall that a Fredholm operator $`T`$ is one with closed range and finite-dimensional kernel and cokernel (denoted $`\text{ker }(T)`$ and $`\text{coker }(T):=\mathrm{ker}(T^{})`$). The index of a Fredholm operator $`T`$ is defined by
$$\text{index }(T):=dim\text{ker }(T)dim\text{coker }T.$$
(26)
A fundamental theorem (e.g., , p. 128, Thm. 5.36) states that the index is invariant under compact perturbations: for any Fredholm operator $`T`$ and compact operator $`C`$, $`T+C`$ is Fredholm, and $`\text{index }(T+C)=\text{index }(T)`$.
Formula (26) makes sense when $`T`$ has finite-dimensional kernel and cokernel even if $`T`$ doesn’t have closed range. However, since the closed range hypothesis is needed to prove the fundamental theorem just mentioned, the term “index” is generally restricted to Fredholm operators. Nevertheless, for purposes of the present exposition, it will be convenient to broaden the definition of $`\text{index }(T)`$ to include cases in which $`T`$ has finite-dimensional kernel and cokernel, but not necessarily closed range.
When told of the curvature formula (11) given by Theorem 4, W. Arveson remarked that it looked something like an operator index and that he had been working on a conjecture that under appropriate hypotheses, the curvature of a $`d`$-contraction would be the index of an associated operator which he calls $`D_+`$, reminiscent of the Dirac operator. Shortly thereafter, he wrote up these results in , which proves this for $`d`$-contractions whose associated Hilbert modules are finite rank, pure, and graded, in the terminology of . It asks if the “graded” hypothesis can be removed, and also if the associated operator $`D_+`$ necessarily has closed range (and so is Fredholm).
For the case of a 1-contraction, the associated operator $`D_+`$ is unitarily equivalent to $`T`$. Corollary 6 below observes that under the hypotheses of Theorem 4, $`T`$ is Fredholm, and its curvature equals $`\text{index }(T)`$. The differences between Corollary 6 and the specialization of Arveson’s result to the single operator case are that the closed range property is proved for $`d=1`$, and the “graded” hypothesis is not needed. This holds out hope that the “graded” and “closed range” hypotheses might be removable for $`d`$-contractions with $`d>1`$.
The interest in identifying the curvature with an index, apart from its evident aesthetic appeal, is that the index is stable under compact perturbations, but the curvature is not known to possess such stability. The strongest result along these lines known as of this writing is , Corollary 1, Stability of Curvature, which proves stability of the curvature under certain special finite rank perturbations. Arveson notes that removing the “closed range” hypothesis would establish a much stronger stability of curvature result, and removing the “graded” hypothesis would strengthen it further.
###### Corollary 6
Let $`T`$ be an operator satisfying the hypotheses of Theorem 4. Then $`T`$ is Fredholm, and
$$K(T)=\mathrm{index}(T).$$
Proof: Let $`T`$ be an operator on a Hilbert space $`H`$ satisfying the hypotheses of Theorem 4. First we sketch the simple proof that $`T`$ must be Fredholm.
The assumed finiteness of the rank of $`1TT^{}`$ implies that $`\text{coker }(T)`$ is finite-dimensional. Since we have already noted that $`\text{rank }(1T^{}T)\text{rank }(1TT^{})`$, also $`\text{ker }(T)`$ is finite-dimensional. That $`T`$ must have closed range under these circumstances can be easily seen by noting that closed range is equivalent to a gap above 0 in the spectrum of $`T^{}T`$. If there were not such a gap, then $`1T^{}T`$ would not have finite rank.
To show that (11) equals $`\text{index }(T)`$, let $`\stackrel{~}{T}`$ be the operator on $`H\text{range }(1TT^{})`$ defined by the operator matrix:
$$\stackrel{~}{T}:=\left[\begin{array}{cc}T\hfill & 0\hfill \\ 0\hfill & 0\hfill \end{array}\right].$$
Let $`Q`$ be the partial isometry
$$Q:=\left[\begin{array}{cc}T\hfill & \sqrt{1TT^{}}\hfill \\ 0\hfill & 0\hfill \end{array}\right]$$
of Proposition 1. Since $`Q`$ is a compact perturbation of $`\stackrel{~}{T}`$,
$$\text{index }(Q)=\text{index }(\stackrel{~}{T})=\text{index }(T).$$
Also, since $`Q`$ is a partial isometry, $`dim\text{ker }(Q)=dim\text{range }(1Q^{}Q)`$ and $`dim\text{coker }(Q)=dim\text{range }(1QQ^{})`$. Hence by Proposition 1,
$`\text{index }(T)=\text{index }(Q)`$ $`:=`$ $`dim\text{ker }(Q)dim\text{coker }(Q)`$
$`=`$ $`dim\text{range }(1Q^{}Q)dim\text{range }(1QQ^{})`$
$`=`$ $`dim\text{range }(1T^{}T)dim\text{range }(1TT^{})`$
$`=`$ $`K(T).`$
Remark: The above proof that $`\text{index }(T)=dim\text{range }(1T^{}T)dim\text{range }(1TT^{})`$ is concise and natural within our context, but may not be the most insightful. A slightly messier but more straightforward proof can be based on the well-known fact that for any operator $`T`$, the restriction of $`T^{}T`$ to its initial space (defined as the orthogonal complement of its nullspace) is unitarily equivalent to the restriction of $`TT^{}`$ to the closure of its initial space. (The equivalence can be implemented by the partial isometry $`U`$ in the polar decomposition $`T=U\sqrt{T^{}T}`$ restricted to its initial space.) From this it follows that any nonzero eigenvalue for $`T^{}T`$ is also an eigenvalue for $`TT^{}`$, with the same multiplicity, so that in the expression $`dim\text{range }(1T^{}T)dim\text{range }(1TT^{})`$, the dimensions of the eigenspaces corresponding to nonzero eigenvalues cancel, leaving the only contribution to this expression as $`dim\text{ker }(T)dim\text{ker }(T^{})=\text{index }(T)`$.
## 4 A contraction with non-integral curvature
Now we apply Theorem 2 to construct a simple example of an operator with non-integral curvature, in fact with arbitrary real curvature $`\kappa 0`$. This shows that Theorem 4’s hypothesis that $`T`$ be pure cannot be omitted. I have been told that the existence of non-pure contractions with non-integral curvatures was implicitly known or expected by experts in the field, so the interest of the example may lie more in its simplicity than novelty.
It is enough to produce a partial isometry $`Q`$ with $`\text{range }\mathrm{\Delta }_Q`$ spanned by a single unit vector $`e`$ satisfying
$$\underset{n\mathrm{}}{lim}Q^ne^2=\kappa .$$
First suppose $`0\kappa 1`$, and set $`\lambda :=\sqrt{1\kappa }`$.
Let $`T`$ be the bilateral weighted shift defined on an orthonormal basis $`\{e_n\}_{n=\mathrm{}}^{\mathrm{}}`$ by:
$$Te_n:=\{\begin{array}{cc}e_{n+1}\hfill & \text{if }n0\hfill \\ \lambda e_1\hfill & \text{if }n=0\hfill \end{array}.$$
Then one routinely computes that $`\mathrm{\Delta }_T:=\sqrt{1TT^{}}`$ is the rank 1 operator whose only non-zero eigenvalue is $`\sqrt{1\lambda ^2}=\sqrt{\kappa }`$, with corresponding eigenvector $`e_1`$.
Let $`Q`$ be the associated partial isometry given by Proposition 1. We may realize $`Q`$ as acting on a space with orthonormal basis $`\{e_{\mathrm{}}\}\{e_n\}_{n=\mathrm{}}^{\mathrm{}}`$ obtained by adjoining a new unit vector named $`e_{\mathrm{}}`$ to the previous orthonormal basis for the space on which $`T`$ was defined. Then $`Q`$ is defined by $`Qe_n:=Te_n`$ for $`n`$ finite, and $`Qe_{\mathrm{}}:=\sqrt{1\lambda ^2}e_1=\sqrt{\kappa }e_1`$.
As in Proposition 1, one routinely computes that $`\mathrm{\Delta }_Q:=\sqrt{1QQ^{}}`$ is the one-dimensional projection with range spanned by $`e_{\mathrm{}}`$. From Theorem 2
$$K(Q)=\underset{n\mathrm{}}{lim}Q^ne_{\mathrm{}}^2=\underset{n\mathrm{}}{lim}T^{n1}\sqrt{\kappa }e_1^2=\kappa .$$
This shows that any $`\kappa `$ with $`0\kappa 1`$ can be the curvature of some contraction. To see that any real number can be the curvature of some contraction, first check that curvature is additive over direct sums: for any two contractions $`T_1,T_2`$, we have
$$K(T_1T_2)=K(T_1)+K(T_2).$$
This follows routinely from the original definition (1) of curvature, or slightly more easily, from Theorem 2. Then any desired non-negative real curvature can be obtained by direct summing appropriate copies of the above example.
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# VLT-ISAAC near-IR Spectroscopy of ISO selected Hubble Deep Field South Galaxies1footnote 11footnote 1Based on observations with ISO, an ESA project with instruments funded by ESA member states (especially the PI countries: France, Germany, the Netherlands, and the United Kingdom) with the participation of ISAS and NASA. 2footnote 22footnote 2Based on observations collected at the European Southern Observatory, Chile, ESO No 63.O-0022
## 1 Introduction
Until recently, most of our knowledge about high-z galaxies has come from optical surveys. The COBE detection of an extragalactic far–infrared$`/`$submm background (Puget et al. 1996), with an integrated intensity similar to or greater than that of optical light (e.g. Hauser et al. 1998), strongly suggests that a significant fraction of the cosmic star formation in the Universe is obscured by dust and thus missed by the various optical surveys.
With the advent of the Infrared Space Observatory (ISO, Kessler et al. 1996) deep mid–IR surveys for distant galaxies, have been successfully carried out for the first time. Operating in the 5 – 18 $`\mu `$m band sensitive to warm dust and emission from Polycyclic Aromatic Hydrocarbons (PAH), ISOCAM on board ISO was more than 1000 times more sensitive than IRAS and thus had the potential to study infrared bright galaxies at redshifts beyond 0.5. A number of cosmological surveys have been performed with ISOCAM especially in the LW3 filter (12 – 18$`\mu `$m). These surveys range from the wide and shallow European Large Area ISO Survey (ELAIS, Oliver et al. 2000), to deep pencil–beam surveys in the Lockman Hole, the Marano field, the Northern and Southern Hubble Deep Fields (HDF–N, HDF–S) and the distant cluster Abell 2390, reaching limiting flux densities of 50 – 100 $`\mu `$Jy (Elbaz et al. 1999).
At the bright end the ISOCAM source counts combined with those of IRAS are in good agreement with no or moderate evolution. At fainter flux densities the counts steepen considerably and at $``$ 200 – 600 $`\mu `$Jy they are about an order of magnitude greater than the predictions of no evolution models. This steepening in the log N–log S plot and a pronounced maximum in the differential number counts at $``$ 400 $`\mu `$Jy suggest that the ISOCAM surveys have revealed a population of strongly evolving galaxies. Elbaz et al. (2000, in preparation) show that this population plausibly accounts for a significant fraction of the far-IR background. The next step is to explore the nature of the ISOCAM population with optical$`/`$near-IR spectroscopy. In this letter we report on the first near-IR (rest-frame R–band) spectroscopic survey of a representative sample of faint ISOCAM galaxies in the HDF–S field.
## 2 Sample Selection
The HDF–S was observed by ISOCAM as part of the ELAIS survey. The observations were carried out at two wavelengths, LW2 ( 6.75 $`\mu `$m) and LW3 (15 $`\mu `$m). Oliver et al. (2000, in preparation) and Aussel et al. (2000, in preparation) analyzed the data independently. Aussel et al. used the PRETI method and detected 63 sources brighter than S<sub>15μm</sub> = 100 $`\mu `$Jy in the LW3 band. We used this source list as input for our observations.
We selected the sample for ISAAC follow up from the HDF–S LW3 sources based on the following criteria: a) a reliable LW3 detection, b) H<sub>α</sub> in the wavelength range of ISAAC and, c) a secure counterpart in the I band image (Dennefeld et al. 2000, in preparation), or a counterpart in the K band image (ESO Imaging Survey (EIS) Deep). We did not apply any selection based on colors. Our reference sample contains 25 galaxies with 15 $`\mu `$m flux densities ranging between 100–800 $`\mu `$Jy. It is thus a fair representation of the strongly evolving ISOCAM population near the peak of the differential source counts (Elbaz et al. 1999). From these 25 optically identified sources we randomly selected 12 sources for ISAAC–follow up. To select the near-IR band (Z, SZ, J, H) for our spectroscopy we used spectroscopic redshifts from optical spectroscopy, where available, for z$`<`$0.7 (Dennefeld et al. 2000, in prep.). Otherwise we used photometric redshift estimates based on the model PEGASE (Fioc and Rocca-Volmerange 1997). Our photometric redshift determinations turned out to be accurate to $`\mathrm{\Delta }`$z$`\pm `$0.1.
## 3 Observations and Results
We collected the spectra during 1999 September 20–24 with the infrared spectrometer ISAAC (Moorwood et al. 1998) on the ANTU–ESO telescope (formerly UT1), on Paranal, Chile. For the observations we used the low resolution grating R$`{}_{s}{}^{}`$ 600 and a 1$`\times `$2 long slit. To maximize the observing efficiency each slit position included on average two galaxies at any given orientation. Most of the targets were first acquired directly from a 1–2 min exposure in the H-band. In the case of the very faint objects (H $``$ 20.0 mag) we offset from a brighter star in the HDF-S field. Observations were made by nodding the telescope $`\pm `$ 20 along the slit to facilitate sky subtraction (always avoiding overlap of the two objects in the slit). Individual exposures ranged from 2–4 minutes. Sky conditions were excellent throughout the acquisition of the spectra, with seeing values in the range 0.4–1.0. For each filter, observations of spectroscopic standard stars were made in order to flux calibrate the galaxy spectra.
The data were reduced using applications from ECLIPSE (Devillard 1998) and IRAF packages. Accurate sky subtraction is critical to the detection of faint lines. Sky was removed by subtracting the pairs of offset frames. In some cases this left a residual signal (due to temporal sky changes) which was then removed by performing a polynomial interpolation along the slit. OH sky emission lines were also carefully removed from the spectra. Spectrum extraction for each galaxy was performed using the APEXTRACT package. Standard wavelength calibration was applied.
The spectra for all 12 galaxies observed with ISAAC are shown in Figure 1. Table 1 contains exposure times, H<sub>AB</sub> magnitudes, measured spectroscopic redshifts (from H<sub>α</sub> detections), H<sub>α</sub> line fluxes and Equivalent Widths (EW) (and where resolved \[NII\] fluxes). We convert the H<sub>α</sub> line fluxes to luminosities using H<sub>o</sub> = 50 km s<sup>-1</sup> Mpc<sup>-1</sup> and $`\mathrm{\Omega }`$ = 0.3. We also list in Table 1 FIR luminosities based on LW3 fluxes (see Section 5 for more details). We note that we have not detected any H<sub>α</sub> Broad Line components.
## 4 The nature of the ISOCAM faint galaxies: Dusty and Luminous Starbursts
Prior to our study no near-infrared (rest-frame R-band) spectroscopy had been carried out for the ISOCAM population, primarily because of the faintness of the galaxies. Aussel et al. (1999) and Flores et al. (1999) have presented optical spectroscopic analysis (rest-frame B-band) for HDF-N and the 1415$`+`$52 field of the Canada-France Redshift Survey (CFRS), respectively. Aussel et al. (1999) cross-correlated the ISOCAM HDF-N galaxies with the optical catalog of Barger et al. (1999) resulting in 38 galaxies with confirmed spectroscopic redshifts. Flores et al. have identified 22 galaxies with confirmed spectroscopic information. In both of these samples the median redshift is about 0.7. Our ISOHDFS sample contains 7 galaxies 0.4$`<`$z$`<`$0.7 and 5 galaxies with 0.7$`<`$z$`<`$1.4 and thus, has a z-distribution very similar to the HDF-N (Aussel et al.) and CFRS (Flores et al.) samples.
Rest-frame B-band spectra host a number of emission and absorption lines related to the properties of the starburst in the galaxy. Based on these features galaxies can be classified according to their starburst history. Strong H<sub>δ</sub>, H<sub>ϵ</sub> Balmer absorption and no emission lines are characteristic of passively evolving k$`+`$A galaxies. The presence of Balmer absorption lines implies the presence of a dominating A-star population formed about 0.1–1 Gyr ago. The simultaneous presence of Balmer absorption and moderate flux \[OII\] and H<sub>β</sub> emission termed as e(a) or S$`+`$A galaxies, indicates that, in addition, there is ongoing star formation. The relative importance of these star formation episodes depends on the extinction (especially of the current star formation component). If the extinction is low the galaxy is primarily a post–starburst system. If the extinction toward the star forming region is high then the galaxy could be a powerful starburst. The majority of the galaxies in the CFRS field ($``$ 70%) display optical spectra characteristic of e(a) galaxies. As evidenced by the detections in Figure 1, the ISOCAM galaxies are in fact powerful starbursts hidden by large amounts of dust extinction.
Remarkably, dusty starbursts such as M82 (L$``$ 10<sup>10</sup> L, Kennicutt et al. 1992), LIRGs (Luminous InfraRed Galaxies, L$``$ 10<sup>11</sup> L, Wu et al. 1998), and many bright ULIRGs (L$``$ 10<sup>12</sup> L, Liu and Kennicutt 1995), show e(a) B–band spectra. Local e(a) galaxies have large H<sub>α</sub> equivalent widths (EW), at the same time demonstrating active current star formation and differential dust extinction. The measured EW ratio (\[OII\]$`/`$H<sub>α</sub>) for e(a) galaxies appears to be somewhat low. Such low ratios have already been observed in the spectra of distant clusters (Dressler et al. 1999), nearby mergers (Poggianti and Wu 2000), the dusty LIRGs studied by Wu et al. (1998) or the interacting$`/`$merging systems studied by Liu and Kennicutt (1995). The behaviour of the (OII)$`/`$(H$`{}_{\alpha }{}^{}+`$NII) ratio is shown in the EW(OII)–EW(H$`{}_{\alpha }{}^{}+`$ NII) diagram of Figure 2: the majority of the points lie below the straight line. For our ISOHDFS sample we use the EW(H<sub>α</sub>) measured from the observations presented here. For the EW(OII) we use a median value of 20$`\pm `$15 Ȧ which was recently measured from FORS2 I-band spectra of a small sample of ISOHDFS galaxies (Franceschini et al. 2000, in preparation). This value is in agreement with the results presented by Flores et al. (1999) for the CFRS galaxies.
It follows from Figure 2 that the ISOHDFS galaxies occupy the same region in the EW(OII)$`/`$EW(H$`{}_{\alpha }{}^{}+`$ NII) diagram as actively starforming galaxies. Intrinsic differential dust extinction is responsible for the somewhat low EW(OII)$`/`$ EW(H$`{}_{\alpha }{}^{}+`$NII) ratio. The \[OII\] emission is affected more than H<sub>α</sub> simply because of its shorter wavelength. The continuum is due to A-stars which come from earlier (0.1–1.0 Gyr) star formation activity that is not energetically dominant and plays a small role once the dusty starburst is dereddened. This scenario implies that these galaxies undergo multiple burst events: the less extincted population is due to an older burst while in the heavily dust enshrouded HII regions there is ongoing star formation. We conclude that ISOCAM galaxies are actively starforming, dust enshrouded galaxies, akin to local LIRGs (e.g. NGC 3256, Rigopoulou et al. 1996).
## 5 Star Formation Rates and Extinction Corrections
The conversion factor between ionizing luminosity and star formation rate (SFR) is usually computed using an evolutionary synthesis model. Only massive stars ($`>`$ 20 M) with short lifetimes ($``$ 10<sup>6</sup> yrs) contribute to the integrated ionizing flux. Using the stellar synthesis code STARS (Sternberg 1998) we create models for solar abundances, a Salpeter IMF (1–100 M) and slowly decaying bursts with ages in the range of a few$`\times `$ 10<sup>7</sup>–10<sup>8</sup> yrs, and SFR decay time-scales in the range 10<sup>7</sup>-10<sup>9</sup> yrs. Averaging, we obtain:
SFR(M$`{}_{}{}^{}/yr`$) = 5 $`\times `$10<sup>-42</sup>L$`_{H_\alpha }`$(erg s<sup>-1</sup>).
(1)
We have used this formula to estimate the SFR rates in Table 2. The SFR estimates based on Eqn.(1) are a factor of 1.6 smaller than the SFR estimates based on the Kennicutt (1998) relationship that refers to stars in the range 0.1–100 M. Averaging over our models, the SFR scales with the FIR luminosity as:
SFR(M$`{}_{}{}^{}/yr`$) = 2.6 $`\times `$10<sup>-44</sup>L<sub>FIR</sub>(erg s<sup>-1</sup>)
(2)
Since extinction is at play, the SFR estimates we quote in the first column of Table 2 are lower limits to the real SFR in these galaxies. We derive the extinction based on V–K color indices (magnitudes taken from the EIS Survey). Using STARS as well as the Starburst99 (Leitherer et al. 1999) codes for various star formation histories (ie bursts of different duration, and continuous star formation) we calculate the range of intrinsic colors. The model predicted intrinsic V–K colors are in the range 1.1 – 1.5. We have applied infrared and optical K-corrections from Poggianti (1997) and Coleman (1980), respectively. Comparing the observed V–K colors to the predicted ones we obtain a median color excess of 2.0 which corresponds to a median A<sub>V</sub> of 1.8 assuming a screen model for the extinction. This A<sub>V</sub> value corresponds to a median correction factor for the SFR(H<sub>α</sub>) of $``$4.
The SFR can also be inferred from far-infrared (FIR) luminosities according to Eqn. (2). The SFR(FIR) estimates in Table 2 are based on the method of Franceschini et al. (2000, in prep.) which makes use of the 15 $`\mu `$m flux and assumes a L$`{}_{FIR}{}^{}/`$L<sub>MIR</sub> ratio of $``$10 (for an M82 like SED, Vigroux et al. 1999). The SFR(FIR) estimates turn out to be a factor of 5 to 50 higher than the SFR estimates inferred from the non-extinction corrected H<sub>α</sub>. However, if we apply the correction factor of $``$4 we deduced for the H<sub>α</sub> then SFR(FIR)$`/`$SFR(H<sub>α</sub>)$``$ 3, confirming that the extinction is much higher than can be predicted using (UV or) optical observations. Thus, ISOCAM galaxies are in fact actively star forming highly dust enshrouded galaxies. One exception is source ISOHDFS 38 for which we have evidence for the presence of a dominant AGN component (both from the LW2$`/`$LW3 ratio and the H$`{}_{\alpha }{}^{}/`$\[NII\] ratio).
Based on our SFR estimates, and the evidence for exinction presented in section 4.1 we conclude that ISOCAM galaxies are indeed dust enshrouded actively star-forming galaxies and not decaying post-starburst systems.
## 6 Conclusions
We have presented NIR spectroscopy, rest frame R-band, of a sample of ISO selected galaxies from the Hubble Deep Field South. We have detected H<sub>α</sub> emission in almost all of them. The detections of the H<sub>α</sub> line combined with the large H<sub>α</sub> EWs are consistent with the idea that these galaxies are ongoing powerful dusty starbursts.
Using the observed H<sub>α</sub> emission lines we estimate that the SFR rate in the ISOHDFS galaxies ranges between 2 and 50 M$`{}_{}{}^{}/`$ yr, far higher than those inferred from local Starbursts (Calzetti 1997) and local spirals (Kennicutt 1992). We have compared these rates of SF with the values estimated from the FIR luminosities, which are typically a factor 5 to 50 larger because of dust obscuration. We estimated the H<sub>α</sub> extinction using standard extinction laws. The H<sub>α</sub> extinction corrected SFR estimates are then higher although still fall short of the SFR estimates based on FIR luminosities ( SFR(FIR):50–400 M$`{}_{}{}^{}/`$ yr). This result demonstrates that it is very dangerous to derive star formation rates from UV or optical data alone since these wavelengths are susceptible to higher extinction. Thus a significant fraction of star formation is missed by optical surveys. We conclude that ISO has detected in the mid-IR the most active, luminous and dust-enshrouded starbursts at z$``$ 0.4-1.4, which would have remained otherwise unnoticed by optical surveys. This population of strongly evolving active dusty starbursts is likely to account for a substantial fraction of the FIR$`/`$submm background (Elbaz et al. 2000, in preparation).
This work is supported by the EC TMR Network “European Large Area ISO Surveys” (contract No. ERBFMRX-CT96-0068). We thank the EC TMR Network “Galaxy Formation and Evolution ” for making redshift information available prior to publication. H.A. is supported by the TMR network “Galaxy formation and evolution”, contract No. ERBFRX-CT96-0086. We thank Amiel Sternberg for fruitful discussions.
FIGURES
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# Uniqueness of solutions to Hamilton-Jacobi equations arising in the Calculus of Variations
## 1 Introduction
Let us consider a Bolza problem of the Calculus of Variations
(1.1)
$$\mathrm{min}\left\{_0^tL(y(s),y^{}(s))𝑑s+\phi (y(t)):yW^{1,1}(0,t;IR^n),y(0)=x\right\},$$
where the final cost $`\phi :IR^nIR_+\{+\mathrm{}\}`$ is lower semicontinuous and the Lagrangian $`L:IR^n\times IR^nIR_+`$ is locally bounded and lower semicontinuous. We assume also that $`L(x,)`$ is convex for every $`xIR^n`$ and that the following Tonelli type coercivity assumption is satisfied: there exists a function $`\mathrm{\Theta }:IR^nIR_+`$ such that
(1.2)
$$\underset{|u|\mathrm{}}{lim}\frac{\mathrm{\Theta }(u)}{|u|}=+\mathrm{},(x,u)IR^n\times IR^n,L(x,u)\mathrm{\Theta }(u).$$
These assumptions guarantee the existence of absolutely continuous minimizers (see, e.g., ). The classical Lagrange problem (with the fixed final condition $`y(t)=z`$) may be reduced to the form (1.1) by simply setting $`\phi (z):=0`$ and $`\phi :=+\mathrm{}`$ elsewhere.
The value function $`V:IR_+\times IR^nIR_+\{+\mathrm{}\}`$ for the Bolza problem (1.1) is defined by
$$V(t,x):=\mathrm{min}\left\{_0^tL(y(s),y^{}(s))𝑑s+\phi (y(t)):yW^{1,1}(0,t;IR^n),y(0)=x\right\}.$$
Under our assumptions on $`L`$ and $`\phi `$ the value function is lower semicontinuous on $`IR_+\times IR^n`$ and locally Lipschitz on $`IR_+^{}\times IR^n`$, where $`IR_+^{}:=IR_+\{0\}`$. Moreover it satisfies the initial condition
$$xIR^n,\underset{h0+,yx}{lim\; inf}V(h,y)=\phi (x).$$
If $`V`$ is smooth and $`L`$ is continuous, then $`V`$ is a classical solution to the Hamilton-Jacobi equation
(1.3)
$$\{\begin{array}{cc}V_t+H(x,V_x)=0\hfill & \text{in }IR_+^{}\times IR^n\text{,}\hfill \\ & \\ V(0,)=\phi \hfill & \text{in }IR^n\text{,}\hfill \end{array}$$
with Hamiltonian $`H:IR^n\times IR^nIR`$ defined by
(1.4)
$$H(x,p):=\underset{uIR^n}{sup}\left(p,uL(x,u)\right),$$
where $`,`$ denotes the scalar product in $`IR^n`$. In other words, $`H(x,)`$ is the Legendre-Fenchel conjugate of $`L(x,)`$.
It is well known that (1.3) may not have smooth solutions, even if $`H`$ and $`\phi `$ are smooth. To overcome this difficulty, different notions of generalized solutions have been proposed. The notion of viscosity solution can be introduced by means of subdifferentials and superdifferentials.
We recall that the subdifferential $`_{}W(x)`$ of a function $`W:IR^nIR\{+\mathrm{}\}`$ at a point $`x\text{dom}(W)`$ is defined by
$$_{}W(x):=\{pIR^n:\underset{yx}{lim\; inf}\frac{W(y)W(x)p,yx}{|yx|}0\}.$$
while the superdifferential $`_+W(x)`$ is defined by $`_+W(x):=_{}(W)(x)`$.
A continuous function $`W:IR_+^{}\times IR^nIR`$ is said to be a viscosity solution to the Hamilton-Jacobi equation (1.3) if the following conditions are satisfied:
(1.5)
$$(t,x)IR_+^{}\times IR^n,(p_t,p_x)_{}W(t,x),p_t+H(x,p_x)0,$$
(1.6)
$$(t,x)IR_+^{}\times IR^n,(p_t,p_x)_+W(t,x),p_t+H(x,p_x)0.$$
In and the uniqueness of a bounded uniformly continuous viscosity solution of (1.3) is proved under some assumptions on $`H`$, which imply in particular that $`H`$ is continuous and $`H(,p)`$ is uniformly continuous for every $`pIR^n`$ . In \[10, Theorem 4.5\] we proved the following result, that can be applied also to unbounded solutions.
###### Theorem 1.1
Assume that $`L`$ is continuous. Let $`W:IR_+^{}\times IR^nIR_+`$ be a locally Lipschitz viscosity solution of (1.3) which satisfies the initial condition
(1.7)
$$xIR^n,\underset{h0+,yx}{lim\; inf}W(h,y)=\phi (x).$$
Then $`W=V`$ on $`IR_+^{}\times IR^n`$.
To describe the new uniqueness results, we introduce the notion of contingent directional derivatives of a function $`W:IR^nIR\{+\mathrm{}\}`$. These are defined, for every $`x\text{dom}(W)`$ and for every $`uIR^n`$, by
(1.8)
$$D_{}W(x)(u):=\underset{{\scriptscriptstyle \genfrac{}{}{0pt}{}{h0+}{vu}}}{lim\; inf}\frac{W(x+hv)W(x)}{h}.$$
The main result of this paper is the following theorem, which shows that the value function $`V`$ is the unique viscosity solution in the larger class of continuous functions whose contingent derivatives satisfy the following very weak assumptions:
(1.9)
$$(t,x)\text{dom}(W),t>0,D_{}W(t,x)(0,0)=0,$$
(1.10)
$$(t,x)\text{dom}(W),t>0,uIR^n,D_{}W(t,x)(1,u)<+\mathrm{}.$$
###### Theorem 1.2
Assume that $`L`$ is continuous. Let $`W:IR_+^{}\times IR^nIR_+`$ be a continuous viscosity solution of (1.3) which satisfies (1.9), (1.10), and the initial condition (1.7). Then $`W=V`$ on $`IR_+^{}\times IR^n`$.
In \[10, Theorem 4.4\] we considered also a different notion of generalized solution and we proved the following uniqueness result in the class of locally Lipschitz functions.
###### Theorem 1.3
Assume that $`L`$ is continuous. Let $`W:IR_+^{}\times IR^nIR_+`$ be a locally Lipschitz function which satisfies the initial condition (1.7) and solves the Hamilton-Jacobi equation (1.3) in the following sense :
(1.11)
$$(t,x)IR_+^{}\times IR^n,(p_t,p_x)_{}W(t,x),p_t+H(x,p_x)=0.$$
Then $`W=V`$ on $`IR_+^{}\times IR^n`$.
In this paper we shall prove the following result, which provides uniqueness in the larger class of lower semicontinuous functions $`W`$ satisfying (1.9) and (1.10).
###### Theorem 1.4
Assume that $`L`$ is continuous. Let $`W:IR_+^{}\times IR^nIR_+\{+\mathrm{}\}`$ be a lower semicontinuous function which satisfies the initial condition (1.7), the technical conditions (1.9) and (1.10), and solves the Hamilton-Jacobi equation (1.3) in the sense of (1.11). Then $`W=V`$ on $`IR_+^{}\times IR^n`$.
In the proof of Theorem 1.4 we use the following comparison result, which follows immediately from \[10, Theorem 4.1, Remark 4.2 and Proposition 4.3\].
###### Theorem 1.5
Assume that $`L`$ is continuous. Let $`W:IR_+^{}\times IR^nIR_+\{+\mathrm{}\}`$ be a lower semicontinuous function which satisfies the initial condition (1.7). Suppose that $`W`$ is a subsolution of the Hamilton-Jacobi equation (1.3) in the following sense:
(1.12)
$$(t,x)IR_+^{}\times IR^n,(p_t,p_x)_{}W(t,x),p_t+H(x,p_x)0$$
Then $`WV`$ on $`IR_+^{}\times IR^n`$.
In the proof of Theorems 1.2 and 1.5 we use also a very general comparison result (Theorem 1.7) for lower semicontinuous viscosity supersolutions of the Hamilton-Jacobi equation (1.3). To our knowledge the strongest result in this direction, dealing with possibly discontinuous Lagrangians, is the following theorem proved in \[10, Theorem 5.1\], where the notion of supersolution is given by using contingent inequalities.
###### Theorem 1.6
Let $`W:IR_+\times IR^nIR_+\{+\mathrm{}\}`$ be a lower semicontinuous function which satisfies the initial condition $`W(0,)=\phi `$. Suppose that $`W`$ is a supersolution of (1.3) in the following sense :
$$(t,x)\text{dom}(W),t>0,uIR^n,D_{}W(t,x)(1,u)L(x,u).$$
Then $`WV`$ on $`IR_+\times IR^n`$.
The comparison result for viscosity supersolutions we are going to prove is the following theorem, where we need the additional assumptions (1.9) and (1.10).
###### Theorem 1.7
Let $`W:IR_+\times IR^nIR_+\{+\mathrm{}\}`$ be a lower semicontinuous function which satisfies (1.9) and (1.10). Suppose that $`W(0,)=\phi `$ and that $`W`$ is a viscosity supersolution of the Hamilton-Jacobi equation (1.3), i.e., $`W`$ satisfies (1.5). Then $`WV`$ on $`IR_+\times IR^n`$.
## 2 Preliminaries
Let $`KIR^n`$ be a nonempty subset and $`xK`$. The contingent cone $`T_K(x)`$ to $`K`$ at $`x`$ is defined by
$$vT_K(x)\underset{h0+}{lim\; inf}\frac{\text{dist}(x+hv,K)}{h}=0.$$
The negative polar cone $`T^{}`$ to a subset $`TIR^n`$ is given by
$$T^{}:=\{vIR^n:wT,v,w0\}.$$
When $`K`$ is convex , then $`\left[T_K(x)\right]^{}`$ coincides with the the usual normal cone $`N_K(x)`$ of convex analysis.
The epigraph $`pi(W)`$ of a function $`W:IR^nIR\{+\mathrm{}\}`$ is defined by
$$pi(W):=\{(x,r)IR^n\times IR:rW(x)\}.$$
We shall need the following version of Rockafellar’s result (see ).
###### Lemma 2.1
Let $`x\text{dom}(W)`$ and let $`(p,0)\left[T_{pi(W)}(x,W(x))\right]^{}`$ be such that $`p0`$. Then there exist $`x_\epsilon `$ converging to $`x`$ (as $`\epsilon 0+`$) and
$$(p_\epsilon ,q_\epsilon )\left[T_{pi(W)}(x_\epsilon ,W(x_\epsilon ))\right]^{}$$
converging to $`(p,0)`$ as $`\epsilon 0+`$ such that $`q_\epsilon <0`$ for every $`\epsilon >0`$.
A closed subset $`K`$ of $`IR^n`$ is called a viability domain of a set-valued map $`G:KIR^n`$ if for every $`xK`$
$$G(x)T_K(x)\mathrm{}.$$
The following statement summarizes several versions of the viability theorem (see ).
###### Theorem 2.2 (Viability)
Let $`KIR^n`$ be a closed set and let $`G:KIR^n`$ be an upper semicontinuous set-valued map with compact convex values. The following conditions are equivalent:
$`K`$ is a viability domain of $`G`$;
$`G(x)\overline{co}T_K(x)\mathrm{}`$ for every $`xK`$;
for every $`xK`$ there exist $`\epsilon >0`$ and a solution $`y:[0,\epsilon [K`$ to the Cauchy problem
(2.1)
$$\{\begin{array}{cc}y^{}(t)G(y(t)),\hfill & \\ y(0)=x.\hfill & \end{array}$$
The equivalence (a) $``$ (b) was proved in . This proof was simplified in \[2, page 85\]. The fact that (a) $``$ (c) was first proved by Bebernes and Schuur in . A proof can be found in or .
The next theorem allows to deal with some unbounded set-valued maps with closed convex values. As usual, $`B_R`$ denotes the closed ball with centre $`0`$ and radius $`R`$.
###### Theorem 2.3
Let $`KIR^n`$ be a closed set and let $`G:KIR^n`$ be an upper semicontinuous set-valued map with closed convex values. We assume that for every $`xK`$ there exists $`R>0`$ such that for all small $`h>0`$
$$\text{dist}(x+hG(x),K)=\text{dist}(x+h(G(x)B_R),K).$$
Then the following statements are equivalent:
$`K`$ is a viability domain of $`G`$;
for every $`xK`$ and for every $`p\left[T_K(x)\right]^{}`$ we have $`\underset{uG(x)}{inf}p,u0`$.
Theorem 2.3 is a direct consequence of the following more technical result, which will be crucial in the proof of Theorem 1.7.
###### Theorem 2.4
Let $`KIR^n`$ be a closed set, let $`xK`$, and let $`G:KIR^n`$ be a set-valued map with non-empty closed convex values. Assume that there exists $`R>0`$ such that for all small $`h>0`$
(2.2)
$$\text{dist}(x+hG(x),K)=\text{dist}(x+h(G(x)B_R),K).$$
Assume also that the support function, defined by
$$\sigma (y,p):=\underset{uG(y)}{sup}p,u,$$
satisfies the following upper semicontinuity condition at $`x`$: for every $`\epsilon >0`$ there exists a neighbourhood $`U`$ of $`x`$ such that
(2.3)
$$\sigma (x,p)+\epsilon |p|>\sigma (y,p)$$
for every $`yUK`$ and for every $`p`$ in the set
$$N_{G(x)}(B_R):=\{pIR^n:uG(x)B_R,pN_{G(x)}(u)\}.$$
Finally, assume that for every $`yK`$ in a suitable neighbourhood of $`x`$ we have
(2.4)
$$\sigma (y,p)0$$
for every $`p\left[T_K(y)\right]^{}`$ such that $`pN_{G(x)}(B_R)`$. Then $`G(x)T_K(x)B_R\mathrm{}`$.
From the proof given below it follows that the same result holds if $`\left[T_K(y)\right]^{}`$ is replaced by the set of proximal normals to $`K`$ at $`y`$.
Proof — Let us define the function $`g:IR_+IR_+`$ by
$$g(h):=\frac{1}{2}\text{dist}(x+hG(x),K)^2.$$
Observe that $`g`$ is locally Lipschitz around zero and $`g(0)=0`$. For all small $`h>0`$ let us consider $`u_hG(x)B_R`$ and $`x_hK`$ such that $`\text{dist}(x+hG(x),K)=|x+hu_hx_h|`$. Then $`x_hx`$ when $`h0+`$.
Consider $`h>0`$ such that $`g^{}(h)`$ does exist and fix any $`uG(x)`$. Since $`G(x)`$ is convex, for all nonnegative $`h,k`$ we have $`(h+k)G(x)=hG(x)+kG(x)`$. Therefore
$$g^{}(h)\frac{1}{2}\underset{k0+}{lim}\frac{|x+hu_h+kux_h|^2|x+hu_hx_h|^2}{k}=p_h,u,$$
where $`p_h:=x+hu_hx_h`$. Consequently
(2.5)
$$g^{}(h)\underset{uG(x)}{inf}p_h,u=\underset{uG(x)}{sup}p_h,u=\sigma (x,p_h).$$
As $`x_h`$ is a point of $`K`$ with minimum distance from $`x+hu_h`$, the vector $`p_h`$ is a proximal normal to $`K`$, therefore it belongs to $`\left[T_K(x_h)\right]^{}`$. On the other hand, $`u_h`$ is the point of $`G(x)`$ with minimum distance from $`(x_hx)/h`$. Thus $`p_hN_{G(x)}(u_h)N_{G(x)}(B_R)`$. By (2.4) we have
(2.6)
$$\sigma (x_h,p_h)0.$$
By the uniform upper semicontinuity (2.3) of $`\sigma `$ for every $`\epsilon >0`$ there exists $`h_\epsilon >0`$ such that for $`0<h<h_\epsilon `$
(2.7)
$$\sigma (x,p_h)\sigma (x_h,p_h)+|p_h|\epsilon \sigma (x_h,p_h)+hR\epsilon ,$$
where the last inequality follows from the fact that $`|p_h|hR`$. From (2.5), (2.6), and (2.7) we obtain that $`g^{}(h)hR\epsilon `$ for every $`h(0,h_\epsilon )`$ at which the derivative $`g^{}(h)`$ exists.
Integrating $`g^{}`$ we deduce that for $`0<h<h_\epsilon `$
$$g(h)\frac{Rh^2}{2}\epsilon .$$
This implies that for $`0<h<h_\epsilon `$
$$\frac{\text{dist}(x+hu_h,K)}{h}=\frac{\text{dist}(x+hG(x),K)}{h}\sqrt{R\epsilon }.$$
Let $`h_i0+`$ be a sequence such that $`u_{h_i}`$ converges to some $`uG(x)`$. From the very definition of contingent cone we deduce that $`uT_K(x)`$. $`\mathrm{}`$
## 3 Proof of the comparison theorem
This section is devoted to the proof of the new comparison theorem for viscosity supersolutions.
Proof of Theorem 1.7 — We first claim that $`H`$ is locally bounded. Indeed for all $`xIR^n`$ we have $`H(x,p)L(x,0)`$, thus $`H`$ is locally bounded from below. On the other hand, $`H(x,p)\mathrm{\Theta }^{}(p)`$, where $`\mathrm{\Theta }^{}`$ denotes the Legendre-Fenchel conjugate of $`\mathrm{\Theta }`$. Since the function $`\mathrm{\Theta }`$ has a superlinear growth, the convex function $`\mathrm{\Theta }^{}`$ takes only finite values, so it is locally bounded. This shows that $`H`$ is also locally bounded from above.
Consequently, the function $`H(x,)`$ is locally Lipschitz with respect to $`p`$, locally uniformly with respect to $`x`$. By (1.2) $`H(,p)`$ is upper semicontinuous with respect to $`x`$. These two properties together imply that for every $`xIR^n`$, for every $`M>0`$, and for every $`\epsilon >0`$ there exists a neighbourhood $`U`$ of $`x`$ such that
(3.1)
$$yU,pB_M,H(x,p)+\epsilon >H(y,p).$$
Let us define the set-valued map $`G:IR^nIR\times IR^n\times IR`$ with closed convex values by
$$G(x):=\{(1,u,L(x,u)\rho ):\rho 0,uIR^n\}.$$
Let $`K:=pi(W)`$, let $`(t,x)\text{dom}(W)`$ with $`t>0`$, and let $`z:=(t,x,W(t,x))`$. We want to show that all assumptions of Theorem 2.4 are satisfied (here $`z`$ plays the role of $`x`$). To prove (2.2) we show that there exists $`R>0`$ such that for all small $`h>0`$
$`uIR^n,|(1,u,L(x,u))|R,`$
$`\text{dist}((t,x,W(t,x))+h(1,u,L(x,u)),K)=\text{dist}((t,x,W(t,x))+hG(x),K).`$
As $`W0`$, it is easy to see that for every $`h>0`$ there exist $`u_hIR^n`$ and $`z_h=(t_h,x_h,r_h)IR\times IR^n\times IR`$, with $`r_hW(t_h,x_h)`$, such that
$`|(t,x,W(t,x))+h(1,u_h,L(x,u_h))(t_h,x_h,r_h)|`$
$`=\text{dist}((t,x,W(t,x))+hG(x),K)h(L(x,0)+1).`$
If $`W(t,x)hL(x,u_h)r_h>0`$, by increasing $`r_h`$ we could make $`W(t,x)hL(x,u_h)r_h=0`$, which would contradict the definition of distance. Therefore $`W(t,x)hL(x,u_h)r_h0`$, hence
$$hL(x,u_h)+W(t_h,x_h)W(t,x)hL(x,u_h)+r_hW(t,x)h(L(x,0)+1).$$
We claim that $`u_h`$ is bounded for all $`h>0`$ small enough. Assume by contradiction that there exist $`h_i0+`$ such that $`|u_{h_i}|\mathrm{}`$.
Case 1. Assume first that for a subsequence, still denoted by $`h_i`$, we have $`h_i|u_{h_i}|c`$ for some $`c>0`$. Since $`W(t_h,x_h)0`$, we have
$$h_iL(x,u_{h_i})W(t,x)h_i(L(x,0)+1);$$
dividing by $`h_i|u_{h_i}|`$ and taking the limit we get
$$\underset{i\mathrm{}}{lim\; sup}\frac{L(x,u_{h_i})}{|u_{h_i}|}<+\mathrm{},$$
which contradicts (1.2).
Case 2. It remains to consider the case $`h_i|u_{h_i}|0`$. Since
$$\mathrm{max}\{|t_h+ht|,|x_hxhu_h|\}h(L(x,0)+1),$$
we deduce that for some $`vIR^n`$ and for a subsequence, still denoted by $`h_i`$, we have
$$\underset{i\mathrm{}}{lim}\frac{t_{h_i}t}{h_i|u_{h_i}|}=0\text{and}\underset{i\mathrm{}}{lim}\frac{x_{h_i}x}{h_i|u_{h_i}|}=v.$$
Furthermore,
$$h_iL(x,u_{h_i})+W(t_{h_i},x_{h_i})W(t,x)h_i(L(x,0)+1);$$
dividing by $`h_i|u_{h_i}|`$ and taking the limit yields
$$\underset{i\mathrm{}}{lim}\frac{W(t_{h_i},x_{h_i})W(t,x)}{h_i|u_{h_i}|}=\mathrm{}.$$
Hence $`D_{}W(t,x)(0,v)=\mathrm{}`$, which implies $`D_{}W(t,x)(0,0)=\mathrm{}`$ (see ). This contradicts (1.9) and completes the proof of our claim.
Thus $`|u_h|`$ is uniformly bounded when $`h>0`$ is small. As $`L`$ is locally bounded, there exists $`R>0`$ such that $`|(1,u_h,L(x,u_h))|R`$ for all small $`h>0`$.
Observe next that, if $`(p_t,p_x,q)N_{G(x)}(B_R)`$, then $`q0`$. Moreover, if $`q=0`$, then $`p_x=0`$; if $`q>0`$, then there exists $`uB_R`$ such that $`p_x/q_uL(x,u)`$, the subdifferential of $`L(x,)`$ at $`u`$. As $`L`$ is locally bounded, this implies that there exists a constant $`M`$ such that $`p_x/qB_M`$ for every $`(p_t,p_x,q)N_{G(x)}(B_R)`$ with $`q>0`$.
To prove (2.3) it is enough to show that for every $`\epsilon >0`$ there exists a neighbourhood $`U`$ of $`x`$ such that
(3.2)
$$\underset{uIR^n}{sup}(p_t+p_x,uqL(x,u))+\epsilon |(p_t,p_x,q)|$$
$$>\underset{uIR^n}{sup}(p_t+p_x,uqL(y,u))$$
for every $`yU`$ and for every $`(p_t,p_x,q)N_{G(x)}(B_R)`$. If $`q=0`$, then $`p_x=0`$ and (3.2) is trivial. If $`q>0`$, then (3.2) can be written as
$$p_t+qH(x,\frac{p_x}{q})+\epsilon |(p_t,p_x,q)|>p_t+qH(y,\frac{p_x}{q}),$$
which follows easily from (3.1).
Let us check that (2.4) holds true. Fix $`(s,y,r)pi(W)`$, with $`s>0`$, and $`(p_t,p_x,q)\left[T_{pi(W)}(s,y,r)\right]^{}`$. Since $`(0,0,1)T_{pi(W)}(s,y,r)`$, we have $`q0`$. Therefore (2.4) is equivalent to
(3.3)
$$\underset{uIR^n}{sup}(p_t+p_x,u+qL(y,u))0.$$
If $`q<0`$, then $`(p_t/|q|,p_x/|q|,1)\left[T_{pi(W)}(s,y,r)\right]^{}`$, hence $`(p_t/|q|,p_x/|q|)_{}W(s,y)`$ (see \[5, page 249\]) and we deduce (3.3) from (1.5). If $`q=0`$ and $`p_x0`$, then the supremum in (3.3) is $`+\mathrm{}`$. If $`q=0`$ and $`p_x=0`$, then $`(p_t,0,0)\left[T_{pi(W)}(s,y,r)\right]^{}`$. By (1.10) there exists $`uIR^n`$ such that $`z:=D_{}W(s,y)(1,u)<+\mathrm{}`$. As $`pi(D_{}W(s,y)(,))=T_{pi(W)}(s,y,W(s,y))`$, the vector $`(1,u,z)`$ belongs to $`T_{pi(W)}(s,y,W(s,y))`$, which is contained in $`T_{pi(W)}(s,y,r)`$. By the definition of $`\left[T_{pi(W)}(s,y,r)\right]^{}`$ we obtain $`p_t0`$, which yields (3.3) when $`q=0`$ and $`p_x=0`$.
From Theorem 2.4 we deduce that
$$G(x)\left(T_{pi(W)}(t,x,W(t,x))\right)\mathrm{}.$$
As $`pi(D_{}W(t,x)(,))=T_{pi(W)}(t,x,W(t,x))`$, we obtain that
$$uIR^n,D_{}W(t,x)(1,u)L(x,u).$$
This and Theorem 1.6 imply that $`WV`$ on $`IR_+\times IR^n`$. $`\mathrm{}`$
###### Remark 3.1
From the proof of Theorem 1.7 we see that condition (1.10) yields, for $`t>0`$,
$$(p_t,0,0)\left[T_{pi(W)}(t,x,W(t,x))\right]^{}p_t0.$$
If, in addition, the subsolution inequality (1.12) is satisfied, and
$$(p_t,0,0)\left[T_{pi(W)}(t,x,W(t,x))\right]^{},$$
then by Rockafellar’s Lemma 2.1 there exist $`(t_i,x_i)(t,x)`$ and $`(p_t^i,p_x^i,q_i)\left[T_{pi(W)}(t_i,x_i,W(t_i,x_i))\right]^{}`$, with $`q_i<0`$, such that $`(p_t^i,p_x^i,q_i)(p_t,0,0)`$. Then $`(p_t^i/q_i,p_x^i/q_i)_{}W(t_i,x_i)`$ and from (1.12) we obtain $`p_t^i/q_iL(x_i,0)0`$, hence $`p_t^i+q_iL(x_i,0)0`$. Taking the limit we get $`p_t0`$. This shows that (1.10) and (1.12) together imply the following geometric condition for $`t>0`$:
$$(p_t,0,0)\left[T_{pi(W)}(t,x,W(t,x))\right]^{}p_t=0.$$
## 4 Proofs of the uniqueness results
We begin with the proof of the uniqueness theorem for viscsity solutions.
Proof of Theorem 1.2 — By Theorem 1.7 we have $`WV`$. Recall that the hypograph of $`W`$ is defined by
$$yp(W):=\{(t,x,r)IR_+^{}\times IR^n\times IR:rW(t,x)\}.$$
Define the closed set
$$K:=yp(W)(IR_{}\times IR^n\times IR).$$
We claim that for all $`(t,x,r)K`$
(4.1)
$$uIR^n,(1,u,L(x,u))\overline{co}T_K(t,x,r).$$
It is enough to prove it in the case $`t>0`$ and $`r=W(t,x)`$, since the other cases are evident. Fix $`uIR^n`$. Then by (1.6)
(4.2)
$$(p_t,p_x)_+W(t,x),p_t+p_x,uL(x,u)0.$$
We want to prove that
(4.3)
$$(p_t,p_x,q)\left[T_{yp(W)}(t,x,W(t,x))\right]^{}p_t+p_x,uqL(x,u)0.$$
When $`q>0`$ we have $`(p_t/q,p_x/q)_+W(t,x)`$, thus (4.3) follows from (4.2). By Lemma 2.1, applied to $`W`$, if $`(0,0,0)(p_t,p_x,0)\left[T_{yp(W)}(t,x,W(t,x))\right]^{}`$, then for some $`(t_i,x_i)(t,x)`$ and $`(p_t^i,p_x^i,q_i)\left[T_{yp(W)}(t_i,x_i,W(t_i,x_i))\right]^{}`$, with $`q_i>0`$, we have $`(p_t^i,p_x^i,q_i)(p_t,p_x,0)`$. So
$$p_t^i+p_x^i,uq_iL(x_i,u)0.$$
Taking the limit we get $`p_t+p_x,u0`$, which concludes the proof of (4.3).
By the separation theorem, (4.1) follows from (4.3). Since the lower set-valued limit of contingent cones is equal to Clarke’s tangent cone (see for instance ), from the continuity of $`L`$ we deduce that, for all $`(t,x)IR_+^{}\times IR^n,`$
$$(1,u,L(x,u))C_{yp(W)}(t,x,W(t,x)).$$
Fix $`\epsilon >0`$. Then it is not difficult to check that
$$uIR^n,(1,u,L(x,u)\epsilon )\text{Int}\left(C_{yp(W)}(t,x,W(t,x))\right).$$
By \[4, Proposition 13, p. 425\] this yields
(4.4)
$$\underset{h0+,vu}{lim\; inf}\frac{W(th,x+hv)W(t,x)}{h}L(x,u)\epsilon .$$
We have to show that for all $`t>0,xIR^n`$, $`V(t,x)W(t,x)`$. Let $`y`$ be a minimizer of the Bolza problem (1.1). Since $`L`$ is continuous, $`y^{}L^{\mathrm{}}(0,t;IR^n)`$ by . Consider a sequence of continuous functions $`u_i:[0,t]IR^n`$ which is bounded in $`L^{\mathrm{}}(0,t;IR^n)`$ and converges to $`y^{}`$ almost everywhere in $`[0,t]`$, and let $`t_i0+`$ and $`x_iy(t)`$ be such that $`W(t_i,x_i)\phi (y(t))`$. Define
$$y_i(s):=x_i_s^{tt_i}u_i(\tau )𝑑\tau ,s[0,tt_i].$$
and $`y_i(s):=x_i`$ for $`s>tt_i`$. Then $`y_i`$ converges to $`y`$ uniformly in $`[0,t]`$. Fix $`i`$ and set $`\psi (s):=W(ts,y_i(s))`$ for $`s[0,tt_i]`$. By (4.4) for every $`s[0,tt_i[`$ we have
(4.5)
$$\underset{h0+}{lim\; sup}\frac{\psi (s+h)\psi (s)}{h}L(y_i(s),u_i(s))\epsilon .$$
Consider the system
$$\{\begin{array}{cc}(\alpha ^{}(s),z^{}(s))=(1,L(y_i(s),u_i(s))\epsilon ),s0\hfill & \\ & \\ (\alpha (0),z(0))=(0,W(t,y_i(0))).\hfill & \end{array}$$
where we have set $`u_i(s):=u_i(t)`$ for all $`stt_i`$. It has the unique solution
$$(\alpha (s),z(s)):=(s,W(t,y_i(0))_0^sL(y_i(\tau ),u_i(\tau ))𝑑\tau \epsilon s).$$
According to Theorem 2.2 and (4.5), this solution is viable in $`yp(\psi )([tt_i,+\mathrm{}[\times IR)`$, i.e., for all $`s[0,tt_i]`$ we have $`(\alpha (s),z(s))yp(\psi )`$. Thus for all $`s[0,tt_i]`$
$$W(ts,y_i(s))W(t,y_i(0))_0^sL(y_i(\tau ),u_i(\tau ))𝑑\tau \epsilon s.$$
In particular
$$W(t,y_i(0))W(t_i,x_i)+_0^{tt_i}L(y_i(\tau ),u_i(\tau ))𝑑\tau +\epsilon (tt_i).$$
Since the functions $`L`$ and $`W`$ are continuous, taking the limit we get
$$W(t,x)\epsilon t\phi (y(t))+_0^tL(y(\tau ),y^{}(\tau ))𝑑\tau V(t,x).$$
Finally, as $`\epsilon 0+`$ we obtain $`W(t,x)V(t,x)`$. $`\mathrm{}`$
Proof of Theorem 1.4 — The inequality $`WV`$ is proved in Theorem 1.7, while the inequality $`WV`$ follows from Theorem 1.5. $`\mathrm{}`$
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# 𝛾^∗𝑁→Δ transition form factors: a new analysis of the JLab data on 𝑝(𝑒,𝑒'𝑝)𝜋⁰ at Q2=2.8 and 4.0 (GeV/c)2
\[
## Abstract
Recent JLab data of the differential cross section for the reaction $`p(e,e^{}p)\pi ^0`$ in the invariant mass region of $`1.1<W<1.4`$ GeV at four-momentum transfer squared $`Q^2=`$ 2.8 and 4.0 (GeV/c)<sup>2</sup> are analyzed with two models, both of which give an excellent description of most of the existing pion electroproduction data below $`W<1.5`$ GeV. We find that at up to $`Q^2=4.0`$(GeV/c)<sup>2</sup>, the extracted helicity amplitudes $`A_{3/2}`$ and $`A_{/2}`$ remain comparable with each other, implying that hadronic helicity is not conserved at this range of $`Q^2`$. The ratios $`E_{1+}/M_{1+}`$ obtained show, starting from a small and negative value at the real photon point, a clear tendency to cross zero, and to become positive with increasing $`Q^2`$. This is a possible indication of a very slow approach toward the pQCD region. Furthermore, we find that the helicity amplitude $`A_{1/2}`$ and $`S_{1/2}`$, but not $`A_{3/2}`$, starts exhibiting the scaling behavior at about $`Q^22.5`$(GeV/c)<sup>2</sup>.
\]
In a recent experiment , electro-excitation of the $`\mathrm{\Delta }`$ was studied at $`Q^2=`$ 2.8 and 4.0 (GeV/c)<sup>2</sup> via the reaction $`p(e,e^{}p)\pi ^0`$. It was motivated by the possibility of determining the range of momentum transfers where perturbative QCD (pQCD) would become applicable. In the limit of $`Q^2\mathrm{}`$, pQCD predicts the dominance of helicity-conserving amplitudes and scaling results . The hadronic helicity conservation should have the consequence that the ratio between magnetic dipole $`M_{1+}^{(3/2)}`$ and electric quadrupole $`E_{1+}^{(3/2)}`$ multipoles, $`R_{EM}=E_{1+}^{(3/2)}/M_{1+}^{(3/2)}`$, approaches 1. The scaling behavior predicted by pQCD for the helicity amplitudes is $`A_{1/2}^\mathrm{\Delta }Q^3`$, $`A_{3/2}^\mathrm{\Delta }Q^5`$, and the Coulomb helicity amplitude $`S_{1/2}^\mathrm{\Delta }Q^3`$, resulting in $`R_{SM}=S_{1+}^{(3/2)}/M_{1+}^{(3/2)}const`$. On the other hand, in symmetric $`SU(6)`$ quark models, the $`\gamma N\mathrm{\Delta }`$ transition can proceed only via the flip of a single quark spin in the nucleon, leading to $`M_{1+}`$ dominance and $`E_{1+}=S_{1+}0`$. Recent experiments give nonvanishing ratios $`R_{EM}`$ lying between $`2.5\%`$ and $`3.0\%`$ at $`Q^2=0`$. This has been widely taken as an indication of a deformed $`\mathrm{\Delta }`$, namely, an admixture of a D state in the $`\mathrm{\Delta }`$. Accordingly, the question of how $`R_{EM}`$ would evolve from a very small negative value at $`Q^2=0`$ to $`+100\%`$ at sufficiently high $`Q^2`$, has attracted great interest both theoretically and experimentally.
In Ref. , the differential cross sections were measured in the invariant mass region of $`\mathrm{\hspace{0.17em}\hspace{0.17em}1.1}<W<1.4`$ GeV. Two methods were used to extract the contributing multipoles. The first one, which is model and energy independent, consisted of making approximate multipole fits to angular distributions independently at each $`W`$, assuming $`M_{1+}`$ dominance, and only $`S`$ and $`P`$ wave contributions . Another extraction of the resonance amplitudes was performed using the effective Lagrangian method . In this model-dependent analysis, the resonant multipoles are expressed as a sum of background and resonance amplitudes, both prescribed by an effective Lagrangian, and unitarized with the K-matrix method. The parameters in the model were fitted to data points with energy $`W`$ only up to $`1.31`$ GeV. The ratios $`R_{EM}`$ and $`R_{SM}`$ extracted with these two methods are both small, negative, and tending to more negative values with increasing $`Q^2`$, indicating that pQCD is not yet applicable in this region of $`Q^2`$. Recently, it was shown that the $`Q^2`$-dependence of the ratios $`R_{EM}`$ and $`R_{SM}`$ extracted in Ref. can be explained in a dynamical model for electromagnetic production of pions, together with a simple scaling assumption for the bare $`\gamma ^{}N\mathrm{\Delta }`$ form factors.
Because of the significance of the physics involved in the $`Q^2`$ evolution of $`R_{EM}`$ and $`R_{SM}`$, it is important to employ the best possible extraction method in the analysis of the data. In fact, the values of $`R_{EM}`$ and $`R_{SM}`$ extracted with the two methods used in Ref. differ from each other by factors of 2 and 1.5 at $`Q^2=`$ 2.8 and 4.0 (GeV/c)<sup>2</sup>, respectively. In this letter, we present the results of a new analysis of the data of Ref. , using a new version (hereafter called MAID) of the unitary isobar model developed at Mainz (hereafter called MAID98) , and the dynamical model developed recently in Ref. , which both give excellent descriptions of most of the existing pion photo- and electroproduction data . Our analysis is similar to the second method used in in the sense that it also makes use of a model. However, we fit all the data points measured up to $`W=`$ 1.4 GeV and obtain smaller values of $`\chi ^2`$ per d.o.f.
In the dynamical approach to pion photo- and electroproduction , the t-matrix can be expressed as $`t_{\gamma \pi }(E)=v_{\gamma \pi }+v_{\gamma \pi }g_0(E)t_{\pi N}(E)`$ and the physical multipoles in channel $`\alpha `$ are given by
$`t_{\gamma \pi }^{(\alpha )}(q_E,k)=\mathrm{exp}(i\delta ^{(\alpha )})\mathrm{cos}\delta ^{(\alpha )}`$ (1)
$`\times `$ $`\left[v_{\gamma \pi }^{(\alpha )}(q_E,k)+P{\displaystyle _0^{\mathrm{}}}𝑑q^{}{\displaystyle \frac{q^2R_{\pi N}^{(\alpha )}(q_E,q^{})v_{\gamma \pi }^{(\alpha )}(q^{},k)}{EE_{\pi N}(q^{})}}\right],`$ (2)
where $`v_{\gamma \pi }`$ is the transition potential for $`\gamma ^{}N\pi N`$, and $`t_{\pi N}`$ and $`g_0`$ denote the $`\pi N`$ t-matrix and free propagator, respectively, with $`EW`$ the total energy in the CM frame. $`\delta ^{(\alpha )}`$ and $`R_{\pi N}^{(\alpha )}`$ are the $`\pi N`$ scattering phase shift and reaction matrix in channel $`\alpha `$, respectively; $`q_E`$ is the pion on-shell momentum and $`k=|𝐤|`$ is the photon momentum.
In a resonant channel like (3,3) in which the $`\mathrm{\Delta }(1232)`$ plays a dominant role, the transition potential $`v_{\gamma \pi }`$ consists of two terms, $`v_{\gamma \pi }(E)=v_{\gamma \pi }^B+v_{\gamma \pi }^\mathrm{\Delta }(E),`$ where $`v_{\gamma \pi }^B`$ is the background transition potential and $`v_{\gamma \pi }^\mathrm{\Delta }(E)`$ corresponds to the contribution of the bare $`\mathrm{\Delta }`$. The resulting t-matrix can be decomposed into two terms $`t_{\gamma \pi }(E)=t_{\gamma \pi }^B+t_{\gamma \pi }^\mathrm{\Delta }(E)`$, where $`t_{\gamma \pi }^B(E)=v_{\gamma \pi }^B+v_{\gamma \pi }^Bg_0(E)t_{\pi N}(E),`$ and $`t_{\gamma \pi }^\mathrm{\Delta }(E)=v_{\gamma \pi }^\mathrm{\Delta }+v_{\gamma \pi }^\mathrm{\Delta }g_0(E)t_{\pi N}(E).`$ Here $`t_{\gamma \pi }^B`$ includes the contributions from the nonresonant background and renormalization of the vertex $`\gamma ^{}N\mathrm{\Delta }`$. The advantage of such a decomposition is that all the processes which start with the excitation of the bare $`\mathrm{\Delta }`$ are summed up in $`t_{\gamma \pi }^\mathrm{\Delta }`$. Note that the multipole decomposition of both $`t_{\gamma \pi }^B`$ and $`t_{\gamma \pi }^\mathrm{\Delta }`$ would take the same form as Eq. (2).
For a correct description of the resonance contributions we need, first of all, a reliable description of the nonresonant part of the amplitude. In MAID98, the background contribution was described by Born terms obtained with an energy dependent mixing of pseudovector-pseudoscalar $`\pi NN`$ coupling and t-channel vector meson exchanges, namely, $`t_{\gamma \pi }^{B,\alpha }(\mathrm{MAID98})=v_{\gamma \pi }^{B,\alpha }(W,Q^2)`$. The mixing parameters and coupling constants were determined from an analysis of nonresonant multipoles in the appropriate energy regions. In the new version of MAID, the $`S`$, $`P`$, $`D`$ and $`F`$ waves of the background contributions are complex numbers defined in accordance with the K-matrix approximation,
$$t_{\gamma \pi }^{B,\alpha }(\mathrm{MAID})=\mathrm{exp}(i\delta ^{(\alpha )})\mathrm{cos}\delta ^{(\alpha )}v_{\gamma \pi }^{B,\alpha }(W,Q^2).$$
(3)
From Eqs. (2) and (3), one finds that the difference between the background terms of MAID and of the dynamical model is that off-shell rescattering contributions (principal value integral) are not included in MAID. To take account of the inelastic effects at the higher energies, we replace $`\mathrm{exp}i(\delta ^{(\alpha )})\mathrm{cos}\delta ^{(\alpha )}=\frac{1}{2}[\mathrm{exp}(2i\delta ^{(\alpha )})+1]`$ in Eqs. (2) and (3) by $`\frac{1}{2}[\eta _\alpha \mathrm{exp}(2i\delta ^{(\alpha )})+1]`$, where $`\eta _\alpha `$ is the inelasticity. In our actual calculations, both the $`\pi N`$ phase shifts $`\delta ^{(\alpha )}`$ and inelasticity parameters $`\eta _\alpha `$ are taken from the analysis of the GWU group . Furthermore, the off-shell rescattering effects in the dynamical model are evaluated with the reaction matrix $`R_{\pi N}^{(\alpha )}(q_E,q^{})`$ as prescribed by a meson exchange model .
Following Ref. , we assume a Breit-Wigner form for the resonance contribution $`𝒜_\alpha ^R(W,Q^2)`$ to the total multipole amplitude,
$$𝒜_\alpha ^R(W,Q^2)=\overline{𝒜}_\alpha ^R(Q^2)\frac{f_{\gamma R}(W)\mathrm{\Gamma }_RM_Rf_{\pi R}(W)}{M_R^2W^2iM_R\mathrm{\Gamma }_R}e^{i\varphi },$$
(4)
where $`f_{\pi R}`$ is the usual Breit-Wigner factor describing the decay of a resonance $`R`$ with total width $`\mathrm{\Gamma }_R(W)`$ and physical mass $`M_R`$. The expressions for $`f_{\gamma R},f_{\pi R}`$ and $`\mathrm{\Gamma }_R`$ are given in Ref. . The phase $`\varphi (W)`$ in Eq. (4) is introduced to adjust the phase of the total multipole to equal the corresponding $`\pi N`$ phase shift $`\delta ^{(\alpha )}`$. Because $`\varphi =0`$ at resonance, $`W=M_R`$, this phase does not affect the $`Q^2`$ dependence of the $`\gamma NR`$ vertex.
We now concentrate on the $`\mathrm{\Delta }(1232)`$. In this case the magnetic dipole $`(\overline{𝒜}_M^\mathrm{\Delta })`$ and the electric quadrupole $`(\overline{𝒜}_E^\mathrm{\Delta })`$ form factors are related to the conventional electromagnetic helicity amplitudes $`A_{1/2}^\mathrm{\Delta }`$, $`A_{3/2}^\mathrm{\Delta }`$ and $`S_{1/2}^\mathrm{\Delta }`$ by
$`\overline{𝒜}_M^\mathrm{\Delta }(Q^2)`$ $`=`$ $`{\displaystyle \frac{1}{2}}(A_{1/2}^\mathrm{\Delta }+\sqrt{3}A_{3/2}^\mathrm{\Delta }),`$ (5)
$`\overline{𝒜}_E^\mathrm{\Delta }(Q^2)`$ $`=`$ $`{\displaystyle \frac{1}{2}}(A_{1/2}^\mathrm{\Delta }+{\displaystyle \frac{1}{\sqrt{3}}}A_{3/2}^\mathrm{\Delta }),`$ (6)
$`\overline{𝒜}_S^\mathrm{\Delta }(Q^2))`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}S_{1/2}^\mathrm{\Delta },`$ (7)
where $`k^2=Q^2+[(W^2m_N^2Q^2)/2W]^2`$. We stress that the physical meaning of these resonant amplitudes in different models is different. In MAID, they contain contributions from the background excitation and describe the so called ”dressed” $`\gamma N\mathrm{\Delta }`$ vertex. However, in the dynamical model the background excitation is included in $`𝒜_\alpha ^B`$ and the electromagnetic vertex $`\overline{𝒜}_\alpha ^\mathrm{\Delta }(Q^2)`$ corresponds to the ”bare” vertex.
In the dynamical model of Ref. , a scaling assumption was made concerning the (bare) form factors $`\overline{𝒜}_\alpha ^\mathrm{\Delta }(Q^2)`$, namely, that all of them have the same $`Q^2`$ dependence. In the present analysis, we do not impose the scaling assumption and write, for electric ($`\alpha =E`$), magnetic ($`\alpha =M`$) and Coulomb ($`\alpha =S`$) multipoles,
$`\overline{𝒜}_\alpha ^\mathrm{\Delta }(Q^2)=X_\alpha ^\mathrm{\Delta }(Q^2)\overline{𝒜}_\alpha ^\mathrm{\Delta }(0){\displaystyle \frac{k}{k_W}}F(Q^2),`$ (8)
where $`k_W=(W^2m_N^2)/2W`$. The form factor $`F`$ is taken to be $`F(Q^2)=(1+\beta Q^2)e^{\gamma Q^2}G_D(Q^2),`$ where $`G_D(Q^2)=1/(1+Q^2/0.71)^2`$ is the usual dipole form factor. The parameters $`\beta `$ and $`\gamma `$ were determined by setting $`X_M^\mathrm{\Delta }=1`$ and fitting $`\overline{𝒜}_M^\mathrm{\Delta }(Q^2)`$ to the data for $`G_M^{}`$ as defined in . The values of $`\overline{𝒜}_M^\mathrm{\Delta }(0)`$ and $`\overline{𝒜}_E^\mathrm{\Delta }(0)`$ were determined by fitting to the multipoles obtained in the recent analyses of the Mainz and GWU groups. Both $`X_E`$ and $`X_S`$ are to be determined by the experiment with $`X_\alpha ^\mathrm{\Delta }(0)=1`$. Note that deviations from $`X_\alpha ^\mathrm{\Delta }=1`$ value will indicate a violation of the scaling law. Similar treatment is also applied to the $`N^{}(1440)`$ resonance with two additional parameters $`X_M^{P11}`$ and $`X_S^{P11}`$ corresponding to the transverse and longitudinal resonance transitions in the isospin 1/2 channel.
The dynamical model and MAID are used to analyze the recent JLab differential cross section data on $`p(e,e^{}p)\pi ^0`$ at high $`Q^2`$. All measured data, 751 points at $`Q^2`$=2.8 and 867 points at $`Q^2`$=4.0 (GeV/c)<sup>2</sup> covering the entire energy range $`1.1<W<1.4`$ GeV, are included in our global fitting procedure. We obtain a very good fit to the measured differential cross sections. In fact, the values of $`\chi ^2/d.o.f.`$ for our two models are smaller than those obtained in Ref. (see Table 1). Our results for the $`G_M^{}`$ form factor are shown in Fig. 1. Here the best fit is obtained with $`\gamma =0.21`$ (GeV/c)<sup>-2</sup> and $`\beta =0`$ in the case of MAID, and $`\gamma =0.40`$ (GeV/c)<sup>-2</sup> and $`\beta =0.52`$(GeV/c)<sup>-2</sup> in the case of the dynamical model.
With the resonance parameters $`X_\alpha ^\mathrm{\Delta }(Q^2)`$ determined from the fit, the ratios $`R_{EM}=ImE_{1+}/ImM_{1+}`$ and $`R_{SM}=ImS_{1+}/ImM_{1+}`$ of the total multipoles and the helicity amplitudes $`A_{1/2}`$ and $`A_{3/2}`$ can then be calculated at resonance. We perform the calculations for both physical $`(p\pi ^0)`$ and isospin 3/2 channels and find them to agree with each other. The extracted $`Q^2`$ dependence of the $`X_\alpha ^\mathrm{\Delta }`$ parameters is: $`X_E^\mathrm{\Delta }(\mathrm{MAID})=1Q^2/3.7,X_E^\mathrm{\Delta }(\mathrm{DM})=1+Q^4/2.4`$, $`X_S^\mathrm{\Delta }(\mathrm{MAID})=1+Q^6/61,X_S^\mathrm{\Delta }(\mathrm{DM})=110Q^2`$, with $`Q^2`$ in units (GeV/c)<sup>2</sup>.
Our extracted values for $`R_{EM}`$ and $`R_{SM}`$ and a comparison with the results of Ref. are presented in Table 1 and shown in Fig. 2. The main difference between our results and those of Ref. is that our values of $`R_{EM}`$ show a clear tendency to cross zero and change sign as $`Q^2`$ increases. This is in contrast with the results obtained in the original analysis of the data which concluded that $`R_{EM}`$ would stay negative and tend toward more negative values with increasing $`Q^2`$. Furthermore, we find that the absolute value of $`R_{SM}`$ is strongly increasing. In terms of helicity amplitudes, our results for a small $`R_{EM}`$ can be understood in that the extracted $`A_{3/2}`$ remains as large as the helicity conserving $`A_{1/2}`$ up to $`Q^2=4.0(GeV/c)^2`$, resulting in a small $`E_{1^+}`$.
Finally, we show our results for $`Q^3A_{1/2}^\mathrm{\Delta },Q^5A_{3/2}^\mathrm{\Delta },`$ and $`Q^3S_{1/2}^\mathrm{\Delta }`$ in Fig. 3. The bare form factors obtained with DM is used since the scaling behavior predicted by pQCD arises from the $`3q`$ Fock states in the nucleon and $`\mathrm{\Delta }`$. It is interesting to see that $`S_{1/2}^\mathrm{\Delta }`$ and $`A_{1/2}^\mathrm{\Delta }`$ clearly starts exhibiting the pQCD scaling behavior at about $`Q^22.5(GeV/c)^2`$. The maximal value for the $`Q^3A_{1/2}^\mathrm{\Delta }`$ which we obtained in this region is about -0.11 GeV<sup>5/2</sup>. This is in between of the asymptotic values for this quantity discussed in Ref., i.e., $`Q^3A_{1/2}^\mathrm{\Delta }=`$ -0.08 GeV<sup>5/2</sup> and -0.17 GeV<sup>5/2</sup>. However, it is difficult to draw any definite conclusion for $`Q^5A_{3/2}^\mathrm{\Delta }`$. From these results, it appears likely that scaling will set in earlier than the helicity conservation as predicted by pQCD.
In summary, we have re-analyzed the recent JLab data for electroproduction of the $`\mathrm{\Delta }(1232)`$ resonance via $`p(e,e^{}p)\pi ^0`$ with two models for pion electroproduction, both of which give excellent descriptions of the existing data. We find that $`A_{3/2}^\mathrm{\Delta }`$ is still as large as $`A_{1/2}^\mathrm{\Delta }`$ at $`Q^2=4`$ (GeV/c)<sup>2</sup>, which implies that hadronic helicity conservation is not yet observed in this region of $`Q^2`$. Accordingly, our extracted values for $`R_{EM}`$ are still far from the pQCD predicted value of $`+100\%`$. However, in contrast to previous results we find that $`R_{EM}`$, starting from a small and negative value at the real photon point, actually exhibits a clear tendency to cross zero and change sign as $`Q^2`$ increases, while the absolute value of $`R_{SM}`$ is strongly increasing. In regard to the scaling, our analysis indicates that $`S_{1/2}^\mathrm{\Delta }`$ and $`A_{1/2}^\mathrm{\Delta }`$, but not $`A_{3/2}^\mathrm{\Delta }`$, starts exhibiting the pQCD scaling behavior at about $`Q^22.5(GeV/c)^2`$. It appears likely that the onset of scaling behavior might take place at a lower momentum transfer than that of hadron helicity conservation.
It will be most interesting to have data at yet higher momentum transfer in order to see the region where the helicity amplitude $`A_{1/2}^\mathrm{\Delta }`$ finally dominates over $`A_{3/2}^\mathrm{\Delta }`$. It is only there that we could expect to see the onset of the asymptotic behavior of $`R_{EM}+100\%`$ and $`R_{SM}const`$.
###### Acknowledgements.
We are grateful to Paul Stoler and Rick Davidson for useful communications. S.S.K. would like to thank the Department of Physics at National Taiwan University for warm hospitality and gratefully acknowledges the financial support of the National Science Council of ROC. This work was supported in part by NSC under Grant No. NSC89-2112-M002-038, by Deutsche Forschungsgemeinschaft (SFB443) and by a joint project NSC/DFG TAI-113/10/0.
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# Regge analysis of diffractive and leading baryon structure functions from DIS
## I Introduction
One of the most striking results obtained at the DESY HERA $`ep`$ collider was the discovery by the H1 and ZEUS collaborations H1-old ; ZEUS-old that deep inelastic scattering (DIS) events tagged with rapidity gaps exhibit mass distributions whose shape resemble very much those observed in hadron-hadron diffraction experiments. More recently, both the H1 and ZEUS collaborations reported H1-Lead ; ZEUS-Lead analyses of another class of DIS events whose pretty flat distribution turned out to be quite similar to the leading particle spectrum, also observed in hadron reactions. These similarities suggest that the Regge pole phenomenology Collins , successfully used to describe diffractive events and the leading particle effect in hadron processes Collins ; Mirian , might also be employed to analyze the corresponding events obtained in DIS.
In a conventional DIS process, $`epeX`$, a high energy electron of four-momentum $`k`$ interacts with a proton of four-momentum $`P`$ through the emission of a photon of virtuality $`Q^2`$. As long as the photon has high enough momentum, it can resolve the internal partonic structure of the proton, interacting with its partons through a hard scattering which breaks up the hadron. In this inclusive reaction only the outgoing electron is detected in the final state (Fig. 1a).
If, besides the electron, one specific kind of hadron is detected in the final state, we have a semi-inclusive DIS process, $`epehX`$. Among processes of this kind there are events for which it is possible to recognize, in the final hadronic state, particles that bear some identity with the original proton, i.e., they are close in rapidity to the original proton and carry a significant fraction of its momentum. In a particular case, events such as these may be characterized by a large rapidity gap between the products of the $`\gamma ^{}p`$ hard scattering and the outgoing proton debris (Fig. 1b). If those debris are identified with a proton, neutron or any other baryon close related to the original proton, we have the above mentioned leading baryon effect, $`\gamma ^{}pXN`$, which, in analogy with the hadron case Mirian , could, in principle, be described by Regge phenomenology in terms of reggeon and pion exchanges Szczurek .
Furthermore, if the detected baryon is carrying more than 90% of the incoming proton momentum and is identified with a proton itself (or, equivalently, if a rapidity gap is detected nearby the proton fragmentation region), then the dominant interaction mechanism is a single diffractive scattering, $`\gamma ^{}pXp`$, in which the virtual photon interacts with the proton through a color singlet exchange with the vacuum quantum numbers, which in Regge phenomenology is known as pomeron exchange Collins .
With the above statements we just intend to make the point that, speaking in terms of theory, diffractive DIS events are part of a wider class of interactions, the semi-inclusive DIS processes, within which the leading particle effect is found. Thus, if one wants to capture the Regge behavior presumably observed by a certain kind of DIS data, one should take into account all available data at once, which, in this case, means to consider simultaneously diffractive and leading particle data in the same analysis. This is the scope of the present paper.
Semi-inclusive processes have been measured by the H1 and ZEUS collaborations in the HERA $`ep`$ colliding machine at DESY, where positrons of $`27.5GeV`$ collide with protons of $`820GeV`$. The H1 Collaboration has made high-statistic measurements of the diffractive structure function $`F_2^D`$ in the process $`epeXY`$, where $`Y`$ represents a hadronic system with mass lower than $`1.6GeV`$ and rapidity closest to that of the incident proton H1-Diff . H1 also measured the leading proton and neutron structure functions, $`F_2^{LP}`$ and $`F_2^{LN}`$ respectively, in the reaction $`epeNX`$, where $`N`$ is the identified nucleon H1-Lead . The ZEUS Collaboration has measurements of the diffractive structure function $`F_2^D`$ in the reaction $`epepX`$ Zeus , and preliminary leading baryon measurements have also been reported ZEUS-Lead .
Now, let us examine these experimental findings through a phenomenological gaze. The first attempts to describe them by the Regge formalism were based on the Ingelman and Schlein model Ingelman by which diffraction in DIS is understood as a two-step process: first the proton emits a pomeron, then the pomeron is hard scattered by the virtual photon. In such a view, the pomeron is a quasi-particle that carries a fraction $`\xi `$ of the proton’s momentum and has its own structure function that could be expressed in terms of $`\beta `$ and $`Q^2`$ (here $`\beta `$ plays the role of the Bjorken variable for the pomeron; see its definition in the next section). Accordingly, the measured structure function $`F_2^{D(4)}(\xi ,t,\beta ,Q^2)`$ would be factorized as
$$F_2^{D(4)}(\xi ,t,\beta ,Q^2)=f_{𝙸𝙿}(\xi ,t)F_2^{𝙸𝙿}(\beta ,Q^2),$$
(1)
where $`f_{𝙸𝙿}(\xi ,t)`$ is the flux of the pomeron out of the proton, which is a function of $`\xi `$ and $`t`$, the squared four-momentum transferred at the proton vertex. $`F_2^{𝙸𝙿}(\beta ,Q^2)`$ represents here the pomeron structure function.
Several analyses were made based on Eq. (1) and on this factorization hypothesis, including those performed by the H1 and ZEUS collaborations H1-old ; ZEUS-old (see also Varios and references quoted therein). In fact, this kind of analyses has been used to establish the pomeron intercept $`\alpha _{𝙸𝙿}`$ from the diffractive DIS data.
Although the preliminary experimental results seemed to confirm the factorization hypothesis H1-old ; ZEUS-old , subsequent high-statistic data measured in an extended kinematical region by the H1 Collaboration proved that such a simple factorized expression is clearly violated H1-Diff . Since then it has been conjectured H1-Diff ; Golec-Biernat that secondary reggeonic exchanges could play an important role in diffractive events, in such a way that the structure function could be written as
$$F_2^{D(4)}(\xi ,t,\beta ,Q^2)=f_{𝙸𝙿}(\xi ,t)F_2^{𝙸𝙿}(\beta ,Q^2)+f_{𝙸𝚁}(\xi ,t)F_2^{𝙸𝚁}(\beta ,Q^2),$$
(2)
where $`f_{𝙸𝚁}(\xi ,t)`$ is the reggeon flux factor, and $`F_2^{𝙸𝚁}(\beta ,Q^2)`$ is the reggeon structure function. Within this approach, the change in the diffractive pattern displayed by the H1 data could be explained without giving up the idea of Regge factorization for each contribution. The H1 Collaboration itself was very successful in describing the bulk of the diffractive structure function data with a fitting expression akin to Eq. (2) (see H1-Diff ).
In fact, not only the diffractive data, but also the H1 leading proton structure function data can be fairly described within the same framework as well by just adding up to Eq. (2) an extra pion contribution as required in such a case (see H1-Lead ). The leading neutron structure function is described by the same scheme, but in that case only pion exchange is necessary H1-Lead .
Since the leading baryon data were obtained some time after the diffractive structure function measurements, these H1 analyses were performed independent of each other. However, as stated previously, it is our belief that both diffractive and leading proton processes should be analyzed together, as two parts of the same semi-inclusive process, in the same fashion as in the hadronic case Mirian . In this way it would be possible to establish more precisely the role of the pomeron and the secondary reggeon exchanges, since the diffractive data are dominated by the former and has the latter only as a background, while the reverse is true for the leading proton data. Therefore, in this work we consider these data sets as complementary ones, i.e., our basic assumption is that the diffractive and leading proton structure functions are parts of one and the same semi-inclusive proton structure function , which can be expressed in a way similar to Eq. (2). Throughout this work we will use the notation $`F_2^{SI}`$ for the semi-inclusive proton structure function, when referring to the diffractive and leading proton structure function data together.
The purpose of this paper is to reach a better understanding about the role of the pomeron and reggeon contributions in the interface between the diffractive and non-diffractive regimes through a global fit of the proton structure function obtained from H1 semi-inclusive DIS data (the ZEUS data were not employed in the fitting procedure, but their diffractive structure function measurements were used for checking our final results). In Sec.II, we define the kinematical variables and cross sections while our fitting procedure is presented in Sec.III. In Sec.IV, we present our fit results and a preliminary discussion, while a procedure to compare diffractive and leading proton data is described in Sec.V. Our main conclusions are summarized in Sec.VI.
## II Kinematics and Cross Sections
The usual variables employed to describe $`ep`$ DIS are depicted in Fig. 1a. One can define the squared energy in the $`ep`$ center of mass system (CMS) in terms of the four-momenta $`P`$ and $`k`$, referring respectively to the incoming proton and electron (or positron), as
$$s=(P+k)^2$$
(3)
and the squared energy in the $`\gamma ^{}p`$ CMS as
$$W^2=(P+q)^2.$$
(4)
The photon virtuality $`Q^2`$, the Bjorken $`x`$ and the variable $`y`$ are given by
$`q^2`$ $`=`$ $`Q^2=(kk^,)^2,`$
$`x`$ $`=`$ $`{\displaystyle \frac{Q^2}{2Pq}}={\displaystyle \frac{Q^2}{W^2+Q^2m_p^2}},`$
$`y`$ $`=`$ $`{\displaystyle \frac{Pq}{Pk}}.`$
If we ignore the proton mass, we have the following relations among these variables:
$$Q^2=xys$$
(5)
and
$$W^2=Q^2\frac{(1x)}{x}\frac{Q^2}{x},$$
(6)
being that $`x<<1`$ has been assumed in the latter expression.
For the case presented in Fig. 1b, where a baryon with four-momentum $`P^{}`$ is detected in the final state, we can also define the variables
$`t`$ $`=`$ $`(PP^{})^2,`$ (7)
$`\xi `$ $`=`$ $`{\displaystyle \frac{Q^2+M_X^2t}{Q^2+W^2}},`$ (8)
$`\beta `$ $`=`$ $`{\displaystyle \frac{Q^2}{Q^2+M_X^2t}}={\displaystyle \frac{x}{\xi }},`$ (9)
where the $`\beta `$ variable represents the fraction of momentum carried by a struck parton in the pomeron (if a pomeron exchange model is assumed).
Also, for leading baryons, it is usual to describe the data in terms of the fraction of momentum carried by the outgoing proton, $`z=P^{}/P`$, where $`z`$ is connected with $`\xi `$ by
$$z=1\xi .$$
(10)
The differential cross section for a semi-inclusive DIS process giving rise to leading baryon behavior is written as
$$\frac{d^3\sigma }{dxdQ^2dz}=\frac{4\pi \alpha _{em}^2}{xQ^4}\left[1y+\frac{y^2}{2(1+R)}\right]F_2^{LB(3)}(z,x,Q^2).$$
(11)
In the case of diffractive events, such a cross section is often expressed in terms of the $`\beta `$ and $`\xi `$ variables,
$$\frac{d^3\sigma }{d\beta dQ^2d\xi }=\frac{4\pi \alpha _{em}^2}{\beta Q^4}\left[1y+\frac{y^2}{2(1+R)}\right]F_2^{D(3)}(\xi ,\beta ,Q^2).$$
(12)
Here $`R=\sigma _L/\sigma _T`$ is the ratio between the cross sections for longitudinally and transversely polarized virtual photons. Under certain conditions, it is possible to assume $`R0`$ and thus the experimental behavior of the cross sections (11) and (12) is expressed in terms of the structure functions $`F_2^{LB(3)}(z,x,Q^2)`$ and $`F_2^{D(3)}(\xi ,\beta ,Q^2)`$. Specifically for the H1 diffractive data, such assumption was applied for those data with $`y<0.45`$ H1-Diff .
Thus, our analysis is directed to study the behavior of both $`F_2^{LB(3)}(z,x,Q^2)`$ and $`F_2^{D(3)}(\xi ,\beta ,Q^2)`$ data. We notice that these data are already integrated over the $`t`$-range corresponding to their respective experiments. In order to compare these data among themselves it is necessary to explicitly introduce the $`t`$-dependence on the structure functions. We discuss that issue in details in Sec. V.
## III Model, Parameters and Fitting Procedure
In the present study we have used the diffractive structure function data $`F_2^D`$ obtained by the H1 Collaboration H1-Diff , together with their measurements of the leading baryon structure functions $`F_2^{LP}`$ for protons and $`F_2^{LN}`$ for neutrons H1-Lead , in the same analysis. The $`F_2^D`$ data cover the kinematical ranges:
$`1.210^4<`$ $`x`$ $`<\mathrm{\hspace{0.33em}2.37}10^2,`$
$`4.5<`$ $`Q^2`$ $`<\mathrm{\hspace{0.33em}75}GeV^2,`$
$`0.04<`$ $`\beta `$ $`<\mathrm{\hspace{0.33em}0.9}.`$
while, for the leading baryon $`F_2^{LB}`$ measurements, the covered kinematical region are:
$`10^4<`$ $`x`$ $`<\mathrm{\hspace{0.33em}3.3}10^3,`$
$`2.5<`$ $`Q^2`$ $`<\mathrm{\hspace{0.33em}28.6}GeV^2,`$
$`3.710^4<`$ $`\beta `$ $`<2.710^2.`$
We notice that, although these data sets are overlapping in terms of $`x`$ and $`Q^2`$ ranges, they are complementary in terms of the $`\beta `$, the Bjorken variable for the presumable pomeron constituents.
As stated before, the H1 diffractive structure function, $`F_2^{D(3)}(\xi ,\beta ,Q^2)`$, can be written as a combination of two Regge exchanges with the quantum numbers of the vacuum, the pomeron and the reggeon ones H1-Diff . The most general expression for such a diffractive structure function reads
$$F_2^{D(3)}(\xi ,\beta ,Q^2)=g_{𝙸𝙿}(\xi )F_2^{𝙸𝙿}(\beta ,Q^2)+g_{𝙸𝚁}(\xi )F_2^{𝙸𝚁}(\beta ,Q^2)+g_𝙸(\xi )F_2^𝙸(\beta ,Q^2).$$
(13)
Here, functions $`g_{𝙸𝙿}(\xi )`$ and $`g_{𝙸𝚁}(\xi )`$ represent, respectively, the pomeron and reggeon flux factors integrated over $`t`$, while $`F_2^{𝙸𝙿}(\beta ,Q^2)`$ and $`F_2^{𝙸𝚁}(\beta ,Q^2)`$ are the pomeron and reggeon structure functions. The last term on the right-hand-side of Eq. (13), $`g_𝙸(\xi )F_2^𝙸(\beta ,Q^2)`$, accounts for a possible interference effect between the pomeron and reggeon exchanges.
The fluxes are taken from the Regge phenomenology of hadronic soft diffraction, and are written as
$$g_{𝙸𝙿}(\xi )=\xi ^{12\alpha _{𝙸𝙿}^0}_{|t_{min}|}^{|t_{max}|}e^{(\alpha _{𝙸𝙿}^,ln\xi )t}F_1^2(t)𝑑t$$
(14)
and
$$g_{𝙸𝚁}(\xi )=\xi ^{12\alpha _{𝙸𝚁}^0}_{|t_{min}|}^{|t_{max}|}e^{(b_{𝙸𝚁}^0\alpha _{𝙸𝚁}^,ln\xi )t}𝑑t,$$
(15)
where $`|t_{min}|`$ and $`|t_{max}|`$ are the minimum and maximum absolute $`t`$ values of the data for each experiment. In these expressions, the parameters $`\alpha _{𝙸𝙿}^0`$, $`\alpha _{𝙸𝚁}^0`$ and $`\alpha _{𝙸𝙿}^,`$, $`\alpha _{𝙸𝚁}^,`$ are, respectively, the intercept and slope of the pomeron and reggeon linear trajectories, that is
$$\alpha _{𝙸𝙿}(t)=\alpha _{𝙸𝙿}^0+\alpha _{𝙸𝙿}^,t\mathrm{and}\alpha _{𝙸𝚁}(t)=\alpha _{𝙸𝚁}^0+\alpha _{𝙸𝚁}^,t,$$
(16)
and $`F_1(t)`$ in Eq. (14) is the Dirac form factor given by
$$F_1(t)=\frac{4m_p^20.28t}{4m_p^2t}\left(\frac{1}{1t/0.71}\right)^2.$$
(17)
The interference term $`g_𝙸(\xi )F_2^𝙸(\beta ,Q^2)`$ is related to the pomeron and reggeon fluxes and structure functions by
$$F_2^𝙸(\beta ,Q^2)=\sqrt{F_2^{𝙸𝙿}(\beta ,Q^2)F_2^{𝙸𝚁}(\beta ,Q^2)}$$
(18)
and
$$g_𝙸(\xi )=2I_{|t_{min}|}^{|t_{max}|}\mathrm{cos}\{\frac{\pi }{2}[\alpha _{𝙸𝙿}(t)\alpha _{𝙸𝚁}(t)]\}\sqrt{e^{b_{𝙸𝚁}t}F_1^2(t)}\xi ^{1\alpha _{𝙸𝙿}(t)\alpha _{𝙸𝚁}(t)}𝑑t.$$
(19)
The expression above is quite similar to the one used by the H1 Collaboration to account for interference contribution in their diffractive structure function analysis H1-Diff . Following their procedure, we introduced a free parameter $`I`$ to account for the degree of interference between the pomeron and reggeon exchanges. Such a parameter is allowed to vary from 0 to 1.
Here we mostly intend to explore the connection between the diffractive and leading proton regimes, although the available data are quite separated in terms $`\beta `$. Therefore, we need a general functional form for the pomeron structure function that could be able to consider both the low $`\beta `$ (leading proton) and high $`\beta `$ (diffractive) regimes. In order to do that, we choose for the pomeron a functional form based on the same phenomenological parameterization as used in the H1 QCD analysis of the diffractive structure function H1-Diff , where a quark flavor singlet distribution $`\beta S_q(\beta ,Q^2)=u+\overline{u}+d+\overline{d}+s+\overline{s}`$ and a gluon distribution $`\beta G(\beta ,Q^2)`$ are parameterized in terms of the coefficients $`C_j^{(S)}`$ and $`C_j^{(G)}`$, according to:
$`\beta S(\beta ,Q^2=Q_0^2)`$ $`=`$ $`\left[{\displaystyle \underset{j=1}{\overset{n}{}}}C_j^{(S)}P_j(2\beta 1)\right]^2exp({\displaystyle \frac{a}{\beta 1}})`$
$`\beta G(\beta ,Q^2=Q_0^2)`$ $`=`$ $`\left[{\displaystyle \underset{j=1}{\overset{n}{}}}C_j^{(G)}P_j(2\beta 1)\right]^2exp({\displaystyle \frac{a}{\beta 1}}).`$ (20)
where $`P_j(\zeta )`$ is the $`j^{th}`$ member in a set of Chebyshev polynomials, with $`P_1=1`$, $`P_2=\zeta `$ and $`P_{j+1}(\zeta )=2\zeta P_j(\zeta )P_{j1}(\zeta )`$. We have summed these terms up to $`n=3`$ and set $`Q_0=2GeV^2`$, in order to contemplate the $`Q^2`$ range of both diffractive and leading proton data. Following H1, we also set $`a=0.01`$. Therefore, Eq. (20) has 6 parameters to be fixed by the fit.
Since it is not possible to totally separate the pomeron structure function from its flux factor, the parameters $`C_j^{(S)}`$ above also set the overall normalization of the pomeron contribution. The gluon and quark distributions above are evolved in leading order (LO) and next-to-leading order (NLO) by using the QCDNUM16 package QCDNUM , and the final pomeron structure function is written in terms of the singlet quark distribution as
$$F_2^{𝙸𝙿}(\beta ,Q^2)=<e^2>(u+\overline{u}+d+\overline{d}+s+\overline{s})$$
(21)
where $`<e^2>`$ is the average charge of the distribution, and for three flavors $`<e^2>=2/9`$.
For the reggeon, we assume the hypothesis of a direct relation between the reggeon structure function and the pion structure function by using
$$F_2^{𝙸𝚁}(\beta ,Q^2)=N_{𝙸𝚁}F_2^\pi (\beta ,Q^2),$$
(22)
where $`N_{𝙸𝚁}`$ is a free normalization parameter, and for the pion structure function we choose the LO GRV parameterization GRV . Such a choice is supported by the good description it provided for the H1 leading baryon data H1-Lead .
In fact, the identification of the reggeon structure function with the pion one is not new, and some authors already have applied it to the analysis of the H1 diffractive structure function data Royon .
Specifically for our case, we also choose to identify the reggeon exchange explicitly with the $`f_2`$ family of resonances, which has the right quantum numbers for the processes analyzed here and is characterized by its high intercept, $`\alpha _{𝙸𝚁}^00.68`$ Dino1 .
For the leading proton structure function, $`F_2^{LP(3)}(\xi ,\beta ,Q^2)`$, besides the pomeron and reggeon contributions, the pion exchange also plays a major role. In fact, the pion contribution is known to have an important role in hadronic leading proton Mirian and seems to work as an effective background for $`\overline{p}p`$ diffractive reactions at small $`t`$ Dino2 , besides its role in DIS H1-Lead . Indeed, pion exchange has a well known phenomenological behavior, so we took the pion flux factor straight out of the literature as being
$$f_\pi (\xi ,t)=\frac{g_{pp}}{4\pi }\frac{1}{4\pi }\frac{|t|}{(t0.02)^2}\xi ^{12\alpha _\pi (t)},$$
(23)
where $`g_{pp}/4\pi =13.6`$ is the coupling constant for $`pppX`$. Note that for the inclusive neutron production, $`ppnX`$, there is an extra factor $`2`$ in the coupling constant due to the Clebsh-Gordan coefficient for such a process.
For the pion structure function, $`F_2^\pi (\beta ,Q^2)`$, we took the LO GRV GRV parameterization. With the flux above and the GRV structure function, we were successful in describing the DIS leading neutron data without any free parameter.
The expression for the leading proton structure function then reads
$$F_2^{LP(3)}(\xi ,\beta ,Q^2)=g_{𝙸𝙿}(\xi )F_2^{𝙸𝙿}(\beta ,Q^2)+g_{𝙸𝚁}(\xi )F_2^{𝙸𝚁}(\beta ,Q^2)+g_\pi (\xi )F^\pi (x,Q^2).$$
(24)
As we said at the beginning, our main assumption is that the diffractive and leading proton structure function are components of one and the same semi-inclusive (SI) structure function, which combines the contribution from both Eq. (13) and Eq. (24) in a single expression that reads
$`F_2^{SI(3)}(\xi ,\beta ,Q^2)=g_{𝙸𝙿}(\xi )F_2^{𝙸𝙿}(\beta ,Q^2)`$ $`+`$ $`g_{𝙸𝚁}(\xi )F_2^{𝙸𝚁}(\beta ,Q^2)`$ (25)
$`+`$ $`g_\pi (\xi )F^\pi (\beta ,Q^2)+g_I(\xi )F_2^I(\beta ,Q^2).`$
It should be noted that in the equation above, the pion contribution is significant only for $`\xi 0.1`$, therefore for the diffractive regime, $`\xi 0.05`$, Eq. (25) reduces to Eq. (13), where no pion exchange is considered.
In overall, we dealing with a maximum of 8 free parameters to be fixed by the fitting procedure. These parameters come from the pomeron structure function, Eq. (20) (6 parameters), reggeon normalization, Eq. (22) (1 parameter), and the interference contribution, Eq. (18) (1 parameter). As mentioned before, the pion contribution (flux factor, Eq. (23), and structure function, given by the GRV parameterization GRV ) is totally fixed by the standard phenomenology having no free parameter left.
The other parameters, such as the pomeron and reggeon trajectories (intercept and slopes), the slope of the reggeon $`t`$-dependence and the $`a`$ parameter from Eq. (20 were kept fixed by their values from the literature, since they are quite well established. In Table 1 we present the values used for these parameters throughout this paper.
It should be mentioned that we excluded from the fit all data lying the the resonance region ($`M_X^22GeV^2`$) and/or with $`y0.45`$. That leave us with a total of 170 diffractive structure function data and 48 leading proton structure function data, with adds to a total of 218 data.
## IV Results and Discussion
In Table 2, we present the results of our first three fits. Fit 1 represents the results of our global LO analysis of diffractive and leading proton structure function data, using Eq. (25) with no interference term included ($`I=0`$). Since we are dealing with two different sets of data, we added the statistic and systematic errors in quadrature. A $`\chi ^2/d.o.f.`$ of $`1.277`$ was obtained.
Fit 2 corresponds to the results of a global NLO analysis of the diffractive and leading proton structure function data, using Eq. (25) as the fitting equation, again with no interference term included ($`I=0`$). Although some of the parameters have significantly changed in comparison to Fit 1, the final result provided a $`\chi ^2/d.o.f.=1.276`$ which is basically the same as the one from the global LO fit.
Fit 3 corresponds to a fit of Eq. (13) to the diffractive structure function data only. The final $`\chi ^2/d.o.f.`$ obtained, with only statistical errors included, was $`\chi ^2/d.o.f.=1.106`$. Although, in this case, the interference component was left free, it was ruled out by the fit. An observation to be made at this point is that one must be careful when comparing this $`\chi ^2/d.o.f.`$ result with the one from the H1 QCD analysis of the same set of data H1-Diff , since our sample includes two sets of data that where not taken into account in the H1 analysis (those for $`Q^2=45GeV^2`$ and $`Q^2=75GeV^2`$ at $`\beta =0.9`$). That gives us a total of $`170`$ data, whereas H1 has only $`161`$. Our choice for the reggeon intercept has also some effect in improving the final $`\chi ^2`$ result.
Table 3 presents the results of global fits when the interference parameter $`I`$ set free. It was bounded to vary in the interval $`0I1`$, but, as can be seen, in both fits it assumed the maximum upper value. Comparing these results respectively to Fits 1 and 2, the $`\chi ^2/d.o.f.`$ improved a little in both the LO fit ($`\chi ^2/d.o.f.=1.16`$) and the NLO fit ($`\chi ^2/d.o.f.=1.18`$).
Table 2 and 3 also present the individual contributions to the $`\chi ^2`$ coming from the diffractive and leading proton data. For three of our global fits, we have a diffractive contribution around $`74\%`$, with the leading proton one around $`26\%`$. The only departure from these values comes from the global NLO fit with the interference parameter $`I`$ set free (Fit 5). For that we have the diffractive data contributing with $`68\%`$ and the leading proton data with $`32\%`$. It is worth to remember that, for the global fits, our data sample is composed of 218 data, 170 coming from diffractive and 48 from leading proton structure function. Therefore, the diffractive data corresponds to $`78\%`$ of our global data set, and the leading proton data to $`22\%`$.
In order to test the parameterization of the pomeron structure function, we compare some of our results for $`F_2^{𝙸𝙿}(\xi ,\beta ,Q^2)`$, Eq. (20), with the independent measurement of $`F_2^{D(3)}(\xi ,\beta ,Q^2)`$ by the ZEUS Collaboration ZEUS-old , where no sign of secondary exchanges was found. As shown in Fig. 2, all of the three fits exhibited are in good agreement with the data (which were not used in the fitting procedure), indicating that the pomeron contribution has been fairly accounted.
In Fig. 3, we plotted the diffractive structure function data from H1 Collaboration in comparison with the results of the same three fits shown in Fig 1. As can be seen, the agreement among the three fits is quite good at small $`\xi `$, but as $`\xi `$ increases Fit 3 grows faster than the other two. The difference between Fit 1 and Fit 2 (not shown in the figure) is quite small over the entire diffractive range of $`\xi `$, which is expected since both fits give close values for the $`\chi ^2`$.
Fig. 4 shows the $`F_2^{LP}`$ data from H1 Collaboration together with the results from Fit 2 to illustrate the description of the leading particle behavior. The leading neutron data, from the same experiment, are also included (these data can be described assuming pion exchange as the only contribution for the reaction and so were not employed in the fitting procedure).
After showing all of these results, some comments are in order. Firstly, from Fits 1 and 2, we see that applying LO or NLO evolution equations produce basically the same result in terms of $`\chi ^2`$, although, as expected, some parameters suffer a little change (the same can be said about Fits 4 and 5). We remind that these parameters reflect the quark and gluon content of the pomeron as obtained from different scenarios.
The comparison between Fit 2 (global) and Fit 3 (only diffractive data) present much more remarkable effects. Not only the parameters change, but in the latter case there is a strong enhancement of the secondary contribution. However, this is a suspicious effect since the diffractive data are quite limited in terms of the $`\xi `$ variable and secondary reggeon contribution are supposed to play an important role only for $`\xi 0.1`$ (see more comments about this aspect in the next section).
When we perform the global fit, but leaving the interference term completely free to be established by the $`\chi ^2`$ minimization, it assumes its maximum value (Fits 4 and 5). Again it is the case of asking whether this outcome reflects a reliable physical effect or is just a fitting artifact. Answering this question is beyond the scope of this paper, but we have strong evidences indicating that the introduction of the interference term makes the corresponding structure functions inadequate to describe the results of diffractive photo- and eletroproduction of dijets by both H1 and ZEUS collaborations. On the other hand, diffractive structure function obtained without interference effects allow a very good description of both dijet production processes Altem .
## V Bringing diffractive and leading proton structure functions together
Now, some words are needed to explain how we handled together both the sets of data displayed in Fig. 5, since it is the central piece of our study. In that figure, we bring together the diffractive and leading proton data and compare the results of our three NLO fits to this combined set of semi-inclusive data.
Here, we are mostly interested in analyzing the behavior of these data in terms of $`\xi `$. Since the $`\beta `$ range for the diffractive and leading proton data are very distinct, the usual procedure of plotting together data with the same values of $`\beta `$ and $`Q^2`$ would not be the best choice. There is, however, a large overlap of these two sets in terms of the variables $`x`$ and $`Q^2`$. Thus, we choose to combine the data in groups with the same (or as close as possible) values of $`x`$ and $`Q^2`$. That is a more proper way to show that the difference between the diffractive and the leading proton regime is due to the $`\xi `$ region where the semi-inclusive process $`epepX`$ is measured, according to our assumption that both sets of data can be embraced by the same semi-inclusive structure function.
Still a problem remains. Besides the different $`\beta `$ range, both the diffractive and leading proton structure functions were measured at different $`t`$ intervals. The diffractive data were measured for the interval $`|t_{min}|<|t|<1GeV^2`$, whereas the leading proton ones where measured for the interval $`|t_{min}|<|t|<|t_0|`$, where
$`t_{min}`$ $`=`$ $`{\displaystyle \frac{m_p^2\xi ^2}{(1\xi )}},`$ (26)
$`t_0`$ $`=`$ $`{\displaystyle \frac{p_{T,max}^2}{(1\xi )}}+t_{min},`$ (27)
with $`p_{T,max}=0.2\mathrm{GeV}`$. Since this last interval corresponds to a range smaller than the diffractive one and since the phenomenological $`t`$ dependence coming from the diffractive region seems to be well established for both hadronic and DIS events, in Fig. 5 we scaled down the diffractive structure function data in order to make them comparable to the leading proton data.
It should be noticed that such a correction is intended only as a visualization device. In our whole fitting analysis, we took the data at their correct measured $`t`$ intervals.
In order to make such a correction as independent of our own analysis as possible, we choose to proceed by the following way. A fit of Eq. (13) to the diffractive structure function data was performed, with the interference parameter $`I`$ set to zero (no interference). The fluxes were those given by Eq. (14) and Eq. (15), with the pomeron and reggeon intercept kept fixed with those values obtained from the H1 analysis H1-Diff ($`1.20\pm 0.01`$ and $`0.57\pm 0.01`$ respectively). For any fixed values of $`\beta `$ and $`Q^2`$, the pomeron and reggeon structure functions were treated as free parameters to be fixed. Once those parameters were determined for each set, it was possible to calculate the ratio
$$R(\xi ,\beta ,Q^2)=\frac{_{|t_{min}|}^{|t_{0,Lead}|}F_2^{D(3)}(t,\xi ,\beta ,Q^2)𝑑t}{_{|t_{min}|}^{|t_{0,Diff}|}F_2^{D(3)}(t^{},\xi ,\beta ,Q^2)𝑑t^{}},$$
(28)
which should be used to correct each measured diffractive structure function data point at a given $`\xi `$, $`\beta `$ and $`Q^2`$.
Such procedure provided a correction factor that is a function of $`\xi `$, going from 0.25 to 0.4. This is reflected in the curves shown in Fig. 5. From that figure it is clear that the $`f_2`$ reggeon contribution coming from Fit 3 overestimates the leading proton data by a factor 2 at least. The only parameter related to this exchange is the normalization $`N_{𝙸𝚁}`$, and from Table 2 it is clear that the fit to the diffractive data alone drives such parameter to a very high value, compared with the one from the global Fits 2 and 5, that are both quite compatible with the combined sets of data.
## VI Conclusions
The analysis in this paper shows that we have to be very careful before drawing conclusions about the role of Regge exchanges in diffractive DIS. If only the H1 high statistic diffractive data were used, as we have done in our Fit 3, an extrapolation of such a result to the leading proton region will overestimate those data by, at least, a factor 2 (Fig. 5). It could be argued that such an extrapolation goes to low $`\beta `$ values beyond the range of the fitted data, and our pomeron structure function would not be valid anymore. That is true, but the point is that the pomeron contribution alone is not important in such extrapolation. It is the secondary reggeon plus the pion contribution that play the major role in the leading proton region. The pion contribution itself is fixed and provides a quite reasonable description of the leading neutron data. The same pion structure function is used by the reggeon exchange, and it has been shown that a such combination provides a good description of the leading baryon data H1-Lead . Therefore our choice of structure functions for the secondary exchanges works well in both regimes, and it is fair to expect that, extrapolating the information about the ratio between pomeron and reggeon from the diffractive SF to the leading proton regime, we should be able to have a decent qualitative description of the leading proton data, but instead we were left with a result that not only does not describe the data, but also lives no room for corrections with extra reggeon exchanges.
The main problem in connecting the diffractive and leading baryon regimes seems to come from the relative weight that the fit put over the reggeon contribution in each case. For instance, the normalization parameter $`N_{𝙸𝚁}`$ changes from $`7.25`$, when only diffractive data are used, to $`2.058`$, when both diffractive and leading proton data are put together.
Although the interference term has some impact over the reggeon contribution, it plays a minor role that does not improve at all the discrepancies discussed above.
The fact that ZEUS Collab. has found no secondary exchange in their diffractive measurements ZEUS-old ; Zeus is also an evidence that the diffractive structure function data alone cannot conclusively provide information concerning the contribution of the secondary reggeon exchange in semi-inclusive $`ep`$ reactions. Therefore, the leading baryon data represent an important constraint that must be taken into account in any analysis based on the Regge picture of diffraction.
The next step following this analysis is to show how these different parameterizations affect the theoretical predictions for the cross sections of diffractive photo- and electro-production of dijets, also measured by ZEUS and H1 Collaborations Zeus-phot ; H1-phot . This is going to be reported in a forthcoming paper Altem .
## Acknowledgments
We would like to thank the Brazilian governmental agencies CNPq and FAPESP for financial support.
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# HIV time hierarchy: Winning the war while, loosing all the battles.
## Introduction:
The natural progression of the Human Immuno-deficiency Virus (HIV) infection varies between individuals, however a general pattern of the progression has been observed (Figure 1).
* Within weeks of infection, a short transient jump of plasma virema (virion concentration in plasma) is seen together with a marked decrease in Immune cell counts (CD4 T Helper cells).
* Partial control of the disease by the immune system ensues, causing variable periods of practically a-symptomatic clinical latency, which can last years. During this period the Immune cell population continues to slowly decline until the immune system is so crippled it can no longer contain the disease.
* A renewed outbreak of the virus which, together with constitutional symptoms and the onslaught by opportunistic diseases, cause death .
The dynamics of HIV was traditionally described using simple homogeneous ODEs (for a review see Perelson ). This method was enlarged to models considering the spatial structure mainly by Zorzenon dos Santos and Coutinho . Such models consider the importance of the localization of the interactions between the HIV virions and immune cells. Taking into consideration the global features of immune response and the high mutation rate of HIV they described the spatial and temporal interactions of infected and healthy cells in lymphoid tissue in the body.
The methodology used in spatially extended models was limited up to now mainly to cellular automata , or to compartmental models . The main advantage of cellular automata is their capacity to emphasize the emergence and the importance of spatial structures. For example the propagation of waves of infection in the lymph nodes, or the creation of immune cell aggregations. Consequently Zorzenon dos Santos and Coutinho, in contradiction with most ODE models, succeeded to reproduce all the stages of HIV evolution, using a single set of rules \[3, figure 13 in 2\].
The present model is inspired by the observed interesting features in the spatially extended model . Our model acts in the shape space rather than in the physical space and exploits the dynamical implications of the HIV virus ’propagation’ in the molecular shape space: The immune system’s need to identify the virus’s form correctly in the shape space of possible viral and immune receptor shapes. The main ingredients in our model are the mutations (represented by propagation of the multiplying virions in the structure-less infinite dimensional shape space) and the confrontation with the immune system cells (resulting in the disappearance of the both cells and virions). Our model does not display any geometrical patterns in either space or shape space. In this model we assume that we can average the spatially distributed reactions (as in the ODE approach).
The model presented in the following pages relies on the basic known facts on the immune defense system. The viruses (antigens) have various characteristic geometrical shapes. In order to act against them, the immune system has to ’identify’ them by producing cells which contain shapes complementary to the geometry of (parts of) these antigens. Since the immune system does not know a priori what is the characteristic shape of every new invading virus, the immune system generates randomly cells with various shapes. If a cell encounters (by chance) an antigen (virion) with complementary shape, then more cells with characteristic shapes identical to it are produced and a mechanism is triggered for the destruction of all the individuals (virions) belonging to this virus strain (and sharing the same shape).
Usually the destruction mechanism is quite efficient and once a virion is ’identified’ by the above random search in the shape space, its fate and the fate of all the virions in the same strain is sealed: they are wiped out by the immune system within days.
With HIV however, the issue is more complicated: Since the virus’s replication mechanism is relatively imprecise, as it multiplies it undergoes a large amount of mutations/changes in shape compared to those found in other kinds of virus . Based on empirical and theoretical results in the research of HIV we propose the following scenario. The immune system cells that are complementary to the old shape are ineffective in dealing wiht the new mutant virus strain. The virions belonging to the strain with the new shape can multiply with impunity until a strain of immune cells which fits the new shape is generated by the immune system. Once the Immune system’s shape generation process succeeds to produce by chance a immune cell carrying a complementary shape to the new virus strain and this cell encounters (by chance) a virion belonging to the new strain, the new strain is wiped out too (with the exception of the eventual new mutants that again can multiply freely until their new shape is detected by the immune system). The process continues indefinitely with the virus loosing every battle but succeeding to produce increasingly many small populations of new shapes (figure 2) .
The process is compounded by the additional fact that virions can kill directly and indirectly immune cells (whether or not they are of complementary shape).
The immune system continues to win every battle until the increase in production of immune cells with shape complementary to a virus strain is overcome by the rate at which the immune cells are destroyed by various other HIV strains. At this point the immune system has effectively lost the war.
In the rest of the paper we provide a detailed microscopic simulation model which supports this scenario and that fits quite well the known phenomenological data on HIV in terms of the following basic mechanisms:
* The local (in shape space) destruction of HIV by the immune system.
* The fast mobility of HIV in shape space (high mutation rate).
* The global destruction of immune system cells by HIV.
## 1 The Model
We represent the shape space of the virus by a random lattice in which each site ($`i`$) has a fixed number of neighbors. Neighboring sites represent shapes that can be reached one from the other by a single base mutation of the virus. The occupation number on each site ($`N_{Vi}`$)represents the number of virions with that shape existing in the organism. The immune system cells that recognize that shape are also represented through an occupation number on the same site($`N_{Ci}`$). Note that the existence of a virus and an immune cell on same lattice site does not imply their proximity in real space: Quite to the contrary they might be located in very distant locations in the organism. There is however a small probability that the virus and the corresponding cell will meet and react in real space. Therefore each pair of virion and immune cell located on the same lattice site has a small, but finite probability to react (according to the rules described in detail below).
One represents the eventual mutations of the virus and the immune cells by their rate of diffusion in the shape space ($`D_V`$,$`D_C`$). More precisely both viruses and immune cells have a certain probability for jumping between neighboring sites.<sup>1</sup><sup>1</sup>1The diffusion rate, and the neighborhood structure implied by the node’s connectivity can be different for the virus and for the immune cells
HIV can replicate in and destroy immune cells. This is irrespective of the cell’s characteristic shape. I.e. the virus can destroy immune cells located on sites arbitrarily far away from the site of the virus). In our model we represent this by:
* A virus proliferation rate proportional to the total immune cells population ($`C_{tot}`$).
* An immune cell death rate proportional to the total viral population ($`V_{tot}`$).
We list bellow the reactions taking place in the model:
1. When an HIV virion and an immune cell reside on the same site the immune cell duplicates with a rate of $`\tau _C`$. However following realistic biological data, we limit the multiplication rate of the immune cells (to a factor of 3 per day).
2. When an HIV virion and an immune cell reside on the same site the virion is destroyed with a probability rate of $`d_V`$.
3. Each HIV virion replicates with a rate ($`\tau _VC_{tot}`$) proportional to the total number of immune cells.
4. Each immune cell is destroyed with a probability rate ($`d_CV_{tot}`$) proportional to the total number of virions.
5. New immune cells, with various shapes are created continuously. We represent this by a probability rate ($`\lambda `$) for an immune cell to appear on a random lattice site.
6. Both the immune cells and the HIV virus diffuse slowly in the shape space with rates $`D_V`$ and $`D_C`$ .
## 2 Results
The reactions listed in the previous section lead to the following scenario for the evolution of the HIV infection. The virus enters the body in a high concentration limited to a restricted number of strains. As long as there is no immune cell in the site corresponding to a certain strain, the virions located in this strain site proliferate exponentially. The rise in virus concentration is accompanied by a destruction of random T cells hosting the virus (according to item 4). These T cells can be located anywhere in the lattice. Eventually one of the immune cells generated by the immune system will fall by chance on the site of this strain (according to item 5) . This immune cell will proliferate very fast, since, according to item 1, the proliferation rate increases with the local viral concentration. The resulting high local immune cells concentration will destroy all the virions of this strain (item 2). As a result, after a short period (1-2 month) all the initial strains will be discovered by the immune system and will be destroyed ending the acute phase of the disease. In the absence of virus diffusion in shape space (i.e. mutations), this would stop the disease (Figure 3).
It is the diffusion rate of the virus in shape space (item 6) which is responsible for the continuation of the infection. Before all the initial strains are destroyed, some of the virions have a chance to mutate and escape to lattice points containing no immune cells.
Each of these new lattice points in the shape space will contain a lower concentration of virus than the original infection. Since the probability for the discovery of a virus strain by the immune cells is proportional to the virus strain concentration (according to item 1) the time it will take for the immune system to discover, and destroy these new strains (item 2) will be longer than in the acute phase. By the time these new strains are destroyed some virions from these strains will have diffused to neighboring sites (undergo mutations), and constitute the germs for a new generation of emerging strains. These strains in turn will have the fate of their predecessors in the previous generation. One sees now that one can describe the long-term evolution of the HIV infection as an iterative process. More precisely the long-term evolution will consist of a chain of small infections, each of which is easily defeated by the organism. However after each such infection the number of strains will grow. Therefore even though the number of virions in each strain is always kept under control by the organism, the total amount of virions will stochastically slowly increase. The increase is stochastic, since it depends on the random time it takes the immune system to discover and destroy each new strain. As the number of virions increases, so does the death rate of the immune cells (item 4) . At a certain stage the death rate will be high enough to impair the capacity of the immune system to react locally to new strains. This constitutes the last stage of the disease (Figure 4).
This typical scenario can vary from person to person:
1. The typical case observed in most hosts is a disease composed of 3 stages (Figure 4). The first acute stage is due to the fast proliferation of the original strains before the appropriate immune response is set. This stage ends when the appropriate immune response to each and every original strains is completed. This stage takes approximately $`\lambda `$ln(Number of original strains) days (see below). This gives an expected number of virions at the peak, which is much higher than in a usual disease (Appendix A) . The latent stage is the stage in which the number of virion strains is limited, and the number of virions in each strain is low. This stage will last as long as the number of strains is much smaller than $`\frac{\tau _C}{d_C}`$ (Appendix B). When the number of strains grows above $`\frac{\tau _C}{d_C}`$ the last stage of the disease occurs. This is the regime in which the local (in shape space) activation of the immune system by the local virus strain (item 2) becomes lower than the global destruction rate of immune cells by the virions (item 4), and the immune system fails to destroy locally (item 4) the existing strains. At this stage many old strains that were kept under control during the previous stages can reappear.
2. If the efficency of immune reaction to HI (items 1 and 2) is high enough, the immune system will manage to defeat the new virus strains fast enough. Thus the average number of strains will stay constant or slowly decrease. In this case the total number of virions will vary around a fixed number, or slowly decrease to 0. This fate can be the one of the long time carriers8. In reality, the total number of virions never decreases to zero, since there are other sources (macrophages, neuronal cells, etc.) that contribute a small number of new virions (Figure 5).
3. When, on the other hand, the efficency of immune reaction is low a large number of new strains will be created before the immune cells will finish destroying the virions from the strains of the first infection. In this case the acute phase will directly lead to the death of the host (Figure 6).
The results of our simulations show that although each new strain that emerges is destroyed within a few weeks (like any ordinary disease), the long term evolution of the infection takes years. The time scale that determines the long-term evolution scale is the diffusion (mutation) rate of the virus. The number of new strains growes exponentially with time , but the unit of time in this exponential is the time it takes for a new strain to establish itself i.e. weeks.
One might ask how the scale hierarchy between the cellular interactions (hours) and the evolution of the infection (years) emerges in this model. The answer is the following. The single strain lifetime is determined by an exponential rate proportional to the local interaction rate (hours). This exponent rises to macroscopic virus concentration within weeks. Only when the value of this exponent is high enough does the long-term mechanism of virus mutation become operative. Thus the time unit for one interaction in the shape space is the time needed for a single strain to establish itself. The time scale of the entire disease is the time necessary for evolving a macroscopic number of strains. Therefore we expect that the time scale of the entire disease will relate to the time scale of new strain creation as the time scale of the strains relate to the individual virion division time. The actual numbers are indeed (13 years/ 2 weeks) = (2 weeks/hours).
## 3 Conclusions
The main success of the present model is the natural emergence of a hierarchy of very different dynamical time scales. The very long-term decrease in immune (CD4+ T) cells count cannot be explained by a simple dynamic system. The transition from the microscopic time scales (hours) to the macroscopic time scales (years) requires a profound explanation. A simple dynamical system would require extreme fine tuning of its parameters in order to achieve such a transition. We propose a mechanism that can bridge between the microscopic and macroscopic time scales, that does not need fine tuning. The transition can be expressed best by representing the evolution of the system in the shape space (the strain of the virus), rather than the real space (the location of the virus in the organism). The relavent unit of time step for the operations in shape space is the time it takes for a strain to reach a macroscopic concentration and therefore have a significant probability to generate a mutant. This time step is of the order of weeks and not of hours. The evolution of the disease takes a few hundred time steps. i.e a few hundreds of weeks.
The mechanisms operating at the short times and at the long times are completely different. At the microscopic level the mechanism is the recognition and the destruction of the virions by the immune cells. The short time scale evolution of the disease is similar to any other disease. The long scale evolution of HIV infection is based on the competition between localized (in shape space) processes and global processes. To be precise the evolution is due to the spread of the virus strains across the shape space. In the initial acute phase most of the viral load is distributed between a small number of localized virus strains. At the last stages of the disease the viral load is distributed between a large variety of many strains. The immune cells manage to destroy locally every particular virus strain within a couple of weeks from its emergence. However as the number of strains at each given moment increases, the virus succeeds to destroy an increasing amount of immune cells. Thus locally the virus looses every battle. Yet in the end the shape space is filled with a multitude of small but numerous strain populations. At that point the virus wins the war by killing more immune cells than it activates.
In short: During the latent stage every virus strain looses every battle since it is activating the cells that can destroy it. Transition to AIDS occurs when the combined contemporary virus population wins the war by killing in total more immune cells than the sum of cells it activates.
## Acknowledgment
We acknowledge Professor Zorzenon dos Santos for sharing with us the content of reference 3 before publication.
## Appendix A - Acute phase
Imagine one has a lattice containing V nodes out of which N are occupied by the strains that constitutes the initial infection. We assume that the population of each initial infection strain is large enough that if the immune system generates randomly an immune cell in this site (A cell with a complementary shape to that virus strain) then that cell will discover this strain with probability 1. The probability rate for generating in a given lattice site an immune cell is $`\lambda `$. Therefore the probability rate for discovering one of the strains is $`N\lambda `$. Thus the average time it takes the immune system to detect the first strain in the initial infection is $`\frac{1}{N\lambda }`$. Once this strain was discovered the probability rate to discover one of the remaining (N-1) strains is $`(N1)\lambda `$. This means that the time it takes to discover a second strain is $`\frac{1}{(N1)\lambda }`$ and so on. Thus the average time it will take for the immune system to discover the N strains is: $`\frac{1}{lambda}\mathrm{\Sigma }_1^N\frac{1}{i}`$ which is approximately $`ln(N)\lambda `$. As long as the virus strain is undetected by the immune system it can proliferate freely. Therefore the strain discovered last had the most time to proliferate. The factor by which it grew more than a single strain infection is :
$$e^{\frac{ln(N)1}{\lambda }\tau _VC_{total}}N^{\frac{\tau _VC_{total}}{\lambda }}$$
(1)
## Appendix B - Transition to AIDS
In this appendix we estimate the conditions for the final collapse of the immune system, when it fails to react appropriately to local virus strains. The dynamics of the population of the immune cells in a given site in the shape space is dominated by 2 parallel processes:
* Proliferation due to the activation of the immune cells by the interaction with the local (in shape space) virus strain ($`\tau _CV_i`$)(item 1).
* Random global destruction by arbitrarily shaped virus strains ($`d_CV_{tot}`$) (item 4).
An immune strain can increase its population if its proliferation rate is higher than its destruction rate. $`\tau _CV_i>d_CV_{tot}`$ . In other words we need the ratio between local virus concentration and global virus concentration to be :$`\frac{V_{tot}}{V_i}<\frac{\tau _C}{d_C}`$. If we have N virus strains proliferating in the system, we can assume that their population are of the same order of magnitude, and estimate $`\frac{V_{tot}}{V_i}=N_{strains}`$. Thus in order for an immune strain to be able to rise its concentration and react appropriately against the corresponding virus strain One needs the number of virus strains not to exceed : $`N_{strains}=\frac{\tau _C}{d_C}`$. If one assumes that the virus population is not equally divided between strains the inequality is even more stringent.
## Appendix C - Discretization vs Continuous Differential Equations
The evolution of a dynamical system can be expressed by 2 types of models:
1. Ordinary Differential Equations (ODE) that simulate the evolution of the average population under the assumption of spatial homogeneity.
2. Microscopic simulation (MS) models, that compute each reaction separately.
The ODEs have the advantage of being cheap in CPU time. They enable us to simulate precisely the system when its concentration is very high, but fail to describe the stochastic aspects of the system. The MS takes into account the stochastic effects and describes precisely the discrete aspects of the agents and of the strains. However MS is very inefficient if the number of agents and the probability for reaction are high.
In the present model most of the sites are basically empty. However there is a relatively small number of sites occupied by a macroscopic number of virions and immune cells. This special situation invalidates both the possibility to use continuity assumption (ODE), and discrete operations (MS).
We solved this problem by using a hybrid model. This model computes the probability for interaction between every 2 agents at every site on the lattice in a given time interval. If this probability is higher than a threshold (30) then the number of agents created or destroyed is computed in a deterministic way using an ODE formalism for this site. If on the other hand the probability for a reaction is lower than the threshold the number of new agents created or destroyed is computed in a discrete stochastic way.
This is a particular application of the hybrid models we developed .
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# Twistor Newsletter 43, pp. 14-21 (1997) Singularities of wavefronts and lightcones in the context of GR via null foliations
##
The null-surface approach to general relativity essentially reformulates general relativity in terms of two real functions on the bundle of null directions over the spacetime manifold (locally $`M^4\times S^2`$ with a metric $`g_{ab}(x^a)`$ on $`M^4`$). These two functions are $`Z(x^a,\zeta ,\overline{\zeta })`$, representing a sphere’s worth of null foliations of the spacetime (i.e., such that $`\stackrel{~}{g}^{ab}Z,_aZ,_b=0`$ for all values of $`\zeta `$ and for members $`\stackrel{~}{g}_{ab}`$ of the conformal class of $`g_{ab}`$), and $`\mathrm{\Omega }(x^a,\zeta ,\overline{\zeta })`$, representing a sphere’s worth of conformal factors with the role of picking the members of the the conformal class that satisfy the Einstein equations. More detail on this formulation can be found in .
Within the context of the null-surface approach to general relativity a family of null coordinate systems $`(\theta ^i,i=0,1,+,)`$ is heavily used to derive dynamical equations for $`Z`$ and $`\mathrm{\Omega }`$. The coordinates $`\theta ^i`$ are defined by derivation from $`Z`$:
$$\theta ^i=(Z,ð\overline{ð}Z,ðZ,\overline{ð}Z)(u,R,\omega ,\overline{\omega }),$$
(1)
which, for fixed $`\zeta `$, represents a transformation between an arbitrary coordinate system $`x^a`$ and $`\theta ^i`$. The coordinates $`\theta ^i`$ are adapted to the null foliations, so that the leaves $`Z(x^a,\zeta ,\overline{\zeta })=const.`$ of the foliation constitute surfaces of constant coordinate $`u`$.
There is some gauge freedom in the theory, due to the fact that there are different foliations which are all null with respect to the same metric, such as, for Minkowski spacetime, a foliation based on outgoing lightcones off a world line as opposed to a foliation by null planes. In the case of asymptotically flat spacetimes, we customarily fix the gauge by requiring that our null foliation consists of surfaces that asymptotically become null planes.
For every fixed value of $`\zeta `$, our special coordinates have a well defined interpretation. The coordinate $`u`$ labels leaves of the null foliation. The coordinates $`(\omega ,\overline{\omega })`$ label null geodesics on a fixed leaf. The remaining coordinate $`R`$ acts as a parameter along every null geodesic in the leaf. Because of our gauge choice of null surfaces becoming planes at null infinity, the null geodesics labeled by $`(\omega ,\overline{\omega })`$ constitute bundles of asymptotically parallel null geodesics.
It is natural to raise the objection that, generically, any vacuum spacetime other than Minkowski focuses non-diverging bundles of null geodesics , and therefore our coordinate systems break down at the point of focusing, by assigning different labels to the same spacetime point. In this respect, although coordinate singularities are irrelevant to the physical content of the dynamical null-surface equations, we feel that the break down of these particular coordinates has a certain appeal since it entails the existence and location of caustics. Caustics and their singularities, as well as the singularities of wavefronts, have been classified by Arnol’d within a sophisticated mathematical context . On the other hand, caustics are increasingly being considered in the field of astrophysical observations .
The null-surface approach provides a dual interpretation to the function $`Z`$. The condition $`Z(x^a,\zeta ,\overline{\zeta })=u`$ for fixed $`x^a`$ picks up the points $`(u,\zeta ,\overline{\zeta })`$ at scri which are connected to $`x^a`$ by null geodesics. These points lie on a two-surface at scri, referred to as the lightcone cut of the point $`x^a`$. Generically, due to focusing in the interior, the lightcone cuts have self-intersections and typical wavefront singularities such as cusps and swallowtails. This appears to pose a technical difficulty regarding the perturbative approach to solving the null-surface equations, since the occurrence of this type of singularity generically entails divergences in the derivatives of the lightcone cuts.
In the following, we examine the occurrence of singularities of wavefronts and lightcones and its relevance to the null-surface approach.
### A Singularities of the null coordinates
Consider a foliation of spacetime by past null cones from points $`(u,\zeta ,\overline{\zeta })`$ at scri along a fixed null generator $`\zeta `$. (This is the compactified version of a foliation by null surfaces that are asymptotically null planes). Every past lightcone in this foliation has singularities, in the sense that the lightcone “folds” and self-intersects, due to focusing in the interior spacetime, as shown in Fig. 1. The points where the past lightcone is singular are points where neighboring null geodesics of the congruence intersect, and are thus conjugate to the point at scri. Translating this into the physical non-compactified spacetime, these points in the interior are “focal points”, such that light rays emitted from them are asymptotically parallel. How can we locate these “focal points” in terms of null-surface variables?
A preliminary answer to this question can be approached quite directly from a consideration of the geodesic deviation vector of the congruence that becomes asymptotically parallel in a direction $`(\zeta ,\overline{\zeta })`$ at scri. Every null geodesic in this congruence is characterized by fixed values of $`(u,\zeta ,\overline{\zeta },\omega ,\overline{\omega })`$. The geodesic deviation vector has been derived earlier in . Along a fixed null geodesic, the cross sectional area of the congruence has the expression:
$$A_p(R;u,\zeta ,\overline{\zeta },\omega ,\overline{\omega })=\frac{1}{\mathrm{\Omega }^2\sqrt{1\mathrm{\Lambda },_1\overline{\mathrm{\Lambda }},_1}},$$
(2)
where $`,_1`$ represents $`/R`$. This is an expression of the area $`A_p`$ in terms of the two null-surface variables
$$\mathrm{\Omega }=\mathrm{\Omega }(x^a,\zeta ,\overline{\zeta })\mathrm{\Lambda }ð^2Z(x^a,\zeta ,\overline{\zeta }).$$
(3)
In (2), the quantities $`\mathrm{\Omega }(x^a,\zeta ,\overline{\zeta })`$ and $`\mathrm{\Lambda }(x^a,\zeta ,\overline{\zeta })`$ are evaluated at fixed $`(\zeta ,\overline{\zeta })`$ and at values of $`x^a`$ along the null geodesic given by $`(u,\zeta ,\overline{\zeta },\omega ,\overline{\omega })`$.
It is relevant to point out that the variable $`\mathrm{\Omega }`$ actually is the product of two factors, one of which is an arbitrary confromal factor for the conformal class (and as such, it does not depend on $`\zeta `$), whereas the other factor carries conformal information via its $`\zeta `$-dependence. Therefore it is not surprising to find that it plays a role in a completely conformally invariant matter such as the determination of conjugate points of null geodesic congruences.
Focusing takes place at the value of $`R`$ such that $`A_p=0`$. However, the only way for the area to vanish is that the denominator become infinite. The square root in the denominator can not diverge before becoming pure imaginary. This leaves us with the previously unsuspected result that along a fixed null geodesic in the past lightcone from a point $`(u,\zeta ,\overline{\zeta })`$ at scri, $`\mathrm{\Omega }`$ must blow up for focusing to take place.
The vanishing of $`A_p`$ is also related to the vanishing of the determinant of the metric, since in these coordinates we have
$$g\text{det}(g_{ij})=(\text{det}(g^{ij}))^1=\frac{1}{\mathrm{\Omega }^8(1\mathrm{\Lambda },_1\overline{\mathrm{\Lambda }},_1)}=\mathrm{\Omega }^4A_p^2$$
(4)
Thus both the vanishing of $`A_p`$ and the fact that $`\mathrm{\Omega }`$ diverges result in the vanishing of the 4-dimensional volume element.
This divergence of $`\mathrm{\Omega }`$ has another significant consequence. Along the null geodesics labeled by $`(\omega ,\overline{\omega })`$ there is a choice of an affine parameter $`s`$, which is related to $`R`$ via
$$\frac{ds}{dR}=\mathrm{\Omega }^2\text{or alternatively}\frac{dR}{ds}=\mathrm{\Omega }^2.$$
(5)
Since the affine parameter is regular, it follows that $`dR/ds`$ blows up as well at the point where $`\mathrm{\Omega }`$ does. Thus $`R`$ is a bad coordinate (as we might have suspected) in the neighborhood of a focal point. Both $`R`$ and $`\mathrm{\Omega }`$ are, however, determined by the function $`Z`$ through the same second derivative (See for details):
$$Rð\overline{ð}Z(x^a,\zeta ,\overline{\zeta }),\mathrm{\Omega }^2g^{ab}Z,_að\overline{ð}Z,_b,$$
(6)
thus it is consistent to attribute both complications to $`ð\overline{ð}Z(x^a,\zeta ,\overline{\zeta })`$ becoming singular, since $`g^{ab}`$ is smooth in a good choice of coordinates $`x^a`$.
### B Singularities of the lightcone cuts
In asymptotically flat spacetimes the function $`Z`$ can also be viewed as describing the intersection of the lightcone of a point $`x^a`$ with scri, via $`u=Z(x^a,\zeta ,\overline{\zeta })`$ where $`(u,\zeta ,\overline{\zeta })`$ are Bondi coordinates on scri. This intersection is a two-surface in a three-dimensional space and can be thought as one member of a series of “wavefronts” obtained by slicing the lightcone of a point $`x^a`$ with a one-parameter family of past lightcones, the last one of them being scri itself. In this respect, the two-surface at scri given by $`u=Z(x^a,\zeta ,\overline{\zeta })`$ can be thought of as a two-dimensional wavefront in a three-dimensional space, for which there is a standard treatment in singularity theory.
Wavefronts are considered as projections of smooth two-dimensional Legendrian manifolds in a five-dimensional space down to a three-dimensional space. The singularities of the wavefront are the places where the projection is singular. The Legendrian manifold itself is obtained from a generating function. More specifically, consider the 1-jet bundle over the two dimensional configuration space $`(x^1,x^2)`$, given in coordinates $`(z,x^1,x^2,p_1,p_2)`$, and a function $`S(p_1,x^2)`$. A Legendrian manifold is a two dimensional subspace of points $`(z,x^1,x^2,p_1,p_2)`$ that can be specified by
$$x^1=\frac{S}{p_1},p_2=\frac{S}{x^2},z=S(p_1,x^2)x^1p_1,\text{ parametrized by }x^2\text{ and }p_1\text{.}$$
(7)
The projection of this Legendrian manifold down to $`(z,x^1,x^2)`$ is a two-surface in three dimensions, representing a wavefront. This wavefront is singular where the projection breaks down, namely at points such that the $`3\times 2`$ Jacobian matrix
$$\left(\begin{array}{cc}\frac{x^1}{x^2}& \frac{x^1}{p_1}\\ & \\ \frac{x^2}{x^2}& \frac{x^2}{p_1}\\ & \\ \frac{z}{x^2}& \frac{z}{p_1}\end{array}\right)$$
(8)
has rank less than 2.
In our case, we take $`(x^1,x^2)`$ as coordinates on the sphere by defining $`\zeta =x^1+ix^2`$ and interpret the projection of $`z`$ as the lightcone cut function $`Z(x^a,\zeta ,\overline{\zeta })`$ (the dependence on the spacetime points $`x^a`$ is considered parametric in the instance of lightcone cuts, being regarded as fixed here). The lightcone cut is regular except at some values of $`\zeta `$ at which the projection breaks down. Carrying through the calculation of the determinant of the Jacobian matrix in these terms, we obtain the result that the projection breaks down at points $`\zeta `$ such that
$$\frac{1}{ð^2Z(x^a,\zeta ,\overline{\zeta })}=0\text{and}\frac{1}{ð\overline{ð}Z(x^a,\zeta ,\overline{\zeta })}=0.$$
(9)
This calculation is parametric in $`x^a`$, where $`x^a`$ represents the apex of the lightcone intersecting scri at the lightcone cut surface. This is equivalent to the statement that both $`\mathrm{\Lambda }`$ and $`R`$ must blow up at particular values of $`\zeta `$ for the lightcone cut of a given spacetime point $`x^a`$ (See Fig.2). For every point on the cut, the quantities $`\mathrm{\Lambda }`$ and $`R`$ have finite values, except at the singular points shown in Fig.2. Given the lightcone cut function $`Z(x^a,\zeta ,\overline{\zeta })`$ for fixed $`x^a`$, the values of $`\zeta `$ for which $`\mathrm{\Lambda }`$ and $`R`$ diverge determine the null geodesics for which $`x^a`$ is a focal point. This is consistent with the result found in the previous section.
### C Singularities of the lightcones
A related issue which we are also concerned with is the location of the singularities of the lightcone of a point in the interior spacetime. Consider the lightcone of a point, namely, the congruence of all the null geodesics through that point, and follow one null geodesic in the congruence out to the future. Generically, due to curvature, there is at least one future point along this null geodesic where neighboring geodesics intersect, referred to as a point conjugate to the apex. The cross sectional area of the congruence vanishes at points which are conjugate to the apex along any null geodesic of the congruence, the locust of all such points being sometimes referred to as the caustic surface (although it would be more accurate to call it the singularity of the lightcone). The lightcone folds and self-intersects beyond the occurrence of the earliest of the points conjugate to the apex, as shown in Fig. 2. How do we formulate the condition for the location of such singular points in terms of null-surface variables?
We can approach this issue from an analysis of the geodesic deviation vector along a fixed null ray $`\zeta `$ of the lightcone of an interior point $`x_0^a`$. The geodesic deviation vector in this case can be derived solely from knowledge of $`Z(x^a,\zeta ,\overline{\zeta })`$ through a procedure that is standard to the null-surface approach. The coordinate transformation (1) can, in principle, be piece-wise inverted, yielding, perhaps with different branches,
$$x^a=f^a(u,\omega ,\overline{\omega },R,\zeta ,\overline{\zeta })$$
(10)
for fixed $`\zeta `$. The lightcone of a point $`x_0^a`$ can be obtained now by substituting
$$u=Z(x_0^a,\zeta ,\overline{\zeta })\omega =ðZ(x_0^a,\zeta ,\overline{\zeta })\overline{\omega }=\overline{ð}Z(x_0^a,\zeta ,\overline{\zeta })$$
(11)
and
$$R=ð\overline{ð}Z(x_0^a,\zeta ,\overline{\zeta })+r$$
(12)
into (10), in this manner obtaining, perhaps with several branches,
$$x^a=F^a(r;x_0^a,\zeta ,\overline{\zeta })$$
(13)
as the lightcone of the point $`x_0^a`$ parametrized by the directions $`(\zeta ,\overline{\zeta })`$ labeling the null geodesics, and the parameter $`r`$ along each null geodesic. The geodesic deviation vector is
$$M^a=ðF^a=\frac{f^a}{\theta ^i}ð\theta ^i+ð^{}f^a$$
(14)
where $`ð^{}`$ is taken keeping $`\theta ^i`$ fixed. This expression can be worked out straightforwardly to yield $`M^a`$ in terms of $`Z`$ and its derivatives, and the cross sectional area of the congruence along the null ray $`\zeta `$ can subsequently be obtained:
$$A_{lc}(r;x_0^a,\zeta ,\overline{\zeta })=\frac{(\mathrm{\Lambda }\mathrm{\Lambda }_0)(\overline{\mathrm{\Lambda }}\overline{\mathrm{\Lambda }}_0)(RR_0)^2}{\mathrm{\Omega }^2\sqrt{1\mathrm{\Lambda },_1\overline{\mathrm{\Lambda }},_1}}$$
(15)
The sublabel $`0`$ indicates evaluation at $`r=0`$ (the apex). The quantities $`\mathrm{\Lambda }_0`$ and $`R_0`$ appearing in (15) can be thought of as referring to the lightcone cut of the apex, whereas $`\mathrm{\Lambda }`$ and $`R`$ correspond to the lightcone cut of the point at $`r`$ along the null geodesic labeled by $`\zeta `$ (See Fig. 3). By comparing with (2) we can see that
$$A_{lc}=A_pH(r;x_0^a,\zeta ,\overline{\zeta })\text{with}H(r;x_0^a,\zeta ,\overline{\zeta })(\mathrm{\Lambda }\mathrm{\Lambda }_0)(\overline{\mathrm{\Lambda }}\overline{\mathrm{\Lambda }}_0)(RR_0)^2.$$
(16)
This equation relates the cross sectional areas of two different congruences containing the same null ray $`\zeta `$, namely the cone of lightrays through $`x_0^a`$ and the congruence of asymptotically parallel rays parallel to the null ray $`\zeta `$. This implies a relationship between the focusing of either congruence and the behavior of the quantities $`\mathrm{\Lambda }`$ and $`R`$. We distinguish three alternatives.
1. $`A_{lc}=0`$ and $`A_p=0`$ at some value $`r`$. Then $`H`$ must not diverge at a rate faster than $`A_p^1`$ at that point.
2. $`A_{lc}=0`$ with $`A_p0`$ at some value $`r`$. Then $`H`$ must vanish at that point.
3. $`A_p=0`$ with $`A_{lc}0`$ at some value $`r`$. Then $`H`$ must blow up at a rate faster than $`A_p^1`$ at that point.
As an example of the first instance, if $`A_p=0`$ at the apex then the apex is a focal point, because $`H=0`$, and is therefore finite, at $`r=0`$ (See Fig 4). So should be any point conjugate to it along the null ray $`\zeta `$; however, at such points $`\mathrm{\Lambda }_0,R_0,\mathrm{\Lambda }`$ and $`R`$ all blow up, therefore $`H`$ becomes quite intractable and its behavior has yet to be verified (See Fig. 5).
In the second case, the apex is not a focal point, and its conjugate point is located at $`r`$ such that
$$H(r;x_0^a,\zeta ,\overline{\zeta })=0.$$
(17)
This is the generic situation illustrated in Fig.2. See also Fig. 6.
In the third case, the point $`r`$ at which $`A_p`$ vanishes is a focal point along the null ray $`\zeta `$, and therefore $`\mathrm{\Lambda }`$ and $`R`$ both blow up, thus again $`H`$ becomes intractable but is not unlikely that it blows up, since $`\mathrm{\Lambda }_0`$ and $`R_0`$ are finite in this case. See Fig. 6.
This analysis applies to every null ray $`\zeta `$ in the lightcone. For some null ray $`\zeta _m`$ the conjugate point occurs closest to the apex, for some other direction $`\zeta _f`$ the apex has a conjugate point at infinity, and finally there are values of $`\zeta `$ for which no focusing takes place at any value of $`r`$, as can be seen from Fig.2.
### D Conclusion
Although the work reported on here is very much in progress, we believe we are finding significant clues as to what consequences the occurrence of focusing of null geodesics has for the null-surface formulation of general relativity. Our ultimate goal is to find a way to integrate the singularity issue with the null-surface dynamical equations, which played no role in this discussion.
This work has been supported by the NSF under grant No. PHY 92-05109.
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# Addemdum to: ”The Mathematical Structure of Quantum Superspace as a Consequence of Time Asymmetry”.
## I A rough coincidence becomes more precise.
In paper one of us reported a rough coincidence between the time where the minimum of the entropy gap $`\mathrm{\Delta }S=S_{act}S_{\mathrm{max}}`$ , takes place and the time where all the stars will exhaust their fuel. The time where the minimum of $`\mathrm{\Delta }S`$ is located was:
$$t_{cr}t_0\left(\frac{2}{3}\frac{\omega _1}{T_0}\frac{t_{NR}}{t_0}\right)^3$$
(1)
The following numerical values were chosen: $`\omega _1=T_{NR},`$ the temperature of the nuclear reactions within the stars (that was considered as the main source of entropy), $`t_{NR}=\gamma ^1`$ the characteristic time of these nuclear reactions, $`t_0`$ the age of the universe, and $`T_0,`$ the cosmic micro-wave background temperature, and making some approximations the rough coincidence was obtained.
Now we have reconsider the problem and conclude that, even if nuclear reactions within the stars are a source of entropy, the parameters $`T_{NR}`$ and $`t_{NR}`$ are not the good ones to define the behavior of the term $`e^{\gamma t/2}\rho _1`$ of equation (100) of paper , since they do not correspond to the main unstable system that we must consider. In fact the main production of entropy in a star is not located in its core, where the temperature is almost constant (and equal to $`T_{NR})`$, but in the photosphere where the star radiates. The energy radiated from the surface of the star is produced in the interior by fusion of light nuclei into heavier nuclei. Most stellar structures are essentially static, so the power radiated is supplied at the same rate by these exothermic nuclear reactions that take place near the center of the star . We can decompose the whole star in two branch systems, as explained in section VII of paper , where a chain of branch systems was introduced. We have two branch systems to study: the core and the photosphere. The core gives energy to the photosphere and in turn the photosphere diffuses this energy to the surroundings of the star, namely in the bath of microwave radiation at temperature $`T_{0\text{ .}}`$In this way, we have two sources of entropy production: the radiation of energy at the surface of the star and the change of composition inside the star (as time passes we have more helium and less hydrogen). Since the core of a star is near thermodynamic equilibrium, we neglect the second and we concentrate on the first: the radiation from the surface of the star (related with the difference between the star and the background temperatures). So the temperature of the photosphere and not the one of the core must be introduced in our formula. Thus it is better to consider the photosphere as the unstable system that defines the term $`e^{\gamma t/2}\rho _1`$ of equation (100) . So we must change $`T_{NR}`$ and $`t_{NR}`$ by $`T_P`$, the temperature of the photosphere and $`t_S`$ the characteristic lifetime of the star. Then we must change eq. (1) to:
$$t_{cr}t_0\left(\frac{2}{3}\frac{T_P}{T_0}\frac{t_S}{t_0}\right)^3$$
(2)
As the 90% of the stars are dwarfs with photosphere temperature $`T_P=10^3K`$ and the characteristic lifetime $`t_S=10^9`$ if we take these values we reach again to.
$$t_{cr}10^4t_0$$
(3)
but now with no approximation. The order of magnitude of $`t_{cr}`$ is a realistic one. In fact, $`10^4t_01.5\times 10^{14}years`$ after the big-bang the conventional star formation will end and it is also considered that all the stars will exhaust their fuel so it is reasonable that this time would be of the same order than the one where the entropy gap stops its decreasing and begins to grow. So the rough coincidence it is now a precise order of magnitude coincidence and therefore the comprehension of paper is improved.
## II Acknowledgments.
We wish to thank Omar Benvenutto for fruitful discussions. This work was partially supported by grants Nos. CI1\*-CT94-0004 of the European Community, PID-0150 and PEI-0126-97 of CONICET (National Research Council of Argentina), EX-198 of the Buenos Aires University and 12217/1 of Fundación Antorchas and the British Council.
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# Fermi-Edge Singularities In AlxGa1-xAs Quantum Wells : Extrinsic Versus Many-Body Scattering Processes
\[
## Abstract
A Fano resonance mechanism is evidenced to control the formation of optical Fermi-edge singularities in multi-subband systems such as remotely doped Al<sub>x</sub>Ga<sub>1-x</sub>As heterostructures. Using Fano parameters, we probe the physical nature of the interaction between Fermi sea electrons and empty conduction subbands. We show that processes of extrinsic origin like alloy-disorder prevail easily at 2D over multiple diffusions from charged valence holes expected by many-body scenarios.
\] Drastic manifestations are expected from the interaction of a magnetic or charged impurity with a Fermi sea of electrons. Since the successfull explanation of the divergent X-ray absorption edges of simple metals by Mahan and Nozières et al. thirty years ago, a Fermi sea of electrons under optical excitation has been recognized as a model system to study these issues experimentally. Indeed a many-body electronic state can develop in presence of the positively charged electronic vacancy (or core hole) involved in optical processes, which multiply scatters conduction electrons throughout and along the Fermi surface. This induces divergences at the Fermi-edge of optical spectra, the so-called Fermi-edge singularities (FES). FES are by essence highly sensitive to phase space restrictions such as hole localization and reduced dimensionalities : they cannot form without either strong hole localization or the dimensionality being unity.
Electron systems embedded in semiconductor heterostructures have therefore attracted much attention to test these predictions. FES were first observed by Skolnick et al. in the low-temperature Photoluminescence (PL) spectrum of a remotely doped In<sub>0.47</sub>Ga<sub>0.53</sub>As/InP quantum well (QW). Emphasis was put on the anomalous temperature dependence of the singularity and on the localization of valence holes by alloy fluctuations. More recently, experiments by Chen et al. in In<sub>0.15</sub>Ga<sub>0.85</sub>As/AlGaAs QWs with weaker hole localization also put forward the tunability of FES, by bringing the first empty QW subband into resonance with the Fermi level. It was then proposed that empty conduction subbands could act as additional scattering channels for Coulomb processes, in qualitative agreement with many-body multi-subband numerical calculations in the infinite hole-mass approximation.
Although this interpretation is widely referred to, no experiments have been undertaken to test many-body schemes beyond qualitative agreement. While attempting to do so with Al<sub>x</sub>Ga<sub>1-x</sub>As QWs, it appeared to us that non-coulombian intersubband scatterings can induce FES as well. The scope of this Letter is to describe this issue quantitatively. Our experiments demonstrate that extrinsic processes like alloy-disorder can prevail easily over multiple Coulomb scatterings of Fermi sea electrons predicted in the framework of many-body scenarios.
The paper is organized as follows. We show that the formation of FES in Al<sub>x</sub>Ga<sub>1-x</sub>As QWs is governed by a Fano resonance mechanism between Fermi-sea electrons and discrete excitonic transitions associated with empty conduction subbands. This model is first validated by the existence of scaling properties of PL spectra when FES are enhanced by reduced intersubband spacings. It is then used to gain insight into the microscopic nature of conduction Intersubband Couplings (ICs) at work. Indeed, the stronger the ICs, the more divergent the FES, so that one can probe directly in experiments the efficiency of extrinsic scattering processes such as alloy-disorder, remote-doping disorder or artificial ICs in lateral superlattices. Experimental results fall in close agreement with microscopic Fano calculations. We demonstrate that alloy-disorder is the dominant contribution to observed FES in Al<sub>x</sub>Ga<sub>1-x</sub>As QWs.
Samples investigated in this work are remotely doped Al<sub>1-x</sub>Ga<sub>x</sub>As/Al<sub>0.33</sub>Ga<sub>0.67</sub>As QWs grown on GaAs substrates by molecular beam epitaxy. They are all of same thickness L<sub>z</sub>=25nm, spacers (in the range 5-8 nm) and sheet density $`n_s8.10^{11}`$cm<sup>-2</sup>, but vary in their aluminium content $`x`$ ($`2.3\%x7.1\%`$). Confined 2D-levels are represented in fig. 1a. The asymmetry of the QW potential originates in the dipole formed by remote ionized dopants ($`z>0`$) and the Degenerate Electron Gas (DEG) partially filling the first conduction subband E<sub>1</sub>. This confers a PL activity to the two first subbands E<sub>1</sub> and E<sub>2</sub> with photocreated valence holes localized on potential fluctuations at the top of the heavy-hole HH<sub>1</sub> subband (fig. 1b). We perform PL spectroscopy at 1.8K. Samples are optically excited by a Ti:Sa laser at 1.7 eV, with a $``$1 W.cm<sup>-2</sup> density so as to avoid inhomogeneous heating of the DEG within the 40$`\mu m^2`$ laser spot. The DEG PL (E<sub>1</sub>HH<sub>1</sub>) extends from lowest wave-vector transitions at E<sub>g</sub> up to Fermi wave-vector transitions at E<sub>g</sub>+E<sub>F</sub> (E$`{}_{F}{}^{}25`$meV). Without any influence of E<sub>2</sub>HH<sub>1</sub>, its oscillator strength decreases monotonously with energy, because of hole localization and indirect optical processes (see spectrum $``$ in fig.1c). The PL of the empty QW subband E<sub>2</sub> exhibits a dominant excitonic feature E$`_{2_X}`$ of high oscillator strength and discrete character, visible either by thermal activation above E<sub>g</sub>+E<sub>F</sub> in PL, or in PL Excitation spectra. We define $`\mathrm{\Delta }`$=E$`_{2_X}`$-E<sub>g</sub>-E<sub>F</sub>. Variations of $`\mathrm{\Delta }`$ are achieved along a given sample by use of the flux gradients of effusion cells in the epitaxy chamber. We stop the wafer rotation during the spacer layer growth between dopants and the QW. The thinner the spacer, the stronger the electric field at the QW interface. This tunes the E<sub>2</sub>-E<sub>1</sub> energy separation, while E<sub>F</sub> hardly changes under illumination. As seen from PL spectra of sample A ($`x`$=7.1$`\%`$) in fig.1c, a FES forms and develops when $`\mathrm{\Delta }`$ is decreased to zero.
To interpret these data, we consider the Fano resonance model depicted in fig.2a. E$`_{2_X}`$ is taken as a discrete level coupled with a matrix-element $`𝒲`$ to the continuum of E<sub>1</sub>HH<sub>1</sub> electron-hole pairs. By assuming an infinite hole mass, all physical parameters simply refer to the conduction band : $`𝒲`$ equals to the IC between E<sub>1</sub> and E<sub>2</sub>, and the E<sub>1</sub>HH<sub>1</sub> continuum is populated by the E<sub>1</sub> Fermi-Dirac electrons. We take a parabolic dispersion and a constant PL oscillator strength for E<sub>1</sub>. Fano gave an analytical description of the spreading of a discrete level coupled to a continuum. Optical FES occur (see fig.2a) due to the partial filling of E<sub>1</sub>HH<sub>1</sub> near E$`_{2_X}`$, and get more divergent as $`\mathrm{\Delta }`$ is reduced. A qualitative agreement with experiments is thus obtained, without any particular assumption on the physical nature of $`𝒲`$.
The validation of the Fano resonance model comes from the scaling property evidenced in fig.2a that all PL spectra for various $`\mathrm{\Delta }`$ should display a common envelope lineshape when plotted with E$`_{2_X}`$ as the origin of energies. The comparison with experimental data is not straightforward because E<sub>1</sub>-E<sub>2</sub> varies experimentally, while the E<sub>1</sub>HH<sub>1</sub> PL oscillator strength intrinsically depends on energy. We compensate for these nominal E<sub>1</sub>HH<sub>1</sub> PL variations by dividing all data of fig.1c by the E<sub>1</sub>HH<sub>1</sub> spectrum with $`\mathrm{\Delta }\mathrm{}`$ ($``$ in fig.1c). Data are then plotted in fig. 2b with E$`_{2_X}`$ as the origin of energies. Remarkably, all spectra superpose on each other in the range of populated electron states. No clear transition exists from convergent to divergent Fermi edges in fig 2b. This means that FES only appear in raw PL data when the intrinsic PL decay of E<sub>1</sub>HH<sub>1</sub> at E<sub>F</sub> gets balanced by the positive slope of the Fano envelope for small $`\mathrm{\Delta }`$ values.
In order to get a quantitative analysis of the envelope profile, we assume a statistical disorder property for the interaction $`𝒲`$ : $`<𝒲>`$= 0 and $`<𝒲^2>0`$. This accounts for alloy disorder, which is later evidenced to dominate FES in sample A. The normalized PL intensity can then be derived analytically from ref. :
$$I(E)=f(E).\frac{q^2+(EE_{2_X})^2/\mathrm{\Gamma }^2}{1+(EE_{2_X})^2/\mathrm{\Gamma }^2}$$
where E is the PL energy ; $`f(E)`$ the occupation number in the E<sub>1</sub> subband ; $`\mathrm{\Gamma }`$ is the observed half width at half maximum of the resonance, only related to the statistical squared interaction average $`<𝒲^2>`$ and the density of states $`𝒟`$ of the E<sub>1</sub> conduction subband by $`\mathrm{\Gamma }=\pi <𝒲^2>𝒟`$ ; $`q^2`$ is the experimental oscillator strength of the excitonic resonance relative to the continuum, inversely proportional to $`<𝒲^2>`$ and $`𝒟^2`$. It depends also strongly on wavefunction overlaps through the ratio of the oscillator strengths of E<sub>2</sub>HH<sub>1</sub> and E<sub>1</sub>HH<sub>1</sub>. $`q^2`$ is therefore quite sensitive to small geometrical fluctuations from sample to sample.
Our data are nicely fitted by this model (fig. 2b), with Fano parameters $`\mathrm{\Gamma }`$=0.60 meV and $`q`$=12.2. All fits use Fermi-Dirac distributions of electrons with an effective temperature T=9.0$`\pm `$1K. This shows that electrons are indeed thermalized, though not with the lattice at 1.8 K due to the incomplete and slow relaxation of photocreated electrons above the E<sub>2</sub> subband edge. The consistence of small $`\mathrm{\Delta }`$ data with our fit indicates that the E$`_{2_X}`$ resonance remains weakly populated enough to stay in the “atom-like” regime. Also, taking $`f(E)`$=1 in the Fano lineshape formula, we can check that it fits the FES “enhancement factor” (FES PL intensity divided by its $`\mathrm{\Delta }\mathrm{}`$ value) measured by Chen et al. in In<sub>0.15</sub>Ga<sub>0.85</sub>As QW structures. Their data are indeed nicely reproduced with $`\mathrm{\Gamma }`$=0.66 meV and $`q=6.6`$ (fig.2c).
Before quitting the phenomenological level, we focus on the temperature dependence of FES in systems with explicit intersubband interaction. Observed thermal quenchings can be understood simply, even though such quenchings are expected from many-body theories, where an actual occupation number discontinuity is required to enhance multiple Coulomb scatterings at E<sub>F</sub>. Here or in ref., the FES disappears, only because the PL relative minimum between E<sub>g</sub>+E<sub>F</sub> and E$`_{2_X}`$ vanishes with raised temperatures.
Up to this point, we successfully assessed the model with respect to $`\mathrm{\Delta }`$ and temperature variations. We now analyse the microscopic origin of the Fano parameter $`\mathrm{\Gamma }`$. On the theoretical side, $`\mathrm{\Gamma }`$ only depends on the strength $`𝒲`$ of intersubband couplings and can be computed for a given microscopic process. Experimentally, the Fano model predicts that the stronger $`\mathrm{\Gamma }`$, the more divergent the Fermi-edges at fixed $`\mathrm{\Delta }`$. By designing appropriate samples, one can thus test : i) whether extrinsic scattering processes like alloy-disorder play an effective role on the formation of FES ; ii) if experimental variations of $`\mathrm{\Gamma }`$ correlate with microscopic calculations.
We display in fig.3a the PL spectra of remotely doped quantum wells C, B and A of with QW aluminium content $`x`$ equal to 0.023, 0.044 and 0.071 respectively, taken at fixed $`\mathrm{\Delta }`$=4.3$`\pm `$0.1 meV. As seen from the global PL lineshapes, the localization of photocreated valence holes remains constant, dominated by the roughness at the QW interface (see fig. 1) rather than by random alloy potential fluctuations. Enhancement of many-body processes due to hole localization can therefore be excluded. Nevertheless, FES get more pronounced with increased alloy concentration $`x`$, while the excitonic resonance broadens. Fitted $`\mathrm{\Gamma }`$ Fano parameters linearly increase with $`x`$ within experimental uncertainty (fig.3b). This explains the quadratic enhancement of FES visible from fig.3a when $`\mathrm{\Gamma }`$ is linearly increased in the regime where $`\mathrm{\Delta }/\mathrm{\Gamma }1`$. It also demonstrates that alloy disorder is the dominant contribution to the IC parameter $`𝒲`$.
To quantify this, we calculate $`\mathrm{\Gamma }`$ microscopically in the infinite hole-mass approximation :
$$\mathrm{\Gamma }_{alloy}=x(1x)m\mathrm{\Omega }_o\delta V^2/L_z\mathrm{}^2$$
This applies for a square quantum well of width $`L_z`$, a conduction electron mass $`m`$, $`\delta V`$ being the conduction band offset between pure AlAs and pure GaAs, and $`\mathrm{\Omega }_o`$ the crystal cell volume. With $`L_z`$=18 nm (representative of the confinement length of E<sub>1</sub> and E<sub>2</sub> wave-functions), $`m`$=0.07 $`m_0`$ and $`x`$=7.1$`\%`$, we obtain $`\mathrm{\Gamma }`$=0.61 meV. This quantitative agreement is striking for such a simple model. The fit from fig.3b is achieved by introducing a residual scattering $`\mathrm{\Gamma }_o`$=0.2 meV in the pure GaAs limit and thus taking a Fano parameter $`\mathrm{\Gamma }^2=\mathrm{\Gamma }_o^2+\mathrm{\Gamma }_{alloy}^2`$. We can also estimate $`\mathrm{\Gamma }_{alloy}`$$``$0.5 meV for the In<sub>0.15</sub>Ga<sub>0.85</sub>As QWs of ref., in close agreement with either our fit (fig.2c) or the empirical two-level coupling (0.6 meV) measured by Chen et al..
We now focus on ICs induced by random positioning of ionized dopants. Strictly speaking, $`𝒲`$ now depends on the occupation of E<sub>1</sub> states. We nonetheless assume $`𝒲`$ equal to its value for Fermi wave-vector (k<sub>F</sub>) electrons. $`\mathrm{\Gamma }`$ can then be computed, and gets proportional to $`\mathrm{exp}2k_F.z_o`$ in the limit of a thick spacer layer $`z_o`$, with a prefactor of $``$60 meV for a square QW structure of width L<sub>z</sub>=20 nm and a sheet density $`n_s10^{12}`$ cm<sup>-2</sup>. Due to high doping ($`k_F2.10^8`$m<sup>-1</sup>), this predicts a poor efficiency even for very shallow spacers. This is already visible from fig. 2b where indirect optical processes \- of same physical origin - only affect small wave-vector conduction states. Also, we measured no increase of $`\mathrm{\Gamma }`$ in a sample similar to C but with a 2.5 nm spacer.
We finally mention the case of FES in tilted lateral superlattices, where artificial intersubband couplings are created by a non separable 1D periodic confinement between the growth ($`z`$) and an in-plane ($`x`$) direction. This has been shown to promote optical FES. In fact, our data also fit to a Fano scheme with parameter $`\mathrm{\Gamma }`$=3.2 meV. The larger strength of ICs compared to 2D QWs explains the formation of pronounced FES even for large $`\mathrm{\Delta }`$ parameters. The experimental $`\mathrm{\Gamma }`$ value matches a microscopic calculation, using a typical value of 30 meV for the peak-to-peak lateral confinement amplitude.
Microscopic calculations of Fano parameters fall therefore in quantitative agreement with our experiments, where ICs have been varied in physical nature and over one decade in amplitude. This demonstrates that non-coulombian scattering processes can efficiently control the formation of FES in multi-subband semiconductor structures, in the simple picture of a band-structure partially filled by a DEG. We stress that Coulomb processes actually depend on the accurate distribution of charged particles with respect to empty subbands and do not quantitatively match Fano lineshapes, even if a qualitative analogy was underlined in early many-body interpretations. The question of extrinsic ICs can also be raised about ref.. A FES enhancement of a factor $``$2 would correspond to $`q\mathrm{\Delta }/\mathrm{\Gamma }`$ for a Fano process involving empty subbands. Estimating $`\mathrm{\Delta }`$30 meV and $`\mathrm{\Gamma }_{alloy}=`$1.0 meV, this criterium discards random alloy-disorder processes provided $`q30`$. An evaluation of actual disorder processes and relevant absorption data on E$`_{2_X}`$ are definitely lacking to give any conclusive answer, but the QW doping on both sides used in ref. should reduce advantageously the E<sub>2</sub> and HH<sub>1</sub> overlap and thus $`q`$. Bringing $`q`$-values down by optical selection rules indeed suppresses excitonic enhancement effects.
In conclusion, this paper shows that extrinsic intersubband scatterings must be carefully evaluated while analysing FES. This applies especially to lower dimensionality issues since disorder gets generally enhanced, and conduction subband spacings lowered. Realistic many-body theories should thus include disorder and multiple subbands. Our work should extend under magnetic field where the interaction of excitons and a DEG is debated. Extremely low-disorder structures seem required for clear electron-electron manifestations.
We thank F. Petit, P. Denk, F. Lelarge, A. Cavanna, R. Planel, C. Tanguy, B. Jusserand and B. Etienne for stimulating discussions. This work was supported by DRET and a SESAME grant of the Région Ile de France.
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# Condensate oscillations, kinetic equations and two-fluid hydrodynamics in a Bose gas
## 1 Introduction
Trapped Bose-condensed atomic gases are remarkable because, in spite of the fact that these are very dilute systems, they exhibit robust coherent dynamic behaviour when perturbed. These quantum “wisps of matter” are a new phase of highly coherent matter. While binary collisions are very infrequent, the coherent mean field associated with the Bose condensate ensures that interactions play a crucial role in determining the collective response of these superfluid gases.
In our discussion of the theory of collective oscillations of atomic condensates, which is the main theme of these lectures, the macroscopic Bose wavefunction $`\mathrm{\Phi }(𝐫,t)`$ will play a central role. This wavefunction is the BEC order parameter. The initial attempts at defining this order parameter began with the pioneering work of Fritz London in 1938, was further developed by Bogoliubov in 1947 and finally formalized in the general quantum field theoretic formalism of Beliaev in 1957. The first extension of these ideas to inhomogeneous Bose condensates was by Pitaevskii and, independently, by Gross in 1961, which led to the now famous Gross-Pitaevskii (GP) equation of motion for $`\mathrm{\Phi }(𝐫,t)`$. Most of this early work was limited to $`T=0`$ where, in a dilute Bose gas, one can assume all of the atoms are in the condensate.
In Section 2, I will first review the dynamics of a pure condensate at $`T=0,`$ based on solving the linearized GP equation. This limit is especially appealing since one can ignore all the complications which arise from the presence of non-condensate atoms (the thermal cloud). We will discuss the normal mode solutions of the $`T=0`$ GP equation of motion using the “quantum hydrodynamic” formalism, which works in terms of the local condensate density $`n_c(𝐫,t)`$ and superfluid velocity $`𝐯_c(𝐫,t)`$. Within the Thomas-Fermi approximation, Stringari has shown that the equations of motion for these two variables can be combined to give a wave equation for oscillations of the condensate.
In Section 3, we derive a generalized form of the GP equation for $`\mathrm{\Phi }(𝐫,t)`$ which is valid at finite temperatures. It involves terms which are coupled to the non-condensate component (the thermal cloud) and thus its solution in general depends on knowing the equations of motion for the dynamics of the non-condensate. In the present lectures, we will restrict ourselves to finite temperatures where the non-condensate can be described by a quantum kinetic equation for the single-particle distribution function $`f(𝐩,𝐫,t).`$ In this Boltzmann equation, the relatively high energy non-condensate atoms are simply free atoms moving in a self-consistent Hartree-Fock mean field.
A unique feature of a Bose-condensed gas is that the kinetic equation for $`f(𝐩,𝐫,t)`$ involves a collision integral (denoted by $`C_{12}[f,\mathrm{\Phi }])`$ describing collisions between condensate and non-condensate atoms. The generalized GP equation for $`\mathrm{\Phi }(𝐫,t)`$ also has a term which is related to the $`C_{12}`$ collisions. This gives rise to damping of condensate fluctuations. In Section 3, we derive a finite $`T`$ Stringari wave equation with damping using the static Popov approximation. This means the thermal cloud is treated as always being in static thermal equilibrium, with $`f(𝐩,𝐫,t)=f_0(𝐩,𝐫)`$ being given by the equilibrium Bose distribution.
In Section 4, we turn to a detailed treatment of the dynamics of the coupled condensate and non-condensate components, starting from the finite $`T`$ generalized GP equation for $`\mathrm{\Phi }(𝐫,t)`$ and our kinetic equation for $`f(𝐩,𝐫,t)`$. We derive a new set of two-fluid hydrodynamic equations, following the recent work by Zaremba, Nikuni and Griffin. This derivation assumes the non-condensate is in local hydrodynamic equilibrium, induced by rapid collisions between the atoms in the thermal cloud. As a result, the non-condensate is completely described in terms of the non-condensate density $`\stackrel{~}{n}(𝐫,t),`$ velocity $`𝐯_n(𝐫,t)`$ and local pressure $`\stackrel{~}{P}(𝐫,t)`$. Both the condensate and non-condensate exhibit coupled coherent oscillations at the same frequency.
In Section 5, we review the famous two-fluid hydrodynamic equations first derived by Landau in 1941 and which form the basis of our understanding of the hydrodynamic behaviour of superfluid <sup>4</sup>He. We discuss first sound and second sound and consider how they differ in superfluid <sup>4</sup>He and in a uniform Bose gas. We also point out the appearance of a new hydrodynamic zero frequency mode which appears in our generalized two-fluid equations. This mode is associated with the relaxation time for the condensate and non-condensate atoms to come into “diffusive” equilibrium with each other. Inclusion of damping due to thermal conductivity of the thermal cloud shows that this mode is the analogue of the well-known thermal diffusion mode in a classical gas.
We close this Introduction with some general references which may be useful to the reader. The 1998 Varenna Summer School lectures on BEC has excellent articles on recent research on BEC in atomic gases, including very detailed reviews of the experimental work carried out by the JILA and MIT groups. For a general introduction to the theory of trapped Bose gases, we recommend the Reviews of Modern Physics article by a leading theoretical group at the University of Trento. This authoritative review concentrates on thermodynamic properties and the collisionless dynamics. A recent long paper by Zaremba, Nikuni and the author discusses the details of the derivation of the coupled hydrodynamic equations for the condensate and the thermal cloud. The present lectures may be viewed as an introduction to this paper, with more emphasis on the general structure of the derivation and the implications of this new set of superfluid equations.
## 2 Pure condensate dynamics (at $`T=0`$)
This section will be largely a review of standard material but will provide the starting point for generalizations in the following sections. As we have mentioned, at $`T=0`$, one can assume that all atoms in a dilute Bose gas are described by the macroscopic wavefunction $`\mathrm{\Phi }(𝐫,t)`$. This obeys the time-dependent Hartree equation of motion first written down 40 years ago by Pitaevskii and Gross,
$$i\mathrm{}\frac{\mathrm{\Phi }}{t}(𝐫,t)=\left[\frac{\mathrm{}^2^2}{2m}+V_{ex}(𝐫)+V_H(𝐫,t)\right]\mathrm{\Phi }(𝐫,t).$$
(1)
Here the trap harmonic potential is
$$V_{ex}(𝐫)=\frac{1}{2}m\omega _0^2r^2\text{(isotropic)}$$
(2)
and the self-consistent condensate Hartree potential is
$$V_H(𝐫,t)=𝑑𝐫^{}v(𝐫𝐫^{})n_c(𝐫,t)=gn_c(𝐫,t).$$
(3)
As usual, since we are interested in extremely low energy atoms, we use the $`s`$wave approximation and approximate the interatomic potential by the pseudopotential
$`v(𝐫𝐫^{})`$ $`=`$ $`{\displaystyle \frac{4\pi a\mathrm{}^2}{m}}\delta (𝐫𝐫^{})`$ (4)
$`=`$ $`g\delta (𝐫𝐫^{}),`$
where $`a`$ is the correct $`s`$wave scattering length. For further discussion of the effective interaction to use in the Gross-Pitaevskii (GP) equation, we refer to the lectures by Burnett in this volume. The condensate density is given by $`n_c(𝐫,t)=|\mathrm{\Phi }(𝐫,t)|^2`$ and hence (1) reduces to a NLSE for $`\mathrm{\Phi }(𝐫,t),`$ namely
$$i\mathrm{}\frac{\mathrm{\Phi }(𝐫,t)}{t}=\left[\frac{\mathrm{}^2_r^2}{2m}+V_{ex}(𝐫)+g|\mathrm{\Phi }(𝐫,t)|^2\right]\mathrm{\Phi }(𝐫,t).$$
(5)
Since all atoms are in the identical quantum state, there is no exchange mean field in (1).
The $`T=0`$ GP equation in (5) has been the subject of literally hundreds of papers since the discovery of BEC in laser-cooled trapped atomic gases - and the equation appears in almost every chapter in this book. As discussed in recent reviews , it gives an excellent quantitative description of both the static and dynamic (linear and non-linear) behaviour in trapped Bose gases below about $`T\stackrel{<}{}\mathrm{\hspace{0.17em}0.5}T_{BEC}.`$ The accuracy can be of the order of a few percent, which is as much as one can expect since the non-condensate fraction of atoms at $`T=0`$ is also estimated to be about $`1\%`$ or so. The $`T=0`$ GP equation has been extended to deal with two-component Bose gases (involving two different atomic hyperfine states) and the effect of perturbations related to laser and rf fields (used to manipulate the Bose condensates) is easily incorporated. For further discussion and references, see the lectures by Ballagh in this book.
Our main purpose in this section is to use the GP equation (5) to discuss collective oscillations of a pure condensate. In Section 3, we will discuss the extension of this equation at finite $`T`$ to deal with the effect of non-condensate atoms. We first briefly review the static equilibrium solution $`\mathrm{\Phi }_0(𝐫)`$ of the GP equation, given by
$$\widehat{\psi }(𝐫,t)\mathrm{\Phi }(𝐫,t)=\mathrm{\Phi }_0(𝐫)e^{i\mu t/\mathrm{}},$$
(6)
where $`\mu `$ is the chemical potential of the condensate. The physics behind this can be seen from
$`N1|\widehat{\psi }(𝐫,t)|N`$ $`=`$ $`e^{iE_{N1}t/\mathrm{}}N1|\widehat{\psi }(𝐫)|Ne^{iE_Nt/\mathrm{}}`$ (7)
$`=`$ $`N1|\sqrt{N}|N1e^{i(E_NE_{N1})t/\mathrm{}}`$
$`=`$ $`\sqrt{N}e^{i\mu t/\mathrm{}}.`$
Using (6) in (5) gives
$$\mu _{c0}\mathrm{\Phi }_0(𝐫)=\left[\frac{\mathrm{}^2^2}{2m}+V_{ex}(𝐫)+g|\mathrm{\Phi }_0(𝐫)|^2\right]\mathrm{\Phi }_0(𝐫).$$
(8)
Assuming a vortex-free ground state, this GP equation for the static condensate wavefunction $`\mathrm{\Phi }_0(𝐫)=\sqrt{n_{c0}(𝐫)}`$ leads to the following expression for the equilibrium condensate chemical potential
$$\mu _{c0}=\frac{\mathrm{}^2^2\sqrt{n_{c0}(𝐫)}}{2m\sqrt{n_{c0}(𝐫)}}+V_{ex}(𝐫)+gn_{c0}(𝐫).$$
(9)
A standard approximation in solving (9) is to ignore the kinetic energy associated with the condensate amplitude $`\sqrt{n_c},`$ ie, neglect the $`\frac{\mathrm{}^2^2}{2m}`$ term. In this “Thomas-Fermi” (TF) approximation, the static GP equation for $`\mathrm{\Phi }_0(𝐫)`$ reduces to
$$\left[V_{ex}(𝐫)+g|\mathrm{\Phi }_0(𝐫)|^2\right]=\mu _{c0},$$
(10)
which is easily inverted to give the condensate density profile
$`n_{c0}(𝐫)`$ $`=`$ $`{\displaystyle \frac{1}{g}}[\mu _{c0}V_{ex}(𝐫)]`$ (11)
$`=`$ $`{\displaystyle \frac{1}{g}}\left[\mu _{c0}{\displaystyle \frac{1}{2}}m\omega _0^2r^2\right].`$
Clearly in the TF approximation, the size of the condensate is $`R_{TF}`$, where
$$\mu _{c0}=\frac{1}{2}m\omega _0^2R_{TF}^2.$$
(12)
One finds $`\mu _{c0}`$ from the condition $`𝑑𝐫n_{c0}(𝐫)=N=N_c,`$ which gives
$$\mu _{c0}=\frac{\mathrm{}\omega _0}{2}\left[15\frac{N_ca}{a_{HO}}\right]^{2/5};a_{HO}(\mathrm{}/m\omega _0)^{1/2}.$$
(13)
We recall that the oscillator length $`a_{HO}`$ is the size of the ground state Gaussian wavefunction of an atom in a parabolic potential. Combining (12) and (13) gives
$`R_{TF}`$ $`=`$ $`a_{HO}\left(15{\displaystyle \frac{N_ca}{a_{HO}}}\right)^{1/5}`$ (14)
$``$ $`a_{HO},\text{if}{\displaystyle \frac{N_ca}{a_{HO}}}1.`$
The TF approximation (11) for $`n_{c0}(𝐫)`$ is very good for large $`N_c`$, except for a small region near the edge of condensate $`(rR_{TF}).`$ In practice, for typical values of $`a_{HO}`$ and $`a`$, one finds the TF approximation is very good for $`N_c\stackrel{>}{}\mathrm{\hspace{0.17em}10}^4`$ atoms.
To discuss condensate fluctuations around the static equilibrium value of $`\mathrm{\Phi }_0(𝐫)`$, we first reformulate the time-dependent GP equation in terms of the condensate density and phase variables
$$\mathrm{\Phi }(𝐫,t)=\sqrt{n_c(𝐫,t)}e^{i\theta (𝐫,t)}.$$
(15)
Inserting this into the GP equation (5), gives
$`i\mathrm{}{\displaystyle \frac{\sqrt{n_c}}{t}}\sqrt{n_c}{\displaystyle \frac{\mathrm{}\theta }{t}}=\mu _c(𝐫,t)\sqrt{n_c}`$ $``$ $`i\mathrm{}^2{\displaystyle \frac{\sqrt{n_c}}{2m}}^2\theta +{\displaystyle \frac{\mathrm{}^2\sqrt{n_c}}{2m}}(\mathbf{}\theta )^2`$ (16)
$``$ $`{\displaystyle \frac{i\mathrm{}^2}{m}}(\mathbf{}\sqrt{n_c})\mathbf{}\theta ,`$
where we have defined
$$\mu _c(𝐫,t)\frac{\mathrm{}^2^2\sqrt{n_c}}{2m\sqrt{n_c}}+V_{ex}(𝐫)+gn_c(𝐫,t).$$
(17)
Separating out the real and imaginary path gives two equations. The real part gives
$$\mathrm{}\frac{\theta }{t}=(\mu _c+\frac{1}{2}m𝐯_c^2),$$
(18)
where the condensate velocity field is defined by
$$m𝐯_c(𝐫,t)\mathrm{}\mathbf{}\theta (𝐫,t).$$
(19)
The imaginary part of (16) gives
$$i\frac{\sqrt{n_c}}{t}=\frac{\sqrt{n_c}}{2}\mathbf{}𝐯_c𝐯_c\mathbf{}\sqrt{n_c},$$
(20)
which can be rewritten in the form
$$\frac{n_c}{t}=n_c\mathbf{}𝐯_c𝐯_c\mathbf{}n_c=\mathbf{}(n_c𝐯_c).$$
(21)
¿From now on, we shall set $`\mathrm{}=1`$ except in some final formulas.
In summary, we have shown that the time-dependent GP equation for the two-component order parameter $`\mathrm{\Phi }(𝐫,t)`$ is completely equivalent to the following two coupled equations for the condensate density $`n_c(𝐫,t)`$ and velocity field $`𝐯_c(𝐫,t)`$:
$`{\displaystyle \frac{n_c}{t}}=\mathbf{}(n_c𝐯_c)`$
$`m\left({\displaystyle \frac{𝐯_c}{t}}+{\displaystyle \frac{1}{2}}\mathbf{}𝐯_c^2\right)=\mathbf{}\mu _c,`$ (22)
where the position and time-dependent generalized condensate chemical potential $`\mu _c`$ is defined in (17). In this formulation, a complete description of the condensate dynamics is given in terms of two variables $`n_c(𝐫,t)`$ and $`𝐯_c(𝐫,t)`$, reminiscent of the hydrodynamic equations for a classical fluid. In the recent BEC literature, the equations in (22) are often referred to as the “hydrodynamic” theory. A better description would be to call it the “quantum hydrodynamic” theory, since it is equivalent to the mean-field GP equation. In later sections, we shall see that an exact GP equation taking into account the dynamics of the non-condensate leads to a generalized set of equations quite analogous to (22). The key equation (19) defining the superfluid velocity field of the condensate will turn out to be quite general. Needless to say, all aspects related to superfluidity of a trapped Bose gas (as compared to Bose condensation) are tied to the fact that the condensate exhibits motion related to the gradient of a phase (which means that, ignoring vortices, the condensate motion is irrotational, $`\mathbf{}\times 𝐯_c(𝐫,t)=0).`$
Stringari first pointed out that within the dynamic TF approximation, one could combine the two equations in (22) into a single condensate wave equation. This approach will be the basis of our development in these lectures. Neglecting the kinetic energy term $`^2\sqrt{n_c}`$ in (17) as small compared to the condensate interaction energy $`gn_c`$, we linearize the resulting equations around the equilibrium values
$`n_c`$ $`=`$ $`n_{c0}+\delta n_c`$
$`𝐯_c`$ $`=`$ $`𝐯_{c0}+\delta 𝐯_c.`$ (23)
We obtain
$`{\displaystyle \frac{\delta n_c}{t}}`$ $`=`$ $`\mathbf{}(n_{c0}\delta 𝐯_c)\mathbf{}(𝐯_{c0}\delta n_c)`$
$`{\displaystyle \frac{\delta 𝐯_c}{t}}`$ $`=`$ $`\mathbf{}\left(\mu _{c0}^{}+g\delta n_c+m𝐯_{c0}\delta 𝐯_c\right),`$ (24)
where $`\mu _{c0}^{}\mu _{c0}+\frac{1}{2}m𝐯_{c0}^2`$ is independent of position. Assuming that there is no vortex in the solution of the static GP equation (ie, $`𝐯_{c0}=0),`$ (24) reduces to the following coupled linearized equations for $`\delta n_c`$ and $`\delta 𝐯_c`$:
$`{\displaystyle \frac{\delta n_c}{t}}`$ $`=`$ $`\mathbf{}(n_{c0}(𝐫)\delta 𝐯_c)`$
$`{\displaystyle \frac{\delta 𝐯_c}{t}}`$ $`=`$ $`{\displaystyle \frac{g}{m}}\mathbf{}\delta n_c.`$ (25)
These can be combined to give the well known $`T=0`$ Stringari wave equation
$$\frac{^2\delta n_c}{t^2}=\frac{g}{m}\mathbf{}\left[n_{c0}(𝐫)\mathbf{}\delta n_c\right].$$
(26)
Since the derivation used the TF approximation (where $`n_{c0}(𝐫)`$ is given by (11) for $`rR_{TF}),`$ one can rewrite (26) in the equivalent form (for an isotropic trap potential)
$$\frac{^2\delta n_c}{t^2}=\frac{\mu _{c0}}{m}\mathbf{}\{[[1\frac{r^2}{R_{TF}^2}]\mathbf{}\delta n_c\},rR_{TF},$$
(27)
where $`\mu _{c0}`$ is given by (13). It turns out that the normal mode solutions $`\delta n_c(𝐫,t)=\delta n_\omega (𝐫)e^{i\omega t}`$ of (27) have frequencies which are independent of the interaction strength $`g`$ or the value of $`N_c`$. This is a feature of the underlying TF approximation, which typically starts to breakdown (as noted earlier) when $`N_c\stackrel{<}{}\mathrm{\hspace{0.17em}10}^4`$ atoms. This is shown by explicit numerical solutions of the coupled Bogoliubov equations of motion which describe the normal mode solutions of the linearized GP equation when we take the kinetic energy of the condensate amplitude $`\sqrt{n_c}`$ fully into account. For values of $`N_c\stackrel{<}{}\mathrm{\hspace{0.17em}10}^4`$, the normal mode frequencies depend significantly on the magnitude of $`N_c`$, as shown in Fig. 1.
We also note that (26) can be equally well rewritten in terms of the superfluid velocity $`\delta 𝐯_c`$ or, equivalently, the phase fluctuations $`\delta \theta `$. This emphasizes that the measured condensate density fluctuations are directly related to the existence of phase fluctuations. The existence of collective modes of a pure condensate may thus be viewed already as “evidence” of superfluidity, the latter being always a consequence of the phase coherence of the macroscopic wavefunction given in (15) which gives rise to the irrotational velocity in (19).
We conclude this section with several examples of condensate normal modes based on solving the $`T=0`$ Stringari equation (26). A wonderful aspect about the collective oscillations of a condensate in a trapped gas is you can “see” them. As Ketterle has remarked, these condensates are robust \- one can kick them, shake them and these “wisps” of Bose-condensed matter keep their integrity. The uniform Bose-condensed gas is especially simple, since $`\delta n_c=\delta n_{k\omega }e^{i(𝐤𝐫\omega t)}.`$ This gives
$$\omega ^2\delta n_{k\omega }=\frac{gn_{c0}}{m}(k^2)\delta n_{k\omega },$$
(28)
or $`\omega ^2=c_0^2k^2,`$ with $`c_0\sqrt{gn_{c0}/m}.`$ This recovers the well-known Bogoliubov phonon oscillations of a uniform Bose condensate. The neglect of the kinetic energy in our TF approximation precludes us from obtaining the particle-like behaviour at large values of the wavevector $`k`$
The Kohn (or sloshing) mode corresponds to the oscillation of the centre-of-mass of the static condensate profile with the trap frequency $`\omega _0.`$ This mode is described by
$$n_c(𝐫,t)=n_{c0}(𝐫𝜼(t)),$$
(29)
where $`\frac{d𝜼(t)}{dt}𝐯_c(t)`$ and
$$\frac{d^2𝜼(t)}{dt^2}=\omega _0^2𝜼(t).$$
(30)
The proof is simple. Linearizing (29), we find
$$n_c(𝐫,t)=n_{c0}(𝐫)𝜼\mathbf{}n_{c0}(𝐫),$$
(31)
which gives the explicit form for the condensate fluctuation,
$$\delta n_c(𝐫,t)=\frac{1}{g}m\omega _0^2𝐫𝜼(t).$$
(32)
Inserting this into the Stringari equation (26), we obtain
$`\mathbf{}n_{c0}{\displaystyle \frac{d^2𝜼(t)}{dt^2}}`$ $`=`$ $`\omega _0^2\mathbf{}\left[n_{c0}\mathbf{}(𝜼(t)𝐫)\right]`$ (33)
$`=`$ $`\omega _0^2\mathbf{}\left[n_{c0}𝜼\right]`$
$`=`$ $`\omega _0^2𝜼\mathbf{}n_{c0}.`$
This confirms that the centre-of-mass position $`𝜼(t)`$ of the static condensate distribution satisfies the SHO equation (30) with frequency $`\omega _0`$.
The breathing (or monopole) condensate normal mode corresponds to a velocity fluctuation of the form
$$\delta 𝐯_c(𝐫,t)=A𝐫e^{i\omega t}(rR_{TF}).$$
(34)
Using $`\mathbf{}𝐫=3`$ and $`\mathbf{}(r^2𝐫)=5r^2,`$ it is easy to verify from the continuity equation in (22) that the associated density fluctuation is
$`i\omega \delta n_c`$ $`=`$ $`\mathbf{}\left[n_{c0}(𝐫)\delta 𝐯_c\right]`$ (35)
$`=`$ $`{\displaystyle \frac{\mu _{c0}}{g}}A\left[35{\displaystyle \frac{r^2}{R_{TF}^2}}\right]e^{i\omega t},`$
or
$$\delta n_\omega (𝐫)B\left(1\frac{5}{3}\frac{r^2}{R_{TF}^2}\right),$$
(36)
where $`\mu _{c0}`$ is given by (12). Inserting this into the Stringari wave equation (27) gives
$`\omega ^2\delta n_\omega (𝐫)`$ $`=`$ $`{\displaystyle \frac{\mu _{c0}}{m}}\mathbf{}\left[\left(1{\displaystyle \frac{r^2}{R_{TF}^2}}\right)B{\displaystyle \frac{5}{3}}{\displaystyle \frac{2𝐫}{R_{TF}^2}}\right]`$ (37)
$`=`$ $`5\omega _0^2\delta n_\omega (𝐫).`$
Thus the breathing mode for an isotropic parabolic trap has a frequency $`\omega =\sqrt{5}\omega _0.`$ As noted earlier, this mode frequency is not explicitly dependent on the interaction strength $`g`$. However the underlying theory is very dependent on mean-field effects. Moreover, we recall that a non-interacting trapped Bose gas has a breathing mode with a frequency $`\omega =2\omega _0`$ at all temperatures.
As a final example, we consider the so-called “surface” modes of a $`T=0`$ condensate, as described by
$$\delta 𝐯_c(𝐫,t)=A\mathbf{}\left[r^lY_{lm}(\theta ,\varphi )\right]e^{i\omega t}.$$
(38)
This mode corresponds to phase fluctuations of the form $`\delta \theta _\omega =mAr^lY_{lm}(\theta ,\varphi )`$. One may easily verify that $`\mathbf{}\delta 𝐯_c=0`$ and hence from (25) we obtain
$$\frac{^2\delta 𝐯_c}{t^2}=\omega _0^2\mathbf{}(\delta 𝐯_c𝐫).$$
(39)
This is equivalent to
$$\frac{^2\delta \theta }{t^2}=\omega _0^2(\mathbf{}\delta \theta )𝐫$$
(40)
or
$$\omega ^2\delta \theta _\omega (𝐫)=\omega _0^2l\delta \theta _\omega (𝐫).$$
(41)
Thus the surface oscillations of the condensate phase have a frequency given by $`\omega =\sqrt{l}\omega _0(l=1,2,3,\mathrm{}).`$
One great advantage of the quantum hydrodynamic formalism (within the TF approximation) is that it is easy to also treat the $`T=0`$ normal modes of an anisotropic parabolic well. We refer to the literature for further discussion. We also note that it is straightforward to derive a Stringari wave equation for the oscillations of a vortex state (this has been done by Svidzinsky and Fetter ).
## 3 Generalized GP equation at finite temperatures
In this section, we generalize the $`T=0`$ GP equation of Section 2 to finite $`T`$, so that it includes the effect of the coupling of the condensate to the non-condensate degrees of freedom. We concentrate in this Section on how the $`T=0`$ condensate modes are renormalized and damped by the thermal cloud. We defer discussion of the appearance of new collective modes mainly associated with the thermal cloud to sections 4 and 5.
Up to the beginning of 2000, the study of the dynamics of a trapped Bose gas at finite temperatures has been largely ignored by experimentalists but actively studied by many theorists (especially in the last year or so). One of the reasons for the lack of finite temperature data is that there are so many interesting phenomena to study at $`T0`$! However, another reason seems to be the implicit belief that the presence of the non-condensate just complicates the behaviour of a pure $`T=0`$ condensate - but is not the source of any interesting new physics. In this and subsequent sections, I hope to argue that this “belief” is quite wrong. The coupling of the condensate and non-condensate degrees of freedom at finite $`T`$ leads to a new two-component system in which both components can exhibit coherent behaviour, analogous to the well-known superfluid macroscopic phenomena appearing in liquid <sup>4</sup>He.
We shall see that, as expected, the finite temperature GP equation of motion for $`\mathrm{\Phi }(𝐫,t)`$ is not closed. It’s general solution involves knowing the equations of motion for the non-condensate. We will work within an approximation where the non-condensate atoms can be described by a quantum Boltzmann transport equation for the single-particle distribution function $`f(𝐩,𝐫,t),`$ with
$$\stackrel{~}{n}(𝐫,t)=\frac{d𝐩}{(2\pi )^3}f(𝐩,𝐫,t).$$
(42)
The details of this kinetic equation will be developed in Section 4. In this section, we concentrate on incorporating the effect of a static thermal cloud into the finite $`T`$ dynamics of the condensate.
The theory of interacting Bose-condensed fluids is most usefully discussed using quantum field operators. This procedure was formalized by Beliaev in 1957 and developed by Bogoliubov , Gavoret and Nozières , Martin and Hohenberg , and others in the early 1960’s. We recall:
$`\widehat{\psi }^{}(𝐫)=\text{creates atom at}𝐫`$
$`\widehat{\psi }(𝐫)=\text{destroys atom at}𝐫.`$ (43)
These quantum field operators satisfy the usual Bose commutation relations, such as
$$[\widehat{\psi }(𝐫),\widehat{\psi }^{}(𝐫^{})]=\delta (𝐫𝐫^{}).$$
(44)
All observables can be written in terms of these quantum field operators, such as the density $`\widehat{n}(𝐫)=\widehat{\psi }^{}(𝐫)\widehat{\psi }(𝐫)`$ and interaction energy
$`\widehat{V}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle 𝑑𝐫𝑑𝐫^{}\widehat{\psi }^{}(𝐫^{})\widehat{\psi }^{}(𝐫)v(𝐫𝐫^{})\widehat{\psi }(𝐫^{})\widehat{\psi }(𝐫)}`$ (45)
$`=`$ $`{\displaystyle \frac{1}{2}}g{\displaystyle 𝑑𝐫\widehat{\psi }^{}(𝐫)\widehat{\psi }^{}(𝐫)\widehat{\psi }(𝐫)\widehat{\psi }(𝐫)}.`$
The crucial idea due to Bogoliubov and later generalized by Beliaev is to separate out the condensate component of the field operators,
$$\widehat{\psi }(𝐫)=\widehat{\psi }(𝐫)+\stackrel{~}{\psi }(𝐫),$$
(46)
where
$$\widehat{\psi }(𝐫)\mathrm{\Phi }(𝐫)=\text{Bose macroscopic wavefunction}.$$
This quantity plays the role of the “order parameter” for the Bose superfluid phase transition:
$`\mathrm{\Phi }(𝐫)`$ $`=0T>T_c`$
$`0T<T_c.`$
We note that $`\mathrm{\Phi }(𝐫)\sqrt{n_c}e^{i\theta }`$ is a 2-component order parameter. Clearly, $`\mathrm{\Phi }(𝐫)`$ is not simply related to the many-particle wavefunctions $`\mathrm{\Psi }(𝐫_1,𝐫_2,\mathrm{}𝐫_N).`$ The thermal average in $`\widehat{\psi }(𝐫)`$ involves a small symmetry-breaking perturbation to allow $`\mathrm{\Phi }(𝐫)`$ to be finite,
$$\widehat{H}_{SB}=\underset{\eta 0}{lim}𝑑𝐫\left[\eta (𝐫)\widehat{\psi }^{}(𝐫)+\eta ^{}(𝐫)\widehat{\psi }(𝐫)\right].$$
(47)
The philosophy behind the concept of symmetry-breaking was extensively discussed by Bogoliubov , for a variety of condensed matter systems, in an article which is still highly recommended.
It is useful to make a few comments on the physics behind $`\mathrm{\Phi }(𝐫,t)`$. $`\mathrm{\Phi }(𝐫,t)`$ is a coherent state, with a “clamped” value of phase - rather than a Fock-state of fixed $`N`$, with no well-defined phase. $`\mathrm{\Phi }(𝐫,t)`$ acts like a classical field, since quantum fluctuations are negligible when $`N_c`$ is large. Probably P.W. Anderson deserves the greatest credit for understanding (in the period 1958-1963) the new physics behind working with a broken-symmetry state $`\mathrm{\Phi }(𝐫,t)`$, both in BCS superconductors and in superfluid <sup>4</sup>He. It captures the physics of the new phase of matter (such as the occurrence of the Josephson effect) and the associated superfluidity. The symmetry-breaking perturbation (47) allows the system to internally set up off-diagonal symmetry-breaking fields, which persist even when the external symmetry-breaking perturbation in (47) is set to zero at the end $`(\eta 0)`$. The same sort of physics is the basis of the BCS theory of superconductors.
One can formulate the GP and Bogoliubov approximations directly in terms of a variational many-particle wavefunction (see the lectures by Leggett in this volume). However, such formulations are limited to simple mean-field approximations. The explicit introduction of the broken-symmetry order parameter $`\mathrm{\Phi }(𝐫,t)`$ gives a systematic way of isolating the role of the Bose condensate in a general treatment of an interacting Bose-condensed fluid. As we shall see, the resulting formalism allows one to deal with questions related to damping as well as superfluidity in both the collisionless and hydrodynamic regions.
The exact Heisenberg equation of motion for the field operator is
$`i{\displaystyle \frac{\widehat{\psi }(𝐫,t)}{t}}`$ $`=`$ $`\left[{\displaystyle \frac{^2}{2m}}+V_{ex}(𝐫)+\delta U(𝐫,t)\right]\widehat{\psi }(𝐫,t)`$ (48)
$`+`$ $`\eta (𝐫)+g\widehat{\psi }^{}(𝐫,t)\widehat{\psi }(𝐫,t)\widehat{\psi }(𝐫,t),`$
where $`\delta U(𝐫,t)`$ is a small time-dependent driving potential. This gives an exact equation of motion for $`\mathrm{\Phi }(𝐫,t)\widehat{\psi }(𝐫,t),`$
$`i{\displaystyle \frac{\mathrm{\Phi }(𝐫,t)}{t}}`$ $`=`$ $`\left[{\displaystyle \frac{^2}{2m}}+V_{ex}(𝐫)+\delta U(𝐫,t)\right]\mathrm{\Phi }(𝐫,t)`$ (49)
$`+`$ $`\eta (𝐫)+g\widehat{\psi }^{}(𝐫,t)\widehat{\psi }(𝐫,t)\widehat{\psi }(𝐫,t),`$
with \[(using the decomposition (46)\]
$$\widehat{\psi }^{}\widehat{\psi }\widehat{\psi }=|\mathrm{\Phi }|^2\mathrm{\Phi }+2|\mathrm{\Phi }|^2\stackrel{~}{\psi }+\mathrm{\Phi }^2\stackrel{~}{\psi }^{}+\mathrm{\Phi }^{}\stackrel{~}{\psi }\stackrel{~}{\psi }+2\mathrm{\Phi }\stackrel{~}{\psi }^{}\stackrel{~}{\psi }+\stackrel{~}{\psi }^{}\stackrel{~}{\psi }\stackrel{~}{\psi }.$$
(50)
Taking the symmetry-breaking average of (50), one finds
$$\widehat{\psi }^{}\widehat{\psi }\widehat{\psi }=n_c\mathrm{\Phi }+\stackrel{~}{m}\mathrm{\Phi }^{}+2\stackrel{~}{n}\mathrm{\Phi }+\stackrel{~}{\psi }^{}\stackrel{~}{\psi }\stackrel{~}{\psi },$$
(51)
where
$`n_c(𝐫,t)`$ $``$ $`|\mathrm{\Phi }(𝐫,t)|^2=\text{condensate density}`$
$`\stackrel{~}{n}(𝐫,t)`$ $``$ $`\stackrel{~}{\psi }^{}(𝐫,t)\stackrel{~}{\psi }(𝐫,t)=\text{non-condensate density}`$
$`\stackrel{~}{m}(𝐫,t)`$ $``$ $`\stackrel{~}{\psi }(𝐫,t)\stackrel{~}{\psi }(𝐫,t)=\text{off-diagonal (anomalous) density}.`$ (52)
Using (51) in (49), our “exact” equation of motion for $`\mathrm{\Phi }(𝐫,t)`$ is
$`i{\displaystyle \frac{\mathrm{\Phi }}{t}}`$ $`=`$ $`\left[{\displaystyle \frac{^2}{2m}}+V_{ex}+gn_c(𝐫,t)+2g\stackrel{~}{n}(𝐫,t)\right]\mathrm{\Phi }`$ (53)
$`+`$ $`g\stackrel{~}{m}(𝐫,t)\mathrm{\Phi }^{}+g\stackrel{~}{\psi }^{}(𝐫,t)\stackrel{~}{\psi }(𝐫,t)\stackrel{~}{\psi }(𝐫,t).`$
We now consider various approximations to the generalized GP equation in (53):
1. The Hartree-Fock-Bogoliubov (HFB) approximation for $`\mathrm{\Phi }`$ corresponds to neglecting the three-field correlation function $`\stackrel{~}{\psi }^{}\stackrel{~}{\psi }\psi `$ but keeping the $`n_c,\stackrel{~}{n}`$ and $`\stackrel{~}{m}`$ fluctuations. This HFB has been exhaustively treated in many recent papers, in conjunction with the separate equations of motion for $`\stackrel{~}{n}`$ and $`\stackrel{~}{m}`$ This HFB approximation can be used to generate the same normal mode spectrum which the Beliaev second-order self-energy approximation gives at finite $`T`$
2. The dynamic Popov approximation corresponds to ignoring both $`\stackrel{~}{\psi }^{}\stackrel{~}{\psi }\stackrel{~}{\psi }`$ and $`\stackrel{~}{m}=\stackrel{~}{\psi }\stackrel{~}{\psi }`$ in (53). Theories of this kind involve coupled equations for $`\mathrm{\Phi }`$ and $`\stackrel{~}{n}`$
3. The static Popov approximation involves even a further simplification, namely it ignores fluctuations in the density $`\stackrel{~}{n}(𝐫,t)`$ of the thermal cloud,
$$\stackrel{~}{n}(𝐫,t)n_{c0}(𝐫).$$
(54)
Using (54) in (53) corresponds to treating the dynamics of the condensate moving in the static mean field of the non-condensate thermal cloud, ie,
$$i\mathrm{}\frac{\mathrm{\Phi }}{t}=\left[\frac{^2}{2m}+V_{ex}(𝐫)+2g\stackrel{~}{n}_0(𝐫)+gn_c(𝐫,t)\right]\mathrm{\Phi }(𝐫,t).$$
(55)
There is a considerable literature based on this static Popov-approximation. In a nutshell, in this section, we will discuss an extension of (55) which involves including the damping associated with the $`\stackrel{~}{\psi }^{}\stackrel{~}{\psi }\stackrel{~}{\psi }`$ term, but again treating the thermal cloud statically.
As in Section 2, we first rewrite (53) in terms of the condensate amplitude and phase. Inserting (15) into (53) and following the same procedure which led to the $`T=0`$ equations in (22), we obtain:
$$\frac{n_c}{t}+\mathbf{}n_c𝐯_c=2gIm\left[\mathrm{\Phi }^2\stackrel{~}{m}+\mathrm{\Phi }^{}\stackrel{~}{\psi }^{}\stackrel{~}{\psi }\stackrel{~}{\psi }\right]$$
(56)
$$\frac{\theta }{t}=\left(\mu _c+\frac{(\mathbf{}\theta )^2}{2m}\right),$$
(57)
where now (compare with (17))
$`\mu _c(𝐫,t)=`$ $``$ $`{\displaystyle \frac{^2\sqrt{n_c}}{2m\sqrt{n_c}}}+V_{ex}+gn_c+2g\stackrel{~}{n}`$ (58)
$`+`$ $`{\displaystyle \frac{g}{n_c}}Re\left[\mathrm{\Phi }^2\stackrel{~}{m}+\mathrm{\Phi }^{}\stackrel{~}{\psi }^{}\stackrel{~}{\psi }\stackrel{~}{\psi }\right].`$
Formally, (53) and hence (56) - (58) are exact results. In these lectures, we will limit ourselves to finite temperatures where the dominant thermal excitations are well approximated by high energy non-condensate atoms moving in a self-consistent dynamic HF mean-field
$`\stackrel{~}{\epsilon }_p(𝐫,t)`$ $`=`$ $`{\displaystyle \frac{p^2}{2m}}+V_{ex}(𝐫)+2g\left[n_c(𝐫,t)+\stackrel{~}{n}(𝐫,t)\right]`$ (59)
$``$ $`{\displaystyle \frac{p^2}{2m}}+U(𝐫,t).`$
In this region, we shall neglect the anomalous pair correlations among the thermal atoms (ie, $`\stackrel{~}{m}=0`$). One can show that both $`\stackrel{~}{m}`$ and $`\stackrel{~}{\psi }^{}\stackrel{~}{\psi }\stackrel{~}{\psi }`$ in (56) and (58) are of order $`g`$. Thus one sees the correction terms in (56) and (58) involving these functions are of order $`g^2`$. In fact, calculation shows that to order $`g`$, $`\stackrel{~}{\psi }^{}\stackrel{~}{\psi }\stackrel{~}{\psi }`$ is imaginary. Our procedure is to ignore the $`O(g^2)`$ terms in (56) and (58) except those which are imaginary and hence are a source of damping of condensate motion.
The end result of this approach is that we are left with
$`{\displaystyle \frac{n_c}{t}}+\mathbf{}n_c𝐯_c`$ $`=`$ $`\mathrm{\Gamma }_{12}[f,\mathrm{\Phi }]`$
$`m\left({\displaystyle \frac{𝐯_c}{t}}+{\displaystyle \frac{1}{2}}\mathbf{}𝐯_c^2\right)`$ $`=`$ $`\mathbf{}\mu _c,`$ (60)
with
$$\mu _c(𝐫,t)=\frac{^2\sqrt{n_c}}{2m\sqrt{n_c}}+V_{ex}(𝐫)+gn_c(𝐫,t)+2g\stackrel{~}{n}(𝐫,t)$$
(61)
and we have defined the new function
$$\mathrm{\Gamma }_{12}[f,\mathrm{\Phi }]2gIm\left[\mathrm{\Phi }^{}\stackrel{~}{\psi }^{}\stackrel{~}{\psi }\stackrel{~}{\psi }\right].$$
(62)
We see that $`\mathrm{\Gamma }_{12}`$ plays the role of a source term in the condensate continuity equation in (60) . It depends on both $`\mathrm{\Phi }(𝐫,t)`$ and the single-particle distribution function $`f(𝐩,𝐫,t).`$ It is useful to note that (60) and (61) are equivalent to the approximate GP equation (see also Refs. )
$$i\mathrm{}\frac{\mathrm{\Phi }}{t}=\left[\frac{\mathrm{}^2^2}{2m}+V_{ex}(𝐫)+gn_c(𝐫,t)+2g\stackrel{~}{n}(𝐫,t)i\mathrm{}R(𝐫,t)\right]\mathrm{\Phi },$$
(63)
where the dissipative function is
$$R(𝐫,t)\frac{\mathrm{\Gamma }_{12}[f,\mathrm{\Phi }]}{2n_c(𝐫,t)}0(g^2).$$
(64)
To proceed, we need to calculate the field correlation $`\stackrel{~}{\psi }^{}\stackrel{~}{\psi }\stackrel{~}{\psi }`$ which determines $`\mathrm{\Gamma }_{12}`$ in (62). This has been done in Appendix A of Ref. , working to lowest order in $`g`$ at finite $`T`$ where the single-particle spectrum in (59) is adequate. The result is
$`\stackrel{~}{\psi }^{}\stackrel{~}{\psi }\stackrel{~}{\psi }`$ $`=`$ $`ig{\displaystyle \frac{\mathrm{\Phi }(𝐫,t)}{(2\pi )^5}}{\displaystyle 𝑑𝐩_1𝑑𝐩_2𝑑𝐩_3}`$ (65)
$`\times `$ $`\delta (𝐩_c+𝐩_1𝐩_2𝐩_3)\delta (ϵ_c+\stackrel{~}{ϵ}_1\stackrel{~}{ϵ}_2\stackrel{~}{ϵ}_3)`$
$`\times `$ $`\left[f_1(1+f_2)(1+f_3)(1+f_1)f_2f_3\right],`$
where $`f_1f(𝐩,𝐫,t),ϵ_c\mu _c(𝐫,t)+\frac{1}{2}m𝐯_c^2(𝐫,t)`$ is the local condensate atom energy and $`𝐩_cm𝐯_c`$ is the condensate atom momentum. In Section 4, we shall see that (65) is closely related to the collision integral $`C_{12}[f,\mathrm{\Phi }]`$ which enters the Boltzmann equation for $`f(𝐩,𝐫,t)`$ and describes collisions with the non-condensate atoms in which one condensate atom is involved. We note that to leading order, the expression in (65) is imaginary. Using this in (62), we find
$`\mathrm{\Gamma }_{12}(𝐫,t)`$ $`=`$ $`2g^2{\displaystyle \frac{n_c(𝐫,t)}{(2\pi )^5}}{\displaystyle 𝑑𝐩_1𝑑𝐩_2𝑑𝐩_3\delta (𝐩_c+𝐩_1𝐩_2𝐩_3)}`$ (66)
$`\times `$ $`\delta \left(ϵ_c+\stackrel{~}{ϵ}_1\stackrel{~}{ϵ}_2\stackrel{~}{ϵ}_3\right)`$
$`\times `$ $`\left[f_1(1+f_2)(1+f_3)(1+f_1)f_2f_3\right].`$
As a first application of our new equations in (60), we will ignore the dynamics of the non-condensate cloud. That is, we will assume that the condensate interacts with a static thermal cloud in thermal equilibrium
$$f(𝐩,𝐫,t)f^0(𝐩,𝐫)=\frac{1}{e^{\beta [p^2/2m+U_0(𝐫)\stackrel{~}{\mu }_0]}1},$$
(67)
where $`\stackrel{~}{\mu }_0`$ is the equilibrium chemical potential of the non-condensate atoms. In Section 4, we will discuss the $`C_{22}`$ collisions which produce this static equilibrium Bose distribution. The value of $`\stackrel{~}{\mu }_0`$ is known since it must equal the condensate equilibrium chemical potential $`\mu _{c0}`$ given by the static solution of the generalized GP equation. Within the TF approximation, (61) reduces to
$$\mu _{c0}=V_{\mathrm{ex}}(𝐫)+gn_{c0}(𝐫)+2g\stackrel{~}{n}_0(𝐫)$$
(68)
and hence we see that (59) reduces to
$`\stackrel{~}{ϵ}_p(𝐫)\stackrel{~}{\mu }_0`$ $`=`$ $`{\displaystyle \frac{p^2}{2m}}+U_0(𝐫)\stackrel{~}{\mu }_0`$ (69)
$`=`$ $`{\displaystyle \frac{p^2}{2m}}+gn_{c0}(𝐫).`$
This result is consistent with the correct Bogoliubov “excitation energy.” The static non-condensate density in (68) is given by
$$\stackrel{~}{n}_0(𝐫)=\frac{1}{\mathrm{\Lambda }^3}g_{3/2}\left(z_0=e^{\beta gn_{c0}(𝐫)}\right).$$
(70)
We can now calculate $`\mathrm{\Gamma }_{12}(𝐫,t)`$ using the equilibrium Bose distribution (67) for the thermal cloud. One finds (for a related calculation, see Ref. ).
$$\mathrm{\Gamma }_{12}^0(𝐫,t)=\frac{n_c(𝐫,t)}{\tau _{12}(𝐫,t)}\left[e^{\beta (\stackrel{~}{\mu }_0\epsilon _c(𝐫,t)\frac{1}{2}mv_c^2)}1\right],$$
(71)
where we have defined a $`C_{12}`$ collision time
$`{\displaystyle \frac{1}{\tau _{12}(𝐫,t)}}`$ $``$ $`{\displaystyle \frac{2g^2}{(2\pi )^5}}{\displaystyle 𝑑𝐩_1𝑑𝐩_2𝑑𝐩_3\delta (𝐩_c+𝐩_1𝐩_2𝐩_3)}`$ (72)
$`\times `$ $`\delta (\epsilon _c+\stackrel{~}{\epsilon }_{p_1}\stackrel{~}{\epsilon }_{p_2}\stackrel{~}{\epsilon }_{p_3})(1+f_1^0)f_2^0f_3^0.`$
$`\mathrm{\Gamma }_{12}[f_0,\mathrm{\Phi }]`$ in (71) still depends on $`n_c(𝐫,t)`$ and $`𝐯_c(𝐫,t)`$ of the condensate through $`\mu _c,𝐩_c`$ and $`ϵ_c`$. We note that when the condensate is in static equilibrium, $`\mu _c(𝐫,t)\mu _{c0}`$ and $`𝐯_{c0}=0`$. In this case, $`\mu _{c0}=\stackrel{~}{\mu }_0`$ and the expression in the square bracket in (71) is seen to vanish. Thus $`\mathrm{\Gamma }_{12}(f_0,\mathrm{\Phi }_0)=0,`$ as it should when both components are in static thermal equilibrium.
Summarizing, at this point we have a closed set of equations which can be used to describe the dynamics of the condensate in a trapped Bose gas which include the interactions with a static equilibrium thermal cloud. These equations are
$$\frac{n_c}{t}+n_c𝐯_c=\mathrm{\Gamma }_{12}[f^0,\mathrm{\Phi }]$$
(73)
$$m\left(\frac{}{t}+𝐯_c\mathbf{}\right)𝐯_c=\mathbf{}\mu _c,$$
(74)
where $`\mathrm{\Gamma }_{12}[f^0,\mathrm{\Phi }]`$ is given explicitly by the expression in (71) and
$$\mu _c(𝐫,t)=\frac{^2\sqrt{n_c}}{2m\sqrt{n_c}}+V_{ex}(𝐫)+2g\stackrel{~}{n}_0(𝐫)+gn_c(𝐫,t).$$
(75)
¿From (73) and (74), we can obtain linearized equations of motion for the condensate fluctuations $`\delta n_c`$ and $`\delta 𝐯_c`$. We use the fact that, to lowest order in the fluctuations from static equilibrium, (71) reduces to
$$\delta \mathrm{\Gamma }_{12}^0=\frac{\beta n_{c0}(𝐫)}{\tau _{12}^0(𝐫)}\delta \mu _c(𝐫,t),$$
(76)
where the “equilibrium” $`C_{12}`$ collision rate \[using (69) in (72)\] is defined by
$`{\displaystyle \frac{1}{\tau _{12}^0(𝐫)}}`$ $``$ $`{\displaystyle \frac{2g^2}{(2\pi )^5}}{\displaystyle 𝑑𝐩_1𝑑𝐩_2𝑑𝐩_3\delta (𝐩_1𝐩_2𝐩_3)}`$ (77)
$`\times `$ $`\delta \left({\displaystyle \frac{p_1^2p_2^2p_3^2}{2m}}gn_{c0}\right)(1+f_1^0)f_2^0f_3^0.`$
In the static TF approximation, we recall that the equilibrium distribution is given by $`f_i^0=[e^{\beta (p_i^2/2m+gn_{c0})}1]^1`$.
In the present discussion, we further restrict ourselves to the dynamic Thomas-Fermi limit valid for large $`N_c`$,
$`{\displaystyle \frac{\delta n_c}{t}}+\mathbf{}(n_{c0}\delta 𝐯_c)`$ $`=`$ $`{\displaystyle \frac{1}{\tau ^{}}}\delta n_c`$ (78)
$`m{\displaystyle \frac{\delta 𝐯_c}{t}}`$ $`=`$ $`g\mathbf{}\delta n_c.`$ (79)
The collision time $`\tau ^{}(𝐫)`$ describes collisions between the condensate and non-condensate atoms when the condensate is perturbed away from equilibrium,
$$\frac{1}{\tau ^{}(𝐫)}=\frac{gn_{c0}(𝐫)}{k_\mathrm{B}T}\frac{1}{\tau _{12}^0(𝐫)}.$$
(80)
The new term on the right-hand side of (78) causes damping of the condensate fluctuations due to the lack of collisional detailed balance between the condensate and the static thermal cloud. We note that this collision time $`\tau ^{}(𝐫)`$ is only a function of the position $`𝐫`$ through its dependence on the static condensate density $`n_{c0}(𝐫)`$.
We can easily combine (78) and (79) to obtain what we shall refer to as the finite $`T`$ Stringari wave equation
$$\frac{^2\delta n_c}{t^2}\frac{g}{m}\mathbf{}(n_{c0}\mathbf{}\delta n_c)=\frac{1}{\tau ^{}}\frac{\delta n_c}{t}.$$
(81)
If we neglect the right-hand side, we obtain the undamped finite $`T`$ Stringari normal modes $`\delta n_c(𝐫,t)=\delta n_i(𝐫)e^{i\omega _it}`$ given by the solution of
$$\frac{g}{m}\mathbf{}\left[n_{c0}(𝐫)\mathbf{}\delta n_i(𝐫)\right]=\omega _i^2\delta n_i(𝐫).$$
(82)
As has been noted by several authors in recent papers , $`n_{c0}(𝐫)`$ at finite $`T`$ can be well approximated by the TF condensate profile at $`T=0`$ but with the number of atoms in the condensate $`N_c(T)`$ now being a function of temperature. This is because the static mean field of the non-condensate plays a minor role. In the regions where the condensate density is finite, we effectively always have $`n_{c0}(𝐫)\stackrel{~}{n}_0(𝐫),`$ as illustrated in Fig. 2. With this approximation for $`n_{c0}(𝐫)`$, the solutions of the finite $`T`$ Stringari equation (81) will be identical to those at $`T=0`$ (see Section 2), since the $`T=0`$ Stringari frequencies do not depend on the magnitude of $`N_c`$. Of course, as shown by calculations solving the coupled Bogoliubov equations , when $`N_c\stackrel{<}{}\mathrm{\hspace{0.17em}10}^4`$ the TF approximation breaks down. Thus the condensate collective mode frequencies will always become temperature dependent close to $`T_{\mathrm{BEC}}`$, where the TF approximation is no longer valid.
We can use the undamped Stringari modes given by (82) as a basis set to solve (81) and find the damping of these modes. Writing $`\delta n_c(𝐫)=_ic_i\delta n_i(𝐫)`$, and using the orthonormality condition $`𝑑𝐫\delta n_i(𝐫)\delta n_j(𝐫)=\delta _{ij}`$, one obtains the following algebraic equations for the coefficients $`c_i`$
$$\omega ^2c_i=\omega _i^2c_ii\omega \underset{j}{}\gamma _{ij}c_j,$$
(83)
where
$$\gamma _{ij}𝑑𝐫\delta n_i(𝐫)\delta n_j(𝐫)/\tau ^{}(𝐫).$$
(84)
Assuming the damping is small (our present discussion is for the collisionless region), (83) is easily solved using perturbation theory by setting $`\gamma _{ij}=0`$ for $`ji`$. This gives the damped Stringari frequency (to lowest order) $`\mathrm{\Omega }_i=\omega _ii\mathrm{\Gamma }_i`$, with
$$\mathrm{\Gamma }_i\frac{\gamma _{ii}}{2}=\frac{1}{2}𝑑𝐫\frac{\delta n_i(𝐫)^2}{\tau ^{}(𝐫)}.$$
(85)
This result for $`\mathrm{\Gamma }_i`$ is reasonable, namely it involves a spatial average over $`1/\tau ^{}(𝐫)`$ weighted with respect to the undamped density fluctuations of the Stringari wave equation (81). Calculation shows that the effect of coupling to other modes ($`\gamma _{ij}0`$) is extremely small.
Before discussing the trapped gas, it is useful to first apply our theory to a homogeneous gas, where $`\tau ^{}`$ is independent of position. In this case, (85) simply reduces to $`\mathrm{\Gamma }_i=1/2\tau ^{}`$. Although our model applies only to the collisionless region, it is useful to compare the inter-component collision time for a uniform gas in the collisionless and hydrodynamic regimes. In Section 5, we show that the inter-component collision time $`\tau _\mu `$ in the hydrodynamic region is given by $`\tau _\mu =\sigma _H\tau ^{}`$, where the temperature-dependent factor $`\sigma _H`$ depends on various thermodynamic functions \[$`\sigma _H`$ is given explicitly by (164) in Section 4\]. In Fig. 3, we compare $`1/\tau ^{}`$ and $`1/\tau _\mu `$ as functions of $`T`$. We see that $`\sigma _H`$ dramatically alters the inter-component relaxation rate $`1/\tau _\mu `$ appropriate to the hydrodynamic regime, as compared to $`1/\tau ^{}`$ involved in the collisionless regime. For completeness, in Fig. 3 we also plot the often-used classical collision time as well as $`\tau _{12}^0`$ defined in (77).
We briefly discuss some numerical calculations of the inter-component damping using the above formalism. We consider <sup>87</sup>Rb atoms in a spherically symmetric trap with frequency $`\nu _0=10`$ Hz and $`N=2\times 10^6`$. In the collisionless limit, we require $`\omega _i\tau _{\mathrm{cl}}1`$. We obtain an upper limit on $`1/\tau _{\mathrm{cl}}`$ by taking the density in the center of the trap $`n(r=0)`$, which gives $`1/\tau _{\mathrm{cl}}=8a^2N\omega _0^3m/(\pi k_\mathrm{B}T)`$. For the parameters we use, $`\omega _{10}\tau _{\mathrm{cl}}19`$ (compared to $`\omega _{02}\tau _{\mathrm{cl}}20`$ for the JILA data and $`\omega _{02}\tau _{\mathrm{cl}}2`$ for the MIT data on collective oscillations at finite $`T`$).
In Fig. 4 we plot the damping rate $`\mathrm{\Gamma }_{10}`$ for the breathing mode ($`n=1`$, $`l=0`$) as a function of temperature up to $`T=0.95T_{\mathrm{BEC}}`$, where $`N_c7\times 10^4`$. At higher temperatures, the Thomas-Fermi approximation will start to break down and the mode frequencies become temperature dependent.
The damping of condensate modes we have discussed in this section is due to the fact that the condensate is out of diffusive equilibrium with the thermal cloud. At finite $`T`$, the collective oscillations of the coupled condensate and non-condensate can be generally split into two classes. One mode mainly involves (out-of-phase) motion of the condensate and is the finite $`T`$ generalization of the analogous $`T=0`$ condensate mode. For the same symmetry, there is another mode which mainly involves the (in-phase) motion of the non-condensate thermal cloud. This mode is naturally viewed as a generalization of the $`T>T_{BEC}`$ thermal cloud oscillation to temperatures below $`T_{BEC}.`$ These two types of modes (for a given symmetry) have been obtained in the collision-dominated hydrodynamic region, when the dynamics of the thermal cloud is fully allowed for. Treating (to first approximation) the thermal cloud (the non-condensate component) as always in static equilibrium in this section means that our theory is only applicable to condensate oscillations which involve motions which are out-of-phase with the thermal cloud.
As we discuss in Section 4, the possibility that the condensate may be out of diffusive equilibrium with the non-condensate is a general feature of the dynamics of trapped Bose-condensed gases. In the hydrodynamic two-fluid region, it leads to a new relaxational mode. This effect is missed in the well-known form of the two-fluid hydrodynamics developed by Landau , where one (implicitly) assumes that the normal and superfluid components are always in local diffusive equilibrium. We will return to this question in Sections 4 and 5.
Of course, in addition to the inter-component relaxation discussed in this section, one also has Landau and Beliaev damping of condensate oscillations. These mechanisms are treated in more detail by Burnett in this book. Landau damping is briefly discussed in Section 4.
Clearly one can extend the formalism of this section to any problem based on the $`T=0`$ GP equation. In particular, it can be used to discuss this kind of inter-component damping at finite $`T`$ of oscillations in two-component Bose gases, vortex dynamics , and Josephson oscillations between two traps.
## 4 Dynamics of the coupled condensate and non-condensate at finite temperatures
In contrast to Section 3, we now grapple with the dynamics of the thermal cloud. As we noted in Section 3, in dealing with this very complicated problem, we will be quite modest and treat the non-condensate using the simplest microscopic model which captures the important physics. We limit ourselves to finite $`T`$, where the non-condensate atoms can be described by the particle-like HF spectrum given by (59). For trapped Bose gases, this spectrum is probably adequate down to quite low temperatures, for reasons discussed in an important paper by the Trento group. To extend our present analysis to very low temperatures is much more involved since then the excitations of the thermal cloud take on a collective aspect (ie, a Bogoliubov-type quasiparticle spectrum must be used).
The single-particle spectrum (59) is appropriate in the semi-classical limit, where we can use the single-particle distribution function $`f(𝐩,𝐫,t)`$ given by the solution of a kinetic equation. This procedure generalizes the approach of Boltzmann (1880’s) for describing binary collisions in a classical gas. Such a quantum Boltzmann equation for a trapped Bose-condensed gas at finite temperatures has been derived and extensively discussed by Zaremba, Nikuni and the author. The conditions of validity are
$$k_BTgn_c,k_BT\mathrm{}\omega _0,$$
(86)
where $`\mathrm{}\omega _0`$ is the spacing of the energy levels of the harmonic trap. More general but less explicit discussions of kinetic equations are given in Refs. . Related work of Gardiner and coworkers is discussed in the lectures by Ballagh in this book. The quantum kinetic equation we use is given by
$`{\displaystyle \frac{f(𝐩,𝐫,t)}{t}}`$ $`+`$ $`{\displaystyle \frac{𝐩}{m}}\mathbf{}_rf(𝐩,𝐫,t)\mathbf{}_rU(𝐫,t)\mathbf{}_pf(𝐩,𝐫,t)`$ (87)
$`=`$ $`C_{22}[f]+C_{12}[f].`$
The right hand side describes how binary collisions effect the value of the single-particle distribution function $`f(𝐩,𝐫,t)`$. The effective time-dependent HF potential $`U(𝐫,t)`$ is defined in (59). The effect of collisions between excited atoms in the non-condensate is described by:
$`C_{22}[f]`$ $`=`$ $`{\displaystyle \frac{2g^2}{(2\pi )^5}}{\displaystyle 𝑑𝐩_2𝑑𝐩_3𝑑𝐩_4\delta (𝐩+𝐩_2𝐩_3𝐩_4)}`$ (88)
$`\times `$ $`\delta \left(\stackrel{~}{\epsilon }_p+\stackrel{~}{\epsilon }_{p_2}\stackrel{~}{\epsilon }_{p_3}\stackrel{~}{\epsilon }_{p_4}\right)`$
$`\times `$ $`\left[(1+f)(1+f_2)f_3f_4ff_2(1+f_3)(1+f_4)\right].`$
We recall that creating a Boson gives a factor $`(1+f)`$ and destroying a Boson gives $`f.`$ In the classical high temperature limit, $`f1`$ and the collision integral $`C_{22}`$ in (88) becomes much simpler.
In addition to $`C_{22}`$ collisions, we also have collisions which involve one condensate atom:
$`C_{12}[f]`$ $`=`$ $`{\displaystyle \frac{2g^2}{(2\pi )^2}}{\displaystyle 𝑑𝐩_1𝑑𝐩_2𝑑𝐩_3\delta (m𝐯_c+𝐩_1𝐩_2𝐩_3)}`$ (89)
$`\times `$ $`\delta \left(\epsilon _c+\stackrel{~}{\epsilon }_{p1}\stackrel{~}{\epsilon }_{p2}\stackrel{~}{\epsilon }_{p3}\right)\left[\delta (𝐩𝐩_1)\delta (𝐩𝐩_2)\delta (𝐩𝐩_3)\right]`$
$`\times `$ $`\left[n_c(1+f_1)f_2f_3n_cf_1(1+f_2)(1+f_3)\right].`$
For convenience, we recall that the condensate atom has
$`\text{energy:}\epsilon _c=\mu _c+{\displaystyle \frac{1}{2}}mv_c^2;\mu _c=V_{ex}+gn_c+2g\stackrel{~}{n}`$
$`\text{momentum:}𝐩_c=m𝐯_c`$ (90)
We note the key difference between $`C_{12}`$ and $`C_{22}`$ collisions:
* $`C_{22}`$ and $`C_{12}`$ both conserve energy and momentum in collisions.
* $`C_{12}`$ does not (but $`C_{22}`$ does) conserve the number of condensate atoms. $`C_{12}`$ describes how atoms are “kicked” in and out of the condensate.
It turns out in the generalized GP equation given in (63), the damping term $`iR(𝐫,t)`$ is closely related to $`C_{12}[f]`$. This makes sense, since the $`C_{12}`$ collisions modify the condensate described by $`\mathrm{\Phi }(𝐫,t).`$ One easily may verify that \[comparing (66) and (89)\]
$$\mathrm{\Gamma }_{12}[f,\mathrm{\Phi }]=\frac{d𝐩}{(2\pi )^3}C_{12}[f(𝐩,𝐫,t)].$$
(91)
More precisely, the three-field correlation function in the exact equation of motions \[see (56) and (58)\] is related to $`C_{12}[f]`$ by
$$\frac{d𝐩}{(2\pi )^3}C_{12}[f]=2g\sqrt{n_c}Im\stackrel{~}{\psi }^{}\stackrel{~}{\psi }\stackrel{~}{\psi }.$$
(92)
We have to solve the equation of motions for $`f(𝐩,𝐫,t)`$ and $`\mathrm{\Phi }(𝐫,t)`$ treating $`C_{12}[f]`$ very carefully. There will be an exchange of atoms between the $`\stackrel{~}{n}(𝐫,t)`$ and $`n_c(𝐫,t)`$ components through the $`C_{12}`$ collisions. We can use these coupled equations for a variety of problems. Recently these coupled equations have been used in two problems:
1. Derivation of a generalized set of two-fluid hydrodynamic equations.
2. Discussion of the rate of growth of a condensate due to a sudden quench in which the high energy spectrum of the thermal cloud distribution is suddenly removed. This work has many points of contact with the formalism reviewed by Ballagh in this book for condensate growth.
In this section, we will linearize these equations and consider the collective oscillations of the combined system composed of condensate and non-condensate. It is useful to introduce two regimes to describe collective modes in interacting systems :
1. Collisionless (produced by mean fields)
$$\omega \tau _R1\text{or}T\tau _R\left(\omega \frac{2\pi }{T}\right)$$
(93)
2. Hydrodynamic (produced by collisions)
$$\omega \tau _R1\text{or}T\tau _R,$$
(94)
where $`\tau _R`$ is some appropriate relaxation time. What should we use for $`\tau _R`$? For a uniform classical gas, this is the collision time
$$\frac{1}{\tau _c}=\sqrt{2}\stackrel{~}{n}\sigma \overline{v}$$
(95)
where
$`\sigma =8\pi a^2\text{(for Bose atoms)};a=s\text{-wave scattering length}.`$
$`\overline{v}\text{average velocity of atoms}\sqrt{{\displaystyle \frac{k_BT}{m}}}.`$
$`\stackrel{~}{n}=\text{density of excited atoms}.`$
Even for a Bose-condensed gas, taking $`\tau _R\tau _c`$ is a reasonable first estimate. To get into the interesting hydrodynamic region $`(\omega \tau _R1)`$, we need a small value of $`\tau _R`$, ie, a large density $`\stackrel{~}{n}`$ or a large collision cross-section $`\sigma `$ (perhaps by working near a Feshbach resonance). We note that current BEC experiments deal with $`N10^6`$ atoms or larger and thus the kind of density of thermal atoms needed to be in the hydrodynamic region now seem achievable.
Let us look at the kinetic equation (87), writing it in the schematic form:
$$\widehat{}f=C_{22}[f]+C_{12}[f].$$
(96)
In the collisionless region, we need only solve $`\widehat{}f=0.`$ For comparison with the collision-dominated hydrodynamic domain, we briefly discuss some features of this collisionless domain, limiting ourselves to a uniform gas for simplicity. This enables us to understand how Landau and Beliaev damping arise from the collisionless Boltzmann equation, in contrast to the inter-component collisional damping (discussed in Section 3) which has its origin in the $`C_{12}[f,\mathrm{\Phi }]`$ collision integral.
We consider linear response theory for the deviations of $`f(𝐩,𝐫,t)`$ from thermal equilibrium
$$f(𝐩,𝐫,t)=f_0(𝐩,𝐫)+\delta f(𝐩,𝐫,t).$$
(97)
For a uniform normal Bose gas to first order in $`\delta f,`$ the collisionless kinetic equation reduces to
$$\frac{\delta f}{t}+\frac{𝐩}{m}\mathbf{}_r\delta f\mathbf{}_r[2g\delta n(𝐫,t)]+\delta U(𝐫,t)]\mathbf{}_pf_0(𝐩)=0.$$
(98)
Assuming $`\delta U(𝐫,t)=\delta U_{k\omega }e^{i(𝐤𝐫\omega t)}`$, then we have $`\delta f(𝐩,𝐫,t)=\delta f_{k\omega }(𝐩)e^{i(𝐤𝐫\omega t)}`$, where $`\delta f_{k\omega }(𝐩)`$ satisfies
$$i\omega \delta f(𝐩)+i\frac{𝐩𝐤}{m}\delta f(𝐩)[2g\delta n_{k\omega }+\delta U_{k\omega }]i𝐤\mathbf{}_pf_0(𝐩)=0,$$
(99)
where
$$f_0(𝐩)=\frac{1}{e^{\beta (\frac{p^2}{2m}+gn_0\mu _0)}1}f_0(ϵ_p)$$
(100)
and
$$\delta n_{k\omega }\frac{d𝐩}{(2\pi )^3}\delta f_{k\omega }(𝐩).$$
(101)
Solving (99), one finds
$$\delta n_{k\omega }=\frac{d𝐩}{(2\pi )^3}\frac{𝐩𝐤}{m}\frac{f_0(ϵ_p)}{ϵ_p}\frac{[2g\delta n_{k\omega }+\delta U_{k\omega }]}{\omega \frac{𝐩𝐤}{m}}$$
(102)
or
$$\delta n_{k\omega }=\frac{\chi _{nn}^0(𝐤,\omega )}{12g\chi _{nn}^0(k,\omega )}\delta U_{k\omega }\chi _{nn}(𝐤,\omega )\delta U_{k\omega }.$$
(103)
Here
$$\chi _{nn}^0(𝐤,\omega )=\frac{d𝐩}{(2\pi )^3}\frac{𝐩𝐤}{m}\frac{f_0(ϵ_p)}{ϵ_p}\frac{1}{\omega \frac{𝐩𝐤}{m}}$$
(104)
is the $`k0`$ limit of the density response function of a non-interacting uniform Bose gas.
This shows how one can use the collisionless kinetic equation to find the density response function $`\chi _{nn}(𝐤,\omega )`$ of a uniform Bose gas. The density fluctuations are given by the poles of $`\chi _{nn}(𝐤,\omega )`$ in (103). In this RPA with exchange, the solutions of
$$12g\chi _{nn}^0(𝐤,\omega )=0$$
(105)
are called “zero sound” modes (following Landau’s terminology in the analogous calculation for an interacting Fermi gas ). It is clear that the solutions $`\omega (k)`$ of (105) will be Landau damped, with a width related to
$``$ $`2gIm\chi _{nn}^0(𝐤,\omega +i0^+)`$ (106)
$`=`$ $`2g\pi {\displaystyle \frac{d𝐩}{(2\pi )^3}[f_0(ϵ_p)f_0(ϵ_{p+k})]\delta [\omega (k)(ϵ_{p+k}ϵ_p)]}.`$
Clearly this damping will come from the zero sound mode absorbing a thermal excitation of energy $`ϵ_p`$ to create a thermal excitation $`ϵ_{p+k}`$,
$$\mathrm{}\omega (k)=ϵ_{p+k}ϵ_p..$$
(107)
We might also remark that if this linear response calculation is extended to deal with a Bose-condensed gas, one finds a new damping mechanism because now $`\chi _{nn}^0(𝐤,\omega )`$ describes non-interacting Bogoliubov excitations $`E_p`$. In addition to the contributions of the kind given by (107), one finds $`\chi _{nn}^0(𝐤,\omega )`$ now has additional poles at
$$\mathrm{}\omega (k)=E_{p+k}+E_p.$$
(108)
These give rise to “Beliaev” damping, where the zero sound mode can decay into two Bogoliubov excitations. In the limit of $`T0`$, Beliaev damping dominates over Landau damping. It was first calculated at $`T=0`$ by Beliaev in 1957 in his classic study of a weakly interacting Bose gas.
Several recent papers have discussed Landau and Beliaev damping of collective modes in a trapped Bose gas. This problem has been nicely formulated directly in terms of response functions by Minguzzi and Tosi in conjunction with mean fields produced by the condensate and non-condensate density fluctuations. Giorgini has given a physically appealing discussion working directly with the coupled mean field equations of motion for $`\mathrm{\Phi }(𝐫,t),\stackrel{~}{n}(𝐫,t)`$ and $`\stackrel{~}{m}(𝐫,t).`$ The approach of the Oxford group is reviewed in Burnett’s lectures in this book.
In this connection, a simple treatment of Landau damping of condensate collective modes can be given starting from our generalized GP equation in (63). In Section 3, we ignored the fluctuations in the HF term $`2g\stackrel{~}{n}(𝐫,t)`$ and concentrated on the damping associated with the $`iR(𝐫,t)`$ term in (69). A simple way of including the $`\delta \stackrel{~}{n}(𝐫,t)`$ fluctuations induced by the mean field associated with the condensate fluctuations is to use (see Refs. )
$$\delta \stackrel{~}{n}(𝐫,\omega )=𝑑𝐫^{}\chi _{\stackrel{~}{n}\stackrel{~}{n}}^0(𝐫,𝐫^{},\omega )2g\delta n_c(𝐫^{},\omega ).$$
(109)
Here $`\chi _{\stackrel{~}{n}\stackrel{~}{n}}^0`$ is the density response function for a non-interacting gas of atoms with the equilibrium HF spectrum given by (69) and the chemical potential $`\stackrel{~}{\mu }_0`$ is equal to $`\mu _{c0}`$ in (68). For a uniform Bose gas, the condensate collective modes given by (63) are found to be
$$\omega ^2=c_0^2k^2[1+4g\chi _{\stackrel{~}{n}\stackrel{~}{n}}^0(𝐤,\omega )]\frac{i\omega }{\tau ^{}}.$$
(110)
Here $`c_0=\sqrt{gn_{c0}/m}`$ is the Bogoliubov phonon velocity \[see (28)\] and $`\tau ^{}`$ is defined in (80). In the long wavelength limit, one finds
$$\underset{k0}{lim}Im\chi _{\stackrel{~}{n}\stackrel{~}{n}}^0(𝐤,\omega =c_0k)=\frac{mk_BT}{2\pi c_0}\frac{1}{3},$$
(111)
where the $`1/3`$ factor arises from the exchange term in the single-particle self-energy. Thus one obtains
$$\omega =c_0ki(\mathrm{\Gamma }_L+\frac{1}{2\tau ^{}}),$$
(112)
where the Landau damping is given by (RPA with exchange)
$$\mathrm{\Gamma }_L=\frac{4}{3}ak_BTk.$$
(113)
The exact result for $`\mathrm{\Gamma }_L`$ obtained from a careful evaluation of the Beliaev self-energies at finite $`T`$ is the same as (113), apart from a slightly smaller numerical coefficient ($`4/3`$ is replaced by $`3\pi /8)`$. This difference is due to our neglect of the fluctuations associated with the anomalous pair density $`\stackrel{~}{m}(𝐫,t)`$. Such terms are in the exact equation for $`\mathrm{\Phi }(𝐫,t)`$ given in (53) but were neglected in deriving (63).
To conclude this discussion of different kinds of damping in the collisionless limit, we can also include the effect of the collision terms in (87) as a perturbative correction to the solution determined by $`\widehat{}f=0.`$ For a Bose gas above $`T_{BEC}`$, one might use the simple “relaxation time” approximation
$$C_{22}[f](ff_0)/\tau _{22}=\delta f/\tau _{22}.$$
(114)
The inclusion of such a term in (98) simply involves the change
$$\frac{\delta f}{t}\frac{\delta f}{t}+\frac{\delta f}{\tau _{22}},$$
(115)
and hence one finds (103) again, with the replacement $`\omega \omega +i/\tau _{22}`$ in the denominator of (104). The collision time $`\tau _{22}`$ is the analogue of the classical expression given in (95). One sees that including $`C_{22}[f]`$ in the simple relaxation time approximation (114) only slightly alters the Landau damping given by (106). This is quite different from the contribution arising from $`C_{12}[f]`$ in so far as this enters directly into the generalized GP equation (see Section 3).
We now leave the discussion of the collisionless domain and turn to the collision-dominated hydrodynamic region. In this case, the collisions between the non-condensate atom are assumed to be sufficiently rapid that the form of $`f(𝐩,𝐫,t)`$ is largely determined by the requirement that
$$C_{22}[f,\mathrm{\Phi }]=0.$$
(116)
That is, the collisions are so strong they force the system to be in local equilibrium. The unique solution (denoted by $`\stackrel{~}{f})`$ of the equation (116) is well-known to be given by
$$\stackrel{~}{f}(𝐩,𝐫,t)=\frac{1}{e^{\beta [\frac{(𝐩m𝐯_n)^2}{2m}+U(𝐫,t)\stackrel{~}{\mu }(r,t)]}1},$$
(117)
where $`𝐯_n`$ is the average local velocity and $`\stackrel{~}{\mu }`$ is the local chemical potential of the thermal atoms. This local equilibrium Bose distribution involves the local variables $`\beta ,𝐯_n,\stackrel{~}{\mu }`$ and $`U`$, all of which depend on $`(𝐫,t)`$. Why must $`\stackrel{~}{f}`$ have the form in (117)? To satisfy $`C_{22}[f_1]=0,`$ we must have \[see (88)\]
$$(1+f_1)(1+f_2)f_3f_4f_1f_2(1+f_3)(1+f_4)=0,$$
(118)
and this requires that $`f`$ be given by the Bose distribution (117). Here we have used the fact that a Bose distribution satisfies
$$f(x)\frac{1}{e^x1}=[f(x)+1]$$
(119)
and energy and momentum conservation
$$\begin{array}{cc}& 𝐩_1+𝐩_2=𝐩_3+𝐩_4\\ & \\ & \stackrel{~}{\epsilon }_{p_1}+\stackrel{~}{\epsilon }_{p_2}=\stackrel{~}{\epsilon }_{p_3}+\stackrel{~}{\epsilon }_{p_4}\end{array}\}.$$
However, while the fact that $`\stackrel{~}{f}`$ is given by the local equilibrium Bose distribution in $`(\text{117})`$ ensures that $`C_{22}[\stackrel{~}{f}]=0,`$ one finds that $`C_{12}[\stackrel{~}{f}]0.`$ More precisely, the term in (89) reduces to
$`\left[(1+\stackrel{~}{f}_1)\stackrel{~}{f}_2\stackrel{~}{f}_3\stackrel{~}{f}_1(1+\stackrel{~}{f}_2)(1+\stackrel{~}{f}_3)\right]`$
$`\left[e^{\beta [\stackrel{~}{\mu }\mu _c\frac{1}{2}m(𝐯_n𝐯_c)^2]}1\right](1+\stackrel{~}{f}_1)\stackrel{~}{f}_2\stackrel{~}{f}_3,`$ (120)
using energy and momentum conservation. The expression in the square bracket in (120) only vanishes if the condensate and non-condensate are in diffusive equilibrium, which requires that the chemical potentials must be equal
$$\stackrel{~}{\mu }=\mu _c+\frac{1}{2}m(𝐯_n𝐯_c)^2.$$
(121)
When we perturb the system, this may not be true, ie, the two components may be out of diffusive equilibrium.
We can now derive hydrodynamic equations for the non-condensate by taking moments of the Boltzmann equation, the standard procedure used in classical gases. The first moment gives a continuity equation with a source term:
$`{\displaystyle 𝑑𝐩\left\{\stackrel{~}{f}=C_{12}[\stackrel{~}{f}]\right\}}{\displaystyle \frac{\stackrel{~}{n}}{t}}=\mathbf{}(\stackrel{~}{n}𝐯_n)+\mathrm{\Gamma }_{12}[\stackrel{~}{f}],`$ (122)
where we have used the fact that $`C_{22}[\stackrel{~}{f}]=0`$ and
$`\stackrel{~}{n}`$ $``$ $`{\displaystyle \frac{d𝐩}{(2\pi )^3}\stackrel{~}{f}(𝐩,𝐫,t)}`$
$`\stackrel{~}{n}𝐯_n`$ $``$ $`{\displaystyle \frac{d𝐩}{(2\pi )^3}\frac{𝐩}{m}\stackrel{~}{f}(𝐩,𝐫,t)}`$
$`\mathrm{\Gamma }_{12}[\stackrel{~}{f},\mathrm{\Phi }]`$ $``$ $`{\displaystyle \frac{d𝐩}{(2\pi )^3}C_{12}[\stackrel{~}{f},\mathrm{\Phi }]}.`$ (123)
More explicitly, we find \[compare with (71)\]
$`\mathrm{\Gamma }_{12}[\stackrel{~}{f}]`$ $`=`$ $`{\displaystyle \frac{2g^2n_c}{(2\pi )^5}}[e^{\beta [\stackrel{~}{\mu }\mu _c\frac{1}{2}m(𝐯_n𝐯_c)^2]}1]`$ (124)
$`\times `$ $`{\displaystyle 𝑑𝐩_1𝑑𝐩_2𝑑𝐩_3\delta (m𝐯_c+𝐩_1𝐩_2𝐩_3)}`$
$`\times `$ $`\delta (\epsilon _c+\stackrel{~}{\epsilon }_1\stackrel{~}{\epsilon }_2\stackrel{~}{\epsilon }_3)(1+\stackrel{~}{f}_1)\stackrel{~}{f}_2\stackrel{~}{f}_3`$
$``$ $`{\displaystyle \frac{n_c}{\tau _{12}}}\left[e^{\beta [\stackrel{~}{\mu }\mu _c\frac{1}{2}m(𝐯_n𝐯_c)^2]}1\right].`$
We note that $`\tau _{12}`$ is a collision time which describes the $`C_{12}`$ collisions between the condensate and non-condensate atoms. Adding (122) to the continuity equation in (60)
$$\frac{n_c}{t}=\mathbf{}(n_c𝐯_c)\mathrm{\Gamma }_{12}[\stackrel{~}{f},\mathrm{\Phi }],$$
(125)
we see that the source term $`\mathrm{\Gamma }_{12}`$ cancels out to give
$$\frac{(n_c+\stackrel{~}{n})}{t}=\mathbf{}(n_c𝐯_c+\stackrel{~}{n}𝐯_n).$$
(126)
Thus our theory gives the correct continuity equation for the total local density $`n=n_c+\stackrel{~}{n}.`$
Similarly, one finds
$`{\displaystyle 𝑑\mathrm{𝐩𝐩}\left\{\widehat{}\stackrel{~}{f}=C_{12}[\stackrel{~}{f}]\right\}}m\stackrel{~}{n}\left({\displaystyle \frac{𝐯_n}{t}}+{\displaystyle \frac{1}{2}}\mathbf{}𝐯_n^2\right)`$
$`=\mathbf{}\stackrel{~}{P}(𝐫,t)\stackrel{~}{n}U(𝐫,t)m(𝐯_n𝐯_c)\mathrm{\Gamma }_{12}[\stackrel{~}{f}],`$ (127)
where the kinetic pressure is given by
$$\stackrel{~}{P}(𝐫,t)=\frac{m}{3}\frac{d𝐩}{(2\pi )^3}(𝐩m𝐯_n)^2\stackrel{~}{f}(𝐩,𝐫,t).$$
(128)
Finally, the second moment gives
$`{\displaystyle 𝑑𝐩p^2\left\{\stackrel{~}{f}=C_{12}[\stackrel{~}{f}]\right\}}{\displaystyle \frac{\stackrel{~}{P}}{t}}+\mathbf{}(\stackrel{~}{P}𝐯_n)`$
$`={\displaystyle \frac{2}{3}}\stackrel{~}{P}\mathbf{}𝐯_n+{\displaystyle \frac{2}{3}}\left[\mu _c+{\displaystyle \frac{1}{2}}m(𝐯_n𝐯_c)^2U\right]\mathrm{\Gamma }_{12}[\stackrel{~}{f}].`$ (129)
The detailed derivation of these equations is not important here. It involves straightforward manipulations using the explicit form of $`\stackrel{~}{f}`$ in (117).
The hydrodynamic equations (122), (127) and (129) describe the non-condensate in terms of three new “coarse-grained” local variables:
$$\stackrel{~}{n}(𝐫,t),𝐯_n(𝐫,t)\text{and}\stackrel{~}{P}(𝐫,t).$$
These are coupled to the two additional local variables which describe the condensate:
$$n_c(𝐫,t),𝐯_c(𝐫,t).$$
We note that the two condensate equations of motion given by (60) are always “hydrodynamic” in form. In contrast, it is only in the collision-dominated region that the non-condensate dynamics can also be described in terms of a few collective variables. Both components exhibit coupled, coherent collective motions at the same frequency. This is the essence of two-fluid superfluid behaviour, familiar in liquid <sup>4</sup>He studies but still an unexplored frontier in trapped Bose gases.
We will now discuss the linearized version of our condensate and non-condensate equations for local equilibrium as given by (60), (121), (127) and (129). We work to first order in the fluctuations around static equilibrium,
$`\stackrel{~}{n}`$ $`=`$ $`\stackrel{~}{n}_0+\delta \stackrel{~}{n},𝐯_n=\delta 𝐯_n,\stackrel{~}{P}=\stackrel{~}{P}_0+\delta \stackrel{~}{P}`$
$`n_c`$ $`=`$ $`n_{c0}+\delta n_c,𝐯_c=\delta 𝐯_c.`$ (130)
What is new about the two-fluid hydrodynamic equations derived above is the role of the source term $`\mathrm{\Gamma }_{12}[\stackrel{~}{f}].`$ In a linearized theory expanded around the static equilibrium Bose distribution $`f_0`$ (where $`\mathrm{\Gamma }_{12}[f_0,\mathrm{\Phi }_0]`$ vanishes), one finds
$$\mathrm{\Gamma }_{12}[\stackrel{~}{f},\mathrm{\Phi }]=\delta \mathrm{\Gamma }_{12}[\stackrel{~}{f},\mathrm{\Phi }]=\frac{\beta _0n_{c0}}{\tau _{12}^0}\delta \mu _{diff},$$
(131)
where
$$\mu _{diff}(𝐫,t)\stackrel{~}{\mu }(𝐫,t)\mu _c(𝐫,t).$$
(132)
Here $`\tau _{12}^0(𝐫)`$ is the $`C_{12}`$ collision time defined in (124) with both components being in static equilibrium (see also (77)). The linearized coupled hydrodynamic equations for the two components are given by the ZGN equations :
$`{\displaystyle \frac{\delta n_c}{t}}`$ $`=`$ $`\mathbf{}(n_{c0}\delta 𝐯_c)\delta \mathrm{\Gamma }_{12}`$
$`m{\displaystyle \frac{\delta 𝐯_c}{t}}`$ $`=`$ $`\mathbf{}\delta \mu _c`$ (133)
and
$`{\displaystyle \frac{\delta \stackrel{~}{n}}{t}}`$ $`=`$ $`\mathbf{}(\stackrel{~}{n}_0\delta 𝐯_n)+\delta \mathrm{\Gamma }_{12}`$
$`m\stackrel{~}{n}_0{\displaystyle \frac{\delta 𝐯_n}{t}}`$ $`=`$ $`\mathbf{}\delta \stackrel{~}{P}\delta \stackrel{~}{n}\mathbf{}U_0(𝐫)2g\stackrel{~}{n}_0[\mathbf{}\delta n_c+\mathbf{}\delta \stackrel{~}{n}]`$
$`{\displaystyle \frac{\delta \stackrel{~}{P}}{t}}`$ $`=`$ $`{\displaystyle \frac{5}{3}}\mathbf{}(\stackrel{~}{P}_0\delta 𝐯_n)+{\displaystyle \frac{2}{3}}\delta 𝐯_n\mathbf{}\stackrel{~}{P}_0{\displaystyle \frac{2}{3}}gn_{c0}\delta \mathrm{\Gamma }_{12},`$ (134)
with $`\delta \mathrm{\Gamma }_{12}`$ given by (131) and (within the TF approximation)
$$\delta \mu _c=g\delta n_c+2g\delta \stackrel{~}{n}.$$
(135)
We note that to lowest order, $`\delta \mathrm{\Gamma }_{12}`$ does not appear in the second equation in (134). In addition, in the last term in the third equation in (134), we have used $`\mu _{c0}U_0=gn_{c0}`$. The condensate couples into the non-condensate equations in (134) directly via the mean field terms $`g\delta n_c`$ and, indirectly, through $`\delta \mathrm{\Gamma }_{12}`$ (see Section 5).
Despite appearances, this coupled set of hydrodynamic equations form a closed set. We have 9 scalar equations in (133) and (134), while they appear to involve 10 fluctuating variables
$$\delta n_c,\delta 𝐯_c;\delta \stackrel{~}{n},\delta 𝐯_n;\delta \stackrel{~}{P},\delta \mu _{diff}.$$
(136)
However one can show that $`\delta \mu _{diff}`$ can be written as a linear combination of the variables $`\delta \stackrel{~}{P},\delta \stackrel{~}{n}`$ and $`\delta n_c`$ and thus we are in fact left with only 9 variables. In Section 5, we prove that (133) and (134) are equivalent to the two-fluid hydrodynamic equations of Landau when applied to a trapped Bose gas, but only when the condensate and non-condensate are in diffusive local equilibrium, ie, $`\mu _c(𝐫,t)=\stackrel{~}{\mu }(𝐫,t).`$ It will turn out that the rate at which $`\delta \mu _{diff}(𝐫,t)`$ relaxes to zero is controlled by a relaxation time $`\tau _\mu `$ related to the $`C_{12}`$ collisions. The Landau two-fluid hydrodynamics is only valid for low frequency oscillations which satisfy $`\omega \tau _\mu 1`$, as we discuss in Section 5.
## 5 Two-fluid hydrodynamics in Bose liquids and gases
To put our new two-fluid equations derived in Section 4 into some sort of context, we first briefly review the more general theory developed by Landau in 1941 to explain superfluidity in liquid <sup>4</sup>He.
The original discovery of superfluidity in liquid <sup>4</sup>He is associated with the famous 1938 papers of Kapitza in Moscow and Allen and Misener at Cambridge. (I cannot resist remarking that Allen and Misener had been graduate students at the University of Toronto where they carried out (with Burton) the pioneering studies on vanishing viscosity in the period 1935-1937). These and subsequent experiments in the next few years showed that superfluid <sup>4</sup>He could exhibit very bizarre behaviour compared to ordinary liquids. This led to the development of a two-fluid theory of the hydrodynamic behaviour of liquid <sup>4</sup>He by Landau (1941). An earlier but less complete version of Landau’s hydrodynamic equations was developed by Tisza in the period 1938-40. For further discussion of this early history, I refer to a recent article of mine.
In this early work, superfluidity (the term was coined in 1938 by Kapitza) was entirely associated with the relative motion of the normal fluid and the superfluid components under a variety of conditions. The main point was that while the normal fluid exhibited finite viscosity and thermal conductivity typical of ordinary fluid, the superfluid component (which exhibited irrotational flow) did not. In more recent times, the aspect of superfluidity which has been emphasized (see Ch.4 of Ref. and Leggett’s lectures in this book, for example) are those most directly tied to the fact that the superfluid velocity is associated with the gradient of the phase of the macroscopic wavefunction $`\mathrm{\Phi }(𝐫,t).`$ Given that this point of view is indeed more fundamental, it is still crucial to understand why superfluidity persists even in the presence of a dissipative normal fluid. This question can be addressed using two-fluid hydrodynamic equations.
In essence, Landau developed his generic two-fluid hydrodynamics by generalizing the standard theory of classical hydrodynamics to include the equations of motion for a new “superfluid” degree of freedom. We recall that classical fluid dynamics was developed well before one knew about the existence of atoms. Since the work of Maxwell and Boltzmann in the 1880’s, we know the “coarse-grained” hydrodynamic description of a fluid in terms of a few quantities like $`n(𝐫,t)`$ and $`𝐯(𝐫,t)`$ is only valid when the collisions between atoms are strong enough to produce local equilibrium. It only describes low frequency phenomena, where the condition in (94) is satisfied.
In his 1941 paper, Landau did not connect the superfluid component with the motion of a “Bose condensate.” Indeed, he rejected the efforts by Tisza and F. London to use a Bose-condensed gas to get some insight into superfluid <sup>4</sup>He. However, since the period 1957-1965, Landau’s superfluid degree of freedom has been understood microscopically in terms of the complex order parameter $`\mathrm{\Phi }(𝐫,t).`$ As noted earlier, the superfluid velocity field $`𝐯_s(𝐫,t)`$ is related to the gradient of the phase of $`\mathrm{\Phi }(𝐫,t)`$, as given by (19). In the same period, it was also realized that there are two distinct kinds of quantum fluids:
1. Bose fluids (associated with a Bose condensate wavefunction $`\mathrm{\Phi }`$)
2. Normal Fermi fluids (associated with the key role of a Fermi surface).
These two kinds of quantum fluids were magnificently described in two books by Nozières and Pines written around 1965, although the one on superfluid Bose liquids was only published in 1990 (it was in wide circulation as a preprint before then). I think the clearest account of the connection between $`\mathrm{\Phi }(𝐫,t)`$ and superfluidity in Bose fluids is still the discussion given in Chapters 4 and 5 of Vol II by Nozières and Pines. I highly recommend it to everyone in the BEC field.
Landau’s pioneering work on superfluid <sup>4</sup>He has two separate aspects which are logically distinct but sometimes confused with each other:
1. The two-fluid equations describing hydrodynamic behaviour. These equations are generic and apply (under certain conditions) to trapped Bose gases as well as to superfluid <sup>4</sup>He. In the period 1947-1950, Landau and Khalatnikov extended these equations to include hydrodynamic damping of the normal fluid (described by various kinds of viscosities, thermal conductivity, etc). This work is all described in the classic 1965 monograph by Khalatnikov.
2. A “microscopic” theory of the elementary excitations describing the normal fluid of liquid <sup>4</sup>He - the famous phonon - roton spectrum. Within his picture of a weakly interacting gas of quasiparticles, Landau and coworkers could calculate the thermodynamic and transport properties of superfluid <sup>4</sup>He. These quantities enter into the expressions given by the two-fluid equations for the first and second sound modes (velocity and damping). Of course, the roton part of the spectrum is not valid for a dilute Bose gas.
As we have noted, Landau’s original formulation of his two-fluid hydrodynamic equations was phenomenological in that the superfluid component was not given an explicit microscopic basis. The first derivation of the Landau hydrodynamic equations starting simply from the existence of the macroscopic order parameter $`\mathrm{\Phi }(𝐫,t)`$ was given by Bogoliubov in 1963. This derivation built on Bogoliubov’s earlier work on deriving hydrodynamic equations for classical liquids without going through the intermediate stage of using Boltzmann-like kinetic equations. While Bogoliubov’s derivation is often viewed as being quite general, buried in his complex analysis is the assumption that the normal fluid and the superfluid are in local equilibrium with each other. Even today, probably the definitive account which formulates the various levels of theory for Bose superfluids is the classic paper by Hohenberg and Martin published in 1965. It clearly shows (to the patient reader!) the central unifying role of $`\mathrm{\Phi }(𝐫,t)`$, summarizes the collisionless and hydrodynamic domains, and finally gives criteria for developing and judging various approximation schemes using Green’s function techniques.
With the preceding discussion as a preamble, I will now summarize the linearized two-fluid hydrodynamic equations of Landau, following the approach given in Ch. 7 of Ref. . As noted, these equations are valid for both superfluid liquids as well as gases. The differences come in only at the last stage when we evaluate the thermodynamic coefficients appearing in these equations, using the appropriate quasiparticle excitation spectrum. The first two (linearized) Landau equations are familiar from ordinary fluid dynamics
$$\frac{\delta n}{t}+\mathbf{}\delta 𝐣=0$$
(137)
$$m\frac{\delta 𝐣}{t}=\mathbf{}\delta P\delta n\mathbf{}V_{ex},$$
(138)
where we have included the effect of an external potential. Landau’s work incorporated two components and hence the total mass density and mass current fluctuations are given by
$$\delta \rho m\delta n=\delta \rho _s+\delta \rho _n$$
(139)
$$m\delta 𝐣\rho _{s0}\delta 𝐯_s+\rho _{n0}\delta 𝐯_n.$$
(140)
The new superfluid component was argued to only exhibit pure potential (irrotational) flow and carry no entropy. Thus Landau’s final two hydrodynamic equations were new,
$$m\frac{\delta 𝐯_s}{t}=\mathbf{}\delta \mu $$
(141)
$$\frac{\delta s}{t}+\mathbf{}(s_0\delta 𝐯_n)=0..$$
(142)
Here the local equilibrium entropy density is $`s(𝐫,t)=s_0+\delta s`$, the local equilibrium pressure is $`P(𝐫,t)=P_0+\delta P`$ and the local chemical potential is $`\mu (𝐫,t)=\mu _0+\delta \mu (𝐫,t).`$ We note that the Landau hydrodynamic equations only involves 8 equations, in contrast to the 9 equations we derived at the end of Section 4. In particular, we see that Landau does not have separate continuity equations for the superfluid $`\rho _s(𝐫,t)`$ and normal fluid $`\rho _n(𝐫,t)`$ densities. Finally, it is assumed that these two components are always in local equilibrium with each other and hence the fluctuations are related by the thermodynamic identity
$$n_0\delta \mu =s_0\delta T+\delta P,$$
(143)
where $`n_0`$ is the equilibrium total density.
For simplicity, we now specialize our analysis to a uniform superfluid, where $`\rho _{s0},\rho _{n0},s_0`$ and other thermodynamic quantities are all position-independent. Multiplying (141) by $`n_{s0}`$ gives \[using (143)\]
$`\rho _{s0}{\displaystyle \frac{\delta 𝐯_s}{t}}`$ $`=`$ $`n_{s0}\mathbf{}\delta \mu `$ (144)
$`=`$ $`{\displaystyle \frac{n_{s0}}{n_0}}\mathbf{}(s_0\delta T+\delta P)`$
$`=`$ $`{\displaystyle \frac{n_{s0}}{n_0}}\mathbf{}\delta P+{\displaystyle \frac{n_{s0}}{n_0}}s_0\mathbf{}\delta T.`$
Using this in (137) gives
$`\mathbf{}\delta P`$ $`=`$ $`\rho _{s0}{\displaystyle \frac{\delta 𝐯_s}{t}}+\rho _{n0}{\displaystyle \frac{\delta 𝐯_n}{t}}`$
$`=`$ $`{\displaystyle \frac{n_{s0}}{n_0}}\mathbf{}\delta P+{\displaystyle \frac{n_{s0}}{n_0}}s_0\mathbf{}\delta T+\rho _{n0}{\displaystyle \frac{𝐯_n}{t}}`$
or
$$\rho _{n0}\frac{\delta 𝐯_n}{t}=\frac{n_{n0}}{n_0}\mathbf{}\delta P\frac{n_{s0}}{n_0}s_0\mathbf{}\delta T.$$
(145)
Combining (137) and (138) gives
$$\frac{^2\delta \rho }{t^2}=m\mathbf{}\frac{\delta 𝐣}{t}=^2\delta P.$$
(146)
Finally, (142) gives \[using (145)\]
$`{\displaystyle \frac{^2\delta s}{t^2}}`$ $`=`$ $`s_0\mathbf{}{\displaystyle \frac{\delta 𝐯_n}{t}}`$ (147)
$`=`$ $`{\displaystyle \frac{s_0}{\rho _0}}^2\delta P+{\displaystyle \frac{s_0^2}{\rho _0}}\left({\displaystyle \frac{\rho _{s0}}{\rho _{n0}}}\right)^2\delta T`$
$`=`$ $`{\displaystyle \frac{s_0}{\rho _0}}{\displaystyle \frac{^2\delta \rho }{t^2}}+{\displaystyle \frac{s_0^2}{\rho _0}}\left({\displaystyle \frac{\rho _{s0}}{\rho _{n0}}}\right)^2\delta T.`$
The last equation gives the local entropy fluctuations in terms of the local mass density $`\delta \rho `$ and temperature $`\delta T`$ fluctuations. Eq.(147) can be re-written in terms of the local entropy per unit mass $`\overline{s}(𝐫,t)s(𝐫,t)/\rho (𝐫,t).`$ Using
$$\delta \overline{s}=\frac{s_0}{\rho _0^2}\delta \rho +\frac{1}{\rho _0}\delta s,$$
(148)
(147) takes on the simpler form
$$\frac{^2\delta \overline{s}}{t^2}=\overline{s}_0^2\left(\frac{\rho _{n0}}{\rho _{s0}}\right)^2\delta T..$$
(149)
Expanding $`\delta P`$ and $`\delta T`$ in terms of $`\delta \rho `$ and $`\delta \overline{s}`$ fluctuations
$`P`$ $`=`$ $`{\displaystyle \frac{P}{\rho }}|_{\overline{s}}\delta \rho +{\displaystyle \frac{P}{\overline{s}}}|_\rho \delta \overline{s}`$
$`\delta T`$ $`=`$ $`{\displaystyle \frac{T}{\rho }}|_{\overline{s}}\delta \rho +{\displaystyle \frac{T}{\overline{s}}}|_\rho \delta \overline{s},`$ (150)
we see that (149) and (146) reduce to two coupled scalar equations for $`\delta \rho `$ and $`\delta \overline{s}`$. Inserting the normal mode solutions
$$\delta \rho ,\delta \overline{s}e^{i(𝐤𝐫\omega t)},$$
(151)
one finds these coupled algebraic equations have two phonon solutions $`\omega ^2=u^2k^2`$, where $`u^2`$ is the solution of the quadratic equation:
$`u^4`$ $``$ $`u^2\left[{\displaystyle \frac{P}{\rho }}|_T+{\displaystyle \frac{T}{\overline{c}_v}}\left({\displaystyle \frac{1}{\rho _0}}{\displaystyle \frac{P}{T}}|_\rho \right)^2+{\displaystyle \frac{\rho _{s0}}{\rho _{n0}}}{\displaystyle \frac{T\overline{s}_0^2}{\overline{c}_v}}\right]`$ (152)
$`+`$ $`{\displaystyle \frac{\rho _{s0}}{\rho _{n0}}}\left({\displaystyle \frac{T\overline{s}_0^2}{\overline{c}_v}}\right){\displaystyle \frac{P}{\rho }}|_T=0.`$
Here $`\overline{c}_v=\frac{\overline{s}}{T}|_\rho `$ is the equilibrium specific heat per unit mass.
The coefficients in (152) are daunting, but only involve equilibrium thermodynamic quantities which can be calculated for any given superfluid. We emphasize that (152) is valid for both a uniform Bose-condensed gas and superfluid <sup>4</sup>He, assuming that the superfluid and normal fluid are in local equilibrium with each other. However, as we shall show, the detailed characteristics and behaviour of the two phonon modes (first and second second) are quite different in a Bose gas and in superfluid <sup>4</sup>He.
A key feature about superfluid <sup>4</sup>He is that (typical of any liquid) $`P/T|_\rho 0.`$ This can be shown to be equivalent to $`C_pC_v,`$ where $`C_{p,v}`$ are the specific heats at constant pressure or constant volume. In this case, (152) reduces to
$$u^4u^2(A+B)+AB=0,$$
(153)
with two solutions $`u^2=A,B:`$
$`u_1^2`$ $`=`$ $`{\displaystyle \frac{P}{\rho }}\text{: first sound}`$ (154)
$`u_2^2`$ $`=`$ $`{\displaystyle \frac{\rho _{s0}}{\rho _{n0}}}\left({\displaystyle \frac{T\overline{s}_0^2}{\overline{c}_v}}\right)\text{: second sound}.`$ (155)
Working out the associated motions, one finds the first sound mode $`(\omega =u_1k)`$ involves the in-phase motion of the superfluid and normal fluid components. Moreover, one can show that first sound is essentially a pressure wave. In contrast, the second sound mode $`(\omega =u_2k)`$ involves the out-of-phase motion of the two components, with $`\delta 𝐣0`$ or
$$\rho _{n0}\delta 𝐯_n=\rho _{s0}\delta 𝐯_s..$$
(156)
This corresponds to an almost pure temperature wave. The successful detection in 1946 of a second sound mode (or temperature wave) was of tremendous significance in low temperature physics. The good agreement of the measured second sound velocity $`u_2`$ with the Landau expression in (155), calculated using the postulated phonon-roton quasiparticle spectrum, vindicated both aspects of the Landau theory of superfluid <sup>4</sup>He.
In contrast, in a (Bose-condensed) gas, the pressure and temperature fluctuations are strongly coupled and hence $`P/T|_\rho `$ in (152) plays a significant role. We recall that in a classical gas, one has $`C_p/C_v=\frac{5}{3}.`$ Evaluating all the thermodynamic derivatives using our finite $`T`$ microscopic model (the single-particle HF spectrum in (59) replaces the roton spectrum of superfluid <sup>4</sup>He) and working to lowest order in $`gn_c/k_BT,`$ one finds (after very lengthy calculations)
$$u_1^2=\frac{5}{3}\frac{k_BT}{m}\frac{g_{5/2}(z_0=1)}{g_{3/2}(z_0=1)}+O(gn)$$
(157)
$$u_2^2=\frac{gn_{c0}}{m}+\mathrm{}$$
(158)
It turns out that the first sound mode in a Bose gas is largely an oscillation of the thermal cloud (the normal fluid). The second sound mode is largely an oscillation of the condensate (the superfluid). Both involve density fluctuations and thus will have significant weight in the dynamic structure factor related to the density response function. The quite different features of first and second sound in Bose gases vs superfluid <sup>4</sup>He should be remembered when reading standard texts about superfluid <sup>4</sup>He. Results equivalent to (157) and (158) are given in Refs. .
Comparing (158) with (28), it is clear that in a uniform Bose-condensed gas, the hydrodynamic second sound mode at finite $`T`$ smoothly extrapolates to the Bogoliubov phonon mode in the collisionless region. The fact that the first sound velocity $`u_1`$ does not depend on the interaction strength $`g`$ to lowest order is typical of ordinary sound waves in any gas. However, one must remember that interactions (collisons) play a crucial indirect role in enforcing dynamic local equilibrium.
We now turn to a discussion of the generalized two-fluid hydrodynamic equations we derived at the end of Section 4. Again, for simplicity, we limit our analysis to a uniform Bose-condensed gas. If we can ignore vorticity in both fluids, we can introduce two velocity potentials
$`\delta 𝐯_c`$ $``$ $`\mathbf{}\varphi _c(𝐫,t)`$
$`\delta 𝐯_n`$ $``$ $`\mathbf{}\varphi _n(𝐫,t),`$ (159)
and reduce the equations in (133) and (134) to three equations for $`\varphi _c,\varphi _n`$ and $`\delta \mu _{diff}`$ \[the latter is defined in (132)\]:
$`m{\displaystyle \frac{^2\varphi _c}{t^2}}`$ $`=`$ $`gn_{c0}^2\varphi _c+2g\stackrel{~}{n}_0^2\varphi _n+{\displaystyle \frac{\sigma _H}{\tau _\mu }}\delta \mu _{diff}`$ (160)
$`m{\displaystyle \frac{^2\varphi _n}{t^2}}`$ $`=`$ $`\left({\displaystyle \frac{5}{3}}{\displaystyle \frac{\stackrel{~}{P}_0}{\stackrel{~}{n}_0}}+2g\stackrel{~}{n}_0\right)^2\varphi _n+2gn_{c0}^2\varphi _c{\displaystyle \frac{2}{3}}{\displaystyle \frac{\sigma _H}{\tau _\mu }}\delta \mu _{diff}`$ (161)
$$\frac{\delta \mu _{diff}}{t}=\frac{2}{3}gn_{c0}^2\varphi _ngn_{c0}^2\varphi _c\frac{\delta \mu _{diff}}{\tau _\mu }.$$
(162)
Here $`\tau _\mu `$ is a new relaxation time governing how $`\delta \mu _{diff}`$ relaxes to $`0`$ (ie, $`\stackrel{~}{\mu }\mu _c)`$. It is found to be related to the $`C_{12}`$ collision time $`\tau ^{}`$ defined in (80) by
$$\frac{1}{\tau _\mu }=\frac{1}{\sigma _H\tau ^{}},$$
(163)
where the dimensionless hydrodynamic renormalization factor $`\sigma _H`$ is given by
$$\sigma _H\frac{\frac{5}{2}\stackrel{~}{\gamma }_0\stackrel{~}{P}_0\frac{3}{2}g\stackrel{~}{n}_0^2}{\frac{5}{2}\stackrel{~}{P}_0(1\stackrel{~}{\gamma }_0)+2g\stackrel{~}{n}_0n_{c0}+\frac{2}{3}g\stackrel{~}{\gamma }_0n_{c0}^2+\frac{3}{2}g\stackrel{~}{n}_0^2}.$$
(164)
For completeness, we recall that
$`\stackrel{~}{P}_0(z_0)`$ $`=`$ $`{\displaystyle \frac{k_BT}{\mathrm{\Lambda }_0^3}}g_{5/2}(z_0=e^{\beta gn_{c0}})`$
$`\stackrel{~}{n}_0(z_0)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Lambda }_0^3}}g_{3/2}(z_0)`$
$`\stackrel{~}{\gamma }_0`$ $`=`$ $`{\displaystyle \frac{g}{k_BT\mathrm{\Lambda }_0^3}}g_{1/2}(z_0),`$ (165)
where the Bose-Einstein functions are $`g_n(z)_{l=1}^{\mathrm{}}z^l/l^n`$ and $`\mathrm{\Lambda }_0`$ is the equilibrium thermal de Broglie wavelength. Calculation shows that $`\sigma _H`$ becomes large as $`TT_{BEC}`$ and thus $`\tau _\mu `$ can become very large near the superfluid transition. This is not in contradiction with the fact that the present discussion is for the collision-dominated region $`[\omega \tau _{22}1,`$ where $`\tau _{22}`$ is some relaxation time associated with the $`C_{22}`$ collision integral in (88)\].
It is clear from the fact that we have 3 coupled equations for the 3 variables $`\varphi _c,\varphi _n`$ and $`\delta \mu _{diff}`$, we will obtain a new mode, in addition to the usual first and second sound oscillations discussed earlier. It is convenient to first eliminate $`\delta \mu _{diff}`$ using the solution of (162),
$$i\omega \delta \mu _{diff}=\frac{2}{3}gn_{c0}(k^2)\varphi _ngn_{c0}(k^2)\varphi _c\frac{\delta \mu _{diff}}{\tau _\mu }$$
or
$$\delta \mu _{diff}=\frac{gn_{c0}\tau _\mu }{1i\omega \tau _\mu }(\varphi _c\frac{2}{3}\varphi _n)k^2.$$
(166)
Using this in (160) and (161) gives
$$m\omega ^2\varphi _c=gn_{c0}\left(1\frac{\sigma _H}{1i\omega \tau _\mu }\right)k^2\varphi _c+2g\stackrel{~}{n}_0\left(1+\frac{\sigma _H}{3(1i\omega \tau _\mu )}\frac{n_{c0}}{\stackrel{~}{n}_0}\right)k^2\varphi _n$$
(167)
$`m\omega ^2\varphi _n`$ $`=`$ $`\left({\displaystyle \frac{5}{3}}{\displaystyle \frac{\stackrel{~}{P}_0}{\stackrel{~}{n}_0}}+2g\stackrel{~}{n}_0\left[1{\displaystyle \frac{2\sigma _H}{9(1i\omega \tau _\mu )}}{\displaystyle \frac{n_{c0}^2}{\stackrel{~}{n}_0^2}}\right]\right)k^2\varphi _n`$ (168)
$`+`$ $`2gn_{c0}\left(1+{\displaystyle \frac{\sigma _H}{3(1i\omega \tau _\mu )}}{\displaystyle \frac{n_{c0}}{\stackrel{~}{n}_0}}\right)k^2\varphi _c.`$
These two coupled equations for $`\varphi _n`$ and $`\varphi _c`$ are easily solved and we obtain (as expected) first and second sound modes. Clearly the velocities $`u_1`$ and $`u_2`$ will now depend on the value of $`\omega \tau _\mu ,`$ although it turns out that this dependence is not very strong. This is shown by the results in Fig. 5.
The $`\omega \tau _\mu 0`$ limit is of special interest since the equations (167) and (168) can then be shown to be completely equivalent to the predictions of the Landau two-fluid equations discussed earlier in this section. To be precise, the first and second sound velocities $`u_{1,2}^2`$ obtained from (167) and (168) in the limit $`\omega \tau _\mu 0`$ agree with those given by (152). This equivalence makes physical sense since in the limit of $`\tau _\mu 0,`$ one sees that $`\delta \mu _{diff}0`$ very rapidly \[see (162) and (166)\]. Thus we have proven that the Landau two-fluid hydrodynamic equations are correct if the condensate and thermal cloud are in diffusive local equilibrium. This, as noted earlier, is an (sometimes implicit) assumption in all previous derivations of the Landau equations.
This Landau limit $`(\omega \tau _\mu 0)`$ is, in fact, very subtle in the context of our microscopic calculation. We see that in this limit, there are still correction terms in (167) and (168) which are proportional to the hydrodynamic renormalization factor $`\sigma _H`$ as defined in (164). These terms are crucial in ensuring that (167) and (168) reproduce the results of the Landau two-fluid equations. Another way of seeing this is that even though $`\delta \mu _{diff}0`$ when $`\tau _\mu 0,`$ we note that $`\delta \mathrm{\Gamma }_{12}`$ in (131) is still finite. Using (166), one finds
$$\delta \mathrm{\Gamma }_{12}=n_{c0}\sigma _H\left(\varphi _c\frac{2}{3}\varphi _n\right)k^2.$$
(169)
In conclusion, one might say that the hydrodynamic renormalization factor $`\sigma _H`$ contains a key part of the physics buried in the hydrodynamic two-fluid equations of Landau.
In the opposite limit $`\omega \tau _\mu 1`$ (which can arise near $`T_{BEC}`$) we see all the terms proportional to $`\sigma _H`$ in (167) and (168) are negligible. This domain is missed in the Landau two-fluid equations. It describes situations in which the condensate and thermal cloud are out of diffusive equilibrium with each other. It would be of great interest to look for this kind of phenomenon in trapped Bose gases.
One can work out the frequency of the new mode associated with the dynamics of $`\delta \mu _{diff}`$ and for a uniform gas, it is well approximated by
$$\omega _Ri/\tau _\mu .$$
(170)
Thus, in general, our two-fluid hydrodynamic equations predict the existence of a relaxational mode peaked at zero frequency. One can improve the theory to include deviations from local equilibrium \[ie, deviations of $`f`$ from $`\stackrel{~}{f}`$ in (116)\]. This involves a Chapman-Enskog kind of calculation familiar in the theory of classical gases and gives rise to hydrodynamic damping. The normal fluid (non-condensate) equations of motion have new terms corresponding to shear viscosity $`(\eta )`$ and thermal conductivity $`(\kappa )`$ transport coefficients , where we recall that $`\kappa `$ and $`\eta `$ are proportional to some $`\tau _{22}`$ collision time and hence go as $`1/g^2`$. The damping of first sound, second sound and the relaxational mode due to small deviations from local equilibrium among the thermal atoms $`(C_{22}[f,\mathrm{\Phi }])`$ has been recently worked out in detail by Nikuni, Zaremba and the author. Of particular interest is the effect on the relaxational mode, which is now described by \[compare with (170)\]
$$\omega _Ri\left[\frac{1}{\tau _\mu }+A\kappa k^2\right].$$
(171)
Effectively the relaxation mode is strongly coupled into thermal conduction processes (or flow of heat). The expression for the coefficient $`A`$ in (171) is somewhat complicated but above $`T_{\mathrm{BEC}}`$ $`(\frac{1}{\tau _\mu }0),`$ we find the mode reduces to the well-known thermal diffusion mode,
$$\omega _R=i\frac{\kappa k^2}{n_0C_p}=iD_Tk^2..$$
(172)
This result strongly suggests that our new relaxational mode below $`T_{\mathrm{BEC}}`$ is the renormalized version of the usual thermal diffusion mode above $`T_{\mathrm{BEC}}`$.
We recall that in the standard two-fluid hydrodynamic theory of Landau, the thermal diffusion mode below $`T_{\mathrm{BEC}}`$ disappears below $`T_{\mathrm{BEC}}`$, with, its spectral weight going into the emerging second sound doublet. In Fig. 6, we schematically illustrate the different hydrodynamic mode spectra predicted for above and below $`T_{\mathrm{BEC}}`$.
To illustrate the essential physics, the detailed discussion in this section has been limited to a uniform Bose-condensed gas. The analysis can, of course, be extended to trapped Bose gas. In particular, the out-of-phase dipole mode in the hydrodynamic limit is of special interest. This corresponds to centre-of-mass oscillation of the equilibrium condensate and non-condensate density profiles,
$`n_c(𝐫,t)`$ $`=`$ $`n_{c0}(𝐫𝜼_c(t))`$
$`\stackrel{~}{n}(𝐫,t)`$ $`=`$ $`\stackrel{~}{n}_0(𝐫𝜼_n(t)),`$ (173)
with
$$𝐯_n=\frac{d𝜼_n}{dt};𝐯_c=\frac{d𝜼_c}{dt}.$$
(174)
One finds that the hydrodynamic equations give (as expected) an undamped in-phase normal mode, with $`𝐯_n=𝐯_c`$ and $`\omega =\omega _0`$, where $`\omega _0`$ is the trap frequency. This is a generalized version of the $`T=0`$ Kohn mode discussed at the end of Section 2. In addition, however, there is an out-of-phase mode satisfying
$$N_c𝐯_c=\stackrel{~}{N}𝐯_n,$$
(175)
with a frequency different from the trap frequency $`\omega _0`$. The frequency of this mode is shown in Fig. 7. Calculation shows that this out-of-phase mode is only damped by its coupling to the relaxational mode given by (170). There is no hydrodynamic-type damping from the finite transport coefficients of the normal gas. Calculations are in reasonable agreement with the observed damping of such a out-of-phase mode. A careful study of this out-of-phase dipole mode as a function of the temperature would be a nice way of probing the unusual hydrodynamics of a Bose-condensed gas.
In these lectures, I have not made any attempt to analyze the available experimental investigations of collective modes at finite temperatures. These pioneering experimental studies are very promising but we need more systematic investigations, especially as a function of temperature and density. A key reason why such studies are perhaps a unique way of probing the many-body dynamics of superfluid gases is that one can measure collective mode frequencies very accurately (at $`T=0`$, with errors of only a few percent).
In this section, we have put emphasis on deriving two-fluid hydrodynamic equations by starting from an approximate but still microscopic model. Such a derivation has been carried out recently by Zaremba, Nikuni and the author for a trapped Bose gas. The pioneering work on deriving two-fluid hydrodynamic equations for a uniform dilute Bose gas was by Kirkpatrick and Dorfman. In this regard, it is useful to point out that the linearized Landau two-fluid hydrodynamic equations \[as given by (137) - (142)\] are expected to be exact as long as the two components are in local equilibrium. Thus, for a uniform system, the exact first and second sound velocities are given by the solutions of (152). The only question is how to calculate the various thermodynamic functions which are involved in this equation. The analogous Landau equations for a trapped gas allows us to go past the simplified microscopic models we have used in Section 4 to derive equations of this kind. In particular, one should be able to use the Landau two-fluid equations as a direct probe of the superfluid density, just as one does in superfluid <sup>4</sup>He. One might be able to detect small differences between the magnitude of the superfluid density $`n_{s0}`$ and the condensate density $`n_{c0}`$ (at the level of our model in Section 4, $`n_{s0}`$ and $`n_{c0}`$ are equal). In this connection, we note that in the formal zero temperature limit, the Landau two-fluid equations reduce to two coupled equations for a pure superfluid $`(\rho _s=\rho ,\rho _n=0`$ at $`T=0)`$. This limit has been used to find a more accurate version of the $`T=0`$ quantum hydrodynamic equations discussed in Section 2.
In conclusion, I hope I have given some insight into why the two-fluid hydrodynamics of trapped Bose gases has so much potential interest. I hope some of the young experimentalists attending this Summer School will take up the challenge to study this new frontier. In a sense, the natural next step after understanding the dynamics of a pure condensate $`(T=0)`$ is to study the two-fluid hydrodynamics of trapped gases (at finite $`T`$) since now one has two components which can execute coupled coherent collective motions. In contrast, the collisionless region at finite $`T`$ as discussed in Section 3 does not seem as interesting since the thermal cloud has always such a low density (see Fig. 2). As a result, mean field effects produced by the non-condensate are never very important , above or below $`T_{\mathrm{BEC}}`$, and thus no new many-body dynamics emerges which is different than already found at $`T=0`$. What is new is the different mechanisms of damping of collisionless condensate motion, as discussed in Sections 3 and 4.
## Acknowledgments
I am grateful to Craig Savage and Mukunda Das for organizing this Australian BEC Summer School, for inviting me to participate, and their hospitality and assistance. Much of the work reviewed in these lectures was carried out with Eugene Zaremba, Tetsuro Nikuni, Jamie Williams and Milena Imamović-Tomasović. I would like to thank Jamie Williams for a critical reading and help with preparing the manuscript. I also thank Helen Iyer for texing the manuscript in record time. My research is supported by NSERC of Canada.
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# Siegel coordinates and moduli spaces for morphisms of Abelian varieties
## 1 Introduction
We describe the moduli spaces of morphisms between polarized abelian varieties over the complex field. Several instances of these modular varieties occur in the literature. The moduli spaces of isogenies between elliptic curves occur as Hecke correspondences on modular curves, cf. , . The moduli spaces of embeddings of elliptic curves into abelian surfaces appear through their images in the moduli spaces of abelian surfaces. The moduli spaces of abelian surfaces containing elliptic curves coincide with certain Humbert modular surfaces, as shown by Kani , and may also be seen as spaces of curves of genus two which cover elliptic curves, as in Murabayashi . Moduli spaces of isogenies on abelian surfaces have been studied by Hulek and Weintraub , and by Birkenhake and Lange .
In this paper we propose a general treatment, in terms of Siegel coordinates. We identify the discrete invariants which correspond to the irreducible components of the moduli space of morphisms, and we describe the irreducible components as quotients of products of Siegel spaces.
The discrete invariants are derived from a Poincaré decomposition of morphisms. A morphism $`VW`$ of polarized abelian varieties, after suitable isogenies $`X\times X^{}V`$ and $`Y\times Y^{}W`$, each a sum of two embeddings of complementary abelian subvarieties, lifts to a morphism $`X\times X^{}Y\times Y^{}`$ which is given by an isogeny $`XY`$ multiplied by the zero map on $`X^{}`$.
The sequence $`\delta `$ which collects the types of polarizations of the six varieties occurring in the decomposition is an invariant of the morphism. A second invariant is defined by taking into account how the morphisms in the decomposition behave with respect to the polarizations.
A morphism of abelian varieties $`XY`$ with respect to symplectic bases gives rise to an integer matrix $`M`$ which represents the induced homomorphism $`H_1(X,)H_1(Y,)`$, and what is only dependent on the polarizations is the equivalence class $`[M]`$ under the natural action of symplectic groups. In the same way a Poincaré decomposition, which consists of a sequence of three morphisms, determines a sequence $`\tau `$ of integer matrices relative to symplectic bases, and what is invariant is an equivalence class $`[\tau ]`$, defined in the natural way. The collective datum $`\delta ,[\tau ]`$ is what we call a type of morphisms between polarized abelian varieties, and we characterize the discrete data which are types of morphisms.
The global coordinates for morphisms are the Siegel coordinates of the varieties occurring in the Poincaré decomposition. As the type is fixed, it turns out that the morphism is determined if the varieties $`X,X^{}`$ and $`Y^{}`$ are given. Thus, for every type of morphisms we construct a complex analytic variety $`\text{F}_\delta [\tau ]`$ which is a coarse moduli space for morphisms of the given type. The moduli variety is irreducible, and is obtained from a product of Siegel spaces modulo the action of a discrete group.
The construction is presented in several steps. In §2 we introduce the invariant $`[M]`$ mentioned before. Using this, in §3 the moduli space of isogenies is constructed. In §4 we describe some known examples, such as the Hecke correspondences. In §5 we introduce a refined invariant, the type of an embedding, derived from classical Poincaré reducibility. In §6 the moduli space of embeddings is constructed. In §7 the Poincaré reducibility for morphisms is presented and in §8 the moduli space of general morphisms is constructed. Finally in §9 we give the proof of the moduli property for these spaces, in a sketched form, and in §10 we shortly discuss the question of classifying the discrete data which are types of morphisms.
## 2 A first invariant
Some notation. Every abelian variety $`X`$ will be endowed with a polarization $`L_X`$, that we identify with the corresponding alternating form on $`H_1(X,)`$. We denote by $`D=(d_1,\mathrm{},d_n)`$ the sequence of elementary divisors, the type of the polarization, where $`n`$ is the dimension of $`X`$. The alternating form will be determined by means of some symplectic basis $`\lambda =(\lambda _1,\mathrm{},\lambda _{2n})`$, such that $`L_X(\lambda _i,\lambda _j)`$ is equal to $`d_i`$ for $`j=i+n`$, to $`d_i`$ for $`i=j+n`$, and 0 otherwise, two symplectic bases for the polarization being related by some element of the symplectic group $`Sp(D,)`$. We use the symbol $`\widehat{D}`$ for the matrix of the alternating form with respect to a symplectic basis $`\lambda `$, and sometimes we write this as $`\lambda ,\lambda =\widehat{D}`$, dropping the symbol $`L_X`$. Note that the datum $`D`$ encodes the dimension $`n`$.
The collection of isomorphism classes of abelian varieties endowed with symplectic basis for a polarization of type $`D`$ will be denoted by $`\text{A}_D^{}`$. The superscript is to remind the choice of symplectic bases, and in this way it will be used again in the following. The bijective map $`\text{A}_D^{}\text{H}_n`$ with the Siegel space introduces the structure of a complex manifold. The collection of isomorphism classes of abelian varieties with polarization of type $`D`$ is bijective to the quotient variety $`\text{H}_n/Sp(D,)`$. In a more precise language it is well known that one should speak of varieties representing the appropriate moduli functors.
We consider morphisms $`f:XY`$ between abelian varieties, both endowed with polarization. As morphisms between abelian varieties we mean in the following the morphisms of analytic groups, the same as the morphisms of analytic varieties up to translations in the range.
An isomorphism $`ff^{}`$ with another morphism $`f^{}:X^{}Y^{}`$ is a couple of isomorphisms $`u:XX^{}`$ and $`v:YY^{}`$, which preserve the polarizations, such that $`f^{}=vfu^1`$. We denote by $`D,E`$ the pair of types of polarizations on $`X,Y`$, of dimensions $`n,m`$.
Definition. Define $`\text{F}_{D,E}^{}`$ to be the collection of morphisms $`f`$ of abelian varieties, both endowed with polarizations of the given types and with symplectic bases for the polarizations, modulo isomorphisms $`ff^{}`$ which preserve (the polarizations and) the bases. Define $`\text{F}_{D,E}`$ to be the collection of morphisms of abelian varieties with polarizations of the given types, modulo isomorphisms (which preserve the polarizations). Define $`\mathrm{\Gamma }_{D,E}=Sp(D,)\times Sp(E,)`$.
The group $`\mathrm{\Gamma }_{D,E}`$ acts on the set $`\text{F}_{D,E}^{}`$ by producing the change of symplectic bases. The element $`(A,B)`$ acts sending the isomorphism class represented by the morphism $`f`$ endowed with symplectic bases $`\gamma ,\lambda `$ into the isomorphism class represented by $`f`$ with the bases $`\gamma A`$ and $`\lambda B`$. There is an identification
$$\text{F}_{D,E}=\begin{array}{c}\text{F}_{D,E}^{}\\ \\ \mathrm{\Gamma }_{D,E}\end{array}$$
To a morphism $`f`$ is associated the homomorphism $`f_{}:H_1(X,)H_1(Y,)`$, called the rational representation of $`f`$. With respect to the symplectic bases $`\gamma ,\lambda `$, it is represented by an integer matrix $`M`$, that we also consider as a homomorphism $`^{2n}\stackrel{M}{}^{2m}`$. This defines a map
$$\text{F}_{D,E}^{}Hom(^{2n},^{2m})$$
sending the isomorphism class of $`(f,\gamma ,\lambda )`$ to the matrix $`M`$. We denote by $`\text{F}_{D,E}^{}(M)`$ the fibre over the matrix $`M`$.
The group $`\mathrm{\Gamma }_{D,E}`$ also acts on the set $`Hom(^{2n},^{2m})`$ by producing the equivalence of matrices that represent a given homomorphism $`f_{}`$ under a change of bases. The element $`(A,B)`$ acts as $`MBMA^1`$. The map above is equivariant with respect to the actions of $`\mathrm{\Gamma }_{D,E}`$, so there are isomorphisms of the fibres over two matrices in the same orbit. Consider the quotient map
$$\text{F}_{D,E}=\begin{array}{c}\text{F}_{D,E}^{}\\ \\ \mathrm{\Gamma }_{D,E}\end{array}\begin{array}{c}Hom(^{2n},^{2m})\\ \\ \mathrm{\Gamma }_{D,E}\end{array}$$
sending the isomorphism class of $`f`$ to the equivalence class $`[M]`$. We denote by $`\text{F}_{D,E}[M]`$ the fibre over the equivalence class $`[M]`$. There is an identification
$$\text{F}_{D,E}[M]=\begin{array}{c}\text{F}_{D,E}^{}(M)\\ \\ \mathrm{\Gamma }_{D,E}(M)\end{array}$$
where $`\mathrm{\Gamma }_{D,E}(M)`$ denotes the stabilizer of $`M`$
The equivalence class $`[M]`$ is a first invariant that may be used in order to separate the components of the space $`\text{F}_{D,E}`$.
## 3 Moduli of isogenies
The invariant defined in the preceding section is indeed sufficient for classifying the isogenies.
Let $`f:XY`$ be a morphism of abelian varieties of the same dimension $`n`$, endowed with polarizations of types $`D,E`$, and with symplectic bases, and let $`M`$ be the matrix of $`f_{}`$ relative to the bases. The morphism is an isogeny if and only if $`detM0`$. Then there is an exact sequence
$$0^{2n}\stackrel{M}{}^{2n}F0$$
(1)
where the group $`F`$ is defined as the cokernel of $`M`$. Moreover there is an isomorphism $`F\text{Ker}(f)`$.
We say that $`f`$ is an isogeny of polarized abelian varieties if $`L_X=f^{}L_Y`$ holds, and this happens if and only if
$$\widehat{D}={}_{}{}^{t}M\widehat{E}M$$
(2)
Note that $`detM0`$ follows from this. For instance one has polarizations of the same type if and only if $`M`$ is symplectic, and then $`F=0`$.
We call the datum $`\delta =(D,E)`$ the polarization type of the isogeny $`f`$. We call the datum $`\delta ,M`$ the type of the isogeny relative to the symplectic bases.
Definition. Let $`\delta =(D,E)`$ be a pair of types of polarizations, of the same dimension $`n`$. Define the subsets $`\text{I}_\delta ^{}\text{F}_{D,E}^{}`$ and $`\text{I}_\delta \text{F}_{D,E}`$ consisting of isomorphism classes of isogenies of polarized abelian varieties, with symplectic bases selected or not selected. Write moreover $`\mathrm{\Gamma }_\delta =\mathrm{\Gamma }_{D,E}`$.
So we have the restricted map
$$\text{I}_\delta ^{}GL(2n,)$$
The image is characterized by equality (2), and the fibres are described as $`\text{I}_\delta ^{}(M)=\text{F}_{D,E}^{}(M)`$. Then we have the quotient map
$$\text{I}_\delta =\frac{\text{I}_\delta ^{}}{\mathrm{\Gamma }_\delta }\frac{GL(2n,)}{\mathrm{\Gamma }_\delta }$$
sending the isomorphism class of the isogeny $`f`$ to the equivalence class $`[M]`$.
We say that the datum $`\delta ,[M]`$ is the type of the isogeny (tout-court, relative to the polarizations, independent of symplectic bases).
The fibre of the quotient map is described as
$$\text{I}_\delta [M]=\frac{\text{I}_\delta ^{}(M)}{\mathrm{\Gamma }_\delta (M)}$$
where $`\mathrm{\Gamma }_\delta (M)`$ denotes the stabilizer of $`M`$.
For a fixed type of isogenies, consider the two projections
$$\begin{array}{ccc}& \text{I}_\delta ^{}(M)& \\ & & \\ \text{A}_D^{}& & \text{A}_E^{}\end{array}$$
in which to $`f:XY`$ are associated the varieties $`X`$ and $`Y`$ respectively.
###### Proposition 3.1.
The two projections are bijective.
###### Proof.
There is a bijective correspondence on the level of comples tori, which restricts to abelian varieties. Given a complex torus $`Y`$ and a basis in $`H_1(Y)`$, i.e. an isomorphism $`H_1(Y)^{2n}`$, from the exact sequence (1) a representation $`H_1(Y)F0`$ is deduced, and this determines a complex torus $`X`$ together with a covering map $`XY`$, unique up to isomorphisms. Then from the exact sequence an isomorphism $`H_1(X)^{2n}`$ is also obtained. Conversely, given a complex torus $`X`$ and an isomorphism $`H_1(X)^{2n}`$, by means of the exact sequence an inclusion $`FX`$ is deduced, hence a quotient torus $`Y=X/F`$, together with an isomorphism $`H_1(Y)^{2n}`$.
Let $`f:XY`$ be a morphism of complex tori constructed as above. The basis in $`H_1(Y)`$ is symplectic for a unique antisymmetric form $`L_Y`$ of type $`E`$. This determines a natural $``$-bilinear form $`\overline{L}_Y`$ on the space $`H^0(\mathrm{\Omega }_Y)^{}`$ such that $`\text{im}(\overline{L}_Y)`$ coincides with $`L_Y`$ on $`H_1(Y)`$, and one has a polarization if and only if $`\overline{L}_Y`$ is a positive hermitian form. Similarly the basis in $`H_1(X)`$ is symplectic for an antisymmetric form $`L_X`$ of type $`D`$, which determines a form $`\overline{L}_X`$. The two bases are related through the matrix $`M`$. The relation (2) between the types is the condition so that there is equality $`L_X=f^{}L_Y`$ of the antisymmetric forms. This implies $`\overline{L}_X=f^{}\overline{L}_Y`$. As the induced isomorphism $`f^{}:H^0(\mathrm{\Omega }_X)^{}H^0(\mathrm{\Omega }_Y)^{}`$ is indeed $``$-linear, it follows that $`\overline{L}_X`$ is a positive hermitian form if and only if the same holds for $`\overline{L}_Y`$. ∎
The sets $`\text{A}_D^{},\text{A}_E^{}`$ both have the structure of an analytic variety, represented by the Siegel space $`\text{H}_n`$. Because of 3.1 there is a diagram of bijections
$$\begin{array}{ccc}& \text{I}_\delta ^{}(M)& \\ & & \\ \text{H}_n& & \text{H}_n\end{array}$$
###### Proposition 3.2.
The bijection in the bottom line is an automorphism of the Siegel space.
###### Proof.
The Siegel space $`\text{H}_n`$ may be viewed as a quotient: the space $`L(D)`$ of lattice bases $`\lambda =(\lambda _1,\mathrm{},\lambda _{2n})`$ in $`^n`$, subject to the locally closed condition that the alternating form defined by $`\lambda ,\lambda =\widehat{D}`$ is the imaginary part of some positive hermitian form (a condition that the Riemann relations express in terms of the Siegel coordinates of $`\lambda `$), divided by the natural action of $`GL(n,)`$. An equivariant isomorphism $`L(D)L(E)`$ induces an automorphism of $`\text{H}_n`$.
If the lattice basis $`\lambda `$ defines a torus $`Y`$ then the lattice basis $`\gamma :=\lambda M`$ determines the complex torus $`X`$ that covers $`Y`$; conversely if $`\gamma `$ defines a torus $`X`$ then $`\lambda :=\gamma M^1`$ defines the quotient torus $`Y`$. This correspondence $`\gamma \lambda `$ is an automorphism of the space of all lattice bases. It follows from the preceding proof that the correspondence restricts to an isomorphism $`L(D)L(E)`$. ∎
The stabilizer subgroup $`\mathrm{\Gamma }_\delta (M)`$ consists of the pairs $`(A,B)\mathrm{\Gamma }_\delta `$ such that $`MA=BM`$. In other words there is a diagram of group homomorphisms
$$\begin{array}{ccc}& \mathrm{\Gamma }_\delta (M)& \\ & {}_{}{}^{\alpha }^\beta & \\ GL(2n,)& \begin{array}{c}\\ _{\varphi _M}\end{array}& GL(2n,)\end{array}$$
where $`\alpha ,\beta `$ are the two projections, and where $`\varphi _M`$ is the inner automorphism such that $`AMAM^1=B`$. Equality (2) implies that $`\varphi _M`$ restricts to an isomorphism $`Sp(D,)Sp(E,)`$ of the rational symplectic groups. So there is an induced isomorphism
$$\begin{array}{ccc}Sp(D,)\varphi _M^1GL(2n,)& & \varphi _MGL(2n,)Sp(E,)\\ _{||}& & _{||}\\ \alpha \mathrm{\Gamma }_\delta (M)& & \beta \mathrm{\Gamma }_\delta (M)\end{array}$$
The actions of the groups $`\alpha \mathrm{\Gamma }_\delta (M)`$ and $`\beta \mathrm{\Gamma }_\delta (M)`$ on $`\text{H}_n`$ are properly discontinuous, because they are restrictions of the actions of symplectic groups. It is well known that on the quotient spaces there are natural structures of analytic varieties. Moreover in the diagram of bijections
$$\begin{array}{ccc}& \text{I}_\delta [M]& \\ & & \\ \begin{array}{c}\text{H}_n\\ \\ \alpha \mathrm{\Gamma }_\delta (M)\end{array}& & \begin{array}{c}\text{H}_n\\ \\ \beta \mathrm{\Gamma }_\delta (M)\end{array}\end{array}$$
the arrow in the bottom line is an isomorphism of varieties. This determines on $`\text{I}_\delta [M]`$ a unique structure of a complex analytic variety. Being a quotient of $`\text{H}_n`$, this is an irreducible variety.
###### Theorem 3.3.
The variety $`\text{I}_\delta [M]`$ is a coarse moduli space for isogenies of the given type.
The proof will be given in section 9.
## 4 Hecke correspondences
We show how the previous description looks for isogenies of elliptic curves, the case $`n=1`$.
Recall that on an elliptic curve the polarizations are just the positive integer multiples of the unique principal polarization, and that an isogeny of elliptic curves acts on polarizations as multiplication by the degree. With respect to the notation in the previous sections, this means that if we take $`E=(1)`$, the principal type, then necessarily $`D=(d)`$, the degree of the isogeny, and condition (2) requires that $`detM=\pm d`$. Recall moreover that the symplectic group is the full $`SL(2,)`$, independent of the polarization degree, and that every $`M`$ is reduced under unimodular transformations to a unique diagonal form with diagonal $`(d_1,d_2)`$ consisting of positive integers such that $`d_1|d_2`$, and necessarily $`d_1d_2=d`$. With this choice of data we simply write $`\text{I}(d_1,d_2)`$ instead of $`\text{I}_\delta [M]`$ for the moduli space of isogenies of the given type.
This implies for the kernel of the isogeny, occurring in (1), an isomorphism $`F_{d_1}\times _{d_2}`$. If $`d_1=1`$ we have a cyclic subgoup of the elliptic curve which dominates in the isogeny, and so we see that the space $`\text{I}(1,p)`$ is identified to the modular curve $`X_0(p)`$ parametrizing elliptic curves with distinguished cyclic subgroup of order $`p`$. If $`d_1>1`$ then the isogenies factor as $`XZY`$, an isogeny of type $`(1,d_1)`$ followed by an isogeny of type $`(1,d_2/d_1)`$.
Let $`a,b,p`$ be positive integers such that $`a|b`$ and $`(b,p)=1`$, and define $`d_1=a`$ and $`d_2=bp`$. Because of the isomorphism $`_{d_1}\times _{d_2}(_a\times _b)\times _p`$, every isogeny $`f`$ of type $`(d_1,d_2)`$ factors as $`X\stackrel{u}{}U\stackrel{g}{}Y`$, where $`u`$ is of type $`(a,b)`$ and $`g`$ is of type $`(1,p)`$, and also factors as $`X\stackrel{h}{}V\stackrel{v}{}Y`$ where $`h`$ is of type $`(1,p)`$ and $`v`$ is of type $`(a,b)`$. In this way two morphisms $`\text{I}(d_1,d_2)X_0(p)`$ are defined, sending $`fg`$ and $`fh`$. The diagram
$$\begin{array}{ccc}& \text{I}(d_1,d_2)& \\ & & \\ X_0(p)& & X_0(p)\end{array}$$
is known as a Hecke correspondence on the modular curve of degree $`p`$. This correspondence is studied in , . We end with recalling two more examples which fall within the scope of the present treatment.
###### Example 4.1.
Level structures. If $`X`$ is a variety with polarization of type $`D`$, there is an isogeny $`\varphi :X\widehat{X}`$ onto the dual variety. A structure on $`X`$ of level $`D`$ is a symplectic basis of $`\text{Ker}(\varphi )`$. This is represented by a symplectic basis $`\lambda `$ of $`H_1(X)`$ such that $`\varphi _1(\lambda )(D^1\times D^1)`$ is a basis of $`H_1(\widehat{X})`$. The basis of $`H_1(\widehat{X})`$ is symplectic for a principal polarization, type $`E=1`$. With respect to the bases the isogeny $`\varphi `$ is represented by the matrix $`M=D\times D`$. In other words the space of polarized varieties with structure of level $`D`$ is a discrete quotient of the space $`\text{I}_{}^{}{}_{(D,1)}{}^{}(D\times D)`$. The quotient is by the discrete group which produces in $`X`$ a change of symplectic basis without changing the symplectic basis in $`\text{Ker}(\varphi )`$. See \[6, Ch. 8, §3\].
###### Example 4.2.
Varieties with isogeny. If $`X`$ is a variety with polarization of type $`D`$, an isogeny $`f:XY`$ such that $`\text{Ker}(f)`$ admits as elementary divisors the elements of the diagonal $`D`$ is called an isogeny of type $`D`$. A canonical basis of $`\text{Ker}(f)`$ is represented by a symplectic basis $`\lambda `$ of $`H_1(X)`$ such that $`f_{}(\lambda )(1\times D^1)`$ is a basis of $`H_1(Y)`$. The basis of $`H_1(Y)`$ is a symplectic basis for a principal polarization on $`Y`$, from which the polarization on $`X`$ is obtained via pullback. With respect to the bases the isogeny $`f`$ is represented by the matrix $`M=1\times D`$. In other words the space of abelian varieties with isogeny of type $`D`$ coincides with the space $`\text{I}_{(D,1)}[1\times D]`$. See \[6, p. 245\]. The modular variety of abelian surfaces with isogeny of type $`D=(1,p)`$ is studied in , .
## 5 Refined invariant from reducibility
For embeddings of abelian varieties a refined invariant is derived from the well known
Poincaré reducibility theorem. Let $`e:XY`$ be an embedding into a polarized abelian variety. There is an abelian subvariety $`X^{}Y`$ such that the sum $`s:X\times X^{}Y`$ is an isogeny. In other words $`XX^{}`$ is a finite subgroup and $`X+X^{}=Y`$. Moreover if $`L_X:=e^{}L_Y`$ and $`L_X^{}:=e_{}^{}{}_{}{}^{}L_Y`$ are the induced polarizations, then $`s^{}L_Y=p^{}L_Xp^{}{}_{}{}^{}L_{X^{}}^{}`$ in the Neron-Severi group of $`X\times X^{}`$, where $`p,p^{}`$ are the two projections. Because of this property the complementary variety $`X^{}`$ is uniquely determined.
Definition. We say that an isogeny $`s:X\times X^{}Y`$ is a sum of embeddings if so are the restrictions $`X,X^{}Y`$. If $`Y`$ is polarized we say that the isogeny is a sum of complementary embeddings if moreover $`s^{}L_Y=p^{}L_Xp^{}{}_{}{}^{}L_{X^{}}^{}`$ holds, where $`L_X,L_X^{}`$ are the induced polarizations.
The reducibility theorem establishes a 1-1 correspondence between embeddings $`e:XY`$ into a polarized abelian variety and isogenies $`s:X\times X^{}Y`$ which are sum of complementary embeddings, and this allows to describe the space of embeddings as the space of isogenies of the special form. Isomorphisms of embeddings $`ee^{}`$ will correspond with the natural definition of isomorphisms $`ss^{}`$ of isogenies of the special form.
Let $`D,D^{}`$ and $`E`$ be the types of the polarizations on $`X,X^{}`$ and $`Y`$. Introduce symplectic bases in $`X,Y`$ and also in $`X^{}`$. The induced homomorphism $`s_{}:H_1(X\times X^{})H_1(Y)`$ is represented by an exact sequence
$`0^{2n}\times ^{2n^{}}\stackrel{M}{}^{2(n+n^{})}F0`$ (3)
where $`M`$ has nonzero determinant, and where $`F`$ is defined as the cokernel. Note that the invariant of section 2 is the restriction $`^{2n}^{2(n+n^{})}`$.
The relation $`p^{}L_Xp^{}{}_{}{}^{}L_{X^{}}^{}=s^{}L_Y`$ of the polarizations is equivalent to
$$\widehat{D}\times \widehat{D}^{}={}_{}{}^{t}M\widehat{E}M$$
(4)
where the left hand side is a diagonal block matrix.
We call the datum $`\delta =(D,D^{},E)`$ the polarization type of the embedding, and we call the full datum $`\delta ,M`$ the type of the embedding relative to the symplectic bases.
Definition. Let $`\delta =(D,D^{},E)`$ be a sequence of polarization types, of dimensions $`n,n^{}`$ and $`n+n^{}`$. Define $`\text{E}_\delta ^{}`$ to be the collection of embeddings associated to isogenies $`s`$ which are sum of complementary embeddings, of the given polarization type, all varieties endowed with symplectic bases, modulo isomorphisms $`ss^{}`$ which preserve polarizations and bases. Define $`\text{E}_\delta `$ to be the collection of embeddings associated to isogenies which are sum of complementary embeddings, of the given polarization type, modulo isomorphisms (which preserve the polarizations). Define $`\mathrm{\Gamma }_\delta =Sp(D,)\times Sp(D^{},)\times Sp(E,)`$.
The group $`\mathrm{\Gamma }_\delta `$ acts on $`\text{E}_\delta ^{}`$ by producing the change of symplectic bases. The element $`(A,A^{},B)`$ acts sending the isomorphism class represented by the isogeny $`s`$ endowed with symplectic bases $`\gamma ,\gamma ^{},\lambda `$ into the isomorphism class represented by $`s`$ with the bases $`\gamma A,\gamma ^{}A^{}`$ and $`\lambda B`$. There is an identification
$$\text{E}_\delta =\begin{array}{c}\text{E}_\delta ^{}\\ \\ \mathrm{\Gamma }_\delta \end{array}$$
We have seen that there is a natural map
$$\text{E}_\delta ^{}GL(2(n+n^{}),)$$
which to the isomorphism class of $`(s,\gamma ,\gamma ^{},\lambda )`$ associates the matrix $`M`$.
The group $`\mathrm{\Gamma }_\delta `$ also acts on the set of matrices, by producing the equivalence of matrices for a given homomorphism $`s_{}`$ under a change of symplectic bases. The element $`(A,A^{},B)`$ acts sending $`MBM(A\times A^{})^1`$. The classifying map above is equivariant with respect to the actions of $`\mathrm{\Gamma }_\delta `$ and therefore induces a map
$$\text{E}_\delta =\begin{array}{c}\text{E}_\delta ^{}\\ \\ \mathrm{\Gamma }_\delta \end{array}\begin{array}{c}GL(2(n+n^{}),)\\ \\ \mathrm{\Gamma }_\delta \end{array}$$
which to the isomorphism class of $`s`$ associates the equivalence class $`[M]`$. The datum $`\delta ,[M]`$ is what we call the type of the embedding.
The fibre over an equivalence class $`[M]`$ is described as
$$\text{E}_\delta [M]=\begin{array}{c}\text{E}_\delta ^{}(M)\\ \\ \mathrm{\Gamma }_\delta (M)\end{array}$$
where $`\mathrm{\Gamma }_\delta (M)`$ denotes the stabilizer of $`M`$.
If the fibre $`\text{E}_\delta ^{}(M)`$ is nonempty the datum $`\delta ,M`$ is an effective type of embeddings. In addition to equality (4), it is clear that the matrix $`M`$ has to satisfy one more property, the property that any isogeny represented by $`M`$ will be a sum of two embeddings. We have the following characterization.
###### Lemma 5.1.
The additional condition so that the datum $`\delta ,M`$ is a type of embeddings is that $`M`$ fits into a diagram
$$\begin{array}{ccccccccc}0& & ^{2n}\times ^{2n^{}}& \stackrel{R\times R^{}}{}& ^{2n}\times ^{2n^{}}& & F\times F& & 0\\ & & ||& & & & & & \\ 0& & ^{2n}\times ^{2n^{}}& \underset{M}{}& ^{2(n+n^{})}& & F& & 0\end{array}$$
where the bottom line is the pullback of the top line under the diagonal homomorphism of $`F`$, and where the top line is a product of two exact sequences.
###### Proof.
The isogeny is a sum of embeddings if and only if the inclusion $`FX\times X^{}`$ is given by a pair of inclusions $`FX,X^{}`$. If the condition is satisfied, as in the beginning of the section, calling $`\overline{X}=X/F,\overline{X}^{}=X^{}/F`$ the quotient varieties, the product isogeny $`X\times X^{}\overline{X}\times \overline{X}^{}`$ factors as
$$\begin{array}{ccc}X\times X^{}& & \overline{X}\times \overline{X}^{}\\ ||& & & & \\ X\times X^{}& & Y\end{array}$$
Introducing bases also in the homology of the varieties $`\overline{X}`$ and $`\overline{X}^{}`$, the diagram is obtained. Conversely if a diagram of integer homomorphisms as above exists for $`M`$ then the inclusion of the subgroup into the product variety is given by a pair of inclusions. ∎
## 6 Moduli of embeddings
The refined invariant introduced in the preceding section is sufficient for classifying the embeddings.
###### Proposition 6.1.
For every type of embeddings there is a natural bijection
$$\text{E}_\delta ^{}(M)\text{A}_D^{}\times \text{A}_D^{}^{}$$
###### Proof.
To an embedding is associated a pair of varieties, as we have seen. Conversely, given two abelian varieties $`X,X^{}`$, with symplectic bases for polarizations of types $`D,D^{}`$, by means of the isomorphism $`H_1(X\times X^{})^{2n}\times ^{2n^{}}`$, from the exact sequence (3) an inclusion $`FX\times X^{}`$ is deduced. Define the torus $`Y=X\times X^{}/F`$ and call $`s:X\times X^{}Y`$ the quotient isogeny. Then one has an isomorphism $`H_1(Y)^{2(n+n^{})}`$, a symplectic basis for an alternating form $`L_Y`$ of type $`E`$. Condition (4) on the types implies that $`p^{}L_Xp^{}{}_{}{}^{}L_{X^{}}^{}=s^{}L_Y`$. Using the same argument as in the proof of Proposition 3.1 it is then seen that $`s^{}L_Y`$ being a polarization implies that $`L_Y`$ is a polarization. Finally because of Lemma 5.1 the isogeny $`s`$ is the sum of two embeddings. In particular the restriction $`XY`$ is an embedding of the given type. ∎
The stabilizer $`\mathrm{\Gamma }_\delta (M)`$ consists of the triplets $`(A,A^{},B)\mathrm{\Gamma }_\delta `$ such that $`M(A\times A^{})=BM`$. In other words there is a diagram of inclusions
$$\begin{array}{ccc}& \mathrm{\Gamma }_\delta (M)& \\ & {}_{}{}^{\alpha }& \\ GL(2n,)\times GL(2n^{},)& \begin{array}{c}\\ _{\varphi _M}\end{array}& GL(2n+2n^{},)\end{array}$$
where the descending arrows are the natural projections, and the horizontal arrow is given by $`(A,A^{})M(A\times A^{})M^1`$. Equality (4) implies that $`Sp(D,)\times Sp(D^{},)=\varphi _M^1Sp(E,)`$. Therefore
$$\alpha \mathrm{\Gamma }_\delta (M)=Sp(D,)\times Sp(D^{},)\varphi _M^1GL(2n+2n^{},)$$
A product of Siegel spaces $`\text{H}_n\times \text{H}_n^{}`$ is the moduli space of embeddings with symplectic bases of polarizations of the given types. The group $`\alpha \mathrm{\Gamma }_\delta (M)`$ acts on the product through the actions of symplectic groups. Hence the action is properly discontinuous and the quotient variety exists. The bijection
$$\text{E}_\delta [M]\begin{array}{c}\text{H}_n\times \text{H}_n^{}\\ \\ \alpha \mathrm{\Gamma }_\delta (M)\end{array}$$
introduces on the left hand side the structure of a complex analytic variety. It is an irreducible variety.
###### Theorem 6.2.
The variety $`\text{E}_\delta [M]`$ is a coarse moduli space for embeddings of the given type.
The proof will be given in section 9.
###### Example 6.3.
Embeddings of elliptic curves into principally polarized abelian surfaces, the case $`n=n^{}=1`$ and $`E=I_2`$. In this situation the pullback polarization is of type $`D=(k)`$ where $`k=X\mathrm{\Theta }_Y`$ is the degree relative to the polarization. Similarly one has $`D^{}=(k^{})`$ and necessarily $`detM=\pm kk^{}`$. This is seen for instance in , where one also finds some first step in the classification of the types of embeddings (Lemma 1) to the effect that the matrix $`M`$ may be reduced to a form
$$\left(\begin{array}{cccc}0& \hfill k& & \\ 0& \hfill 0& & \\ 1& \hfill 0& & \\ 0& \hfill 1& & \end{array}\right)$$
The moduli space $`\text{E}_\delta (M)`$ is a surface, a discrete quotient of $`\text{H}_1\times \text{H}_1`$. Consider the natural map $`\text{E}_\delta (M)\text{A}_2`$. The space of principally polarized abelian surfaces which contain some elliptic curve of degree $`k`$ is known to be irreducible, and is known to coincide with the so called Humbert modular surface of invariant $`\mathrm{\Delta }=k^2`$. See for instance , which also contains the historical references. This means that for fixed $`k`$ all $`\text{E}_\delta (M)`$ have the same image in $`\text{A}_2`$.
## 7 Poincaré decomposition of morphisms
The treatment of general morphisms is done by patching together the arguments developed so far for isogenies and embeddings. The starting point is a Poincaré decomposition of morphisms.
###### Proposition 7.1.
Let $`f:VW`$ be a morphism of abelian varieties, endowed with polarizations. There are complementary abelian subvarieties $`X,X^{}V`$ and $`Y,Y^{}W`$, and an isogeny $`g:XY`$ such that the following is a commutative diagram
$$\begin{array}{ccc}X\times X^{}& & V\\ g\times 0& & f& & \\ Y\times Y^{}& & W\end{array}$$
(5)
where the horizontal arrows are isogenies, sums of the given embeddings. The diagram is uniquely determined from the polarizations.
###### Proof.
Take $`X^{}=\text{Ker}_0(f)`$, the connected component of 0 in the kernel, take $`Y=f(V)`$, let $`X`$ and $`Y^{}`$ be the complementary abelian subvarieties, and let $`g`$ be the restriction of $`f`$. Uniqueness of the choice is a consequence of uniqueness in the classical reducibility theorem. ∎
Definition. In the decomposition above if $`X^{}0`$ it is not possible to have on $`V`$ a pullback polarization. However it makes sense to point out the case in which the isogeny $`g:XY`$ preserves the polarizations. In this situation we say that $`f`$ is compatible with the polarizations, or simply a morphism of polarized abelian varieties.
The proposition establishes a 1-1 correspondence between morphisms $`f:VW`$ which are compatible with the polarizations, and sequences of the form
$$\begin{array}{ccc}X\times X^{}& & V\\ g\times 0& & \\ Y\times Y^{}& & W\end{array}$$
where the horizontal arrows are isogenies, each a sum of complementary embeddings, and $`g`$ is an isogeny which preserves the polarizations, and moreover the sequence is of some special type, such that a morphism $`f`$ actually exists which fills in a commutative diagram (5). This allows to describe the space of morphisms as the space of sequences of the special type.
Let $`D,D^{},E`$ and $`H,H^{},K`$ be the polarization types of $`X,X^{},V`$ and $`Y,Y^{},W`$. Introducing symplectic bases in all these varieties, the associated diagram of homology groups is represented by a diagram
$$\begin{array}{ccccccccc}0& & ^{2n}\times ^{2n^{}}& \stackrel{M}{}& ^{2(n+n^{})}& & F& & 0\\ & & P\times 0& & & & & & \\ 0& & ^{2m}\times ^{2m^{}}& \underset{N}{}& ^{2(m+m^{})}& & G& & 0\end{array}$$
(6)
where $`n=m`$ and where the matrices $`M,N`$ and $`P`$ have nonzero determinant.
The cokernels $`F,G`$ are isomorphic to the kernels of the horizontal isogenies in the decomposition, i.e. the subgroups $`XX^{}`$ and $`YY^{}`$. Thus the morphism $`f`$ is isomorphic to a morphism
$$\begin{array}{c}X\times X^{}\\ \\ F\end{array}\begin{array}{c}Y\times Y^{}\\ \\ G\end{array}$$
which is a quotient of the product $`g\times 0`$ in the diagram.
We denote by $`\delta =(D,D^{},E,H,H^{},K)`$ the sequence of types of polarizations, and call it the polarization type of the morphism. We denote by $`\tau =(M,N,P)`$ the sequence of matrices which represent the decomposition of the morphism, and we call the datum $`\delta ,\tau `$ the type of the morphism relative to the symplectic bases.
The types $`D,D^{},E`$ are related to $`M`$ in equality (4), and similarly the types $`H,H^{},K`$ are related to $`N`$. As the morphism is compatible with the polarizations then $`\widehat{D}={}_{}{}^{t}P\widehat{H}P`$ also holds as in (2). The matrices $`M,N,P`$ satisfy a number of conditions. The product $`N(0\times P)M^1`$ is an integer matrix, associated to $`f`$. Both matrices $`M,N`$ satisfy the condition of Lemma 5.1. Finally it is quite clear that one more property holds. It is the property that, if a sequence of isogenies of the special form is given, and if the sequence is of the given type $`\delta ,\tau `$, then a morphism $`f`$ actually exists which fills in the commutative diagram (5). A characterization is given in the following lemma. If all these properties are verified then the datum $`\delta ,\tau `$ is an effective type of morphisms.
###### Lemma 7.2.
The additional condition for the datum $`\delta ,\tau `$ as above to be a type of morphisms is that the matrix $`P`$ fits in some diagram
$$\begin{array}{ccccccccc}0& & ^{2n}& \stackrel{R}{}& ^{2n}& & F& & 0\\ & & ||& & \overline{P}& & \\ & & ^{2n}& \underset{P}{}& ^{2n}\end{array}$$
in other words that $`P=\overline{P}R`$ with $`\text{Coker}(R)=F`$.
###### Proof.
By hypothesis the matrices of the isogenies satisfy the condition of Lemma 5.1 and there are inclusions $`FX,X^{}`$ and $`GY,Y^{}`$ as is seen in the proof of the lemma. Using these inclusions it is easily seen that there exists $`f:VW`$ if and only if $`g\times 0`$ sends $`FG`$ and this happens if and only if $`g:XY`$ sends $`F0`$. If this happens then there is a factorization $`X\overline{X}Y`$ where $`\overline{X}=X/F`$. Introducing bases in homology as in the proof of the lemma, the associated matrices satisfy $`P=\overline{P}R`$ where $`\overline{P}`$ is the matrix of $`\overline{X}Y`$ and where $`R`$ is the matrix of $`X\overline{X}`$. Conversely if the matrix $`\overline{P}`$ exists then $`F0`$ under $`XY`$. ∎
## 8 General moduli spaces
Finally we show that the moduli spaces of morphisms of given types are irreducible varieties.
Definition. Let $`\delta `$ be a sequence of 6 polarization types, whose dimensions are in the relation as in the preceding section, diagram (6). Define $`\text{F}_\delta ^{}`$ to be the collection of morphisms $`f`$ admitting a decomposition (5) of the given polarization type, all varieties endowed with symplectic bases, modulo isomorphisms which preserve polarizations and bases. Define $`\text{F}_\delta `$ to be the collection of morphisms admitting a decomposition of the given polarization type, modulo isomorphisms (which preserve the polarizations). Define $`\mathrm{\Gamma }_\delta `$ to be the product of the 6 symplectic groups for the types in $`\delta `$.
The group $`\mathrm{\Gamma }_\delta `$ acts on $`\text{F}_\delta ^{}`$ by producing the change of symplectic bases, and there is an identification
$$\begin{array}{ccc}\text{F}_\delta & =& \begin{array}{c}\text{F}_\delta ^{}\\ \\ \mathrm{\Gamma }_\delta \end{array}\end{array}$$
For the natural map on $`\text{F}_\delta ^{}`$ which to the isomorphism class of a morphism $`f`$ endowed with bases associates the type $`\tau `$ relative to the bases, we denote by $`\text{F}_\delta ^{}(\tau )`$ the fibre over $`\tau `$. It is nonempty if the datum $`\delta ,\tau `$ is a type of morphisms.
The group $`\mathrm{\Gamma }_\delta `$ also acts on the set of sequences $`\tau `$ by producing the equivalence under the change of bases. The classifying map on $`\text{F}_\delta ^{}`$ is equivariant, and therefore defines a map on $`\text{F}_\delta `$ which to the isomorphism class of $`f`$ associates the equivalence class $`[\tau ]`$. The datum $`\delta ,[\tau ]`$ is what we call the type of the morphism.
If $`\text{F}_\delta [\tau ]`$ is the fibre over $`[\tau ]`$, we have the natural identification
$$\begin{array}{ccc}\text{F}_\delta [\tau ]& =& \begin{array}{c}\text{F}_\delta ^{}(\tau )\\ \\ \mathrm{\Gamma }_\delta (\tau )\end{array}\end{array}$$
where $`\mathrm{\Gamma }_\delta (\tau )`$ is the subgroup of $`\mathrm{\Gamma }_\delta `$ stabilizer of $`\tau `$.
###### Proposition 8.1.
For every type of morphisms there is a natural bijection
$$\text{F}_\delta ^{}(\tau )\text{I}_{D,H}^{}(P)\times \text{A}_D^{}^{}\times \text{A}_H^{}^{}$$
###### Proof.
Assume that we are given abelian varieties $`X,X^{},Y,Y^{}`$ with symplectic bases for polarizations of types $`D,D^{},H,H^{}`$, and an isogeny $`g:XY`$ of type $`(D,H),P`$. Because of Lemma 5.1 for $`M`$ there are inclusions $`FX,X^{}`$. Then define the torus $`V=X\times X^{}/F`$. By means of diagram (6) an isomorphism $`H_1(V)^{2(n+n^{})}`$ is obtained. Because the types $`D,D^{},E`$ are related, this gives on $`V`$ a basis for a polarization of type $`E`$, as follows from Proposition 6.1. Similarly there are inclusions $`GY,Y^{}`$ and defining $`W=Y\times Y^{}/G`$ one has an isomorphism $`H_1(W)^{2(m+m^{})}`$, a basis for a polarization of type $`K`$. From Lemma 7.2 it follows that the morphism $`g\times 0:X\times X^{}Y\times Y^{}`$ sends $`FG`$ and induces a morphism $`f:VW`$. With respect to the bases this morphism is of the given type $`\delta ,\tau `$. ∎
It was seen in 3.2 that the family $`\text{I}_{D,H}^{}(P)`$ is represented by the Siegel space $`\text{H}_n`$. It follows from the proposition above that $`\text{F}_\delta ^{}(\tau )`$ is represented by the product $`\text{H}_n\times \text{H}_n^{}\times \text{H}_m^{}`$, and therefore there is a bijective map
$$\text{F}_\delta [\tau ]\frac{\text{H}_n\times \text{H}_n^{}\times \text{H}_m^{}}{\mathrm{\Gamma }_\delta (\tau )}$$
It is seen in the usual way that the stabilizer $`\mathrm{\Gamma }_\delta (\tau )`$ acts properly discontinuously, so the quotient variety exists, and this introduces on the left hand side the structure of a complex analytic variety, and an irreducible one.
###### Theorem 8.2.
The variety $`\text{F}_\delta [\tau ]`$ is a coarse moduli space for morphisms of the given type.
The proof will be given in section 9.
## 9 The moduli property
In this section we give the proof of the moduli property, in a sketched form. The proof is based on an analytic treatment of the moduli property of the Siegel modular variety.
The moduli property means two things. First, each modular variety parametrizes a family of morphisms. This is constructed in a natural way from the universal family of abelian varieties over the Siegel space. The details are lenghty, but natural, and will be omitted. Second, every flat family of morphisms determines a morphism to the suitable variety of moduli. This part of the proof will be sketched in the various cases.
Proof of 3.3. A flat family of isogenies is a surjective morphism $`f:\text{X}\text{Y}`$ between two flat families of polarized abelian varieties of the same dimension over some scheme $`T`$. Locally over open subschemes $`T^{}T`$ the two families can be endowed with symplectic bases and this determines two liftings $`T^{}\text{A}_D^{},\text{A}_E^{}`$. With respect to the bases all isogenies in the family are represented by some constant integer matrix, which we may assume is the given $`M`$. Then the two liftings above are two morphisms $`T^{}\text{H}_n`$ which are identified under the automorphism $`\text{H}_n\text{H}_n`$ of Proposition 3.2, and so both represent on $`T^{}`$ the same morphism $`T\text{I}_\delta [M]`$. ∎
Proof of 6.2. A flat family of embeddings $`\text{X}\text{Y}`$ into a polarized family of abelian varieties determines a flat family of complementary varieties $`\text{X}^{}`$, with the induced polarizations. Locally over open subschemes $`T^{}`$ of the parameter scheme $`T`$ the families can be endowed with symplectic bases, and this determines two liftings $`T^{}\text{A}_D^{},\text{A}_D^{}^{}`$ and so a morphism $`T^{}\text{H}_n\times \text{H}_n^{}`$ which locally represents a morphism $`T\text{E}_\delta [M]`$. ∎
Proof of 8.2. Let $`f:\text{V}\text{W}`$ be a flat family of morphisms of the given type, parametrized by $`T`$. By hypothesis all morphisms in the family are of the same rank. Because of this, there is a flat subfamily $`\text{X}^{}\text{V}`$ which parametrizes the connected kernels $`\text{Ker}_0(f_t)`$, and the image $`\text{Y}=f(\text{V})`$ also is a flat family. Then the families $`\text{X},\text{Y}^{}`$ of complementary varieties are determined relative to the polarizations, and finally is defined a family of isogenies $`g:\text{X}\text{Y}`$. They are all flat families. Locally over open subschemes $`T^{}T`$ the families of abelian varieties can be endowed with symplectic bases for the type $`\delta `$, with respect to which the morphisms occurring in the factorization of $`f_t`$ are represented by a sequence of integer matrices, independent of $`t`$, belonging to the given type, we may assume it is precisely the sequence $`\tau `$. From the universal property of the Siegel space and from Theorem 1.1 it follows that there is a morphism $`T^{}\text{I}_{D,H}^{}(P)\times \text{A}_D^{}^{}\times \text{A}_H^{}^{}`$ which locally represents a morphism $`T\text{F}_\delta [\tau ]`$. ∎
Finally we may define the total space $`\text{F}=\text{F}_\delta [\tau ]`$, disjoint sum of countably many analytic varieties. This will be a coarse moduli space for all morphisms of polarized abelian varieties. What is only required in addition is the observation that in a family of morphisms $`f:\text{V}\text{W}`$ the image $`f(\text{V})`$ has constant fibre dimension over the parameter $`T`$. This is a general statement of rigidity for morphisms of abelian varieties, and will be considered elsewhere. Then as in the preceding proof the decomposition is constructed, with respect to which the family is of some well determined type $`\delta ,[\tau ]`$, and finally a morphism $`T\text{F}_\delta [\tau ]`$ is obtained.
## 10 Some questions in integral linear algebra
We have seen that the components of the moduli space of morphisms between polarized abelian varieties correspond to equivalence classes of certain rather complicated discrete data. Thus the question arises of finding canonical forms for these data. All what we know is contained in the examples of section 4 and in example 6.3. Here we recollect some results and formulate some natural questions.
The types of isogenies are the data of the form $`\delta ,M`$ where $`\delta =(D,E)`$ is a pair of polarization types of the same dimension $`n`$, and $`M`$ is an integer square matrix of order $`2n`$, such that $`\widehat{D}={}_{}{}^{t}M\widehat{E}M`$ holds (equality (2)).
For given $`\delta `$ we ask whether some $`M`$ exists. A necessary condition is that $`E`$ divides $`D`$, as follows from \[2, Ch. 9, p. 85, ex. 1a\]. Moreover, a solution $`M`$ is equivalent to $`BMA^1`$ for every element $`(A,B)`$ of the group $`\mathrm{\Gamma }_\delta `$, a pair of symplectic matrices for the types $`D,E`$. It would be useful to find canonical forms, i.e. the simplest possible forms in the equivalence classes $`[M]`$. Finally, it is clear that equality (2) says that the smaller datum $`E,M`$ determines the datum $`D`$. It might be useful to rephrase the whole treatment in terms of the smaller datum.
The types of embeddings are the data of the form $`\delta ,M`$ where $`\delta =(D,D^{},E)`$ is a sequence of polarization types, of dimensions $`n,n^{},n+n^{}`$, and $`M`$ is an integer square matrix of order $`2(n+n^{})`$, which is obtained from some product $`R\times R^{}`$ in the way described in Lemma 5.1, and such that $`\widehat{D}\times \widehat{D^{}}={}_{}{}^{t}M\widehat{E}M`$ holds (equality (4)).
For the same reason as before we find that a necessary condition is that $`E`$ divides the type of the matrix $`\widehat{D}\times \widehat{D^{}}`$, which is determined by the pair $`D,D^{}`$, cf. \[2, Ch. 7, p. 94, rem. 3\]. Moreover, we may replace $`M`$ with any equivalent matrix of the form $`BM(A\times A^{})^1`$, for any element $`(A,A^{},B)`$ of the group $`\mathrm{\Gamma }_\delta `$, and it would be useful to find some simpler equivalent form. It is possible that this problem, for embeddings, might become easier to work with once we have some solution of the previous problem, for isogenies.
The analogous discussion can be done about the classification of the general types of morphisms.
address: Dipartimento di Matematica e Informatica, Università di Perugia, via Vanvitelli 1, 06123 Perugia, Italy
e-mail: guerra@unipg.it
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# Contents
## 1 Introduction
Questions of how to compare hadronic observables to the apparent underlying fundamental theory of QCD lie at the heart of understanding the nature of strong interactions. Thirty years after its inception, QCD in $`D=4`$ spacetime dimensions still stubbornly refuses to admit a global solution. The asymptotic freedom property of the theory permits the perturbative calculation of (Euclidean) Green functions involving large values of momentum transfer or energy release in terms of quarks and gluons, the fundamental objects of QCD. But at lower scales one enters the nonperturbative regime, which not only invalidates (or at least complicates) the standard perturbative methods of field theory developed in QED, but also leads to a dramatic change in the physical spectrum of the theory. Instead of quarks and gluons, only colorless hadrons are produced as asymptotic states in any process, even at arbitrarily large energy.
Many nontrivial theoretical techniques respecting QCD first principles have been developed to study nonperturbative features of the theory. Yet despite numerous advances, no one has been able to compute the masses, wavefunctions, or transition amplitudes of hadrons in terms of quark masses and couplings directly from the QCD Lagrangian. Moreover, many existing theoretical tools are expressed through various expansions in certain small parameters; the actual range of each parameter where the expansions are applicable is often not well known. In such a situation, it is clearly advantageous to build a soluble toy field theory that incorporates as many features of the QCD Lagrangian as possible.
Such a theory does indeed exist, the famous ’t Hooft model , which is defined by the Yang-Mills Lagrangian in $`D=2`$ spacetime dimensions in the limit of a large number of colors $`N_c`$. As was shown in the original paper, the quark-antiquark sector of the theory admits an infinite tower of confined, color-singlet solutions that can be obtained, in principle, to an arbitrary degree of numerical accuracy. The reason for this solubility lies precisely in the defining features of the model. Large $`N_c`$ eliminates all Feynman diagrams with internal $`q\overline{q}`$ loops and nonplanar gluons. On the other hand, $`D=2`$ allows gluon self-couplings to be eliminated by gauging away one component of the gauge potential $`A^\mu `$. Since only two components are initially present, the commutator term $`[A^\mu ,A^\nu ]`$ in the covariant derivative, which gives gluon self-coupling, vanishes identically in such gauges. Then the only remaining Feynman diagrams to be summed for the quark-antiquark Green function are “rainbow” and “ladder” diagrams, whose Schwinger-Dyson equations can be solved, giving rise to an integral expression called the ’t Hooft equation (discussed in Sec. 2).
The ’t Hooft model provides an excellent laboratory for testing various approaches to strong interaction physics. After all, the ’t Hooft equation provides a means to compute hadronic masses, wavefunctions, and transition amplitudes in terms of the underlying partonic degrees of freedom.
In this work we are specifically interested in questions of local quark-hadron duality in the inclusive decays of heavy quarks. The notion of duality in general terms was first introduced in the early days of QCD in Ref. but not pursued for quite some time. A more detailed consideration was given a few years ago by Shifman and later reiterated in a number of papers (see, e.g., Refs. ), with applications relevant to Minkowskian observables amenable to study via an operator product expansion (OPE). This allows the formulation of the concept of local duality in a more quantitative way, including nontrivial nonperturbative effects; we refer the reader to these recent publications for the theoretical aspects. Here the question of duality is studied concretely by comparing the weak decay width of a meson containing a heavy quark computed in two ways. In terms of partonic degrees of freedom, one has an OPE depending upon the free quark diagram (with perturbative corrections) and a number of nonperturbative matrix elements suppressed by powers of the heavy quark mass. In terms of the hadronic degrees of freedom, one simply computes the weak decay amplitude for each allowed exclusive channel, and adds them up one by one. This comparison is especially instructive since one may consider the behavior of solution as the mass $`m_Q`$ of the heavy decaying quark is varied.
Such a problem was first considered in Ref. , where the main elements in numerical computations of exclusive decay rates were annunciated. The hadronic result was compared to the Born-level free-partonic diagram as a function of $`m_Q`$. In terms of the OPE, the latter is the tree-level piece of the Wilson coefficient corresponding to the unit operator. The numerical agreement was seen to be remarkable, in that the onset of the asymptotic agreement was clearly visible already for relatively small values of $`m_Q`$. The intrinsically limited numerical accuracy for sufficiently heavy quarks, however, prohibited drawing a definite conclusion about the size of nonperturbative corrections for asymptotically large $`m_Q`$. Additional numerical studies considered similar questions for weak decay topologies other than the simple spectator tree diagram, in particular weak annihilation (WA).
The validity of the OPE was addressed analytically in Refs. , which considered on one hand the nature of the OPE for heavy quark decays, and on the other an explicit $`1/m_Q`$ expansion of the decay amplitudes, which allows an analytical summation of the individual decay rates in the asymptotic regime. The agreement of the two approaches through relative order $`1/m_Q^4`$ was obtained by means of a number of sum rules derived directly from the ’t Hooft equation, the archetype of which first appeared in Ref. .
While adequate to illustrate the theoretical validity of the OPE for the inclusive decay widths of heavy flavors, the analytic methods per se cannot help in answering the practical question relevant to phenomenology of beauty and charm quarks: Namely, how accurately do the OPE-improved parton computations describe the true weak decay width of a heavy flavor meson with finite mass, only a few times larger than the typical strong interaction scale? A purely analytic expansion can hardly be used for this purpose, since it is a priori unknown how small an expansion parameter must be for the expansion to start yielding a reasonable approximation, not to mention achieving the necessary precision. To obtain insights into the size of deviations between the actual decay widths and the expressions obtained from the OPE for quarks in the intermediate mass range, one must employ real numerical computations.
In this paper we focus on semileptonic decays of heavy quarks. In the contexts of both real QCD and the ’t Hooft model, they are technically simpler than nonleptonic decays. Moreover, the magnitude of local duality violation is phenomenologically most important in semileptonic decays when one extracts $`|V_{cb}|`$ and $`|V_{ub}|`$. We use the techniques developed in Ref. to evaluate the required decay rates, and confront the total decay width with the expansion in terms of a power series in $`1/m_Q`$ of Ref. . Moreover, by making use of a number of relations derived in the large-$`m_Q`$ limit of the ’t Hooft equation , members of the set of nonperturbative matrix elements involved can be related to each other, providing an economical description of the nonperturbative physics. These are the tools that allow us to study the onset of quark-hadron duality.
As explained in Appendix A, we use a scheme based on the modified Multhopp method, by which the ’t Hooft equation is converted into an infinite-dimension eigenvector system that for practical reasons must be truncated at some number $`N`$ of eigenvector modes. The asymptotic convergence of this approach has not been rigorously studied, although it apparently must yield unlimited accuracy when the number of the Multhopp modes $`N`$ goes to infinity. Yet the rate of convergence at large $`N`$ is not well known. Additionally, large quark masses turn out to require one to use a larger $`N`$ for sufficient numerical accuracy, as discussed in Sec. 2. It therefore seems mandatory to make an independent cross-check of the numerical accuracy. We investigate this problem by comparing the numerical values of a number of static properties of heavy mesons at different values of $`m_Q`$, with the results of their $`1/m_Q`$ expansions obtained analytically from the ’t Hooft equation; this is the topic of Sec. 3. We find that our solutions have sufficient numerical accuracy for masses $`m_Q`$ corresponding to physical values (in the sense explained in Sec. 3) as large as $`20`$ GeV.
The duality of the inclusive widths of heavy-flavor hadrons to the parton-level widths, including the power corrections from the OPE, emerges through a set of sum rules that equate sums of weighted transition probabilities to possible final states and expectation values of the local heavy quark operators. Since our main interest lies in $`bc`$ transitions, which carry in practice a limited energy release, the most relevant are the so-called small velocity (SV) sum rules, which we study here in the heavy quark limit. The behavior of these sum rules not only shapes the semileptonic $`bc`$ decays in actual QCD, but is also important for the determination of the basic parameters of the heavy quark expansion.
An additional advantage of the heavy quark limit for our investigation is that we are able to compute the SV amplitudes semi-analytically, using the exact relations derived from the ’t Hooft equations and relying for input only on a few static parameters, which can be computed with a high precision. A discussion of these relations appears in Sec. 4. We find that the SV sum rules in the ’t Hooft model are saturated to an unexpectedly high degree by the first excitation above the ground state (which we henceforth call the “$`P`$-wave” excitation, despite the fact that in $`D=2`$ only radial excitations occur). Its contributions to even the Darwin ($`\rho _D^3`$) and kinetic ($`\mu _\pi ^2`$) expectation values constitute over $`90\%`$ and $`96\%`$ of the totals, respectively, while it saturates the “optical” sum rule for $`M_Bm_b`$ to a $`1.5\%`$ accuracy. This appears to be an intriguing dynamical feature of the model. A similar high-saturation effect has been observed in a quark flux-tube model , for the contribution from the “valence” quarkonium states.
We study the size of violations of local duality in the semileptonic decays $`bc\mathrm{}\overline{\nu }`$ assuming vectorlike weak currents and massless leptons. These assumptions are important for comparison with QCD far beyond the obvious parallel of closely resembling the actual world: The strength of the resonance-related duality violation crucially depends on the threshold behavior in the decay probabilities, which is completely different in two and four dimensions. The two-body phase space, while $`|\stackrel{}{p}|`$ in $`D=4`$, is $`1/|\stackrel{}{p}|`$, that is, infinite at threshold, in $`D=2`$. On the other hand, the situation is special for massless leptons: Their invariant mass is always zero if they are produced by a vectorlike source, and the weak vertex is then proportional to the momentum. As a result, in this case the threshold behavior of the decay rate becomes $`|\stackrel{}{p}|`$ much in the same way as in real QCD. This is a crucial detail if one tries to draw practical lessons from the ’t Hooft model. The need for a vectorlike coupling in $`D=2`$ is even more stark for the parton-level calculation. There one finds that the integrated three-body phase space actually diverges for massless leptons, and only the behavior of the weak decay amplitude renders the width finite. We provide more arguments in favor of such a choice in Sec. 5, which is dedicated to the inclusive decay widths.
In Sec. 6 we briefly illustrate how well the duality works for the vacuum correlator of light quarks in the timelike domain. In the context of the heavy quark expansion this is relevant for the nonleptonic decay widths, including spectator-dependent effects like WA.
Section 7 summarizes our investigation and discusses the conclusions that can be drawn for actual QCD.
Appendices describe the computational technique employed and contain a number of relations for the heavy quark limit of the ’t Hooft equation employed in these numerical studies.
## 2 The ’t Hooft Equation and Its Solutions
We first review some well-known properties of the ’t Hooft model both as a reminder and to establish notation. Confinement is manifest in 1+1 spacetime dimensions with large $`N_c`$, and the quark($`m_1`$)-antiquark($`\overline{m}_2`$) two-particle irreducible Green function, i.e., the meson wavefunction $`\phi (x)`$, is given by the ’t Hooft equation:
$$M_n^2\phi _n(x)=\left(\frac{m_1^2\beta ^2}{x}+\frac{m_2^2\beta ^2}{1x}\right)\phi _n(x)\beta ^2_0^1dy\phi _n(y)\mathrm{P}\frac{1}{(yx)^2},$$
(1)
where $`x`$ is the momentum fraction in light-cone coordinates carried by the quark, and
$$\beta ^2\frac{g_s^2}{2\pi }(N_c1/N_c).$$
(2)
Since $`\beta `$ is finite in the large-$`N_c`$ limit, it provides a natural unit of mass. Thus, all masses in this paper are understood as multiples of $`\beta `$. Indeed, as pointed out in Ref. , $`\beta `$ fills the role in 1+1 dimensions of served by $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ in 3+1. We discuss the estimation of $`\beta `$ as a particular number in Sec. 3.
The singularity of the QCD Coulomb interaction in Eq. (1) is regularized using a principal value prescription, indicated by P in Eq. (1).
Solutions $`n=0,1,\mathrm{}`$ of the ’t Hooft equation alternate in parity, with the lowest being a pseudoscalar. The general analytic solution in closed form is not known. As the eigenvalue index $`n`$ increases, the eigenvalues $`M_n^2`$ asymptotically approach $`\beta ^2[\pi ^2n+O(\mathrm{ln}n)]`$.
The static limit $`m_1m_Q\mathrm{}`$ is most easily studied by employing the “nonrelativistic” variables $`M_n=m_Q+ϵ_n`$, $`t=(1x)m_Q`$ and $`\mathrm{\Psi }_n(t)=\frac{1}{\sqrt{m_Q}}\phi _n\left(1\frac{t}{m_Q}\right)`$, in terms of which Eq. (1) assumes the form
$$ϵ_n\mathrm{\Psi }_n(t)=\frac{m_2^2\beta ^2}{2t}\mathrm{\Psi }_n(t)+\frac{t}{2}\mathrm{\Psi }_n(t)\frac{\beta ^2}{2}_0^{\mathrm{}}ds\frac{\mathrm{\Psi }_n(s)}{(ts)^2}.$$
(3)
We solve the finite-mass ’t Hooft equation using a numerical method called the Multhopp technique , a venerable system for solving integral equations with singular kernels. It was first applied to the ’t Hooft equation in Ref. . The idea is to expand the wavefunction in a series of modes, not unlike Fourier analysis, and then turn the equations for the mode coefficients into an equivalent infinite-dimension eigenvector problem. In practice, one then truncates at some point where the higher modes are deemed to have little effect upon the wavefunction solutions, which is of course strongly dependent on the highest value of $`n`$ used. The detailed formulas for applying the standard Multhopp technique to mesons with unequal quark masses in the ’t Hooft model appear in Appendix A of Ref. .
Intrinsic to the original Multhopp technique is the evaluation of the wavefunction at a discrete set of points called “Multhopp angles,” which in the current problem are equivalent to
$$x_k=\frac{1}{2}\left[1+\mathrm{cos}\left(\frac{k\pi }{N+1}\right)\right],k=1,\mathrm{},N,$$
(4)
where $`N`$ is the number of modes retained in the numerical solution. The mode coefficients are then obtained by the use of a discrete inversion formula \[(A7) in \]. However, the Multhopp solutions can be seen to vanish as $`\sqrt{x}`$ and $`\sqrt{1x}`$ at the endpoints $`x=0`$ and $`x=1`$, respectively \[see (A10)–(A11) in \], while the exact solutions are known to vanish as $`x^{\gamma _1}`$ and $`(1x)^{\gamma _2}`$, respectively, where
$$\frac{m_i^2}{\beta ^2}+\pi \gamma _i\mathrm{cot}\pi \gamma _i=\mathrm{\hspace{0.25em}1},$$
(5)
leading to a type of Gibbs phenomenon in the Multhopp solutions. Since the Multhopp angles cease to sample the wavefunction at some finite distance from the endpoints, it may be expected that the wavefunctions thus obtained are numerically inaccurate there. This shortcoming led Brower, Spence, and Weis to improve the Multhopp technique by eliminating the Multhopp angles and using instead a continuous inversion formula. The algebraic details are presented in Appendix A, and it is this improved numerical technique that is used in obtaining our results.
## 3 Heavy Quark Expansion and Cross Check of the Algorithm
Let us first establish a bit of notation. The mass of a heavy quark of flavor $`Q`$ is labeled as $`m_1m_Q`$; in the weak transitions considered in subsequent sections, the final-state quark $`q`$ is assigned the mass $`m_q`$. The spectator antiquark mass $`m_2`$ is labeled by $`m`$, or $`m_{\mathrm{sp}}`$ if there is any chance of confusion.
As explained in the previous section and Appendix A, we use the modified Multhopp technique to find numerical solutions of the ’t Hooft eigenstate problem. Since the heavy meson wavefunctions are peaked near the end of the interval, the accuracy deteriorates with increasing $`m_Q`$. The same, in principle, applies to the high excitations of light hadrons. A more appropriate strategy for heavy quarks is to start with a solution of the infinite-mass (static) equation. This has been done analytically , and full consistency<sup>3</sup><sup>3</sup>3In the case that the fermions $`f`$ created by the weak current have $`m_f=0`$, this agreement was shown up to and including $`O(1/m_Q^4)`$ terms in the weak decay width in , while terms up to and including $`O(m_f^2/m_Q)`$ were shown to coincide with those in the OPE in . In the current work we take $`m_f=0`$. with the OPE was demonstrated.
However, our practical interest lies in the properties of heavy hadrons with $`m_Q`$ lying in the intermediate domain, specifically for $`m_Q`$ one order of magnitude larger than $`\beta `$. The convergence of the $`1/m_Q`$ expansion in this case is too difficult to quantify analytically. This is just the situation where the numerical computations are best employed.
Therefore, an important element of the analysis is to check the accuracy of the numerical computations of both the heavy hadron masses and wavefunctions at different values of $`m_Q`$. To this end, we compute the masses and certain moments for the ground and first excited states, and compare them to the analytic $`1/m_Q`$ expansion. In general, the terms in the $`1/m_Q`$ expansion depend on a number of expectation values in the static limit, like the kinetic one $`\mu _\pi ^2=\overline{Q}(i\stackrel{}{D})^2Q`$, etc. However, one can show that the parameters appearing here through high order in $`1/m_Q`$ can be expressed in terms of just the asymptotic value $`\overline{\mathrm{\Lambda }}=M_{H_Q}m_Q`$ and the corresponding decay constant. These quantities are the ones most accessible to numerical evaluation; in particular, $`\overline{\mathrm{\Lambda }}`$ is expected to be the most accurately determined quantity.
Our main computations refer to the case of $`m_{\mathrm{sp}}=0.56\beta `$, as chosen in Ref. . It corresponds (see Sec. 4) to a mass of the strange quark in QCD. The choice of a noticeable light quark mass may be motivated by an attempt to mimic the effect of the transverse gluons absent in $`D=2`$, which in a certain respect supply some effective mass to the light quark. Clearly, this can be only a rather crude approximation, since the bare quark mass breaks chiral invariance. One can suppose, nevertheless, that this side effect is not too important for our purposes. The chiral symmetry is spontaneously broken anyway, and the presence of a massless versus a massive pion does not seem to be essential for the range of problems we address here. On the other hand, the effect of the transverse degrees of freedom is known to soften the $`x1`$ singularity of the heavy quark distribution function , similar to the impact of the light quark mass in the ’t Hooft model. The behavior of the distribution function affects the inclusive decays of the heavy quarks in an essential way.
We also present some results for $`m_{\mathrm{sp}}=0.26\beta `$, partly to explore light quark dependences of matrix elements and partly to investigate the beginnings of failure of the numerical solutions as $`m_{\mathrm{sp}}0`$. The number $`N`$ of Multhopp modes used is $`500`$; we considered smaller $`N`$ as well to study this dependence, but since the behavior was found to be stable, we do not dwell on it further here.
The masses of heavy hadrons obey
$$M_{H_Q}m_Q=\overline{\mathrm{\Lambda }}+\frac{\mu _\pi ^2\beta ^2}{2m_Q}+\frac{\rho _D^3\rho _{\pi \pi }^3}{4m_Q^2}+O\left(\frac{\beta ^4}{m_Q^3}\right).$$
(6)
where
$`\mu _\pi ^2`$ $`=`$ $`\overline{Q}(i\stackrel{}{D})^2Q,\rho _D^3={\displaystyle \frac{1}{2}}\overline{Q}(\stackrel{}{D}\stackrel{}{E})Q,`$
$`\rho _{\pi \pi }^3`$ $`=`$ $`{\displaystyle \frac{1}{2}}iT\{\overline{Q}(i\stackrel{}{D})^2Q(x),\overline{Q}(i\stackrel{}{D})^2Q(0)\}_{q=0},`$ (7)
and these expectation values refer to the infinite mass limit. In the $`\rho _{\pi \pi }^3`$ expression, $`q`$ is the momentum variable conjugate to $`x`$, and diagonal transitions within the correlator have been removed.
Since $`\overline{\mathrm{\Lambda }}`$ in QCD traditionally denotes the mass difference between a ground-state pseudoscalar meson and its corresponding heavy quark in the large $`m_Q`$ limit \[as it is defined in Eq. (6)\], and we need it for a number of the excited states $`H_Q^{(n)}`$ as well, we assign the notation
$$\overline{\mathrm{\Lambda }}^{(n)}ϵ^{(n)},$$
(8)
and use $`ϵ`$ and $`\overline{\mathrm{\Lambda }}`$ throughout the paper on equal footing. Equations (6)–(3) hold for each state $`H_Q^{(n)}`$ with $`n=0,1,\mathrm{}`$, so that an implicit superscript $`(n)`$ is to be understood in these expressions.
According to Ref. , the following relations hold in the ’t Hooft model:
$$\mu _\pi ^2=\frac{\overline{\mathrm{\Lambda }}^2m^2+\beta ^2}{3},\rho _D^3=\frac{\beta ^2F^2}{4},\rho _{\pi \pi }^3=\frac{1}{36}\left[8\overline{\mathrm{\Lambda }}(\overline{\mathrm{\Lambda }}^2m^2+\beta ^2)+3\beta ^2F^2\right].$$
(9)
Here $`F`$ is the scaled decay constant in the heavy quark limit, i.e.,
$$F^{(n)}=_0^{\mathrm{}}dt\mathrm{\Psi }_n(t)=\underset{m_Q\mathrm{}}{lim}_0^{m_Q}dt\frac{1}{\sqrt{m_Q}}\varphi _n\left(1\frac{t}{m_Q}\right)=\underset{m_Q\mathrm{}}{lim}c_n\sqrt{m_Q},$$
(10)
where the superscript is suppressed if there is no ambiguity,
$$c_n=_0^1dx\phi _n(x),$$
(11)
and the exact relation between $`c_n`$ and the decay constant of the $`n`$th excitation is given in Eq. (68). In the heavy quark limit one has
$$\overline{\mathrm{\Lambda }}=m_Q1x,\mu _\pi ^2=m_Q^2\left(x^2x^2\right),$$
(12)
but there are $`O(1/m_Q)`$ corrections to these relations. For further applications to the decay widths we also consider the scalar expectation value
$$\frac{1}{2M_{H_Q}}\overline{Q}Q=\frac{m_Q}{M_{H_Q}}\frac{1}{x}.$$
(13)
Then the following expansions hold:
$`\sqrt{m_Q}c_n`$ $`=`$ $`\left(1{\displaystyle \frac{2[2ϵ^{(n)}m(1)^n]}{3m_Q}}\right)F^{(n)}+O\left({\displaystyle \frac{\beta ^{5/2}}{m_Q^2}}\right),`$
$`m_Q1x`$ $`=`$ $`\overline{\mathrm{\Lambda }}{\displaystyle \frac{\overline{\mathrm{\Lambda }}^2+\mu _\pi ^2}{m_Q}}+{\displaystyle \frac{4\overline{\mathrm{\Lambda }}(6\overline{\mathrm{\Lambda }}^2+8\mu _\pi ^2+3\beta ^2)+\beta ^2F^2}{24m_Q^2}}+O\left({\displaystyle \frac{\beta ^4}{m_Q^3}}\right),`$
$`m_Q^2\left(x^2x^2\right)`$ $`=`$ $`\mu _\pi ^2{\displaystyle \frac{1}{3m_Q}}\left(8\overline{\mathrm{\Lambda }}\mu _\pi ^2+\beta ^2F^2\right)+O\left({\displaystyle \frac{\beta ^4}{m_Q^2}}\right),`$
$`{\displaystyle \frac{m_Q}{M_{H_Q}}}{\displaystyle \frac{1}{x}}`$ $`=`$ $`1{\displaystyle \frac{\mu _\pi ^2\beta ^2}{2m_Q^2}}{\displaystyle \frac{\rho _D^3\rho _{\pi \pi }^3}{2m_Q^3}}+O\left({\displaystyle \frac{\beta ^4}{m_Q^4}}\right).`$ (14)
We note that values of $`\mu _\pi ^2`$, $`\rho _D^3`$, or $`\rho _{\pi \pi }^3`$ determined from the expansions Eqs. (3) suffer degraded numerical accuracy compared to those taken directly from Eqs. (9) since $`\overline{\mathrm{\Lambda }}`$ and $`F`$ are determined from more stable expansions (in particular, they do not depend upon close numerical cancellations). Therefore, we use Eqs. (9) as primary information and relegate Eqs. (3) to numerical checks. Our method of determining $`\overline{\mathrm{\Lambda }}`$ from the $`M_{H_Q}m_Q`$ expression, designed to minimize the influence of potentially large uncertainties at large $`m_Q`$, is described in Appendix C.
Values of $`M_{H_Q}m_Q`$ and and the averages in Eqs. (3) as functions of $`m_Q`$ from $`m=0.56\beta `$ to $`50\beta `$ are presented in Table 1 for both the ground and first excited states. Similar results for just the ground state with $`m=0.26\beta `$ are presented in Table 2. Based upon the 10 data points presented in Table 1 for the ground state, one may fit to a polynomial in $`1/m_Q`$, obtaining
$`{\displaystyle \frac{1}{\beta }}(M_{H_Q}m_Q)`$ $`=`$ $`1.3170.086{\displaystyle \frac{\beta }{m_Q}}0.050{\displaystyle \frac{\beta ^2}{m_Q^2}}+O\left({\displaystyle \frac{\beta ^3}{m_Q^3}}\right),`$
$`c_0\sqrt{{\displaystyle \frac{m_Q}{\beta }}}`$ $`=`$ $`2.0322.775{\displaystyle \frac{\beta }{m_Q}}+O\left({\displaystyle \frac{\beta ^2}{m_Q^2}}\right),`$
$`{\displaystyle \frac{m_Q}{\beta }}1x`$ $`=`$ $`1.3162.491{\displaystyle \frac{\beta }{m_Q}}+3.789{\displaystyle \frac{\beta ^2}{m_Q^2}}+O\left({\displaystyle \frac{\beta ^3}{m_Q^3}}\right),`$
$`{\displaystyle \frac{m_Q^2}{\beta ^2}}\left(x^2x^2\right)`$ $`=`$ $`0.80744.050{\displaystyle \frac{\beta }{m_Q}}+O\left({\displaystyle \frac{\beta ^2}{m_Q^2}}\right),`$
$`{\displaystyle \frac{m_Q}{M_{H_Q}}}{\displaystyle \frac{1}{x}}`$ $`=`$ $`1+0.099{\displaystyle \frac{\beta ^2}{m_Q^2}}0.044{\displaystyle \frac{\beta ^3}{m_Q^3}}+O\left({\displaystyle \frac{\beta ^4}{m_Q^4}}\right).`$ (15)
The corresponding expressions using the approach of Appendix C (neglecting the one for $`M_{H_Q}m_Q`$, which is used as input and hence is identical through $`O(\beta ^2/m_Q^2)`$) read
$`c_0\sqrt{{\displaystyle \frac{m_Q}{\beta }}}`$ $`=`$ $`2.0352.816{\displaystyle \frac{\beta }{m_Q}}+O\left({\displaystyle \frac{\beta ^2}{m_Q^2}}\right),`$
$`{\displaystyle \frac{m_Q}{\beta }}1x`$ $`=`$ $`1.3182.544{\displaystyle \frac{\beta }{m_Q}}+4.512{\displaystyle \frac{\beta ^2}{m_Q^2}}+O\left({\displaystyle \frac{\beta ^3}{m_Q^3}}\right),`$
$`{\displaystyle \frac{m_Q^2}{\beta ^2}}\left(x^2x^2\right)`$ $`=`$ $`0.80783.996{\displaystyle \frac{\beta }{m_Q}}+O\left({\displaystyle \frac{\beta ^2}{m_Q^2}}\right),`$
$`{\displaystyle \frac{m_Q}{M_{H_Q}}}{\displaystyle \frac{1}{x}}`$ $`=`$ $`1+0.096{\displaystyle \frac{\beta ^2}{m_Q^2}}0.066{\displaystyle \frac{\beta ^3}{m_Q^3}}+O\left({\displaystyle \frac{\beta ^4}{m_Q^4}}\right).`$ (16)
This agreement between the two approaches is quite excellent and is exhibited in Figs. 14 for $`M_{H_Q}m_Q`$ and the quantities in Eqs. (3); in general, the exact results are presented as points on a solid line, while each fit using Eqs. (3) is presented as a dashed line. In Fig. 6 the analogous expression $`M_{H_Q}m_Q`$ for the $`m=0.56\beta `$ first excited state is presented, while Fig. 7 uses the same methods and values from Table 2 to present $`M_{H_Q}m_Q`$ for the $`m=0.26\beta `$ ground state. In Fig. 1 and especially in Fig. 7, the quality of numerical results is seen (as expected) to begin breaking down at large $`m_Q`$ and small $`m`$, since $`N=500`$ is fixed. We conclude that the numerical routine we rely upon is sufficiently accurate for $`N=500`$ up to $`m_Q(25÷30)\beta `$. The critical value of $`m_Q`$ also depends, however, on the meson’s light quark mass, decreasing for small $`m`$. This is expected since at small $`m`$ the sharpness of the wavefunction as $`x1`$ becomes stronger, and more Multhopp functions are required to approximate it: Each Multhopp function vanishes like $`\sqrt{1x}`$. Likewise, the required $`N`$ increases for the excited states. Still, one can check that it is possible to go as high as $`m_Q=15\beta `$ even for $`m`$ as small as $`0.1\beta `$.
It turns out that a numerically significant cancellation occurs in the value of $`\mu _\pi ^2\beta ^2`$ in $`1/m_Q^2`$ corrections and, in particular, at the $`1/m_Q^3`$ level between $`\rho _D^3`$ and $`\rho _{\pi \pi }^3`$, for the ground state just around our primary value $`m=0.56\beta `$. Such a numerical suppression of the power corrections is accidental and does not occur for the excited states, nor for $`m=0.26\beta `$.
Let us note that the expectation value of the light-quark scalar density in the heavy meson turns out very close to unity for the ground state, which may be seen by taking $`m_Qm`$ and $`x1x`$ in Eq. (13) and referring to Table 3; this is a characteristic feature of a nonrelativistic (with respect to the light quark) bound-state system. It implies an almost simple additive dependence of $`\overline{\mathrm{\Lambda }}`$ on the light quark mass $`m`$,
$$\overline{\mathrm{\Lambda }}\overline{\mathrm{\Lambda }}_{|m=0}+m,$$
(17)
and indeed one can verify this feature by comparing $`M_{H_Q}m_Q`$ values between Table 1 ($`m=0.56\beta `$) and Table 2 ($`m=0.26\beta `$). While this pattern is expected when the spectator quark is heavy, it a priori needs not hold when it is light. This supports the naive expectation that the chiral symmetry breaking may lead to a description in some aspects resembling the nonrelativistic constituent quark model. The above expectation value, however, decreases for the excited states, as expected from such a picture.
Drawing semi-quantitative conclusions for QCD requires a translation rule between the mass parameters in the two theories, that is, an estimate of the value of $`\beta `$ in GeV. Different dimensionful quantities can be taken as the yardstick; since the theories are not identical, this translation rule must be introduced with some care. As follows from the heavy quark sum rules, the physics of duality in the decay widths of heavy flavors crucially depends on the properties of the lowest excited heavy-quark states, in particular the $`P`$-wave excitations with opposite parity to the ground-state multiplet. It will become evident from the next section that they are of primary importance for the $`1/m_Q`$ expansion of static properties as well. Therefore, we choose the mass difference between the lowest parity-even ($`P`$-wave) state and the parity-odd ground-state meson to gauge the translation between the mass scales.
In the ’t Hooft model the mass difference $`ϵ_1ϵ_0`$ for light spectators amounts to about $`1.3\beta `$. Real charm spectroscopy suggests that the first $`P`$-wave excitations are between $`400`$ and $`500\text{MeV}`$ above the ground state. Taking the larger value for sake of illustration, we arrive at the estimate
$$\beta \mathrm{\hspace{0.33em}400}\text{MeV},$$
(18)
which is adopted in our analysis. This falls rather close to the estimate of Ref. , which relied on a quite different type of effects in the light-quark systems.
Assuming a value for the “bare” $`b`$ quark mass in QCD (normalized at the appropriate scale $`\text{ }>m_b`$) of about $`4\text{GeV}`$, we conclude that mesons with quarks of masses
$$m_b10\beta ,m_c(2.5÷3.5)\beta ,$$
(19)
represent in the ’t Hooft model the actual beauty and charm mesons.
The value of $`\overline{\mathrm{\Lambda }}1.3\beta 500\text{MeV}`$ seems to be in a reasonable correspondence with the size of this difference in QCD when it is normalized at a low hadronic scale, $`\overline{\mathrm{\Lambda }}(1\text{GeV})(600\pm 60)\text{MeV}`$ .
It should be noted, however, that the kinetic expectation value in the ’t Hooft model turns out to be rather small, $`\mu _\pi ^20.8\beta ^20.12\text{GeV}^2`$. This is not surprising, since the chromomagnetic field is absent in two dimensions, while it was shown to be crucial in the real case. Indeed, the comparison is better justified for the difference $`\mu _\pi ^2\mu _G^2`$ in actual QCD versus the value of $`\mu _\pi ^2`$ in the ’t Hooft model. These questions were discussed in detail in Ref. , and can be easily understood using the sum rule representation. Due to the absence of spin in two dimensions, there is no difference between the would-be spin-$`1/2`$ and spin-$`3/2`$ light degrees of freedom. In particular, labeling the “oscillator strengths” $`\tau `$ defined in the next section by spin rather than excitation number, $`\tau _{1/2}=\tau _{3/2}`$ and $`ϵ_{1/2}=ϵ_{3/2}`$ effectively hold. Then the sum
$$\mu _\pi ^2=\mathrm{\hspace{0.33em}3}\underset{n}{}ϵ_n^2|\tau _{1/2}|^2+\mathrm{\hspace{0.25em}6}\underset{n}{}ϵ_n^2|\tau _{3/2}|^2\mathrm{\hspace{0.33em}9}\underset{n}{}ϵ_n^2|\tau _{1/2}|^2,$$
(20)
and the latter sum is just the general expression for $`\mu _\pi ^2\mu _G^2`$:
$$\mu _\pi ^2\mu _G^2=\mathrm{\hspace{0.33em}9}\underset{n}{}ϵ_n^2|\tau _{1/2}|^2.$$
(21)
Accepting such an identification suggested in Ref. and the estimate $`\mu _\pi ^2\mu _G^2(0.15\pm 0.1)\text{GeV}^2`$, we again observe a reasonable agreement with the findings of the ’t Hooft model.
## 4 Duality in the SV Sum Rules
A useful theoretical limit—the so-called small velocity (SV) regime—was suggested in the mid 80’s as a theoretical tool for studying semileptonic heavy quark decays. This refers to kinematics where both $`b`$ and $`c`$ quarks are heavy, but the energy release is limited, so that the velocity of the final charm hadron is small. At large energy release the OPE for the width must converge rapidly to the actual hadronic width. Still, at fixed energy release the deviations, although $`1/m_Q`$ suppressed, are present regardless of the absolute values of masses.
In the SV regime the semileptonic decays proceed either to the ground-state charm final state, $`D`$ or $`D^{}`$ (the semi-elastic transitions), or to excited “$`P`$-wave” states of the opposite parity. Other decays are suppressed by higher powers of velocity, or by heavy quark masses.
The equality of the sum of partial decay widths and its OPE expansion is achieved through the sum rules that relate the sums of the $`P`$-wave transition probabilities, weighted with powers of the excitation energies, to the static characteristics of the decaying heavy hadron. The onset of convergence of the OPE expansion for the widths is then directly related to the pattern of saturation of the sum rules by the lowest excitations. If higher states contribute significantly, they delay the onset of duality, while their absence leads to a tight quark-hadron duality after the first $`P`$-wave channel is open.
Knowledge of degree of saturation of the heavy quark sum rules is also important for another reason: It determines the hadronic scale above which one can apply the perturbative treatment to compute corrections or account for evolution of the effective operators. The lower this scale, the more predictive in turn is the treatment of the nonperturbative effects in the OPE.
A recent review of the SV sum rules can be found in Ref. (the perturbative aspects are discussed in more detail in Ref. ). For most practical purposes addressed here, one can consider the perturbative effects to be absent in the ’t Hooft model. In particular, the heavy quark parameters do not depend perturbatively on the normalization point, and there is no need in the explicit ultraviolet cutoff to introduce a normalization point. The sum rules we address are
$`\rho _k^2{\displaystyle \frac{1}{4}}`$ $`=`$ $`{\displaystyle \underset{n}{}}\tau _{nk}^2,`$ (22)
$`{\displaystyle \frac{1}{2}}\overline{\mathrm{\Lambda }}_k`$ $`=`$ $`{\displaystyle \underset{n}{}}(ϵ_nϵ_k)\tau _{nk}^2,`$ (23)
$`\left(\mu _\pi ^2\right)_k`$ $`=`$ $`{\displaystyle \underset{n}{}}(ϵ_nϵ_k)^2\tau _{nk}^2,`$ (24)
$`\left(\rho _D^3\right)_k`$ $`=`$ $`{\displaystyle \underset{n}{}}(ϵ_nϵ_k)^3\tau _{nk}^2.`$ (25)
Here $`k`$ and $`n`$ denote excitation indices for the initial and final states, respectively (in practice only transitions from the ground state are interesting, so we limit ourselves to $`k=0`$; in this case the index $`k`$ is omitted). The so-called “oscillator strengths” $`\tau `$ parameterize the transition amplitudes into the opposite-parity states in the SV limit,
$$\frac{1}{2m_Q}n\left|\overline{Q}\gamma _\mu Q\right|k=\tau _{nk}ϵ_{\mu \nu }v^\nu +O(\stackrel{}{v}^{\mathrm{\hspace{0.17em}3}}),$$
(26)
where $`\stackrel{}{v}`$ is the velocity of the final state hadron. In the diagonal transition $`\rho _k^2`$ is the slope of the Isgur-Wise (IW) function of state $`|k`$:
$$\frac{1}{2m_Q}k(\stackrel{}{v})\left|\overline{Q}\gamma _0Q\right|k(0)=\mathrm{\hspace{0.33em}1}\rho _k^2\frac{\stackrel{}{v}^{\mathrm{\hspace{0.17em}2}}}{2}+O(\stackrel{}{v}^{\mathrm{\hspace{0.17em}4}}).$$
(27)
The expressions for $`\tau _{nk}`$ and $`\rho _k`$ in terms of the light-cone wavefunctions are
$`\tau _{nk}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}dt\mathrm{\Psi }_n(t)t{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}\mathrm{\Psi }_k(t)=\underset{m_Q\mathrm{}}{lim}{\displaystyle _0^1}dx\phi _n(x)(1x){\displaystyle \frac{\mathrm{d}}{\mathrm{d}x}}\phi _k(x),`$
$`\rho _k^2`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}dt\left|\left(t{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}+{\displaystyle \frac{1}{2}}\right)\mathrm{\Psi }_k(t)\right|^2=\underset{m_Q\mathrm{}}{lim}{\displaystyle _0^1}dx\left|\left[(1x){\displaystyle \frac{\mathrm{d}}{\mathrm{d}x}}{\displaystyle \frac{1}{2}}\right]\phi _k(x)\right|^2.`$
The finite-$`m_Q`$ corrections to the $`x`$-integral forms turn out to be rather significant, leading to significant problems in precision numerical studies. To avoid this problem we use the analytic expression for the inelastic amplitudes obtained in Ref. :
$$\tau _{nk}=n\left|t\frac{\mathrm{d}}{\mathrm{d}t}\right|k=\frac{\beta ^2}{2(ϵ_nϵ_k)^3}F^{(n)}F^{(k)}\left(\frac{1(1)^{nk}}{2}\right),$$
(29)
where $`F^{(n)}`$ are the asymptotic values of the decay constants $`c_n`$ scaled up by the factor $`\sqrt{m_Q}`$, as in Eq. (10). The constants $`F^{(n)}`$ are computed as the values of $`c_n\sqrt{m_Q}`$ at $`m_Q=15\beta `$ (see Table 1) augmented by the $`1/m_Q`$ corrections detailed in the first of Eqs. (3), while values of $`ϵ_n`$ are computed using the procedure described in Appendix C.
The results of the computations for the case $`m_Q=15\beta `$, $`m=0.56\beta `$ are presented in Table 4. Our central result is a surprisingly good saturation of the sum rules: The first $`(n=1)`$ excitation generates $`99.4\%`$ of $`\rho ^2`$, $`98.5\%`$ of $`\overline{\mathrm{\Lambda }}`$, $`96\%`$ of $`\mu _\pi ^2`$, and even $`91\%`$ of $`\rho _D^3`$. The rest is almost completely saturated by the second $`P`$-wave state ($`n=3`$), where the cumulative values for the same quantities read $`99.92\%`$, $`99.73\%`$, $`99.1\%`$, and $`96.7\%`$, respectively.
In terms of absolute numbers, the sum rules Eqs. (22)–(25) would give $`\rho ^21/4=0.529`$, $`\overline{\mathrm{\Lambda }}=1.278\beta `$, $`\mu _\pi ^2=0.782\beta ^2`$, and $`\rho _D^3=0.99\beta ^3`$, the last of which gives $`F^{(0)}=1.99\sqrt{\beta }`$, in fine agreement with the values obtained from the values obtained in the previous section via the methods described in Appendix C. The few-percent discrepancy corresponds to the accuracy in determinations of squared decay constants.
The level of saturation by the lowest open channels is extraordinary. The explicit reason for such a perfect saturation of the sum rules involving even rather high, $`ϵ^3`$ powers of the excitation energy can be read off Eq. (29)—$`\tau `$’s are inversely proportional to the third power of the excitation energy. With the asymptotics $`ϵ_n\sqrt{n}`$, $`F^{(n)}n^{1/4}`$, the first excitation energy $`ϵ_1ϵ_0`$ is notably smaller than the next one $`ϵ_3ϵ_0`$ including three energy gaps. The general peculiarity of the ’t Hooft model leading to such a saturation is not understood completely.
With this pattern of saturation of the SV sum rules for the ground-state meson, one expects an early onset of the accurate duality for the inclusive widths in the $`bc`$ transitions, only slightly above the threshold of the first excitation. Demonstrating this result through direct evaluation of the decay widths is one of the purposes of the next section.
## 5 Local Duality in the Decay Widths
The semileptonic widths described in this work were considered in detail in Ref. . Here we recapitulate a few basic points. The weak decay Lagrangian is
$$_{\mathrm{weak}}=\frac{G}{\sqrt{2}}(\overline{c}\gamma _\mu b)(\overline{e}\gamma ^\mu \nu ).$$
(30)
In terms of the previous notation, $`Qb`$, $`qc`$ (or, later in this section, $`u`$), and $`H_QB`$. The key property of all $`D=2`$ vectorlike currents is that for $`m_e=m_\nu =0`$, the invariant mass $`q^2`$ of the lepton pair is always zero. For all computational purposes decays into this massless lepton pair are equivalent to decays into a single massless pseudoscalar particle $`\varphi `$ weakly coupled to quarks according to
$$\stackrel{~}{}_{\mathrm{weak}}=\frac{G}{\sqrt{2\pi }}\overline{c}\gamma _\mu bϵ^{\mu \nu }_\nu \varphi .$$
(31)
Several arguments favor our choice of a vectorlike weak decay interaction in the ’t Hooft model. One is of course the simplicity of Eq. (31). Another is that for $`q^2=0`$ some difficult problems of renormalization are absent, as we now discuss. The central problem in applying the OPE in practice is disentangling perturbative and nonperturbative effects. More precisely, this refers to the separation of short-distance effects attributed to the coefficient functions from long-distance effects residing in the matrix elements of the effective heavy-quark operators.
The perturbative corrections, for example those that renormalize the weak quark current, are generally rather nontrivial, even in the ’t Hooft model. However, according to the nonrenormalization theorem of Ref. , such vertex corrections are absent from the decays with $`q^2=0`$. This allows one to isolate the problem of renormalization of the underlying current from the question of interest in our study: possible deviations of the full decay widths due to the presence of thresholds in the production of the hadronic resonances.
In reality, from the OPE viewpoint some short-distance corrections still remain even in this special kinematic region due to the high-momentum tails in the meson wavefunctions. These tails come from the hard gluon exchanges between the constituents. In principle, these “hard” components can also be separated from the “soft” bound-state dynamics explicitly. However, in practice this is not necessary: These effects are completely contained in the meson wavefunctions.
Another advantage of vectorlike currents is apparent when one notes that the $`D=2`$ three-body “semileptonic” phase space diverges logarithmically for massless leptons. Explicitly, for the decay $`Mm+m_{\mathrm{}}+m_{\mathrm{}}`$ (equal lepton masses are assumed to render the expressions simpler), the three-body phase space turns out to be
$`\mathrm{\Phi }_3(M;m,m_{\mathrm{}},m_{\mathrm{}})`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^3(Mm)\sqrt{(M+m)^24m_{\mathrm{}}^2}}}`$ (32)
$`\times K\left[{\displaystyle \frac{(M+m)^2\left[(Mm)^24m_{\mathrm{}}^2\right]}{(Mm)^2\left[(M+m)^24m_{\mathrm{}}^2\right]}}\right],`$
where $`K`$ is the complete elliptic integral of the first kind. As $`m_{\mathrm{}}m`$, one regains Eqs. (4.1)–(4.3) of , while as $`m_{\mathrm{}}0`$, the argument of the elliptic integral goes to unity, and $`K(1ϵ)\mathrm{ln}(8/ϵ)/2`$. This is a manifestation of the logarithmic infrared divergence of the massless scalar Green function at large distance in $`D=2`$. A detailed calculation shows that the vector nature of the weak coupling regularizes the phase space integral, preventing the partonic rate from diverging in the limit of massless leptons. Furthermore, as discussed in the Introduction, this also removes the $`1/|\stackrel{}{p}|`$ singularity in the threshold behavior for hadronic two-body decays.
As a final advantage of vectorlike currents and the special kinematic point $`q^2=0`$, note that at $`q^2=0`$ the $`BD^{(n)}`$ transition amplitudes are directly expressed in terms of the overlap between the initial and the final wavefunctions:
$$q_\mu \frac{1}{2M_B}n|ϵ^{\mu \nu }J_\nu |B=q_z_0^1dx\phi _n(x)\phi _B(x),$$
(33)
where $`q_z=|\stackrel{}{p}|=(M_B^2M_n^2)/2M_B`$, so that the partial decay width for $`BD^{(n)}\mathrm{}\overline{\nu }`$ is given by
$$\mathrm{\Gamma }_n=\frac{G^2}{4\pi }\frac{M_B^2M_n^2}{M_B}\left|_0^1dx\phi _n(x)\phi _B(x)\right|^2\theta (M_BM_n).$$
(34)
The threshold suppression mentioned above is manifested in the explicit factor $`(M_B^2M_n^2)`$: The reciprocal of this factor in the phase space is removed by $`q_z^2`$ from the matrix element. It is also possible to derive this result directly using the methods of Ref. ; note, however, that these expressions are much simpler than those of Ref. , because the vectorlike current with massless leptons restricts $`q^2`$ to 0. The sum of these widths over all open channels is to be compared to the OPE prediction. The remarkable speed of saturation in $`n`$, anticipated in the last section, is illustrated for one sample case in Table 5.
Turning to the OPE, we mention one more problem associated with an accurate understanding of local duality violation. Apart from the purely theoretical aspect that OPE power series are generally only asymptotic and, thus have a formally zero radius of convergence in $`1/m_Q`$, one normally has additional practical limitations. Only a limited number of the terms, as well as the associated expectation values, are usually known, which places additional theoretical uncertainties that dominate in practice at sufficiently large $`m_Q`$.
This feature can be naturally incorporated in the analysis of our concrete model. We account completely only through terms that scale like $`1/m_Q^4`$, the highest order that emerges from the OPE free from the four-fermion operators . The rest, although calculable in principle term-by-term in the ’t Hooft model, are taken to represent the OPE “tails” discarded by the unavoidable truncation.
Using the sum rules of the ’t Hooft model, Ref. established the following exact representation for the total decay width:
$$\mathrm{\Gamma }_B=\frac{G^2}{4\pi }\frac{m_b^2m_c^2}{M_B}_0^1\frac{\mathrm{d}x}{x}\phi _B^2(x)\underset{M_n>M_B}{}\mathrm{\Gamma }_n,$$
(35)
where $`\mathrm{\Gamma }_n`$ at $`M_n>M_B`$ are understood as given by Eq. (34) without the explicit $`\theta `$-function singling out the open channels; such $`\mathrm{\Gamma }_n`$ are therefore all negative. On the other hand, the OPE yields the result
$$\mathrm{\Gamma }_B=\frac{G^2}{4\pi }\frac{m_b^2m_c^2}{m_b}\left[\frac{m_b}{M_B}_0^1\frac{\mathrm{d}x}{x}\phi _B^2(x)+O\left(\frac{\beta ^5}{M^5}\right)\right],$$
(36)
with $`M`$ generically denoting the OPE expansion parameter; we do not specify here if it is $`m_b`$ or $`m_bm_c`$, or some other combination. It was shown in Ref. that the $`\mathrm{\Gamma }_n`$ term in Eq. (35) is dual to the order term in Eq. (36); however we do not use this here and rather treat the latter as an intrinsic uncertainty in the “practical” version of the OPE.
Thus, our strategy is to compare the exact width
$$\mathrm{\Gamma }_B=\frac{G^2}{4\pi }\underset{M_n<M_B}{}\frac{M_B^2M_n^2}{M_B}\left|_0^1dx\phi _n(x)\phi _B(x)\right|^2,$$
(37)
to
$$\mathrm{\Gamma }_{\mathrm{OPE}}=\frac{G^2}{4\pi }\frac{m_b^2m_c^2}{m_b}\frac{m_b}{M_B}_0^1\frac{\mathrm{d}x}{x}\phi _B^2(x).$$
(38)
The expectation value $`\frac{m_b}{M_B}\frac{1}{x}`$ above can either be evaluated numerically, or in the spirit of the OPE, computed in the form of a $`1/m_b`$ expansion, the last of Eqs. (3). It turns out that the expansion converges very rapidly to the exact result, so that this does not significantly affect the observed pattern of local duality at the quantitative level. The Born-term partonic rate is simply given by $`\mathrm{\Gamma }_{\mathrm{OPE}}`$ with this expectation value set to unity,
$$\mathrm{\Gamma }_b=\frac{G^2}{4\pi }\frac{m_b^2m_c^2}{m_b}.$$
(39)
The main practical interest of these calculations lies in the $`bc`$ width with its limited energy release $`E_r`$. In general, $`E_r`$ can be small either if $`m_b`$ is not large enough, or even at large $`m_b`$ if $`E_r=m_bm_c`$ (or $`m_bm_c\sqrt{q^2}`$ if $`q^2`$ is nonzero) is insufficient due to a significant $`c`$ quark mass. The latter case falls into the SV category, and the violations of duality are suppressed here even at the maximal $`q^2`$ by heavy quark symmetry, as was pointed out in the mid-80’s . Therefore, one a priori expects a different pattern in the two cases. We try to separate the possible effects by considering different choices for $`m_b`$ and $`m_c`$ rather than by only taking them close to their realistic values.
With these arguments in mind, one can expect to find significant effects of duality violation in the cases where $`1/m_c`$ or $`1/m_b`$ effects are important. As suggested in Ref. , in this case it is advantageous to fix $`m_b`$ close to its actual value, and vary $`m_c`$ from near $`m_b`$ down to smaller values, changing in this way the energy release. At one end of the interval the local duality is supported by the heavy quark symmetry with large quark masses and SV kinematics, while at another end it rests on the large energy release.
We start from the SV case when $`m_b`$ is fixed and large and $`m_c`$ is large as well, varying the energy release by increasing $`m_c`$ towards $`m_b`$. Since the violation of local duality is expected to be suppressed for all values of $`m_c`$, high numerical accuracy is vital. We fix $`m_b=15\beta `$ ($`6\text{GeV}`$), and vary $`m_c`$ from $`5\beta `$ up to $`m_b`$. The results are given in Table 6 and Fig. 8. We note that the difference between the two widths is so small that one must plot $`\mathrm{ln}(\mathrm{\Gamma }_B/\mathrm{\Gamma }_{\mathrm{OPE}}1)`$ rather than the widths themselves. This is expected since the SV sum rules are very well saturated, as detailed in the previous section—the higher thresholds are then strongly suppressed numerically at finite energy release. But for $`m_c`$ approaching $`m_b`$, where they could be noticeable, the heavy quark symmetry works efficiently since both quarks are very heavy. In fact, the only prominent features on the plot occur when thresholds to the first few $`D`$ states of opposite parity to the ground-state $`B`$ meson are crossed, for example between $`m_c=13.5`$ and $`14\beta `$. The deviation is extremely small also for smaller $`m_c`$ where the $`c`$ quark velocity is rather large—yet there the energy release is significant, and a large number of excited states (up to $`18`$ at $`m_c=5\beta 2\text{GeV}`$) are produced. Table 7 and Fig. 9 show analogous results for $`m_b=10\beta `$, $`m=0.56\beta `$.
To render the duality violation more apparent, we consider (Table 8 and Fig. 10) the same decay widths for a $`b`$ quark with half the mass, $`m_b=5\beta 2\text{GeV}`$. Even here the deviation is below per mill as soon as the first excitation can appear with sufficient phase space. The duality-violating component at last exhibits the proper oscillating behavior (note the decrease between $`m_c=3`$ and $`3.5\beta `$ or 1 and $`1.5\beta `$), but this effect is too small to be extracted reliably at larger energy release where this property becomes an asymptotic rule.
As follows from our computations, local duality is violated at a tiny level in the $`bc`$ decays in the ’t Hooft model whenever it is a priori meaningful to apply OPE. A possible reason behind this might be that for unidentified reasons the heavy quark symmetry works for the inclusive widths too effectively, down to relatively low masses and velocities of order $`1`$. This was conjectured in the early papers on the subject .
Therefore, our final attempt in the quest for the sizeable duality violation in beauty is considering the $`(bu)`$-type transitions, where the heavy quark symmetry per se does not constrain the individual transition form factors. We fix in our expressions $`m_c=m=0.56\beta `$ or $`0.26\beta `$ (but still keep the two quarks flavor-distinguished) and vary $`m_b`$ from $`1\beta 0.4\text{GeV}`$ to $`12\beta 4.5\text{GeV}`$. The results are shown in Table 9 and Fig. 11, and Table 10 and Fig. 12, respectively. Although the difference between the actual width and its OPE approximation is larger, it still is very small and approaches a percent level for $`m_b`$ as low as $`2\beta 0.8\text{GeV}`$. The total decay width is no longer saturated to such a high degree by transitions to the ground state, especially for larger $`m_b`$. Nevertheless, the duality is amazingly well satisfied when just the first few open channels are summed. Again, the only prominent features in the plots appear when crossing kinematic thresholds due to the lightest $`D`$ mesons of opposite parity to the ground-state $`B`$.
The extraordinary agreement between $`\mathrm{\Gamma }_B`$ and $`\mathrm{\Gamma }_{\mathrm{OPE}}`$ may be underscored by instead plotting (Fig. 13, final column of Table 9) the difference between $`\mathrm{\Gamma }_B`$ and the Born-term partonic rate $`\mathrm{\Gamma }_b`$ given in Eq. (39). From an algebraic point of view, $`\mathrm{\Gamma }_B`$ and $`\mathrm{\Gamma }_{\mathrm{OPE}}`$ differ generically at $`O(1/M^5)`$, while $`\mathrm{\Gamma }_B`$ and $`\mathrm{\Gamma }_b`$ begin to differ already at $`O(1/M^2)`$.
Thus, we find local duality between the actual semileptonic decay width and its OPE expansion to be very well satisfied in all cases.
Before concluding this section, let us briefly address duality in the differential distribution $`\mathrm{\Gamma }^1\mathrm{d}\mathrm{\Gamma }/\mathrm{d}E`$. In the heavy quark limit the shape of the final-state hadronic mass distribution follows the heavy quark distribution function in the decaying meson; for the $`bu`$ decays under consideration, this is the light-cone distribution function $`F(x)`$. In decays with $`q^2=0`$ the recoil energy of the lepton pair $`E`$ is directly related to the final state mass $`M_h`$:
$$E=\frac{M_B^2M_h^2}{2M_B}.$$
(40)
Since $`q^2=0`$, these decays are analogous to $`bs\gamma `$ in the Standard Model. In the large-$`m_b`$ limit one has
$$\frac{1}{\mathrm{\Gamma }}\frac{\mathrm{d}\mathrm{\Gamma }}{\mathrm{d}E}=\frac{2}{\overline{\mathrm{\Lambda }}}F\left(\frac{2Em_b}{\overline{\mathrm{\Lambda }}}\right).$$
(41)
At finite $`m_b`$ in a theory with narrow resonances the actual distribution is given by the comb of $`\delta `$-functions with spacing in the argument of Eq. (41) of order $`\overline{\mathrm{\Lambda }}/m_b`$. In order to obtain a continuous result, we adopt the simple ansatz of averaging over the peaks. Using Eq. (40) to define the energy $`E_n`$ of the $`n`$th state $`M_h`$, we integrate the $`\delta `$-function for the $`n`$th state evenly over the energy range $`(E_n+E_{n+1})/2`$ to $`(E_n+E_{n1})/2`$, i.e., the midpoints between energy eigenvalues. Letting $`N`$ be the maximum number of kinematically allowed $`M_h`$ values, we establish the endpoint bins by defining $`E_1=E_{\mathrm{max}}=M_B/2`$ and $`E_{N+1}=E_{\mathrm{min}}=0`$.
We find that our numerical computations yield a distribution resembling the light-cone distribution function $`\phi ^2`$; specifically,
$$F(y)=\overline{\mathrm{\Lambda }}\mathrm{\Psi }^2\left((1y)\overline{\mathrm{\Lambda }}\right)=\underset{m_Q\mathrm{}}{lim}\frac{\overline{\mathrm{\Lambda }}}{m_Q}\phi ^2\left(1(1y)\frac{\overline{\mathrm{\Lambda }}}{m_Q}\right).$$
(42)
Recalling that $`M_B=m_b+\overline{\mathrm{\Lambda }}+O(1/m_b)`$ and combining Eqs. (41) and (42) yields
$$\frac{1}{\mathrm{\Gamma }}\frac{\mathrm{d}\mathrm{\Gamma }}{\mathrm{d}E}\underset{m_Q\mathrm{}}{lim}\frac{2}{m_Q}\phi ^2\left(1\frac{M_B2E}{m_Q}\right).$$
(43)
The two sides of this expression are plotted in Fig. 14, using $`m_Q=25\beta `$ to represent the limit $`m_Q\mathrm{}`$, while the actual distribution is considered at $`m_b=10\beta `$. The agreement is quite remarkable. The continuous distribution appears to pass approximately through the midpoint of each bin; owing to the near-equal spacing of ’t Hooft model eigenvalues in $`M_n^2`$, Eq. (40) shows that these bin midpoints are very close to the values $`E_n`$ themselves.
It is also interesting to consider integration over a range of $`E`$. In particular, define $`\mathrm{\Phi }(12E/M_B)`$ as the cumulative fractional width from maximum energy $`M_B/2`$ down to the given $`E`$; then $`\mathrm{\Phi }(0)=0`$ and $`\mathrm{\Phi }(1)=1`$. While the exact result for $`\mathrm{\Delta }\mathrm{\Gamma }/\mathrm{\Gamma }`$ amounts to an integration of the $`\delta `$-function differential widths renormalized so that the cumulative result approaches unity, the integral of the continuous distribution gives
$$\mathrm{\Phi }(y)=\underset{m_Q\mathrm{}}{lim}\frac{2}{m_Q}_0^ydz\phi ^2\left(1\frac{M_{H_Q}}{m_Q}z\right).$$
(44)
These two curves are presented in Fig. 15. Two features particularly stand out in this plot. First, even for $`m_b`$ as large as $`10\beta 4\text{GeV}`$, the overwhelming part of the decay probability falls into the transitions to at most four lowest states. Second, the continuous curve seems to provide a nearly optimal description possible for the step-like exact distribution. The point-to-point deviation for all plotted values with $`12E/M_B>0.04`$ does not exceed half of the contribution of the nearest threshold.
## 6 Duality in the Vacuum Current Correlator
In this section we briefly illustrate the onset of duality for the absorptive part of the vector current correlator with light quarks, of the type that determines the normalized cross section $`R(e^+e^{}\text{hadrons})`$ as a function of energy. In the context of the heavy quark decays this is relevant in nonleptonic decay widths in two kinds of processes: in spectator-independent decays, where $`R(q^2)`$ determines the weight with which the semileptonic width at given $`q^2`$ must be integrated over $`q^2`$ (see Ref. ), and in the effects of WA decays.
In either case, at $`N_c\mathrm{}`$ the cross section appears as a comb-like collection of $`\delta `$-functions:
$$R(q^2)=\underset{n}{}c_n^2\delta (q^2M_n^2);c_n=_0^1dx\phi _n(x).$$
(45)
The above expression for the residues refers to the case where a vector current is considered. We suppress here the factor of $`\sqrt{N_c/\pi }`$ relating $`c_n`$ to $`f_n`$ \[Eq. (68)\]. We also assume in what follows that the light quark masses are equal, $`m_u=m_d`$, and are $`O(\beta )`$ or less, in order to reach asymptotic $`q^2`$ more quickly.
In the extreme situation of infinitely narrow resonances one cannot, of course, discuss a point-to-point equality of the cross section $`R(q^2)`$ with its OPE in the form of $`1/q^2`$ expansion. A meaningful comparison is possible if each resonant peak is somehow averaged over an interval no smaller than the distance between adjacent peaks, the latter being approximately given by $`\mathrm{\Delta }q^2\pi ^2\beta ^2`$ . It is worth recalling that $`R(q^2)`$ is proportional to $`m^2/q^4`$, so one must consider nonvanishing masses for the vector current, and address the OPE terms formally suppressed by $`m^2/q^2`$.
This question was first addressed in the context of nonleptonic decays in Ref. using the numerical approach. Duality for the average cross section in the same manner as above, i.e., using sum rules derived from the ’t Hooft equation and analytically matching terms in the $`1/m_Q`$ expansion, was obtained in Ref. . Yet establishing the asymptotics per se cannot tell us beforehand how early one can expect the onset of duality. Here we study this question numerically, in the domain of intermediate $`q^2`$.
The concrete amount of the deviation between $`R(q^2)`$ and $`R^{\mathrm{OPE}}(q^2)`$ in the case of direct resonances may depend in an essential way on the chosen smearing procedure. Interested in the qualitative features only, we choose a rather simplified, crude method: We spread the integral of $`R(q^2)`$ evenly over the interval between the successive resonances. More precisely, we put
$$\overline{R}(q^2)=\frac{1}{M_{2n+1}^2M_{2n1}^2}_{M_{2n1}^2}^{M_{2n+1}^2}dq^2R(q^2)=\frac{c_{2n}^2}{M_{2n+1}^2M_{2n1}^2},$$
(46)
for $`M_{2n1}^2<q^2<M_{2n+1}^2`$, with $`M_1^2=4m^2`$, the partonic pair production threshold. Here we use the fact that $`c_n`$ vanish for odd $`n`$ when $`m_u=m_dm`$. This smearing is very similar to that described for the differential width in the last section, except that averaging is performed in $`q^2`$ rather than $`E`$. The free quark loop $`R(q^2)`$, which is of course the leading term of the OPE, is given by
$$R_0(q^2)=\frac{2m^2}{q^4}\frac{1}{\sqrt{14m^2/q^2}}.$$
(47)
Table 11 and Fig. 16 show the results for our reference case $`m=0.56\beta `$. The agreement of the average hadronic cross section with the parton-computed probability again turns out to be very good. Apparently, this can be related to two facts: the heavy suppression of power corrections to $`R(q^2)`$ in the OPE (see Eqs. (34)–(35) in ), and an early onset of the asymptotics in the spectrum,
$$M_{n+1}^2M_n^2\pi ^2\beta ^2,$$
(48)
which even at $`n=0`$ is satisfied to about $`15\%`$.
## 7 Discussion and Summary
The main motivation behind the present study has been to assess the magnitude of local duality violations in the inclusive semileptonic decays of beauty particles. We considered this question using the ’t Hooft model as a toy theory in which all relevant decay amplitudes can be evaluated numerically. The ’t Hooft model, while retaining certain key features of full $`D=4`$ QCD that shape the spectrum of hadrons (quark confinement, chiral symmetry breaking), still differs from $`D=4`$ in many respects. Yet using it as a lab for exploration carries an important advantage—it allows no “wiggle room” for interpretation of the results. There are no ad hoc parameters to choose or adjust, and as soon as the underlying weak decay Lagrangian is fixed, the numerical results are unambiguous and must be accepted at face value. This positively distinguishes this approach from various models where often the conclusions, even qualitatively, depend on the arbitrary choice of parameters according to one’s preferences. The question of a particular model being compatible with the general dynamical properties of QCD underlying the OPE approach, often quite problematic in simplified quark models, does not arise for the ’t Hooft model.
Although the simplest illustration of the asymptotic nature of the decay width $`1/m_Q`$ expansion and related violations of local duality follows just from the existence of hadronic thresholds (see, e.g., ), violation of local duality is a more universal phenomenon that is not directly related to existence of hadronic resonances nor even confinement itself. This has been illustrated in Ref. by the example of soft instanton effects that do not lead, at least at small density, to quark confinement—but do indeed generate computable oscillating duality-violating contributions to the total decay rates.
Nevertheless, there is a widespread opinion that decays with manifest resonance structure in the final state are most difficult for—if compatible at all with—the standard OPE. Even the possibility that the OPE does not fully apply in the case of “hard” confinement has been occasionally voiced in the literature. The analytic studies performed in Refs. , which explicitly demonstrate in the ’t Hooft model the applicability of the OPE to the total widths, should help to allay such conceptual concerns. Nevertheless, the intuition remains that resonance dominance is not “favorable” for the OPE, and problems might show up, for instance, through a delayed numerical onset of duality, in that the approximate equality of the OPE predictions and the actual decay widths may set in only after a significant number of thresholds has been passed. To address such issues, the ’t-Hooft model seems to represent the most certain testing ground for local duality in the domain of decays of moderately heavy quarks.
Contrary to naive expectations, we found surprisingly accurate duality between the (truncated) OPE series for $`\mathrm{\Gamma }_{\mathrm{sl}}`$ and the actual decay widths. The deviations are suppressed to a very high degree almost immediately after the threshold for the first excited final state hadron is passed. No suspected delay in the onset of duality was found.
The key property that governs the onset of the $`1/m_Q`$ expansion for the semileptonic widths is the pattern of saturation of the heavy quark sum rules. We examined a particular class, the SV sum rules in the heavy quark limit, that has the most transparent quantum mechanical meaning. We found them saturated to an amazing degree by the very first excitation. The contribution of the remaining, higher states to the slope of the IW function, $`\overline{\mathrm{\Lambda }}`$, and $`\mu _\pi ^2`$ does not exceed a few percent. Even in the Darwin operator sum rule, the first excitation accounts for $`90\%`$ of the whole expectation value, despite the fast-growing weight, $`(ϵ_kϵ_0)^3`$ of higher-order contributions. This peculiarity underlies the early onset of duality for the case when initial- and final-state quarks are both heavy.
Some of the duality-violating features observed in these studies have natural explanations. At fixed energy release $`m_Qm_q`$ the magnitude of the deviations is smaller if $`m_Q`$, $`m_q`$ are both large (as in $`bc`$) than if they are both small. This is expected, since in the former case the heavy quark symmetry for the elastic amplitude additionally enforces approximate duality even when no expansion in large energy release can be applied.
It is interesting, however, that at fixed $`m_b`$ the duality violation decreases rapidly as $`m_c`$ decreases, in full accord with the OPE where the higher order terms are generally suppressed by powers of $`1/(m_bm_c)`$. This is clearly a dynamical feature that goes beyond heavy quark symmetry per se, the quality of which deteriorates as $`m_c`$ decreases.
It is also instructive to note that including the calculated power-suppressed OPE terms significantly reduces the difference between the actual decay width and its purely partonic evaluation. Moreover, the seeds of oscillations inherent to duality violation (as functions of quark masses), can be seen. Since we adopted the truncated OPE expansion to mirror the existing implementation of the OPE in QCD, the deviations do not average to zero but rather oscillate around the (rapidly dissipating) contributions attributed to discarded higher-order terms.
The numerical effects of duality violation we study turn out to be typically quite small. Partially this can be attributed to moderate size of the corresponding expectation values multiplying $`1/m_Q^k`$ corrections in the OPE. Yet certainly not all power corrections in heavy quarks are suppressed in the model. It is well known from ordinary quantum mechanics that masses (eigenvalues) typically are much more robust against perturbations than wavefunctions themselves (or transition amplitudes). We observe a similar pattern in the ’t Hooft model. For example, $`1/m_Q`$ corrections to the meson decay constants turn out very significant even at the scale of the $`b`$ quark mass. Apparently, the inclusive decay rates fall into the class of “robust” observables, although, as explained above, this was difficult to anticipate beforehand.
We note here another “fragile” observable, the light-cone heavy quark distribution function, which can be measured in decays of the type $`bs\gamma `$. In $`D=2`$ the scaled spread $`m_Q^2\left(x^2x^2\right)`$ of the $`x`$ distribution approaches $`\mu _\pi ^2`$ at large $`m_Q`$. Yet, as seen in Fig. 4, even at the $`b`$ quark mass one would obtain from this distribution only about $`60\%`$ of the actual value of $`\mu _\pi ^2`$, due to significant $`1/m_Q`$ corrections. This caveat may be important for existing analyses of the decay distributions in $`B`$ decays, where such effects routinely are not included.
We also briefly addressed the inclusive differential decay distributions in the analogues of $`bu\mathrm{}\overline{\nu }`$ or $`bs\gamma `$ decays. Generally, we find good agreement (at the scale corresponding to the physical $`b`$ mass) with the parton-based prediction incorporating effects of the “Fermi motion,” and in particular for the partially integrated probability
$$\mathrm{\Phi }(x)=\frac{1}{\mathrm{\Gamma }_{\mathrm{sl}}}_0^{xM_B^2}dM_h^2\frac{\mathrm{d}\mathrm{\Gamma }_{\mathrm{sl}}}{\mathrm{d}M_h^2}.$$
(49)
This distribution, following Refs. , is examined in real $`B`$ decays in the quest for $`|V_{ub}|`$ . However, the point-to-point deviations are clearly still significant, for the decays to only the 4 or 5 lowest final states saturate the overwhelming fraction of the total decay probability. It is quite conceivable, though, that such deviations are less pronounced in actual QCD owing to the significant resonance widths and to a richer resonance structure.
The vacuum current correlator also turns out to be especially robust; even neglecting all OPE corrections except the leading partonic contribution leads to excellent agreement with the hadronic result.
Turning to the direct phenomenological conclusions that can be inferred from our studies, we see that, to the extent our findings can be transferred to real QCD, violation of local duality in the total semileptonic widths of $`B`$ mesons is not an issue. The scale of duality violation lies far below the phenomenologically accessible limits, and cannot affect the credibility of $`|V_{cb}|`$ or $`|V_{ub}|`$ extractions.
In reality there are, of course, essential conceptual differences between the two theories, including those aspects that are expected to be essential for local duality (for a discussion, see Ref. ). Although many seem to optimistically suggest that duality violation is more pronounced in the ’t Hooft model than for actual heavy flavor hadrons, some differences may still work in the opposite direction. In $`D=2`$ there are no dynamical gluons, nor a chromomagnetic field that in $`D=4`$ provides a significant scale of nonperturbative effects in heavy flavor hadrons. Likewise, there is no spin in $`D=2`$, and no corresponding $`P`$-wave excitations of the light degrees of freedom (the so-called $`j=3/2`$ states), which seem to play an important role in $`D=4`$.
Two-dimensional QCD neither has long perturbative “tails” of actual strong interactions suppressed weakly (by only powers of $`\mathrm{log}`$s of the energy scale). In $`D=2`$ the perturbative corrections are generally power-suppressed, as follows from the dimension of the gauge coupling. As discussed in Ref. , it is conceivable that the characteristic mass scale for freezing out the transverse gluonic degrees of freedom is higher than in the “valence” quark channels. This would imply a possibly higher scale for onset of duality in $`\alpha _s/\pi `$ corrections to various observables.
Regardless of these differences, we conclude that presence of resonance structure per se is not an obstacle for fine local quark-hadron duality tested in the context of the OPE. As we see in the ’t Hooft model, resonances themselves do not seem to demand a larger duality interval. As soon as the mass scale of the states saturating the sum rules in a particular channel (quark or hybrid) has been passed, the decay width can be well approximated numerically by the expansion stemming from the OPE.
The ground states of heavy mesons in the ’t Hooft model exhibit relatively small expectation values of nonperturbative operators ($`\mu _\pi ^2`$, $`\rho _D^3`$, but not $`\overline{\mathrm{\Lambda }}`$) compared to real QCD, if our identification $`\beta 400\text{MeV}`$ is adopted. This may be regarded as a reason for small duality violation for $`\mathrm{\Gamma }_{\mathrm{sl}}`$ in the model. However, even if we scale $`\beta `$ up to $`700`$$`800\text{MeV}`$ to make up for smallness of the nonperturbative OPE effects, the duality violation is still very small, and superficially rather insignificant even in charm.
We note, however, that the specific choice Eq. (30) of the weak interaction effectively requires decays to occur only at $`q^2=0`$, and therefore the effects of four-fermion operators of the type $`(\overline{Q}\mathrm{\Gamma }q)(\overline{q}\mathrm{\Gamma }Q)`$ are totally absent, at least in the lowest orders of perturbation theory (cf. Ref. , Sec. III.B.3). As was suggested in Ref. , it is conceivable that the apparent excess in $`\mathrm{\Gamma }_{\mathrm{sl}}(D)`$ is simply related to a noticeable magnitude of the non-valence (nonfactorizable) expectation values $`D\left|(\overline{c}\mathrm{\Gamma }s)(\overline{s}\mathrm{\Gamma }c)\right|D`$. If this conjecture is true, similar effects in $`\mathrm{\Gamma }_{\mathrm{sl}}(BX_u\mathrm{}\nu )`$ are still suppressed but possibly detectable in future precision experiments. In the context of the present study, it suffices to say that this would be a legitimate OPE effect rather than a manifestation of a significant local duality violation in the strict sense.
Acknowledgments: R.L. thanks the Department of Energy for support under Contract No. DE-AC05-84ER40150; N.U. acknowledges the support of the NSF under grant number PHY96-05080, by NATO under the reference PST.CLG 974745, and by RFFI under grant No. 99-02-18355. We are grateful to N. Isgur for inspiring interest and discussions, and to M. Burkardt for invaluable insights. N.U. also thanks I. Bigi, M. Shifman and A. Vainshtein for encouraging interest and collaboration on related issues, and A. Zhitnitsky for useful comments. N.U. enjoyed the hospitality of Physics Department of the Technion and the support of the Lady Davis grant during completion of this paper.
## Appendix A The BSW Improvement of the Multhopp Technique
The Brower-Spence-Weis (BSW) improvement of the Multhopp technique avoids the need for evaluating the wavefunction at a discrete set of points called “Multhopp angles,” thus improving the behavior of the solutions in the endpoint regions, as described in Sec. 2. Here we exhibit the expressions used by BSW, correcting along the way some minor typographical errors in their work.
Starting with the ’t Hooft equation (1) with bare quark masses $`m_1`$ and $`m_2`$, one converts the kinematic variables $`x,y`$ to angular variables:
$$x=\frac{1+\mathrm{cos}\theta }{2},y=\frac{1+\mathrm{cos}\theta ^{}}{2},$$
(50)
in terms of which the ’t Hooft equation reads
$`{\displaystyle \frac{M_p^2}{2}}\phi _p(\theta )`$ $`=`$ $`\left[{\displaystyle \frac{m_1^2}{1+\mathrm{cos}\theta }}+{\displaystyle \frac{m_2^2}{1\mathrm{cos}\theta }}\right]\phi _p(\theta )`$ (51)
$`+{\displaystyle _0^\pi }d\theta ^{}\phi _p(\theta ^{})\mathrm{P}{\displaystyle \frac{1}{(\mathrm{cos}\theta \mathrm{cos}\theta ^{})^2}}.`$
Expanding
$$\phi _p(\theta )=\underset{n=1}{\overset{\mathrm{}}{}}a_n^{(p)}\mathrm{sin}n\theta ,$$
(52)
and using the continuous inversion identity (contrast with Eq. (A7) of Ref. )
$$_0^\pi d\theta \mathrm{sin}m\theta \mathrm{sin}n\theta =\frac{\pi }{2}\left(\delta _{mn}\delta _{m,n}\right),$$
(53)
one obtains the infinite-dimensional eigenvector system
$$M_p^2a_n^{(p)}=\left(H_0+V\right)_{nm}a_m^{(p)},$$
(54)
where
$`\left(H_0\right)_{nm}`$ $`=`$ $`+{\displaystyle \frac{4}{\pi }}{\displaystyle _0^\pi }d\theta \left[{\displaystyle \frac{m_1^2}{1+\mathrm{cos}\theta }}+{\displaystyle \frac{m_2^2}{1\mathrm{cos}\theta }}\right]\mathrm{sin}n\theta \mathrm{sin}m\theta ,`$ (55)
$`V_{nm}`$ $`=`$ $`{\displaystyle \frac{4}{\pi }}\beta ^2{\displaystyle _0^\pi }d\theta \mathrm{sin}n\theta {\displaystyle _0^\pi }d\theta ^{}\mathrm{sin}\theta ^{}\mathrm{sin}m\theta ^{}\mathrm{P}{\displaystyle \frac{1}{(\mathrm{cos}\theta \mathrm{cos}\theta ^{})^2}}.`$ (56)
Both of these integrals can be evaluated, with the result
$`\left(H_0\right)_{nm}`$ $`=`$ $`4\mathrm{min}(n,m)\left[(1)^{m+n}m_1^2+m_2^2\right],`$ (57)
$`V_{nm}`$ $`=`$ $`V_{n1,m1}\left({\displaystyle \frac{m}{m1}}\right)+{\displaystyle \frac{8m}{n+m1}}\left[{\displaystyle \frac{1+(1)^{n+m}}{2}}\right],`$ (58)
where $`V_{n0}=V_{0m}=0`$ for $`m,n0`$. This recursive form for $`V_{nm}`$ is most convenient for numerical calculations; however, one may also write the closed-form solution,
$$V_{nm}=4m\left[\frac{1+(1)^{n+m}}{2}\right]\left[\psi \left(\frac{1nm^{}}{2}\right)\psi \left(\frac{1|nm|}{2}\right)\right].$$
(59)
Note that the “potential” $`V`$ in Eq. (59) is real but not symmetric, owing to the extra $`\mathrm{sin}\theta ^{}`$ in Eq. (56); therefore, the “Hamiltonian” $`H_0+V`$ is not Hermitian, and the eigenvectors $`a^{(p)}`$ are not orthogonal. This is a direct result of converting the exact wavefunctions, which are eigenfunctions of a Hermitian Hamiltonian when written in terms of the variable $`x`$ (and therefore orthogonal in $`x`$), into orthogonal functions of the variable $`\theta `$. This transformation is nonunitary because the number of modes used is not infinite; therefore, the overlap of different eigenvector solutions should be small when a large number of modes are used. Indeed, this turns out to be empirically true; nevertheless, we take the further step of orthogonalizing the numerical eigenvector solutions recursively by means of the standard Gram-Schmidt procedure, i.e.,
$$|\phi _{\mathrm{orth}}^{(p)}=\frac{|\phi ^{(p)}_{j=0}^{p1}|\phi _{\mathrm{orth}}^{(j)}\phi _{\mathrm{orth}}^{(j)}|\phi ^{(p)}}{\sqrt{\phi _{}^{(p)}|\phi _{}^{(p)}_{j=0}^{p1}\phi _{\mathrm{orth}}^{(j)}|\phi _{}^{(p)}}}.$$
(60)
For $`N=500`$ modes, this typically changes expectation values by one part in $`10^5`$.
The expressions for these overlaps and other matrix elements in terms of the mode coefficients $`a_n`$ are presented in Appendix B.
## Appendix B Matrix Elements
A number of useful overlaps and other integrals are straightforward to evaluate in terms of the mode coefficients, using the expressions (52). Solving them amounts to evaluating a number of trigonometric integrals. Such expressions are especially convenient since they permit a number of integrations that introduce no numerical uncertainties (except due to machine precision) beyond those of solving the original Multhopp-BSW eigenvector equation (54).
In particular, denote the $`p`$th eigenstate wavefunction presented in Eq. (52) by $`\phi _p^{(a)}`$ and that for some other set of masses in the $`q`$th eigenstate by $`\phi _q^{(b)}`$; the latter wavefunction then has an expansion like (52) with mode coefficients $`b_n^{(q)}`$. Truncating after $`N`$ modes, one then finds
$`\phi _p^{(a)}|\phi _q^{(b)}`$ $`=`$ $`{\displaystyle _0^1}dx\phi _p^{(a)}(x)\phi _q^{(b)}(x)`$
$`=`$ $`2{\displaystyle \underset{m=1}{\overset{N}{}}}ma_m^{(p)}{\displaystyle \underset{n=1}{\overset{N}{}}}nb_n^{(q)}\left[{\displaystyle \frac{1+(1)^{m+n}}{2}}\right]{\displaystyle \frac{1}{\left[1(mn)^2\right]\left[1(m+n)^2\right]}}.`$
Indeed, the normalization integral $`_0^1dx\phi (x)^2=1`$ is just the case $`a=b`$ and $`p=q`$, in agreement with Eq. (A9) of Ref. .
Other useful expectation values include
$`x{\displaystyle \frac{1}{2}}_p`$ $`=`$ $`{\displaystyle _0^1}dx\left(x{\displaystyle \frac{1}{2}}\right)\left[\phi _p^{(a)}(x)\right]^2`$
$`=`$ $`{\displaystyle \underset{m=1}{\overset{N}{}}}ma_m^{(p)}{\displaystyle \underset{n=1}{\overset{N}{}}}na_n^{(p)}\left[{\displaystyle \frac{1(1)^{m+n}}{2}}\right]{\displaystyle \frac{1}{\left[4(mn)^2\right]\left[4(m+n)^2\right]}},`$
$`\left(x{\displaystyle \frac{1}{2}}\right)^2_p={\displaystyle _0^1}dx\left(x{\displaystyle \frac{1}{2}}\right)^2\left[\phi _p^{(a)}(x)\right]^2`$ (63)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{m=1}{\overset{N}{}}}ma_m^{(p)}{\displaystyle \underset{n=1}{\overset{N}{}}}na_n^{(p)}\left[{\displaystyle \frac{1+(1)^{m+n}}{2}}\right]\left[216(m^2+n^2)+(m^2n^2)^2\right]`$
$`\times \left[\left(1(mn)^2\right)\left(1(m+n)^2\right)\left(9(mn)^2\right)\left(9(m+n)^2\right)\right]^1.`$
Note that the spread of the wavefunction may be computed about any convenient point in $`x`$, viz.,
$$\left(ax+b\right)^2\left(ax+b\right)^{}^2=a^2\left(x^2x^{}^2\right),$$
(64)
so that the additive constants of $`1/2`$ above are irrelevant. Also,
$$\frac{1}{x}_p=_0^1dx\frac{1}{x}\left[\phi _p^{(a)}(x)\right]^2=\underset{m=1}{\overset{N}{}}a_m^{(p)}\underset{n=1}{\overset{N}{}}a_n^{(p)}I_{mn},$$
(65)
where
$`I_{mn}`$ $`=`$ $`2{\displaystyle \underset{j=|mn|/2}{\overset{(m+n)/21}{}}}{\displaystyle \frac{1}{2j+1}}`$
$`=`$ $`\psi \left({\displaystyle \frac{m+n^{}+1}{2}}\right)\psi \left({\displaystyle \frac{|mn|+1}{2}}\right),mn\mathrm{even};`$
$`I_{mn}`$ $`=`$ $`{\displaystyle \frac{1}{|mn|}}{\displaystyle \frac{1}{m+n}}2{\displaystyle \underset{j=(|mn|1)/2}{\overset{(m+n1)/21}{}}}{\displaystyle \frac{1}{2j+1}}`$ (66)
$`=`$ $`{\displaystyle \frac{1}{|mn|}}{\displaystyle \frac{1}{m+n}}+\psi \left({\displaystyle \frac{|mn|}{2}}\right)\psi \left({\displaystyle \frac{m+n^{}}{2}}\right),mn\mathrm{odd}.`$
One also finds
$$\frac{1}{1x}_p=_0^1dx\frac{1}{1x}\left[\phi _p^{(a)}(x)\right]^2=\underset{m=1}{\overset{N}{}}a_m^{(p)}\underset{n=1}{\overset{N}{}}a_n^{(p)}J_{mn},$$
(67)
where, using the notation of Eq. (66), one finds $`J_{mn}=+I_{mn}`$ for $`mn`$ even, and $`J_{mn}=I_{mn}`$ for $`mn`$ odd.
Finally, the decay constant of the $`p`$th excitation \[cf. Eqs. (10)–(11)\] is given by
$$f_p^{(a)}=\sqrt{\frac{N_c}{\pi }}_0^1dx\phi _p^{(a)}(x)=\sqrt{\frac{N_c}{\pi }}c_p=\sqrt{\frac{N_c}{\pi }}\times \frac{\pi }{4}a_1^{(p)}.$$
(68)
## Appendix C Additional Relations Used in the Analysis
The numerical calculation of large-$`m_Q`$ matrix elements with acceptable accuracy relies on achieving a balance between competing effects.
On one hand, Multhopp solutions to the ’t Hooft equation with $`m_Q\beta `$ tend to suffer degraded numerical accuracy since they are highly concentrated into the small kinematic region $`1x1`$. As discussed in Sec. 2, the endpoint regions $`x0`$ and $`1`$ are where the Multhopp solutions—or more precisely, their derivatives—tend to break down. This effect is compounded when $`m\beta `$, since lighter quark masses force sharper endpoint behavior in the wavefunction. Although the BSW solution ameliorates this behavior, as $`m_Q`$ is increased one eventually faces the problem of attempting to represent a function with only a very small region of support in $`x`$ by a finite number of modes with support over the full range $`x[0,1]`$. In practice, we gauge the errors committed through such “lattice spacing” effects by computing a given quantity with $`N=500`$ and noting the amount by which its value shifts if one uses instead $`N=100`$, and as expected, such errors become substantial (as much as a few percent) by the time one reaches $`m_Q>25\beta `$ or $`m<0.4\beta `$.
On the other hand, although numerical solutions with $`m_Q,mO(\beta )`$ have the highest numerical accuracy, they also have substantial $`O(1/m_Q)`$, $`O(1/m_Q^2)`$, etc. corrections that are difficult to disentangle.
We adopt an intermediate strategy of employing certain exact relations that hold for the ’t Hooft solutions. To determine the relevant static expectation values, we solve the finite-$`m_Q`$ heavy hadron mass expansion for $`\overline{\mathrm{\Lambda }}`$ \[Eq. (6)\]:
$$M_{H_Q}m_Q=\overline{\mathrm{\Lambda }}+\frac{\mu _\pi ^2\beta ^2}{2m_Q}+\frac{\rho _D^3\rho _{\pi \pi }^3}{4m_Q^2}+O\left(\frac{\beta ^4}{m_Q^3}\right).$$
(69)
Neglecting the order term and using the relations \[Eq. (9)\]
$$\mu _\pi ^2=\frac{\overline{\mathrm{\Lambda }}^2m^2+\beta ^2}{3},\rho _D^3=\frac{\beta ^2F^2}{4},\rho _{\pi \pi }^3=\frac{1}{36}\left[8\overline{\mathrm{\Lambda }}(\overline{\mathrm{\Lambda }}^2m^2+\beta ^2)+3\beta ^2F^2\right],$$
(70)
we thus arrive at an equation cubic in $`\overline{\mathrm{\Lambda }}`$ that depends on $`F^2`$. We solve it at $`m_Q=15\beta `$.
The asymptotic value of the scaled decay constant $`F^{(n)}=\sqrt{m_Q}c_n`$ must also be evaluated at a finite value of $`m_Q`$, thus including $`1/m_Q`$-suppressed pieces. We account for them explicitly using the expansion \[the first of Eqs. (3)\]
$$\sqrt{m_Q}c_n=\left(1\frac{2[2\overline{\mathrm{\Lambda }}^{(n)}m(1)^n]}{3m_Q}\right)F^{(n)}+O\left(\frac{\beta ^{5/2}}{m_Q^2}\right).$$
(71)
We likewise solve this equation for $`F^{(n)}`$ at $`m_Q=15\beta `$.
Turning to the analysis of the SV sum rules Eqs. (22)–(25) in Sec. 4, we note that their rapid saturation demands an exceptionally high precision in evaluating both the oscillation strengths $`\tau `$ in the r.h.s. and the expectation values in the l.h.s. Reaching such an accuracy through direct computation seems impossible. Therefore, we use a number of identities to get meaningful results. First, we employ the expression for $`\tau _{nk}`$ in terms of $`ϵ_n`$, $`ϵ_k`$, and the corresponding decay constants:
$$\tau _{nk}=\frac{\beta ^2}{2(ϵ_nϵ_k)^3}F^{(n)}F^{(k)}\left(\frac{1(1)^{nk}}{2}\right).$$
(72)
Then we make use of the fact that the discussed sum rules, being completeness sums, are exact when summation includes all excitations (see Ref. ). Therefore, one has
$`{\displaystyle \frac{1}{\rho _k^2\frac{1}{4}}}\left[\rho _k^2{\displaystyle \frac{1}{4}}{\displaystyle \underset{\mathrm{}=1}{\overset{n}{}}}\tau _\mathrm{}k^2\right]`$ $`=`$ $`{\displaystyle \frac{1}{\rho _k^2\frac{1}{4}}}{\displaystyle \underset{\mathrm{}=n+1}{\overset{\mathrm{}}{}}}\tau _\mathrm{}k^2,`$ (73)
$`{\displaystyle \frac{2}{\overline{\mathrm{\Lambda }}_k}}\left[{\displaystyle \frac{1}{2}}\overline{\mathrm{\Lambda }}_k{\displaystyle \underset{\mathrm{}=1}{\overset{n}{}}}(ϵ_{\mathrm{}}ϵ_k)\tau _\mathrm{}k^2\right]`$ $`=`$ $`{\displaystyle \frac{2}{\overline{\mathrm{\Lambda }}_k}}{\displaystyle \underset{\mathrm{}=n+1}{\overset{\mathrm{}}{}}}(ϵ_{\mathrm{}}ϵ_k)\tau _\mathrm{}k^2,`$ (74)
$`{\displaystyle \frac{1}{\left(\mu _\pi ^2\right)_k}}\left[\left(\mu _\pi ^2\right)_k{\displaystyle \underset{\mathrm{}=1}{\overset{n}{}}}(ϵ_{\mathrm{}}ϵ_k)^2\tau _\mathrm{}k^2\right]`$ $`=`$ $`{\displaystyle \frac{1}{\left(\mu _\pi ^2\right)_k}}{\displaystyle \underset{\mathrm{}=n+1}{\overset{\mathrm{}}{}}}(ϵ_{\mathrm{}}ϵ_k)^2\tau _\mathrm{}k^2,`$ (75)
$`{\displaystyle \frac{1}{\left(\rho _D^3\right)_k}}\left[\left(\rho _D^3\right)_k{\displaystyle \underset{\mathrm{}=1}{\overset{n}{}}}(ϵ_{\mathrm{}}ϵ_k)^3\tau _\mathrm{}k^2\right]`$ $`=`$ $`{\displaystyle \frac{1}{\left(\rho _D^3\right)_k}}{\displaystyle \underset{\mathrm{}=n+1}{\overset{\mathrm{}}{}}}(ϵ_{\mathrm{}}ϵ_k)^3\tau _\mathrm{}k^2.`$ (76)
The sums on the r.h.s. can be accurately evaluated since the higher contributions fall off in magnitude very quickly. In practice, we truncate the sum at $`\mathrm{}=20`$.
A similar approach was used to evaluate the duality-violating difference $`\mathrm{\Gamma }_B\mathrm{\Gamma }_{\mathrm{OPE}}`$ as a function of $`m_Q`$. We use the exact relation \[Eqs. (34), (35), (38)\]
$$\mathrm{\Gamma }_B=\frac{G^2}{4\pi }\frac{m_b^2m_c^2}{M_B}_0^1\frac{\mathrm{d}x}{x}\phi _B^2(x)\frac{G^2}{4\pi }\underset{M_n>M_B}{}\frac{M_B^2M_n^2}{M_B}\left|_0^1dx\phi _n(x)\phi _B(x)\right|^2,$$
(77)
and therefore,
$$\frac{\mathrm{\Gamma }_B\mathrm{\Gamma }_{\mathrm{OPE}}}{\mathrm{\Gamma }_{\mathrm{OPE}}}=\left(_0^1\frac{\mathrm{d}x}{x}\phi _B^2(x)\right)^1\underset{n}{}\frac{M_n^2M_B^2}{m_b^2m_c^2}\left|_0^1dx\phi _n(x)\phi _B(x)\right|^2\theta (M_nM_B).$$
(78)
The summation runs over all final excited states kinematically forbidden in the decay. Once again, the sum converges rapidly and is dominated by the lowest couple of states.
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# Assisted inflation in Friedmann-Robertson-Walker and Bianchi spacetimes
## 1 Introduction
In many inflationary models the effective potential energy density of a scalar field is responsible for an epoch of accelerated inflationary expansion . Very often one assumes that inflation is driven by a scalar field of the Liouville form, i.e., a exponential potential, because this kind of potential arises in various higher-dimensional supergravity and superstring models .
Although there are many scalar fields in superstring theories, in the past it was often assumed that typically only one scalar field was responsible for the inflation, while those having higher exponents were quickly redshifted away. However, it has been found that the so-called assisted inflation may occur when several scalar fields are present, even if each individual field is too steep to drive the inflation, provided that the fields are uncoupled and interact only through the geometry. On the other hand, if the fields interact directly with each other, the opposite effect may happen and the presence of cross couplings beteween fields may hinder inflation .
Assisted inflation has been mainly studied in power-law solutions ($`at^p`$) for the spatially flat Friedmann-Robertson-Walker (FRW) cosmology, which can be shown to be the late-time attractor for the evolution of this kind of model. (Recently, Green and Lidsey discussed in the context of assisted inflation the late-time evolution in a general geometry.)
The purpose of this work is to extend previous studies on multi-scalar field cosmologies in two directions: firstly, we will use general solutions (instead of the special power-law ones) to analyze if the presence of several fields generically helps or impedes inflation, and secondly, going beyond the aforementioned FRW cosmology, we will consider the anisotropic inhomogeneous generalization given by Bianchi type I models. Thus in section 2.1 we deal with $`n`$ interacting scalar fields in a FRW spacetime and use the general solution to show that the larger is the number of interacting scalar fields the less likely is inflation. Section 2.2 makes plausible for more general scalar field potentials the results obtained in section 2.1 for exponential potentials. We do not know the solution for $`n`$ non-interacting scalar fields in a FRW cosmology in the general case, but we use the discussion of the late-time attractor of section 2.3 to restrict the analysis of uncoupled fields in section 2.4 to the particular case in which all fields are assumed to be equal, for which the general solution can be found. We show that non-interacting fields generically cooperate to assist inflation. The density fluctuations corresponding to the last case are discussed in section 2.5. General solutions of anisotropic Bianchi I cosmologies with interacting and uncoupled fields are used in the first part of section 3.1 and subsection 3.2, respectively, to check that also in those cases interacting fields make inflation more difficult while uncoupled fields assit it. The stability of power-law solutions is dicussed in 3.3. Finally in section 4 we turn to power-law solutions of the Bianchi VI<sub>0</sub> model, reaching again the same conclusions.
## 2 The $`n`$-scalar field problem in a flat FRW spacetime
In the following we will consider two kinds of problems in flat FRW spacetimes in which there are $`n`$ homogeneous scalar fields driven by exponential potentials. First of all, we will assume that the scalar fields are interacting through a product of exponential potentials. Then we will consider the case in which the scalar fields are uncoupled because the potential is a sum of potentials involving a single field.
### 2.1 The interacting $`n`$-scalar field problem in flat FRW spacetime
The problem of $`n`$ interacting homogeneous scalar fields, $`\varphi _i`$, driven by a product of $`n`$ exponential potentials $`V_i=V_{0i}\mathrm{e}^{k_i\varphi _i}`$ minimally coupled to gravity in a flat Robertson-Walker spacetime, with metric
$$ds^2=dt^2+a^2(t)\left[dx^2+dy^2+dz^2\right],$$
(1)
is formulated by the system of equations
$$3H^2=\frac{1}{2}\dot{\varphi }^2+V,$$
(2)
$$\ddot{\stackrel{}{\varphi }}+3H\dot{\stackrel{}{\varphi }}V\stackrel{}{k}=0,$$
(3)
where $`H=\dot{a}/a`$ and the potential
$$V(\stackrel{}{\varphi })=V_0\mathrm{e}^{\stackrel{}{k}\stackrel{}{\varphi }},$$
(4)
allows for interactions between the fields. $`V_0`$ is the constant $`V_{01}V_{02}\mathrm{}V_{0n}`$, and $`\stackrel{}{k}=(k_1,k_2,\mathrm{},k_n)`$ is an $`n`$ -component constant vector with respect to an orthonormal basis in the $`n`$ -dimensional Euclidean internal space to which the vector $`\stackrel{}{\varphi }=(\varphi _1,\varphi _2,\mathrm{},\varphi _n)`$, built with the $`n`$ minimally coupled scalar fields, also belongs. In the remaining of this section $`k`$, $`k^2`$ and $`\dot{\varphi }^2`$ stand for $`\left|\stackrel{}{k}\right|`$, $`\stackrel{}{k}\stackrel{}{k}`$, and $`\dot{\stackrel{}{\varphi }}\dot{\stackrel{}{\varphi }}`$ respectively. Potentials of this type are of interest because they may be considered just as an approximation to a more complex potential. In fact, in higher-dimensional superstring theories, the scalar field is like one of the matter fields that contribute to the action and effective potential of the theory. Loop expansion , or expansion in the number of interacting particles , of the action leads to a perturbative expression of the potential which is a summation of exponential terms . From (2)–(3), and discarding the trivial static metric solution $`H=0`$, $`\frac{1}{2}\dot{\varphi }^2+V=0`$ , we get
$$\dot{H}=\frac{1}{2}\dot{\varphi }^2.$$
(5)
Using this equation and the system (2)–(3) one finds that the Klein-Gordon equations (3) have the first integrals
$$\dot{\stackrel{}{\varphi }}=H\stackrel{}{k}+\frac{\stackrel{}{c}}{a^3},$$
(6)
where $`\stackrel{}{c}=(c_1,c_2,\mathrm{},c_n)`$ is an arbitrary vector integration constant. As (6) involves only geometrical quantities the Einstein-Klein-Gordon equations uncouple and the general parametric solution can be obtained .
One can easily verify that in this kind of models the Einstein-Klein-Gordon equations have power-law solutions
$$at^{\frac{2}{k^2}},$$
(7)
so that they inflate at all times when $`k^2<2`$. We will show below that the general solution has the same kind of behavior when the scale factor $`a`$ is large enough.
Inserting (6) in (5) we obtain the following second-order equation for the scale factor $`a(t)`$
$$\ddot{s}+s^m\dot{s}+\frac{1}{4\mathrm{cos}^2\sigma }s^{2m+1}=0,(\sigma \pi /2),$$
(8)
where
$$\mathrm{cos}\sigma =\frac{\stackrel{}{c}\stackrel{}{k}}{ck},$$
(9)
$`m=6/k^2<0`$, $`c=\left|\stackrel{}{c}\right|`$, $`a`$ and $`t`$ have been replaced by the new variables $`s`$ and $`\tau `$ defined by
$$a=s^{\frac{m}{3}},\tau =ckt\mathrm{cos}\sigma ,$$
(10)
and the dot stands for derivation with respect to $`\tau `$. Once $`s(\tau )`$ is known one can compute, in principle, the scale factor $`a(\tau )`$ from (10 ), and the fields $`\stackrel{}{\varphi }(t)`$ from equations (6).
Equation (8) is a particular case of the second order nonlinear ordinary differential equation
$$\ddot{s}+\alpha f(s)\dot{s}+\beta f(s)f(s)𝑑s+\gamma f(s)=0,$$
(11)
where $`f(s)`$ is some real function and $`\alpha `$, $`\beta `$ and $`\gamma `$ are constant parameters. Depending on the values of $`k^2`$ our equation (8) corresponds to the next two cases:
1. $`k^26`$ and
$$f(s)=s^m,\alpha =1,\beta =\frac{m+1}{4\mathrm{cos}^2\sigma },\gamma =0.$$
(12)
2. $`k^2=6`$ and
$$f(s)=\frac{1}{s},\alpha =1,\beta =0,\gamma =\frac{1}{4\mathrm{cos}^2\sigma }.$$
(13)
Inserting the first integrals (6) in the Einstein equation (2 ), we get a quadratic equation in the expansion rate $`H`$ which has real solutions only when its discriminant is non-negative. This condition leads to
$$\beta \frac{1}{4\left[1+\frac{2a^6V}{c^2}\right]}<\frac{1}{4},$$
(14)
discarding the oscillatory solutions of equation (11), which would be obtained for a negative potential or $`\beta >1/4`$. It follows that, if $`k^26`$, real solutions exist always for
$$k^2<k_0^2=\frac{6}{\mathrm{sin}^2\sigma }.$$
(15)
This sets a restriction on the integration constants and the exponent of the exponential-potential.
The general solution of (11) can be obtained making the nonlocal transformation of variables
$$z=f(s)𝑑s,\eta =f(s)𝑑\tau .$$
(16)
Under this transformation (8) becomes a linear inhomogeneous ordinary differential equation with constant coefficients
$$z^{\prime \prime }+\alpha z^{}+\beta z+\gamma =0,$$
(17)
which for $`\alpha >0`$ and $`\beta >0`$ is the equation of a damped harmonic oscillator in a constant external field. Here the indicates differentiation with respect to $`\eta `$. In our case the nonlocal transformation of variables (16) is
$$z=\frac{s^{m+1}}{m+1},\eta =s^m𝑑\tau ,$$
(18)
for $`k^26`$, and
$$z=\mathrm{ln}s,\eta =\frac{d\tau }{s},$$
(19)
for $`k^2=6`$.
Taking $`\alpha =1`$ in (17) we get
$$z=\{\begin{array}{cc}b_1\mathrm{e}^{\lambda _+\eta }+b_2\mathrm{e}^{\lambda _{}\eta },\hfill & \text{for }k^26;\hfill \\ b_1+b_2\mathrm{e}^\eta \gamma \eta ,\hfill & \text{for }k^2=6;\hfill \end{array}$$
(20)
where
$$\lambda _\pm =\frac{1\pm \sqrt{14\beta }}{2},$$
(21)
and $`b_1`$, $`b_2`$ are real integration constants.
We see from (10) and (18) that
$$az^{\frac{m}{3(m+1)}},$$
(22)
so that, while for $`k^2>6`$ the scale factor goes to $`0`$ as $`z0`$, one has that for $`k^2<6`$ the scale factor goes to infinity when $`z0`$. On account of equations (20) one may assume that the latter occurs as $`\eta \eta _0`$, for an appropriate value of $`\eta _0`$. Let us now write $`\eta =\eta _0+\delta \eta `$ and expand (20) and (22) about $`z=0`$ and $`\eta =\eta _0`$, keeping only terms linear in $`\delta \eta `$, so that $`\delta z=\delta \eta `$ and
$$a=\delta z^{\frac{m}{3(m+1)}}=\delta \eta ^{\frac{m}{3(m+1)}}.$$
(23)
Now, we see from (10) and (18) that
$$\delta \eta =a^3\delta \tau a^3\delta t.$$
(24)
If we use this in (23), we get successively
$$aa^{\frac{m}{m+1}}\delta t^{\frac{m}{3(m+1)}},$$
(25)
and
$$a\delta t^{\frac{m}{3}}=\delta t^{\frac{2}{k^2}}.$$
(26)
We see then that if $`k^2<2`$ the solution does inflate at least along some time interval. Now, since $`k^2=k_1^2+k_2^2+\mathrm{}+k_n^2`$, one concludes from our analysis of the general solution that the larger is the number of the interacting scalar fields the less likely will be $`k^2<2`$ as well that inflation take place. On the other hand, for $`6<k^2<k_0^2`$, one has $`0<\beta <1/4`$ that implies that both $`\lambda _\pm <0`$; thus (18) and ( 20) tell us that the scale factor goes to infinity when $`z\mathrm{}`$ and $`\eta \mathrm{}`$ (with no loss of generality we are taking both constants $`b_1,b_2>0`$). By keeping only the dominant terms we have for very large negative $`\eta `$ that
$$a\mathrm{e}^{\frac{m\lambda _{}\eta }{3(m+1)}}.$$
(27)
This relation and (24), that is generally valid, yield
$$\frac{\delta a}{a}a^3\delta t$$
(28)
whose general solution is $`at^{1/3}`$. In this case the solution does not inflate and it approaches to the free scalar field solution. These results are in accordance with those obtained with particular solutions by other authors .
It would be surprising if for $`k^2=6`$ the scale factor had a behavior drastically departing from that suggested by the previous analysis when $`k^2(k_0^2,6)(6,2)`$. In fact, we see from (10) and (19) that $`z=3\mathrm{log}a`$ so that, as a consequence of (20), $`a`$ goes to infinity as $`\eta \mathrm{}`$. From (20) we see that $`\mathrm{log}z\eta `$, when $`\eta \mathrm{}`$ and, by using (24),
$$\frac{\delta a}{a\mathrm{log}a}\frac{\delta z}{z}\delta \eta \frac{\delta t}{a^3},$$
(29)
which can be written as
$$\frac{\delta a}{a}a^3\delta t\mathrm{log}a.$$
(30)
In the asymptotic regime in which $`a\mathrm{}`$ we have $`a\delta a<a^2\delta a/\mathrm{log}a<a^2\delta a`$, so that from (30) we get $`a^2\stackrel{<}{}t\stackrel{<}{}a^3`$ and, finally,
$$t^{1/3}\stackrel{<}{}a\stackrel{<}{}t^{1/2},\text{when }a\mathrm{}.$$
(31)
We see that the solution does not inflate.
### 2.2 More general potentials
The results of the above section show that we can introduce an $`n`$ -dimensional Euclidean internal vector space containing the $`n`$-component vector $`\stackrel{}{\varphi }=(\varphi _1,\varphi _2,\mathrm{},\varphi _n)`$ built with $`n`$ minimally coupled scalar fields. Let us assume now that the potential has the general form
$$V=V(\mathrm{\Phi }),\mathrm{\Phi }=\stackrel{}{k}\stackrel{}{\varphi }.$$
(32)
In this case the Einstein equation (2) remains unchanged, however the Klein-Gordon equation (3) becomes
$$\ddot{\stackrel{}{\varphi }}+3H\dot{\stackrel{}{\varphi }}+V^{}\stackrel{}{k}=0.$$
(33)
where the indicates derivatives with respect to the variable $`\mathrm{\Phi }`$. Since (2), (32) and (33) are invariant under rotations of the orthogonal axes in the Euclidean space above mentioned, we may choose the first axis of this internal space along the vector $`\stackrel{}{k}`$. Then the Klein-Gordon equation splits into one equation for $`\varphi _1=\varphi `$
$$\ddot{\varphi }+3H\dot{\varphi }+kV^{}=0,$$
(34)
and $`n1`$ free field Klein-Gordon equations for $`\varphi _2=\mathrm{}=\varphi _n=\psi `$
$$\ddot{\psi }+3H\dot{\psi }=0.$$
(35)
From (35) we obtain the first integral $`\dot{\psi }=c_0/a^3`$, where $`c_0`$ is an arbitrary integration constant. Hence the original $`n`$-scalar field problem is equivalent to considering a self-interacting scalar field with stiff matter. In fact, the Einstein equation (2) now reads
$$3H^2=\frac{1}{2}\dot{\varphi }^2+V+\frac{c_\psi ^2}{2a^6},$$
(36)
where $`c_\psi ^2`$ is the sum of $`n1`$ positive defined integration constants.
Taking into account that the scalar field $`\varphi `$ depends only on $`t`$, its energy-momentum tensor may be written in the perfect fluid form
$$T_{ik}=(p_\varphi +\rho _\varphi )u_iu_k+p_\varphi g_{ik},$$
(37)
where
$`\rho _\varphi `$ $`=`$ $`{\displaystyle \frac{1}{2}}\dot{\varphi }^2+V(\varphi ),`$
$`p_\varphi `$ $`=`$ $`{\displaystyle \frac{1}{2}}\dot{\varphi }^2V(\varphi ).`$ (38)
The fluid interpretation of the scalar field has proven very useful in the study of the inflationary and Q-matter scenarios . In particular it leads to consider its equation of state $`p_\varphi =\left(\gamma _\varphi 1\right)\rho _\varphi `$. On the other hand, the state equation for fluid representing stiff matter is $`p_f=\left(\gamma _f1\right)\rho _f`$ with $`\gamma _f=2`$. Because of the additivity of the stress-energy tensor it makes sense to consider an effective perfect fluid description with equation of state $`p=\left(\gamma 1\right)\rho `$ where $`p=p_f+p_\varphi `$, $`\rho =\rho _f+\rho _\varphi `$ and
$$\gamma =\frac{\gamma _f\rho _f+\gamma _\varphi \rho _\varphi }{\rho _f+\rho _\varphi }$$
(39)
is the overall (i.e. effective) adiabatic index. For this effective perfect fluid the dynamical equations are
$$3\frac{\dot{a}^2}{a^2}=\rho $$
(40)
and
$$\frac{\ddot{a}}{a}=\frac{1}{6}\left[\rho +3p\right],$$
(41)
where $`p`$ and $`\rho `$ are the density and pressure of the effective perfect fluid. They involve the self-interacting scalar field and the $`n`$-1 free scalar fields. Inflationary solutions occur when $`\ddot{a}>0`$; this means that the expansion is dominated by a gravitationally repulsive stress that violates the strong energy condition $`\rho +3p<0`$ or equivalently $`\gamma <2/3`$ . When we impose this condition on (39), we obtain that
$$V(\mathrm{\Phi })>\dot{\varphi }^2+\frac{c_\psi ^2}{a^6}.$$
(42)
Hence, when there is only one scalar field the solution inflates for $`V(\varphi )>\dot{\varphi }^2`$. However, the $`n`$ interacting scalar fields produce a desassisted inflation because the solutions inflate only if (42) holds. Using the effective perfect fluid description we can generalize this analysis to an arbitrary potential $`V(\varphi _1,\varphi _2,\mathrm{},\varphi _n)`$. In this case we obtain that $`V(\varphi _1,\varphi _2,\mathrm{},\varphi _n)>\dot{\varphi }_1^2+\dot{\varphi }_2^2+\mathrm{}..+\dot{\varphi }_n^2`$, so that the interacting $`n`$ scalar fields might make inflation more unlikely in a FRW spacetime.
### 2.3 The $`n`$-scalar field attractor problem in the FRW spacetime
Let us now assume that the $`n`$ homogeneous scalar fields, $`\varphi _i`$, in the FRW spacetime are driven by a general potential $`V=V(\varphi _i)`$. In that case the Einstein-Klein-Gordon equations are
$$3H^2=\frac{1}{2}\underset{i=1}{\overset{n}{}}\dot{\varphi }_i^2+V,$$
(43)
$$\ddot{\varphi _i}+3H\dot{\varphi _i}+V_{,\varphi _i}=0.$$
(44)
where $`V_{,\varphi _i}`$ stand for $`V/\varphi _i`$. From these equations we get
$$\dot{H}=\frac{1}{2}\underset{i=1}{\overset{n}{}}\dot{\varphi _i}^2.$$
(45)
In order to investigate the stable scalar field configurations, we introduce the quantity
$$\omega =\frac{_{i=1}^n\dot{\varphi _i}^2}{n\dot{\varphi }_\alpha ^2},$$
(46)
which reduces to $`\omega =1`$ for the configuration $`\varphi _1=\varphi _2=\mathrm{}=\varphi _n`$. Using (43)–(46) we find the differential equation for $`\omega `$:
$$\dot{\omega }=2\frac{nV_{,\varphi _\alpha }\dot{\varphi }_\alpha \omega \dot{V}}{n\dot{\varphi }_\alpha ^2}.$$
(47)
(In this section no summation convention applies to repeated Greek indexes.) If we further assume that the potential satisfies the condition
$$\dot{V}=nV_{,\varphi _\alpha }\dot{\varphi }_\alpha ,$$
(48)
equation (47) becomes
$$\dot{\omega }=2\frac{V_{,\varphi _\alpha }}{\dot{\varphi }_\alpha }(\omega 1)$$
(49)
which has a fixed point solution: $`\omega =1`$. Furthermore, the general solution of (49) can be found using (44):
$$\omega =1+\frac{c}{a^6\dot{\varphi }_\alpha ^2}$$
(50)
where $`c`$ is an arbitrary integration constant. However, it is useful to express the last solution in terms of geometrical quantities with the aid of (45) and (46). The final result is
$$\omega =\left(1+\frac{nc}{2a^6\dot{H}}\right)^1.$$
(51)
Evaluating (51) in the asymptotic regime, it can be easily shown that the particular solution $`\omega =1`$ is an attractor for evolutions that behave asymptotically as $`at^\nu `$ with $`\nu >1/3`$.
This result strongly suggests that the special case in which all scalar fields are equal may be the late-time attractor of more general problems. In fact, this has been proved in for potential (4) with $`k_1=\mathrm{}=k_n`$ and $`V_{01}=\mathrm{}=V_{0n}`$, which guarantee assumption (48). We did not use this fact in the previous section because we were able to use the general solution with no additional assumption on $`k_i`$, $`V_{0i}`$ and $`\varphi _i`$, but it will be useful to simplify somewhat the problem analyzed in the following section.
### 2.4 The non-interacting $`n`$-scalar field problem in the FRW spacetime
Let us now assume that the $`n`$ homogeneous scalar fields, $`\varphi _i`$, in the spacetime given by (1) do not interact directly, but are driven by a sum of $`n`$ exponential potentials $`V_i=V_{0i}\mathrm{e}^{k_i\varphi _i}`$. In that case the Einstein-Klein-Gordon equations are
$$3H^2=\underset{i=1}{\overset{n}{}}\left[\frac{1}{2}\dot{\varphi }_i^2+V_i\right],$$
(52)
$$\ddot{\varphi _i}+3H\dot{\varphi _i}k_iV_i=0.$$
(53)
From now on we will consider a simplified problem in which $`k_1=k_2=\mathrm{}=k_nk`$ and $`V_{01}=V_{02}=\mathrm{}=V_{0n}V_0`$. As discussed in the previous section, we can expect in this particular case that in the asymptotic evolution all scalar fields tend to a common limit. That this is actually the case has been proved in . In consequence, in the remaining of this section we will take $`\varphi _1=\varphi _2=\mathrm{}=\varphi _n\varphi `$, so that $`V_1=V_2=\mathrm{}=V_nV=V_0\mathrm{e}^{k\varphi }`$ and equations (52)–(53) reduce to
$$3H^2=\frac{n}{2}\dot{\varphi }^2+nV,$$
(54)
$$\ddot{\varphi }+3H\dot{\varphi }kV=0.$$
(55)
One can easily get from (54)–(55)
$$\dot{H}=\frac{n}{2}\dot{\varphi }^2$$
(56)
and the first integral of the Klein-Gordon equation (55)
$$\dot{\varphi }=\frac{k}{n}H+\frac{c}{a^3},$$
(57)
where $`c`$ is an arbitrary integration constant.
As in the previous section, the Einstein-Klein-Gordon equations have power-law solutions
$$at^{\frac{2n}{k^2}},$$
(58)
so that this type of solutions inflates at all times when $`k^2<2n`$. We will now show that the general solution of the Einstein-Klein-Gordon equations ( 54) and (55) has the same kind of behavior when the scale factor $`a`$ is large enough.
Inserting (57) in (56) we obtain the following second-order equation for the scale factor $`a(t)`$
$$\ddot{s}+s^m\dot{s}+\frac{1}{4}s^{2m+1}=0,$$
(59)
where $`m=6n/k^2<0`$, the dot means derivative with respect to $`\tau `$ and we have used, instead of $`a`$ and $`t`$, the new variables $`s`$ and $`\tau `$ defined by
$$a=s^{\frac{m}{3}},\tau =ckt.$$
(60)
We have again a particular case of equation (11) with $`\alpha =1`$ and to find the general exact solution of (59) one has to consider the following two possibilities:
1. $`k^26n`$ and
$$f(s)=s^m,\beta =\frac{m+1}{4},\gamma =0.$$
(61)
2. $`k^2=6n`$ and
$$f(s)=\frac{1}{s},\beta =0,\gamma =\frac{1}{4}.$$
(62)
The general solution of (59) can be obtained by performing the nonlocal transformation of variables given in (18)–(19) with the new value of $`m`$, which reduces (59) to the linear inhomogeneous ordinary differential equation with constant coefficients (17).
If one now repeats the analysis of the previous section by systematically taking $`\sigma =0`$ and using the new $`m`$, the same final results for $`z`$ and $`a`$ are obtained: just replace $`k^2=k_1^2+\mathrm{}k_n^2`$ by $`k^2/n`$ in the quantities involved in (20) and (22). One readily concludes that, when $`a\mathrm{}`$, these models inflate if $`k^2<2n`$, so that the fields cooperate to make inflation more likely in the so-called “assisted inflation”, which was first discussed —but only for power-law solutions— in .
This result could have been anticipated since the present case is included in the mathematical problem set in section 2.1 by equations (8) and (9) by just taking $`\mathrm{cos}\sigma =\pm 1`$ (that corresponds to the simplified problem under consideration) and $`m=6n/k^2`$.
### 2.5 Density fluctuations
The fact that the contributions of the density fluctuations differ significantly in different inflationary universe models, has motivated a detailed study of all the alternatives. In this context, it is interesting to derive the spectral indices for the perturbations that would be created during the periods of inflation described by the solutions given in the last section. It is well known that, for multi-scalar field models, the spectrum of the curvature perturbation reads ,
$$P_S=\left(\frac{H}{2\pi }\right)^2\frac{N}{\varphi _i}\frac{N}{\varphi _j}\delta _{ij},$$
(63)
where $`N`$ is the number of $`e`$-foldings of inflationary expansion remaining, and there is a summation over $`i`$ and $`j.`$ In the case considered in the previous section,where all the scalar fields are equal, (63) yields
$$P_S(\stackrel{~}{k})=\left(\frac{H}{2\pi }\right)^2\frac{1}{n}\frac{H^2}{\dot{\varphi }^2}_{aH=\stackrel{~}{k}}.$$
(64)
where $`H`$ and $`\dot{\varphi }`$ have to be evaluated at the time when the wave number of interest $`\stackrel{~}{k}`$ leaves the horizon during inflation. Also in this case, the spectral index $`n_S(\stackrel{~}{k})`$ defined as
$$n_S(\stackrel{~}{k})=1+\frac{\mathrm{d}\mathrm{ln}P_S}{\mathrm{d}\mathrm{ln}\stackrel{~}{k}}$$
is given by
$$1n_S=2\frac{\dot{H}}{H^2}.$$
(65)
The availability of exact solutions allows us to express the relevant quantities as functions of the variable $`\eta `$ introduced in equation (16). The scale factor reads (see (22), which is satisfied for uncoupled scalar fields with $`m=6n/k^2`$, as one can easily see)
$$a(\eta )=\left[\left(m+1\right)z(\eta )\right]^{\frac{m}{3\left(m+1\right)}}$$
(66)
where $`z(\eta )`$ is given by (20)–(21). The spectrum $`P_S(\stackrel{~}{k})`$ obtained from (64) is shown in Figure 1, for $`k=3`$, $`c=1`$, $`b_1=1`$, $`b_2=0.001`$ and different values of $`n`$. Here, inflation is possible for $`n5`$, and one can see that the peak of the spectral distribution moves towards the high frequency region as $`n`$ increases. The corresponding spectral index $`n_S`$ is shown in Figure 2. From (65) and the general solutions (20), (21) and (66), it can be shown that the value of $`n_S1`$ in the asymptotic region $`a1`$ as the number $`n`$ of present fields increases. This feature is exemplified in Figure 2, where we can see how the larger is the value of $`n`$, the closer is the spectrum to the scale invariance .
## 3 The $`n`$-scalar field problem in Bianchi type I models
Now we turn to the general Bianchi type I model with $`n`$ homogeneous scalar fields driven by exponential potentials. As we proceeded in the case of FRW spacetimes, we will first assume that the scalar fields are interacting through a product of exponential potentials, an then we will consider the case in which the scalar fields are uncoupled because the potential is a sum of potentials involving a single field.
### 3.1 The interacting $`n`$-scalar field problem in the anisotropic Bianchi type I model
The general Bianchi type I model is the anisotropic generalization of the spatially flat FRW universe expanding differently in the $`x`$, $`y`$, and $`z`$ directions. In the usual synchronous form its line element is given by
$$ds^2=dT^2+a_1^2(T)dx^2+a_2^2(T)dy^2+a_3^2(T)dz^2.$$
(67)
For convenience we use the semiconformal coordinates
$$dt\frac{dT}{a_3},\mathrm{e}^fa_3^2,Ga_1a_2,\mathrm{e}^p\frac{a_1}{a_2},$$
(68)
to cast the metric (67) into the form
$$ds^2=\mathrm{e}^{f(t)}\left(dt^2+dz^2\right)+G(t)\left(\mathrm{e}^{p(t)}dx^2+\mathrm{e}^{p(t)}dy^2\right).$$
(69)
We first consider, as in section 2.1, $`n`$ scalar fields $`\varphi _i`$ interacting directly through the exponential potential (4). The problem of $`n`$ interacting homogeneous scalar fields, $`\varphi _i`$, driven by a product of $`n`$ exponential potentials $`V_i=V_{0i}e^{k_i\varphi _i}`$, minimally coupled to gravity in the Bianchi I spacetime (69), is formulated by the following system of Einstein-Klein-Gordon equations
$$\dot{p}=\frac{a}{G},$$
(70)
$$\mathrm{e}^f=\frac{\ddot{G}}{2VG},$$
(71)
$$\frac{\ddot{G}}{G}\frac{1}{2}\left(\frac{\dot{G}}{G}\right)^2\frac{\dot{G}}{G}\dot{f}+\frac{1}{2}\dot{p}^2=\dot{\varphi }^2,$$
(72)
$$\ddot{\stackrel{}{\varphi }}+\frac{\dot{G}}{G}\dot{\stackrel{}{\varphi }}\mathrm{e}^fV\stackrel{}{k}=0,$$
(73)
where $`a`$ is an arbitrary integration constant. It can be easily seen that the vector
$$\dot{\stackrel{}{\varphi }}=\frac{\dot{G}}{G}\frac{\stackrel{}{k}}{2}+\frac{\stackrel{}{m}}{G},$$
(74)
(where $`\stackrel{}{m}`$ is a $`n`$-dimensional vector whose components are integration constants) is a first integral of the Klein-Gordon equation set (73). Inserting (70) in equation (74) the general solution of the Klein-Gordon equations is found:
$$\stackrel{}{\varphi }=\stackrel{}{\varphi }_0+p\frac{\stackrel{}{m}}{a}+\frac{\stackrel{}{k}}{2}\mathrm{ln}G,$$
(75)
where $`\stackrel{}{\varphi }_0`$ is an arbitrary constant vector. Equations (70 )–(72) along with (75) uncouple and their solutions can be obtained if one is able to solve the following third-order equation for $`G`$
$$G\ddot{G}^2\stackrel{\mathrm{}}{G}\dot{G}G+\left(\frac{1}{2}\frac{k^2}{4}\right)\ddot{G}\dot{G}^2+\left(m^2+\frac{a^2}{2}\right)\ddot{G}=0.$$
(76)
Once $`G(t)`$ is known, in principle one can compute $`p(t)`$ and $`\stackrel{}{\varphi }(t)`$ from equations (70) and (75), respectively; $`f(t)`$ is then obtained from (71).
The Einstein-Klein-Gordon equations admit power-law solutions, $`G=t^\alpha `$, but they happen to be isotropic and, thus, equal to the ones discussed in section 2.1. Thus we will analyze, instead, the general solution of (76).
Equation (76) has the first integral
$$G\frac{\ddot{G}}{\dot{G}}+(K1)\dot{G}+\frac{M^2}{\dot{G}}=C,$$
(77)
where $`C`$ is an arbitrary constant and
$$K\frac{k^2}{4}\frac{1}{2},M^2m^2+\frac{a^2}{2}.$$
(78)
If instead of $`t`$ and $`G`$ we use the new variables $`z`$ and $`\tau `$ defined, for $`C0`$, in
$$G=z^{1/K},t=\frac{\tau }{C},$$
(79)
then, equation (77) becomes
$$z^{\prime \prime }+z^{1/K}z^{}+\frac{KM^2}{C^2}z^{12/K}=0,$$
(80)
where a prime denotes the derivative with respect to $`\tau `$. This equation is, once more, a particular case of (11) and can be linearized by using the non-local transformation (16), which in this case is
$$yz^{1/K}𝑑z=K\frac{z^{11/K}}{K1},\eta z^{1/K}𝑑\tau =\frac{C}{a}p$$
(81)
for $`K1`$, and
$$yz^1𝑑z=\mathrm{ln}z,\eta z^1𝑑\tau =\frac{C}{a}p,$$
(82)
for $`K=1`$. If we take
$$\beta (K1)\frac{M^2}{C^2},\gamma =0,$$
(83)
for $`K1`$ and
$$\beta =0,\gamma =\frac{M^2}{C^2},$$
(84)
for $`K=1`$, equation (80) reduces to two particular cases of equation (17) for $`\alpha =1`$. The trivial solution of this equation gives the implicit general solution of (76) which can be written, for arbitrary $`a`$, $`M`$ and non-vanishing $`C`$, as
$$G=\left[\mathrm{e}^{\eta /2}\left(C_1\mathrm{e}^{\lambda \eta }+C_2\mathrm{e}^{\lambda \eta }\right)\right]^{\frac{1}{K1}}$$
(85)
for $`K1`$, and as
$$G=C_1e^{\gamma \eta +C_2e^\eta }$$
(86)
for $`K=1`$. $`C_1`$ and $`C_2`$ are integration constants and $`\lambda =\sqrt{14\beta }/2`$.
To check whether a model inflates, we will look at the sign of the deceleration parameter $`q=\theta ^2\left(3\dot{\theta }+\theta ^2\right)`$, where $`\theta =u_{;a}^a`$ is the expansion and $`\dot{\theta }=\theta _{,a}u^a`$, $`u^a`$ being the four-velocity of the cosmic fluid. Since in this case we are dealing with comoving coordinates, $`u^a=(e^{f/2},0,0,0)`$, one can see that, apart from a positive factor, the deceleration is
$$q9K\dot{G}^23(2km+C)+(kmC)^2+9M.$$
(87)
We see from (85) that when $`K<1`$ (i.e., when $`k^2<6`$) $`G`$ blows up for some value $`\eta =\eta _0=\frac{1}{2\lambda }\mathrm{log}\left(C_2/C_1\right)`$, provided that $`C_1C_2<0`$. If we expand (85) around this value ($`\eta =\eta _0+\delta \eta `$) we get
$$G\delta \eta ^{\frac{1}{K1}},$$
(88)
and from (79) and (81)
$$d\eta =\frac{d\tau }{G}=\frac{C}{G}dt,$$
(89)
so that
$$\dot{G}=\frac{C}{G}G^{}\frac{1}{K1}\delta \eta ^1.$$
(90)
In consequence, when $`G\mathrm{}`$ the deceleration parameter (87) is
$$q9K\frac{C^2}{(K1)^2}\frac{1}{\delta \eta ^2},(\text{when }\delta \eta 0),$$
(91)
and there is inflation if $`K<0`$, i.e., if $`k^2=k_1^2+k_2^2+\mathrm{}+k_n^2<2`$. We conclude that in these anisotropic universes also a greater number of interacting scalar fields makes inflation less likely.
### 3.2 The non-interacting $`n`$-scalar field problem in the Bianchi type I model
We will now assume that the $`n`$ homogeneous scalar fields $`\varphi _i`$ in the metric (67) do not interact directly, but are driven by a sum of $`n`$ exponential potentials $`V_i=V_{0i}\mathrm{e}^{k_i\varphi _i}`$. To simplify the task of finding exact solutions of the Einstein-Klein-Gordon equations, we will further assume that $`k_1=k_2=\mathrm{}=k_nk`$, $`\varphi _1=\varphi _2=\mathrm{}=\varphi _n\varphi `$ and $`V_1=V_2=\mathrm{}=V_nV=V_0\mathrm{e}^{k\varphi }`$, so that aforementioned equations can be written as
$$\dot{p}=\frac{a}{G},$$
(92)
$$\mathrm{e}^f=\frac{\ddot{G}}{2nVG},$$
(93)
$$\frac{\ddot{G}}{G}\frac{1}{2}\left(\frac{\dot{G}}{G}\right)^2\frac{\dot{G}}{G}\dot{f}+\frac{1}{2}\dot{p}^2=n\dot{\varphi }^2,$$
(94)
$$\ddot{\varphi }+\frac{\dot{G}}{G}\dot{\varphi }k\mathrm{e}^fV=0,$$
(95)
where $`a`$ is an arbitrary integration constant. It is easy to check that
$$\dot{\varphi }=\frac{k}{2n}\frac{\dot{G}}{G}+\frac{m}{G}$$
(96)
is a first integral of the Klein-Gordon equation (95), in terms of the new integration constant $`m`$. Inserting (92) in equation ( 96) the general solution of the Klein-Gordon equations is found:
$$\varphi =\varphi _0+p\frac{m}{a}\frac{k}{2n}\mathrm{ln}G,$$
(97)
where $`\varphi _0`$ is an arbitrary constant. The equations (92)–( 94) along with (97) uncouple and their solutions can be obtained if one is able to solve the following third-order equation for $`G`$
$$G\ddot{G}^2\stackrel{\mathrm{}}{G}\dot{G}G+\left(\frac{1}{2}\frac{k^2}{4n}\right)\ddot{G}\dot{G}^2+\left(m^2+\frac{a^2}{2}\right)\ddot{G}=0.$$
(98)
Since this equation is the same as (76) once one replaces $`k^2=k_1^2+k_2^2+\mathrm{}+k_n^2`$ by $`k^2/n`$, we may repeat the calculations of the previous section to reach the opposite conclusion: if several non-interacting scalar fields are present, they will cooperate to “assist” the inflation, which will be more likely and occurs for $`k^2<2n`$.
### 3.3 Stability of power-law solutions in Bianchi type I model
For many purposes it is interesting to investigate the stability of the solutions of (76). In particular, we hope that the solution representing an accelerated expansion of the universe, and the solutions that correspond to the assisted inflation, be stable. To this end we introduce the variable
$$\mathrm{\Omega }=\frac{\dot{h}}{h^2},$$
(99)
where $`h=\frac{\dot{G}}{G}`$, in equation (76):
$$\dot{\mathrm{\Omega }}+\left[\mathrm{\Omega }+K\frac{M^2}{h^2G^2}\right](\mathrm{\Omega }+1)h=0.$$
(100)
This equation has the fixed point solution $`\mathrm{\Omega }=1`$. Note that equation (100) has also the fixed point solution $`\mathrm{\Omega }=K`$ if $`\dot{G}\mathrm{}`$ asymptotically. The corresponding asymptotic limits of these solutions can be obtained by solving (99) for them. The final result is $`Gt`$ and $`Gt^{1/K}`$ respectively. Let us investigate the stability of these solutions when $`G`$ blows up. From (100) it is easy to see that $`\mathrm{\Omega }=1`$ is unstable because expanding the solutions about it, $`\mathrm{\Omega }=1+ϵ`$ with $`ϵ1`$, the sign of $`\dot{ϵ}`$ depend on the slope of potential and the initial conditions. In fact, the corresponding solution $`Gt`$ does not satisfies Einstein equations (cf. (93)) and was introduced when multiplying (94) with $`G^2\ddot{G}`$ to obtain (98). On the other hand, the asymptotic solution $`\mathrm{\Omega }=K`$ is stable because the dynamical equation for the perturbation $`ϵ`$
$$\dot{ϵ}=\frac{1K}{Kt}ϵ$$
(101)
near the attractor indicates that $`ϵ`$ decreases for $`K<1`$. In particular, the inflationary solutions, that occur for $`K<0`$, are stable.
## 4 The $`n`$-scalar field problem in a Bianchi VI<sub>0</sub> model
The Bianchi VI<sub>0</sub> model can be written as follows:
$$ds^2=\mathrm{e}^{f(t)}\left(dt^2+dz^2\right)+G(t)\left(\mathrm{e}^zdx^2+\mathrm{e}^zdy^2\right).$$
(102)
We first consider, as in previous sections, $`n`$ scalar fields $`\varphi _i`$ interacting through the exponential potential (4). The corresponding Einstein-Klein-Gordon equations are
$$\mathrm{e}^f=\frac{\ddot{G}}{2VG},$$
(103)
$$\frac{\ddot{G}}{G}\frac{1}{2}\left(\frac{\dot{G}}{G}\right)^2\frac{\dot{G}}{G}\dot{f}+\frac{1}{2}=\dot{\varphi }^2,$$
(104)
$$\ddot{\stackrel{}{\varphi }}+\frac{\dot{G}}{G}\dot{\stackrel{}{\varphi }}\mathrm{e}^fV\stackrel{}{k}=0.$$
(105)
As formerly, one can check that the vector
$$\dot{\stackrel{}{\varphi }}=\frac{\dot{G}}{G}\frac{\stackrel{}{k}}{2}+\frac{\stackrel{}{m}}{G},$$
(106)
where $`\stackrel{}{m}`$ is a $`n`$-dimensional arbitrary constant vector, is a first integral of the Klein-Gordon equation set (105). By using this result and the value of $`\dot{f}`$ one obtains from (103), we get
$$G\ddot{G}^2\stackrel{\mathrm{}}{G}\dot{G}G+\left(\frac{1}{2}\frac{k^2}{4}\right)\ddot{G}\dot{G}^2+\frac{1}{2}\ddot{G}G^2+m^2\ddot{G}=0.$$
(107)
The solutions of this equation are not known. However, investigating the stability of its fixed points, the asymptotic behavior of the general solution can be obtained in a simple way. In order to see whether assisted inflation works in Bianchi type VI<sub>0</sub> metrics it is sufficient to analyze the special case $`m^2=0`$. In terms of the variable $`\mathrm{\Omega }`$ defined by (99), equation (107) becomes
$$\dot{\mathrm{\Omega }}+\left[\mathrm{\Omega }+K\frac{1}{2h^2}\right](\mathrm{\Omega }+1)h=0.$$
(108)
This equation has three fixed points: $`\mathrm{\Omega }_1=1`$, $`\mathrm{\Omega }_2=K`$ if $`h\mathrm{}`$ asymptotically, and $`\mathrm{\Omega }_3=0`$, which correspond to $`Gt`$, $`Gt^{1/K}`$ and $`G\mathrm{e}^{t/\sqrt{2K}}`$ respectively. Now, we investigate the stability of these solutions when $`G`$ blows up. Expanding the solution about fixed points, that is, making $`\mathrm{\Omega }=\mathrm{\Omega }_{1,2,3}+ϵ`$ with $`ϵ1`$, we get
$$\dot{ϵ}=\frac{1}{2h}ϵ$$
(109)
for $`\mathrm{\Omega }_1`$, which shows it is unstable, as in the case discussed in the previous section. On the other hand, one obtains equation (101) for the linear approximation around $`\mathrm{\Omega }_2`$, and
$$\dot{ϵ}=\sqrt{\frac{1}{2K}}ϵ$$
(110)
in the case of $`\mathrm{\Omega }_3`$. We conclude that the asymptotic solution $`\mathrm{\Omega }_2`$ is stable for $`K<0`$, which means $`k^2<2`$ (for this set of potential slopes we have inflation), and $`\mathrm{\Omega }_3`$ is stable for $`K>0`$. Note that $`Gt^{1/K}`$ is only an asymptotic solution of (107), with $`m=0`$; however, it acts as an attractor for all solutions that are close to it.
In this special case in which $`m^2=0`$, one can readily see that the the deceleration parameter for the solutions $`Gt^{1/K}`$ is (after recovering the implicit absolute value around $`t`$)
$$q\frac{K}{1K}|t|^{21/K},$$
(111)
which is negative for $`1/2K<0`$. Again, the more interacting scalar fields the less likely is inflation, which ensues when $`k^2=k_1^2+k_2^2+\mathrm{}+k_n^2<2`$.
If one assumes now that the scalar fields do not interact directly and one sets all the fields equal, as in sections 2.4 and 3.2, it is easily seen that the attracting power-law solutions inflate when $`k^2<2n`$, so that the scalar fields cooperate to assist inflation.
## 5 Conclusions
We have studied the effects of the appearance of more than one scalar field both in FRW spacetimes and in anisotropic Bianchi I cosmologies. Instead of using the important but particular power-law solutions, we have taken advantage of the general solution to analyze the generic behavior in FRW. In Bianchi I the power-law solutions are isotropic and, thus, very particular, so that the use of the general solution is even more illustrative. In all cases we have found, in agreement with calculations made by other authors with power-law solutions in FRW, that the existence of more than one scalar field assists inflation provided that they are uncoupled and interact only through expansion. Also, in this case, the spectrum of density perturbations becomes closer to scale invariance as the number of fields increases. If, on the contrary, the fields interact directly with each other, inflation is less likely to occur. The same behavior has been obtained in Bianchi VI<sub>0</sub> universes, but (not having available the general solution) only for power-law solutions, which have been shown to be attractors when inflation arises. These results reinforce our belief that the presence of several uncoupled scalar fields in more general cosmologies fosters inflation, but that, in contradistinction, mutually interacting scalar fields tend to hinder the inflationary process.
## Acknowledgments
This work was partially supported by the University of the Basque Country through the Research Project UPV172.310-EB150/98 and the General Research Grant UPV172.310-G02/99, as well as by the University of Buenos Aires under Project TX93. We thank the anonymous referee for his helpful comments.
## Figures
Figure 1. Spectrum of the curvature perturbations for $`k=3`$, $`c=1`$, $`b_1=1`$ and $`b_2=0.001`$. As $`n`$ increases (inflation is possible when $`n5`$) the spectral distribution peak is shifted to high frequencies.
Figure 2. The spectral index $`n_S(\stackrel{~}{k})`$ for the general solution approaches $`1`$ in the asymptotic region $`a1`$ as the number of fields increases.
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# Trapping of single atoms with single photons in cavity QED
## I Introduction
An exciting advance in recent years has been the increasing ability to observe and manipulate the dynamical processes of individual quantum systems. In this endeavor, an important physical system has been a single atom strongly coupled to the electromagnetic field of a high-$`Q`$ (optical or microwave) cavity within the setting of cavity quantum electrodynamics (cavity QED). Here the coupling frequency of one atom to a single mode of an optical resonator is denoted by $`g_0`$ (i.e., $`2g_0`$ is the one-photon Rabi frequency), with the regime of strong coupling defined by the requirement that $`g_0(\gamma ,\kappa )`$, where $`\gamma `$ is the atomic decay rate to modes other than the cavity mode and $`\kappa `$ is the decay rate of the cavity mode itself. In this circumstance, the number of photons required to saturate an intracavity atom is $`n_0\gamma ^2/g_0^21`$ and the number of atoms required to have an appreciable effect on the intracavity field is $`N_0\kappa \gamma /g_0^21`$ .
Although there have been numerous laboratory advances which demonstrate the effect of strong coupling on the internal degrees of freedom of an atomic dipole coupled to the quantized cavity field (i.e., $`g_0`$ $`\kappa ,`$ $`\gamma `$), the consequences of strong coupling for the external, atomic center-of-mass motion with kinetic energy $`E_k`$ have only recently been explored experimentally . In a regime of strong coupling for the external degrees of freedom, $`g_0>E_k/\mathrm{}`$, a single quantum is sufficient to profoundly alter the atomic center-of-mass (CM) motion, as an atom moves through a region of spatially varying coupling coefficient $`g(\stackrel{}{r})=g_0\psi (\stackrel{}{r})`$ \[e.g., as arises in the Gaussian mode of a Fabry-Perot cavity, $`\psi (\stackrel{}{r})`$\].
Perhaps most strikingly, the spatial variation of the cavity mode can lead to a confining potential sufficient to trap an atom within the cavity mode even for a single quantum of excitation of the atom-cavity system, as first discussed in the work of Refs.. This is illustrated in Fig. 1, which shows the the possibility for trapping by excitation to the lower component $`|`$ in the Jaynes-Cummings manifold of eigenstates. Modifications of the atomic CM dynamics can in turn significantly alter the cavity field. This situation is very different from the usual case for trapped atoms or ions in fixed external potentials, in that here the confining field and the atomic motion can be strongly interacting, in which case the overall state of the system must be determined in a self-consistent fashion.
The experimental requirements to investigate strong coupling for both the internal and external degrees of freedom are stringent \[namely, $`g>(E_k/\mathrm{},\gamma ,\kappa )`$\], and have required the integration of the techniques of laser cooling and trapping with those of cavity QED, as was initially achieved in 1996 and as illustrated in Fig. 2. Mechanical effects due to strong coupling with single quanta were first observed in 1998 , in an experiment with peak coupling energy $`\mathrm{}g_05`$ mK and with initial atomic kinetic energy $`E_k400\mu `$ K.
Following this theme, two groups recently reported trapping of single atoms with intracavity fields at the single-photon level, beginning with the work of Ref. and culminating in that of Refs.. That such trapping might be possible in these experiments is indicated by the fact that the ratio $`R`$ of initial atomic kinetic energy $`E_k`$ to the coherent coupling energy $`\mathrm{}g_0`$, $`RE_k/\mathrm{}g_0,`$ is less than unity. For the work in Refs., $`R0.06`$, while for that in Ref. $`R0.27`$. Although these ratios are indicative of the possibility of trapping with single quanta in cavity QED, the actual forces and confining potentials are somewhat more complex to analyze, as we shall see. Moreover, beyond providing single-quantum forces sufficient for atomic localization, strong coupling also means that the presence of one atom can significantly modify the intracavity field, thereby providing a means to track atomic motion by way of the light emerging from the cavity.
To understand the basic scheme for trapping of single atoms with single quanta in cavity QED, consider the energies $`\mathrm{}\beta _\pm `$ for the first excited states $`|\pm `$ of the atom-cavity system. Along the radial direction $`\rho =\sqrt{y^2+z^2}`$ and for optimal $`x`$ (standing-wave) position, $`\beta _\pm (\rho )`$ has the spatial dependence indicated in Fig. 1, which neglects dissipation. The ground state of the atom-cavity system is $`|a,0`$; the atom is in its ground state $`a`$, and there are no photons in the cavity. For weak coupling (atom far from the cavity mode center), the first two excited states are that of one photon in the cavity and the atom in the ground state, $`|a,1`$, and of the atom in the excited state $`e`$ with no photons in the cavity, $`|e,0`$. These two states are separated by an energy $`\mathrm{}\mathrm{\Delta }_{ac}`$, where $`\mathrm{\Delta }_{ac}\omega _{cavity}\omega _{atom}`$ is the detuning between the “bare” (uncoupled) atom and cavity resonances.
As an atom enters the cavity along $`\rho `$ it encounters the spatially varying mode of the cavity field, and hence a spatially varying interaction energy $`\mathrm{}g(\stackrel{}{r})`$, given by $`g(\stackrel{}{r})=g_0\mathrm{cos}(kx/)\mathrm{exp}((y^2+z^2)/w_0^2)`$ ($`k=2\pi /\lambda `$). The bare states map via this coupling to the dressed states $`|\pm `$ shown in the figure, with energies
$$\beta _\pm =\frac{\omega _{atom}+\omega _{cavity}}{2}\pm [g(\stackrel{}{r})^2+\frac{\mathrm{\Delta }_{ac}^2}{4}]^{1/2}.$$
(1)
Our interest is in the state $`|`$; the spatial dependence of the energy $`\mathrm{}\beta _{}(\stackrel{}{r})`$ represents a pseudopotential well that can be selectively populated by our choice of driving field $`_{probe}(t)`$ and $`\mathrm{\Delta }_{probe}`$ to trap the atom, as first suggested by Parkins. The system is monitored with a weak probe beam as an atom enters the cavity mode; detection of an atom transit signal triggers an increase in driving strength to populate the state $`|`$ and trap the atom. Because the experiments in the optical domain have atomic and cavity decay times ($`\kappa ^1,\gamma ^1`$) that are small compared to the time $`\tau `$ for motion through the cavity field, the atom-cavity system must be continually re-excited by way of $`_{probe}`$, thereby providing an effective pseudopotential on time scales $`\delta t`$ such that ($`\kappa ^1,\gamma ^1`$)$`\delta t\tau `$.
Although a full theory based on the preceding discussion is sufficient to provide detailed agreement with the experimental observations of Refs. (as we shall show in subsequent sections), it is reasonable to ask to what extent such a theory based on the interactions in cavity QED is necessary. In particular, it might well be that the well-established theory of laser cooling and trapping in free space could provide an adequate description of the potentials and heating rates, with the cavity merely providing a convenient means for attaining a strong drive field. With respect to the experimental results of Pinkse et al. (Ref. ), we find that this is in fact largely the case; there are only small quantitative distinctions between the free-space theory and the appropriate quantum theory. One interesting feature to note in this experiment is enhanced cooling of the atomic motion relative to the parameters of Hood et al.. This effect, which enables trapping in this parameter regime, arises through cavity-mediated cooling . For these parameters, the average localization time from simulations is extended by 75% relative to the equivalent free-atom signal; both these times are shorter than the time for an atom to transit freely through the cavity.
By contrast, in the regime of the experiment of Hood et al. (Ref.), the cavity QED interactions result in a strong suppression of dipole heating along the cavity axis relative to the free-space theory, which has a strong effect on both the duration and character of the observed atom transits. In the cavity QED setting it becomes possible to create a potential deep enough to trap an atom without simultaneously introducing heating rates that cause rapid escape from that potential. For these parameters, the average experimentally observed localization time is a factor of 3.5 longer than the equivalent free-atom average. The results of extensive numerical simulations of trapping times and radial oscillation frequencies, and their validation by way of comparisons to experimentally measured distributions, demonstrate the essential role of the single-photon trapping mechanism in the experiment of Ref.. At root is the distinction between the nonlinear response of an atom in free space and one strongly coupled to an optical cavity. For these experimental parameters, the eigenvalue structure of Fig. 1 leads to profound differences between the standard theory of laser cooling and trapping, and the extension of this theory to the regime of strong coupling in cavity QED.
Note that prior experiments in our group have confirmed that the full quantum treatment of the one-atom master equation in cavity QED is required for a description of the dynamics associated with the internal degrees of freedom for a single atom in an optical cavity in the regime $`g>(\gamma ,\kappa )`$. These experimental confirmations come by way of measurements of the nonlinear susceptibility for the coupled system in settings close to that for the experiment of Ref.. A principal goal of this paper is to investigate the extent to which a theory of atomic motion within the setting of cavity QED is likewise a necessary component in describing the center-of-mass dynamics for the experiments of Refs..
A second goal is to examine the related question of the extent to which inferences about atomic motion within the cavity can be drawn from real-time observations of the cavity field, either via photon counting or heterodyne detection of the cavity output. The interactions in cavity QED bring an in principle enhancement in the ability to sense atomic motion beyond that which is otherwise possible in free space. Stated more quantitatively, the ability to sense atomic motion within an optical cavity by way of the transmitted field can be characterized by the optical information $`I=\alpha g_0^2\mathrm{\Delta }t/\kappa \alpha \mathrm{\Delta }t`$, which, roughly speaking, is the maximum possible number of photons that can be collected as signal in time $`\mathrm{\Delta }t`$ with efficiency $`\alpha `$ as an atom transits between a region of optimal coupling $`g_0`$ and one with $`g(\stackrel{}{r})g_0`$. A key enabling aspect of the experiments in Refs. is that $`=g_0^2/\kappa (\kappa ,\gamma )`$, leading to information about atomic motion at a rate that far exceeds that from either cavity or spontaneous decay (as in fluorescence imaging). In practice, for detection strategies employed experimentally, information is extracted at a somewhat lower rate. For example, in the experiment of Hood et al. , the photon count rate would be $`(2.7\times 10^7/`$s$`)`$ (including the overall escape and detection efficiency $`\alpha 0.15`$), while for the experiment of Pinkse et al. it is $`(2.2\times 10^6/`$s$`)`$ (including an estimated overall escape and detection efficiency $`\alpha 0.11`$). For time scales $`\mathrm{\Delta }t10\mu `$ s as relevant to the following discussion, atomic motion through the spatially varying cavity mode leads to variations in the transmitted field that can be recorded with a high signal-to-noise ratio, namely, a signal of $`2.7\times 10^2`$ photons for the experiment of Hood et al. and $`2.2\times 10^1`$ for that of Pinkse et al., where each is calculated for an intracavity field strength of one photon.
The value of the optical information itself does not tell the complete story. For cavity QED experiments like those considered here, one records either the sequence of photoelectric counts or the heterodyne current versus time, from which necessarily only limited inferences about atomic motion can be drawn. However, if center-of-mass dynamics (i.e., axial and radial motions) occur on well-separated time scales, then it is reasonable to suggest that appropriate signal processing techniques could extract information about these motions from the single time sequence of the photocurrent $`i(t)`$. Such processing could presumably occur in real time if $`\alpha `$ is much faster than the rates for radial and axial motion \[e.g., the oscillation frequencies $`(f_r,f_a)`$ in a potential well, with $`f_rf_a`$\]. Unfortunately, in neither experiment is $`\alpha `$ large enough to resolve the axial dynamics directly, so the task of disentangling the radial and axial motion signals becomes more difficult, and theoretical simulations of the experiment become useful in understanding the nature of the observed transmission signals.
This difficulty arises in the experimental regime of Pinkse et al.. For these parameters, axial heating leads to frequent bursts of large-amplitude motion along the cavity axis, with envelopes extending over time scales comparable to those for radial motion. Consequently, at experimental bandwidths (averaging times), both types of motion give rise to qualitatively similar modulations in the measured transmission signal. Furthermore, motion in the radial direction has a strong diffusive component, giving rise to a wide spread of time scales for radial motion. Our simulations discussed in Sec. V suggest that for these parameters, short-time-scale modulations ($`300`$ $`\mu `$s) tend to be mostly due to bandwidth averaging over axial motion, while longer ($`500`$ $`\mu `$ s) variations such as presented in Fig. 2 of Ref. typically reflect radial motion, though these long-time-scale variations are generally modified in amplitude by the presence of axial motion. Modulations on intermediate time scales appear ambiguous in their dynamical origin.
By contrast, as shown in Ref., for the parameters of Hood et al. atoms are well localized along the standing-wave direction throughout most of the trapping interval, with axial motion giving rise to negligible signal until finally rapid axial heating leads to atomic escape. Consequently, observed variations in the photocurrent $`i(t)`$ are simpler than those of Ref. , and directly yield the radial atomic position. Furthermore, in this experiment the radial oscillation frequency is large compared to the spontaneous emission heating rate, meaning that the resulting atomic motion is largely conservative (rather than diffusive) in nature, taking place in a known potential (as demonstrated both experimentally and by way of numerical simulation). Hence, from $`i(t)`$ it becomes possible to make detailed inferences about the radial motion, even to the point of real-time observations of the anharmonic motion of a single atom and of the reconstruction of actual atomic trajectories.
The structure of the paper is as follows. Following this introduction, in Sec. II we present a detailed description of our theoretical model and its use for the implementation of numerical simulations. Section III compares effective potentials and momentum diffusion rates derived for the two experiments, along with their analogs for the hypothetical case of an equal-intensity free-space trap. These calculations explore the distinction between quantum and classical, and also give insight into the nature of atomic motion expected in both experiments. Sample simulated trajectories are presented for both cases. In Sec. IV we present experimental and simulation results for the case of Hood et al., which serve both to verify the simulations and also to demonstrate important features of the resulting motion. Sec. V gives the application of the same tools to analyze the experiment of Pinkse et al.; we see that standing-wave motion and diffusive radial motion complicate the correlation between atomic position and detected field in this case. Finally, axial motion is explored in more depth, and Fourier analysis of our simulations show that oscillations of comparable amplitude and frequency should be visible for both atoms confined (but heated) within a well, and atoms skipping along the standing wave.
### Principal findings
The theoretical treatment and numerical simulation of the motion of a single atom strongly coupled to an optical cavity, as described in Sec. II, lead to a surprisingly rich range of often qualitatively different dynamics. The motion may be essentially conservative and tightly confined around antinodes of the standing wave, or essentially dissipative and diffusive and involve interesting flights between different potential wells of the standing wave. Indeed we find that the existing experimental results of Hood et al. and Pinkse et al. exemplify these very different dynamical regimes. Key features of the atomic motion in both experimental regimes are addressed as follows:
Figures 3–5 and their associated discussion in Sec. III elucidate the nature of the trapping potential and momentum diffusion in an optical cavity as opposed to a free space standing wave. In particular we find that, even when the atom-cavity system is strongly coupled and driven such that it has a mean intracavity photon number of roughly $`1`$, the trapping potential and momentum diffusion may be only slightly different from those in a free-space standing wave, and in fact this is the case for the parameters of Pinkse et al. On the other hand, for the parameters of Hood et al. the usual fluctuations of the dipole force along the standing wave are suppressed by an order of magnitude, which to our knowledge represents qualitatively new physics for optical forces at the single-photon level within the context of cavity QED. We show that in the parameter regime of Pinkse et al. the heating rates are such that the atom could be expected to gain energy equal to a significant fraction of the total trapping potential during a single motional oscillation period for both axial and radial motion. By this measure the heating rates in the experiment of Hood et al. are much slower, indicating more nearly conservative motion, and this could be expected to have a profound effect on the qualitative nature of the dynamics in the two experiments.
Figures 6 and 7 and the corresponding text in Sec. III present simulated transits for both experiments, and discuss the qualitative features of atomic dynamics in both cases. For the parameter regime of Hood al., conservative radial motion dominates diffusion and standing-wave motion, with atomic trajectories localized at peaks of a single standing-wave antinode. Atoms trapped with the mean trapping time execute several radial orbits. The eventual escape is typically due to heating along the cavity axis. By contrast, for the experiment of Pinkse et al., a trajectory of typical duration, as in Fig. 7(a), does not experience a complete radial orbit and in fact resembles a scattering event, with a large contribution from radial diffusion as well. For these events the observed localization time is comparable to the time for free flight through the cavity. Axially the simulations show that in longer duration transits the atom frequently skips between wells of the standing-wave potential due to repeated heating and recooling.
Section IV, with Figs. 8–10, presents a more detailed and quantitative investigation of trapping and motional dynamics for the experiment of Hood et al. The ability of our simulations to closely reproduce the mean trapping times observed in the experiment provides evidence of their accuracy and utility. As illustrated in Fig. 9, the triggering strategy leads to significant modifications of the distribution of residence times within the cavity. The essentially conservative nature of the dynamics and the strong axial confinement make it possible to confidently ascribe oscillations in the transmitted intensity to radial motion of the atom. As shown in Fig. 10, the experimentally observed oscillations are consistent with the calculated potential. The conservative nature of the motion is further confirmed by the separation of orbital periods by angular momentum that is also apparent in this figure.
Section V, with Figs. 11–15, presents a detailed analysis of trapping and motional dynamics for the experiment of Pinkse et al. Again, our simulations are sufficient to reproduce the reported mean localization time. In this case, the triggering strategy leads to relatively minor modifications of the distribution of residence times for an atom within the cavity. In this case the dissipative nature of the evolution is significant; essentially no long-term localization is observed if the sign of the friction coefficient is reversed, whereas this has little effect in the parameter regime of Hood et al. These largely dissipative and diffusive motional dynamics are found to have significant effect on the information about the motion that is available in the transmitted field. For those events with a long localization time, the axial motion of the atom is repeatedly heated and cooled, resulting in slow variations in envelope of the amplitude of the rapid oscillations of the transmitted light. The time scale of these variations is comparable to that for radial motion of the atom. There are thus no unambiguous signatures for radial motion and only longer time scale excursions of the atom in the radial potential lead to variations of the output field that may be confidently ascribed to the radial motion. Likewise, although information about axial motion is also available in the output light, we find that it is in general difficult to distinguish large oscillations in a single well of the axial potential from free flight over several wells as attempted in Ref. .
## II Theoretical Model and Numerical Simulations
In this section we outline the derivation from the full quantum-mechanical master equation of the “semiclassical model” for the atomic motion used in Ref. . It turns out that this model is able to reproduce the experimental observations very accurately. Note that here the term “semiclassical” refers to approximations with respect to the atomic center-of-mass motion, and not to the internal degrees of freedom, for which the full quantum character is retained. This situation should not be confused with the semiclassical theory of cavity QED for which expectation values of field operators $`\widehat{O}_{field}`$ and atomic operators $`\widehat{O}_{field}`$ are assumed to factorize, $`\widehat{O}_{field}\widehat{O}_{atom}=\widehat{O}_{field}\widehat{O}_{atom}`$; no such approximation is made here. To distinguish these two cases, we introduce the term quasiclassical for the case of atomic motion.
The validity of the quasiclassical model depends on a separation of time scales between the atomic motion and the cavity and internal atomic dynamics. We adapt the work of Dalibard and Cohen-Tannoudji to the situation of a quantized cavity mode. A similar derivation in the bad-cavity limit appears in . The details of the derivation are essentially unchanged from free space, since the terms of the master equation which refer to the dynamics of the cavity have no explicit dependence on the operators describing the atomic motion. However, we do find conditions for the validity of the approximation for this system which depend on the properties of the cavity. Finally, we describe in more detail the numerical simulations of the resulting model first presented in Ref. . These simulations are of the kind discussed in Refs..
An analytical calculation of force, momentum diffusion and friction coefficients for the quasiclassical model of atomic motion in the low driving limit was derived by Horak and co-workers , who found a regime in which the steady-state temperature scaled as the cavity decay rate. This allows cooling of the atom below the Doppler limit, so long as the cavity can be made to have lower loss than the atom. However, the parameters of Refs. are very far from this low driving limit. Hence we employ numerical techniques based on solving the appropriate master equations by expansions in terms of Fock states of the cavity field . Note that a very early contribution developed a different theoretical framework and numerical scheme for calculating the force and friction (but not the momentum diffusion) of an atom in a cavity (or “colored vacuum”) . Very recently, Vuletic and Chu found cavity-mediated cooling in a slightly different regime to that considered by Horak et al.
### A Model of atom-light interaction in a cavity
The Hamiltonian for a two-level atom interacting with a single mode of the electromagnetic field in an optical cavity using the electric dipole and rotating-wave approximations (in the interaction picture with respect to the laser frequency) is
$`H`$ $`=`$ $`{\displaystyle \frac{\stackrel{}{p}^2}{2m}}+\mathrm{}(\omega _{atom}\omega _{probe})\sigma ^{}\sigma +\mathrm{}(\omega _{cavity}\omega _{probe})a^{}a`$ (3)
$`+\mathrm{}g(\stackrel{}{r})(a^{}\sigma +\sigma ^{}a)+\mathrm{}\left(a^{}+^{}a\right).`$
This is the familiar Jaynes-Cummings Hamiltonian modified to take into account the external degrees of freedom of the atom and the spatial variation of the cavity mode. The first term is the kinetic energy of the atom, and the next two terms are the energy in the internal state of the atom and the cavity excitation. The fourth term describes the position-dependent interaction of the cavity mode and the atomic dipole. It is important to note that $`\stackrel{}{r}`$ and $`\stackrel{}{p}`$ are operators. Thus, for example, the exact strength of the coupling between the atomic internal state and the cavity field depends on the shape of the atomic wave packet, which is in turn determined by the mechanical effects of the cavity field. Some implications of this Hamiltonian are considered in detail by Vernooy and Kimble . The Hamiltonian has been written in terms of cavity and dipole operators that rotate at the frequency of the probe field $`\omega _{probe}`$. The real atomic transition (cesium in Ref. and rubidium in Ref. ) in fact involves several degenerate magnetic sublevels, but we assume that the cavity is driven by circularly polarized light and that the atom is optically pumped such that it occupies an effective two-level system described by the dipole operator $`\sigma `$ with the quantization axis along $`x.`$
Dissipation in the system is due to cavity losses and spontaneous emission. By treating modes external to the cavity as a heat reservoir at zero temperature in the Born, Markov, and rotating-wave approximations, it is possible to derive the standard master equation for the density operator $`\rho `$ of the system as
$`{\displaystyle \frac{d\rho }{dt}}`$ $`=`$ $`{\displaystyle \frac{i}{\mathrm{}}}[H,\rho ]+\kappa (2a\rho a^{}a^{}a\rho \rho a^{}a)`$ (6)
$`+{\displaystyle \frac{3\gamma }{4\pi }}{\displaystyle d^2\widehat{𝐤}S(\widehat{𝐤}\widehat{𝐱})\mathrm{exp}(ik\widehat{𝐤}𝐫)\sigma \rho \sigma ^{}\mathrm{exp}(ik\widehat{𝐤}𝐫)}`$
$`\gamma (\sigma ^{}\sigma \rho +\rho \sigma ^{}\sigma ).`$
The third and fourth terms describe the effect of spontaneous emission on the atomic motion including the momentum kick experienced by the atom as a result of the spontaneous emission. The unit vector $`\widehat{𝐤}`$ is the direction of an emitted photon. The pattern of dipole radiation is accounted for by the angular factor $`S(\widehat{𝐤}\widehat{𝐱})=\left[1+(\widehat{𝐤}\widehat{𝐱})^2\right]/2`$ .
### B Quasiclassical motion of the center of mass
It is possible to eliminate the internal and cavity dynamics adiabatically in favor of the slower dynamics of the motional state in parameter regimes of direct relevance to current experiments. Intuitively, for the quasiclassical approximation to work, the state of the atom needs to be sufficiently localized in position and momentum on the scales important to the problem so that it can be thought of as a classical particle. The conditions for adiabatically eliminating the internal and cavity dynamics roughly correspond to this idea. It turns out that it is necessary first that exchanges of momentum with either the cavity field or by spontaneous emission into free space should result in momentum kicks that are small compared with the momentum spread $`\mathrm{\Delta }p`$ of the atomic Wigner function, thus
$$\epsilon _1\mathrm{}/\mathrm{\Delta }p1.$$
(7)
For an atom which is in a minimum uncertainty state with respect to the position-momentum Heisenberg inequality this requires that the state is localized to better than a wavelength. The atomic motional state will in general be a mixture allowing the position spread to be broader. However, this requirement means that the motional state can be thought of as a probabilistic mixture of pure states localized to within a wavelength, and so places a limit on the coherence length of the motional state . Second it is important that the range of Doppler shifts of the atom due to its momentum spread is small compared to the atomic and cavity linewidths, thus
$$\epsilon _2k\mathrm{\Delta }p/m\gamma k\mathrm{\Delta }p/m\kappa 1.$$
(8)
In this paper it will be assumed that the root-mean-square atomic momentum obeys this inequality, thus making a low velocity approximation, but the arguments here can in fact be generalized to arbitrary mean velocities of the atom . The Heisenberg inequality means that this also requires a minimum position spread of the atom
$$\mathrm{\Delta }r\mathrm{}k/m\gamma ,\mathrm{}k/m\kappa .$$
(9)
These criteria are a simple generalization of the situation for laser cooling in free space which can be imagined as the situation $`\kappa \mathrm{}`$. The consistency of these conditions, which effectively put lower and upper limits on the atomic momentum spread, requires that
$$\frac{^{\mathrm{}^2k^2/2m}}{\mathrm{}\gamma }1,\frac{^{\mathrm{}^2k^2/2m}}{\mathrm{}\kappa }1.$$
(10)
The first of these conditions is well known for laser cooling in free space—the requirement that the recoil energy of the atomic transition be much lower than the Doppler energy, which effectively controls the limiting temperature of the laser cooling. This condition is well satisfied for heavy atoms such as cesium and rubidium and the optical transitions employed in cavity QED experiments considered here. The analogous condition brought about by the cavity dynamics requires that the recoil energy associated with exchanging excitation with the cavity field is much smaller than the energy width of the cavity resonance. Just as the first criterion implies that the atom still be in resonance with a driving field at its transition frequency after spontaneously emitting, the second criterion implies that absorbing or emitting a photon from the cavity will leave the atom near the cavity resonance. In the experiments of Refs. , $`\kappa \gamma `$, so that this second criterion does not place a stronger restriction on the validity of the approximations than the free-space limit. However, it is important to note that the design of the cavity, as well as the atom and transition that are chosen, now has an effect on the validity of the approximation. It would be possible, for example, to change the cavity length in such a way that the system moves from a regime in which the quasiclassical treatment is appropriate into one in which it is not. In practice for cold atoms cooled to roughly the Doppler limit ($`\mathrm{\Delta }p^2/2m\mathrm{}\gamma ,\mathrm{}\kappa `$) it will be the case that $`\epsilon _1\epsilon _2\sqrt{\left(\mathrm{}^2k^2/2m\right)/\mathrm{}\gamma },\sqrt{\left(\mathrm{}^2k^2/2m\right)/\mathrm{}\kappa }`$, and so a consistent expansion should be to equal order in these small parameters.
The derivation of Ref. may be applied to our problem, and proceeds by transforming the master equation \[Eq. 6\] into an evolution equation for a Wigner operator,
$$W(\stackrel{}{r},\stackrel{}{p},t)=\frac{1}{h^3}d^3\stackrel{}{u}\stackrel{}{r}+\frac{1}{2}\stackrel{}{u}|\rho |\stackrel{}{r}\frac{1}{2}\stackrel{}{u}\mathrm{exp}(i\stackrel{}{p}\stackrel{}{u}/\mathrm{}),$$
(11)
describing the complete state of the system. An approximate Fokker-Planck equation for the Wigner function describing the motional degrees of freedom alone is found by writing this equation as a Taylor expansion in terms of the small parameters $`\epsilon _1`$ and $`\epsilon _2`$, and truncating that expansion at third order. The force operator is defined as the gradient of the atom-cavity coupling
$$\stackrel{}{F}(\stackrel{}{r})=\mathrm{}g_0\stackrel{}{}\psi (\stackrel{}{r})(a^{}\sigma +\sigma ^{}a).$$
(12)
It is possible to show that the Fokker-Planck equation for the atomic Wigner function $`f`$ takes the form
$$\frac{}{t}f+\frac{\stackrel{}{p}}{m}\frac{}{\stackrel{}{r}}f=\stackrel{}{\varphi }(\stackrel{}{r})\frac{}{\stackrel{}{p}}f+\underset{ij}{}D_{ij}\frac{^2}{p_ip_j}f+\mathrm{}^2k^2\gamma \sigma ^{}\sigma _{\rho _s}\underset{ij}{}E_{ij}\frac{^2}{p_ip_j}f+\underset{ij}{}\eta _{ij}\frac{^2}{p_ir_j}f+\underset{ij}{}\mathrm{\Gamma }_{ij}\frac{}{p_i}\left(p_jf\right).$$
(13)
The quantities appearing in the Fokker-Planck equation can be calculated from the master equation for the internal and cavity degrees of freedom alone that is obtained by setting $`\stackrel{}{r}`$ to some real number value $`\stackrel{}{r}_0`$, and disregarding the kinetic-energy term. We define $`\rho _s(\stackrel{}{r})`$ as the steady state of this master equation, with the steady-state expectation value of the arbitrary operator $`c`$ given by $`c_{\rho _s}=`$Tr$`(c\rho _s(\stackrel{}{r}))`$. The parameters appearing in the Fokker-Planck equation can then be expressed as follows:
$`\stackrel{}{\varphi }(\stackrel{}{r})`$ $`=`$ $`\text{Tr}\left[\stackrel{}{F}(\stackrel{}{r})\rho _s(\stackrel{}{r})\right],`$
$`D_{ij}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑\tau \left[{\displaystyle \frac{1}{2}}F_i\left(\tau \right)F_j\left(0\right)+F_j\left(0\right)F_j\left(\tau \right)_{\rho _s}\varphi _i\varphi _j\right],`$
$`E_{ij}`$ $`=`$ $`{\displaystyle \frac{3}{8\pi }}{\displaystyle d^2\widehat{k}S(\widehat{k}\widehat{x})\widehat{k}_i\widehat{k}_j},`$
$`\eta _{ij}`$ $`=`$ $`{\displaystyle \frac{1}{m}}{\displaystyle _0^{\mathrm{}}}𝑑\tau \tau \left[{\displaystyle \frac{1}{2}}F_i\left(\tau \right)F_j\left(0\right)+F_j\left(0\right)F_j\left(\tau \right)_{\rho _s}\varphi _i\varphi _j\right],`$
$`\mathrm{\Gamma }_{ij}`$ $`=`$ $`{\displaystyle \frac{i}{m\mathrm{}}}{\displaystyle _0^{\mathrm{}}}𝑑\tau \tau [F_i\left(\tau \right),F_j\left(0\right)]_{\rho _s}.`$
Simple integrations give $`E_{xx}=2/5`$ and $`E_{yy}=3/10=E_{zz}`$ and all other components of $`E`$ are zero. Excepting the different definition of the force operator $`\stackrel{}{F}`$, these are the expressions that can be derived in case of a free-space light field . However, it is important to bear in mind the extra conditions on the validity of the adiabatic elimination. The master equation \[Eq. 6\] means that the force expectation values and correlation functions can be very different from those that are calculated in free space. In practice, the contribution from the parametric tensor $`\eta `$ is often smaller than that from the diffusion tensor $`D`$ by a factor of order $`\epsilon `$, and is usually disregarded in treatments of free-space laser cooling.
Thus, as assumed in earlier work, calculating the quasiclassical motion of the atom in a cavity field only requires that the force and its correlation function be evaluated for the full atom-cavity master equation. Such prior treatments assumed that the atom is motionless; however, they can be extended to atoms moving at some velocity under the same conditions . The diffusion coefficients may be found by first calculating the correlation functions via the quantum regression theorem and numerical integration, or directly via matrix-continued fraction techniques . A matrix-continued fraction calculation requires that the field mode be periodic, and as such it only works along the standing-wave axis of the cavity mode. In directions perpendicular to this, the calculation of correlations from the master equation is essentially the only option if the atom is not slowly moving.
### C Stochastic simulations of the quasiclassical model
It is possible to recast the Fokker-Planck equation of Eq. 13 into a simple set of stochastic equations which describe atomic trajectories in the cavity field. These equations can be used to gain intuition about the atomic motion and how it is affected by mechanical forces. The diffusion and friction tensors can be rewritten using the definition of the force operator \[Eq. (12)\]
$`D`$ $`=`$ $`\mathrm{}^2g_0^2\left[\stackrel{}{}\psi (\stackrel{}{r})\right]\left[\stackrel{}{}\psi (\stackrel{}{r})\right]^T{\displaystyle _0^{\mathrm{}}}𝑑\tau \left[{\displaystyle \frac{1}{2}}\mathrm{\Phi }\left[\tau \right]\mathrm{\Phi }\left[0\right]+\mathrm{\Phi }\left[0\right]\mathrm{\Phi }\left[\tau \right]_{\rho _s}\mathrm{\Phi }_{\rho _s}^2\right]`$ (15)
$`=`$ $`\mathrm{}^2g_0^2\xi (\stackrel{}{r})\left[\stackrel{}{}\psi (\stackrel{}{r})\right]\left[\stackrel{}{}\psi (\stackrel{}{r})\right]^T,`$ (16)
$`\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle \frac{i}{m}}\mathrm{}g_0^2\left[\stackrel{}{}\psi (\stackrel{}{r})\right]\left[\stackrel{}{}\psi (\stackrel{}{r})\right]^T{\displaystyle _0^{\mathrm{}}}𝑑\tau \tau [\mathrm{\Phi }\left[\tau \right],\mathrm{\Phi }\left[0\right]]_{\rho _s}`$ (17)
$`=`$ $`{\displaystyle \frac{\mathrm{}g_0^2}{m}}\chi (\stackrel{}{r})\left[\stackrel{}{}\psi (\stackrel{}{r})\right]\left[\stackrel{}{}\psi (\stackrel{}{r})\right]^T,`$ (18)
where $`\mathrm{\Phi }=a^{}\sigma +\sigma ^{}a`$. Writing the parameters of the quasiclassical model in this form relies on the approximation that the atom is slowly moving, namely, that it does not move a significant fraction of a wavelength during a cavity or atomic lifetime. Note that the functions $`\xi `$ and $`\chi `$ depend on position only through the coupling $`g=g_0\psi `$. They can be calculated efficiently by finding $`D_{xx}`$ and $`\mathrm{\Gamma }_{xx}`$ using matrix-continued fractions, and then dividing off the gradient factors. A matrix-continued fraction technique cannot be used to find the other components of the momentum diffusion or the friction tensors directly, since the field mode is not periodic across the Gaussian profile of the mode.
It is now straightforward to convert the Fokker-Planck equation for the Wigner function into an equivalent set of Itô stochastic differential equations. The resulting (Itô) equations are
$`d\stackrel{}{x}`$ $`=`$ $`{\displaystyle \frac{1}{m}}\stackrel{}{p}dt,`$ (20)
$`d\stackrel{}{p}`$ $`=`$ $`\mathrm{}g_0\mathrm{\Phi }\stackrel{}{}\psi dt{\displaystyle \frac{\mathrm{}g_0^2}{m}}\chi (\stackrel{}{r})\left(\stackrel{}{p}\stackrel{}{}\psi \right)\stackrel{}{}\psi `$ (22)
$`+2\mathrm{}g_0\sqrt{\xi (\stackrel{}{r})}\stackrel{}{}\psi dW_1+2\mathrm{}k\gamma \sqrt{\sigma ^{}\sigma }\sqrt{E}d\stackrel{}{W.}`$
The Wiener increment $`dW_1`$ has the usual properties, in particular $`dW_1^2=dt`$. The vector $`d\stackrel{}{W}`$ is a vector of three such increments. The terms in the equation for the momentum are the mean radiative force, its first-order dependence on momentum, and its fluctuations due to the atom-cavity system and due to the coupling to free space, respectively. These equations depend on the quantities $`\mathrm{\Phi },\chi ,\xi `$ and $`\sigma ^{}\sigma `$, which are functions of position through $`g`$ only. A straightforward simulation of these equations only needs to store ordered look-up tables of these quantities for given values of $`g`$, rather than for all possible values of $`\stackrel{}{r}`$. All of the other quantities that appear, including $`g`$, are simple functions of $`\stackrel{}{r}`$ and $`\stackrel{}{p}`$. At each time step the algorithm searches the look-up table for the current value of $`g`$, starting from the previous value, and reads off the current values of $`\mathrm{\Phi },\chi ,\xi ,\sigma ^{}\sigma `$. A linear interpolation for the two closest values of $`g`$ was used, but more sophisticated interpolation schemes could be implemented. Since $`g`$ will not change by a large amount in any one timestep the search can be very efficient; a routine from Ref. was used for this. In the low-velocity limit of the quasiclassical theory, these stochastic differential equations describe all the motional dynamics of the atom inside the cavity. The term proportional to $`\eta _{ij}`$ leads to correlations between the atomic position and momentum. The effect of $`\eta _{ij}`$ is typically small compared to friction and diffusion and has been ignored for the moment as is common practice in free-space standing waves. Terms in the SDE corresponding to the $`\eta `$ term in the Fokker-Planck equation could easily be added. This would mean adding a new noise source which would affect the evolution of the position as well as the momentum.
## III Application of the Model to Experimental Regimes
### A Potentials and heating rates for atomic motion
The “quasiclassical” model discussed in the previous Sec. II can give us a great deal of information about the nature of the dynamics that may be expected in the parameter regimes relevant to the experiments of Hood et al. and Pinkse et al. . In particular we are interested in whether quantization of the cavity field leads to any significant change in the dynamics, in the sense of asking whether the atomic motion is very different in the cavity from what it would be in a free-space standing wave of the same intensity and geometry as the cavity mode. Second, we can investigate the nature of the resulting atomic motion in the cavity field, which can be either predominantly conservative or significantly diffusive and dissipative, depending on the particular parameters of interest.
To obtain a feel for the type of atom dynamics expected, effective potentials and heating rates were calculated for both axial and radial directions of motion. The effective potential of the atom in the cavity field may be calculated from the force by
$`U(\stackrel{}{r})={\displaystyle _0^\stackrel{}{r}}\stackrel{}{F}(\stackrel{}{r}^{})𝑑\stackrel{}{r}^{}.`$
The heating rates represent the average increase in the motional energy due to the momentum diffusion at a given position $`\stackrel{}{r}`$, and may be calculated from the diffusion tensor according to
$`{\displaystyle \frac{dE}{dt}}(\stackrel{}{r})=\text{Tr}[D(\stackrel{}{r})]/m.`$
Thus the axial potential at the center of the mode is $`U(0,x)=_0^x\stackrel{}{F}(0,x^{})𝑑x^{}`$ and the associated axial heating rate is $`dE(0,x)/dt=D_{xx}(0,x)/m.`$ These quantities along with their radial equivalents $`U(\rho ,0)`$ and $`dE(\rho ,0)/dt`$ are plotted in Fig. 3 for the parameters of Hood et al.. The force and momentum diffusion coefficient for the cavity system were calculated according to the formulas described above by numerical techniques based on Ref. . The field state is expanded in terms of number states, and truncated at an appropriate level and a matrix-continued fraction algorithm is used to calculate $`D`$. The axial potentials and heating rates have $`\lambda /2=426`$ nm periodicity inherited from the standing-wave field strength. Observe that the axial heating rates have minima at both field antinodes and field nodes.
The first thing to note is that the axial and radial heating rates are very different. In the radial direction, heating is dominated by diffusion due to spontaneous-emission recoils. Axially, however, the reactive or dipole fluctuation component of the diffusion dominates. This is because the reactive component is proportional to the gradient of the field squared, which is much larger for the axial direction where variations are greater (by a factor of $`2\pi w_0/\lambda `$). This contribution also has the property that it does not saturate with the atomic response.
It is already clear that it should be possible to trap individual atoms, since the potential depth of roughly 2.5 mK is greater than the initial energy of the atoms in the experiment (around 0.46 mK) and the heating rate in the radial potential is relatively slow. Over 50 $`\mu `$s (a time scale over which the atomic motion is strongly affected by the potential) the total heating will typically still be small compared to the depth of the potential. However, the importance of the quantum character of the relevant fields or phenomena is not ensured by the statement that trapping occurs with a mean field strength of $`\overline{m}1`$ photon, since this is trivially the case in an equivalent free-space volume for a field of the same intensity as that inside the cavity.
In order to see whether a full quantum description of the atom-cavity is necessary in order explain observed effects, Fig. 3 also shows the values calculated for an atom in an equivalent free-space standing wave, calculated by standard techniques. This free-space standing wave has the same geometry as the cavity mode, and the same peak field strength $`g_0|a|^2(0,0).`$ The detuning between the free-space field and the atom is chosen to be $`\mathrm{\Delta }_{probe}`$. Perhaps surprisingly, the only large difference between the two models is in the axial heating rate, where a strong suppression of the axial heating is seen in the quantum calculation. This suppression is an effect of the quantized nature of the intracavity field. The self-consistent coupling of the cavity field and atomic position (in a semiclassical sense) cannot explain this suppression; in fact, by itself this coupling would lead to an increase in diffusion over the free-space case, since the atomic motion within the cavity induces steeper gradients in the field. The suppression of diffusion is then evidence that it is necessary to use a fully quantum description, and speak of single photons rather than classical fields for these experimental parameters. As discussed in Ref. , this suppression of the axial heating was essential for the trapping of atoms in the cavity. Thus for these experimental parameters, the eigenvalue structure of Fig. 1 leads to profound differences between the standard theory of laser cooling and trapping and the extension of this theory to the regime of strong coupling in cavity QED.
By way of comparison, the same quantities are plotted for the parameters relevant to Pinkse et al. in Fig. 4.. The smaller value of $`g_0`$ in this experiment leads to a smaller effective potential, since the spatial gradients of the dressed state energy levels (which lead to the potential) are proportional to $`g_0`$. More importantly, the diffusion values calculated from the full quantum model discussed above are now little different from those of the equivalent free-space standing wave. This lack of a clear difference in potentials or diffusion indicates that the quantized nature of the field is not required to explain the radial trapping observed in Ref. . Note that the resulting axial heating rates are essentially the same as those of Ref. in absolute magnitude; however, in Ref. the potential was made deeper without the expected corresponding increase in diffusion. For the parameters of Ref. one additional interesting feature appears—enhanced cooling of the atom motion relative to the parameters of Ref. . This arises through cavity-mediated cooling , and as we shall see, has an important effect on the axial dynamics of atoms in the experiment of .
We now wish to use these potentials and heating rates to gain an intuitive understanding of the character of atomic motion that we would expect to observe in each case. In particular, we are interested in exploring the degree to which the atomic motion in the potential can be close to conservative motion, or likewise the degree to which it could be dominated by diffusion.
The time scales of relevance to the conservative motion may be characterized by the period associated with small-amplitude oscillations in the bottom of the axial ($`\tau _a=1/f_a`$) and radial ($`\tau _r=1/f_r`$) potential wells. If the energy changes only by a small fraction (relative to the the total well depth $`U_0`$) on this time scale, motion will be nearly conservative. Fig. 5 plots the potentials and heating rates for the two cases in this new set of scaled units; heating rates are expressed as an energy increase per oscillation period, as a fraction of $`U_0`$ (note as the atom heats and explores the anharmonicity of the potential, this only lengthens the period of oscillation). Interestingly, we see a clear qualitative difference in the nature of the atomic motional dynamics. For the parameters of Hood et al., in the radial plane spontaneous emission only gives small perturbations to the energy over the time scale of single orbits, and motion is nearly conservative. We note that this low level of diffusion enabled the reconstructions of single-atom trajectories in Ref. , for which the small changes in angular momentum could be accurately tracked. A quite different regime is found for the parameters of Pinkse et al., where the radial atomic motion is strongly affected by heating from spontaneous emission kicks. Here an average atom gains an energy of nearly half the well depth in what would be a radial orbit time, adding a large diffusive component to the motion. This same scaling shows that the axial heating rate is also much more rapid on the scale of the potential in Ref. , which suggests that the atom will more quickly escape its confinement near an antinode and begin to skip along the standing wave. The qualitative understanding of the atomic motion gained here is borne out by the simulations of Refs. and , and is explored in more detail in the simulations to follow.
### B Simulated transits
Simulations of the kind described in Sec. II were performed for the parameters of the two experiments, and individual instances of these simulations give insight into the dynamics of the motion—for example, the relative significance of conservative or dissipative dynamics—and the correlation between atomic motion and the cavity field state, which is in turn measured by detection of the output field. Ensembles of these trajectories provide the statistics of the motion described by the Fokker-Planck equation \[Eq. 13\] which may then be used to provide histograms of transit times to compare to the experimental data or to test reconstruction algorithms for the motion. In order to approximate the experiment as closely as possible, some effort was made to match the detailed experimental conditions. The two general considerations were to reasonably accurately estimate the initial distribution of atomic positions and momenta for atoms and to consider detection noise and bandwidth when simulating the feedback switching of the probe laser power.
For each trajectory in the simulations, initial atomic position and momentum values were drawn from a probability distribution, which was chosen to correspond to the cloud of atoms following laser cooling and then free fall or launching by an atomic fountain to the cavity mirrors. In the simulations, all the atoms started in a horizontal plane $`1`$ $`\frac{3}{4}`$ mode waists above or below the center of the cavity mode, where mechanical effects on the atom are negligible. Since the MOT from which the atoms are falling or rising has dimensions much larger than the cavity mode, the initial position in the axial direction was chosen from a flat distribution over the cavity mode, and the initial position along the $`y`$ axis was also chosen from a flat distribution over 1 $`\frac{1}{2}`$ mode waists on either side of the mode center—this distance could be modified but atoms that are far out in the mode radially do not typically cause large increases in the cavity transmission, and therefore do not trigger the feedback. The velocity of the atom along the cavity axis is limited by the fact that it must not hit one of the mirrors while falling toward the cavity, and this was also chosen from a flat distribution where the speed was not more than 0.46 cm/s for the cavity of Hood et al.. Although the two experiments have rather different geometries, we estimate that this consideration leads to a very similar limiting velocity for motion along the axis. In the experiment of Pinkse et al., we used 0.4 cm/s. The velocity along the $`z`$ axis was chosen from a Gaussian distribution appropriate to the temperature of the MOT ($`20`$ $`\mu K`$) after polarization gradient cooling. For the velocities in the vertical direction were chosen by calculating as appropriate for an atom falling freely from the MOT (the MOT is situated 3.2mm above the mode with a spatial extent of standard deviation 0.6mm). Thus atoms arriving at the cavity axis have a mean vertical velocity $`\overline{v}=25`$cm$`/`$s. Some of these parameters such as the height, size and temperature of the initial MOT are not precisely known for the experiment, so that some consideration of the variation of the histograms and other features of the resulting simulations has been made although no systematic optimization in order to obtain the best agreement has been undertaken. In Ref. the mean initial vertical velocity of atoms entering the cavity is $`20`$ cm/s. This speed is very much less than the mean velocity imparted to the atoms by the pushing beam which launches them from the MOT $`25`$ cm below, and as a result the atoms are all near the top of their trajectories. Simple kinematical calculations show that the resulting distribution of velocities should be rather broad compared to the mean. In the absence of more detailed information about the MOT temperature and spatial size and the strength of the pushing beam we choose the initial vertical velocity distribution to be a Gaussian of mean 20 cm/s and standard deviation 10 cm/s—this leads to a distribution of trapping times with a mean that matches the mean reported in Ref. . Each trajectory proceeds until the atom is either a greater radial distance from the center of the mode than it started from, or it has moved sufficiently far in the axial direction that it would hit one of the cavity mirrors.
The detection and triggering are modeled as follows. In the parameter range in which the “quasiclassical” model is valid, the cavity field comes to equilibrium with the atomic position on a time scale much faster than the atomic motion itself, and thus the light transmitted through the cavity (over bandwidths of the order of tens to hundreds of kilohertz) is associated with the atomic motion. At each point in the simulation the intracavity field and intensity expectation values are stored in order to record for each trajectory a noiseless and infinite-bandwidth trace. In practice, experimental traces will look like filtered and noisy versions of these traces. As an atom enters the cavity mode, a weak driving field is present for probing. In order to model the triggering step, the field intensity $`a^{}a`$ or field amplitude modulus squared $`|a|^2`$ is averaged over a time equal to the bandwidth of the detection in the case of heterodyne detection as in Ref. , or over the time windows in which photocounts are binned in the case of direct photodetection as in Ref. . A random number with the appropriate variance to represent the shot noise is added and the total is compared with some predecided level—if the transmission exceeds this level the probe laser beam is increased in strength in order to attempt to trap the atom. In the case of Ref. the trigger level is $`|a|^2=`$ 0.32, the averaging time is 9 $`\mu `$s, and there is a 2-$`\mu `$s delay between triggering and changing the driving laser power. For the experimental bandwidth of 100 kHz, the appropriate noise has standard deviation 0.05 at a transmitted signal of 0.32. These parameters are chosen so as to match as closely as possible the conditions of the experiment. The same procedure is followed for simulations of the parameters . Although the exact triggering protocol is not described there, we assumed that counts over a period of $`10`$ $`\mu `$s were used to decide whether or not to trigger and the noise was chosen to be consistent with the reported photon count rate of $`2\times 10^6`$s<sup>-1</sup> .
Examples of such trajectories are plotted for the parameters of Ref. in Fig. 6 and for those of Ref. in Fig. 7. The chosen trajectories range in length from the experimentally reported mean transit time upward, and are chosen because they show typical features of the dynamics in each case. It is clear that the two experiments are in quite different parameter regimes, as already indicated by the relative sizes of the potentials and heating rates.
For the parameters of Ref. , the atoms orbit in a radial plane; some have nearly circular and some very eccentric orbits. The motion along the axial direction is usually well localized near an antinode of the standing wave, where the axial heating rate is small. This localization occurs because atoms are channeled into the antinodes by the weak potential associated with the initial probing field, which slowly begins to affect an atom as it falls across the mode waist during the detection stage of the experiment. However the strong axial heating that is present away from the antinodes means that once an atom begins to heat axially, it suffers a burst of heating (over several hundred microseconds), which leads to its loss from the potential well associated with a single antinode of the field. Frequently the atom leaves an axial potential well when it is radially far from the center of the cavity mode, since in this case the axial potential becomes weaker. Note that the mean transit time in Ref. corresponds to $`3.5`$ radial orbits around the center of the cavity mode, so transits with multiple oscillations are frequently observed. In Ref. the radial oscillation frequency is slower, so an atom of mean transit time does not in fact make a complete rotation about the mode center. The radial motion in this case is also visibly more stochastic in nature, as a result of the relatively faster spontaneous emission momentum diffusion discussed above.
Another interesting difference between the two parameter regimes is, as suggested in Ref. , the relative importance of atomic motion along the standing wave as opposed to oscillations around a single antinode. In the case of Ref. , long, strongly trapped transits almost always involve intervals when an atom is skipping along the standing wave, as well as intervals when it is oscillating in an individual well. By contrast, for the parameters of Ref. , only a few percent of trajectories involve skipping during times in which the atom is trapped, and this is usually associated with movement over one or two wells with the atom falling back into the adjacent or a nearby well. This happens so quickly that it does not affect the radial motion in practice, or lead to a detectable signal in the output light, so that these rare events of skipping do not affect the reconstructions of Ref. .
As noted in Ref. , the axial motion often becomes more significant at the end of a transit and as the atom is leaving the mode, which leads to atoms skipping a well in perhaps as many as one in five cases at the end of the transit. We find from the simulations that in Ref. , the first escape time from an axial potential well for an atom initially localized near an antinode is sufficiently short compared to the mean trapping time that skipping along the wells almost always takes place. On the other hand, the first escape time is of the order of several times the mean trapping time for the parameters of Ref. , so skipping between standing wells is correspondingly rare.
It is interesting to note that the friction coefficient for the parameters of Pinkse et al. is much more significant than for the experiment of Hood et al., and plays an important role in the axial motion of the atom. As in the trajectories shown here it is a feature of essentially every trajectory for the parameters of Ref. that the atom spends time in potential wells associated with several different antinodes of the field. However, we performed simulations with the sign of the friction coefficient reversed, and found that no more than a few percent of trajectories were recaptured in a second well after having begun to skip along the standing wave. Clearly the dissipative nature of the motion is an integral feature of the dynamics in this regime, and in particular it enables the atoms to fall back into axial potential wells after escape due to the rapid heating in that dimension.
## IV Simulation Results for the Experiment of Hood et al.
Having presented the theoretical basis underlying the simulated atom trajectories, in this section we present results of these simulations and their comparison with experimental results as reported in Ref.. We generate a set of simulated trajectories for the parameters $`(g_0,\gamma ,\kappa )=2\pi (110,2.6,14.2)`$ MHz with detuning parameters $`\mathrm{\Delta }_{ac}=\omega _{cavity}\omega _{atom}=2\pi \times 47`$ MHz and $`\mathrm{\Delta }_{probe}=\omega _{probe}\omega _{atom}=2\pi \times 125`$ MHz. In correspondence with the experimental protocol, the initial pretriggering level of the driving laser gives a 0.05-photon mean-field strength in the empty cavity; when this level rises to 0.32 photons indicating the presence of an atom, we trigger a sixfold increase in the driving strength to a trapping level of a 0.3-photon empty-cavity mean field strength. A close correspondence between theory and experiment is obtained for these results, demonstrating the relevance of this theoretical model to the physics of the actual experiment. In addition, both theoretical and experimental results exhibit features which are relevant to building up a picture of the nature of the single-atom, single-photon trapping and atomic dynamics, both qualitatively and quantitatively.
We begin by presenting the qualitative similarity of experimental and simulated atom transit signals, as observed via detection of cavity transmission as a function of time. Fig. 8 shows two sample experimental transits \[(a) and (b)\] and two sample simulated transits \[(c) and (d)\]. For the simulated transits, traces of the corresponding radial and axial motion are also shown. Transmission is shown here as $`\overline{m}=|a|^2`$, as is appropriate for the balanced heterodyne detection of Ref.. In the case of the simulated results, the simulated transmission signal has been filtered down to the experimental detection bandwidth of 100 kHz, and both technical noise and shot noise have been added. The transmission signal thus processed can be seen to lose some of the clarity with which it reflects the full atomic dynamics, in comparison to the transmission traces of Fig. 6. In particular, the experimental detection bandwidth is much slower than the time scale for axial oscillation in the confining potential, so that observed transmission signals are averaged over the fast variation in $`g`$ caused by these axial oscillations. The observed maximum transmission should therefore be lowered relative to theoretical predictions, by an amount dependent on the amplitude of typical axial motion. Thus this finite-bandwidth effect allows for an experimental estimation of the axial confinement of a typical transit. Such a procedure gives an estimate of confinement within $`70`$ $`nm`$ of an antinode, in good agreement with simulation results which suggest typical confinement within $`50`$ $`nm`$. It is important to note that while such tight confinement appears typical over the duration of a trajectory, atoms commonly undergo rapid diffusive heating near the end of their confinement lifetime, which leads to their escape in a majority of cases.
### A Trapping lifetimes
From the entire set of experimental and simulated trajectories like those of Fig. 8, it is possible to investigate some quantitative aspects of the trapping dynamics. First we focus on the trap lifetimes produced by the triggered-trapping scheme. Figures 9(a) and 9(b) show histograms of experimental transit times for untrapped atoms and for atoms trapped by means of the triggered-trapping strategy. Transit durations are determined from the experimental data by recording the time interval during which the transmission signal is clearly distinguishable from the empty-cavity transmission level, in the presence of experimental noise. Since the signal-to-noise ratio for observing transits depends on the specific probe parameters, one must be careful to compare untriggered and triggered transits observed with the same detunings and intracavity field strengths. The sole difference must be that in the untriggered case, the empty cavity field is set at a constant strength so that the atom falls through the effective potential, whereas in the triggered case the field begins at a lower level and is only turned up once the atom enters the cavity, thus confining the atom. For example, Fig. 8(a) shows sample untriggered (dashed) and triggered (solid) transit signals which correspond to one another in this way.
In Fig. 9 the difference in transit lifetimes between triggered (b) and untriggered (a) cases is immediately striking. For their initial fall velocity of $`\overline{v}=25`$ cm/s, atoms have a free-fall time of $`110`$ $`\mu `$s across the cavity waist $`2w_0=2(14.06`$ $`\mu `$m). As discussed above, the duration of observed transits is limited by the signal-to-noise ratio, which provides a slightly more restrictive cut on transit durations, so the untriggered data set shows a mean duration of $`92`$ $`\mu `$s. In contrast, when the triggered-trapping strategy is employed, the mean trapping lifetime is $`340`$ $`\mu `$s. The dispersion about the mean likewise changes drastically from $`75`$ $`\mu `$s in the untriggered case to $`240`$ $`\mu `$s in the triggered case. These results represent a clear signature of the trapping of single atoms with single photons via this method. In this setting, atoms have been observed to remain trapped in the cavity field for as long as $`1.9`$ ms.
The corresponding theoretical histograms are shown in Figs. 9(c) and 9(d) for the untriggered and triggered cases. The start of the transit is taken to be the time at which an atom could be distinguished in the cavity given the signal to noise, and the final time is taken to be the last point at which the transmission dropped to within the noise of the transmission with no atom. This definition accounts for the fact that as atoms move out in the radial direction the transmission often drops to around the free space value, but returns again to some large value over the time scales of the atomic motion. These levels were chosen to duplicate as closely as possible the protocol for deciding transit times for the experimental data.
The simulated transit set shows a mean trapping time of $`96`$ $`\mu `$s in the untriggered case and $`383`$ $`\mu `$s in the triggered case and dispersions of $`84`$ and $`240`$ $`\mu `$s, respectively. This result is in good agreement with the experimental results when statistical errors and uncertainties in the initial MOT parameters are taken into account. The agreement between experimental and simulated trap lifetimes, in both mean and distribution, gives an indication of the validity of the theoretically calculated trapping potential and diffusive forces on the atom. The 3.5-fold increase in observed lifetimes due to trapping is made possible by the cavity QED interaction, which allows creation of a deep trapping potential without correspondingly large diffusion as in the free-space case.
### B Oscillations and radial motion
We now turn to a more detailed investigation of the dynamics of motion experienced by a trapped atom. As we have seen, the transmission signal for a single trapped atom exhibits large variations over time which may be tentatively identified with atomic motion in the radial (Gaussian) dimensions of the cavity field. Thus, for example, the highest transmission occurs when the atom passes closest to the cavity axis, $`\rho =0`$. To determine the validity of such an identification, we examine the periods of observed oscillation in the transmission signal. The calculated effective potential is approximately Gaussian in the radial dimension, so a one-dimensional conservative-motion model predicts periods as a function of oscillation amplitude in this anharmonic effective potential well. Referring to the sample transits of Fig. 8, one does indeed note a trend toward large modulations with long periods and smaller modulations with shorter periods. To quantify this observation, we plot period $`P`$ versus the amplitude $`A`$ for individual oscillations, where $`A2[(H_1+H_2)/2H_c]/(H_1+H_2)`$, with $`\{H_1,H_2,H_c\}`$ as indicated in Fig. 8.
Figure 10(a) shows the experimental data plotted along with the calculated curve for one-dimensional motion in the effective potential $`U(\rho ,0)`$ (see Fig. 3), for the same parameters as Fig. 8. (This is a different data set from that presented in Fig. 4. of Ref .) Note that since an atom approaches the cavity axis $`\rho =0`$ twice over the course of one orbital period, the predicted period for oscillations in the transmission signal is half the period of the underlying atomic motion. Experimental data clearly map out this calculated curve for radial atomic motion, demonstrating that oscillations in the observed cavity transmission do indeed reflect radial position of an atom as it varies over time within the trap. The agreement also indicates the quantitative correctness of the theoretical model for the radial potential depth and spatial profile. Note that the comparison is absolute with no adjustable parameters.
The same analysis may be performed for transmission oscillations in the set of simulated transits, yielding the plot of Fig. 10(b). This plot again shows agreement with the calculated curve, with some spread away from the line. For simulated transits, it is possible to turn to the underlying atomic position record to determine an angular momentum for the atom during a given oscillation. Thus the oscillation data of Fig. 10(b) are plotted by atomic angular momentum, where lower angular momentum data points are shown with circles. A separation by angular momentum is clearly evident, with lower angular momentum points most closely following the calculated one-dimensional (and thus zero angular momentum) curve. This separation, while it may seem expected, is in fact a non-trivial indication that angular momentum is a valid quantity for the atomic motion over the course of an oscillation period. Since the atomic motion is not in fact conservative, but is also influenced by random (diffusive) forces, a separation by angular momentum can only be expected to occur if the effect of diffusive forces is sufficiently small over the time scale of an orbit in the conservative potential. The plots of Fig. 3 provide an initial indication that this is indeed the case for these parameters, and this idea is borne out by the current investigation. Confidence in the relatively small effect of diffusion over a single orbital period is crucial in the reconstruction of two-dimensional atomic trajectories as in Ref..
## V Simulation Results for the Experiment of Pinkse et al.
Having provided a validation of our capabilities for numerical simulation by way of the results of Sec. IV, we next apply this formalism to the experiment reported in Ref. . At the outset, we note that the various approximations discussed in Sec. II related to the derivation of this quasiclassical model are satisfied to a better degree for this experiment than for the experiment of Ref. . Hence we expect that the correspondence between the simulations and experiment should be at least of the quality as in the preceding section.
Our starting point is the generation of a large set of simulated trajectories for the parameters reported in Ref. , namely, $`(g_0,\gamma ,\kappa )=2\pi (16,3,1.4)`$ MHz with detuning parameters $`\mathrm{\Delta }_{ac}=\omega _{cavity}\omega _{atom}=2\pi \times 35`$ MHz and $`\mathrm{\Delta }_{probe}=\omega _{probe}\omega _{atom}=2\pi \times 40`$ MHz. The initial pretriggering level of the driving laser gives a 0.15-photon mean intensity in the empty cavity; when this level rises to 0.85 photons, indicating the presence of an atom, we trigger an increase in the driving strength to a trapping level of 0.9-photon empty-cavity intensity. These criteria are intended to follow the parameters indicated in Figs. 2 and 3 of Ref.. Note that for the cavity geometry of this experiment, the time for an atom to transit freely through the cavity mode in the absence of any light forces is $`\tau _0=2w_0/\overline{v}=290`$ $`\mu `$s, where as before we take twice the cavity waist $`w_0`$ as a measure of the transverse dimension of the cavity.
### A Histograms of transit durations
From the set of such simulated trajectories ($`400`$ in this particular case), we can construct histograms for the number of events as a function of total transit signal duration. Following the experimental protocol of Ref., which employed photon counting, we base this analysis upon the intracavity photon number $`\overline{n}=a^{}a`$ rather than$`|a|^2`$ as in Ref., although this distinction is not critical to any of the following considerations.
The resulting histograms for the experiment of Ref. are displayed in Fig. 11 for the cases of untriggered and triggered trajectories. As in the discussion of Fig. 9, the external drive strengths are set to be equal for this comparison to provide equal detectability for an atom passing through the cavity mode. Detection with lower external drive strength gives a lower signal-to-noise ratio for atom detection, which results in detected transit durations much shorter than the actual passage time through the cavity (which is of order $`\tau _0=2w_0/\overline{v}`$), as for example in Fig. 2(a) of Ref..
In support of the validity of our simulations for the experiment of Pinkse et al. (including the initial atomic velocity and position distribution and the triggering conditions), note that the mean of $`280`$ $`\mu `$s for the histogram in the triggered case of Fig. 11(b) corresponds quite well with that quoted in Ref., namely, $`\overline{\tau }_{\mathrm{exp}}=250\pm 50`$ $`\mu `$s. Further, the histograms in Fig. 11 exhibit an extension of the mean transit duration from $`160`$ $`\mu `$s for the case of no triggering in (a) to $`280`$ $`\mu `$s with triggering in (b), in support of the claim of trapping in Ref.. The dispersion of events around the mean is quite large in both cases, $`161`$ $`\mu `$ s in the untriggered set and $`282`$ $`\mu `$s in the triggered set. The increase in the mean is largely associated with an increase in the number of events in the range 200-300 $`\mu `$s, as well as in the number of rare events much longer than the mean duration. Once again we note that the dissipative nature of the dynamics plays a crucial role in the observed motion for the experiment of Pinkse et al. A histogram of transit durations calculated with the sign of the friction coefficient reversed has a lower mean than that of transits with no triggering.
However, it is certainly worth noting that the observed “average trapping time” $`\overline{\tau }_{\mathrm{exp}}=250\pm 50`$ $`\mu `$s quoted in Ref. , as well as the corresponding mean time from our simulations, are smaller than the time $`\tau _0=290`$ $`\mu `$s for an atom to transit freely through the cavity mode. Additionally, even in the case of no triggering, there is already a significant number of events with similar long duration to those in (b) with triggering. Such events arise from the relatively large contribution of diffusion-driven fluctuations whereby an atom randomly loses a large fraction of its initial kinetic energy as it enters the cavity. That such fluctuations play a critical role should already be clear from the plots of the confining potentials and diffusion coefficients in Fig. 4.
### B Radial motion
Trapping dynamics can also be explored if atomic oscillation in the trapping potential can be directly observed. Certainly the observations presented in Fig. 10 make this case for the experiment of Ref., with the observed oscillation frequencies found to be in good quantitative agreement with those computed directly from the anharmonic potential without adjustable parameters and with the results of the numerical simulations.
Towards the goal of constructing a similar plot for the parameters of Ref. , consider a long-duration transit event such as that in Fig. 7(c). Recall that the output flux from the cavity is given by the cavity decay rate $`2\kappa _d`$ into the relevant detection channel times the intracavity photon number, or $`I=2\kappa _d\overline{n}=2\kappa _da^{}a`$, with then the detected count rate found from the overall propagation and detection efficiency as $`R=\xi I`$. Of course, in any actual experiment the full information displayed for the intracavity photon number $`\overline{n}`$ is not available because of finite detection efficiencies $`(\xi <1)`$ and the requirement to average over many cavity lifetimes in order to achieve an acceptable signal-to-noise ratio (roughly for a time such that $`\sqrt{R\delta t}>1)`$.
Rather than attempt a detailed analysis of such effects for the experiment of Ref., here we wish to illustrate several generic effects that hinder definitive observation of radial oscillations in this regime. We therefore take the full ideal signal $`\overline{n}(t)`$ with no degradation due to cavity escape efficiency or subsequent system losses (which we estimate to be $`\kappa _d/\kappa 0.17`$ and $`\xi 0.6`$for an overall efficiency of 0.11). As shown in Fig. 12, to this ideal signal we apply a low-pass filter with cutoff $`f_c=10`$ kHz intended to optimize the visibility of any radial oscillations for frequencies $`f5`$ kHz, where $`f_0^{(r)}=2.6`$ kHz is the orbital frequency for small-amplitude oscillation near the bottom of the radial potential. As before, recall that a periodic variation in the radial coordinate at frequency $`f`$ results in a variation in $`\overline{n}`$ at $`2f`$. Precisely such a filtering protocol was implemented for the analysis in Fig. 10, there with $`f_c=25`$ kHz in correspondence to the larger radial oscillation frequencies ($`f_0^{(r)}=9.4`$ kHz for Ref.) .
Not surprisingly, the frequent and large bursts of axial heating evident for the simulated trajectories of Fig. 7 result in large variations in the intracavity photon number on time scales set by twice the axial oscillation frequency $`f_0^{(a)}430`$ kHz. While these axial oscillations cannot be directly resolved in the detected counting signal $`R(t)`$, their envelope nonetheless leads to variations in $`\overline{n}(t)`$ and hence $`R(t)`$ on time scales comparable to that associated with radial motion (i.e., $`1/2f_0^{(r)}`$), as is apparent in Fig. 12. Consequently, the low-pass filtering \[or, equivalently, the time averaging over segments in $`R(t)`$\] that is required experimentally to obtain an acceptable signal-to-noise ratio gives rise to observed variations in $`\overline{n}(t)`$ that can arise from either axial or radial atomic motion. In the particular transit shown in Fig. 12, two apparent variations on time scales $`200`$ $`\mu `$s are introduced by a filtering of the axial motion, whereas the longer modulation ($`600\mu `$s duration) does reflect the radial position of the atom. This is something of a generic feature of the several hundred simulated transits examined; shorter-time-scale modulations ($`300`$ $`\mu `$s) can reflect either a genuine radial excursion or a filtering of axial motion, whereas very long period variations (500–600 $`\mu `$s) are indicative of radial atomic motions. This simply reflects the fact that the bursts of axial motion tend to have time scales limited to a few $`100`$ $`\mu `$s.
To illustrate these points further, we have constructed a plot of period versus normalized amplitude of transmission oscillations from our simulations of the experiment of Pinkse et al., with the result given in Fig. 13. We emphasize that the protocol followed is precisely as for the analysis that led to Fig. 10(b) for the experiment of Hood et al. (see also Fig. 4 of Ref.), with the exception of the aforementioned reduction in the low-pass cutoff frequency. In marked contrast to that case, here there is a poor correspondence between the distribution of orbital periods from the ensemble of simulated trajectories and the prediction from the potential obtained from Eq. (12). Referring to the discussion of Fig. 12 above, we note that about 2/3 of the points in the 100–300 $`\mu `$s range result from averaging over axial motion, whereas for longer-period (P¿300 $`\mu `$s) modulations, 80% of the observed points reflect changes in the radial motion, but with associated transmission amplitude typically modified by the presence of axial motion. The results of Fig. 13 \[which are for the ideal case of $`\overline{n}(t)`$ without signal degradation due to finite escape and detection efficiency\] suggest that only in restricted cases can temporal variations in $`R(t)`$ be attributed to radial motion, and not instead of (or in addition to) the envelopes of axial heating processes. Indeed, such effects are well known in the literature, having been previously discussed for the case of individual atoms falling through the cavity mode (albeit without triggering or trapping) . A similar conclusion was reached, namely, that axial heating processes contaminate the frequency band associated with radial motion, thereby precluding inferences about radial motion. For the data presented by Pinkse et al. , the long ($`500`$ $`\mu `$s) time scale of the modulations suggests an assignment of these signals to radial motion; however, a more detailed characterization of the atom dynamics over a larger ensemble of transits should yield this more definitively.
It is also worth noting that the quoted average trapping time $`\overline{\tau }_{\mathrm{exp}}=250\pm 50`$ $`\mu `$s in Ref. is itself less than $`1/f_0^{(r)}=390`$ $`\mu `$s, which is shortest time for a full radial orbit. Hence any conclusion about motion in the radial plane must necessarily be based upon rare events in the tail of the histograms of Fig. 11. The rare occurrence of these long events is reflected in the small number of data points in Fig. 13, which was constructed from the same number of simulated transits as Fig. 10(b).
### C Axial Motion
We next turn to analyze motion along the axial direction, and to the statement of Pinkse et al. that Fig. 4 of Ref. “is direct evidence for the atom moving along the cavity axis,” as opposed to instances of localization around an antinode for which “hardly any periodic structure is visible.” In their analysis, Pinkse et al. employed a function $`g^{(4)}(ϵ,\tau ,ϵ)`$, whose intention is to pick out two-time correlations in intensity, with an enhanced signal-to-noise ratio of intensity fluctuations by measuring coincidences of photon pairs. Here we attempt to investigate manifestations of the axial motion independent of the details of any specific such function by analyzing $`\overline{n}(t)`$ directly by way of a windowed fast-Fourier transform (FFT). More specifically, for each trajectory from a large ensemble from our simulations, we apply a FFT to the record $`\overline{n}(t)`$ with a Hanning window centered at time $`t_i`$ and of total width $`25`$ $`\mu `$s, with the window then offset sequentially to $`t_{i+1}=t_i+5`$ $`\mu `$s to cover the whole range of a given atomic trajectory. The window width $`25`$ $`\mu `$s is chosen to be in close correspondence to the record length of $`20`$ $`\mu `$s employed by Pinkse et al. Longer window widths do not qualitatively change the results of our analysis, while a substantially shorter-duration window leads to a loss of requisite frequency resolution.
Two examples from an extended set of such transforms are given in Figs. 14 and 15. Parts (a) of each of these figures show the mean intracavity photon number $`\overline{n}(t)`$, the axial coordinate $`x(t)`$, and a contour plot of the windowed FFT $`𝒩_{t_i}(\mathrm{\Omega })`$ for a single atomic trajectory for the parameters of Ref.. Here $`𝒩_{t_i}(\mathrm{\Omega })`$ is the windowed FFT of $`\overline{n}(t)`$ over the entire duration of the trajectory, with $`t_i=t_0+i\times 5`$ $`\mu `$s. Parts (b) of Figs. 14 and 15 compare $`𝒩_{t_i}(\mathrm{\Omega })`$ for two particular values of $`t_i`$, namely, at a time $`t_{flight}`$ corresponding to the midst of a flight of the atom over several antinodes of the intracavity standing wave (i.e., variations in axial coordinate $`x`$ by several units of $`1\lambda /2`$) and at a time $`t_{localized}`$ for which there is appreciable heating along the axial direction but for which there is no flight (i.e., the atom remains localized within the same axial well). The times $`(t_{flight},t_{localized})`$ are indicated by the arrows in the top two panels of parts (a).
Perhaps the most striking aspect of the comparison of the spectral distributions $`\{𝒩_{flight}(\mathrm{\Omega }),𝒩_{localized}(\mathrm{\Omega })\}`$ for the cases with and without flight is their remarkable similarity \[in (b) of Figs. 14 and 15\]. Both display prominent peaks near $`\mathrm{\Omega }_p/2\pi =f_p`$500–600 kHz, which is in accord with the expected frequency for large-amplitude oscillation in the axial potential, for which the harmonic frequency $`f_0^{(a)}430`$ kHz (recall that frequency of atomic dynamics is half the frequency of the associated variations in $`\overline{n}(t)`$). This result is also in accord with that from Fig. 4(b) of Pinkse et al., for which their simulation leads to $`1/\tau _p`$550 kHz for variations in the function $`g^{(4)}`$.
However, our analysis, as in the comparison of $`\{𝒩_{flight}(\mathrm{\Omega }),𝒩_{localized}(\mathrm{\Omega })\}`$ above, indicates that neither the observation of a peak in $`𝒩(\mathrm{\Omega })`$ around $`\mathrm{\Omega }_p`$ nor of oscillatory structure in $`g^{(4)}(ϵ,\tau ,ϵ)`$ around $`\tau _p2\pi /\mathrm{\Omega }_p`$ is sufficient to justify direct evidence for the atom moving along the cavity axis. Rather, peaks in $`𝒩_{t_i}(\mathrm{\Omega })`$ are ubiquitous around frequencies $`\mathrm{\Omega }_p/2\pi `$500–600, and appear whether the atom’s motion is localized (but heated) within a given axial well or whether the atom is in flight across several wells. This feature follows from an analysis of the full record of $`\overline{n}(t)`$ without the deleterious effects of finite escape and detection efficiency, or of finite detection bandwidth. Such a result suggests that the measurements of Fig. 4 in Ref are not in and of themselves sufficient to establish unambiguous observation of atomic motion across several wells of the cavity field standing wave.
Our analysis does suggest that it may still be possible to distinguish between axial motion confined within a well and flight along the cavity axis through a more careful quantitative analysis of the respective spectral distributions $`\{𝒩_{flight}(\mathrm{\Omega }),𝒩_{localized}(\mathrm{\Omega })\}`$. With reference to Figs. 14 and 15, note that a principal distinction between these cases is that in the case of flight there is a large decrease of spectral content in the lowest frequency components around $`\mathrm{\Omega }=0`$. This decrease reflects the fact that axial skipping causes full-range variation in $`g`$, and thereby pulls down the time-averaged value of transmission $`\overline{n}(t)`$. In addition, we note an increase in $`𝒩_{flight}(\mathrm{\Omega })`$ as compared to $`𝒩_{localized}(\mathrm{\Omega })`$ for Fourier components in a broad range around $`\mathrm{\Omega }_p/2`$ and up to $`\mathrm{\Omega }_p.`$ The increase appears to reflect atomic motion that, during skipping, explores the full nonlinear (anharmonic) range of the axial potential. These characteristics of the overall spectral distributions seem to discriminate more reliably between flight and localized heating than does a single-frequency peak criterion; they may still offer an avenue for observing atomic skips across the standing wave.
## VI Conclusions
A principal objective of this paper has been to investigate the extent to which light-induced forces in cavity QED are distinct from their free-space counterparts. Our perspective has been to seek qualitatively new manifestations of optical forces at the single-photon level within the setting of cavity QED. Note that the importance of a quantum character for the relevant fields or phenomena is not ensured by the statement that the mean photon number $`\overline{n}1`$, since this is trivially the case in an equivalent free-space volume for a field of the same intensity as that inside the cavity.
As a starting point, we have presented comparisons between the effective potential $`U_{eff}(\rho ,x)`$ in cavity QED and the corresponding free-space potential, as well as of the diffusion coefficients in both contexts (Figs. 3 and 4). Perhaps surprisingly, even in a regime of strong coupling as in Ref., there are only small differences between the cavity QED and free-space potentials and diffusion coefficients. Note that the comparison of Fig. 4 includes “the back action of the atom on the cavity field”, and yet there are nonetheless no substantive differences between the cavity QED and free-space cases for the experiment of Pinkse et al. Hence, although the cavity QED interactions do bring a substantial advantage for atomic detection within the cavity volume, we conclude that the claim of trapping an atom with single photons in Ref. involves no new characteristics unique to the cavity QED environment, with the conservative forces and diffusion largely described by the well-known free-space theory (Fig. 4). Friction which enhances trapping in this regime can be ascribed to cavity-mediated cooling effects, which are in themselves not uniquely features of the quantized-field treatment. However, more analysis is required to determine if the observed effects of friction do indeed rely on the cavity-field quantization.
By contrast, for the experiment of Hood et al., a comparison of the free-space theory and its cavity QED counterpart demonstrates that the usual fluctuations associated with the dipole force along the standing wave are suppressed by an order of magnitude. A semiclassical treatment of the cavity field yields large diffusions like those calculated for the free-space trap. Indeed, if it were not for the reduction of heating in the quantized cavity QED case, an atom would be trapped for less than the period of a single radial orbit before being heated out of the well for the parameters of Ref. . Our calculations support the conclusion that the suppression in dipole-force heating is based upon the Jaynes-Cummings ladder of eigenstates for the atom-cavity system, which to our knowledge represents qualitatively new physics for optical forces at the single-photon level within the setting of cavity QED.
In terms of a more complete analysis, the effective potential $`U_{eff}(\rho ,x)`$ and the diffusion coefficient $`D(\rho ,x)`$ are important ingredients in the quasiclassical theory that we have developed for atomic motion in cavity QED. By way of detailed, quantitative comparisons with the experiment of Hood et al. in Sec. IV, we have validated the accuracy and utility of our numerical simulations based upon the quasiclassical theory. As part of this comparison, we have demonstrated agreement between experiment and simulation for histograms of the duration of transit events, with mean $`\overline{\tau }_t=340`$ $`\mu `$s for the histogram in the triggered case of Fig. 9b extended well beyond the mean $`\overline{\tau }_u=92`$ $`\mu `$s for the untriggered case. Furthermore, $`\overline{\tau }_t`$ exceeds the transit time $`\tau _0=110`$ $`\mu `$s for an atom to transit freely through the cavity mode. The simulated trajectories of Fig. 6 together with the comparison of Fig. 10 for the experiment of Hood et al. strongly support the conclusion that atomic motion is largely conservative in nature, with only smaller contributions from fluctuating and velocity-dependent forces. Atomic motion is predominantly in radial orbits transverse to the cavity axis. The (suppressed) axial heating is important, but only towards the end of a given trajectory leading to ejection from the trap. Knowledge of the time dependence $`\rho (t)`$ for the radial coordinate (by way of the detected field emerging from the cavity and the solution of the master equation) as well of the confining potential $`U(\rho ,0)`$ allow an algorithm to be implemented for inference of the actual atomic trajectory, as demonstrated in Ref. and discussed in greater detail in Ref..
In the case of Ref., numerical simulations for the parameters appropriate to this experiment lead to histograms with mean $`280`$ $`\mu `$s in the triggered case of Fig. 11(a) and $`160`$ $`\mu `$s for the untriggered case of Fig. 11(b), which should be compared to the time $`\tau _0=290`$ $`\mu `$s for an atom to transit freely through the cavity mode in this experiment. The simulated transits of Fig. 7 indicate that atomic motion in this case is dominated by diffusion-driven fluctuations in both the radial and axial dimensions with friction playing an important role in the axial direction. The character of the motion hampers inference of atomic motion from the record of intracavity photon number. Axial heating leads to repeated large bursts of axial excursions during an atomic transit, and hence to large oscillations in the intracavity photon number $`\overline{n}(t)`$. The envelopes of these oscillations have appreciable Fourier content in the range of interest for observation of radial motion, so that there is not an unambiguous signature for the radial motion in the record of $`\overline{n}(t)`$ on short time scales, such as those presented in Ref.. Similarly, the result by Pinkse et al. for hopping or flights over the antinodes of the cavity standing wave is not substantiated by a closer inspection of the Fourier content of the relevant signals. As documented in Figs. 14 and 15, similar signals can be observed for an atom localized (but heated) within a single standing-wave well. We emphasize that these conclusions concerning the work of Ref. are based upon the analysis of several hundred simulated trajectories, apparently well beyond the few cases presented in that paper.
Beyond these comments directed to the prior work of Refs., we suggest that the capability for numerical simulation of the quasiclassical model of atom motion in cavity QED should have diverse applications. For example, we are currently applying the simulations to the problem of feedback control of atomic motion. Given the capability to infer an atomic trajectory in real time, it should be possible to apply active feedback to cool the motion to the bottom of the effective potential $`U_{eff}(\stackrel{}{r})`$.
## Acknowledgments
We gratefully acknowledge the contributions of K. Birnbaum, J. Buck, H. Mabuchi, S. Tan, and S. J. van Enk to the current research. This work was supported by DARPA via the QUIC Institute administered by ARO, by the NSF, and by the ONR.
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# Topological Solitons in Discrete Space-Time as the Model of Fermions
## Abstract
In the present paper we discuss arguments, favouring the view that massive fermions represent dislocations (i.e. topological solitons) in discrete space-time with Burgers vectors, parallel to an axis of time. If to put symmetrical parts of tensors of distortions ( i.e. derivatives of atomic displacements on coordinates) and mechanical stresses equal zero, the equations of the field theory of dislocations get the form of the Maxwell equations. If to consider these tensors as symmetrical, we shall receive the equations of the theory of gravitation, and it turns out that the sum of tensor of distortions and pseudo-Euclidean metrical tensor is the analogue of metrical tensor. It is shown that we can also get Dirac equation with four-fermion interaction in the framework of the field theory of dislocations. This model explains quantization of electrical charge: it is proportional to the topological charge of dislocation, and this charge accepts quantized values because of discrete structure of the 4-dimensional lattice.
The concept of discrete (in other words quantized) space-time has been discussed in physics of elementary particles already for a long time. One of variants of this approach is the lattice quantum field theory, within the framework of which significant successes were already achieved. It is based on approximation of space-time by a discrete lattice, which is actually 4-dimensional analogue of crystal lattices. But in the theory of solid state it is well known, that the ideal crystal lattices do not exist in a nature. The real lattices always contain defects, in particular, dislocations. From the geometrical point of view the dislocations are the one-dimensional topological solitons. The first homotopy group which classifies them is the group of translations. Therefore it is quite natural to assume, that the similar defects should exist in 4-dimensional space-time lattice. Dislocation in (2+1)-dimensional discrete space-time with a Burgers vector, parallel to an axis of time, is shown in Figure 1. The dislocation line in the figure is parallel to a Burgers vector.
As is known, solitons are particle-like objects; various soliton models of elementary particles were offered by many authors. Ross has assumed, that fermions represent topological defects in space - time, namely, dislocations with Burgers vectors, parallel to an axis of time and taking only quantized values. Unzicker offered other topological model of electron: a disclination in space-time. The author of paper also used analogy between elementary particles and dislocations in crystals. Our approach is closer to the hypothesis by Ross . But the consideration by Ross was based solely on relation between fermionic spin density and space-time torsion. In the present paper we are bringing to your notice a number of other arguments, favouring the view that fermions represent dislocations in discrete space-time with Burgers vectors, parallel to an axis of time. We assume that the discrete space-time represents a 4-dimensional lattice. There are some objects ( we shall call them ”primary atoms” ), capable to be displaced from equilibrium positions, in lattice sites. We suppose that the movement of these ”primary atoms” obeys only the laws of classical Newton mechanics. Certainly, this naive assumption is not absolutely justified. It is conceivable that the movement of these ”primary atoms” can be described by the much more complex laws. But it turns out that this ”naive” assumption allows by a rather simple way to receive a number of the important formulae of electrodynamics, gravitation theory and quantum field theory. Probably, account of deviations of the laws of ”primary atoms” behaviour from the laws of the classical mechanics will allow to describe other phenomena (for example, strong or weak interactions).
We shall remind, that torsion is, by definition, an antisymmetric part of connection multiplied by two:
$$T^\lambda {}_{\mu \nu }{}^{}=\mathrm{\Gamma }^\lambda {}_{\mu \nu }{}^{}\mathrm{\Gamma }^\lambda {}_{\nu \mu }{}^{}.$$
The integral of this function over any 2-dimensional surface, which intersects a dislocation, is the value referred to as Burgers vector in solid state physics:
$$b^\lambda =T^\lambda {}_{\mu \nu }{}^{}df^{\mu \nu },$$
where $`df^{\mu \nu }dx^\mu dx^\nu dx^\nu dx^\mu `$ is the element of area. From the geometrical point of view the space-time, containing dislocations with Burgers vectors $`b^i`$ is the manifold with torsion $`T^i{}_{ah}{}^{}(x_\zeta )=\frac{1}{2}e_{achd}\tau ^cb^iV^d\delta (x_\zeta x_\zeta ^{\mathrm{\hspace{0.33em}0}})`$, where $`e_{cahd}`$ is the 4-dimensional Levi-Civita symbol, $`\tau ^c`$ is the 4-dimensional unit vector tangential to the dislocation line, $`V^d`$ is the 4-dimensional velocity of dislocations, $`\delta (x)`$ is the Dirac delta function, $`x_\zeta ^{\mathrm{\hspace{0.33em}0}}`$ are the coordinates of the dislocation line. On the other hand in the theory of gravitation it was shown , that the fermions are also sources of torsion. The relationship between a spin density tensor $`S^\lambda _{\mu \nu }`$ and a torsion is described by the formula
$$T^\lambda {}_{\mu \nu }{}^{}=\frac{8\pi G}{c^3}(S^\lambda {}_{\mu \nu }{}^{}+\frac{1}{2}\delta _\mu ^\lambda S_\nu \frac{1}{2}\delta _\nu ^\lambda S_\mu ).$$
$`\left(1\right)`$
Here $`G`$ is Newton gravitational constant, $`c`$ is velocity of light, $`S^\lambda {}_{\mu \nu }{}^{}=v^\lambda S_{\mu \nu }`$, $`S_\mu =S^\lambda _{\mu \lambda }`$, $`v^\lambda `$ is the 4-vector of the particle velocity, $`S_{\mu \nu }`$ is the antisymmetric tensor. Its spatial components form a 3-vector $`s_k=(S^{\mathrm{\hspace{0.33em}23}},S^{\mathrm{\hspace{0.33em}31}},S^{\mathrm{\hspace{0.33em}12}})`$ which is equal to 3-dimensional density of the particle spin in reference frame of rest of this particle. In equations (1) and hereinafter all the indices, unless otherwise specified, assume the values from 0 to 3. Therefore, non-moving particle having spin $`s_k`$, $`k=1,2,3`$ create the same torsion, as non-moving dislocations with tangential vector $`\tau _c=(\tau _{\mathrm{\hspace{0.33em}0}},\tau _k)`$ ( vectors $`s_k`$ and $`\tau _k`$ are parallel) and Burgers vector, parallel to an axis of time. This dislocation represents a line in the 4-dimensional space-time which is coincident with a world line of fermion. At any moment of time this line intersects 3-dimensional physical space only in one point, therefore fermions are observed as particle objects. Continuity of a dislocation line is provided by the topological laws . This fact also was proved with use of the Noether theorem . In continuous space-time torsion, created by fixed particles, can accept continuous set of values and in quantized space-time it accepts discrete one. It is well known that spins of elementary particles are quantized. This fact results, in accordance with the formula (1), in discreteness of torsion, created by fixed particles. It confirms hypothesis by Ross of the discrete structure of space-time and the dislocation nature of fermions.
To the present time the gauge theory of dislocations \[4-6\] is developed very well. Many authors paid attention on analogy between this theory and both electrodynamics and theory of gravitation . There are the equations of incompatibility in continuous theory of dislocations
$$_\gamma _\epsilon u_n\left(x_\zeta \right)_\epsilon _\gamma u_n\left(x_\zeta \right)=\frac{2}{c}T_{\gamma n\epsilon }\left(x_\zeta \right),$$
$`\left(2\right)`$
where $`u_n`$ is the vector of displacement of the 4-dimensional continuum particles (in continuous theory the 4-dimensional lattice is approximated by continuous medium ), $`_\gamma =\frac{}{x^\gamma }`$. Equations (2) are a definition of dislocations and a statement that there are no disclinations in the given lattice (since the term descriptive of the contribution from disclinations is absent in the right-hand side of these equations). In other words, (2) are purely geometrical equations.
Tensor $`\beta _{\epsilon n}=_\epsilon u_n`$ refer to as tensor of distortions. Let us introduce tensor $`F_{\epsilon n}`$ which equals an antisymmetric part of a tensor of distortions multiplied by two:
$$F_{\epsilon n}=\beta _{\epsilon n}\beta _{n\epsilon }.$$
$`\left(3\right)`$
By summing each three equations from equations (2), we obtain the following system of the equations
$$_\mu F_{\alpha \beta }+_\alpha F_{\beta \mu }+_\beta F_{\mu \alpha }=\frac{2}{c}\left(T_{\mu \beta \alpha }+T_{\beta \alpha \mu }+T_{\alpha \mu \beta }\right).$$
$`\left(4\right)`$
Let the vector tangential to a dislocation line is equal to $`\tau _\beta =(\tau _{\mathrm{\hspace{0.33em}0}},\tau _{\mathrm{\hspace{0.33em}1}},0,0).`$ We suppose that the quantities $`b_{\mathrm{\hspace{0.33em}0}}`$ and $`\tau _{\mathrm{\hspace{0.33em}1}}`$ are very small. So modern experimental techniques do not allow nonzero members to be detected in the right side of (4). Then the equations (4) get the form of the first pair of the Maxwell equations
$$e^{\alpha \beta \gamma \delta }_\beta F_{\gamma \delta }=0.$$
$`\left(5\right)`$
Newton second law for particles of the continuum containing dislocations has the form
$$^\gamma \sigma _{g\gamma }=\frac{2}{c}C_{g\epsilon i\lambda }T^{\epsilon i\lambda }.$$
$`\left(6\right)`$
The 4-dimensional tensor of mechanical stresses $`\sigma _{g\epsilon }`$ by definition is equal to
$$\sigma _{g\epsilon }=C_{g\epsilon i\lambda }^\lambda u^i.$$
$`\left(7\right)`$
Here $`C_{g\epsilon i\lambda }`$ is the 4-tensor of elastic modules of 4-dimensional continuum. In other words, Lagrangian of ”elastic” waves in 4-dimensional medium under consideration has the form
$$_{\mathrm{\hspace{0.33em}0}}=\frac{1}{2}C^{i\gamma j\epsilon }\beta _{\gamma i}\beta _{\epsilon j}.$$
$`\left(8\right)`$
By virtue of the fact that the tensor $`C_{g\epsilon i\lambda }`$ is sufficiently large the distinction from zero of the right side of equations (6) can be detected in experiments. If it is possible to neglect a symmetric part of the tensor of mechanical stresses, the equations (6) get the form of the second pair of Maxwell equations. If to consider tensors of distortions and mechanical stresses as symmetric, we shall receive the equations of the gauge translation theory of gravitation (that is theory of gravitation in space-time with zero curvature and nonvanishing torsion ) . It turns out that the sum of tensor of distortions and the Minkowskian (pseudo-Euclidean) metrical tensor is the analogue of metrical tensor in gravitation theory. In a general case of dislocations in the 4-dimensional lattice these tensors are neither symmetrical nor antisymmetrical and then we receive as a consequence of our model variant of the uniform theory of electrodynamics and gravitation, offered by Einstein : tensor of electromagnetic field is an antisymmetric part of metrical tensor. It is simultaneously possible to solve a problem, which considered by Einstein as the main lack of the theory with the asymmetrical metric: to give geometrical definition of particles.
In the theory of dislocations the Frenkel-Kontorova model of a dislocation, based on the account of only one nonlinear member in Lagrangian of an elastic field, is well-accepted . In the one-dimensional case it results in the dislocation equation of motion having the form of sine-Gordon equation. The Lagrangian of an elastic field in this case has the form
$$_{SG}=\frac{1}{2}_\mu \varphi ^\mu \varphi +\frac{m^{\mathrm{\hspace{0.33em}2}}}{\beta ^{\mathrm{\hspace{0.33em}2}}}\left[\mathrm{cos}\left(\beta \varphi \right)1\right].$$
$`\left(9\right)`$
Coleman has rigorously proved equivalence of the sine-Gordon soliton and the fundamental fermion of the massive Thirring model in (1+1) dimensions. The Lagrangian of the massive Thirring model has the form
$$_T=i\overline{\psi }\gamma _\mu ^\mu \psi m_f\overline{\psi }\psi \frac{1}{2}g\left(\overline{\psi }\gamma ^\mu \psi \right)\left(\overline{\psi }\gamma _\mu \psi \right),$$
$`\left(10\right)`$
where $`\psi `$ is Fermi field, $`\gamma ^\mu `$ are Dirac matrices in (1 + 1) dimensions. Lagrangians (9) and (10) are equivalent on condition that
$$\frac{4\pi }{\beta ^{\mathrm{\hspace{0.33em}2}}}1=\frac{g}{\pi }.$$
$`\left(11\right)`$
The correspondence between these two models is established by bosonization relations:
$$\frac{m^{\mathrm{\hspace{0.33em}2}}}{\beta ^{\mathrm{\hspace{0.33em}2}}}\mathrm{cos}\left(\beta \varphi \right)=m_f\overline{\psi }\psi ,$$
$$\frac{\beta }{2\pi }e^{\mu \nu }_\nu \varphi =\overline{\psi }\gamma ^\mu \psi j^\mu .$$
$`\left(12\right)`$
The issues of bosonization are considered in more detail in the book by Rajaraman . In particular, there was shown that the topological charge of the sine-Gordon soliton is equivalent to the fermionic charge of the particle of the massive Thirring model. Thus, the discreteness of the fermionic charge in our approach is a consequence of discreteness of the topological charge of solitons (i. e. dislocations), which in turn directly follows from discrete structure of space-time.
Many authors consider incompatibility of a space-time lattice with a condition of Lorentz invariance as traditional lack of models of discrete space-time. In offered model this contradiction is eliminated. The Lorentz transformations arise in this model by a natural way as a consequence of finiteness of velocity of light. But only values, relating to properties of particles: fields, forces of interaction, and so on depend on velocity by Lorentz law. This dependence is a consequence of occurrence of Lorentz roots in expression for classical Green function of the equations (6). Therefore quantities, expressing through Green function: fields, created by particles, force of interaction between particles, and so on only depend on velocity by Lorentz law. All other quantities, including the parameters of the lattice, are not exposed to any transformations.
In this connection we shall notice that relativistic expressions always occur in soliton theories, in particular, in the theory of dislocations in crystals. In this theory instead of velocity of light velocities of sound appear in the formulae. As in solids even in isotropic case not only transversal, but also longitudinal sonic waves can exist, in the theory of dislocations in certain cases expressions, containing Lorentz roots of different kinds: $`\sqrt{1v^{\mathrm{\hspace{0.33em}2}}/c_\lambda ^{\mathrm{\hspace{0.33em}2}}}`$, $`\lambda =1,2`$ can occur. For example, in case of straight dislocation in isotropic medium parallel to axis z with Burgers vector $`b_i=(b,0,0)`$, moving at the velocity v parallel to axis x, the displacements of particles of continuum are described by the following formulas :
$$\begin{array}{c}u_{\mathrm{\hspace{0.33em}1}}(x,y,t)=bc_{\mathrm{\hspace{0.33em}1}}^{\mathrm{\hspace{0.33em}2}}/\left(\pi v^{\mathrm{\hspace{0.33em}2}}\right)[\mathrm{arctg}(y(1v^{\mathrm{\hspace{0.33em}2}}/c_{\mathrm{\hspace{0.33em}2}}^{\mathrm{\hspace{0.33em}2}})^{\mathrm{\hspace{0.33em}1}/2}/(xvt))+\hfill \\ \\ (v^{\mathrm{\hspace{0.33em}2}}/\left(2c_{\mathrm{\hspace{0.33em}1}}^{\mathrm{\hspace{0.33em}2}}\right)1)\mathrm{arctg}(y(1v^{\mathrm{\hspace{0.33em}2}}/c_{\mathrm{\hspace{0.33em}1}}^{\mathrm{\hspace{0.33em}2}})^{\mathrm{\hspace{0.33em}1}/2}/(xvt))],\hfill \\ \\ u_{\mathrm{\hspace{0.33em}2}}(x,y,t)=bc_{\mathrm{\hspace{0.33em}1}}^{\mathrm{\hspace{0.33em}2}}/\left(2\pi v^{\mathrm{\hspace{0.33em}2}}\right)[(v^{\mathrm{\hspace{0.33em}2}}/\left(2c_{\mathrm{\hspace{0.33em}1}}^{\mathrm{\hspace{0.33em}2}}\right)1)(1v^{\mathrm{\hspace{0.33em}2}}/c_{\mathrm{\hspace{0.33em}1}}^{\mathrm{\hspace{0.33em}2}})^{1/2}\times \hfill \\ \\ \mathrm{ln}((xvt)^{\mathrm{\hspace{0.33em}2}}+(1v^{\mathrm{\hspace{0.33em}2}}/c_{\mathrm{\hspace{0.33em}1}}^{\mathrm{\hspace{0.33em}2}})y^{\mathrm{\hspace{0.33em}2}})+(1v^{\mathrm{\hspace{0.33em}2}}/c_{\mathrm{\hspace{0.33em}2}}^{\mathrm{\hspace{0.33em}2}})^{\mathrm{\hspace{0.33em}1}/2}\mathrm{ln}((xvt)^{\mathrm{\hspace{0.33em}2}}+\hfill \\ \\ (1v^{\mathrm{\hspace{0.33em}2}}/c_{\mathrm{\hspace{0.33em}2}}^{\mathrm{\hspace{0.33em}2}})y^{\mathrm{\hspace{0.33em}2}})].\hfill \end{array}$$
Here $`c_{\mathrm{\hspace{0.33em}1}}=(\mu /\rho )^{\mathrm{\hspace{0.33em}1}/2}`$ is the speed of transversal sound waves, $`c_{\mathrm{\hspace{0.33em}2}}=[(\lambda +2\mu )/\rho ]^{\mathrm{\hspace{0.33em}1}/2}`$ is the speed of longitudinal sound waves, $`\lambda `$ and $`\mu `$ are Lame constants. Such relations already are not Lorentz-covariant in traditional sense. Formulae of the dislocation theory can be Lorentz-covariant only in some special cases: for example, in case of straight dislocation with a Burgers vector parallel to a dislocation line in isotropic solid. Therefore under certain conditions Lorentz invariance violation in offered model is possible. May be, it will allow to explain occurrence in the last years of a number of field theories, not satisfying to a condition of Lorentz invariance.
Thus, in the present paper model, describing electromagnetic and gravitational properties of electrons and positrons, is offered. The inclusion of strong and weak interactions in the model is, apparently, future problem. But it already has allowed to unify advantages of soliton models and theories with discrete (quantized) space-time by removing at the same time their lacks. Model explains quantization of electrical charge: it is proportional to the topological charge of dislocation, and this charge accepts quantized values because of discrete structure of the 4-dimensional lattice. At the same time existence of the lattice allows to avoid occurrence of divergences, in particular, results in finite mass of a particle. Within the framework of the given model the phenomenon of annihilation of particle-antiparticle pair is easily explained. The similar process is annihilation of dislocation pairs in solids. From topological reasons follows, that at a meeting of two dislocations with Burgers vectors, which are equal in magnitude and opposite in direction, ideal structure of a lattice restore. Both dislocations (that is both particles) thus disappear, and their energy is radiated as elastic waves. As follows from above-stated, the waves of antisymmetric distortions are perceived by us as electromagnetic, and symmetric as gravitational. It is important to note, that the relativistic properties of particles occur in this model as consequences of the classical Newton mechanics of a 4-dimensional deformable solid.
ACKNOWLEDGEMENTS
Useful discussions with I.V. Barashenkov, M.B. Mineev-Weinstein and O.K. Pashaev are kindly acknowledged.
REFERENCES
D.K. Ross, “Planck’s Constant, Torsion, and Space-Time Defects”, Int. J. Theor. Phys., 28, 1333-1340 (1989).
A. Unzicker, “Teleparallel space–time with torsion yields geometrization of electrodynamics with quantized charges.” preprint gr-qc/9612061
F.W. Hehl, P. von der Heyde, G.D. Kerlick, J.M. Nester, “General relativity with spin and torsion: Foundations and prospects.” Rev. Mod. Phys., 48, 393-416 (1976).
A.M. Kosevich, “The dynamical theory of dislocations.” UFN, 84, 579-609 \[ Sov. Phys. Uspekhi, 7, 837-854 \] (1964).
A.I. Musienko, V.A. Koptsik, “A New Variant of the Gauge Theory of Linear Defects in Crystals.” Kristallografija, 40, 438-445 \[ Crystallography Reports, 40, 398-405 \] (1995).
A.I. Moussienko, V.A. Koptsik, “The Gauge Theory of Dislocations and Disclinations in Crystals with Multiatomic Lattices.” Kristallografija, 41, 586-590 \[ Crystallography Reports, 41, 550-554 \] (1996).
F.W. Hehl, J.D. McCrea, “Bianchi identities and the automatic conservation of energy-momentum and angular momentum in general-relativistic field theories.” Found. Phys., 16, 267-293 (1986).
M.O. Katanaev, I.V. Volovich, “Theory of defects in solids and three-dimensional gravity.” Ann. Phys. (N.Y.), 216, 1-28 (1992).
A. Einstein, “Einheitliche Feldtheorie von Gravitation und Elektrizität.” Sitzunsber. Preuss. Akad. Wiss., phys.-math. Kl., 414-419 (1925). \[ There is a Russian translation: A. Einstein, Collection of Scientific Papers, Vol. 2, Nauka, Moscow, 171-177 (1966)\].
J.P. Hirth and J. Lothe, Theory of Dislocations, (McGraw-Hill, New York, 1968).
S. Coleman, “Quantum sine-Gordon equation as the massive Thirring model.” Phys. Rev. D, 11, 2088-2097 (1975).
R. Rajaraman, Solitons and instantons, (North-Holland, Amsterdam, 1982).
J. Weertman, J.R. Weertman, “Moving Dislocations.” In: Dislocations in Solids, F.R.N. Nabarro, ed., Vol. 3, Moving Dislocations. (North-Holland, Amsterdam, 1980). p. 1-59.
FIGURE
FIG. 1. Dislocation in (2+1)-dimensional discrete space-time with a Burgers vector, parallel to an axis of time. The dislocation line on the figure is parallel to a Burgers vector.
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# Double Inflation in Supergravity and the Boomerang Observations
## I Introduction
After the discovery by COBE/DMR , anisotropies of Cosmic Microwave Background radiation (CMB) become one of the most important targets of modern cosmology. Theoretical works reveal that they contain rich information, i.e., geometry of the universe, the baryon density, the total matter density, the Hubble parameter, thermal history and so on . However, it turns out that the angular resolution of COBE/DMR was too crude to obtain above information. Precise measurements of so-called acoustic peaks in the angular power spectrum on arcminute scales are necessary to determine these cosmological parameters.
There are many attempts to observe CMB anisotropies on arcminute scales after COBE/DMR discovery. Recently, the boomerang team has reported very clear evidence of the first acoustic peak in the angular power spectrum of CMB anisotropies . The location of the peak suggests the flatness of the universe. However, there remain some mysteries in their result. One of them is a relatively low second acoustic peak if there exists. The most natural way to explain such a low peak is to increase the value of $`\mathrm{\Omega }_Bh^2`$, where $`\mathrm{\Omega }_B`$ is the baryon density parameter and $`h`$ is the non-dimensional Hubble constant normalized by $`100\mathrm{k}\mathrm{m}/\mathrm{s}/\mathrm{Mpc}`$, since increasing $`\mathrm{\Omega }_Bh^2`$ boosts only odd number peaks . In order to fit the data, however, $`\mathrm{\Omega }_Bh^2`$ needs to be larger than the value constrained by Big Bang Nucleosynthesis (BBN) .
There are several possibilities to explain the low amplitude of the second acoustic peak. Among them, a tilted initial power spectrum, increasing the diffusion length , and degenerated neutrinos can be considered. However, it seems that none of them provides a small enough second-first peak height ratio without increasing $`\mathrm{\Omega }_Bh^2`$. Here we propose double inflation which breaks a coherent feature of the initial power spectrum as a candidate for solving this mystery.
Recently we studied a double inflation model with hybrid and new inflations and its cosmological implication . It was found that both inflations could produce cosmologically relevant density fluctuations if the total $`e`$-fold number of new inflation is small enough. In this case, there appears a breaking scale in the density power spectrum which corresponds to the horizon scale of the transit epoch from hybrid to new inflation. The fluctuations on scales larger (smaller) than the breaking are produced by hybrid (new) inflation. Assuming the ’standard’ cold dark matter model with $`\mathrm{\Omega }_0=1`$, where $`\mathrm{\Omega }_0`$ is the density parameter of a matter component, we can fit both cluster abundances and galaxy distributions with the COBE/DMR normalization when we set the amplitude of perturbations produced by new inflation is smaller than the one by hybrid inflation. Accordingly we generally obtain smaller acoustic peaks in the CMB angular power spectrum since the peaks are controlled by new inflation and the Sachs-Wolfe tail on larger scales whose amplitude is fixed by the COBE normalization is determined by hybrid inflation.
Therefore, we may be able to explain cluster abundances, galaxy distributions and the low second acoustic peak at once by introducing double inflation if the braking scale is in between first and second peaks by chance. In this paper, we first explain the double inflation model and compare the resultant density power spectrum and CMB angular power spectrum with observational data. We take $`\mathrm{\Omega }_0+\lambda _0=1`$, i.e., a spatially flat universe. Here $`\lambda _0`$ is the density parameter of a cosmological constant. We also take $`\mathrm{\Omega }_Bh^2=0.02`$ which is the best value from BBN analysis .
## II Double inflation model
We adopt a double inflation model proposed in Ref.. Here we briefly describe the model in order to show the relation between model parameters and the power spectrum of the density fluctuations. The model consists of two inflationary stages; the first one is called preinflation. Here we employ hybrid inflation (see also Ref.) as preinflation. We also assume that the second inflationary stage is realized by a new inflation model and its $`e`$-fold number is smaller than $`60`$. Thus, the density fluctuations on large scales are produced during preinflation and their amplitudes should be normalized to the COBE data . On the other hand, new inflation produces fluctuations on small scales. Thus, this power spectrum has a break on the scale corresponding to a turning epoch from preinflation to new inflation. As for the detailed argument of dynamics of our model, see Refs. .
### A First inflationary stage
First, let us briefly discuss a hybrid inflation model . The hybrid inflation model contains two kinds of superfields: one is $`S(x,\theta )`$ and the others are a pair of $`\mathrm{\Psi }(x,\theta )`$ and $`\overline{\mathrm{\Psi }}(x,\theta )`$. Here $`\theta `$ is the Grassmann number denoting superspace. The model is based on the U$`(1)_R`$ symmetry under which $`S(\theta )e^{2i\alpha }S(\theta e^{i\alpha })`$ and $`\mathrm{\Psi }(\theta )\overline{\mathrm{\Psi }}(\theta )\mathrm{\Psi }(\theta e^{i\alpha })\overline{\mathrm{\Psi }}(\theta e^{i\alpha })`$. The superpotential is given by
$$W(S,\mathrm{\Psi },\overline{\mathrm{\Psi }})=\mu ^2S+\lambda S\overline{\mathrm{\Psi }}\mathrm{\Psi }.$$
(1)
The $`R`$-invariant Kähler potential is given by
$$K(S,\mathrm{\Psi },\overline{\mathrm{\Psi }})=|S|^2+|\mathrm{\Psi }|^2+|\overline{\mathrm{\Psi }}|^2+\mathrm{},$$
(2)
where the ellipsis denotes higher-order terms which we neglect in the present analysis for simplicity. We gauge the U$`(1)`$ phase rotation:$`\mathrm{\Psi }e^{i\delta }\mathrm{\Psi }`$ and $`\overline{\mathrm{\Psi }}e^{i\delta }\overline{\mathrm{\Psi }}`$. To satisfy the $`D`$-term flatness condition we take always $`\mathrm{\Psi }=\overline{\mathrm{\Psi }}`$ in our analysis. For $`|S|>|S_c|=\mu /\sqrt{\kappa }`$, the effective potential $`V`$ has a minimum at $`\mathrm{\Psi }=\overline{\mathrm{\Psi }}=0`$. That is, for $`|S|>|S_c|`$, the energy density is dominated by the false vacuum energy density $`\mu ^4`$ and inflation takes place. We identify the inflaton field $`\sigma /\sqrt{2}`$ with the real part of the field $`S`$.
We define $`N_{\mathrm{COBE}}`$ as the $`e`$-fold number corresponding to the COBE scale and the COBE normalization leads to a condition for the inflaton potential,
$$\left|\frac{V^{3/2}}{V^{}}\right|_{N_{\mathrm{COBE}}}\frac{4\pi \mu ^2\sqrt{N_{\mathrm{COBE}}}}{\lambda }5.3\times 10^4,$$
(3)
where $`V`$ is the inflaton potential obtained from Eqs.(1) and (2) including one-loop corrections. In the hybrid inflation model, density fluctuations are almost scale invariant, $`n_{\mathrm{pre}}1`$, where $`n_{\mathrm{pre}}`$ is a spectral index for a power spectrum of density fluctuations.
### B Second inflationary stage
Now, we consider a new inflation model. We adopt an inflation model proposed in Ref. . The inflaton superfield $`\varphi (x,\theta )`$ is assumed to have an $`R`$ charge $`2/(n+1)`$ and U$`(1)_R`$ is dynamically broken down to a discrete $`Z_{2nR}`$ at a scale $`v`$, which generates an effective superpotential ,
$$W(\varphi )=v^2\varphi \frac{g}{n+1}\varphi ^{n+1}.$$
(4)
The $`R`$-invariant effective Kähler potential is given by
$$K(\varphi )=|\varphi |^2+\frac{\kappa }{4}|\varphi |^4+\mathrm{},$$
(5)
where $`\kappa `$ is a constant of order $`1`$. Hereafter we take $`n=4`$ and $`g=1`$ for simplicity.
An important point on the density fluctuations produced by new inflation is that it results in a tilted spectrum with spectral index $`n_{\mathrm{new}}`$ given by
$$n_{\mathrm{new}}12\kappa .$$
(6)
### C Initial value and fluctuations of the inflaton $`\phi `$
The crucial point observed in Ref. is that preinflation sets dynamically the initial condition for new inflation. We identify the inflaton field $`\phi (x)/\sqrt{2}`$ with the real part of the field $`\varphi (x)`$. Then, the value of $`\phi `$ at the beginning of new inflation is given by
$$\phi _b\frac{\sqrt{2}}{\sqrt{\lambda }}v\left(\frac{v}{\mu }\right)^2.$$
(7)
Therefore, in our model, $`\phi _b`$, the $`e`$-fold number of the second inflation ($`N_{\mathrm{new}}`$), and $`N_{\mathrm{COBE}}`$ are determined by only model parameters. On the contrary, in the other double inflation models $`\phi _b`$ should be put by hand. In our model we have four model parameters $`(\mu ,\lambda ,v,\kappa )`$ among which $`\mu `$ is expressed by the other parameters with use of Eq. (3). Thus there are three free parameters.
### D Numerical results
We estimate density fluctuations in double inflation by calculating the evolution of $`\phi `$ and $`\sigma `$ numerically. For given parameters $`\kappa `$ and $`\lambda `$, we obtain the breaking scale $`k_b`$ and the amplitude of the density fluctuations $`\delta _b`$ produced at the beginning of new inflation. Here, $`k_b^1`$ is the comoving breaking scale corresponding to the Hubble radius at the beginning of new inflation. We can understand the qualitative dependence of $`(k_b,\delta _b)`$ on $`(\kappa ,\lambda )`$ as follows: When $`\kappa `$ is large, the slope of the potential for new inflation is too steep, and new inflation cannot last for a long time. Therefore, the break occurs at smaller scales. In fact, we can express $`k_b`$ as
$$k_b\frac{1}{3000}h\mathrm{Mpc}^1\mathrm{exp}\left[50\frac{1}{\kappa }\mathrm{ln}\left(\sqrt{\frac{\lambda (1\kappa )}{12}}\frac{\mu ^2}{v^2}\right)\right].$$
(8)
As for $`\delta _b`$, we can see from Eq.(3) that as $`\lambda `$ becomes larger, $`\mu `$ also must become larger. In addition, we can show that
$$\delta _b\left(\frac{\delta \rho }{\rho }\right)_{\mathrm{new},k_b}\frac{1}{5\sqrt{6}\pi }\frac{\sqrt{\lambda }\mu ^2}{\kappa v},$$
(9)
for a given $`v`$ (see Ref. ). Thus, we have larger $`\delta _b`$ for larger $`\lambda `$.
In our previous work , we have shown that if $`\lambda 𝒪(10^410^3)\mathrm{and}0.1\kappa 0.2`$ , $`k_b`$ is at a cosmological scale ($`10^3h\mathrm{Mpc}^1k_b1h\mathrm{Mpc}^1`$), and density fluctuations produced during new inflation are not too far from those of preinflation ($`0.1P_{\mathrm{new}}/P_{\mathrm{pre}}10`$). Here $`P_{\mathrm{new}}`$ and $`P_{\mathrm{pre}}`$ refer to the amplitude of the power spectrum of the density fluctuations at $`k_b`$, produced by new inflation and preinflation, respectively:
$$P(k)=\{\begin{array}{c}P_{\mathrm{pre}}\left(\frac{k}{k_b}\right)^1T^2(k)(k<k_b),\hfill \\ P_{\mathrm{new}}\left(\frac{k}{k_b}\right)^{n_{\mathrm{new}}}T^2(k)(k>k_b),\hfill \end{array}$$
(10)
where $`T(k)`$ is a matter transfer function.
## III Comparison with observations
### A Second acoustic peak
The spectral index of new inflation $`n_{\mathrm{new}}`$ is $`n_{\mathrm{new}}12\kappa <1`$ \[see Eq.(6)\]. Also, the amplitude of the density fluctuations on smaller scales, which are produced during new inflation, can be smaller than that on larger scales, which is normalized to the COBE/DMR data. Thus, in our double inflation model, if the breaking scale $`k_b`$ is in between the first and the second peaks of the CMB angular power spectrum, there is a possibility to explain the lower second acoustic peak of the boomerang results.
Here we take $`n_{\mathrm{new}}0.8(\kappa 0.1)`$ as an example. Since the location of the first acoustic peak (the multipole moment $`\mathrm{}200`$) corresponds to the comoving wave number $`k𝒪(10^2)h`$Mpc<sup>-1</sup>, we have searched parameter sets within the parameter range of $`0.001h\mathrm{Mpc}^1k_b0.04h\mathrm{Mpc}^1`$, and $`0.71`$, which can produce a lower second acoustic peak of the CMB angular power spectrum.
In Table I, samples of these parameters are listed. From the recent observations of Type Ia supernovae, $`\mathrm{\Omega }_0`$ is estimated as $`\mathrm{\Omega }_00.5`$ . Therefore, we take $`\mathrm{\Omega }_0=0.4,0.5,`$ and $`0.6`$, for example. Also, in Fig.1, we plot the angular power spectrum of the CMB anisotropies for these parameters. From this figure, we can see that our double inflation model has a parameter region which can explain the lower second acoustic peak of the boomerang observations as expected.
### B Cluster abundances and galaxy distribution
As we have seen in the previous subsection, our double inflation model can explain the low second acoustic peak of the boomerang data. However, we need to check whether it is consistent with other observations. In this subsection we compare the result of our double inflation model with the observations of the cluster abundances and galaxy distributions .
Usually the constraint on the power spectrum from observations of the cluster abundances is expressed in terms of $`\sigma _8`$, the specific mass fluctuations within a sphere of a radius of $`8h^1\mathrm{Mpc}`$. Since the power spectrum of the density fluctuations shows a break on the cosmological scale in our double inflation model, we cannot simply employ the value of $`\sigma _8`$ quoted by previous works . We need to calculate the cluster abundances by using the Press-Schechter theory .
When we determine the breaking scale $`k_b`$, the power spectrum ratio $`P_{\mathrm{new}}/P_{\mathrm{pre}}`$, and the spectral index for new inflation $`n_{\mathrm{new}}`$, we can get the power spectrum up to normalization $`A_{\mathrm{cl}}`$. Using this power spectrum we can calculate the comoving abundance of the clusters as
$$n(>M_{\mathrm{min}};A_{\mathrm{cl}})=_{M_{\mathrm{min}}}^{\mathrm{}}\frac{dn(M)}{dM}dM,$$
(11)
where mass distribution $`dn/dM`$ is obtained by the Press-Schechter formula.
Many clusters of galaxies are observed with use of x-ray fluxes. Under the assumption that clusters are hydrostatic, we can obtain the mass-temperature relations as
$$T_{\mathrm{gas}}=\frac{9.37\mathrm{keV}}{\beta (5X+3)}\left(\frac{M}{10^{15}h^1M_{}}\right)^{2/3}(1+z)\mathrm{\Omega }_0^{1/3}\mathrm{\Delta }_c^{1/3},$$
(12)
where $`\mathrm{\Delta }_c`$ is the ratio of the density of a cluster to the background mean density at that redshift, $`\beta `$ is the ratio of specific galaxy kinetic energy to specific gas thermal energy, and $`X`$ is the hydrogen mass fraction. Following Ref. , we take $`X=0.76`$, $`\beta =1`$. Also, $`\mathrm{\Delta }_c`$ can be approximated as $`\mathrm{\Delta }_c18\pi ^2\left[1+0.4093\left(1/\mathrm{\Omega }_01\right)^{0.9052}\right]`$ .
The observed cluster abundance as a function of x-ray temperature can be translated into a function of mass using Eq. (12). Accumulating the observations, Henry and Arnaud gave the fitting formula as
$$\left(\frac{dn(T)/dT}{h^3\mathrm{Mpc}^3\mathrm{keV}^1}\right)=1.8\left\{\begin{array}{c}+0.8\hfill \\ 0.5\hfill \end{array}\right\}\times 10^3\left(\frac{kT}{1\mathrm{k}\mathrm{e}\mathrm{V}}\right)^{4.7\pm 0.5}.$$
(13)
Ref. also gave a table of cluster observations whose temperatures are larger than $`2.5`$ keV, which determines the lower limit $`M_{\mathrm{min}}`$ from Eq. (12). Therefore, by integrating Eq. (13) we obtain
$$6.6\times 10^6n(>M_{\mathrm{min}})4.3\times 10^5.$$
(14)
Matching these abundances, Eq. (11) calculated from the Press-Schechter theory, and Eq. (14) inferred from the x-ray cluster observations, we can determine the normalization (amplitude) of power spectrum, $`A_{\mathrm{cl}}`$. Using this normalization, we can obtain “cluster abundance normalized” $`\sigma _8`$, $`\sigma _{8,\mathrm{cl}}`$, as
$$\sigma _{8,\mathrm{cl}}^2_0^{\mathrm{}}\frac{k^3}{2\pi ^2}P(k;A_{\mathrm{cl}})W^2(kr_0)\frac{dk}{k}|_{r_0=8h^1\mathrm{Mpc}}.$$
(15)
where $`P(k;A)`$ is a present matter density fluctuation power spectrum with a normalization $`A`$, and $`W(x)`$ is a window function. Because of errors in observations, we have some range for allowed $`\sigma _{8,\mathrm{cl}}`$.
On the other hand, we can normalize the power spectrum by COBE data . Therefore, we have “COBE normalized” $`\sigma _8`$, $`\sigma _{8,\mathrm{COBE}}`$ together with $`\sigma _{8,\mathrm{cl}}`$. Bunn and White estimates one standard deviation error of COBE normalization to be $`7\%`$ which is much smaller than the one of cluster normalization. We assume that, therefore, if $`\sigma _{8,\mathrm{COBE}}`$ lies in an allowed $`\sigma _{8,\mathrm{cl}}`$ range, the parameter region of $`k_b`$, $``$, and $`n_{\mathrm{new}}`$ is consistent with the cluster abundance observations. As for the parameter sets for models (a) to (c) in Table I, we have confirmed that they all satisfy the cluster abundance constraint.
We also have to investigate whether our parameter sets are consistent with the observations of galaxy distributions. There are many observations which measure the density fluctuations from galaxy distributions. Among them we use the data sets compiled by Vogeley from Refs. in this paper.
Employing the COBE normalization, we can determine the power spectrum with its overall amplitude if we fix the breaking scale $`k_b`$, the power spectrum ratio $``$, and the spectral index for new inflation $`n_{\mathrm{new}}`$. One might make direct comparison of this power spectrum with above observations of galaxy distributions. However, distribution of luminous objects such as galaxies could differ from underlying mass distribution because of so-called bias. There is even no guarantee that each observational sample has the same bias factor. Therefore, we only consider the shape of the power spectrum here. We change the overall amplitude of each set of observations arbitrarily. Thus, we estimate the goodness of fitting by calculating $`\chi ^2`$ of this power spectrum with fixing $`k_b,`$, and $`n_{\mathrm{new}}`$.
For each parameter set we have chosen in the previous subsection, they fit the observations of galaxy distributions well (reduced $`\chi _{\mathrm{gal}}^21`$), except for model (a) \[see Table I\]. In Fig. 2, we plot the power spectrum of the density fluctuations for the model (b) in Table I as an example.
## IV Conclusions and discussions
The boomerang team has reported that there is a low second acoustic peak in the angular power spectrum of CMB anisotropies. Although there are some explanations to this lower peak, they seem to need higher baryon density than predicted by the Big Bang Nucleosynthesis. In this paper, we have considered the double inflation model in supergravity, and shown that the density fluctuations produced by this inflation model can produce this low second acoustic peak.
Since the density fluctuations in our model has a nontrivial spectrum, we have checked that it is consistent with the observations of the cluster abundances and the galaxy distributions. We have found that the fit to the data in our model is very good if we take $`\mathrm{\Omega }_00.5`$. We can conclude that the double inflation model can account for the boomerang data without conflicting other observations. In particular, we stress that our model does not require high baryon density and hence is perfectly consistent with BBN.
T. K. is grateful to K. Sato for his continuous encouragement. A part of this work is supported by Grant-in-Aid of the Ministry of Education and by Grant-in-Aid, Priority Area “Supersymmetry and Unified Theory of Elementary Particles” (#707).
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# Gamow-Teller Strengths of the Inverse-Beta Transition 176Yb→¹⁷⁶Lu for Spectroscopy of Proton-Proton and other sub-MeV Solar Neutrinos
## Abstract
Discrete Gamow-Teller (GT) transitions, <sup>176</sup>Yb$`^{176}`$Lu at low excitation energies have been measured via the (<sup>3</sup>He,$`t`$) reaction at 450 MeV and at 0. For <sup>176</sup>Yb, two low-lying states are observed, setting low thresholds Q($`\nu `$)=301 and 445 keV for neutrino ($`\nu `$) capture. Capture rates estimated from the measured GT strengths, the simple two-state excitation structure, and the low Q($`\nu `$) in Yb–Lu indicate that Yb-based $`\nu `$-detectors are well suited for a direct measurement of the complete sub-MeV solar electron-neutrino ($`\nu _e`$) spectrum (including $`pp`$ neutrinos) where definitive effects of flavor conversion are expected.
26.65.+t, 25.55.Kr, 27.70.+q
The major result of experimental solar-neutrino research to date shows that the observed fluxes of solar neutrinos are much reduced compared with theoretical predictions based on the standard solar model (SSM) . Particularly sharp questions are posed by the results from the Gallex and Sage experiments at the low-energy end of the spectrum (from the $`p+p`$, <sup>7</sup>Be and CNO reactions in the sun) that are in contrast with those from Super-Kamiokande at high energies (from the decay of <sup>8</sup>B) . The low-energy results imply a negligible flux of <sup>7</sup>Be neutrinos which is inconsistent with the observed sizeable flux of high-energy <sup>8</sup>B neutrinos because the <sup>8</sup>B activity in the sun cannot arise without the precursor <sup>7</sup>Be in the reaction <sup>7</sup>Be+$`p^8`$B. These results suggest mechanisms beyond possible astrophysical shortcomings of the SSM. The only other possibility is non-standard neutrino physics, viz. the conversion of the original electron-flavor solar neutrinos ($`\nu _e`$) to undetected $`\mu `$ and $`\tau `$ flavors. Evidence for $`\nu _\mu `$$`\nu _\tau `$ flavor conversion has recently been observed by Super-Kamiokande .
The conceivable conversion mechanisms produce $`\nu _e`$ flux deficits that are energy dependent (typically strongest at low energies) and result in large deficits of either the $`pp`$ and/or the <sup>7</sup>Be $`\nu _e`$ fluxes, the two dominant features of the sub-MeV spectrum. A deficit in the model-independent $`pp`$-neutrino flux and/or a <sup>7</sup>Be $`\nu _e`$ flux measurably smaller than the total all-flavor flux (expected from the Borexino experiment in particular), would provide the most direct proof of flavor conversion. A real-time measurement of the complete $`\nu _e`$ spectrum from the sun and source-specific fluxes is, thus, of central interest for solving the solar-$`\nu `$ problem. Such data are not available yet since the operating low-energy Ga detectors yield only the integral signal rate above a threshold, not the fluxes from specific solar $`\nu `$ sources.
Recently, Raghavan suggested a possible way to construct a low-threshold, flavor-specific scheme for real-time detection of solar neutrinos via neutrino captures based on the charged current mediated Gamow-Teller (GT) transitions $`\nu _e+^{176}`$Yb$`e^{}`$ \+ <sup>176</sup>Lu and $`\nu _e+^{160}`$Gd$`e^{}`$ \+ <sup>160</sup>Tb. The basic idea is to identify absorption events of $`pp`$ and <sup>7</sup>Be neutrinos by a delayed coincidence between a prompt $`e^{}`$ event and a cascade $`\gamma `$-ray via isomeric states in <sup>176</sup>Lu or <sup>160</sup>Tb. The coincidence tag (with time gates $``$10<sup>-7</sup> s) allows suppression of background events by a factor of 10<sup>7</sup>. For the first time, this scheme offers the tools for practical real-time spectroscopy even of $`pp`$–neutrinos despite the formidable backgrounds that have precluded low-energy $`\nu `$ spectroscopy so far. Feasibility studies are in progress for constructing such a solar $`\nu `$ detector LENS (Low-Energy Neutrino Spectroscopy) .
The basic input data for designing LENS are the cross sections of the GT transitions in Yb and Gd for which the weak matrix elements B(GT) must be determined for each of the low-lying 1<sup>+</sup> states excited by neutrino capture. The low energy-excitations in Yb-Lu and Gd-Tb consist of several close-lying known 1<sup>+</sup> states. Clarification of the states relevant to LENS thus requires high-resolution spectra. We have therefore measured high resolution (<sup>3</sup>He,$`t`$) spectra at 0 using a <sup>3</sup>He beam at 450 MeV. This work complements a different approach to find B(GT) via ($`p,n`$) measurements on the same targets.
The (<sup>3</sup>He,$`t`$) reaction at 450 MeV has already been demonstrated to be useful in getting the B(GT) values. Since the full Fermi transition strength B(F)=N–Z is exhausted by the isobaric analog state (IAS), the B(GT) values can be obtained from the singles (<sup>3</sup>He,$`t`$) spectra by using the ratio of the cross sections to the 1<sup>+</sup> states and the IAS . The measurements described below show that the B(GT) values obtained here can be used reliably to estimate solar neutrino rates in LENS.
The (<sup>3</sup>He,$`t`$) experiment was performed at E(<sup>3</sup>He)=450 MeV at the RCNP cyclotron facility of Osaka University using the spectrometer Grand Raiden which was set at $`\theta `$ = 0 with a solid angle of $`\mathrm{\Delta }\mathrm{\Omega }`$ = 1.6 msr. Both the incident <sup>3</sup>He<sup>++</sup> beam and the reaction products entered into the spectrometer set at 0. A metallic target of <sup>176</sup>Yb (96.7 % enrichment, 3.3 mg/cm<sup>2</sup>) was used. The <sup>3</sup>He<sup>++</sup> beam, which has a much small magnetic rigidity than the tritons,was stopped at the internal Faraday cup in the first dipole (D1) magnet. Details of the experimental setup are described in Ref. . Particles analyzed by the spectrometer were detected in the focal-plane detection system consisting of two multi–wire drift chambers (MWDC’s) and two $`\mathrm{\Delta }`$E plastic scintillation counters, whose signals were used for particle identification as well as for triggering the events. The MWDC’s gave information on the arrival positions at the focal plane and the scattering angles of the tritons.
Very good identification with the focal-plane detector of the spectrograph was required to overcome the huge background due to the <sup>3</sup>He<sup>++</sup> $``$ <sup>3</sup>He<sup>+</sup> atomic charge-exchange process , since the Q-values for the (<sup>3</sup>He,$`t`$) reactions on <sup>176</sup>Yb leading to the lowest 1<sup>+</sup> states in the residual nuclei is $``$ 0.3 MeV. For this purpose, we added a plastic scintillation counter with a thickness of 1 mm in front of the two $`\mathrm{\Delta }E`$ counters of 10 mm thickness for improved particle identification. This front counter should be very thin to avoid creating tritons induced by (<sup>3</sup>He,$`t`$) reactions in the scintillator itself. With this technique, we could obtain a spectrum free from the <sup>3</sup>He<sup>++</sup> $``$ <sup>3</sup>He<sup>+</sup> atomic charge-exchange process.
To achieve a high energy resolution of around 100 keV, we installed an energy defining slit in the beam-transport system from the K=400 MeV ring cyclotron to the target position. The image-defining slit was set at the first beam-focusing position of the beam-line. The technique of dispersion matching between the beam-line and the spectrometer was employed to obtain the required high resolution. An overall energy resolution of $`\mathrm{\Delta }E`$ = 100$``$140 keV was obtained, which was sufficient to resolve the low-lying 1<sup>+</sup> states in <sup>176</sup>Lu.
To confirm the validity of the particle identification, we performed an additional experiment where the <sup>176</sup>Yb target was backed with a thin mylar foil (0.086 mg/cm<sup>2</sup>) to reduce the huge <sup>3</sup>He<sup>++</sup> $``$ <sup>3</sup>He<sup>+</sup> yields. In this run, a high-resolution measurement was not possible because of the achromatic beam-transport mode. However, it was confirmed that the unresolved yield for the 1<sup>+</sup> states at 195 keV and 339 keV has the same relative intensity as the summed yield for these two states in the 0 <sup>176</sup>Yb(<sup>3</sup>He,$`t`$) spectrum taken in the dispersion-matching high-resolution mode.
To gauge the B(GT) values obtained in the (<sup>3</sup>He,$`t`$) analysis, a calibration measurement was performed using the neighboring <sup>164</sup>Dy(<sup>3</sup>He,$`t`$) reaction, leading to the 1<sup>+</sup> ground state in <sup>164</sup>Ho. The <sup>164</sup>Ho ground state is known to $`\beta `$-decay to the <sup>164</sup>Dy ground state with a log ft value of 4.6 (Ref. ). This yields a B(GT) value of 0.293$`\pm `$0.006 for the Dy-Ho transition. The <sup>164</sup>Dy result can thus be used to independently calibrate the strong-interaction part of the (<sup>3</sup>He,$`t`$) reaction cross-sections for measurements on the 1<sup>+</sup> states.
Fig. 1 shows the (<sup>3</sup>He,$`t`$) spectrum on <sup>176</sup>Yb at $`\theta `$=0. We deduced the relative transition strengths of the 194.5 keV to the 338.9 keV states, by fitting the peaks with a Gaussian shape (see Fig. 1), to be 1.0:0.55. The transition strengths to the excited states were compared to that of the IAS. The excitation strengths to GT states and the IAS depend, however, on the volume integrals of the central parts of the effective interaction, $`J_\tau `$ and $`J_{\sigma \tau }`$, and distortion effects parameterized by $`N_\tau ^D`$ and $`N_{\sigma \tau }^D`$ . Using the observed strengths to the ground 1<sup>+</sup> state and the IAS in the <sup>164</sup>Dy(<sup>3</sup>He,$`t`$)<sup>164</sup>Ho reaction at the same bombarding energy of 450 MeV, we calibrated the ratio of the interaction strengths including the distortion effects as
$$\frac{N_{\sigma \tau }^DJ_{\sigma \tau }^2}{N_\tau ^DJ_\tau ^2}=\frac{(\frac{d\sigma }{d\mathrm{\Omega }})_{GT}B(F)}{(\frac{d\sigma }{d\mathrm{\Omega }})_{IAS}B(GT)}=7.44\pm 0.79.$$
(1)
Here, we used B(GT)=0.293$`\pm `$0.006, B(F)=N–Z=32, and $`(\frac{d\sigma }{d\mathrm{\Omega }})_{GT}/(\frac{d\sigma }{d\mathrm{\Omega }})_{IAS}`$=0.068$`\pm `$0.007 for the transition to the ground state of <sup>164</sup>Ho. The value of $`7.44\pm 0.79`$ for the (<sup>3</sup>He,$`t`$) reaction at 150 A$``$MeV agrees nicely with the values obtained by the scaling relation $`[(E/A)/55MeV]^2`$ which fits the systematics of (p,n) reaction data , and has been verified specifically for the <sup>176</sup>Yb(p,n) reactions with values of 4.76 obtained for E<sub>p</sub>=120 MeV and 8.46 for 160 MeV .
The present calibration method is applicable since the distortion effect and the ratio of the effective interactions $`J_{\sigma \tau }/J_\tau `$ are expected to be the same as those for the <sup>176</sup>Yb(<sup>3</sup>He,$`t`$) reaction and the kinematic factors are cancelled out. Empirically, the volume integrals of the effective interaction and the distortion effect are not significantly different among target nuclei in the same mass region.
The B(GT) values for the 1<sup>+</sup> states in <sup>176</sup>Lu were obtained using the equation
$$B(GT)=B(F)\frac{1}{7.44}\frac{(\frac{d\sigma }{d\mathrm{\Omega }})_{GT}}{(\frac{d\sigma }{d\mathrm{\Omega }})_{IAS}},$$
(2)
where B(F) is (N–Z), and the ratio of $`(\frac{d\sigma }{d\mathrm{\Omega }})_{GT}/(\frac{d\sigma }{d\mathrm{\Omega }})_{IAS}`$ is obtained from the (<sup>3</sup>He,$`t`$) spectra at 0. The B(GT) values deduced from the 0 cross sections measured in the (<sup>3</sup>He,$`t`$) experiment at 450 MeV are listed in Table I. The B(GT) values obtained for the possible low-lying 1<sup>+</sup> levels in <sup>160</sup>Tb from <sup>160</sup>Gd(<sup>3</sup>He,$`t`$)<sup>160</sup>Tb are also listed for reference. The B(GT) values for <sup>176</sup>Yb obtained here are in agreement with the (p,n) results .
The reliability of the B(GT) values obtained here is supported by the following considerations: 1) scaling of reaction cross sections for GT resonances relative to the Fermi strength observed in the IAS, 2) direct use of the known <sup>164</sup>Dy weak matrix element to calibrate the strong-interaction factor, 3) agreement of the deduced strong interaction factor with the energy scaling systematics of (p,n) reactions in general and those for <sup>176</sup>Yb in particular, and 4) agreement of the B(GT) values obtained in the (<sup>3</sup>He,$`t`$) and (p,n) works for <sup>176</sup>Yb and <sup>160</sup>Gd.
The impact of the B(GT) values for <sup>176</sup>Yb may now be examined. The complete level scheme involved in an Yb-based detector is shown in Fig. 2. These data and those of Table I show that only the two 1<sup>+</sup> levels at 194.5 keV and 338.9 keV in the final nucleus <sup>176</sup>Lu are populated by $`\nu `$ capture below 3 MeV. The thresholds for $`\nu `$ capture are determined by the level energies and the <sup>176</sup>Yb-<sup>176</sup>Lu mass difference as Q($`\nu `$) = 301 keV and 445 keV for the above two states. The Q value of 301 keV lies below 426 keV, the endpoint of the $`pp`$ $`\nu `$ continuum. The B(GT) value for this transition is the larger of the two.
The above data on the $`\nu `$-capture level structure and the capture strengths establish the fundamental suitability of Yb-LENS for the long-standing problem of not only the detection but also the spectroscopy of solar neutrinos. The response of Yb-LENS to solar neutrinos, calculated for a 20 ton natural Yb target with fluxes given by the SSM is shown in Fig. 3a. The observable spectrum is that of the prompt electron $`e^{}`$ with the energy $`E_e^{}=E_\nu `$ \- Q<sub>ν</sub> that follows $`\nu `$-capture. Consequently, the incident $`\nu _e`$ spectrum is recorded directly. The principal features of the $`pp`$, <sup>7</sup>Be, $`pep`$ and the underlying continuum from CNO reactions in the sun are clearly resolved assuming a conservative 1 $`\sigma `$ energy resolution of $`\mathrm{\Delta }E/E7\%/\sqrt{E(MeV)}`$.
Because of the importance of absolute fluxes for the three major $`\nu `$ sources \[$`pp`$(E<sub>max</sub>=420 keV), <sup>7</sup>Be (862 keV) and $`pep`$ (1442 keV)\] B(GT) values must be known individually for the two states in <sup>176</sup>Lu. The simplicity of the level structure of <sup>176</sup>Lu (only 2 states below 3 MeV) is valuable in contrast to the case of <sup>160</sup>Gd (see Table I) in which 5 levels below 1 MeV participate and more than 13 levels below $``$2 MeV. Therefore, the use of Gd is less favorable than Yb since interpretation of the experimental data is more complicated besides presenting severe problems in the practical operation of the detector.
The merits of <sup>176</sup>Yb are reinforced by plans to measure the $`\nu `$-absorption cross sections directly using mono-energetic $`\nu `$’s from megaCurie (MCi) radioactive sources. Such a calibration requires a series of sources, each with a single neutrino line, selected to discriminate level thresholds so that cross sections to individual levels can be deduced at least from the combined data. In principle, as many sources (lines) are needed as target levels revealed in this work. Clearly, such measurements are practically ruled out for Gd with so many levels to characterize. For the two levels in <sup>176</sup>Yb, two sources are proposed. The first, <sup>51</sup>Cr produces a 752 keV line which populates both the Yb transitions (the weaker line at 426 keV produces a much weaker signal). The second source, <sup>75</sup>Se produces a line at 465 keV which is sensitive only to the lower level. Figs. 3b and c depict the LENS response for these two calibration lines. They show that the neutrino response of an Yb-LENS can be calibrated experimentally without interference from a “background” of solar neutrinos.
In summary, we report experimental results on the Gamow-Teller transitions from <sup>176</sup>Yb to <sup>176</sup>Lu using the (<sup>3</sup>He,$`t`$) reaction at $`E`$(<sup>3</sup>He)=450 MeV. A high-resolution measurement enabled us to resolve the low-lying $`1^+`$ states in the residual nucleus. By taking the ratios of the excitation strengths for the IAS and the Gamow-Teller states, we deduce the B(GT) values for low-lying $`1^+`$ states after calibrating the reaction mechanism. The data establish the foundations for a practical real-time <sup>176</sup>Yb detector for sub-MeV solar $`\nu `$’s including $`pp`$ $`\nu `$’s. Solar $`\nu `$ absorption rates for a 20-ton Yb detector are given together with those for the <sup>51</sup>Cr and <sup>75</sup>Se $`\nu `$ calibrations.
The authors acknowledge the RCNP cyclotron staff for their support during the experiment. The present work has been supported by the Ministry of Education, Science, Sports and Culture (Monbusho) with Grant No. 09041108 and by the Japan Society for the Promotion of Science (JSPS).
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# References
INFN-GEF-TH-4/2000
The general QCD parametrization and the large $`N_c`$ description:
Some remarks
G.Morpurgo Università di Genova and Istituto Nazionale di Fisica Nucleare, Genova, Italy.
Abstract. Stimulated by a recent paper of Buchmann and Lebed,a comparison is presented of the two methods mentioned in the title for treating hadron properties in QCD.Doubts arise on the equivalence of the large $`N_c`$ description to real QCD.
(PACS: 12.38.Aw; 11.15.Pg; 13.40.Dk)
1. Introduction.
A recent paper by Buchmann and Lebed (Large $`N_c`$, Constituents quarks and $`N`$, $`\mathrm{\Delta }`$ charge radii) compares, in a specific case, the $`1/N_c`$ description and the general parametrization (GP) method of QCD . Ref. seems to imply that, for the case at hand, the general parametrization should be looked as a very good (though not fully exact) approximation to the $`1/N_c`$ method, which is regarded as more fundamental.
I recall that both the $`1/N_c`$ and the GP methods are parametrizations to describe hadronic properties, but that the GP method, although down to earth, is founded on QCD, while the same is not so clear for $`1/N_c`$. One reason for my doubts on $`1/N_c`$ is that expressed concisely in : “The basis for the large $`N_c`$ approach is the assumption that $`N_c=3`$ QCD is similar to QCD in the limit $`N_c=\mathrm{}`$. In particular it is assumed that there are no phase transitions as we go from $`N_c=3`$ to $`N_c\mathrm{}`$. Currently the status of these assumptions is not clear, because not much is known about QCD($`N_c=\mathrm{}`$)”.
Note, incidentally, that the factor $`g^2/3`$ producing the hierarchy in the $`1/N_c`$ method is, approximately, the same factor that empirically emerges in depressing the diagrams with one added gluon in the GP method; so that the two approaches are characterized, in practice,by a similar hierarchy.
I will exemplify the general parametrization in a few cases, to clarify the situation. But, before doing this, I note two points:
1. The GP method is an exact consequence of QCD, based only on few general properties of the QCD Lagrangian. For many physical quantities of the lowest multiplet of hadrons (e.g. masses, magnetic moments, electromagnetic and semileptonic matrix elements, e.m. form factors etc.) it leads to an exact spin-flavor parametrization, independent of the choice of the renormalization point of the quark masses in the QCD Lagrangian. It turns out that, for a given quantity, the number of terms in this exact QCD parametrization is rather small, indeed smaller than one might have anticipated. The GP method -which, even if not covariant, is fully relativistic- was developed originally to explain the unexpected semiquantitative success of the non relativistic quark model (NRQM) ; it did this \[2a\] long before the $`1/N_c`$ treatment, and much more directly. It emerged that the structure of the terms in the GP is similar to that of the NRQM. Because terms of increasing complexity in the GP have decreasing coefficients, few terms usually suffice to reproduce the data reasonably well, explaining why the NRQM works already in its most naive form.
2. Although $`SU_6`$ was important in suggesting the NRQM , it does not play a role after that. For baryons the essential point in the construction of the NRQM was that the space part of the octet and decuplet wave function has an overall zero orbital angular momentum: $`L=0`$. This implies the factorizability of the baryon (octet or decuplet) NRQM model state as:
$$\varphi _B=X_{L=0}(𝐫_\mathrm{𝟏},𝐫_\mathrm{𝟐},𝐫_\mathrm{𝟑})W_B(s,f)$$
(1)
where $`X`$ is the space part and $`W_B(s,f)`$ is the spin-flavor factor. (Color is understood.) The $`W_B`$’s are symmetric in the three quark variables and have necessarily $`J=1/2`$ and $`J=3/2`$ for the octet and decuplet, so that, automatically, the $`W_B`$ spin-flavor part of $`\varphi _B`$ has the form prescribed by $`SU_6`$, without the need of invoking $`SU_6`$ at all. The factorizability of $`\varphi _B`$ (1) into a space and spin-flavor factor is essential to derive the simple structure of the general parametrization. In the GP there is no need to relate the states to $`SU_6`$ representations as in the $`1/N_c`$ method; nor to rename constituent quarks as “representation quarks”.
Although I will not derive here the GP method -this was done repeatedly \[2a,3a,4\]- I recall some notation, in order to compare GP and $`1/N_c`$ in a few cases. The symbol $`|\varphi _B`$ indicates, in the quark-gluon Fock space, the state corresponding to no gluons and three quarks with wave function $`\varphi _B`$. The exact eigenstate of the QCD hamiltonian $`H_{QCD}`$ for the baryon $`B`$ (with mass $`M_B`$) at rest is written $`|\psi _B`$. It is $`H_{QCD}|\psi _B=M_B|\psi _B`$. A unitary transformation $`V`$ defined in \[2a\], acting on the auxiliary state $`|\varphi _B`$, transforms it into the exact eigenstate $`|\psi _B`$ of $`H_B`$, so that:
$$|\psi _B=|qqq+|qqq\overline{q}q+|qqq,Gluons+\mathrm{}$$
(2)
where the last form of (2) recalls that $`V|\varphi _B`$ is a superposition of all possible quark-antiquark-gluon states with the correct quantum numbers. In particular, configuration mixing is automatically included in $`V|\varphi _B`$. The mass of a baryon is:
$`M_B=\psi _B|H_{QCD}|\psi _B=\varphi _B|V^{}H_{QCD}V|\varphi _B=`$
$`=W_B|\mathrm{`}\mathrm{`}parametrizedmass\mathrm{"}|W_B`$ (3)
The last step (eliminating the space variables) is due to the factorizability of $`\varphi _B`$ (eq.(1)). In the next section I discuss the “parametrized mass” in (S0.Ex1).
2. The parametrization of the baryon masses in the GP method.
The “parametrized mass” in (S0.Ex1) following from the GP method is \[2e,3a\]:
$`\mathrm{`}\mathrm{`}parametrizedmass^{\prime \prime }=M_0+B{\displaystyle \underset{i}{}}P_i^s+C{\displaystyle \underset{i>k}{}}(𝝈_i𝝈_k)+`$
$`+D{\displaystyle \underset{i>k}{}}(𝝈_i𝝈_k)(P_i^s+P_k^s)+E{\displaystyle \underset{\begin{array}{c}ikj\\ (i>k)\end{array}}{}}(𝝈_i𝝈_k)P_j^s+a{\displaystyle \underset{i>k}{}}P_i^sP_k^s+`$ (6)
$`+b{\displaystyle \underset{i>k}{}}(𝝈_i𝝈_k)P_i^sP_k^s+c{\displaystyle \underset{\begin{array}{c}ikj\\ (i>k)\end{array}}{}}(𝝈_i𝝈_k)(P_i^s+P_k^s)P_j^s+dP_1^sP_2^sP_3^s`$ (9)
where the notation is defined in \[2e\]; $`P_i^s`$’s are the projectors on the strange quarks; $`M_0,B,C,\mathrm{},d`$ are parameters. Of the two parameters $`a`$ and $`b`$ only the combination $`(a+b)`$ intervenes.
A comment on (S0.Ex2): Because the different masses of the lowest octet and decuplet baryons are 8 (barring e.m. and isospin corrections), Eq.(S0.Ex2), with 8 parameters $`(M_0,B,C,D,E,a+b,c,d)`$, is certainly true, no matter what is the underlying theory. Yet the general parametrization (S0.Ex2) is not trivial: The values of the above 8 parameters are seen to decrease strongly on moving to terms with increasing number of indices (Eq.(10)). In deriving (S0.Ex2) from QCD, the term $`\mathrm{\Delta }m\overline{\psi }P^s\psi `$ in the QCD Lagrangian is treated exactly; Eq. (S0.Ex2) is correct to all orders in flavor breaking and the derivation takes into account all possible closed loops. In (S0.Ex2) the parameters (in MeV) are -Ref.\[3a\]:
$$\begin{array}{ccccccc}M_0=1076\hfill & ,& B=192\hfill & ,& C=45.6\hfill & ,& D=13.8\pm 0.3\hfill \\ (a+b)=16\pm 1.4\hfill & ,& E=5.1\pm 0.3\hfill & ,& c=1.1\pm 0.7\hfill & ,& d=4\pm 3\hfill \end{array}$$
(10)
The hierarchy of these numbers is evident and, as shown in \[3a\], it corresponds to a reduction factor $`1/3`$ for an additional pair of indices and $`1/3`$ for each flavor breaking factor $`P_i^s`$. The values (10) decrease strongly with increasing complexity of the accompanying spin-flavor structure. Barring $`c`$ and $`d`$, the following mass formula results \[2e\], a generalization of the Gell-Mann Okubo formula that includes octet and decuplet:
$$\frac{1}{2}(p+\mathrm{\Xi }^0)+T=\frac{1}{4}(3\mathrm{\Lambda }+2\mathrm{\Sigma }^+\mathrm{\Sigma }^0)$$
(11)
The symbols stay for the masses and $`T`$ is the following combination of decuplet masses:
$$T=\mathrm{\Xi }^{}\frac{1}{2}(\mathrm{\Omega }+\mathrm{\Sigma }^{})$$
(12)
Because of the level of accuracy reached in comparing Eq.(11) with the data, we wrote (11) so as to be free of electromagnetic effects. (It can be easily checked the combinations in (11) are independent of electromagnetic and isospin effects, to zero order in flavor breaking.) The data satisfy (11) as follows:
$$l.h.s.=1133.1\pm 1.0r.h.s.=1133.3\pm 0.04$$
(13)
an impressive agreement confirming the smallness of the terms neglected in (S0.Ex2).
One more remark \[3a\]: A QCD calculation, if feasible, would express each ($`M_0,B\mathrm{}c,d`$) in (S0.Ex2) in terms of the quantities in the QCD Lagrangian, the running quark masses -normalized at any $`q`$ that we like to select- and the dimensional (mass) parameter $`\mathrm{\Lambda }\mathrm{\Lambda }_{QCD}`$; for instance, setting for simplicity $`m_u=m_d=m`$:
$$M_0\mathrm{\Lambda }\widehat{M}_0(m(q)/\mathrm{\Lambda },m_s(q)/\mathrm{\Lambda })$$
(14)
where $`\widehat{M}_0`$ is some function. Similarly for $`B,C,D,E,a,b,c,d`$ . The numerical value of the coefficients should be seen as the result of a QCD exact calculation performed with an arbitrary choice of the renormalization point of the running quark masses.
3. A comparison with the large $`N_c`$ method.
We now compare the parametrized baryon mass (S0.Ex2), with that obtained in the $`1/N_c`$ method. There (compare ref., Eq.3.4) the parametrization of the baryon masses is also expressed in terms of 8 parameters (from $`c_{(0)}^{1,0}`$ to $`c_{(3)}^{64,0}`$), but, note, the quantities they multiply are collective rather than individual quark variables. It is again true that, setting to zero the smaller coefficients, one finds a relation (Eq.(4.6) in ) between octet and decuplet baryon masses, which is equivalent to Eq. (11).
Neither in Ref. nor in other papers it was stated that this relation coincides -except for the notation- with (11), published long before. I note only, here, the following: The general QCD parametrization can reproduce the good results of $`1/N_c`$ simply using, as we did, the conjecture that the empirical hierarchy apparent in Eqs. (S0.Ex2,10) for the baryon masses, applies to many or all properties, at least for the lowest baryons with a factorizable $`\varphi _B`$. It is, I repeat, what we always did in (see \[2g\], fig.1). We explained in this way a variety of facts about the magnetic moments, $`\mathrm{\Delta }^+p\gamma `$ , semileptonic decays and many other quantities. Note that the GP method in principle includes all diagrams, not only the planar ones; the closed loops, related in the GP to the Trace terms (see in \[3a\], the ref.14) are also taken into account; their contribution is depressed or not depending on the number of additional gluons that are necessary, due to the Furry theorem (see \[3c\], in particular fig.1). For instance the Trace terms that were written in Ref.\[3f\] are depressed by the Furry theorem and can be neglected as we did. The phrasing of Buchmann and Lebed did not clearly express this.
By the way, on reading Ref. it looks as if by neglecting, as we did, from the start, terms proportional to $`m_um_d`$ (which are of the order $`|m_um_d|/(3\mathrm{\Lambda }_{QCD})510^3`$) we had imposed a “mild physical constraint”. This is not so. But, except for these points of language, Buchmann and Lebed seem to state that the relationship between the radii of $`N`$ and $`\mathrm{\Delta }`$ implied by the GP and $`1/N_c`$ methods is the same. One may ask: Is it then really necessary to start from $`N_c=\mathrm{}`$?
Finally I comment on the Coleman-Glashow (CG) relation considered in a ref.\[3h\](see also ). In ref.\[3h\], we recalled the GP result of ref.\[2f\] (only three index flavor breaking terms violate the CG relation) and showed that neither the $`ud`$ mass difference, nor the Trace terms modify this conclusion. This explains the “miracolous” precision of the CG relation, which neglects entirely flavor breaking in its original derivation; such a precision is much better tested after a recent measurement of the $`\mathrm{\Xi }^0`$ mass .
After the appearance of \[3h\](as hep-ph/004198,20 apr 2000) a preprint by Jenkins and Lebed implied by its title that in the large $`N_c`$ description it is quite natural (not ”miracolous”) that the CG relation is so beautifully verified. It is asserted in that the neglected terms are naturally expected to be of an order in $`1/N_c`$ sufficiently high to guarantee their smallness.
This confidence, however, is not supported by their theory. E.g. the Trace terms present in the general parametrization (corresponding to closed loops) are many \[3h\]. Their negligible or vanishing global contribution cannot be established using only the order in $`1/N_c`$ of a typical term. This is another reason,in addition to the doubts raised in , confirming that it is not established that the $`1/N_c`$ expansion can make predictions having a real QCD foundation.
Acknowledgement. I am very indebted to G.Dillon for frequent discussions.
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# REPORT OF THE QCD WORKING GROUP
## 1 INTRODUCTION
### 1.1 Objectives of the working group
Fully hadronic multi-jet topologies play an important role at Lep2, in the contexts both of physics measurements and of searches for new phenomena. For example four and more hadronic jet topologies dominate the statistics both in the measurements of $`W`$ boson pairs and in the searches for Higgs bosons, because of the large hadronic decay branching ratios of all heavy bosons involved. Improving our understanding of the physics of QCD processes and of the modelling provided by our main generators is relevant at Lep2 for two main reasons:
* In contrast to the other two main decay topologies occurring in boson pair production and studied at Lep2 (the semi- or fully leptonic ones), four-quark production processes leading to fully hadronic topologies must be analysed in the presence of large backgrounds from two-quark production, which can lead to similar multi-jet topologies via hard QCD processes.
* The reconstruction of basic event observables such as for instance boson masses is intrinsically more difficult in fully hadronic channels because of soft QCD processes, which broaden the jets, create ambiguities in assigning the jets, and can also result in cross-talk between the produced bosons (if they are short-lived) which may be large enough to be noticeable in precision measurements such as that of the $`W`$ mass.
This working group on QCD generators has focussed its activity on the first of the two items above, dealing mainly with hard QCD processes. The second item (physics and modeling of soft QCD), has been and still is pursued in the framework of the WWMM-2000 (previously called Crete) workshop .
The work described here was originally motivated by the desire to assess the performance of the various QCD generators used to model QCD backgrounds at Lep2, as well as the expected corresponding theoretical uncertainties. The point of view taken was that final publications at Lep2 should be based on the best possible Monte Carlo programs, and that we should be able to specify corrections when needed, and to quote uncertainties, in a reliable way, particularly when fully satisfactory treatments are not yet available.
In addition to serving the Lep2 community, the improvements of the programs and of the basic understanding also benefits a number of other genuine QCD studies.
In the following section the programs available and investigated by the working group are described by their authors. In the case of standard programs commonly used in the community, only those aspects relevant to the topics studied, and the related improvements stimulated by the working group, are covered. Also several new approaches and options are described.
Then follow five sections where the investigations of the main physics features considered are reported :
* Inclusive (all flavour) jet rates are not extremely well modelled and can result in significant discrepancies, even at Lep2, when four-jet events are selected. The different Monte-Carlo approaches available, and the tuning strategies adopted by the different collaborations, are compared, and a procedure to extrapolate the uncertainty to Lep2 energies, based on the quality of the description achieved at Lep1, is outlined.
* Mass effects in 3- and 4-jet rates were not previously considered in detail by the modellers, but are relevant to analyses in which $`b`$-tagging is used as a tool, such as the Higgs searches at Lep2. In addition several features of the modelling result in uncertainties in basic QCD measurements at Lep1, such as that of the $`b`$-quark mass. A consistent method to quantify the theoretical uncertainty is presented, and the performance of the different Monte-Carlo programs available, including recent improvements, described. Additional uncertainties from gluon splitting processes into $`b\overline{b}`$ (see below as well) in the case of the 4-jet rate are also considered.
* Genuine four-jet observables, particularly angular distributions, are not well described by Monte-Carlo programs based on parton shower approaches matched to matrix elements at the level of three partons. This can result in biases when methods based on topological information are used to select (or anti-select) the events. An additional basic motivation for improving the description in this respect lies in the use of four-jet events to measure the strong interaction coupling constant $`\alpha _\mathrm{S}`$. The emphasis of the work was to estimate uncertainties, and to evaluate new Monte-Carlo programs in which matching of the parton shower approach with matrix elements is attempted beyond three partons.
* The $`b`$-quark fragmentation function is relevant to a number of topics involving $`b`$ quarks, at both Lep1 and Lep2 energies, as it affects for instance the lifetime of $`B`$-hadrons and selection efficiencies of $`b`$-tagging algorithms. Although this topic was not a central one in this working group, it was felt important to report as much as possible the present status and recent results on this topic.
* Processes involving gluon splitting into $`b\overline{b}`$ are poorly known, both theoretically and experimentally, and become more important at Lep2 energies. Several new options exist in the different Monte-Carlo programs, which enable one to alter the rate and kinematics of the production. These are considered in the light both of analytical results and of measurements at Lep1.
The evaluations were based on comparisons of the different Monte-Carlo programs, with analytical results when available, and with data at Lep1. An effort was made to define dedicated observables enabling meaningful comparisons, and to estimate the theoretical uncertainties quantitatively. In several cases the calculations, the Monte-Carlo simulations and the evaluations of systematic uncertainties were extrapolated to Lep2 energies as well. In some cases discrepancies were found between the theoretical expectations, the data, and Monte-Carlo results. An attempt to quantify such discrepancies was then made, and the results served to stimulate improvements by the model builders. Several such improvements were actually achieved in the course of the workshop, and evaluations of the resulting new Monte-Carlo versions was carried out as well.
In the final section, overall conclusions are presented. Although in some instances real progress was achieved thanks to this working group, clearly in many cases still more work and checks are needed. Such additional investigations and developments are mentioned, based on the present knowledge. General recommendations on the use of the present programs are formulated in each of the relevant contexts.
### 1.2 Jet clustering algorithms
The jet clustering algorithms used in this report are those in most common use in $`e^+e^{}`$ experiments: the Jade , Durham and Cambridge algorithms. They are used to define the jets at parton level in the theoretical calculations, and for grouping the selected charged and neutral particles into jets at the experimental level.
The Jade algorithm was the earliest of these and established the method of successive binary clustering that has been adopted in later algorithms. For all pairs of final-state particles $`(i,j)`$, a test variable $`y_{ij}`$ is defined as indicated in Table 1. The minimum of all $`y_{ij}`$ is compared with the so-called jet resolution parameter, $`y_c`$ (often called $`y_{cut}`$). If it is smaller, the two particles are recombined into a new pseudo-particle with four-momentum $`p_k=p_i+p_j`$.<sup>1</sup><sup>1</sup>1Other possible recombination schemes are discussed in The algorithm can be applied again to the new group of pseudo-particles until all pairs satisfy $`y_{ij}>y_c`$. The number of jets in the event is then the number of pseudo-particles one has at the end. In perturbative theoretical calculations, this procedure leads to infrared-finite quantities because one excludes the regions of phase-space that cause trouble. For the same reason, sensitivity to non-perturbative physics is limited and hadronization corrections can be estimated from Monte-Carlo models.
The Jade algorithm was nevertheless found to have some unpleasant theoretical and experimental features, which arise from the fact that its resolution criterion is approximately one of invariant mass, $`M_{ij}^22E_iE_j(1\mathrm{cos}\theta _{ij})>y_cE_{vis}^2`$. This means that particles at widely different angles can be combined into the same jet, leading to theoretical predictions with large higher-order corrections that cannot be resummed, and to the possibility of “ghost jets” (jets in directions where no particles are observed) at the experimental level.
The problems of the Jade algorithm are largely alleviated by replacing the test variable by one that measures the relative transverse momentum of pairs of particles rather than their invariant mass. This led to the formulation of the Durham algorithm, the most widely used for Lep physics, in which $`\mathrm{min}(E_i^2,E_j^2)`$ simply replaces $`E_iE_j`$ in the Jade formula (see Table 1). The resolution criterion then becomes $`k_{Ti}^2>y_cE_{vis}^2`$ at small angles, where $`k_{Ti}`$ is the transverse momentum of a particle/jet relative to the direction of any other in the event.
The Cambridge algorithm has been introduced to cure some remaining defects of the Durham algorithm at low values of the jet resolution $`y_c`$, with a better understanding of the processes involving soft gluon radiation, allowing one to explore regions of smaller $`y_c`$, where furthermore the experimental error of three-jet ratios is expected to be smaller. It uses the same recombination procedure and test variable as Durham but with the new ingredients of angular ordering and soft freezing.
The selection of the first pair of particles to be compared with the resolution parameter is now made according to the ordering variable $`v_{ij}=2(1\mathrm{cos}\theta _{ij})`$ (see Table 1). Then, for the pair of particles with the smallest $`v_{ij}`$, one computes $`y_{ij}`$ and if $`y_{ij}<y_c`$ the two particles are recombined. If not, the soft freezing mechanism comes into the game by considering the softer particle as a resolved jet and by bringing back the other one into the binary procedure. The net effect of the new definition is that NLO corrections to the three-jet fraction become smaller .
In the Durham algorithm one can always define a transition value of $`y_c`$, $`y^{nn+1}`$, in which an $`(n+1)`$-jet configuration event becomes one with $`n`$ (or fewer) jets. Furthermore, the number of jets is monotonically decreasing for increasing $`y_c`$. However, in Cambridge, this property is lost due to the fact that the sequence of clustering depends on the external $`y_c`$ and in some circumstances certain jet topologies are not present for a specific event. In the case of three jets this affects $`1\%`$ of the events in the range $`y_c0.01`$.
For a more thorough discussion of these and other $`e^+e^{}`$ jet algorithms in current use, see .
#### 1.2.1 Jet rates
Having chosen a jet algorithm one may define the n-jet rate, $`R_n`$, by the fraction of hadronic final states that are clustered into precisely $`n`$ jets at jet resolution $`y_c`$:
$$R_n(y_c)=\frac{\sigma _n(y_c)}{\sigma _{had}}$$
(1)
where $`\sigma _n`$ and $`\sigma _{had}`$ are the $`n`$-jet and the total hadronic cross sections, respectively. Here we assume that all processes other than the direct QCD one, $`e^+e^{}Z^0/\gamma ^{}q\overline{q}`$ hadrons, have been eliminated by suitable cuts. For some purposes it will be useful to define jet rates for a particular primary quark flavour:
$$R_n^q(y_c)=\frac{\sigma _{q\overline{q}njets}(y_c)}{\sigma _{q\overline{q}had}}$$
(2)
where $`q=\mathrm{}`$, $`c`$ or $`b`$, with $`\mathrm{}`$ representing a light ($`u,d,s`$) quark.
## 2 MONTE CARLO GENERATORS
This Section gives brief descriptions of the main QCD event generators for two-fermion processes at Lep2, with emphasis on the features relevant to multi-jet and $`b`$-jet fragmentation.
### 2.1 PYTHIA
Pythia is a general-purpose generator . The current version, Pythia 6.1, combines and extends the previous generation of programs, Pythia 5.7, Jetset 7.4 and Spythia . Here we concentrate on those aspects of the program that have been modified as a consequence of the current workshop, or are of specific interest to this working group. Program code, manuals and sample main programs are obtainable from http://www.thep.lu.se/$``$torbjorn/Pythia.html.
#### 2.1.1 Gluon radiation off heavy quarks
The Pythia final-state shower consists of an evolution in the squared mass $`m^2`$ of a parton. That is, emissions are ordered in decreasing mass of the radiating parton, and the Sudakov form factor is defined as the no-emission rate in the relevant mass range. Such a choice is not as sophisticated as the angular one in Herwig or the transverse momentum one in Ariadne, but usually the three tend to give similar results. (An exception, where small but significant differences were found, is the emission of photons in the shower .) One of the advantages is that a mapping between the parton-shower and matrix-element variables is rather straightforward to $`𝒪(\alpha _\mathrm{S})`$ for massless quarks, and that already the basic shower populates the full phase space region very closely the same way as the matrix element. It is therefore possible to introduce a simple correction to the shower to bring the two into agreement.
The other main variable in the shower is $`z`$, as used in the splitting kernels. It is defined as the energy fraction in the CM frame of the event. That is, in a branching $`ab+c`$, $`E_b=zE_a`$ and $`E_c=(1z)E_a`$. In the original choice of $`z`$, which is done at the same time as $`m_a`$ is selected, the $`b`$ and $`c`$ masses are not yet known. A cut-off scale $`Q_01`$ GeV is used to constrain the allowed phase space, by assigning fictitious $`b`$ and $`c`$ masses $`Q_0/2`$ so that $`a`$ can only branch if $`m_a>Q_0`$, but kinematics is constructed as if $`b`$ and $`c`$ were massless. At a later stage, when $`m_b`$ and $`m_c`$ are being selected, possibly well above $`Q_0`$, the previously found $`z`$ may be incompatible with these. The solution is to take into account mass effects by reducing the magnitude of the three-momenta $`𝐩_b=𝐩_c`$ in the rest frame of $`a`$. Expressed in four-momenta in an arbitrary frame, this is equivalent to
$`p_b`$ $`=`$ $`(1k_b)p_b^{(0)}+k_cp_c^{(0)},`$
$`p_c`$ $`=`$ $`(1k_c)p_c^{(0)}+k_bp_b^{(0)},`$ (3)
where $`p_b^{(0)}`$ and $`p_c^{(0)}`$ are the original massless momenta and $`p_b`$ and $`p_c`$ the modified massive ones. The parameters $`k_b`$ and $`k_c`$ are found from the constraints $`p_b^2=m_b^2`$ and $`p_c^2=m_c^2`$.
Angular ordering is not automatic, but is implemented by vetoing emissions that don’t correspond to decreasing opening angles. The opening angle of a branching $`ab+c`$ is calculated approximately as
$$\theta \frac{p_b}{E_b}+\frac{p_c}{E_c}\sqrt{z(1z)}m_a\left(\frac{1}{zE_a}+\frac{1}{(1z)E_a}\right)=\frac{1}{\sqrt{z(1z)}}\frac{m_a}{E_a}.$$
(4)
The procedure thus is the following. In the $`\gamma ^{}/Z^0`$ decay, the two original partons 1 and 2 are produced, back-to-back in the rest frame of the pair. In a first step, they are evolved downwards from a maximal mass equal to the CM energy, with the restriction that the two masses together should be below this CM energy. When the two branchings are found, they define $`m_1`$ and $`m_2`$ and the $`z`$ values of $`13+4`$ and $`25+6`$. These latter branchings obviously have smaller opening angles than the $`180^{}`$ one between 1 and 2, so no angular-ordering constraints appear here. The matching procedure to the matrix element is used to correct the branchings, however, as will be described below. In subsequent steps, a pair of partons like 3 and 4 are evolved in parallel, from maximum masses given by the smaller of the mother (1) mass and the respective daughter (3 or 4) energy. Here angular ordering restricts the region of allowed $`z`$ values in their branchings, but there are no matrix-element corrections. Once $`m_3`$ and $`m_4`$ are fixed, the kinematics of the $`13+4`$ branching needs to be modified according to eq. (3).
Let us now compare the parton-shower (PS) population of three-jet phase space with the matrix-element (ME) one. With the conventional numbering $`q(1)\overline{q}(2)g(3)`$, and $`x_j=2E_j/E_{CM}`$, the matrix element is of the form
$$\frac{1}{\sigma _0}\frac{\mathrm{d}\sigma _{\mathrm{ME}}}{\mathrm{d}x_1\mathrm{d}x_2}=\frac{\alpha _\mathrm{S}}{2\pi }\frac{4}{3}\frac{M(x_1,x_2,r_q)}{(1x_1)(1x_2)}.$$
(5)
For massless quarks
$$M(x_1,x_2,0)=x_1^2+x_2^2,$$
(6)
while for massive ones
$$M\left(x_1,x_2,r_q=\frac{m_q^2}{E_{CM}^2}\right)=x_1^2+x_2^24r_qx_38r_q^2(2r_q+4r_q^2)\left(\frac{1x_2}{1x_1}+\frac{1x_1}{1x_2}\right).$$
(7)
There are two shower histories that could give a three-jet event. One is $`\gamma ^{}/Z^0(0)q(i)\overline{q}(2)q(1)\overline{q}(2)g(3)`$, i.e. with an intermediate ($`i`$) quark branching $`q(i)q(1)g(3)`$. For massless quarks this gives
$`Q^2`$ $`=`$ $`m_i^2=(p_0p_2)^2=(1x_2)E_{CM}^2,`$ (8)
$`z`$ $`=`$ $`{\displaystyle \frac{p_0p_1}{p_0p_i}}={\displaystyle \frac{E_1}{E_i}}={\displaystyle \frac{x_1}{x_1+x_3}}={\displaystyle \frac{x_1}{2x_2}},`$ (9)
$``$ $`{\displaystyle \frac{\mathrm{d}Q^2}{Q^2}}\mathrm{d}z={\displaystyle \frac{\mathrm{d}x_2}{1x_2}}{\displaystyle \frac{\mathrm{d}x_1}{2x_2}}.`$ (10)
The parton-shower probability for such a branching is
$$\frac{\alpha _\mathrm{S}}{2\pi }\frac{4}{3}\frac{1+z^2}{1z}\mathrm{d}z\frac{\mathrm{d}Q^2}{Q^2}=\frac{\alpha _\mathrm{S}}{2\pi }\frac{4}{3}\frac{1x_1}{x_3}\left[1+\left(\frac{x_1}{2x_2}\right)^2\right]\frac{\mathrm{d}x_1\mathrm{d}x_2}{(1x_1)(1x_2)}.$$
(11)
There also is a second history, where the rôles of $`q`$ and $`\overline{q}`$ are interchanged, i.e. $`x_1x_2`$. (On the Feynman diagram level, this is the same set as for the matrix element, except that the shower does not include any interference between the two diagrams.) Adding the two, one arrives at a form
$$\frac{1}{\sigma _0}\frac{\mathrm{d}\sigma _{\mathrm{PS}}}{\mathrm{d}x_1\mathrm{d}x_2}=\frac{\alpha _\mathrm{S}}{2\pi }\frac{4}{3}\frac{S(x_1,x_2,r_q)}{(1x_1)(1x_2)},$$
(12)
with
$$S(x_1,x_2,0)=1+\frac{1x_1}{x_3}\left(\frac{x_1}{2x_2}\right)^2+\frac{1x_2}{x_3}\left(\frac{x_2}{2x_1}\right)^2.$$
(13)
In spite of the apparent complexity of $`S(x_1,x_2,0)`$ relative to $`M(x_1,x_2,0)`$, it turns out that $`S(x_1,x_2,0)M(x_1,x_2,0)`$ everywhere but also that $`S(x_1,x_2,0)>M(x_1,x_2,0)`$. It is therefore straightforward and efficient to use the ratio
$$\frac{\mathrm{d}\sigma _{\mathrm{ME}}}{\mathrm{d}\sigma _{\mathrm{PS}}}=\frac{M(x_1,x_2,0)}{S(x_1,x_2,0)}$$
(14)
as an acceptance factor inside the shower evolution, in order to correct the first emission of the quark and antiquark to give a sum in agreement with the matrix element.
Clearly, the shower will contain further branchings that modify the simple result, e.g. by the emission both from the $`q`$ and the $`\overline{q}`$, but these effects are formally of $`𝒪(\alpha _\mathrm{S}^2)`$ and thus beyond the accuracy we strive to match. One should also note that the shower modifies the distribution in three-jet phase space by the appearance of Sudakov form factors, and by using a running $`\alpha _\mathrm{S}(p_{}^2)`$ rather than a fixed one. In both these respects, however, the shower should be an improvement over the fixed-order result.
The prescription of correcting the first branchings by a factor $`M(x_1,x_2,0)/S(x_1,x_2,0)`$ was the original one, used up until Jetset 7.3. In 7.4 an intermediate “improvement” was introduced, in that masses were used in the matrix-element numerator, i.e. an acceptance factor $`M(x_1,x_2,r_q)/S(x_1,x_2,0)`$. (The older behaviour remained as an option.) The experimental problems found with this procedure has prompted new studies as part of this workshop. Starting with Pythia 6.130, therefore also masses have been introduced in the shower expression, i.e. an acceptance factor $`M(x_1,x_2,r_q)/S(x_1,x_2,r_q)`$ is now used.
In the derivation $`S(x_1,x_2,r_q)`$, one can start from the ansatz
$`x_2`$ $`=`$ $`1{\displaystyle \frac{m_i^2m_q^2}{E_{CM}^2}},`$
$`x_1`$ $`=`$ $`\left(1+{\displaystyle \frac{m_i^2m_q^2}{E_{CM}^2}}\right)\left((1k_1)z+k_3(1z)\right),`$ (15)
$`x_3`$ $`=`$ $`\left(1+{\displaystyle \frac{m_i^2m_q^2}{E_{CM}^2}}\right)\left((1k_3)(1z)+k_1z\right).`$
The quark mass enters both in the energy splitting between the intermediate quark $`i`$ and the antiquark 2, and in the correction procedure of eq. (3) for the sharing of energy in the branching $`q(i)q(1)g(3)`$. The constraints $`p_1^2=m_q^2`$ and $`p_3^2=0`$ give $`k_1=0`$ and $`k_3=m_q^2/m_i^2`$. One then obtains
$`Q^2`$ $`=`$ $`m_i^2=(1x_2+r_q)E_{CM}^2,`$
$`z`$ $`=`$ $`{\displaystyle \frac{1}{2x_2}}\left(x_1r_q{\displaystyle \frac{2x_1x_2}{1x_2}}\right).`$ (16)
By a fortuitous cancellation of mass terms, $`\mathrm{d}Q^2/Q^2\mathrm{d}z`$ is the same as in eq. (10), but the $`(1+z^2)/(1z)`$ factor is no longer simple. Therefore one obtains
$$S(x_1,x_2,r_q)=\frac{1x_1}{x_3}\frac{1x_2}{1x_2+r_q}[1+\frac{1}{(2x_2)^2}(x_1r_q\frac{x_3}{1x_2})^2]+\{x_1x_2\},$$
(17)
where the second term comes from the graph where the antiquark radiates.
The mass effects go in the “right” direction, $`S(x_1,x_2,r_q)<S(x_1,x_2,0)`$, but actually so much so that $`S(x_1,x_2,r_q)<M(x_1,x_2,r_q)`$ in major regions of phase space. This is illustrated in Figure 1. The dashed curve here shows how well the PS and ME expressions agree in the massless case. The dash-dotted one is the well-known “dead cone effect” in the matrix element , and the full the corresponding suppression in the shower. Very crudely, one could say that the massive shower exaggerates the angle of the dead cone by about a factor of two (in this rather typical example).
Thus the amount of gluon emission off massive quarks is underestimated already in the original prescription, where masses entered in the kinematics but not in the ME/PS correction factor. If instead the ratio $`M(x_1,x_2,r_q)/S(x_1,x_2,0)`$ is applied, the net result is a distribution even more off from the correct one, by a factor $`S(x_1,x_2,r_q)/S(x_1,x_2,0)`$. Thus it would have been better not to introduce the mass correction in Jetset 7.4.
Armed with our new knowledge, we can now instead use the correct factor, namely the ratio $`M(x_1,x_2,r_q)/S(x_1,x_2,r_q)`$. A technical problem is that this ratio can exceed unity, in the example of Figure 1 by up to almost a factor of two. This could be solved e.g. by enhancing the raw rate of emissions by this factor. However, another trick was applied, based on the fact that the accessible $`z`$ range is smaller for a massive quark than a massless one. Therefore, without any loss of phase space, $`z`$ can be rescaled to a $`z^{}`$ according to
$$(1z^{})=(1z)^k,\mathrm{with}k=\frac{\mathrm{ln}(m_q^2/E_{CM}^2)}{\mathrm{ln}(Q_0^2/E_{CM}^2)}<1.$$
(18)
The ME/PS correction factor then has to be compensated by $`k`$, and thereby comes below unity almost everywhere — the remaining weighting errors are too small to be relevant.
In Sec. 4.4 of this report it is shown that the corrected procedure now does a good job of describing mass effects in the amount of three-jet events. Problems still remain in the four-jet sector, however, where the emission off heavy quarks is reduced more in Pythia than in the data. These four-jets come in several categories in the Monte Carlo simulation. If one resolved gluon is emitted from the quark and another from the antiquark, or if a gluon branches into two resolved partons, the mass effects should now be included. If the quark emits both resolved gluons, however, the second emission involves no correction procedure. Instead the dead cone effect is exaggerated, similarly to what was shown in Figure 1. That might then explain the discrepancies noted above.
The intention is to find an alternative algorithm that better can take into account mass effects at all steps of the shower. For instance, if the evolution is performed in terms of the variable $`Q^2=m^2m_q^2`$ rather than $`Q^2=m^2`$, then the dead-cone effect is underestimated rather than overestimated. A suppression factor could therefore be implemented to correct down to the desired level. The technical details have yet to be worked out.
#### 2.1.2 The total four-jet rate
The above modifications partly address the four-jet rate off heavy quarks relative to light quarks, but not the shortfall in the overall four-jet rate in Pythia relative to the data. Currently the matrix-element correction procedure is used in the first branching of both sides of the event, i.e. both the quark and the antiquark ones. Thus not only the three-jet but also the four-jet rate is affected. If the correction procedure is only used on the side with the harder emission, here defined as the one occuring at the largest mass, one might hope to increase the four-jet rate relative to the three-jet one. This possibility was studied, for simplicity only for massless quarks. The result was disappointing, however. To the extent that the four-jet rate is at all changed, it is below the 1% level. In retrospect, this is maybe not so surprising, considering how close the matrix-element correction factor is to unity, cf. Figure 1. A solution to the four-jet rate problem therefore remains to be found.
#### 2.1.3 Gluon splitting to heavy quarks
A few new options have been included in Pythia, that allow studies of the gluon splitting rate under varying assumptions. These developments are described in Sec. 7.3.
#### 2.1.4 Fragmentation of low-mass strings
The Lund string fragmentation algorithm has remained essentially unchanged over the years, and generally does a good job of describing data. Some improvements have recently been made (in Pythia 6.135 onwards) in the description of low-mass strings , however.
Whereas gluon emission only adds kinks on the string stretched between a quark end and an antiquark one, a gluon splitting $`gq\overline{q}`$ splits an existing string into two. In this process, one of the new strings can obtain a small invariant mass, so that it can only produce one or two primary hadrons. Such a low-mass system is called a cluster, and is handled separately from ordinary strings. If only one hadron is produced, “cluster collapse”, its flavour is completely specified by the string endpoints.
In fixed-target $`\pi p`$ collisions, strings are often stretched between a produced central charm quark and a beam remnant antiquark or diquark. Thus the cluster collapse mechanism favours the production of charm hadrons that share a valence flavour content with the incoming beam particles. This was predicted in Pythia, but the measurements have shown that production asymmetries are smaller in data than in the model. The new data have therefore been used to tune some aspects of the cluster treatment, and some other improvements were included at the same time. The ones relevant for $`e^+e^{}`$ physics are summarized below.
The quark masses assigned to “on-shell” quarks, e.g. in the event listing, have been changed to $`m_u=m_d=0.33`$ GeV, $`m_s=0.5`$ GeV, $`m_c=1.5`$ GeV and $`m_b=4.8`$ GeV. In previous program versions, lower “current-algebra” masses were used to comply with requirements e.g. for Higgs physics, but these latter needs are now covered by the new running-mass function PYMRUN. The change in masses has consequences in several places, e.g. for the rate of $`gq\overline{q}`$ branchings. In this Section, the main point is the change in the string mass spectrum, and thereby in the fate of strings. For a string $`q_1\overline{q}_2`$, the cluster treatment is applied whenever $`m(q_1\overline{q}_2)<m(q_1)+m(\overline{q}_2)+1`$ GeV, while the normal string routine is used above that.
A cluster can produce either one or two primary hadrons. The choice is made dynamically, as follows. The cluster is assumed to break into two hadrons $`h_1=q_1\overline{q}_3`$ and $`h_2=q_3\overline{q}_2`$ by the production of a new $`q_3\overline{q}_3`$ pair. The composition of the new flavour and the spin multiplet assignment of the hadrons is determined by standard string fragmentation parameters. If $`m(h_1)+m(h_2)<m(q_1\overline{q}_2)`$, an allowed two-body decay of the cluster has been found. Even in case of failure, a subsequent new try might succeed, with another $`q_3`$ or another spin assignment. Therefore a very large number of tries would make each cluster decay to two hadrons if at all possible, while only one try gives a more gradual transition between one and two hadrons as the various two-body thresholds are passed. As a compromize between the extremes, up to two tries are made. If neither succeeds, the cluster collapses to one hadron
In a cluster collapse, it is not possible to conserve energy and momentum within the cluster. Instead other parts of the events have to receive or donate energy to put the hadron on mass shell. The algorithm handling this has now been made more physically appealing, by performing the shuffling to/from the parts of the event that are most closely moving in the same general direction as the collapsing cluster. The technical details are not described here, but one may note that differences are small relative to the previous simpler algorithm (still available as an option and as a last resort, should the more sophisticated one fail to find a sensible solution).
The treatment of a two-body cluster decay has been improved to provide a smoother match to the string description in the overlapping mass region. At a first step, the cluster decay is isotropic. The decay is accepted with a weight $`\mathrm{exp}(p_{}^2/2\sigma ^2)`$, where the $`p_{}`$ is defined relative to the $`q_1\overline{q}_2`$ axis in the cluster rest frame. This agrees with the standard Gaussian string fragmentation $`p_{}`$ spectrum well above threshold, but reverts to isotropic decay near the threshold. Even with $`p_{}`$ fixed, two “mirror” solutions exist for the longitudinal momenta of the hadrons. The relative probabilities are well-defined in the string model, and are here used to make the choice. Near threshold both are equally likely, while further above threshold the $`q_1\overline{q}_3`$ hadron is preferentially moving in the $`q_1`$ direction and vice versa.
#### 2.1.5 A shower interface to four-jet events (massless ME)
A few years ago, an algorithm was developed to allow the Pythia shower to start from a given four-jet configuration, $`q\overline{q}gg`$ or $`q\overline{q}q^{}\overline{q}^{}`$ . This was intended to allow comparisons e.g. of four-jet topologies between matrix-element calculations and data, with a realistic account of showering and hadronization effects not covered by the matrix-element calculations. The standard Pythia shower does not do this well, since it does not include any matching procedure to four-jet matrix elements and therefore does not do e.g. the azimuthal angles in branchings fully correctly.
A problem is that the standard shower routine is really set up only to handle systems of two showering partons, not three or more. (Actually an option does exist for three, but it is primitive and hardly used by anybody.) The trick therefore is to try to guess the “prehistory” of shower branchings that gave the specified four-parton configuration, and thereafter to run a normal shower starting from two partons. Here two of the subsequent branchings already have their kinematics defined, while the rest are chosen freely as in a normal shower. Benefits of having a prehistory include (i) the availability of the standard machinery to take into account recoils when masses are assigned to partons massless in the matrix elements, (ii) a knowledge of angular-ordering constraints on subsequent emissions and azimuthal anisotropies in them, and (iii) information on the colour flow as required for the subsequent string description.
The choice among possible shower histories is based on a weight obtained from the mass poles and splitting kernels. As an example, consider a $`q(1)\overline{q}(2)g(3)g(4)`$ configuration, which could come e.g. from an initial $`q(i)\overline{q}(2)`$ configuration followed by branchings $`q(i)q(1)g(j)`$ and $`g(j)g(3)g(4)`$. The relative weight is then
$$𝒫=𝒫_{i1j}𝒫_{j34}=\frac{1}{m_i^2}\frac{4}{3}\frac{1+z_{i1j}^2}{1z_{i1j}}\frac{1}{m_j^2}\mathrm{\hspace{0.17em}3}\frac{(1z_{j34}(1z_{j34}))^2}{z_{j34}(1z_{j34})}.$$
(19)
Of course, one could imagine including further information, e.g. on azimuthal angles or on a scale-dependent $`\alpha _\text{S}`$.
The original routines were not set up to handle massive quarks, e.g. to correct the $`z`$ definition for the rescaling of eq. (3). This has now been included, and also the interface has been simplified. The re-implementation originally contained a bug, that was fixed in Pythia 6.137.
Users can now CALL PY4JET(PMAX,IRAD,ICOM) to shower and fragment a four-parton configuration. If ICOM is 0 or 1 the configuration is picked up either from the HEPEVT or the PYJETS commonblock. The partons have to be stored in the order $`q\overline{q}gg`$ or $`q\overline{q}q^{}\overline{q}^{}`$, where $`q^{}\overline{q}^{}`$ is assumed to be the secondary quark pair. (Interference terms make the primary/secondary pair distinction nontrivial in a matrix element, but pragmatic recipes should work well.) Initial-state photons can be interspersed anywhere in the given initial state, and final-state photon radiation in the shower is off or on for IRAD 0 or 1. PMAX sets the maximum mass scale allowed in the shower. In an exclusive description, i.e. where one wants four-jet only and not five or more jets, the logical choice would be to put PMAX equal to the mass cutoff applied to the matrix elements. An inclusive picture, where all emissions are allowed below the lowest mass scale of the reconstructed shower, is obtained for PMAX$`=0`$ (or, more precisely, PMAX$`<Q_0`$).
#### 2.1.6 Interfacing 4 parton LO massive ME: Fourjphact.
As already explained in the preceding Sections, complete matrix elements calculations are expected to give a good description of multijet events when large separations among jets are involved and in particular when angular variables are considered. On the other hand, pure ME differential cross sections lack PS and hadronization and cannot reproduce collinear and soft radiation. It is therefore important to have the possibility to start with pure ME calculations and complement them with these additional features. The results obtained in this way (ME + PS + hadronization) can be compared with pure parton level ones as well as with those from dedicated QCD MC’s.
If one takes for example topologies with four or more jets, one expects that a reasonable description for not too small values of the jet resolution $`y_{cut}`$ may be obtained starting with four jet ME at a much lower $`y_{cut}`$ and adding to it PS and hadronization. One must however be aware of the fact that when starting with four parton ME, all events described by two or three parton ME + PS + hadronization are not taken into account. In this respect QCD MC’s, like Herwig or Pythia, surely give a more complete description, as they start PS from two parton ME and match 3 parton production with the respective ME results. The above mentioned approach of starting from 4 parton ME can however be considered as a complementary approach for some studies and a way to check MC results when for instance angular variables or mass effects are involved.
Fourjphact is a Fortran code which has been written to provide a tool for this kind of studies and comparisons. It computes exact LO massive ME for all $`e^+e^{}q\overline{q}q^{}\overline{q}^{}`$ and $`e^+e^{}q\overline{q}gg`$ final states and it interfaces them with the Pythia routine PY4JET described in the preceding Section. It can therefore be used to compute total or differential four jet cross sections at parton level or to study fully hadronic events initiated by 4 partons.
The program, together with instructions and examples, can be found in
http://www.to.infn.it/$``$ballestr/qcd/ .
Here we limit ourselves to a brief description of the main features of the program.
Fourjphact computes all $`ee4q`$ ME’s with the method of ref. while for $`ee2q2g`$ it makes use of the routine of ref . Numerical integration over phase space is performed with VEGAS . Unweighted event generation and distributions at parton level are implemented as in the four fermion program WPHACT . Initial State Radiation is included, when requested, via the structure function approach .
When using the program, one starts by computing some cross section. Unweighted events may be generated during this step, or in a second run in order to obtain a predetermined number of events. These may be passed to Pythia which provides PS and hadronization.
In the cross section computation one may choose between fixed or running $`\alpha _\text{S}(M)`$. In the second case, the scale $`M`$ for $`4q`$ diagrams is chosen to be the invariant mass of the gluon propagator, while for $`2q2g`$ the invariant mass of the two gluons is used.
An inventory of cuts at parton level are already defined in Fourjphact: to implement them one has only to specify the numerical values for minima and maxima of energies, transverse momenta, angles among partons and invariant masses. Jade or Durham or Cambridge $`y_{cut}`$ at parton level can be requested in a similar way. Any other cut can be easily defined in an include file. It must be noticed in this connection that massive LO ME for $`2q2g`$ cannot be computed without any cut or $`y_{cut}`$. $`4q`$ final states can in principle be computed without any cut, as quark masses are exactly accounted for. It is however wiser to use also in this case realistic cuts, in order to avoid regions which are computationally demanding and of dubious physical interpretation at this level of approximation
Parton level distributions can be easily defined in the include file. Corresponding values for each bin will be given after cross section computation in output .dat file. This feature might be useful when one wishes to compare partonic distributions with hadron level ones obtained after the call to PY4JET.
Fourjphact can compute or generate events for one final state at a time ( eg. $`u\overline{u}gg`$ or $`b\overline{b}c\overline{c}`$), or for all 20 final states with quarks (not top) and gluons at the same time. In this last case, the corresponding probability of every channel is determined or read from a file, and the generated events will have the correct fraction of all final states. This “one shot” option is often used when hadronization is required.
In the call to PY4JET(PMAX,IRAD,ICOM) the parameters PMAX, IRAD, ICOM are set respectively to 0.d0, 0, 0 in a data statement. Their meaning is explained in the previous Section and they can of course be changed if needed.
The partons have to be stored in the proper order before the call to PY4JET: this is unambiguous for $`q\overline{q}gg`$ while for $`4q`$ one has somehow to decide which of the two $`q\overline{q}`$ pair corresponds to the secondary emission. Such a distinction between first and second pair is not well defined in the case of ME. As the highest contribution comes, event by event, from the diagrams which have the lower $`q\overline{q}`$ invariant mass as secondary emission, we choose this configuration for giving the proper order to quarks. This we do also in the case of two identical flavours (e.g. $`u\overline{u}u\overline{u}`$).
Examples of results obtained with Fourjphact+Pythia and comparisons with other methods can be found in this report in Sec. 4.4.2, Sec. 5.3.3 and Sec. 5.3.4.
### 2.2 HERWIG
Like Pythia, Herwig is a general-purpose event generator which uses parton showering to simulate higher-order QCD effects. The main differences are the variables used in the parton showers, which are chosen to simplify the treatment of soft gluon coherence, and the hadronization model, which is based on cluster rather than string fragmentation. The current version, described here, is Herwig 6.1 . The program and documentation are available at
http://hepwww.rl.ac.uk/theory/seymour/herwig/
#### 2.2.1 Parton showers
The Herwig parton shower evolution is done in terms of the parton energy fraction $`z`$ and an angular variable $`\xi `$. In the parton splitting $`ijk`$, $`z_j=E_j/E_i`$ and $`\xi _{jk}=2(p_jp_k)/(E_jE_k)`$. Thus $`\xi _{jk}\theta _{jk}^2`$ for massless partons at small angles.
The values of $`z`$ are chosen according to the relevant DGLAP splitting functions and the distribution of $`\xi `$’s is determined by the Sudakov form factors. See e.g. for technical details. Coherence of soft gluon emission is simulated by angular ordering: each $`\xi `$ value must be smaller than the one for the previous branching of the parent parton.
The initial conditions for each parton cascade are determined by the configuration and colour structure of the primary hard process. The initial value of $`\xi `$ for the showering of parton $`j`$ is $`\xi _{jk}`$ where $`k`$ is the parton that is colour-connected to $`j`$. For example, in $`e^+e^{}q\overline{q}g`$ the gluon has a colour that is connected to the antiquark and an anticolour connected to the quark. Therefore the initial angle for the quark jet is the angle between the quark and the gluon. For the gluon jet, the initial angle is either the gluon-quark or gluon-antiquark angle, with equal probability.
In general, the hard process may involve several possible colour flows, which are unique and distinct only in the limit of an infinite number of colours, $`N_c\mathrm{}`$. For example in $`e^+e^{}q\overline{q}g_1g_2`$ either gluon 1 or gluon 2 may be connected to the quark. In the limit $`N_c\mathrm{}`$ these colour flows have distinct matrix elements-squared, $`|_1|^2`$ and $`|_2|^2`$. In Herwig colour flow 1 is chosen with probability $`|_1|^2/(|_1|^2+|_2|^2)`$ and flow 2 with probability $`|_2|^2/(|_1|^2+|_2|^2)`$, after using the full ($`N_c=3`$) matrix element to generate the momentum configuration. In this approximation, each final state has a unique colour flow which tells us how to limit the angles in each parton shower.
The parton showers are terminated as follows. For partons of mass $`m_i`$ there is a cutoff of the form $`Q_i=m_i+Q_0`$, and showering from any parton stops when a value of $`\xi `$ below $`Q_i^2/E_i^2`$ is selected for the next branching. The condition $`\xi >Q_i^2/E_i^2`$ corresponds to the “dead cone” for heavy quarks . Then the parton is put on mass-shell, or given a small non-zero effective mass in the case of gluons. Working backwards from these on-shell partons, one can now construct the virtual masses of all the internal lines of the shower, and the overall jet mass, from the energies and opening angles of the branchings. Finally one can assign the azimuthal angles of the branchings, including EPR-type correlations, and deduce all the 4-momenta in the shower.
Next the parton showers are used to replace the (on mass-shell) partons that were generated in the original hard process. This is done in such a way that the jet 3-momenta have the same directions as the original partons in the c.m. frame of the hard process, but they are boosted to conserve 4-momentum taking into account their extra masses.
We see that combining any tree-level hard process matrix element with parton showers is quite straightforward in Herwig. Double-counting is avoided, or at least suppressed, by angular ordering, which limits the showers to cones defined by the hard process and its colour structure. The price for this simplicity is that one must know both the overall ($`N_c=3`$) matrix element-squared and the separate ones ($`|_1|^2`$ etc) for all the possible colour flows in the limit $`N_c\mathrm{}`$.
One must bear in mind that results from combined matrix elements and parton showers are only likely to make sense if all the energy scales in the hard process being modelled by the matrix element are bigger, or at least not much smaller, than those in the parton showers. Otherwise, the structure of the final state will be determined mainly by the showers and the details of the matrix element become irrelevant. This is ensured in Herwig by a variable EMSCA, set by the hard process subroutine, which acts as an upper limit on the relative transverse momentum of any branching in the associated parton showers. For example, in the $`e^+e^{}`$ 4-jets matrix element option, discussed below in Sec. 2.2.4, EMSCA is (the square root of) the smallest of the invariant quantities $`s_{ij}=2p_ip_j`$ for the 4 partons generated in the hard process.
While the above procedure of attaching parton showers to a hard process generated by a tree-level matrix element may be straightforward, the problem of matching matrix elements and showers beyond tree level is certainly not. So far, this has only been done up to order $`\alpha _\text{S}`$ in Herwig (as in Jetset), for a limited class of processes including $`e^+e^{}q\overline{q}(g)`$. In Herwig the problem separates into two parts. First (“hard” matrix element corrections) there is a region of phase space that $`e^+e^{}q\overline{q}`$ \+ parton showers does not populate at all to order $`\alpha _\text{S}`$. That region can easily be filled by generating a gluon according to the matrix element. Second, there are the (“soft”) matrix element corrections that have to be applied inside the parton showers. As shown in Ref. , the right way to do this is to apply a correction not only to the first branching in each shower but also to every branching that is the “hardest so far”. This is especially important in Herwig where the evolution in $`\xi `$ means that several relatively soft (i.e. low $`p_t`$) wide-angle branchings can precede a harder one with a smaller angle.
To provide full matrix-element matching for 2-, 3- and 4-jets would mean extending the above procedure to next-to-next-to-leading order. There will be unpopulated regions of 4-parton phase space to be filled using the hard 4-jet matrix element, and “semi-hard” regions in which the 3-jet matrix element should be used in combination with order $`\alpha _\text{S}`$ “soft” corrections within a shower. However the bulk of the cross section will be in regions where order $`\alpha _\text{S}^2`$ corrections within the showers must be computed and applied – a daunting prospect.
One may, however, implement a less ambitious procedure for ‘combining’ 2, 3 and 4-jets so as to describe multijet distributions to leading order, which is discussed in Sec. 2.2.5.
#### 2.2.2 Hadronization
Hadronization in Herwig is done using a cluster model. First of all, any “on mass-shell” gluons at the ends of the parton showers are split into light quark-antiquark pairs. As mentioned above, a unique colour flow is generated for each final state, so that each final-state quark is uniquely colour-connected to an antiquark and vice-versa. These connected pairs can therefore form colour-singlet clusters carrying the combined flavour and 4-momentum of the pair. In the simplest case these clusters decay directly into pairs of hadrons according to the density of states for possible pairs of the right flavour. The transverse momentum $`300`$ MeV generated in hadronization is a reflection of the typical momentum release in cluster decay, which is determined by the cutoff $`Q_0`$, the quark masses and the QCD intrinsic scale $`\mathrm{\Lambda }`$.
If a cluster is too light to decay into two hadrons, it is converted into a single hadron of that flavour by donating some 4-momentum to a neighbouring cluster. If its mass is above a flavour-dependent value set by the parameter CLMAX (default value 3.35 GeV),
$$M_{jk}>[\mathrm{𝙲𝙻𝙼𝙰𝚇}^p+(Q_j+Q_k)^p]^{1/p}$$
where the power $`p`$ is given by a parameter CLPOW (default 2.0), it is split collinearly into two lighter clusters. A further parameter PSPLT (default 1.0) specifies the mass distribution of the resulting lighter clusters, which is taken to be proportional to $`M^{\mathrm{𝙿𝚂𝙿𝙻𝚃}}`$.
The cluster mass spectrum falls rapidly at high masses and its peak lies below the threshold for cluster splitting. One can show that these features are asymptotically independent of the energy scale of the hard process. However, there is always a finite probability of producing a very massive cluster. In this case sequential collinear splitting is invoked, leading to string-like hadronization.
#### 2.2.3 $`b`$-jet fragmentation
We concentrate here on primary $`b`$-quark showering and hadronization, leaving discussion of gluon $`b\overline{b}`$ to Sec. 7.4
The main point to note in connection with $`b`$-quark showering is the treatment of quark masses in Herwig parton showers. In the basic algorithm, the quantity $`m_i`$ appears only in the shower cutoff $`Q_i=m_i+Q_0`$, but this affects the distributions of $`\xi `$ and $`z`$ throughout the shower via the constraint
$$Q_j/(E_i\sqrt{\xi _{jk}})<z_j<1Q_k/(E_i\sqrt{\xi _{jk}})$$
at each branching $`ijk`$. Since this is always a low-energy cutoff it seems clear that the relevant value of $`m_i`$ is the pole or constituent mass. On the other hand a running mass might well be more appropriate in evaluating the hard process matrix element and the corresponding matrix element corrections.
In the process of $`b`$-quark hadronization, the input value of $`m_b`$ clearly affects the fraction of $`b`$-flavoured clusters that become a single B meson, the fractions that decay into a B meson and another meson, or into a B baryon and an antibaryon, and the fraction that are split into more clusters. Thus the properties of $`b`$-jets depend on the parameters $`m_b`$, CLMAX, CLPOW and PSPLT in a rather complicated way.
In practice the parameters CLMAX, CLPOW and PSPLT are tuned to global final-state properties and one needs extra parameters to describe $`b`$-jets. A parameter B1LIM has been introduced to allow clusters somewhat above the B$`\pi `$ threshold mass $`M_{th}`$ to form a single B meson if
$$M<M_{lim}=(1+\mathrm{𝙱𝟷𝙻𝙸𝙼})M_{th}.$$
The probability of such single-meson clustering is assumed to decrease linearly for $`M_{th}<M<M_{lim}`$. This has the effect of hardening the B spectrum if B1LIM is increased from the default value of zero.
Finally one should note that the properties of $`b`$-jets in Herwig are also affected by the parameters CLDIR and CLSMR, which control the decay angular distribution of clusters containing a perturbative quark (as opposed to the quark-antiquark pairs produced by the non-perturbative gluon splitting at the end of the parton showers – see above). If CLDIR=0, the decay of such a cluster is taken to be isotropic in its rest frame, as for other clusters. But if CLDIR=1 (the default value), the decay hadron carrying the flavour of the perturbative quark is assumed to continue in the same direction as that quark in the cluster rest-frame. This is suggested by the observation that the leading hadron in a quark jet preferentially carries the quark flavour. The value of CLSMR determines the amount of smearing \[exponential in $`(1\mathrm{cos}\theta )`$\] of this angular correlation. The default value of zero corresponds to perfect correlation. Thus increasing CLSMR tends to soften and broaden the B-hadron distribution in $`b`$-jets. In practice, the predicted spectrum tends to be too soft and CLSMR=0 is preferred.
In Herwig version 6.1, the parameters PSPLT, CLDIR and CLSMR have been converted into two-dimensional arrays, with the first element controlling clusters that do not contain a $`b`$-quark and the second those that do. Thus tuning of $`b`$-fragmentation can now be performed separately from other flavours, by setting CLDIR(2)=1 and varying PSPLT(2) and CLSMR(2). By reducing the value of PSPLT(2), a harder B-hadron spectrum can be achieved.
#### 2.2.4 4-jet matrix element + parton shower option (massless ME)
A new option available in Herwig version 6.1 is to generate events starting from the 4-parton processes $`e^+e^{}q\overline{q}gg`$ and $`e^+e^{}q\overline{q}q\overline{q}`$. The relevant process code is IPROC = 600 +IQ for primary quark flavour IQ or 600 for a sum over all flavours. The matrix elements used are those of Ellis Ross and Terrano and Catani and Seymour , which include the relative orientation of initial and final states but not quark masses. As explained in Sec. 2.2.1, the kinematic effects of quark masses are taken into account in the subsequent parton showers and in matching the showers to the momentum configurations generated according to the matrix elements. As also explained there, the variable $`\mathrm{𝙴𝙼𝚂𝙲𝙰}=\mathrm{min}\{\sqrt{s_{ij}}\}`$ sets a limit on the transverse momenta in the showers and is also used as the scale for $`\alpha _\text{S}`$. The latter feature has the effect of enhancing the regions of small $`s_{ij}`$ relative to matrix element calculations with $`\alpha _\text{S}`$ fixed.
To avoid soft and collinear divergences in the matrix elements, an internal parton resolution parameter Y4JT (default value 0.01) must be set. The interparton distance is calculated using either the Durham or Jade metric. This choice is governed by the logical parameter DURHAM (default .TRUE.). For reliability of the results, one should use the same metric for parton and final-state jet resolution, with a value of Y4JT smaller than the $`y_{cut}`$ value to be used for jet resolution.
#### 2.2.5 Combined 2,3 and 4-jet matrix element + parton shower option
As a result of discussions in the working group, a preliminary version of a combined 2,3 and 4-jet option based on Herwig 6.1 was developed. The strategy for combining matrix elements and parton showers follows that of , with some simplifications, as follows.
The program first generates conventional Herwig $`e^+e^{}`$ hadronic events starting from matched $`q\overline{q}`$ and $`q\overline{q}g`$ matrix elements (process code IPROC=100). After parton showering, the Durham clustering algorithm is applied, and those events with precisely four jets at resolution scale $`y_1\mathrm{𝚈𝟺𝙹𝚃}`$ (default value 0.008) are replaced by events generated using the (massless LO) 4-parton matrix element (IPROC=600), with Durham cutoff $`y_{ij}>\mathrm{𝚈𝟺𝙹𝚃}`$. The 4 parton momenta are distributed according to the matrix element multiplied by a weight factor, which for $`q\overline{q}gg`$ is
$$𝒲(y_1,y_3,y_4)=\frac{\alpha _\text{S}(y_3s)}{\alpha _\text{S}(y_1s)}\frac{\alpha _\text{S}(y_4s)}{\alpha _\text{S}(y_1s)}\mathrm{\Delta }_g(y_1s,y_3s)\mathrm{\Delta }_g(y_1s,y_4s)$$
(20)
where $`y_{3,4}`$ are the jet resolution values at which the partons are just resolved into 3,4 jets, and $`\mathrm{\Delta }_g`$ is the Sudakov form factor of the gluon (see e.g. ).
As explained in , the extra weight factor (20) is necessary to ensure smooth matching to the parton showers at small values of $`y_1`$ — more specifically, to cancel leading and next-to-leading logarithms of $`y_1`$. Since this factor is always less than unity, reweighting is simply achieved by rejecting configurations with $`𝒲(y_1,y_3,y_4)<`$ where $``$ is a random number.
After a 4-parton configuration has been generated, parton showers are generated in the usual way except that (for the 4-parton events only) parton branchings that would lead to sub-jets resolvable at resolution $`y_1`$ are vetoed. This means they are not allowed, but the evolution scale for subsequent branching is reduced as if they had occurred. Again, this is necessary to cancel LL and NLL $`y_1`$-dependence between ME and PS. In Herwig it is simply ensured by resetting $`\mathrm{𝙴𝙼𝚂𝙲𝙰}=\sqrt{y_1s}`$ after the 4-parton hard process.
Combining 2,3 and 4-jet events in Herwig 6.1 according to the above “replacement” algorithm is done by the (Fortran) program hw234jet.f. A prerelease version and some further discussion can be found at http://home.cern.ch/webber/ . To run the program one must link the slightly revised Herwig version 6.103, also available there.
### 2.3 ARIADNE
The Ariadne program is based on the Colour Dipole model where the QCD cascade is described in terms of gluon emissions from independent colour-dipoles between colour-connected partons. The program is described in detail elsewhere , and the following will mainly discuss issues related to gluon radiation off heavy quarks. Gluon splitting into heavy quarks in Ariadne is discussed in Sec. 7.5.
One of the main advantages of the dipole model is that, since gluons are emitted by the dipoles between partons, the interference between diagrams where a gluon is emitted by either of two partons is automatically taken into account, and there is no need to introduce explicit angular ordering. Another related advantage is that, since the first gluon emitted in an e<sup>+</sup>e$`{}_{}{}^{}q\overline{q}`$ event, again is emitted coherently by the $`q`$ and $`\overline{q}`$, the full leading order matrix element can be used explicitly in this emission, and correction procedures necessary in conventional parton shower models are not needed.
#### 2.3.1 Gluon radiation off heavy quarks
For heavy quarks, the default current version of the program uses an approximate extra suppression to suppress gluon radiation close to the direction of the quark to account for the dead-cone effect. This extra suppression can be switched off<sup>2</sup><sup>2</sup>2By setting the switch MSTA(19)=0 in the /ARDAT1/ common block. and, as discussed in Sec. 4.4.1, it seems that this actually improves the description of the heavy-to-light jet-rate measurement somewhat. Recently, the full massive leading order matrix element was implemented in Ariadne<sup>3</sup><sup>3</sup>3Not yet released. A prerelease can be obtained on request to leif@thep.lu.se for the first gluon emission, and it seems that this also describe jet rates a bit better than the approximate dead-cone suppression, although excluding mass effects still seems to give the best desctiption. This needs to be studied further.
### 2.4 APACIC++
Paradigm of the program:
Employ matrix elements to describe the production of jets,
model the evolution of jets with the parton shower.
#### 2.4.1 Introduction
As stated already in the introduction, due to various reasons, the modelling of multijet events in high–energy reactions becomes increasingly important with rising energies.
With emphasis on this modelling of multijet events, the program package Apacic++/Amegic++ has been developed only recently. The philosophy of the new approach presented here is to use matrix elements (ME) and parton showers (PS) in the corresponding regimes of their reliability : matrix elements are employed to describe the production of jets, and parton showers to model their evolution. A general algorithm to match them has been proposed and implemented in Apacic++ , the PS part of the package. The algorithm is based on the paradigm above, namely to restrict the validity of the ME’s for the description of particle emission to the regions of jet–production, i.e. to regions of comparably large angles and energies – or to large $`y_{\mathrm{cut}}`$ of the corresponding jet–clustering scheme. In contrast, the PS is restricted to the disjunct region of jet–evolution, i.e. small angles and low energies – or low $`y_{\mathrm{cut}}`$, respectively.
However, in its current state, the package is capable to deal with multijet production in $`e^+e^{}`$–annihilations only, where the jet–configurations available are determined by the ME generator. In addition to the generic ME part of the package, Amegic++ , interfaces to Debrecen and Excalibur are provided as well. The hadronization of the partons is left to well–established schemes. At the moment, an interface only to the hadronization in the Lund–string picture as implemented in Pythia is supplied.
The short description of the package follows closely the steps of event generation, namely
1. Initialization of matrix elements and jet rates,
2. Choice of jet structure of the single event,
3. Evolution of the jets with the parton shower, and
4. Hadronization.
#### 2.4.2 Initialization of matrix elements and jet rates
The use of matrix elements for the determination of the large–scale jet structure of the single events enforces their initialization and the calculation of the corresponding jet rates before the generation of single events. Since the description of the two other ME–generators can be found elsewhere, only the ME–part Amegic++ of the package will be discussed here. At the present stage, it is capable to deal with the following processes
$`e^+e^{}`$ $``$ $`\gamma ,Z(5)\text{QCD–jets}`$
$`e^+e^{}`$ $``$ $`(4)\text{fermions}`$ (21)
at tree–level in the Standard Model. All particles can be taken massless or massive, which allows for the inclusion of Higgs interactions. Effects due to photonic initial state radiation off the incoming electron pair can be included in the structure function approach.
Amegic++ constructs and integrates the matrix elements fully automatically. It proceeds in the following steps,
1. Building of topologies with unspecified internal lines and specified external legs in all combinations. Mapping of predefined Feynman–rules onto the topologies.
2. Construction of helicity amplitudes corresponding to the Feynman–amplitudes. Gauge test and transformation into a word–string, which is stored in a library.
3. Integration over the phasespace of the outgoing particles. Here, a significant acceleration is gained by using the compiled and linked word–strings out of the library.
As a result of this procedure, Amegic++’s source code of roughly 13 000 lines grows considerably to up to 200 000 lines when libraries for all possible processes are added.
However, as a well–known fact, the integration over the phasespace is plagued with real divergencies related to the soft and collinear emission of massless particles. To handle them, usually the phasespace is cut to avoid the dangerous regions. Then, outgoing particles are identified with jets, which are well–separated in phasespace with a measure $`y_{\mathrm{cut}}`$ depending on the jet–scheme. The package provides different jet–clustering schemes. Consequently, the cross–sections for $`n_j3`$ jets depend sensitively on the choice of the scheme and the corresponding $`y_{\mathrm{cut}}`$. Concentrating for the moment on pure QCD events and defining
$`\sigma _{\mathrm{QCD}}={\displaystyle \underset{q}{}}\sigma _{eeq\overline{q}}`$ (22)
the package provides three different schemes to determine jet rates. With $`\stackrel{~}{\sigma }`$ the various cross sections with the appropriate powers of $`\alpha _\text{S}`$ pulled out, the jet rates in the “direct” scheme read
$`_{n_j}^{\mathrm{dir}.}=\alpha _\text{S}^{n_j2}(\kappa _{n_j}s_{ee}){\displaystyle \frac{\stackrel{~}{\sigma }_{n_j}}{\sigma _{\mathrm{QCD}}}}\text{and}_2^{\mathrm{dir}.}=1{\displaystyle \underset{n_j=3}{}}_{n_j}^{\mathrm{dir}.},`$ (23)
whereas in the two “rescaled” schemes they are defined by
$`_{n_j}^{\mathrm{res}.1}`$ $`=`$ $`_{n_j}^{\mathrm{dir}.}_{n_j+1}^{\mathrm{dir}.}`$
$`_{n_j}^{\mathrm{res}.2}`$ $`=`$ $`_{n_j}^{\mathrm{dir}.}{\displaystyle \underset{m>n_j}{}}\left(1_m^{\mathrm{res}.2}\right)\text{and}_2^{\mathrm{res}.2}=1{\displaystyle \underset{n_j=3}{}}_{n_j}^{\mathrm{res}.2}.`$ (24)
To account for the effect of higher order corrections, the package supplies scale factors $`\kappa _S^{n_j}`$ for the corresponding $`n_j`$–jet rate. They enter in the form of $`\alpha _\text{S}=\alpha _\text{S}(\kappa _Ss_{ee})`$.
Going beyond pure QCD–events, the final states are divided into two ensembles, namely an electroweak one and the QCD ensemble. The former consists of all events with at least four fermions in the final state, where the normalization is given by the appropriate sum of the cross sections taken into account. This division obviously assigns a small amount of QCD events, namely the ones containing four quarks, to the electroweak ensemble. However, it should be noted, that so far this issue of electroweak events is still under further investigation.
#### 2.4.3 Choice of jet structure of the single event
The jet structure of the single events is now determined following the steps below,
1. The number of jets and their flavour structure are chosen according to the rates given above, Eqs.,23, 2.4.2.
2. The kinematical configuration is chosen. An appropriate number of equally distributed four vectors for the outgoing on–shell partons is produced. Their minimal $`y_{\mathrm{cut}}`$, $`y_{\mathrm{min}}`$, is forced to be larger than the $`y_{\mathrm{cut}}=y_0`$ used for the initialization of the jet rates. These fourvectors are then reweighted to reproduce the kinematical configurations as determined by the matrix element, potentially including the effect of higher order corrections. Defining $`|\stackrel{~}{}|^2(\mathrm{max})`$ the largest matrix element squared which can be obtained, Apacic++ again offers three choices for the corresponding weights, namely
$`𝒲^{\mathrm{L}.\mathrm{O}.}`$ $`=`$ $`{\displaystyle \frac{|\stackrel{~}{}|^2(p_i)}{|\stackrel{~}{}|^2(\mathrm{max})}}`$
$`𝒲^\alpha `$ $`=`$ $`\left[{\displaystyle \frac{\alpha _\text{S}(y_{\mathrm{min}}s_{ee})}{\alpha _\text{S}(y_0s_{ee})}}\right]^{n_j2}{\displaystyle \frac{|\stackrel{~}{}|^2(p_i)}{|\stackrel{~}{}|^2(\mathrm{max})}}`$ (25)
for leading order– and $`\alpha _\text{S}`$–corrected weights and a more involved one, which follows closely the reasoning of resummed jet rates. For example, in this scheme , the weight for three–jet events reads
$`𝒲^{\mathrm{NLL}}={\displaystyle \frac{\alpha _\text{S}(y_{\mathrm{min}}^{}s_{ee})}{\alpha _\text{S}(y_0s_{ee})}}{\displaystyle \frac{|\stackrel{~}{}|^2(p_i)}{|\stackrel{~}{}|^2(\mathrm{max})}}{\displaystyle \frac{\mathrm{\Delta }_g(y_{\mathrm{min}}^{}s_{ee})}{\mathrm{\Delta }_g(y_0s_{ee})}},`$ (26)
where $`y_{\mathrm{min}}^{}=\mathrm{min}\{y_{qg},y_{g\overline{q}}\}`$ the minimal $`y_{\mathrm{cut}}`$ related to the gluon of the $`q\overline{q}g`$–configuration of the three–jet event and the Sudakov form factors for the gluon, $`\mathrm{\Delta }_g`$ to be found in .
Note, that at this stage, the outgoing momenta are still on their mass–shell.
3. The colour configuration is determined. This is achieved by constructing relative probabilities $`𝒫_i`$ of the different parton histories. Here, Apacic++ provides three schemes, two of them based on the corresponding amplitudes $`_i`$ related to the single diagrams $`i`$ in the form
$`𝒫_i`$ $`=`$ $`{\displaystyle \frac{1}{_j|_j|^2}}|_i|^2\text{or}`$
$`𝒫_i`$ $`=`$ $`{\displaystyle \frac{1}{\underset{j}{}|_j\underset{l}{}_l^{}|}}|_i{\displaystyle \underset{k}{}}_k^{}|.`$ (27)
The third scheme relies on a shower oriented picture and was proposed in a similar fashion already in and discussed in Sec. 2.1.5. To illustrate this scheme, consider the diagram displayed in Figure 2. Its relative probability reads up to a suitable normalization
$`𝒫_i={\displaystyle \frac{1}{t_1}}P_{qqg}(z_{134}){\displaystyle \frac{1}{t_4}}P_{ggg}(z_{456}),`$ (28)
where the $`t_{ij}=(p_i+p_j)^2`$ and the $`z`$ are the energy fractions related to the various emissions encountered. The parton history and equivalently the colour configuration of the parton ensemble are then chosen according to their relative probabilities.
4. The task left now, is to use the parton history to supply the outgoing on–shell particles with virtual masses to allow them to experience a jet evolution via multiple emission of secondary partons. This is achieved by means of the appropriate Sudakov form factors for each of the outgoing legs, where the starting scale for the form factor is given by the internal $`t_i`$, like $`t_1`$ for the virtual mass of leg $`3`$ and $`t_4`$ for $`5`$ and $`6`$, in the exemplary diagram. To ensure, that no additional jet under the jurisdiction of the initial $`y_0`$ is produced by means of the parton shower, an appropriate veto is introduced into the subsequent shower algorithm. To account for local four momentum conservation during the change from on–shell to off–shell particles the kinematics are slightly rearranged, resulting in slight changes in the energy fractions and the opening angles of the outgoing partons.
#### 2.4.4 Evolution of the jets
The evolution of the jets proceeds in the standard way employing the Sudakov form factors . In addition to the usual switches allowing for the optional inclusion of prompt photons and azimuthal contributions, Apacic++ provides the possibility to use either the ordering by virtualities (LLA) or the ordering by angles (MLLA). However, in contrast to the pure MLLA–parton shower as performed in Herwig , Apacic++ uses a hybrid solution when switching to ordering by angles. Anticipating, that the MLLA–scheme is valid only in the domain of small angles, the first branching in each jet is done using the LLA–prescription, i.e. using the proper virtual mass as evolution parameter in the Sudakov form factor. After that, Apacic++ continues with the scaled angles as evolution parameters. In both cases, the Sudakov form factor yields the probability for no observable branch between scales $`t_+`$ and $`t_{}`$ and has the following form
$`\mathrm{\Delta }(t_+,t_{})=\mathrm{exp}\left\{{\displaystyle \underset{t_{}}{\overset{t_+}{}}}{\displaystyle \frac{dt}{t}}{\displaystyle \underset{z_{}(t)}{\overset{z_+(t)}{}}}𝑑z\alpha _\text{S}\left[p_{}^2(z,t)\right]P(z)\right\}.`$ (29)
The boundaries for the $`z`$ integration are given by
$`z_\pm ^{\mathrm{LLA}}(t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\pm {\displaystyle \frac{1}{2}}\sqrt{1{\displaystyle \frac{4t_0}{t}}}\text{and}`$
$`\sqrt{{\displaystyle \frac{t_0}{t}}}<`$ $`z^{\mathrm{MLLA}}`$ $`<1\sqrt{{\displaystyle \frac{t_0}{t}}},`$ (30)
where $`t_0`$ is the minimal virtual mass allowed in the parton shower. Within Apacic++, however, all partons leaving the jet–evolution have virtual mass $`t_{\mathrm{fin}}`$
$`t_{\mathrm{fin}}=\text{min}\{t_0,m_f^2\}.`$ (31)
The transversal momentum squared for the decays is
$`p_{}^2\stackrel{\mathrm{LLA}}{}z(1z)t\stackrel{\mathrm{MLLA}}{}z^2(1z)^2t`$ (32)
reflecting the interpretation of the $`t`$ in LLA and MLLA as the virtual mass of the decaying particle and the scaled opening angle of the branching, respectively.
#### 2.4.5 Treatment of mass effects
Within the framework of LLA parton showers, Apacic++ treats non–vanishing quark masses in the following way:
1. For the Sudakov form factors of Eq. 29, the minimal virtual mass $`t_0^f`$ is flavour dependent, see Eq. 31. Consequently, the minimal virtual mass for decaying quarks and gluons are changed, too.
2. The boundaries of the $`z`$–integration are determined by Eq. 2.4.4 but with the corresponding replacement
$`4t_0\left(\sqrt{t_0^b}+\sqrt{t_0^c}\right)^2`$ (33)
for branchings of the form $`abc`$. For example, this results in $`4t_0\left(m_b+\sqrt{t_0}\right)^2`$ for $`bbg`$ branchings and $`4t_04m_b^2`$ for $`gb\overline{b}`$ branchings, respectively.
3. Accordingly, for $`gq\overline{q}`$ splittings individual Sudakov form factors are summed. The $`tz`$–pair are chosen according to the sum, when picking the resulting quark flavour forbidden flavours are respected. For example, if $`t<4m_b^2`$, and equivalently, if the $`z`$ value results in quark energies $`E_q<m_q`$, $`b`$ quarks are not picked any more.
#### 2.4.6 Hadronization
At its present stage, Apacic++ performs the hadronization of the outgoing particles with the help of the Lund–string as provided by Pythia. For this purpose, an appropriate interface has been written and included. A similar interface for the cluster–hadronization of Herwig is planned. However, at this place, it should again be noted, that all particles leaving the parton shower of Apacic++ have a non–vanishing mass as given by Eq. 31. Therefore, before entering the Lund–string all particles have to be set on their mass–shell resulting in a small rescaling of their four–vectors.
#### 2.4.7 Summary : physics and computer features
Physics features :
1. The program package Apacic++/Amegic++ is designed for the modelling of multijet events. It is capable to produce and evaluate matrix elements for the production of up to five massive partons in QCD and at least all electroweak processes of the type $`e^+e^{}`$ four fermions allowed in the Standard Model. Additional interfaces to various different M.E. generators describing the production of multijet topologies are available, too.
2. The MEs are matched to the parton shower (PS) via a generically new matching algorithm. This algorithm is capable to deal with – in principle – any number of jets produced via the strong, weak or electromagnetic interaction on equal footing.
3. The hadronization is modelled with the LUND–string approach as provided by Jetset, the corresponding interface is provided, an similar interface to the cluster–hadronization of Herwig is planned.
Computer features:
1. The programming language is C++, allowing for a transparent and user–friendly programming style.
2. The package Apacic++/Amegic++ has been developed under Linux with the GNU–compilers. In addition, it has been tested under AIX, Digital Unix and IRIX.
3. Size of the package is :
Source code : Apacic++ $``$ 7 000 lines Amegic++ $``$ 13 000 lines Own libs : Amegic++ up to 200 000 lines
### 2.5 Tuning and tests of Apacic++ to reproduce event shape data
#### 2.5.1 Introduction
High precision measurements of event shape distributions and inclusive particle spectra, based on Lep1 data taken with the Delphi detector at Lep, are used to determine Apacic++ parameters. Extensive studies are performed to compare predictions of Apacic++ with Delphi data.
Definitions of used observables and a description of their measurement, together with a short description of the Delphi detector, can be found in
#### 2.5.2 The tuning procedure
The tuning procedure is based on a simultaneous fit of Monte Carlo parameters to physical observables, taking correlations between parameters into account. The fit is based on the minimisation of the variable
$$\chi ^2(\stackrel{}{p}):=\underset{\text{observables}}{}\underset{\text{bins}}{}\left(\frac{X_{meas.}X_{MC}(\stackrel{}{p})}{\sigma _{meas.}}\right)^2$$
(34)
The sum extends over all bins of all physical observables included in the fit, $`\sigma _{meas.}`$ being the total (statistical and systematic) error on the measured value $`X_{meas.}`$, $`X_{MC}(\stackrel{}{p})`$ being the Monte Carlo prediction of bin $`X`$ for the parameter setting $`\stackrel{}{p}`$. To perform the fit a fast prediction of $`X_{MC}`$ for any parameter setting $`\stackrel{}{p}`$ is needed. This is approximated by a Taylor expansion:
$$X_{MC}(p_1,p_2,\mathrm{},p_n)=A_0+\underset{i=1}{\overset{n}{}}B_ip_i+\underset{i=1}{\overset{n}{}}C_ip_i^2+\underset{i=1}{\overset{n1}{}}\underset{j=i+1}{\overset{n}{}}D_{ij}p_ip_j+\mathrm{}$$
(35)
The coefficients $`A_0,B_i,C_i`$ and $`D_{ij}`$ are extracted from a systematic parameter variation by applying a singular value decomposition.
For a detailed description of the tuning procedure, see .
#### 2.5.3 Tuning of Apacic++
The generator Apacic++ together with Amegic++, restricted to at most five massless jets from matrix element calculation, followed by a LLA parton shower, is chosen to be tuned. The initial jet-finder is the Durham algorithm, and fragmentation is achieved by the Lund string model.
The tuning of a new Monte Carlo generator is an iterative process. Each tuning triggers a learning process, resulting in improvements in the program, followed by a re-tuning. This process has not yet finally converged, but the quality of the program has reached a competitive state.
Each iteration starts with the selection of parameters to be tuned and the definition of their variation ranges. Since the fit result is difficult to predict, the variation ranges have to be chosen generously, degrading the precision of the fit. A second tuning around the optimal values provides further improvements to the generator.
The fit is performed to a sample of observables, sensitive to the varied parameters. Exchanges in the composition of the observables give hints to systematic uncertainties of the tuning result.
#### 2.5.4 Apacic++ parameters
Within Apacic++ there are parameters describing the matrix element calculation, the parton shower evolution and the Lund fragmentation.
* Matrix element
+ $`y_{cut}^{ini}`$
Emissions of colour charged partons are restricted to resolution parameters $`y_{cut}>y_{cut}^{ini}`$. In previous tunings of Apacic++ an adequate agreement to the reference data could only be achieved for large values of $`y_{cut}^{ini}`$ $`(0.05)`$, resulting in predictions for hard QCD processes dominated by the parton shower. For that reason $`y_{cut}^{ini}`$ was fixed to some sensible value ($`y_{cut}^{ini}=0.005`$). Improvements by re-weighting the kinematic distribution of jets cured this problem. One of the future projects will be to restore $`y_{cut}^{ini}`$ to the list of tuning parameters.
+ $`\kappa _s^{3,4,5}`$
Due to the truncation of the perturbative expansion, matrix element calculations show a significant dependence on the QCD renormalisation scale. Apacic++ accounts for these dependences by a scale parameter $`\kappa _s^{3,4,5}`$ for each $`n`$-jet configuration: $`\alpha _\text{S}=\alpha _\text{S}(\kappa _s^ns)`$
* parton shower
+ $`\alpha _\text{S}(M_\mathrm{Z}^2)`$
The strong coupling constant $`\alpha _\text{S}(M_\mathrm{Z}^2)`$ is responsible for the parton shower evolution
+ cutoff PS
The parton shower ends at a given energy scale, where fragmentation starts<sup>4</sup><sup>4</sup>4The parameter “cutoff PS” in Apacic++ is different from the cutoff parameter $`q_0^2`$ in Pythia: $`4\text{cutoff}=q_0^2`$.
* fragmentation
+ Lund A,B
Lund A and B enter the Lund fragmentation function. Due to the strong anticorrelation between A and B it is sufficient to tune one and keep the other fixed.
+ $`\sigma _q`$
The width of the Gaussian distribution of transverse momentum for fragmentation quarks is given by $`\sigma _q`$.
Table 2 summarises the parameters considered, their variation ranges and an illustrative tuning result.
#### 2.5.5 Data distributions
Measurements of event shape distributions and inclusive particle spectra are taken from ; for definitions of the observables used see there.
Apacic++ parameters from Table 2 are simultaneously fitted to a set of data distributions. For the fit result depending on the composition of the data set some systematic checks are performed to estimate the stability of the fit. The strategy followed within the composition is to include at least one distribution that is sensitive to the parameters being fitted. Within this constraint systematic exchanges in the composition are performed to study uncertainties in the fit result. Table 3 gives a summary of 15 different compositions used.
#### 2.5.6 Results
Predictions of Apacic++ event shape distributions, jet rates and inclusive particle spectra are compared to established Monte Carlo generators (like Pythia, Herwig, Ariadne) and to Delphi data. Figures 3,4,5,6 and 7 give an overview of the behaviour and the relative (dis-)advantages of Apacic++.
#### 2.5.7 Conclusion and outlook
Apacic++ parameters have been tuned to various sets of Delphi event shape distributions, jet rates and inclusive particle spectra.
The fits converged, the tuned parameters came out to be basically reasonable: The parameter $`y_{cut}^{ini}`$ has been fixed in the latest tuning. The parameter for the cutoff of the parton shower is high, giving large weight to the fragmentation and minor to the parton shower. This has to be investigated.
Apacic++ is able to predict all examined observables reasonably well. Still none of the examined Monte Carlo generators is able to predict the tail of the $`p_t^{out}`$ distribution (see however Sec. 3.2).
## 3 INCLUSIVE (ALL FLAVOUR) JET RATES
### 3.1 Tuning issues
During the Lep1 phase a qualitative improvement of the description of the hadronic final state by parton shower fragmentation models has been reached, mainly due to the possibility to precisely tune the models to a vast amount of high quality data . For this task flexible tuning procedures were used allowing interpolation between model responses generated with different parameter settings .
The effects on the model response of the individual parameters of the two major aspects of the models – the parton shower and the actual hadronisation phase – turn out to be strongly correlated. This requires one to determine the most important model parameters in global fits to high statistic event shape and inclusive charged particle spectra and to identified particle data. A recent example for such a fit is discussed in .
### 3.2 Model performance and multi-jet rates
It turns out that the string as well as the cluster hadronisation model are able to represent the major features of particle production, especially the identified particle rates, reasonably well. More detailed discussions can be found in .
Distributions depending mainly on the parton shower phase of the models are in general very well represented. Especially for most of the event shape distributions, data and models agree within a few percent. There are two important exceptions to this rule:
Firstly the tail of the transverse momentum distribution of particles out of the event plane is underestimated by about 30% by most models . A possible explanation for this deficiency is that the parton shower models account for part of the angular structure of multi-jet events by tracing the polarization of the emitted gluons (see e.g. ) to further splittings. This approximation cannot account for interference effects like a full matrix-element calculation. It should be emphasized, however, that the most recent tuning of the latest version of Herwig shows a remarkable improvement of the $`p_t^{out}`$-description. This distribution (see Figure 8) now seems almost perfectly reproduced.
The second exception concerns the inability of Herwig and Pythia/Jetset to simultaneously describe different multi-jet rates with the precision desired by the experiments. This can already be seen from the $`y_{cut}`$ dependence of the multi-jet rates shown in Figure 9 but is more clearly evident from the direct data/model comparisons in Figures 10, 11 and 12. For a well represented three-jet rate, as was perhaps required in the tunings of the Delphi Collaboration , the (differential) four and more jet-rates are systematically overestimated (underestimated) by Herwig or Pythia/Jetset, respectively. This is already observed for the four-jet rate, which is of special importance at Lep2, but is even more so for the five-jet rate. This general trend remains valid even for the aforementioned latest version/tune of Herwig . Depending on the strategy followed by the experiments, this misrepresentation can be distributed differently among the individual jet-rates. For example, the Opal tuning mediates between the rates (see Figure 12). The general discrepancy between multi-jet-rate data and the corresponding model predictions has been reported during the workshop by all experiments.
Only Ariadne so far is able to well represent all jet-rates simultaneously (see Figures 10, 11 and 12). The likely explanation for this difference between Ariadne and the other models lies in the matching between the parton shower and the first order $`\mathrm{q}\overline{\mathrm{q}}\mathrm{g}`$ matrix element simulations performed in Pythia and Herwig in order to well represent the initial hard gluon radiation. This matching is not needed in the dipole model implemented in Ariadne as here the splitting probability for all splittings is given by the lowest order matrix element expression. The “opposite” behavior observed for Pythia and Herwig may indicate, however, that a better agreement may be reachable by suitably improving the matching procedure.
### 3.3 Residual uncertainties
Residual uncertainties due to the imperfect description of multi-jet rates are difficult to review globally as they will depend critically on the individual analyses as well as on the tuning chosen by the different experiments.
An incorrect 4-jet rate at high energies may require a reweighting of the Monte Carlo to properly account for the QCD background in W or Higgs analyses. Due to the correlation of the number of jets of an event with other properties, e.g. the charged multiplicity or the momentum spectrum, this is likely to have unwanted side effects. An incorrect value of the strong coupling (which in the tunings is often fixed by the 3-jet rate) may cause an incorrect energy extrapolation of the models which is hard to control at high energy because of the limited data statistics.
A possible strategy for a determination of systematic error for a QCD type observable such as the four jet rate at Lep2 energies may be the following: The quality of the description of the observable is checked at the Z<sup>0</sup>. A possible misrepresentation at the Z<sup>0</sup> and at high energy is corrected using the same correction factor. A large fraction of the deviation of the correction factor from unity has to be taken as systematic uncertainty of the correction factor at high energy, since the reason for the bad data description, and consequently a possible energy dependence, is unknown.
The additional error for the uncertainty of the energy evolution of the model will in general be small. In the case of the four-jet cross section at high $`y_c`$, it will be dominated by the uncertainty of the strong coupling. From the expected QCD evolution of the four-jet rate this uncertainty is at $`\sqrt{s}=200`$ GeV:
$$\frac{\delta R_4(\sqrt{s})}{R_4(M_\mathrm{Z})}=4b\mathrm{ln}\left(\frac{\sqrt{s}}{M_\mathrm{Z}}\right)\delta \alpha _\text{S}(M_\mathrm{Z})1.9\delta \alpha _\text{S}(M_\mathrm{Z})$$
Here $`b=(332n_f)/12\pi 0.61`$. For an optimistically reachable error of $`\alpha _\text{S}`$ in the models of $`\delta \alpha _\text{S}(M_\mathrm{Z})=0.003`$ this yields $`0.6\%`$.
Employing alternative models and alternative model tunings will provide an important cross check of the above error estimate. A model which correctly represents the four-jet cross section at the Z<sup>0</sup>, but overestimates the three-jet cross section ($`\alpha _\text{S}`$) by about 10% (compare Figure 9 at $`y_{cut}0.01\mathrm{to}0.02`$) may in fact lead to a more optimistic error estimate. The error for the correction factor will vanish in this case at the expense of an increased error of the model extrapolation. This error, however, is still small ($`1.9\delta \alpha _\text{S}=1.9\times 0.0122\%`$).
For some analyses already today the abovementioned deficiency of the multi-jet description of the models leads to important contributions to the systematic error. An example is the Delphi measurement of the W pair production cross section with a fully hadronic final state. A systematic error of 5% (including a possible misrepresentation of the jet angular distributions) is here assigned to the major background of QCD events. This error was estimated by comparing different (uncorrected) models and dominates the overall systematics. With increasing statistics this systematic error will be of similar size to the statistical error. Delphi therefore starts to employ Ariadne which certainly in terms of the jet rates provides the best description of the data, as an alternative model for the full simulation.
## 4 STUDY OF MASS EFFECTS IN 3- AND 4-JET RATES
### 4.1 Introduction
The aim of this Section is to study the theoretical precision in the modeling of the rate of QCD processes leading to 3 and 4 final state jets, at Lep1 and Lep2, and involving $`b`$ quarks. This is important both to help understanding how to treat mass effects in our phenomenological QCD models, and to ensure precise enough control of backgrounds to new particle searches with $`b`$ quarks in the final state, such as for instance the Higgs search. With this aim in mind we compare mass effects on jet rates in the different MC approaches both to data and to analytic calculations.
It is natural to consider that because of their higher mass, $`b`$ quarks must from kinematics radiate fewer gluons than light quarks. More generally such a suppression enters in what is often refered to as the dead cone effect. What we want to know is how well the magnitude of this suppression is modeled in our Monte-Carlo approaches, and what is the related theoretical uncertainty. In some sense this question can also be formulated as that of specifying the appropriate mass which should be used for the $`b`$ quark.
From basic kinematics arguments, it can be shown that the magnitude of the suppression of gluon radiation from mass effects should scale as $`m_b^2/(s.y)`$, where $`m_b,s`$ and $`y`$ are the $`b`$ quark mass, the collision energy squared and jet resolution parameter, respectively. From this scaling law can be anticipated that the effects on jet rates are reduced at Lep2 energies as compared to Lep1, but that they can still remain substantial if jets are defined with a very small $`y`$ parameter.
In practice, from the results obtained, we find that the above scaling law is very approximate, and can really only be used to give a rough indication of the magnitude of effects. In addition to the purely kinematic effects, resulting from the more limited phase-space, significant corrections - often referred to as dynamic mass effects - arise when taking into account properly the $`b`$ quark mass in the matrix elements for 3 or 4 final state partons. Moreover, when comparing the data to the behaviour of the analytic calculations or to that of the MC, one must be careful to appropriately take into account in a consistent way effects resulting in final state jets with $`b`$ quarks other than the radiation of gluons off $`b`$ quarks, such as processes with a gluon splitting into $`b\overline{b}`$.
In what follows, we first describe the method used for the evaluation. This involves describing briefly the experimental analysis and the procedure enabling meaningful comparisons with analytic calculations and with predictions from Monte-Carlo generators. We then describe the analytic calculations which are available, show how these can be parametrised as a function of $`m_b,s`$ and $`y`$ to enable the above mentioned comparisons, as well as extrapolations. A proposed definition for the theoretical uncertainty to be quoted is also presented and discussed. Then follows a Section presenting the results of the comparisons of the different Monte-Carlo approaches studied (Ariadne, Pythia, Herwig and Fourjphact) with the analytic calculations and with the Delphi data (some of which is still preliminary). The comparisons with data are performed using the high statistics data available from Delphi at Lep1. Several of the used Monte-Carlo programs provide different options for the treatment of mass effects. Also, some have recently benefited from improvements. The main features of the different results and behaviours are described. Finally, remaining issues are discussed. One of these concerns the contribution to the observables used in the evaluation from processes involving gluon splitting into $`b\overline{b}`$. A quantitative measure of the precision of the description provided by the different programs for the 3 and 4 jet cases as a function of $`y`$, at both Lep1 and Lep2 energies, is given.
### 4.2 Procedure used for evaluation
#### 4.2.1 Appropriate choice of experimental observable
Any observable with an explicit mass dependence can be computed either using the quark pole mass ($`M_q`$) or the running mass ($`m_q(\mu )`$). Both methods are equally valid, and the results have to be the same if computed to all orders in perturbation theory. However at any fixed finite order the predictions using both schemes do not necessarily agree. At a fixed order, calculations have a dependence on the scale at which the observable is computed. This dependence reveals the size of the higher order terms, as the sum of all the series must have no dependence on the scale.
Aleph has performed a study of the $`b`$ mass effects on several event shape observables comparing their Leading Order (LO) and NLO terms and hadronization effects. Some observables have quite large mass dependence as shown in that study. However as the raw measurements have to be corrected accounting for the hadronization effects, it so happens that some observables have a correction even larger than the size of the effect to be measured.
Therefore the ideal observable to study quark mass effects will exhibit a large mass dependence, low higher-order corrections and small hadronization effects.
#### 4.2.2 $`b`$ quark mass effects on the 3-jet ratio
In order to study the $`b`$ mass effects, both Aleph and Delphi have chosen an observable which fully complies with the above requirements. The observable is the following:
$$R_3^b\mathrm{}(y_c)=\frac{R_3^b(y_c)}{R_3^{\mathrm{}}(y_c)}$$
(36)
where $`R_3^b(y_c)`$ and $`R_3^{\mathrm{}}(y_c)`$ are the $`b`$-quark and light-quark 3-jet rates as defined in Eq. 2.
#### 4.2.3 Data analysis
The jets of the hadronic events can be reconstructed by means of a jet finding algorithm (e.g. Durham or Cambridge – see Sec. 1.2). By choosing the appropriate jet resolution parameter $`y`$, it is possible to force the reconstruction of just three jets in every event. Then, a set of quality cuts are applied over each jet (e.g. minimum charged multiplicity, enough visible energy, …).
#### 4.2.4 Flavour definition
The criterion adopted in and is such that the flavour of the event is defined as that of the quark coupled to the $`Z`$ in the $`Zq\overline{q}`$ vertex. By convention, the production of $`b`$ quarks via the splitting of a bremsstrahlung gluon is ignored in this definition. It was shown that for this particular 3-jet observable, the corresponding effect is practically negligible. As will be seen in Sec. 4.4.4, this is no longer the case in the case of 4-jets.
#### 4.2.5 Flavour tagging technique
The sample of hadronic events is split in two categories, with both events strongly enriched in $`b\overline{b}`$, by using $`b`$-tagging, and events strongly enriched in light quarks, by using $`antib`$-tagging.
The tagging procedure normally uses the impact parameter information of all charged particles in the event. Lep experiments are equipped with silicon vertex detectors allowing accurate track reconstruction. $`b\overline{b}`$ events can be selected by considering the presence of particles with large impact parameter significance (impact parameter over its error), while $`\mathrm{}\overline{\mathrm{}}`$ events can be tagged just by the lack of those tracks. In addition the presence or absence of a reconstructed secondary vertex can be used in the $`b`$-tagging criterion.
#### 4.2.6 Measuring and correcting $`R_3^b\mathrm{}`$
The measured value of our observable can be computed from the ratio of the reconstructed three-jet rates for $`b`$ and light quark tagged flavours ($`R_{3q}^{\mathrm{obs}}(y_\mathrm{c}),q=b,\mathrm{}`$):
$$R_3^{b\mathrm{}\mathrm{obs}}(y_\mathrm{c})=R_{3b}^{\mathrm{obs}}(y_\mathrm{c})/R_3\mathrm{}^{\mathrm{obs}}(y_\mathrm{c})$$
This raw value of the observable must be converted into a parton level one which may be compared with the LO and NLO calculations . The correction method accounts for the detector effects, biases introduced in the flavour tagging, and hadronization effects.
The contribution of each flavour, $`R_{3q}^i`$, to the observed three-jet cross section is given by:
$$R_{3q}^{\mathrm{obs}}=c_q^bR_{3q}^b+c_q^cR_{3q}^c+c_q^{\mathrm{}}R_{3q}^{\mathrm{}}$$
where $`c_q^i`$ represents the flavour content for $`i=b,c,\mathrm{}`$ in each of the experimentally tagged samples.
The reconstruction level and parton level three-jet rates for each flavour and defined sample ($`R_{3q}^i`$ and $`R_{3q}^{\mathrm{par}}`$, respectively) are related by:
$$R_{3q}^i(y_\mathrm{c})=d_{3q}^i(y_\mathrm{c})h_{3i}(y_\mathrm{c})R_{3i}^{\mathrm{par}}(y_\mathrm{c})$$
where $`d_{3q}^i`$ and $`h_{3i}`$ are correction factors, accounting, respectively, for detector acceptance and tagging effects (deduced from the modeling of the detectors response to hadronic events), and for hadronization effects (estimated by comparing the hadron and parton level distributions obtained from MC programs).
By taking the jet rates from $`c`$ quarks equal to those from light quarks ($`R_{3c}^{\mathrm{par}}R_3\mathrm{}^{\mathrm{par}}`$), the measured jet rates can be expressed as:
$$R_{3b}^{\mathrm{obs}}(y_\mathrm{c})=A_b(y_\mathrm{c})R_{3b}^{\mathrm{par}}(y_\mathrm{c})+B_b(y_\mathrm{c})R_3\mathrm{}^{\mathrm{par}}(y_\mathrm{c})$$
$$R_3\mathrm{}^{\mathrm{obs}}(y_\mathrm{c})=A_{\mathrm{}}(y_\mathrm{c})R_{3b}^{\mathrm{par}}(y_\mathrm{c})+B_{\mathrm{}}(y_\mathrm{c})R_3\mathrm{}^{\mathrm{par}}(y_\mathrm{c})$$
where $`A_q`$ and $`B_q`$ are a redefinition of the original set of parameters: $`c_q^i,d_{3q}^i`$ and $`h_{3i}`$. This parametrization allows expressing the corrected observable as:
$$R_3^b\mathrm{}(y_\mathrm{c})=\frac{R_{3b}^{\mathrm{par}}}{R_3\mathrm{}^{\mathrm{par}}}=\frac{B_bB_{\mathrm{}}R_3^{b\mathrm{}\mathrm{obs}}}{A_{\mathrm{}}R_3^{b\mathrm{}\mathrm{obs}}A_b}$$
The results are described in Sec. 4.4. At $`\sqrt{s}=M_\mathrm{Z}`$ the total correction to the raw $`R_3^b\mathrm{}`$ is typically about 10%, the bulk of which corresponds to the detector and tagging effects, while hadronization corrections are of the order of 1%.
Uncertainties in the detector modeling were studied. The main uncertainty results from the limited statistics of fully simulated events with detector response, resulting in limited knowledge of the factors $`c_q^i`$ describing the flavour content of each of the tagged samples. This error was estimated to be roughly 0.3%, but has a strong dependence on the jet resolution parameter $`y_c`$ as it is directly related to the statistics of the three-jet simulated sample. Uncertainties from the modeling of the hadronization were also studied, and are described in Sec. 4.5.
#### 4.2.7 $`b`$ quark mass effects on the 4-jet ratio
The study of the $`b`$ quark mass in the 4-jet rate was performed by Delphi in . The analysis uses an observable defined in a similar manner to $`R_3^b\mathrm{}`$:
$$R_4^b\mathrm{}(y_c)=\frac{R_4^b(y_c)}{R_4^{\mathrm{}}(y_c)}$$
(37)
where $`R_4^b(y_c)`$ and $`R_4^{\mathrm{}}(y_c)`$ are the $`b`$-quark and light-quark 4-jet rates as defined in Eq. 2.
The analysis was similar to the $`R_3^b\mathrm{}`$ and also reveals a clear dependence of the 4-jet rate on the quark mass. However $`R_4^b\mathrm{}`$ has only been computed at LO level (see Sec. 4.3). The result of the analysis are shown in Sec. 4.4.
### 4.3 Analytical calculations
The next-to-leading order (NLO) matrix element (ME) calculation for the process $`e^+e^{}3`$ jets, with complete quark mass effects, has been performed independently by three groups . These predictions are in agreement with each other and were successfully used in the measurements of the bottom quark mass far above threshold and in the precision tests of the universality of the strong interaction at the $`Z`$-pole. Instead, only leading order (LO) predictions for 4-jet final states with heavy quarks are available at present.
In this Section we discuss in detail the ME calculation for the process $`e^+e^{}4`$ jets with quark mass corrections. A procedure to estimate the theoretical uncertainty of the LO calculation, i.e. of the expected higher order corrections, is also described. First, we present and test this procedure in the 3-jet case, where the recent available NLO corrections can be compared with the LO calculation. Then, these results are extrapolated to the 4-jet case for which the higher order corrections are estimated.
Since quarks are not free particles it seems natural to treat their mass like a coupling constant. In other words, we can work with quark masses defined in several renormalization schemes. As was already pointed out in Sec. 4.2.1, the physical result cannot depend on which mass definition is used in the ME calculation, but if at a fixed order in perturbation theory the ME calculation gives different results for different quark mass definitions, this difference should come from higher order corrections. Therefore, at a given order, the spread of the results for different mass definitions can be taken as an estimate of the theoretical uncertainty of the calculation, i.e. as an estimate of higher order corrections.
In the following we consider the ME calculation for the production of bottom quarks through the processes $`e^+e^{}`$ 3-jets, 4-jets in two different schemes: the perturbative pole scheme and the running scheme. In the former, the calculation is performed with the perturbative pole mass $`M_b5`$ GeV. In the later, we use the $`\overline{MS}`$ scheme running mass $`m_b(\mu )`$, normalized at the center of mass energy of the collision, $`m_b(M_\mathrm{Z})3`$ GeV at Lep1 energies. The two results define a band which can be taken as an estimate of the higher order corrections.
At the NLO the bottom quark 3-jet cross section receives contributions from one-loop corrected three parton final states, $`e^+e^{}Z,\gamma ^{}b\overline{b}g`$, and tree level four parton final states, $`e^+e^{}Z,\gamma ^{}b\overline{b}gg,b\overline{b}b\overline{b},b\overline{b}q\overline{q}`$, with $`qb`$. These contributions can be handled and classified through the different cuts to the bubble diagrams of Figure 13 giving rise to three and four parton final states.
We first evaluate the ratio of three-jet rates defined by Eq. 36, which can be interpreted as the measure of the suppression of gluon radiation off $`b`$-quarks with respect to gluon radiation off light quarks, $`\mathrm{}=u,d,s`$.In this ratio, most of the electroweak corrections cancel out.
In Figure 14, the $`R_3^b\mathrm{}`$ observable is presented at NLO for the Durham and the Cambridge algorithms at the $`Z`$-peak energies, $`\sqrt{s}=M_\mathrm{Z}`$, in the running mass (NLO-$`m_b(M_\mathrm{Z})`$) and the pole mass (NLO-$`M_b`$) schemes. For comparison, the LO results – LO-$`m_b(M_\mathrm{Z})`$ and LO-$`M_b`$ – are also plotted. At a fixed order, the band defined by the results in both schemes is taken as our estimate of the theoretical uncertainty of the calculation at this order. As one would naturally expect, the width of this band is reduced at the NLO with respect to the LO result, roughly by a factor two in Durham or even more in Cambridge. Furthermore, we found that the two LO predictions bound the NLO results in both algorithms. This suggests that higher order contributions cannot be too large and may be bounded by the lower order results.
Next we consider the process $`e^+e^{}4`$ jets. The ratio of four-jet rates is defined in Eq. 37. At LO, only four parton final state cuts to the bubble diagrams of Figure 13 have to be considered. Let’s interpret for the moment, $`R_4^b\mathrm{}`$ as the suppression of gluon radiation from primary $`b`$-quarks in four-jet events. By primary quarks we mean that the flavour of an event is defined by the flavour of the quark directly coupled to the $`Z`$ or $`\gamma ^{}`$ bosons. This implies that events where a bottom-antibottom pair is produced from gluon radiation off a light quark pair (the so-called gluon splitting into $`b\overline{b}`$) are considered as light events, even though they actually involve $`b`$ quarks, and their contribution is added to the denominator of Eq. 37). In $`R_3^b\mathrm{}`$, the contribution from gluon splitting into $`b\overline{b}`$ was calculated to be small, resulting in effects of the order of a few permil. This is not the case for four-jet events, where larger contributions can justify, as will be explained in Sec. 4.4.4, the consideration of a different convention for the jet rate ratio defined above, in which the numerator receives contributions from all events involving $`b`$-quarks.
At LO and for the definition with the primary quark convention, the $`R_4^b\mathrm{}`$ ratio can be parameterized in the following way
$$R_4^b\mathrm{}=1+\frac{m_b^2}{s}\left[\sigma _V(s)H_V(\frac{m_b^2}{s},y_c)+\sigma _A(s)H_A(\frac{m_b^2}{s},y_c)\right],$$
(38)
were $`m_b`$ can be either the pole mass, $`M_b`$, or the running mass at some renormalization scale, $`m_b(\mu )`$, typically $`\mu =\sqrt{s}`$, of the bottom quark, $`\sigma _{V,A}`$ is a function of the vector (axial-vector) couplings of the quarks to the $`Z`$-boson and the photon, and $`H_{V,A}`$ gives the behaviour as a function of the resolution parameter $`y_c`$, and can also contain a small residual dependence on the ratio $`m_b^2/s`$.
In Figure 15, the $`R_4^b\mathrm{}`$ ratio is shown at LO in the Durham algorithm at the center of mass energies $`\sqrt{s}=M_\mathrm{Z}`$ and $`\sqrt{s}=189`$ GeV. For the same center of mass energy, the suppression of gluon radiation from $`b`$-quarks is a larger effect in four-jet events than in three-jet events. At $`\sqrt{s}=M_\mathrm{Z}`$, it amounts to roughly $`10\%`$, which also corresponds to the difference between the two LO predictions, LO-$`m_b(M_\mathrm{Z})`$ and LO-$`M_b`$, taken as the theoretical uncertainty of the LO prediction. If the behaviour for three-jet events really could be extrapolated to the four-jet case, we would expect that this difference be reduced hopefully by half, if the still uncalculated NLO corrections were included. At $`\sqrt{s}=189`$ GeV, we get a plot that is roughly scaled by a factor $`4`$ with respect to the result at the Lep1 energies. But even in this case, the theoretical uncertainty, i.e. the difference between the two LO predictions, can be as large as $`5\%`$ for small values of the jet resolution parameter $`y_c`$.
### 4.4 Comparisons for $`R_3^b\mathrm{}`$ and $`R_4^b\mathrm{}`$ at LEP1 and LEP2
Following the procedure described above, we compare at both Lep1 and Lep2 energies, the double jet rate ratios $`R_3^b\mathrm{}`$ and $`R_4^b\mathrm{}`$ obtained with the different Monte-Carlo approaches studied (Ariadne, Pythia, Herwig and Fourjphact) with the analytic calculations, and, in the case of Lep1, with the data. All these comparisons are done at the level of partons (hadronization and non-perturbative effects are considered in Sec. 4.5). In Figures 16, 17 and 18 are indicated both the LO matrix element results for $`m_b=`$ 3 and 5 GeV and the NLO results for $`m_b=`$ 3 GeV, in the case of the 3-jet rates at Lep1. In Figures 19, 20, 21 and 22 are indicated only the LO results for $`m_b=`$ 3 and 5 GeV for the 4-jet rates at Lep1. In Figures 23, 24 and 25 are indicated the LO results for $`m_b=`$ 3 and 5 GeV for the 4-jet rates at Lep2. The band defined by the pair of LO curves is in each case taken, conservatively, to represent the theoretical uncertainty in the ratios $`R_3^b\mathrm{}`$ and $`R_4^b\mathrm{}`$, as was explained in Sec. 4.3.
#### 4.4.1 Three jet rates at $`\sqrt{s}=M_\mathrm{Z}`$
In Figures 16, 17 and 18 are shown the comparisons for $`R_3^b\mathrm{}`$ at $`\sqrt{s}=M_\mathrm{Z}`$ between the analytic calculations, the data and the Monte-Carlo generators Pythia, Herwig, Ariadne and Apacic++, respectively.
The Pythia comparison in Figure 16 shows three curves corresponding to:
* The initial treatment, with mass effects only present in the limitation of the phase-space available in $`bbg`$ branchings in the parton shower, but not in the kinematics of these branchings, and with massless matrix elements used in the matching procedure applied to the 3 jets generated (Jetset versions $``$ 7.3, or any present version with the switch MSTJ(47) set to 1 to turn off subsequent additional mass effects).
* Mass effects present in the limitation of the phase-space available in $`bbg`$ branchings in the parton shower, but not in the kinematics of these branchings, and with massive matrix elements used in the matching procedure applied to the 3 jets generated (Pythia versions 5.7-6.125).
* Mass effects present in the limitation of the phase-space available in $`bbg`$ branchings in the parton shower, in the kinematics of these branchings, and with massive matrix elements used in the matching procedure applied to the 3 jets generated (Pythia versions $``$ 6.130).
As can be seen the recent changes consisting in introducing mass effects at all stages in the treatment (see Sec. 2.1.1) result in a very good behaviour. In the prior versions which are presently still used by a majority of Lep experiments, a behaviour almost as good is obtained by using massless expressions for the matching procedure of the generated 3 jets (switch MSTJ(47) set to 1).
The Herwig comparison in Figure 17 shows a single curves corresponding to the massive options. A reasonable description is evident, even if the predicted rate is slightly lower than the data and NLO result.
The Ariadne comparison in Figure 18 shows three curves corresponding to using the optional extra dead cone suppression available (MSTA(19)=1), and to the new treatment of heavy masses described in Sec. 2.3.1, in which the full leading order massive matrix element was introduced to describe the branching of the first gluon emitted in $`q\overline{q}`$ events. As can be seen the rate is too low when no optional extra dead cone suppression is used, and then gets even worse when it is used. The new treatment of heavy masses is a clear improvement compared to the old treatment with the dead cone suppression option turned on. But the best behaviour is still achieved when no mass corrections are used.
#### 4.4.2 Four jet rates at $`\sqrt{s}=M_\mathrm{Z}`$
The Pythia comparison in Figure 19 shows three curves corresponding to the three cases described above. As can be seen the recent changes consisting in introducing mass effects at all stages in the treatment also improve the description, as for the three jet case, but still results in a rate which is too low by about 5-7% with respect to data and to the LO matrix element predictions. The same trend is seen as for the three jet case that in the prior versions still presently used by a majority of Lep experiments, the best behaviour is obtained by using massless expressions for the matching procedure of the generated 3 jets (switch MSTJ(47) set to 1).
The Herwig comparison is shown in Figure 20. A description compatible with the data and with the LO matrix element prediction can be seen.
The Ariadne comparison in Figure 21 shows three curves corresponding to using the optional extra dead cone suppression available (MSTA(19)=1), and to the new treatment of heavy masses described in Sec. 2.3.1. As can be seen, contrary to the three jet case, the overall behaviour is quite reasonable when no optional extra dead cone suppression is used, but somewhat too low when it is used. The behaviour with the newly changed treatment of heavy masses does not improve the situation significantly: the best behaviour is still achieved when no mass corrections are used, as in the case of $`R_3^b\mathrm{}`$.
Finally the result of an initial investigation of $`R_4^b\mathrm{}`$ with the new Fourjphact program is shown in Figure 22. The LO ME curves calculated within Fourjphact are identical to those shown in Figure 19, Figure 20 and Figure 21. The results after only the subsequent parton shower are also shown, as calculated starting from the LO ME with 4.6 GeV, and indicate as expected an exageration of the suppression of gluon radiation from the $`b`$ quark mass. Presumably, if a way could be found in this program to start the matching procedure from the LO ME with 2.8 GeV (rather than 4.6 GeV) while preserving the corresponding jet angles, one could perhaps contemplate getting the LO ME + parton shower only results to lie in the middle of the band of uncertainty defined in Sec. 4.3, where the NLO results are expected. The result after both parton shower and hadronisation are also shown, although here it is fair to say that the large effect seen from the hadronisation is not understood and should be studied more.
#### 4.4.3 Four jet rates at $`\sqrt{s}=189`$ GeV
The Pythia 4 jet rate comparison was repeated at $`\sqrt{s}=189`$ GeV. This is shown in Figures 23 and 24. A trend similar to that at $`\sqrt{s}=M_\mathrm{Z}`$ can be seen. For values of the jet resolution parameter in the range $`y=0.0020.008`$, even with the most recent treatment consisting in introducing mass effects at all stages in the handling of the three jets, a residual deficit of about 1-3% with respect to the LO matrix element predictions. Also, in the prior versions still presently used by a majority of Lep experiments, the best behaviour is obtained by using massless expressions for the matching procedure of the generated three jets (switch MSTJ(47) set to 1). If the massive treatment is used in these prior versions is used, a deficit as large as about 2-6% results in the range $`y=0.0020.008`$ (see Figure 23).
The Ariadne 4.08 4 jet-rate comparison was also repeated at $`\sqrt{s}=189`$ GeV. This is shown in Figure 25. The same trend is seen as for the 4-jet rate at Lep1 energies: a quite reasonable behaviour is observed when no optional extra dead cone suppression is used, but the rate is somewhat too low when it is used. The behaviour with the recently changed treatment of heavy masses described Sec. 2.3.1 is also shown. As can been the new treatment does not improve the description at Lep2 energies.
#### 4.4.4 Effects of gluon splitting on $`R_4^b\mathrm{}`$
An additional issue was raised during the meetings of this working group concerning the impact on the evaluation of $`R_4^b\mathrm{}`$ arising from uncertainties in processes involving gluon splittings into $`b\overline{b}`$. In the case of the three-jet rate ratio $`R_3^b\mathrm{}`$, the effect was investigated and found to be small. On the contrary, for four jets, because of the lower rate, and because the two $`b`$ quarks emerging from gluon splittings are often resolved, effects are larger.
The standard definition of $`R_4^b\mathrm{}`$ presented in Sec. 4.2.7 considers primary quarks, and is advantageous from the theoretical point of view, but not from the experimental one, where such a distinction is obviously not straightforward. Since the data are extrapolated to the parton level using a Monte-Carlo to compare with the calculations, a wrong assumption on the gluon splitting into $`b\overline{b}`$ translates directly into a bias. In order to estimate this bias, a new definition was proposed, labeled $`R_4^b\mathrm{}`$(NEW) in which are counted in the numerator any event containing $`b`$ quarks, irrespective whether they originate from primary or secondary production, and at the denominator only events with light quarks, excluding those with a gluon splitting into $`b\overline{b}`$. This new observable is closer to the experimental situation, but is known to carry larger NLO corrections, and is hence more uncertain theoretically. It was nonetheless evaluated both at $`\sqrt{s}=M_\mathrm{Z}`$ and at at $`\sqrt{s}=189`$ GeV, analytically and using Pythia 6.131, with the different settings corresponding to different recent treatments of the gluon splitting process into $`b\overline{b}`$ developed in the framework of the working group (see Sec. 7.3). As an example the normalised difference $`(R_4^b\mathrm{}\text{(NEW)}R_4^b\mathrm{})/R_4^b\mathrm{}`$ is shown for $`\sqrt{s}=M_\mathrm{Z}`$ in Figures 26 and 27, respectively for the present default settings (MSTJ(44)=MSTJ(42)=2), and for one of the proposed set of new settings described in Sec. 7.3 (MSTJ(44))=MSTJ(42)=3), corresponding to a $`gb\overline{b}`$ rate which is roughly doubled. The same comparison is shown for $`\sqrt{s}=189`$ GeV in Figures 28 and 29.
A general feature of these plots is that at small $`y_c`$ Pythia is always lower. This arises because in Pythia, contrary to the analytic calculation, the four-jet rates in both the denominator and numerator of $`R_4^b\mathrm{}`$ receive contributions also from two and three jets, which reduce the relative impact from the fraction of events containing a gluon splitting into $`b\overline{b}`$. Considering firstly the results at $`\sqrt{s}=M_\mathrm{Z}`$, one can see that indeed doubling the gluon splitting rate into $`b\overline{b}`$ in Pythia results in a difference between the two definitions which at large $`y_c`$ becomes similar to that obtained in the analytical calculation. Since doubling the gluon splitting rate into $`b\overline{b}`$ tends to be favoured by both experimental results and theoretical work (see Sec. 7.3), one can in principle take this found consistency as evidence confirming that it indeed needs to be doubled in Pythia. To check this further, the fraction of events containing a $`gb\overline{b}`$ splitting in which the two $`b`$ quarks are clustered into separate jets (case of resolved gluon splittings) was evaluated and found to be largely dominant. This indicated that large NLO corrections in the new definition $`R_4^b\mathrm{}`$(NEW) are not expected to arise from gluon splittings into $`b\overline{b}`$, and that most of the effect does occur at LO. From this can be concluded that a procedure consisting in correcting the measured four-jet events in data to extrapolate to the parton level using a Monte-Carlo induces a sensitivity to the correct rate of gluon splittings into $`b\overline{b}`$, at the level of these discrepancies between the differences. A second conclusion is that indeed at $`\sqrt{s}=M_\mathrm{Z}`$ doubling the rate of gluon splittings into $`b\overline{b}`$ would seem to be justified from the found consistency of the comparison with the analytical results.
The picture changes however at $`\sqrt{s}=189`$ GeV. As can be seen from Figure 29, the Pythia curve with doubled gluon splitting into $`b\overline{b}`$ rate now overshoots significantly at large $`y_c`$ To understand the origin of this behaviour, the same study was performed as at $`\sqrt{s}=M_\mathrm{Z}`$ to evaluate the fraction of events containing a gluon splitting into $`b\overline{b}`$ in which the gluon splitting is resolved, and it was found that at 189 GeV only about half are. Hence in this case the origin of the overshooting could be traced to the fact that the new definition receives large contributions at NLO from gluon splittings into $`b\overline{b}`$. Such a behaviour is in fact expected from the scaling with energy of this last contribution, which grows like $`\mathrm{log}(m_b^2/s)`$.
None of these results prevent one from evaluating the performance of Monte-Carlo programs using the standard definition for $`R_4^b\mathrm{}`$. However one must be careful as soon as one wants to compare with data. Moreover, the results at 189 GeV should not be taken as evidence against a larger rate of gluon splittings into $`b\overline{b}`$ in Pythia, which may be needed as explained in Sec. 7.3, but just that for the case of the observable $`R_4^b\mathrm{}`$(NEW) the comparison is not meaningful, because of the large NLO contributions affecting it. A full calculation at NLO would be helpful to study this further.
### 4.5 Discussion of hadronization corrections to $`R_3^b\mathrm{}`$
In most of the experimental analyses that involve hadronic final states, the data need to be corrected to the parton level in order to be compared with the theoretical predictions. In particular, this is true for the studies considered in this report, for instance the determination of the $`b`$-quark mass, or the studies of the flavour independence of the strong coupling constant and the multi-jet production rates. This procedure necessarily implies unfolding the data for detector and acceptance effects as well as for the hadronization process. This introduces biases and uncertainties which need to be carefully studied and quantified to extract reliable measured values within the quoted errors. The detector and acceptance corrections depend on each particular experiment and consequently will not be discussed here. Only the hadronization correction will be the subject of this section. As shown in references the uncertainty arising from the lack of precise knowledge of how the hadronization process takes place limits the experimental precision of the experimental quantities and QCD tests. Any progress leading to a better understanding of the transition from partons to hadrons or finding new observables with better behaved properties will immediately result in improving these measurements.
In particular, let us consider the hadronization corrections associated with the $`R_3^b\mathrm{}`$ observable introduced in Eq. 34 (though on a qualitative basis the same procedure can be easily applied and generalized to the $`R_4^b\mathrm{}`$ observable of Eq. 35 or other event shape variables). This observable summarizes most of the features commented on in previous sections and at the same time has the advantage of being calculated up to NLO (see section 4.3) and of having relatively small fragmentation corrections (of roughly $``$1%). The fragmentation models considered in this exercise are Pythia and Herwig. The analysis is also performed using two jet clustering algorithms: Durham and Cambridge, and the potential results of using one or the other are compared and discussed. The determination of $`b`$-quark mass or the test of the flavour independence of the strong coupling constant can be derived from this observable and the implications are obvious: a better understanding of $`R_3^b\mathrm{}`$ with a smaller error leads to a more precise determination of the $`b`$-quark mass or a more stringent test of the flavour independence of $`\alpha _\text{S}`$.
In general there is no prescription to unambiguously define the fragmentation correction factor to be applied to the $`R_3^b\mathrm{}`$ observable but the procedure described below can be regarded as a reasonable approach. It is based on the Delphi procedure , though others methods could also be envisaged for this purpose . It seems appropiate then to consider all models which give a good description of the data and calculate the correction factor corresponding to each model. The average of the correction factors obtained is taken as the best estimate of the correction factor and the distribution of these values defines the uncertainty. This also means that the models or generators considered in the analysis should be tuned in order to properly describe the data. The overall fragmentation uncertainty can then be quantified by adding in quadrature the two different source of errors: $`\sigma _{mod}`$, the uncertainty due to the dependence of the hadronization correction factors on the two models considered, Pythia and Herwig, and $`\sigma _{tun}`$, uncertainty due to the possible variation of the main fragmentation parameters in Pythia. Hence, the total uncertainty is expressed as:
$$\sigma _{had}(y_c)=\sqrt{\sigma _{mod}^2(y_c)+\sigma _{tun}^2(y_c)}.$$
(39)
Following the Delphi procedure the $`\sigma _{tun}`$ uncertainty is obtained by varying the most relevant parameters of the string fragmentation model incorporated in Pythia ($`Q_0`$, $`\sigma _q`$, $`ϵ_b`$, a and b) within an interval of $`\pm 2\sigma `$ from their central tuned values and assuming that the individual parameter errors are all independent. Figure 30 shows, for the Cambridge algorithm, the $`\sigma _{tun}`$ uncertainty as a function of the jet resolution parameter as well as the contribution of each individual parameter error. For large enough $`y_c`$ values the overall $`\sigma _{tun}`$ uncertainty is seen to be around 3$`0`$$`/`$$`00`$. The different tuned versions of Jetset or Pythia have also been tested and cross-checked to give the same correction factors to the observable $`R_3^b\mathrm{}`$ within 1$`0`$$`/`$$`00`$.
In this exercise, the two fragmentation schemes considered are Pythia and Herwig, therefore the average of the two correction factors obtained is considered as the fragmentation correction factor to $`R_3^b\mathrm{}`$ and the uncertainty $`\sigma _{mod}`$ is taken to be half of their difference. The generators differ not only in the fragmentation process (cluster fragmentation in Herwig and Lund string fragmentation in Pythia) but also in the way the particle decays are implemented. Therefore $`\sigma _{mod}`$ has two contributions, one from the fragmentation scheme itself, $`\sigma _{modfrag}`$, and the other one from the decay tables used, $`\sigma _{moddec}`$, so it can be written as:
$$\sigma _{mod}(y_c)=\sqrt{\sigma _{modfrag}^2(y_c)+\sigma _{moddec}^2(y_c)}.$$
(40)
We present here the results of the $`\sigma _{mod}`$ uncertainty obtained with Herwig version 5.8 as tuned by Delphi and version 6.1 as tuned by Aleph. For Pythia, the Delphi tuning is used. Presently all Lep experiments are working on the tuning of new versions, therefore new and better sets are expected soon. Figure 31 shows $`\sigma _{mod}`$ and $`\sigma _{modfrag}`$ at different $`y_c`$ values as calculated using the Durham jet clustering algorithm. For $`y_c0.02`$ the contribution of $`\sigma _{moddec}`$ to the total model uncertainty becomes small. Also a better agreement between the latest generator versions of Herwig and Pythia is observed over the entire $`y_c`$ region. This result probably is a consequence of the fact that in the latest versions of both generators the $`b`$-fragmentation functions are now similar. How well these models describe this distribution in data is, however, discussed elsewhere (section 6) and new analyses are still emerging from the various LEP and SLC collaborations on this subject.
Apart from improving the performance of the generators by including more precise calculations and better modelling of the various processes taking part, the search and use of observables having better theoretical or experimental properties is also worthwhile to reduce the total uncertainty. Still in the context of measuring the $`b`$-quark mass or of testing the flavour independence of $`\alpha _\text{S}`$, the use of different algorithms to reconstruct jets has extensively been studied and compared in references and for different event shape variables in references . The main conclusion derived from these studies is that not all observables are equally suited to make the above measurements because they are influenced by different higher order corrections which in some cases can be large. This can explain some of the spread the $`b`$-mass values measured by the various experiments. Therefore studying each observable property in both its theoretical and hadronization aspects is mandatory for making precise measurements.
Following the spirit of this section, Figure 32 presents the size of the total hadronization correction uncertainty ($`\sigma _{had}(y_c)`$) for $`R_3^b\mathrm{}`$ using either the Durham or Cambridge algorithms for the jet reconstruction. The use of the Cambridge algorithm reduces the theoretical error on the $`b`$-quark mass determination though the hadronization error is about the same for $`y_c0.02`$. The use of Durham is however limited to the $`y_c`$-region above 0.015 in order to keep the four jet rate below 5% and the hadronization correction small and flat with respect to $`y_c`$. The same arguments applied to Cambridge allows the extension of the $`y_c`$ region down to 0.004, which, although it increases the sensitivity to the $`b`$-quark mass marginally, does increase the sensitivity to the difference of the LO and the NLO predictions, enabling a better experimental distinction between the $`\overline{MS}`$ and pole mass schemes. The curves shown in Figure 32 have been taken from where Pythia 6.131 and Herwig 5.8 were used and, therefore, promising further reductions in those uncertainties can probably still be obtained with the latest versions of these generators.
### 4.6 Conclusions and remaining issues
In this section effects from the $`b`$ quark mass on three and four-jet rates have been analysed, both in real data, and via analytic calculations and Monte-Carlo approaches. A method for evaluating the relevant observables and for estimating the theoretical uncertainty in the predictions in a consistent way has been presented, and the predictions from several Monte-Carlo programs (Ariadne, Pythia, Herwig and Fourjphact) has been evaluated, including in some cases (Ariadne, Pythia) recently improved versions. Moreover results from on-going studies aiming at controlling the additional uncertainties arising from the modeling of non-perturbative effects in the main Monte-Carlo programs used have been reported.
Here below we first quantify the theoretical uncertainty which is appropriate for each of the observables studied, and then summarize the performance of each Monte-Carlo program. Finally we mention remaining issues relevant to the four-jet case, and improvements which are still needed.
#### 4.6.1 Theoretical uncertainty affecting $`R_3^b\mathrm{}`$ and $`R_4^b\mathrm{}`$
Following the prescription described in Sec. 4.3, we quote the theoretical uncertainty in the $`R_3^b\mathrm{}`$ and $`R_4^b\mathrm{}`$ ratios as $`\pm `$ half the difference between the LO ME results corresponding to the pole mass $`M_b5`$ GeV and to the running mass $`m_b(\mu )`$, with $`m_b(M_\mathrm{Z})3`$ GeV. This definition is also supported by the analysis of the Delphi data at $`\sqrt{s}=M_\mathrm{Z}`$, which for both $`R_3^b\mathrm{}`$ and $`R_4^b\mathrm{}`$ lie within these two calculations. As expected the uncertainties are found to be larger for $`R_4^b\mathrm{}`$ than for $`R_3^b\mathrm{}`$, and are reduced for large values of $`y_c`$ and of $`\sqrt{s}`$.
At $`\sqrt{s}=M_\mathrm{Z}`$, for $`R_3^b\mathrm{}`$, from the difference between the LO ME results the error may be quoted as ranging from $`\pm `$ 1 to 2% for values of $`y_c`$ between 0.02 and 0.1 (see Figure 16). Arguably, this estimate is conservative since in this case an NLO calculation exists. For $`R_4^b\mathrm{}`$, from the difference between the LO ME results the error may be quoted as ranging from $`\pm `$ 2 to 4% for values of $`y_c`$ between 0.012 and 0.03 (see Figure 19).
At $`\sqrt{s}=189`$ GeV, for $`R_4^b\mathrm{}`$, from the difference between the LO ME results the error may be quoted as ranging from $`\pm `$ 1 to 3% for values of $`y_c`$ between 0.002 and 0.04 (see Figure 23).
#### 4.6.2 Performance of the different Monte-Carlo programs
Ariadne underestimates $`R_3^b\mathrm{}`$ at $`\sqrt{s}=M_\mathrm{Z}`$ in all of its versions (see Figure 18). The recently improved one does however provide slightly better values than the version with the dead cone suppression switched on. The best description is nonetheless achieved by switching off all mass treatments altogether. In this case the result is within the band of uncertainty defined by the LO ME curves, but is about 1.5% lower than the NLO curve and the data. With the default setting of the present official version of the program (dead cone suppression turned on) the shift with respect to the NLO ME curve ranges from 2 to 4%. On the other hand $`R_4^b\mathrm{}`$ at $`\sqrt{s}=M_\mathrm{Z}`$ is reasonably described in the version with all mass treatments switched off. For this observable both the new improved version and the old mass treatment of the dead cone give similarly low results, A similar qualitative behaviour is observed at $`\sqrt{s}=189`$ GeV (see Figure 25). For both of the versions of Ariadne with either the default treatment of the $`b`$ mass via the dead cone suppression or with the new improved treatment, the underestimation reaches about 3% for $`y_c0.0020.3`$.
Pythia results for $`R_3^b\mathrm{}`$ at $`\sqrt{s}=M_\mathrm{Z}`$ (see Figure 16) show a strong underestimation for the old version (prior to 6.130) with mass effects turned on as per the default of that version. The bias is about 1 to 4% for $`y_c`$ ranging between 0.01 and 0.12. A better behaviour is obtained by turning off all mass effects. The best description is however obtained thanks to the recently improved treatment of mass effects in versions following 6.131. In this case the MC prediction overlaps nicely with the NLO results and with data. On the other hand for $`R_4^b\mathrm{}`$ at $`\sqrt{s}=M_\mathrm{Z}`$ all versions underestimate the rate (see Fig. 19). The recently improved one (following version 6.131) does nonetheless provide the best description. The bias is in this case about 2.5 to 7% for $`y_c`$ ranging between 0.012 and 0.03. With the old default mass treatment (prior to version 6.130) the bias becomes as large as 10 to 12%. The discrepancies are somewhat less pronounced at $`\sqrt{s}=189`$ GeV (see Figure 23 and 24). In this case also the newest version (following 6.131) is the best, and appears to be about 1 to 3% too low. The old default mass treatment (versions prior to 6.130) is the worst, with the suppression exaggerated by 1.5 to 4% for $`y_c`$ ranging from 0.002 to 0.03.
Herwig was studied for the $`R_3^b\mathrm{}`$ and $`R_4^b\mathrm{}`$ observables only at $`\sqrt{s}=M_\mathrm{Z}`$ (see Figure 17 and 20). A fair agreement is seen in both cases although with some slight underestimation.
The new Fourjphact program with massive four-parton matrix elements matched to the parton shower algorithm of Pythia was investigated as well, and an initial preliminary result was shown. More work is needed here to study whether this approach to matching can provide a solution to the description of $`R_4^b\mathrm{}`$ once the full parton shower and hadronisation treatements are implemented.
#### 4.6.3 Remaining issues and improvements needed
As described above it has been found that all programs tend to exaggerate the suppression which arises from the $`b`$ quark mass, by varying amounts, either in the three-jet rate, or in the four-jet rate, or in both. The best global behaviour is seen for Herwig, although in the version 5.8 of this program which was studied, hadronization corrections were quite a bit larger than for example in Pythia. In the most recent versions (version 5.9 and 6.1), hadronization corrections have become closer to those in Pythia. So from the particular point of view discussed in this section, Herwig would seem to be best. More studies are of course needed to confirm that this is also the case at higher energies. Moreover, additional work towards improving further the remaining discrepancies in Ariadne and Pythia is still on-going at the time of this writing, and will hopefully also bring these two programs in line in the near future. Finally the new Fourjphact program looks quite promising and needs to be investigated further in this context.
On the theoretical side work towards carrying out a massive four-parton matrix element calculation at NLO would enable the estimate of the uncertainty in the $`R_4^b\mathrm{}`$ observable described in this section to be checked explicitly and refined. It could then also be used experimentally to probe the running of the $`b`$-quark mass, as has been done with $`R_3^b\mathrm{}`$.
An additional issue – which would also benefit from the availability of an NLO 4-parton massive calculation – concerns the impact on $`R_4^b\mathrm{}`$ from uncertainties in processes involving gluon splitting into $`b\overline{b}`$. Monte Carlo studies indicate that enhancing the rate of $`gb\overline{b}`$ splitting in Pythia, an option discussed in Sec. 7.3, gives consistency with the LO estimate of the impact at $`\sqrt{s}=M_\mathrm{Z}`$ but overshoots it by up to a factor of two at high $`y_c`$ at Lep2 energies. This could be taken as an estimate of higher-order uncertainties resulting from $`gb\overline{b}`$ splitting processes at high energy.
## 5 STUDY OF FOUR JET OBSERVABLES
### 5.1 Introduction
The study of 4-jet final states is of great interest for Lep1 as well as Lep2 physics analyses, and reliable predictions of the properties of such final states by the various Monte Carlo programs are mandatory. From a QCD standpoint of view, 4-jet final states have their origin in the processes $`\mathrm{Z}q\overline{q}gg`$ and $`\mathrm{Z}q\overline{q}q^{}\overline{q}^{}`$, with the secondary partons coming from double gluon Bremsstrahlung and gluon splitting into gluon or quark pairs.
At Lep1 these processes have been employed for the tests of the structure of the underlying gauge group ( and references therein), which is SU(3) in the case of QCD. In order to get sensitivity to the gauge structure of the theory, a specific class of observables has been employed, namely angular distributions of jets in 4-jet events. The perturbative expansion for the differential distributions of these observables starts at $`𝒪(\alpha _\text{S}^2)`$, and only leading-order (LO) predictions were available until recently. However, now the next-to-leading order (NLO) corrections have been computed -, which allows refined studies of 4-jet observables, such as improved tests of the gauge structure or measurements of the strong coupling constant with variables for which the perturbative predictions start at $`𝒪(\alpha _\text{S}^2)`$, only.
At Lep2 the interest in 4-jet final states is more related to background studies in analysis of fully hadronic W decays and searches, such as the 4-jet channel in the search for the Higgs boson. As an example for variables which enter the selection algorithms of those analyses, the sum of the six interjet angles and the angle between the second and third most energetic jets can be mentioned in case of the W cross section measurements . The variable $`y_{34}`$, which will be explained in the next section, enters in a typical preselection of Higgs searches , and a further background rejection is obtained by looking at functions of interjet angles. Therefore it is clear that a good description of 4-jet observables by the Monte Carlo programs is necessary in order to obtain reliable estimates of the QCD backgrounds.
In the following first the observables are defined which are used for the studies of this section, then the predictions of the various models are compared. These comparisons are first made for the leading order matrix element predictions for four-jet observables, then the effects of next-to-leading order contributions and mass corrections are investigated, and finally the Monte Carlo models are compared to each other and to the data for quantities computed from hadrons instead of partons.
### 5.2 Observables
The observables which will be studied in detail, are described in the following. For those events where exactly four jets are found by the Durham jet clustering algorithm with the $`E`$ recombination scheme and a cut-off value of $`y_{cut}=0.008`$, the energy-ordered jet momenta are used to compute the four-jet angular variables listed below :
* the Bengtsson-Zerwas angle : $`\chi _{BZ}=\mathrm{}[(\text{p}_1\times \text{p}_2),(\text{p}_3\times \text{p}_4)]`$
* the Körner-Schierholz-Willrodt angle :
$`\mathrm{\Phi }_{KSW}=1/2,\{\mathrm{}[(\text{p}_1\times \text{p}_4)(\text{p}_2\times \text{p}_3)]+\mathrm{}[(\text{p}_1\times \text{p}_3),(\text{p}_2\times \text{p}_4)]\}`$
* the (modified) Nachtmann-Reiter angle : $`\theta _{NR}^{}=\mathrm{}[(\text{p}_1\text{p}_2),(\text{p}_3\text{p}_4)]`$
* the angle between the two lowest energy jets : $`\alpha _{34}=\mathrm{}[\text{p}_3,\text{p}_4]`$
These variables have already been used extensively for the measurements of the QCD colour factors because the shape of these distributions is sensitive to the group structure.
For all hadronic events, the following event shape variables have been considered :
* D-parameter $`D`$ , which is defined as the product $`D=27\lambda _1\lambda _2\lambda _3`$, with $`\lambda _i`$ being the three eigenvalues of the infrared safe momentum tensor
$$\mathrm{\Theta }^{ij}=\underset{a}{}\frac{p_a^ip_a^j}{|\text{p}_a|}/\underset{a}{}|\text{p}_a|.$$
(41)
The sum on $`a`$ runs over all the final state particles (partons or hadrons), and $`p_a^i`$ is the $`i`$th component of the three-momentum of the particle $`a`$ in the centre-of-mass system.
* $`y_{34}`$ (Durham, $`E`$ recombination scheme), which is the jet resolution parameter when going from four to three jets, i.e., the event is clustered into jets until only four jets are left, and then $`y_{34}=\mathrm{min}y_{ij}`$, where the minimum is taken over all distance measures (defined by the Durham metric) between the remaining jets.
Since at the end of this section a comparison with corrected data will be given, a short description of a typical data analysis is in place. Hadronic events are selected by requiring a minimum number of charged tracks and a minimum charged energy per event. This reduces backgrounds from $`\tau ^+\tau ^{}`$ and two-photon events to negligible levels.
Then the observables have to be corrected for detector effects such as finite acceptance and resolution. This is done by computing the observables from a Monte Carlo before and after the detector simulation and imposing the same track and event selection cuts as for the data. Then bin-by-bin correction factors are computed for every bin $`i`$ of the distribution,
$$C_i^{det}=\frac{N_i^{had}}{N_i^{det}},$$
(42)
where $`N_i^{had}(N_i^{det})`$ denotes the number of entries in the distribution at the hadron (detector) level. The hadron level distributions are obtained by switching off any photon radiation in the initial and final state (ISR, FSR), with all particles having mean lifetimes less than $`10^9`$ s required to decay, and all other particles being treated as stable. So from a measured distribution $`D_i^{meas}`$ a corrected distribution $`D_i^{corr}`$ is obtained according to
$$D_i^{corr}=C_i^{det}D_i^{meas}.$$
(43)
The detector correction factors are typically found within the 5-10% range, increasing at the edges of the phase space.
These corrected distributions can be compared to the predictions from perturbative QCD, which have to be corrected for hadronization effects, i.e., long-distance non-perturbative effects. This is achieved by computing the relevant observables at parton and at hadron level, which allows to define bin-by-bin correction factors similarly to the detector corrections,
$$C_i^{had}=\frac{N_i^{had}}{N_i^{part}}.$$
(44)
The superscript part refers to the parton level. So from a purely perturbative prediction $`D_i^{pert}`$ a corrected QCD prediction $`D_i^{QCD}`$ is obtained according to
$$D_i^{QCD}=C_i^{had}D_i^{pert},$$
(45)
which is to be compared to the corrected data $`D_i^{corr}`$.
The Monte Carlo simulations which most frequently are employed for the computation of these correction factors, as well as for the simulations of QCD backgrounds to other physics channels, are the parton shower models as implemented in Pythia or Herwig , together with the string fragmentation for the former and a cluster fragmentation in case of the latter. It should be considered that the basic idea of the parton shower is to describe well the structure of jets in two-jet like events, since it is based on a collinear approximation of the matrix elements for gluon radiation off quarks. Because of the matching of the first parton branchings to the exact LO matrix elements, also three-jet like quantities are described rather well. However, it can not really be expected that the parton shower approach gives a good description of four-jet quantities. In fact, rather large discrepancies have been observed in the past .
A different approach can be tested by using the matrix element option in the Pythia program, where at the parton level two-, three- and four-parton final states are generated according to the exact NLO matrix elements, and then the hadronization step is performed via the string fragmentation scheme. This model should give a better description of four-jet related quantities. However, it is known not to describe well the energy evolution of basic quantities such as the charged multiplicity .
Therefore new approaches are tried, based on the idea of matching matrix element calculations to parton shower evolutions, as described in the previous sections. In the following these models will be discussed in detail with respect to their description of 4-jet quantities.
### 5.3 Comparison of model predictions
#### 5.3.1 Leading order predictions (parton level)
There is quite a large variety of programs featuring the production of four jets via QCD at the tree–level. Here, the performance of four of them, namely Herwig, Pythia, Debrecen and the package Apacic++/Amegic++ (denoted as Apacic++ in the following), is compared.
Note that according to the corresponding manuals, the four jet expressions within Pythia and Herwig are for massless partons (apart some mass effects which are built in for Pythia) and they contain only the structures to be found for the exchange of virtual photons . However, the claim is, that the additional terms related to intermediate $`Z`$–bosons have only a minor effect . At least for the observables studied here this claim has been verified.
The focus is on the observables defined in Sec. 5.2, specifically the four jet angles $`\alpha _{34}`$, $`\chi _{BZ}`$, $`\varphi _{KSW}`$ and $`\theta _{NR}`$, and $`y_{34}^D`$, the $`y_c`$–value according to the Durham-scheme, where four–jet events turn to three resolvable jets. All results shown and discussed here are on the primary parton level, i.e., results obtained by the appropriate matrix elements squared, and at a centre-of-mass energy of $`91.2`$ GeV, with the argument of $`\alpha _\text{S}`$ kept fixed.
The Monte–Carlo points were produced adopting the following strategy :
1. For Pythia a sample of four–jet events was generated with $`y_{34}^Jy_{cut}^J=0.008`$. This is due to the fact that in Pythia only the Jade scheme is available. Over a large region of phase space, as a rule of thumb, $`y_{cut}^J4y_{cut}^D`$ for the same kinematical configurations.
2. Out of this first sample, only events with $`y_{34}^Dy_{cut}^D=0.004`$ have been selected. For the other three generators, Herwig, Debrecen and Apacic++ the events were directly generated in the Durham-scheme with $`y_{cut}^D=0.004`$.
3. For the four jet angles, jets were defined according to $`y_{cut}^D=0.008`$, thus reducing the sample of step 2 by roughly $`50\%`$. For the $`y_{34}`$–distribution no additional cuts have been applied.
The resulting distributions of $`\mathrm{cos}\alpha _{34}`$, $`|\mathrm{cos}\chi _{BZ}|`$, $`\mathrm{cos}\varphi _{KSW}`$ and $`|\mathrm{cos}\theta _{NR}|`$ can be found in Figure 33. Here, the upper plots exhibit the total number of events per corresponding bin normalized to the total number of events with $`y_{cut}^D=0.008`$, and in the lower plots the relative deviations from the Apacic++ results are displayed. With the exception of the last bin in $`\mathrm{cos}\alpha _{34}`$, the relative (statistical) errors on each distribution are of the size of the symbols. Clearly, the results show a satisfying coincidence with no sizeable relative deviations.
This situation changes when considering the $`y_{34}^D`$–distribution, see Figure 34. Again, the upper plot shows the normalized number of events per bin, and the lower plots depicts the relative deviations from the Apacic++ results. Here, the spread of the statistical errors covers a region from barely visible in the left bins up to three times the size of the symbols in the right bins. However, the deviations of the generators from each other are larger than their individual relative errors and reach up to $`15\%`$. Seemingly, Herwig, Debrecen and Apacic++ coincide. The results obtained by Pythia are somewhat – $`𝒪(15\%)`$ – higher, with the first bin as the only significant exception. Here, Pythia is well below ($`25\%`$) the other generators. However, it should be noted here, that this is probably due to the way the Pythia sample was produced. Since for the production of the Pythia sample in the first step the intrinsic Jade scheme was employed, deviations can be expected especially in the regions where the phase space is cut, i.e., for low $`y_{34}`$. Normalising in the region $`y_{34}>0.01`$, for example, would remove the discrepancy. Turning to the $`D`$–parameter, the different generators agree very well with each other. The relative errors reach roughly the size of the symbols for $`D0.2`$ and are of the order of $`10\%`$ in the last bin. Note that for the $`D`$–parameter as well as for the four jet angles, any difference seen in the $`y_{34}`$–distribution is washed out.
Herwig and Apacic++ provide additional options to supplement the pure matrix elements with running $`\alpha _\text{S}`$ instead of the fixed one with a scale depending on the specific kinematical situation (Herwig and Apacic++) or with some appropriate Sudakov weights (Apacic++), which depend on the flavours and the kinematics of the individual event. These two options are meant to model some aspects of higher order corrections to the pure LO matrix elements and result basically in a shift of events from regions with large $`y_{34}`$ to region with small $`y_{34}`$, see Figure 36. Here, the upper plot shows the number of events per bin in the $`y_{34}`$–distribution normalized to the total number of events and in the lower plot the ratio of the numbers per bin in the uncorrected and the corrected versions of the generators is depicted. The full and empty triangles correspond to Herwig without and with the running $`\alpha _\text{S}`$ option, the diamonds refer to Apacic++ without and with the Sudakov weights (“NLL”), respectively. Obviously, these options “soften” the $`y_{34}`$–distribution of the samples. On the other hand, their effect on the angular distributions is only minor in most of the phase space, see Figure 37. In the four plots the ratios of the corrected (corr) versus the uncorrected (uncorr) options for the four jet angles are displayed. It can be read off, that over the dominant region of phase space available, the inclusion of these corrections does not alter the angular distributions significantly. Rather, their effect is of the order of roughly $`5\%`$ with the only exception of the last bins for small $`\alpha _{34}`$, where the additional weights induce a drastical decrease of up to $`15\%`$. Note, however, that this region is strongly disfavoured, see the corresponding plot in Figure 33, thus, there are comparably large errors on the results.
#### 5.3.2 Next-to-leading order corrections
For four-jet observables large differences were observed between LO order perturbative predictions and corrected experimental data. For instance, the LO perturbative result for the mean value of the D parameter is $`D^{LO}=0.0216`$, while the experimental value is $`0.0618\pm 0.0024`$ , and even after including the large hadronization corrections, a significant discrepancy remains. Such a discrepancy and the large renormalization scale dependence of the LO preturbative result shows clearly that if QCD is to work for four-jet observables, then there have to be large higher-order or non-perturbative corrections. On the other hand, the normalized angular distributions at LO are independent of the strong coupling, and therefore were expected to be insensitive to the renormalization scale, indicating small higher order corrections. For the clarification of the situation the NLO calculations were indispensable. During the second phase of Lep, four parton level NLO programs using different regularization methods were published, Menlo parc , Debrecen , Eerad2 and Mercutio . The one-loop matrix elements implemented in these codes were also different: Eerad2 uses the one-loop matrix elements of Ref. , while the other three programs implemented the one-loop matrix elements published in Ref. . The results of the four programs were compared for distributions of many four-jet observables and very good agreement was found as exemplified here in Table 4, where the NLO four-jet fractions at three different $`y_{\mathrm{cut}}`$ values for the Durham clustering algorithm are compared. In the following, we shall present results obtained using the Debrecen code. This program is the only one of the four that gives the results in a colour decomposed form , which is useful for colour charge measurements.
The general form of the NLO differential cross section for a four-jet observable $`O_4`$ (for instance, D parameter, $`O_4=D`$, or Bengtsson-Zerwas angle, $`O_4=\chi _{BZ}`$) is given by the following equation:
$$\frac{1}{\sigma _0}\frac{d\sigma }{\mathrm{d}O_4}(O_4)=\eta (\mu )^2B_{O_4}(O_4)+\eta (\mu )^3\left[C_{O_4}(O_4)+B_{O_4}(O_4)\beta _0\mathrm{ln}(x_\mu ^2)\right],$$
(46)
where $`\sigma _0`$ denotes the Born cross section for the process $`e^+e^{}\overline{q}q`$, $`\eta (\mu )=\alpha _\text{S}(\mu )C_F/(2\pi )`$, $`x_\mu `$ is the ratio of the renormalization scale to the total centre-of-mass energy, and $`B_{O_4}(O_4)`$, $`C_{O_4}(O_4)`$ are the perturbatively calculable coefficient functions in the Born approximation and the radiative correction, respectively, which are independent of the renormalization scale.
We have performed a high statistics calculation of the $`B`$ and $`C`$ functions for the two event shape variables and for the four angular correlations defined in Sec. 5.2. Figure 38 shows the LO and NLO perturbative cross sections for the two event shapes. We observe that the inclusion of the radiative corrections increases the overall event rate substantially (by 70–130%). For instance, the mean value of the D parameter at NLO is $`D^{NLO}=0.0383(2)`$, which is 77% larger than the LO value, but still far from the measured result. The NLO predictions exhibit significant renormalization scale dependence indicating the importance of even higher orders, which are not likely to be known in the foreseeable future. Thus, for event shapes the perturbative description remains unsatisfactory unless one attempts to use an optimized scale choice . The only possible exception is the four-jet rate for the Durham algorithm, where the relative size of the NLO correction is around 60%, and the resummation of large logarithms exists . Indeed, after matching the next-to-leading logarithmic and NLO results one finds a small renormalization scale dependence and a remarkably good description of the corrected experimental data .
The perturbative result is much more convincing in the case of the angular correlations, although the NLO calculations have brought some surprises, too. In Figure 39 we plot the LO and NLO perturbative predictions for the distributions of the Bengtsson-Zerwas and modified Nachtmann-Reiter angles. We see that the shapes of the distributions change very little when going from LO to NLO. However, when these predictions are used for measuring the QCD color factors, then, at NLO, one uses a quadratic form of the color charge ratios instead of the linear form used in a LO analysis . The different functional form of the fitted function leads to different fitted parameters, even if the shape of the distribution is only slightly changed. In particular, the coefficient of the $`T_R/C_F`$ colour factor ratio receives a very large negative contribution leading to a significant shift in the measured value of this ratio if the NLO prediction is used instead of the LO one .
#### 5.3.3 Mass corrections
For the investigation of mass effects Fourjphact and the package Apacic++/Amegic++ were used with the mass parameters of Table 5. Note that Apacic++ adds the value of a cutoff to the mass parameters, see Sec. 2.4.5.
In general, the inclusion of masses has only a minor effect on the angular distributions. This result can be read off Figure 40. In the upper plots the massive distributions of Fourjphact and Apacic++ are confronted with the massless result of Apacic++, and in the corresponding lower plots the ratios massive/massless are depicted. In most of the bins the effect of masses is of the order of $`2\%`$.
This result is not too surprising, however, since only the $`b`$–quark mass induces any sizeable effect. In Figure 41, the results of different channels as given by Apacic++ are compared, namely $`b\overline{b}b\overline{b}`$, $`b\overline{b}gg`$, $`d\overline{d}d\overline{d}`$, and $`d\overline{d}gg`$. Again, in the upper plots the appropriately normalized number of events per bin is displayed for each channel, and the lower plots depict the corresponding ratios $`(b\overline{b}b\overline{b})/(d\overline{d}d\overline{d})`$ and $`(b\overline{b}gg)/(d\overline{d}gg)`$. Closer inspection of this figure reveals that the by far dominant four jet $`b`$ channel, namely $`b\overline{b}gg`$, is affected on the level well below $`10\%`$ in most of the phase space by mass effects. On the other hand, large effects, best seen in the comparison of $`b\overline{b}b\overline{b}`$ versus $`d\overline{d}d\overline{d}`$, are suppressed by the small relative rates of the corresponding channels.
In the $`y_{34}`$–distribution, mass effects are completely negligible, see Figure 42. In the upper plot, the $`y_{34}`$ distribution as given by Apacic++ with and without the inclusion of masses are depicted. The lower plot shows the ratios of the massive and the massless distributions. As expected, the inclusion of masses results in a slight shift from relatively soft (small $`y_{34}`$) to comparably hard (large $`y_{34}`$) events, since masses shield the collinear regions of particle emission. However, this effect is only of minor size, and in most bins the ratios massive/massless are quite close to 1. Large deviations can be seen in bins with limited statistics only.
#### 5.3.4 Comparison of shower models
In this section we want to compare the predictions of various parton shower models. When working with the new MC options which allow for generating parton showers starting from a four-parton configuration, care has to be taken for particular aspects of these models. In order to avoid singularities in the four-parton generation according to the LO matrix elements, some intrinsic cut-off has to be applied for these programs, for example a Durham resolution criterion in case of Herwig. Therefore the resolution criterion with which jets are selected at the analysis level has to be larger than this intrinsic one. In case of Herwig, the intrinsic cut-off chosen is $`y_{cut}^{intr}=0.004`$, whereas jets are selected with $`y_{cut}=0.008`$. The analysis level can be the final state after the parton shower or after the hadronization step.
Attention has also to be paid to the edges of certain phase space regions, such as the end points in the angular distributions. These regions could be sensitive to the details in the implementations of the models, such as the handling of the masses and virtualities of the partons from which the parton showers start.
The hadronization parameters for the various models are taken from the tuned parameter sets as used by Aleph. No study with respect to variations in these parameters has been performed.
A study of the changes of the shapes for the angular distributions when going from parton to hadron level is outlined below in Sec. 5.4 for the case of Herwig. Here we concentrate on the comparison of the model predictions at hadron level, i.e., after parton showering and hadronization. Also shown is a comparison to data by Aleph. These data are preliminary, with statistical errors only, since the main purpose is to give a qualitative benchmark for the MC predictions. The data have been corrected for detector acceptance and resolution effects by means of bin-to-bin correction factors, as described in Sec. 5.2. The Monte Carlo program employed for this purpose is based on a standard parton shower approach, starting from a $`q\overline{q}`$ pair.
In case of the angular distributions, the normalization is with respect to the total number of four-jet events found, in order to concentrate on the shape. The event shape distributions are normalized to the total number of hadronic events generated.
From Figure 43 the following observations are made : The new option in Pythia 6.1, which interfaces a four-parton event to a parton shower, gives generally a very good description of the angular distributions, whereas the standard parton shower option shows deviations. The two other four-jet MC programs differ from the Pythia four-jet option by about 5-10%, with larger discrepancies seen only at the high end of the $`\mathrm{cos}\alpha _{34}`$ distribution, which is sensitive to mass effects and details of the implementation of the interface, since there two soft jets at close angles are probed. Quark mass effects, which are implemented in Fourjphact, do not have a sizeable impact on the shape of the distributions for the measured sample, which is a normal flavour mixture. However, when studying event samples enriched in heavy quarks, the mass effects should definitely be taken care of.
In Figure 44 similar comparisons are shown for the event shape distributions $`y_{34}`$ and D parameter. Here a rather different picture arises. The four-jet MCs definitely fail to describe these distributions, whereas the standard parton shower approach gives good predictions apart from the very tails. This can be understood from the fact that the low end of the event shape distributions is rather sensitive to contributions from two- and three-jet events, which vanish at leading order in perturbative QCD, but can become sizeable after the parton shower simulation and after hadronization. These contributions are not implemented in the four-jet MC programs. Furthermore, the predictions for the low regions of the event shapes depend strongly on the intrinsic four-jet resolution parameter.
Therefore from these observations we conclude that the new four-jet MC options are well suited for the description of the shape of angular distributions in four-jet events, with remaining uncertainties of the order of 5%, if critical phase space regions are avoided. However, they cannot be used for observables which are sensitive to contributions from two- and three-jet events, which is the case for event shape distributions. Also the relative jet rates cannot be predicted correctly by these programs.
### 5.4 Hadronization corrections
A preliminary study of hadronization effects on the four-jet angular distributions was made using the Herwig four-jet matrix element \+ parton shower option, described in Sec. 2.2.4.
Figures 45 and 46 show the distributions obtained at $`s=M_\mathrm{Z}^2`$ for light primary quarks (IPROC = 601) and primary $`b`$-quarks (IPROC = 605), respectively. The Durham jet metric was used with matrix-element cutoff Y4JT = $`0.004`$ and $`y_{cut}=0.008`$ for the actual resolution of jets. The dotted histograms show the matrix element distributions, with the argument of $`\alpha _\text{S}`$ running as explained in Sec. 2.2.4, while the dashed and solid ones are the reconstructed jet distributions at the parton level (after showering) and hadron level (after decays), respectively. In each case the distributions are normalized to the number of 4-jet events found (with $`y_{cut}=0.008`$) at the relevant level.
One sees that hadronization effects on the shapes of the 4-jet angular distributions are generally not large. The hadron/parton level ratios are, within the limited statistics of the present study, broadly similar to those obtained using the Herwig 2+3 jet ME+PS option (IPROC = 101,105). There are, however, indications that hadronization effects may be overestimated by the 2+3 jet ME+PS option at the 5-10% level in certain places, e.g. the central regions of $`\theta _{NR}`$ and $`\chi _{BZ}`$. Thus the use of correction factors obtained from the 2+3 jet ME+PS option needs to be treated with caution if precision better than 10% is required.
Close study of Figures 45 and 46 reveals the rather surprising fact that the hadron-level results are often closer than the parton-level ones to the (massless) matrix-element distributions. This is particularly true for primary $`b`$-quarks, suggesting that in Herwig the kinematic effects of the $`b`$-quark mass during parton showering are largely cancelled by the effects of hadronization and B-hadron decays.
Thus first results from the Herwig 4-jet ME+PS option suggest that hadronization effects are similar to those obtained using the old 2+3 jet option at the 10% level, and that effects of parton showering and hadronization tend to cancel in the 4-jet angular distributions. This should be investigated further as a function of c.m. energy and jet resolution. Further studies of 4-jet hadronization using the Pythia generator as well as the new combined 2,3 and 4 jet ME+PS Herwig option (see Sec. 2.2.5) are also needed.
### 5.5 Conclusions
A study of the description of 4-jet final states by various Monte Carlo models has been presented. Particular emphasis has been put on the comparison of new Monte Carlo generators, which produce 4-parton final states according to the leading order QCD matrix elements, and then add a parton shower and hadronization. In general good agreement between the different generators has been found. These new models give a better description of 4-jet angular distributions than the standard parton shower models, where the parton shower starts from a quark-antiquark initial state, only. However, jet rates as well as event shape distributions sensitive to 4-jet production, such as $`y_{34}`$, cannot be described by these models, since here also the contributions from 2- and 3-jet events are important. These observations should be taken into account when using the 4-jet MCs for background studies.
Quark mass effects are small for the distributions under consideration, for samples with normal flavour mixture. NLO contributions have minor impact on the shape of 4-jet angular distributions, but they change considerably the 4-jet rate and event shape distributions such as the D parameter. Possible large non-perturbative power corrections to observables such as the mean value of the D parameter have not been studied here.
A preliminary study of hadronization effects on 4-jet angular distributions, using the 4-jet ME+PS option of the Herwig generator only, showed with limited statistics that such effects are broadly similar to those seen with the 2+3 jet ME+PS option, although differences at the 5-10% level are apparent in certain configurations.
## 6 B QUARK FRAGMENTATION FUNCTION
Heavy quark fragmentation functions are a powerful tool in testing the predictivity of perturbative QCD (pQCD), since effects of non-perturbative origin are much more limited in size than in the light-flavour case. At the origin of this behaviour lies the fact that the mass $`m`$ of the heavy quark is much larger than the QCD scale $`\mathrm{\Lambda }`$.
Indeed, on one side the large mass acts as an infrared cutoff for the mass singularities which would appear in the perturbative calculation, ensuring a finite result. The energy distribution of the $`b`$ quark prior to hadronization can therefore be calculated perturbatively. On the other side, hadronization effects have to be phenomenologically modelled, but happen to be small: a heavy quark only loses a momentum fraction of order $`\mathrm{\Lambda }/m`$ when binding with a light one to form a heavy-light meson .
### 6.1 Experimental results
Results for the normalized energy distribution of $`B`$ hadrons, i.e.
$$D(Q,x_E)\frac{1}{\sigma }\frac{d\sigma }{dx_E},$$
(47)
in $`e^+e^{}`$ collisions are given by the Lep collaborations and by the Sld experiment at SLC, at $`Q=M_\mathrm{Z}`$. The scaling variable $`x_E`$ is given by the ratio of the observed $`B`$ particle energy to the beam energy $`E_{beam}`$.
A typical observable measured by experiments is the mean scaled energy fraction $`x_E`$. Table 6 shows some of the most recent determinations of this quantity. In this table the second column identifies the kind of $`B`$ particle observed in the final state, be it the “leading” (also called “primary”) (L) or the “weakly decaying” one (wd). Of course, the average energy of the latter is lower than that of the former, since it has undergone further decaying processes. It should also be noted that the precise details of what the observed final state actually is will at least slightly vary from experiment to experiment. The numbers quoted in the table under the same label “wd” are therefore not exactly comparable, though probably homogeneous enough within the experimental uncertainties so that one can average them. Needless to say, it would be useful if all analyses at some point finally agreed on a single definition for this final state.
The most recent analyses also report fairly accurate data for the full fragmentation function eq. (47), with $`x_E`$ ranging from near zero to one. As expected, these distributions peak very close to one, around $`x_E`$ 0.8-0.9.
### 6.2 Theoretical predictions
The challenge for the theoretical calculations is of course to reproduce not only the mean scaled energy but also, as far as possible, the full fragmentation distribution. A certain degree of phenomenological modelling will be necessary, as perturbative calculations cannot of course describe the hadronization of the $`b`$ quark into $`B`$ mesons and/or baryons. The full fragmentation function is therefore usually described in terms of a convolution between a calculable perturbative component and a phenomenological one:
$$D(Q,x;m,\mathrm{\Lambda };ϵ_1,\mathrm{},ϵ_n)=D^{pert}(Q,x;m,\mathrm{\Lambda })D^{np}(x;ϵ_1,\mathrm{},ϵ_n).$$
(48)
In this equation the perturbative component depends on the centre-of-mass energy $`Q`$, the QCD coupling and the heavy quark pole mass $`m`$, while the non-perturbative one is assumed to depend only on some given set of phenomenological parameters $`(ϵ_1,\mathrm{},ϵ_n)`$, to be determined by fitting the experimental data.
The perturbative component can be either calculated by analytical means or extracted from Monte Carlo simulations of the emission of radiation by the fragmenting heavy quark. In the latter case the theoretical accuracy will of course be lower.
As far as fixed order analytical calculations are concerned, today’s state of the art is the work of ref. . It is accurate up to order $`\alpha _\text{S}^2`$ and also fully includes finite mass terms of the form $`(m/Q)^p`$ with $`p1`$.
Fixed order results do however display two classes of large logarithmic terms: collinear logs, of the form $`\mathrm{log}(Q^2/m^2)`$, and Sudakov logs, $`\mathrm{log}(1x)`$. These terms become large when the centre-of-mass energy is much larger than the heavy quark mass, a fact certainly true at Lep and Sld, and at the $`x1`$ endpoint respectively. All-order resummations for such terms, very important for producing a reliable result, have and are being considered , and are now available at next-to-leading log (NLL) accuracy for both classes of logarithms. Ref. also provides a merging between the fixed order calculation and the collinear-resummed one.
Once pQCD has produced a reliable prediction for the $`b`$ quark fragmentation, one has to “dress” it with some phenomenological modelling in order to describe the observed $`B`$ hadrons distribution, as discussed before. It is important to realize that, since only the convolution of the two factors in Eq. 48 has physical meaning, the parameters fitted in the non-perturbative part will strictly depend on the kind of description adopted for the perturbative term. Different descriptions and/or different parameters in $`D^{pert}(Q,x;m,\mathrm{\Lambda })`$ will lead to different values for the fitted $`(ϵ_1,\mathrm{},ϵ_n)`$ set. Once fitted, such a set will therefore not be usable in conjunction with perturbative descriptions other than the one it has been fitted with.
As an example, Table 7 shows the average scaled energy value of the $`b`$ quark after fragmentation, $`x_E_{pert}`$, as predicted by the perturbative term only, for different values of $`\mathrm{\Lambda }`$ and $`m`$. The calculation used in this example resums to NLL accuracy both collinear and Sudakov logs, but neglects the non-logarithmic finite mass terms. One can clearly see that, on one hand, the purely perturbative predictions are too large to directly describe the experimental results, unless probably unrealistic value for $`\mathrm{\Lambda }_5`$ and $`m_b`$ are used. On the other hand, each different $`(m_b,\mathrm{\Lambda }_5)`$ choice will of course imply a different fitted value for the non-perturbative parameters if the same measured $`x_E`$ is to be correctly described.
While the calculation of the perturbative component does of course follow the rules dictated by pQCD, much more freedom is instead available when choosing the form of the non-perturbative contribution $`D^{np}(x)`$. A lot of different parametrizations have been employed and can be found in the literature: some of them possibly more physically motivated, some chosen only because of practical advantages like an easy Mellin transform or enough parameters to ensure that the shape of the data can be properly described. A list, probably incomplete, of these functional forms can be found in , though it should be noted that some of them appear now to be disfavoured by comparisons with experimental data.
As an example of such a comparison, let us consider one of the best-sellers of these non-perturbative forms, namely the Peterson et al. fragmentation function. Such a function is meant to be one of the physically motivated ones, and has the attractive feature of depending on only one phenomenological parameter $`ϵ`$, which can moreover be roughly related to more fundamental quantities via the relation $`ϵ\mathrm{\Lambda }^2/m^2`$. This function reads
$$D^{np}(x;ϵ)=N(ϵ)\frac{1}{x}\left(1\frac{1}{x}\frac{ϵ}{1x}\right)^2$$
(49)
where $`N(ϵ)`$ is the normalization factor.
The plot in Figure 47 shows the result of attempting a point-by-point fit of the Peterson form, convoluted with the same perturbative calculation used in Table 7, to the latest data from Sld. It can be clearly appreciated how the Peterson model, coupled to this perturbative description, does not seem to offer a valid description of the data. The same conclusion was reached by the Sld Collaboration by coupling this model to the Jetset Monte Carlo description of the perturbative component.
It should however be noted that Eq. 49 was derived under the assumption of describing the hadronization of a heavy quark into a heavy-light meson by picking up a light quark from the vacuum. No attempt was made to include the description of the subsequent decays transforming the leading $`B`$ particle into the weakly decaying ones, which are the ones observed here. Such decays will modify the shape of the fragmentation function, and might at least partially explain the observed discrepancy.
Figure 47 also shows a fit performed with a different non-perturbative function, namely
$$D^{np}(x;\alpha ,\beta )=N(\alpha ,\beta )(1x)^\alpha x^\beta .$$
(50)
This particular form has no immediate physical origin, but it is often used because it has a very simple Mellin transform and is flexible enough to describe the data well. One can indeed see from the plot that it allows for a very good fit of the experimental data.
### 6.3 Monte Carlo predictions
An important issue that remains to be addressed is the performance of the main Monte Carlo event generators in comparison with the latest data on $`b`$-quark fragmentation. Figure 48 shows the results of combining Jetset (version 7.4) parton showers with various models for $`b`$-quark fragmentation into a weakly-decaying B-hadron, compared to the recent Sld data . Jetset plus Lund fragmentation gives a good description of the data whereas, as was the case in the analytical calculations discussed above, the Peterson model does not.
The prediction from Herwig (version 5.7) is seen to be too soft in comparison to the Sld data. As already remarked in Sec. 2.2.3, a harder B-hadron spectrum can be obtained in Herwig 6.1 by varying the $`b`$-quark fragmentation parameters separately. However, detailed tuning of version 6.1 to these data has not yet been attempted.
### 6.4 Concluding remarks
Accurate theoretical preditions exist for the perturbative part of the heavy quark fragmentation function. Collinear and Sudakov logarithms can be resummed to next-to-leading accuracy, and the finite mass terms are known up to order $`\alpha _\text{S}^2`$. All the various contributions can be merged into a single result.
On top of this perturbative result a non-perturbative contribution will always have to be included, and precise experimental data can help identifying the proper shape for such a function.
Predictions for heavy quark fragmentation from the latest versions of Monte Carlo generators have yet to be compared and tuned to the most recent data.
## 7 GLUON SPLITTING INTO BOTTOM QUARKS
### 7.1 Experimental data
The current experimental results on the quantity
$$g_{bb}=\frac{\mathrm{\Gamma }(Z^0q\overline{q}g,gb\overline{b})}{\mathrm{\Gamma }(Z^0\text{hadrons})}$$
are summarized in Table 8. The first three results are based on extracting a secondary $`b\overline{b}`$ signal from a 4-jet sample, whereas the latest Delphi value is obtained from a measurement of the 4$`b`$ rate $`\mathrm{\Gamma }(Z^0b\overline{b}b\overline{b})`$, which should be less model-dependent and give better phase-space coverage.
### 7.2 Analytical predictions
References give leading-order \[$`𝒪(\alpha _\text{S}^2)`$\] predictions as well as the results of resummation of leading \[$`\alpha _\text{S}^n\mathrm{log}^{2n1}(s/m_b^2)`$\] and next-to-leading \[$`\alpha _\text{S}^n\mathrm{log}^{2n2}(s/m_b^2)`$\] logarithms (NLL) to all orders in $`\alpha _\text{S}`$, matched with leading order. The results, for $`\alpha _\text{S}=0.118`$, $`s=M_\mathrm{Z}^2`$ and $`m_b=5.0`$ GeV, are summarized in Table 9.
We see that resummation of logarithms gives a substantial enhancement. Bearing in mind that $`\mathrm{log}(M_\mathrm{Z}^2/m_b^2)6`$ and $`\alpha _\text{S}\mathrm{log}^2(M_\mathrm{Z}^2/m_b^2)4`$, this is not surprising. Still missing are higher-order terms of the form $`\alpha _\text{S}^n\mathrm{log}^{2nm}(s/m_b^2)`$ with $`n,m>2`$, which could well become comparable with the leading term at high orders. The difference between the predictions of and is due to their different treatment of NNLL ($`m=3`$) terms.
In conclusion, the theoretical prediction for $`g_{bb}`$ must still be regarded as quite uncertain, and not in serious disagreement with the data.
Reference also contains predictions of secondary heavy quark production as a function of an event shape variable (the heavy jet mass). There are no data available yet on this, so the authors compare with Monte Carlo results. Their predictions are similar to those of the main event generators discussed below.
Monte Carlo predictions of $`g_{bb}`$ are in principle even more unreliable than the theoretical results presented above, since they do not fully include next-to-leading logarithms or matching to fixed order. Nevertheless they do include some real effects absent from the analytical calculations, such as the effects of phase-space limitations. Different options for the treatment of subleading terms, such as the choice of argument for $`\alpha _\text{S}`$, can easily be explored by providing suitable switches in the programs. Also of course they provide a complete model of the final state, which allows the effects of experimental cuts to be simulated. Relevant developments in the main Monte Carlo programs are described in the following three subsections.
### 7.3 Monte Carlo developments: PYTHIA
#### 7.3.1 Strong coupling argument and kinematics
The default behaviour in Pythia is to let $`\alpha _\text{S}`$ have $`p_T^2`$ as argument. Actually, since the exact kinematics has not yet been reconstructed when $`\alpha _\text{S}`$ is needed, the squared transverse momentum is represented by the approximate expression $`z(1z)m^2`$, where $`z`$ is the longitudinal splitting variable and $`m`$ the mass of the branching parton. Since $`\alpha _\text{S}`$ blows up when its argument approaches $`\mathrm{\Lambda }_{QCD}`$, this translates into a requirement on $`p_T^2`$ or on $`z`$ and $`m`$, restricting allowed emissions to $`p_T>Q_0/2`$, where $`Q_0`$ is the shower cutoff scale. Also when full kinematics is reconstructed, this is reflected in a suppression of branchings with small $`p_T`$. Therefore, if the angular distribution of the $`g`$ decay is plotted in its rest frame, the quarks do not come out with the $`1+\mathrm{cos}^2\theta `$ angular distribution one might expect, but rather something peaked at $`90^\text{o}`$ and dying out at $`0^\text{o}`$ and $`180^\text{o}`$.
For $`gq\overline{q}`$ branchings, the soft-gluon results that lead to the choice of $`p_T^2`$ as scale are no longer compelling, however. One could instead use some other scale that does not depend on $`z`$ but only on $`m`$. A reasonable, but not unique, choice is to use $`m^2/4`$, where the factor 4 ensures continuity with $`p_T^2`$ for $`z=1/2`$. This possibility has been added as new option MSTJ(44)=3. In order for this new option to be fully helpful, a few details in the treatment of the kinematics have also been changed for the $`gq\overline{q}`$ branchings. These changes are not completely unimportant, but small on the scale of the other effects discussed here.
Actually, the change of $`\alpha _\text{S}`$ argument in itself leads to a reduced $`gq\overline{q}`$ splitting rate, while the removal of the $`p_T>Q_0/2`$ requirement increases it. The net result is an essentially unchanged rate, actually decreased by about 10% for charm and maybe 20% for bottom, based on not overwhelming statistics. The kinematics of the events is changed, so experimental consequences would have to be better quantified. However, the changes are not as big as might have been expected – see the following.
#### 7.3.2 Coherence
In the above subsection, it appears as if the $`1+\mathrm{cos}^2\theta `$ distribution would be recovered in the new option MSTJ(44)=3. However, this neglects the coherence condition, which is imposed as a requirement in the shower that successive opening angles in branchings become smaller. Such a condition actually disfavours branchings with $`z`$ close to 0 or 1, since the opening angle becomes large in this limit. It should be noted that the opening angle discussed here is not the true one, but the one based on approximate kinematics, including neglect of masses. One may question whether the coherence arguments are really watertight for these branchings, especially if one considers $`gq\overline{q}`$ close to threshold, where the actual kinematics is quite different from the one assumed in the massless limit used in the normal coherence derivation.
As a means to exploring consequences, two new coherence level options MSTJ(42)=3 and 4 have thus been introduced. In the first, the $`p_T^2`$ of a $`gq\overline{q}`$ branching is reduced by the correct mass-dependent term, $`14m_q^2/m_g^2`$, while the massless approximation is kept for the longitudinal momentum. This is fully within the uncertainty of the game, and no less reasonable than the default MSTJ(42)=2. In the second, no angular ordering at all is imposed on $`gq\overline{q}`$ branchings. This is certainly an extreme scenario, and should be used with caution. However, it is still interesting to see what it leads to.
It turns out that the decay angle distribution of the gluon is much more distorted by the coherence than by the $`\alpha _\text{S}`$ and kinematics considerations described earlier. Both modifications are required if one would like to have a $`1+\mathrm{cos}^2\theta `$ shape, however. Also other distributions, like gluon mass and energy, are affected by the choice of options.
The most dramatic effect appears in the total gluon branching rate, however. Already the introduction of the mass-dependent factor in the angular ordering requirement can boost the $`gb\overline{b}`$ rate by about a factor of two. The effects are even bigger without any angular ordering constraints at all. It is difficult to know what to make of these big effects. The options described here would not have been explored had it not been for the Lep data that seem to indicate a very high secondary charm and bottom production rate. Experimental information on the angular distribution of secondary $`c\overline{c}`$ pairs might help understand what is going on better, but probably that is not possible experimentally.
#### 7.3.3 Summary
In order to study uncertainties in the $`gb\overline{b}`$ rate, some new Pythia options have been introduced, MSTJ(44)=3 and MSTJ(42)=3 and 4, none of them as default (yet). Taken together, they can raise the $`gb\overline{b}`$ rate by a significant factor, as summarized in Table 10. A study of the effects of these options on the 3-jet rate ratio $`R_4^b\mathrm{}`$ is described in Sec. 4.4.4.
The MSTJ(42)=4 option is clearly extreme, and to be used with caution, whereas the others are within the (considerable) range of uncertainty.
The corrections and new options are available starting with Pythia 6.130, obtainable from
www.thep.lu.se/$``$torbjorn/Pythia.html.
### 7.4 Monte Carlo developments: HERWIG
#### 7.4.1 Angular distribution in $`gq\overline{q}`$
In Herwig, the angular-ordering constraint, which is derived for soft gluon emission, is applied to all parton shower vertices, including $`gq\overline{q}`$. In versions before 6.1, this resulted in a severe suppression (an absence in fact) of configurations in which the gluon energy is very unevenly shared between the quarks. For light quarks this is irrelevant, because in this region one is dominated by gluon emissions, which are correctly treated. However for heavy quarks, this energy sharing (or equivalently the quarks’ angular distribution in their rest frame) is a directly measurable quantity, and was badly described.
Related to this was an inconsistency in the calculation of the Sudakov form factor for $`gq\overline{q}`$. This was calculated using the entire allowed kinematic range (with massless kinematics) for the energy fraction $`x`$, $`0x1`$, while the $`x`$ distribution generated was actually confined to the angular-ordered region, $`x,1xm/\sqrt{E^2\xi }`$ (see Sec. 2.2.1).
In Herwig version 6.1, these defects are corrected as follows. We generate the $`E^2\xi `$ and $`x`$ values for the shower as before. We then apply an a posteriori adjustment to the kinematics of the $`gq\overline{q}`$ vertex during the kinematic reconstruction. At this stage, the masses of the $`q`$ and $`\overline{q}`$ showers are known. We can therefore guarantee to stay within the kinematically allowed region. In fact, the adjustment we perform is purely of the angular distribution of the $`q`$ and $`\overline{q}`$ showers in the $`g`$ rest frame, preserving all the masses and the gluon four-momentum. Therefore we do not disturb the kinematics of the rest of the shower at all.
Although this cures the inconsistency above, it actually introduces a new one: the upper limit for subsequent emission is calculated from the generated $`E^2\xi `$ and $`x`$ values, rather than from the finally-used kinematics. This correlation is of NNL importance, so we can formally neglect it. It would be manifested in an incorrect correlation between the masses and directions of the produced $`q`$ and $`\overline{q}`$ jets. This is, in principle, physically measurable, but it seems less important than getting the angular distribution itself right. In fact the solution we propose maps the old angular distribution smoothly onto the new, so the sign of the correlation will still be preserved, even if the magnitude is wrong.
Even with this modification, the Herwig kinematic reconstruction can only cope with particles that are emitted into the forward hemisphere in the showering frame. Thus one cannot populate the whole of kinematically-allowed phase space. Nevertheless, we find that this is usually a rather weak condition, and that most of phase space is actually populated.
Using this procedure, we find that the predicted angular distribution for secondary $`b`$ quarks at Lep energies is well-behaved, i.e. it looks reasonably similar to the leading-order result ($`1+\mathrm{cos}^2\theta ^{}`$), and has relatively small hadronization corrections.
#### 7.4.2 Predictions for $`g_{bb}`$
Reference contains comparisons between analytical predictions for $`g_{bb}`$ and those of Herwig. One result of the analytical calculation is that, to NLL accuracy, one can use the massless formula for the splitting $`gb\overline{b}`$, provided one also sets a cutoff on the virtual gluon mass of $`m_g>e^{5/6}m_b=2.3m_b`$ instead of the kinematic cutoff $`m_g>2m_b`$. Somewhat fortuitously, this is similar to the Herwig method, which uses the massless formula with a cutoff $`m_g>2(m_b+Q_0)`$ with $`Q_00.5`$ GeV. The comparisons in show that the resummed and Herwig predictions are quite similar at Lep1 and Lep2 energies.
Herwig results for $`Z^0`$ decay are summarized in Table 11, for the version used in the original comparisons (5.7) and the latest version, 6.1 . The main difference between the two versions in this context is a change in the default $`b`$-quark mass from $`m_b=5.2`$ GeV to 4.95 GeV, which is justified by the approximate relation $`m_B=m_b+m_l`$ where $`m_l`$ is the light quark mass. We see that the Herwig results are somewhat higher than the resummed predictions in Table 9 and in better agreement with the data in Table 8.
### 7.5 Monte Carlo developments: ARIADNE
The splitting of gluons into a $`q\overline{q}`$ pair does not fit into the dipole picture in an obvious way, since this splitting is related directly to a single gluon rather than to any dipole between two partons. Also, all gluons in emitted in the cascade are massless, and to be able to split into massive quarks, energy has to be required from somewhere. The way the process is included in Ariadne is described in ref. The splitting probability of a gluon is simply divided in two equal parts, each of which is associated to each of the two connecting dipoles. The splitting process can then again be treated as a two-to-three process, where a spectator parton is used to conserve energy and momentum. It can be shown that this is equivalent to standard parton shower approaches in the limit of strongly ordered emssions. But the differences when extrapolating away from that limit can become large, and Ariadne typically gives twice as many secondary b$`\overline{\text{b}}`$ pairs as compared to eg. Jetset.
But this treatment of secondary heavy quarks may lead to rather strange situations (as noted in ref. ). Since transverse momentum of the $`q\overline{q}`$ splitting can become small even for heavy quarks, it is possible to split a gluon so that the mass $`m_{q\overline{q}}^2`$ is larger than the transverse momentum scale – $`p_g^2`$ – at which the gluon was emitted – although the ordering of the emissions, $`p_{q\overline{q}}^2<p_g^2`$, is still respected. To avoid such situations there is an option in Ariadne<sup>5</sup><sup>5</sup>5MSTA(28)$``$0 in the /ARDAT1/ common block. which introduces an extra limit, $`m_{q\overline{q}}^2<p_g^2`$, on gluon splitting.
As discussed already in Sec. 7.3, it is not quite clear if or how the ordering of emissions should be enforced in the case of gluon splitting into massive quarks and, for that reason, Ariadne also includes an option where these splittings are allowed to be non-ordered, ie. $`p_{q\overline{q}}`$ is allowed to be larger than the transverse momentum of the preceeding emission.<sup>6</sup><sup>6</sup>6MSTA(28)$`<`$0 in the /ARDAT1/ common block. The current version (4.10) contains a bug for this option. A bug fix can be obtained on request to leif@thep.lu.se. The corresponding rates of gluon splittings are given in table 12
## 8 OVERALL CONCLUSIONS AND RECOMMENDATIONS
Here we summarise the main results from each Section above and the recommendations that follow from them. The term ‘jet rates’ always refers to the Durham algorithm unless it is explicitly stated otherwise.
### 8.1 Monte-Carlo developments
The relevant features of the main event generators, Pythia, Herwig and Ariadne, were reviewed with emphasis on relevant new developments. In many cases, important modifications and new options were introduced as a result of discussions in the working group.
In Pythia, improved treatment of quark masses in the parton shower permits a better description of the overall 3-jet rate, but there remains a problem of underestimation of the 4-jet rate. Modification of the way in which matrix element corrections are applied has little effect on this.
In Herwig, parameters for the cluster hadronization of $`b`$-quarks have been separated from those for ligher quarks, so that improved tuning to $`b`$-quark fragmentation data will be possible.
In both Pythia and Herwig there are now options to interface parton showers to the massless 4-parton matrix elements. In addition, Fourjphact interfaces the massive 4-parton matrix elements to Pythia parton showers. These options provide an improved description of 4-jet final states, but are not suitable for describing features that receive important 2- and 3-jet contributions. A very recent development in Herwig is an option to combine 2,3 and (massless) 4-parton matrix elements together with parton showers in a way that aims to avoid the worst aspects of double counting. Comparison and tuning of this option to Lep1 data is in progress.
In Ariadne, options exist for switching on and off the dead-cone effect in QCD radiation from heavy quarks. A massive leading-order matrix element correction has also been introduced, as exists in Pythia and Herwig. Overall, Ariadne gives a better description of jet rates than either Pythia or Herwig. However, the description of mass effects in jet rates in Ariadne is not so satisfactory, and in fact turning off the treatment of quark masses altogether appears to provide a better description.
A major new development is the introduction of the new event generator Apacic++, which for the first time interfaces $`n`$-parton matrix elements and parton showers for $`n=2,\mathrm{},5`$. A number of options are available for choosing the relative jet rates, and for initialising and evolving the parton showers. As in Pythia and Ariadne, the Jetset string hadronisation model is used. A first attempt at tuning to Lep1 data gives encouraging results, with fits to most event shape and single-particle distributions of a quality similar to the established generators. Differential jet rates show some features which may be associated with merging the different parton multiplicities. The tuned shower cutoff is high, so that the parton showers have little phase space for evolution and final-state structure is mostly determined by matrix elements and hadronization.
### 8.2 Jet rates (inclusive)
Both Pythia and Herwig have problems with fitting the 3- and 4-jet rates simultaneously as functions of the jet resolution $`y_c`$. For a given tuning, one can describe e.g. the 3-jet rate well, but then the rates for higher jet multiplicities are overestimated by Herwig and underestimated by Pythia. None of the modifications tried was able to eliminate this problem. For analyses at Lep2 energies using Pythia or Herwig, we recommend tuning to the relevant jet rate at Lep1 in order to minimize the associated systematic error, which then results only from the change in that jet rate from Lep1 to Lep2. If this is done, then a systematic error of 2% in the 4-jet rate at Lep2 could be achieved. On the other hand, if only a general tuning at Lep1 is performed, the systematic error could be as large as 5%.
Ariadne gives the best overall description of jet rates and should therefore be considered as the generator of choice for estimating multi-jet backgrounds, e.g. in hadronic WW decay studies. Using Ariadne could lead to a further reduction of systematic errors.
### 8.3 Jet rates (mass effects)
Full NLO massive matrix element calculations of the 3-jet rate are now available. They were used to study the effect of the $`b`$-quark mass on the ratio of $`b`$-quark to light-quark rates, $`R_3^b\mathrm{}(y_c)`$. This is an observable in which many systematic uncertainties tend to cancel. The difference between the running-mass and pole-mass schemes was used as an estimate of higher-order contributions. This difference was indeed reduced relative to the LO calculation, with the predicted NLO band lying within the LO one.
In the case of the 4-jet rate ratio, $`R_4^b\mathrm{}(y_c)`$, only LO massive predictions are available. Therefore, to be cautious, theoretical uncertainties on both $`R_3^b\mathrm{}`$ and $`R_4^b\mathrm{}`$ were estimated at $`\pm `$ half the difference between the LO pole-mass and running-mass predictions. The Delphi Lep1 data do fall within this band over the range measured ($`0.01y_c0.06`$). For detailed numerical estimates of the uncertainties, see Sec. 4.6.
The performance of the Monte Carlo event generators was judged against the theoretical predictions with the estimated uncertainties. Generally speaking, the generators tend to overestimate mass effects, i.e. they underestimate $`R_3^b\mathrm{}`$ and $`R_4^b\mathrm{}`$. Overall, Herwig gave the best agreement at $`\sqrt{s}=M_\mathrm{Z}`$, although the new mass treatment in Pythia describes $`R_3^b\mathrm{}`$ better. Ariadne underestimates more severely, with dead-cone effects and the new massive matrix element correction tending to worsen agreement.
A full NLO massive matrix element calculations of the 4-jet rate would undoubtedly be helpful in reducing the systematic uncertainty on $`b`$-quark mass effects in jet rates, and in testing and improving the performance of event generators.
A study of hadronization effects showed that use of the Cambridge jet algorithm can considerably reduce hadronization corrections to $`R_3^b\mathrm{}`$, which should be helpful in determinations of the $`b`$-quark mass. Comparison of the latest versions of Pythia and Herwig showed that they give closer estimates of hadronization corrections following improvements in Herwig, and that decay effects are small for sufficiently large values of $`y_c`$.
### 8.4 Four-jet observables
The Monte Carlo generators with specific 4-jet options (Pythia, Herwig and Apacic++) were compared with each other and with matrix element calculations, for the standard set of 4-jet angular distributions as well as the differential 4-jet rate and the D-parameter distribution. No significant differences were found at the matrix element level, except for the differential jet rate in Pythia, which was thought to be due to an intrinsic Jade (mass) cut on the 4-parton configurations generated by that program. Good agreement between the programs was found after parton showering and hadronization. Quark mass effects were found to be small (2% level) for the Lep flavour mixture.
NLO corrections to the 4-jet angular distributions are small but can have a significant effect on the extracted colour factors, owing to their different functional dependence on these quantities. On the other hand NLO effects are very large (70-130%) in the 4-jet rate and D-parameter distributions. This indicates that resummation of large higher-order corrections is required. The 4-jet options of the event generators are not able to describe these distributions owing to the lack of 2- and 3-jet contributions. The default 2+3 jet + parton shower options in Pythia and Herwig are more successful here. The new combined multijet + shower options in Herwig and Apacic++, developed during the workshop, may provide a better simultaneous description of these distributions and of the 4-jet angular distributions, provided successful tuning to the Lep1 data can be achieved.
### 8.5 B fragmentation
The data, theory and models for $`b`$-quark fragmentation into B-hadrons were reviewed. New theoretical calculations with resummation of large higher-order terms suggest that non-perturbative effects are small but significantly different from the conventional Peterson parametrization.
No new work could be undertaken on the important topic of comparing the performance of Monte Carlo generators with the data and with theoretical calculations in this area. Comparisons with new data presented recently by the Sld collaboration suggest that the Pythia/Jetset description with the original Lund parametrization of fragmentation is satisfactory at $`\sqrt{s}=M_\mathrm{Z}`$.
### 8.6 $`gb\overline{b}`$ splitting
The experimental results on the rate of gluon splitting into $`b`$-quark pairs (around 3$`0`$$`/`$$`00`$ at $`\sqrt{s}=M_\mathrm{Z}`$) are somewhat higher than the best theoretical estimates (around 2$`0`$$`/`$$`00`$). However the theoretical uncertainties due to unknown sub-leading logarithmic corrections easily cover the discrepancy. This point is emphasised by the sensitivity of Monte Carlo generator predictions to the treatment of subleading and kinematic effects. In Pythia a number of new options have been introduced to vary the treatment of such effects, and particular choices can readily bring the $`gb\overline{b}`$ rate up to the observed value. We provisionally recommend the setting MSTJ(42)=MSTJ(44)=3. In Herwig and Ariadne, the default settings in the latest versions already give adequate agreement with the data. However, an estimated uncertainty as large as 30% remains appropriate.
## REFERENCES
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# Temperature dependence of the interlayer magnetoresistance of quasi-one-dimensional Fermi liquids at the magic angles
## I Introduction
In spite of intensive research over the past decade the nature of the metallic state in low-dimensional strongly correlated materials is still poorly understood. Widely studied materials include cuprate and organic superconductors. Many of the properties of the cuprates cannot be understood within the Fermi liquid picture that has so successfully described conventional metals . Although some properties of the quasi-two-dimensional molecular crystals, $`\kappa `$-(BEDT-TTF)<sub>2</sub>X (Reference) and the quasi-one-dimensional Bechgaard salts (TMTSF)<sub>2</sub>X can be explained within a Fermi liquid framework others cannot. A particular challenge is understanding the dependence of the magnetoresistance of the Bechgaard salts on the direction of the magnetic field, especially (TMTSF)<sub>2</sub>PF<sub>6</sub> under pressures of about 10 kbar . The different angular-dependent magnetoresistance effects in quasi-one-dimensional metals are known as the Danner , magic angle (or Lebed) , and third angular effects, depending on whether the magnetic field is rotated in the $`𝐚𝐜`$, $`𝐛𝐜`$, or $`𝐚𝐛`$ plane, respectively. (The most- and least-conducting directions are the $`𝐚`$ and $`𝐜`$ axes, respectively). The magic angle effect is the most poorly understood of these effects and is the focus of this paper. If $`\theta `$ is the angle between the magnetic field and the $`𝐜`$ axis then at the “magic angles” given by
$$\mathrm{tan}\theta =\frac{b}{c}\frac{p}{q}\pm p,q=1,2,3,\mathrm{}.$$
(1)
where $`b`$ and $`c`$ are the lattice constants in the $`𝐛`$ and $`𝐜`$ directions, Lebed predicted dips in the threshold field for formation of a field-induced spin-density-wave . Although these dips are not observed , features are seen, mostly at $`p/q=1,2`$ in the torque and in all components of the resistance .
A wide range of physical mechanisms have been proposed to explain these effects including commensurability effects changing the electron-electron scattering rate , semi-classical transport , complicated band structures , hot spots on the Fermi surface , cold spots on the Fermi surface , electron-electron interactions , non-Fermi liquid effects , and magnetic field induced changes in effective dimensionality .
The properties of (TMTSF)<sub>2</sub>PF<sub>6</sub> at 10 kbar are particularly difficult to understand. For example, when the magnetic field is perpendicular to the current direction the magnetoresistance is smaller than when it is parallel, the opposite of what one observes in (TMTSF)<sub>2</sub>ClO<sub>4</sub> and in conventional metals. Recently, the temperature dependence of the magnetoresistance when the field direction was fixed at the first magic angle was measured . It was found to be non-monotonic: as the temperature decreased down to $`T_{min}`$, the magnetoresistance decreased, it increased until $`T_{max}`$, and then decreased. It has recently been proposed that these two temperatures actually represent phase transitions between metallic and insulating phases . The magnetoresistance of the quasi-two-dimensional metal $`\alpha `$-(BEDT-TTF)<sub>2</sub>MHg(SCN)<sub>4</sub> \[M = K,Rb,Tl\] also exhibits unusual temperature and angular dependence .
The purpose of this paper is to clarify what properties of the magic angle effects can only be explained within a non-Fermi liquid framework by seeing what effects can be explained within Fermi liquid theory. The interlayer magnetoresistance is calculated within the framework of semi-classical transport theory. It is found that if one takes into account the finite bandwith along the most-conducting direction then dips in the magnetoresistance are observed for $`p/q=1,2,3,\mathrm{}`$. Furthermore, if one assumes a simple Fermi liquid form for the temperature dependence of the scattering rate then at the magic angles the interlayer magnetoresistance does have a non-monotonic temperature dependence. Hence, one should be cautious about associating maxima and minima in the temperature dependence with metal-insulator transitions. However, the results obtained gives a poor description of the observed properties when the field is close to the $`𝐛`$-axis and of the resistivity within the layers.
## II Calculation of the interlayer conductivity
### A Semi-classical transport theory
If the electronic dispersion relation is $`ϵ(\stackrel{}{k})`$ then the electronic group velocity perpendicular to the layers is $`v_z=\frac{1}{\mathrm{}}\frac{ϵ(\stackrel{}{k})}{k_z}.`$ The interlayer conductivity can be calculated by solving the Boltzmann equation in the relaxation time approximation leading to Chambers’ formula
$$\sigma _{zz}=\frac{e^2\tau }{4\pi ^3}v_z(\stackrel{}{k})\overline{v}_z(\stackrel{}{k})\left(\frac{f(ϵ)}{ϵ}\right)d^3\stackrel{}{k},$$
(2)
where $`f(ϵ)`$ is the Fermi function and $`\tau `$ is the scattering time which is assumed to be the same at all points on the Fermi surface. $`\overline{v}_z(\stackrel{}{k})`$ is the electron velocity averaged over its trajectories on the Fermi surface,
$$\overline{v}_z(\stackrel{}{k})=\frac{1}{\tau }_{\mathrm{}}^0\mathrm{exp}\left(\frac{t}{\tau }\right)v_z(\stackrel{}{k}(t))𝑑t$$
(3)
where $`\stackrel{}{k}(0)=\stackrel{}{k}`$. The time dependence of the wave vector $`\stackrel{}{k}(t)`$ is found by integrating the semi-classical equation of motion
$$\frac{d\stackrel{}{k}}{dt}=\frac{e}{\mathrm{}^2}\stackrel{}{}_kϵ\times \stackrel{}{B}.$$
(4)
If the temperature is sufficiently low that $`T<<E_F`$ then $`\frac{f}{ϵ}`$ in Equation (2) can be replaced by a delta function at the Fermi energy and Equation (2) becomes
$$\sigma _{zz}=\frac{e^2\tau }{4\pi ^3}v_z(\stackrel{}{k})\overline{v}_z(\stackrel{}{k})\delta (E_Fϵ(\stackrel{}{k}))d^3\stackrel{}{k}.$$
(5)
### B Dispersion relation along the chains
In the tight binding approximation the dispersion relation in an orthorhombic crystal can be written as
$$ϵ(\stackrel{}{k})=2t_a\mathrm{cos}(k_xa)2t_b\mathrm{cos}(k_yb)2t_c\mathrm{cos}(k_zc),$$
(6)
where $`t_a,t_b,`$ and $`t_a`$ are the inter-site hopping integrals along the different crystal axes. In the Bechgaard salts, $`t_at_b,t_c`$, the dispersion along the chains can be linearized giving $`ϵ(\stackrel{}{k})=\mathrm{}v_F(|k_x|k_F)2t_b\mathrm{cos}(bk_y)2t_c\mathrm{cos}(bk_z)`$, where $`v_F=2t_aa\mathrm{sin}(ak_F)/\mathrm{}`$ is the Fermi velocity and $`k_F`$ is the Fermi wave vector. This linear dispersion has been used in a number of papers on the magic angle effect . If one solves for the interlayer conductivity within semi-classical transport theory one obtains
$$\frac{\sigma _{zz}(\theta )}{\sigma _{zz}^0}=\frac{1}{1+(\omega _{c0}\tau \mathrm{sin}\theta )^2},$$
(7)
where $`\omega _{c0}=ev_FcB/\mathrm{}`$ is the frequency at which an electron traverses the Brillouin zone in the $`𝐜`$-direction when the field is parallel to the $`𝐛`$-axis. Clearly this is a smoothly varying function of $`\theta `$ and does not exhibit any magic angle effects.
We now show that if the full nonlinear dispersion (6) is used then one does obtain magic angle effects. We will re-derive a result obtained earlier by Maki with a view to elucidating the physics in the process.
### C Solution of the semi-classical equations of motion
The group velocity for the dispersion relation (6) is
$$\stackrel{}{v}(\stackrel{}{k})=\frac{1}{\mathrm{}}\stackrel{}{}_kϵ=\frac{1}{\mathrm{}}\left(\begin{array}{c}2at_a\mathrm{sin}(ak_x)\\ 2bt_b\mathrm{sin}(bk_y)\\ 2ct_c\mathrm{sin}(ck_z)\end{array}\right).$$
(8)
The rate of change of the wave vector $`\stackrel{}{k}(t)`$, in a magnetic field in the $`𝐛`$-$`𝐜`$ plane, $`\stackrel{}{B}=(0,B\mathrm{sin}\theta ,B\mathrm{cos}\theta )`$, is given by (4),
$$\frac{d\stackrel{}{k}}{dt}=\frac{1}{\mathrm{}^2}\left(\begin{array}{c}2beBt_b\mathrm{cos}\theta \mathrm{sin}(bk_y)\\ 2aeBt_a\mathrm{cos}\theta \mathrm{sin}(ak_x)\\ 2aeBt_a\mathrm{sin}\theta \mathrm{sin}(ak_x)\end{array}\right)$$
(9)
where terms involving $`t_c`$ have been neglected. This is valid provided $`t_c\mathrm{sin}\theta t_b\mathrm{cos}\theta `$. Hence, the results below will not be valid as $`\theta 90^0`$.
In order to calculate the $`z`$-component of the velocity one needs to obtain $`\stackrel{}{k}_z(t)`$, which Equation (9c) shows is determined by $`k_x(t)`$. To zero-th order in $`t_b`$, $`k_x(t)=k_F`$. Integrating Equation (9b) then gives
$$k_y(t)=k_y(0)+\frac{\omega _b}{b}t,$$
(10)
where
$$\omega _b=v_FeBb\mathrm{cos}\theta /\mathrm{}\omega _{b0}\mathrm{cos}\theta $$
(11)
is the frequency at which an the electron traverses the Brillouin zone in the direction of the $`𝐛`$-axis.
Substituting this into Equation (9a) and integrating gives, to first order in $`t_b/t_a`$,
$$k_x(t)=k_F+\frac{2t_b}{\mathrm{}v_F}\mathrm{cos}(bk_y(0)+\omega _bt).$$
(12)
We obtain $`k_z(t)`$ by using Equation (9c) and substituting in (12), giving
$$\frac{dk_z}{dt}=\frac{2aeBt_a\mathrm{sin}\theta }{\mathrm{}^2}\left[\mathrm{sin}(k_F)\mathrm{cos}\left(\frac{2at_b}{\mathrm{}v_F}\mathrm{cos}(bk_y(0)+\omega _bt)\right)+\mathrm{cos}(ak_F)\mathrm{sin}\left(\frac{2at_b}{\mathrm{}v_F}\mathrm{cos}(bk_y(0)+\omega _bt)\right)\right],$$
(13)
where we have used trigonometric identities to expand $`\mathrm{sin}(ak_x(t))`$. If we take a linear dispersion relation, the second term in (13) will equal zero and we are left with $`\frac{dk_z}{dt}=`$ $`\frac{Bev_F\mathrm{sin}\theta }{\mathrm{}}`$ , where we have assumed that at $`t=0`$, the wave vector in the x-direction ($`k_x`$) is equal to $`k_F`$.
Now to first order in $`t_b/t_a,`$
$$\frac{dk_z}{dt}=\frac{2aeBt_a\mathrm{sin}\theta }{\mathrm{}^2}\left[\mathrm{sin}(ak_F)+\mathrm{cos}(ak_F)\mathrm{sin}\left(\frac{2at_b}{\mathrm{}v_F}\mathrm{cos}(bk_y(0)+\omega _bt)\right)\right].$$
(14)
Integrating this we obtain
$$k_z(t)c=k_z(0)c\omega _ct\gamma _0\mathrm{tan}\theta \mathrm{sin}(bk_y(0)+\omega _bt),$$
(15)
where
$$\omega _c=\omega _{c0}\mathrm{sin}\theta $$
(16)
and
$$\gamma _0=\frac{2ct_b}{\mathrm{}v_F}\frac{a}{b}\mathrm{cot}(ak_F).$$
(17)
### D Evaluation of the interlayer conductivity
Substitution of (15) into the $`z`$-component of the velocity gives
$$v_z(k_z(0),\varphi ,\varphi ^{^{}})=\frac{2ct_c}{\mathrm{}}\mathrm{sin}\left(ck_z(0)+\frac{\omega _c}{\omega _b}\varphi ^{^{}}\gamma _0\mathrm{tan}\theta \mathrm{sin}(\varphi \varphi ^{^{}})\right),$$
(18)
where $`\varphi ^{^{}}=\omega _bt`$, $`\varphi =bk_y(0)`$.
The interlayer conductivity given by (5) can then be written in the form
$$\sigma _{zz}=\frac{e^2}{4\pi ^3b\mathrm{}v_F}_{\pi /c}^{\pi /c}𝑑k_z(0)_0^{2\pi }𝑑\varphi v_z(k_z(0),\varphi )_0^{\mathrm{}}\frac{d\varphi ^{^{}}}{\omega _b}\mathrm{exp}\left(\frac{\varphi ^{^{}}}{\tau \omega _b}\right)v_z(k_z(0),\varphi ,\varphi ^{^{}}).$$
(19)
We now expand Equation (18) using trigonometric identies and substitute the Bessel generating functions to obtain
$`v_z(k_z(0),\varphi ,\varphi ^{^{}})`$ $`=`$ $`{\displaystyle \frac{2ct_c}{\mathrm{}}}[\mathrm{sin}(ck_z(0)+{\displaystyle \frac{\omega _c}{\omega _b}}\varphi ^{^{}})[J_0(\gamma _0\mathrm{tan}\theta )+2{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}J_{2k}(\gamma _0\mathrm{tan}\theta )\mathrm{cos}((2k)(\varphi \varphi ^{^{}}))]`$ (21)
$`+\mathrm{cos}(ck_z(0)+{\displaystyle \frac{\omega _C}{\omega _b}}\varphi ^{^{}})\left[2{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}J_{2k+1}(\gamma _0\mathrm{tan}\theta )\mathrm{sin}((2k+1)(\varphi \varphi ^{^{}}))\right]]`$
and $`v_z(k_z(0),\varphi )`$ is obtained by setting $`\varphi ^{^{}}=0`$. Substituting these expressions into Equation (19) and performing the integrals over $`\varphi ^{^{}}`$, $`\varphi `$ and $`dk_z(0)`$ the final expression for the conductivity becomes
$$\sigma _{zz}(\theta )=\sigma _{zz}^0\left[\frac{J_0(\gamma _0\mathrm{tan}\theta )^2}{1+(\omega _c\tau )^2}+\underset{\nu =1}{\overset{\mathrm{}}{}}J_\nu (\gamma _0\mathrm{tan}\theta )^2\left(\frac{1}{1+\tau ^2(\omega _c\omega _b\nu )^2}+\frac{1}{1+\tau ^2(\omega _c+\omega _b\nu )^2}\right)\right],$$
(22)
where $`\sigma _{zz}^0=\frac{2e^2\tau ct_c^2}{\pi b\mathrm{}^3v_F}`$ is the interlayer conductivity in zero field. Note that for fixed $`\nu `$ and $`z1`$
$$J_\nu (z)\frac{(\frac{z}{2})^\nu }{\mathrm{\Gamma }(\nu +1)}.$$
(23)
Maki obtained a similar result, although he included the corrections to $`\omega _b`$ and $`\omega _c`$ to next order in $`\left(\frac{t_b}{t_a}\right)^2`$. This raises the general question of to what order in $`t_b/t_a`$ is the above expression for $`\sigma _{zz}`$ valid. We only calculated $`k_z(t)`$ to first order in $`t_b/t_a`$. Strictly speaking, this means that (22) is valid to second order in $`t_b/t_a`$. However, we anticipate that a general solution for $`v_z(t)`$ will be of the form $`v_z(t)\mathrm{sin}(\omega _ct)_na_n\mathrm{sin}(\omega _bt)`$ where $`a_n`$ is of order $`(t_b/t_a)^n`$. This means that the coefficients in (22) for $`\nu 2`$ will change but be of the same order.
## III Magic angles
The angular dependence of the interlayer resistivity given by Equation (22) is shown in Figure 1 for several parameter values. Dips occur at the “magic angles” given by $`\omega _c=\nu \omega _b`$ or
$$\mathrm{tan}\theta =\frac{b}{c}\nu \pm \nu =1,2,3,\mathrm{}.$$
(24)
where $`b`$ and $`c`$ are lattice constants. The size of the dip at the $`\nu `$-th magic angle, compared the background magnetoresistance will be of order
$$\left(\frac{\gamma _0}{2}\frac{b}{c}\nu \right)^{2\nu }\left(\frac{\omega _c\tau }{\nu !}\right)^2.$$
(25)
The size of the dips is determined by the parameter $`\gamma _0`$, defined by (17), which is determined by the geometry of the Fermi surface. Note that if $`\gamma _00`$, (22) reduces to (7). This is because the limit $`\gamma _00`$ corresponds to taking a linear dispersion relation. If $`\gamma _01`$ then the dips will decrease in magnitude rapidly with increasing $`\nu `$. For example if $`\gamma _00.1`$ then the $`\nu =1`$ feature will be five orders of magnitude smaller than the $`\nu =3`$ feature. Note that when $`\nu `$ becomes sufficiently large this will no longer be valid because $`\gamma _0\mathrm{tan}\theta 1`$.
We now consider what is a realistic value for $`\gamma _0`$ for the (TMTSF)<sub>2</sub>X materials. If we look at the the form of $`\gamma _0`$ in Equation (17) we note that the factor $`\frac{2ct_b}{\mathrm{}v_F}`$ equals the parameter $`\gamma `$ which determines the periodicity of the Danner oscillations . For (TMTSF)<sub>2</sub>ClO<sub>4</sub> it was estimated to be 0.24. The lattice constants for (TMTSF)<sub>2</sub>PF<sub>6</sub> are $`a=7.3\AA `$, $`b=7.7\AA `$ and $`c=13.5\AA `$, while for (TMTSF)<sub>2</sub>ClO<sub>4</sub>, $`b`$ is twice as large due to anion ordering . The $`\mathrm{cot}(ak_F)`$ term depends on the band filling. At three-quarter filling $`k_F=\frac{3\pi }{4a}`$, and $`\mathrm{cot}(ak_F)=1`$. This gives a value for $`\gamma _0(PF_6)=0.24`$ and $`\gamma _0(ClO_4)=0.12`$. Note that for half-filling $`k_F=\frac{\pi }{2a}`$, and thus $`\gamma _0=0`$ and there will be no magic angle effects unless we solve the semi-classical equations to higher order in $`t_b`$.
Figure 1 is qualitatively similar to experimental results for (TMTSF)<sub>2</sub>ClO<sub>4</sub> at ambient pressure and at 6.0 kbar and (TMTSF)<sub>2</sub>PF<sub>6</sub> at 6.0 kbar (0.3 K and 4 T) . A small difference is that the experimental data shows a small dip near $`90^0`$, whereas the theoretical curve shows no such dip. It is quite possible that the small dips can be explained within semi-classical transport theory if one includes the effect of a finite $`t_c`$ in the solution of the semi-classical transport equations. An analogous effect occurs when the field is rotated in the $`𝐚𝐜`$ plane: for coherent interlayer transport with finite $`t_c`$ a peak in the angular-dependent magnetoresistance occurs when the field is parallel to the $`𝐚`$ axis .
The angular dependence of the interlayer magnetoresistance of (TMTSF)<sub>2</sub>PF<sub>6</sub> at pressures of about 10 kbar is quite different from that shown in Figure 1 At fields less than one tesla the angular dependence is similar to that given by Eqn. (7). However, above one tesla, $`\rho _{zz}(B\mathrm{cos}\theta )^{1.3}`$ and so a large dip is observed near $`90^0`$ and at $`90^0`$ the in-field resistance is comparable to the zero-field resistance .
A number of theoretical papers have predicted effects when $`\mathrm{tan}\theta =\frac{b}{c}\frac{p}{q}`$ where $`p/q`$ is fraction. In contrast, the model considered here only gives effects for $`q=1`$. A review of the experimental literature shows that the only reproducible fractional features seen have been in (TMTSF)<sub>2</sub>ClO<sub>4</sub> at $`p/q=`$ 3/2 and 5/2 . This can be explained within the framework considered here. If (TMTSF)<sub>2</sub>ClO<sub>4</sub> is slowly cooled anion ordering occurs and the lattice constant in the $`𝐛`$ direction doubles so in (24) $`b`$ should be replaced by $`2b`$. However, if a sample which is not completely anion ordered it will produce features at angles corresponding to half-integers for a fully anion-ordered sample.
## IV Temperature dependence of the interlayer magnetoresistance at the magic angles
Suppose that the field direction is fixed at a magic angle and the temperature (and thus the scattering time $`\tau `$) is varied.
The first magic angle $`(\nu =1)`$. Setting $`\omega _c=\omega _b`$ and using the fact that $`\gamma 1`$ to take just the first term in the series, i.e. $`\nu =1`$, the conductivity is
$$\sigma _{zz}(\theta _1)A\tau \left[\frac{1}{1+(\omega _{c0}\mathrm{sin}\theta _1\tau )^2}+\left(\frac{\gamma _0\mathrm{tan}\theta _1}{2}\right)^2\right],$$
(26)
where $`\theta _1`$ represents the first magic angle and $`\sigma _{zz}^0=A\tau `$, where $`A`$ is the ratio of the zero field conductivity to $`\tau `$. A plot of interlayer resistivity verses $`\frac{1}{\sqrt{\tau }}`$ is shown in Figure 2 for different values of $`\omega _0`$. The interlayer resistivity is a non-monotonic function of $`\tau .`$ It will be seen below that this leads to non-monotonic temperature dependence.
We now find for what values of $`\tau `$ the maxima and minima seen in Figure 2 occur. Finding the extrema of Equation (26) as a function of $`\tau `$ gives that the minimun occurs when
$$\omega _{c0}\tau \frac{1}{\mathrm{sin}\theta _1}$$
(27)
and the maximum occurs when
$$\omega _{c0}\tau \frac{2}{\gamma _0\mathrm{sin}\theta _1}.$$
(28)
If $`\gamma _01`$, then the maximum will only be observed at sufficiently high fields and in high purity samples.
The second magic angle $`(\nu =2)`$. To obtain the interlayer conductivity for the second magic angle then a similar argument to that given above leads to
$$\sigma _{zz}(\theta _2)A\tau \left[\frac{1}{1+(\tau \omega _c)^2}+\left(\frac{\gamma _0\mathrm{tan}\theta _2}{2}\right)^2\left(\frac{1}{1+\left(\frac{\tau \omega _c}{2}\right)^2}+\frac{1}{1+\left(\frac{3\tau \omega _c}{2}\right)^2}\right)+\frac{\left(\gamma _0\mathrm{tan}\theta _2\right)^4}{64}\right],$$
(29)
where we have set $`\omega _c=2\omega _b`$ and $`\theta _2`$ is the $`\nu =2`$ magic angle. For small $`\gamma _0`$, the minima is again given by $`\omega _c\tau 1`$. To find the maxima we expand the first two terms in (29) to fourth order in $`\frac{1}{(\tau \omega _c)^2}.`$ The maxima occurs when
$$\omega _{c0}\tau \frac{4}{\mathrm{sin}\theta _2(\gamma _0\mathrm{tan}\theta _2)^2}.$$
(30)
Since this is smaller than (28) by a factor of $`1/\gamma _0`$, in order to see this maximum even higher fields and lower temperatures will be required than for the maximum associated with the first magic angle.
Conductivity as $`\theta 90^0`$: We can expand the term $`\left(\frac{1}{1+\tau ^2(\omega _c\omega _b\nu )^2}+\frac{1}{1+\tau ^2(\omega _c+\omega _b\nu )^2}\right)`$ in the summation in (22) to second order in $`\nu \mathrm{cos}\theta `$ to obtain
$$\frac{2}{1+a^2}(1+\nu ^2\mathrm{cos}\theta ^2A+\mathrm{}\mathrm{}),$$
(31)
where $`a=\omega _0\tau `$ and $`A=\frac{a^2(3a^21)}{(1+a^2)^2}3`$ for $`\omega _0\tau 1`$. Substitution of this into the conductivity gives
$$\frac{\sigma _{zz}(\theta 90^0)}{\sigma _{zz}^0}\frac{1}{(\omega _0\tau )^2}\left[J_0(\gamma _0\mathrm{tan}\theta )^2+2\underset{\nu =1}{\overset{\mathrm{}}{}}J_\nu (\gamma _0\mathrm{tan}\theta )^2(1+3\nu ^2\mathrm{cos}\theta ^2+\mathrm{})\right].$$
(32)
This can be simplified using the identities
$$\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}J_n(z)^2=1\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}n^2J_n(z)^2=z^2/2$$
(33)
to give
$$\frac{\rho _{zz}(\theta =90^0)}{\rho _{zz}^0}\frac{(\omega _{c0}\tau )^2}{1+3\gamma _0^2}.$$
(34)
The resistivity is quadratic in field as for the case of a linear dispersion, but the co-efficient is smaller. This is consistent with Figure 1 which shows that the resistivity near $`90^0`$ does decrease with increasing $`\gamma _0`$. This field dependence is quite different to what is observed in the (TMTSF)<sub>2</sub>X materials. For (TMTSF)<sub>2</sub>ClO<sub>4</sub> at 6.0 kbar a linear field dependence is observed at high fields. In (TMTSF)<sub>2</sub>PF<sub>6</sub> at pressures from 6 to 10 kbar the resistivity saturates as the field increases . However, caution is in order because derivation of the quadratic dependence involved assuming that $`t_c\mathrm{tan}\theta t_b`$ and so we only expect a quadratic dependence slightly away from $`90^0`$ or in the limit $`t_c0`$.
Fermi liquid model for the temperature dependence: We now consider a specific model for the temperature dependence of the scattering time $`\tau `$. In a Fermi liquid the scattering rate, at temperatures much less than the Fermi temperature, has a temperature dependence of the form
$$\frac{1}{\tau }=\frac{1}{\tau _0}+\beta T^2,$$
(35)
where the first term is due to impurity scattering and the second is due to electron-electron scattering. Using this expression for $`\tau `$ in Equation (26) we can now plot the temperature dependence of the resistivity. This is shown in Figure 3 for various values of $`\omega _{c0}\tau _0`$. The resistivity is not a monotonic function of temperature but has a minimum when $`\omega _c\tau (T)1`$. If the field is sufficiently high there is also a maximum. Using (35) we see that the minimum occurs at a temperature
$$T_{min}\left(\frac{\omega _{c0}\mathrm{sin}\theta }{\beta }\right)^{1/2}.$$
(36)
If $`T_{max}^21/(\beta \tau _0)`$ then (28) implies
$$T_{max}\left(\frac{\omega _{c0}\gamma _0\mathrm{sin}\theta \mathrm{tan}\theta }{2\beta }\right)^{1/2}.$$
(37)
We are unaware of any measurements of the temperature dependence of the interlayer magnetoresistance of (TMTSF)<sub>2</sub>ClO<sub>4</sub> at the magic angles. Although the temperature dependence at the $`\nu =1`$ magic angle shown in Figure 3 is similar to that reported in Ref. for the intralayer resistance of (TMTSF)<sub>2</sub>PF<sub>6</sub> at 9 kbar, the observed temperature dependence of the interlayer magnetoresistance is different . It depends weakly on the temperature from 15 K down to about 3 K and then decreases.
The temperature dependence shown in Figure 3 for $`\theta =90^0`$ i.e., as the magnetic field is aligned with the $`𝐛`$ axis the temperature dependence is qualitatively similar to that observed in (TMTSF)<sub>2</sub>ClO<sub>4</sub> at ambient pressure . Again, qualitatively very different behavior was observed in (TMTSF)<sub>2</sub>PF<sub>6</sub> at 10 kbar. There it was found that the in-field resistance had a temperature dependence similar to the zero-field resistance.
Zheleznyak and Yakovenko have given a heuristic argument as to the origin of the maximum and minimum temperatures seen in the intralayer resistance in Ref. , suggesting that metal-insulator and insulator-metal phase transitions occur as the temperature passes through these values . They argue that $`T_{min}\omega _bB`$ and $`T_{max}t_c`$ which is independent of field. Above $`T_{min}`$ the system is a two-dimensional metal. Below $`T_{min}`$, a magnetic field causes the electron motion in the $`𝐛`$ direction to be quantized resulting in a one-dimensional dispersion and correlations producing insulating behavior. Below $`T_{max}`$, the interlayer coupling becomes important and metallic behavior is recovered. In contrast we find that $`T_{min}`$ and $`T_{max}`$ are given by (36) and (37), respectively, and both scale with $`\sqrt{B}`$. Careful measurements should be able to distinguish between these two different field dependences.
## V Conclusions
This paper only considers the interlayer resistivity $`\rho _{zz}`$, whereas magic angle effects are also seen experimentally in the intralayer resistivity $`\rho _{xx}`$. Maki pointed out that the semi-classical theory will only give resonances in $`\rho _{xx}`$ of order $`(t_c/t_a)^2`$ whereas they are observed to be much larger. A possible way around this problem is that experiments that are meant to measure $`\rho _{xx}`$ may actually be measuring some of $`\rho _{zz}`$. This is because in highly anisotropic metals it is difficult to arrange the contacts and current path so it lies completely within the layers. This potential problem increases the motivation to make thin film samples of these metals.
It has been shown that within semi-classical transport theory a nonlinear dispersion parallel to the chains is necessary to produce dips in the interlayer magnetoresistance at integer magic angles. If the field direction is fixed at one of the magic angles then one observes both minima and maxima in the temperature dependence of the interlayer magnetoresistance. This arises from the temperature dependence of the scattering rate. Since these maxima and minima can exist within a Fermi liquid model one should be cautious about associating them with non-Fermi liquid behaviour or metal-insulator transitions. On the other hand, although the Fermi liquid model considered here gives a good description of many of the properties of (TMTSF)<sub>2</sub>PF<sub>6</sub> at pressures from 6 to 8 kbar and of (TMTSF)<sub>2</sub>ClO<sub>4</sub> it gives a poor description of their properties when the field is parallel to the layers and of the intralayer transport. Qualitatively very different behavior is observed in (TMTSF)<sub>2</sub>PF<sub>6</sub> at pressures of about 10 kbar; explaining it remains a considerable theoretical challenge.
###### Acknowledgements.
This work was supported by the Australian Research Council. We thank P.M. Chaikin and M.J. Naughton for very helpful discussions. We also thank P.M. Chaikin for sending us unpublished experimental data.
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